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warning/0506/astro-ph0506537.html | ar5iv | text | # 2-D models of layered protoplanetary discs: I. The ring instability
## 1 Introduction
The layered disc model was proposed by Gammie (1996) to account for accretion-related phenomena in T Tauri stars. He assumed that the angular momentum is transported by the magneto-rotational instability, commonly referred to as the MRI (Balbus & Hawley, 1991). However, in the outer disc (beyond $`0.1`$ AU) the temperature and the ionization degree is so low that the gas is not well coupled to the magnetic field and the MRI decays. There, the only parts of the disc in which the MRI can operate are the surface layers that are ionized by cosmic rays (ionization due to x-ray quanta emitted by the central star was also considered; see Glassgold, Najita & Igea 1997). Sandwiched between the active surface layers is an MRI-free, and, consequently, non-viscous area near the mid-plane of the disc, commonly referred to as the dead zone.
An interesting property of layered discs is that, in general, the accretion rate $`\dot{M}`$ is a function of the radius $`r`$. The specific form of this function depends on the mass-weighted opacity. However, $`\dot{M}`$ increases with $`r`$ if the opacity does not depend on the density and increases with the temperature not faster than $`T^2`$. This is true in the range $`203`$ K $`<`$ $`T`$ $`<`$ $`2290\rho ^{2/49}`$ K, where the opacity is dominated by ice-free grains, and also between $`167`$ K and $`203`$ K, where the opacity drops due to the sublimation of ices as $`T`$ grows. At temperatures below $`167`$ K the opacity varies as $`T^2`$, and the accretion rate is constant as a function of $`r`$ (in this paper we use opacities of gas and dust mixture according to Bell & Lin 1994).
When $`\dot{M}`$ increases with $`r`$, an annulus centred at $`r_0`$ receives more mass per unit time from the outer disc $`(r>r_0)`$ than that it loses to the inner disc $`(r<r_0)`$. As a result, the mass accumulates in the dead zone. Eventually, at one or more locations in the dead zone the surface density becomes so large that the heat released by the accretion of the accumulated matter pushes the temperature to a level at which the collisional ionization can restore the coupling between the gas and the magnetic field. In such case a triggering event, e.g. perturbation due to the passage of the companion star, heat flux from the inner active disc or gravitational instability of the dead zone (Armitage, Livio & Pringle, 2001), could start the accretion and ”ignite” the dead zone, making at least part of it active. Consequently, the rate at which mass is accreted by the central star could increase dramatically. This mechanism was suggested by Gammie (1996) to explain the FU Ori type outbursts (Hartmann & Kenyon, 1996).
The layered disc model was further developed by Huré (2002), who noted that the non-zero thickness of the dead zone is an important factor influencing the structure of the active surface layers (this is because their structure depends on the vertical component of gravity, which in turn depends on the distance from the mid-plane). On the other hand, detailed MHD simulations of surface layers performed in a 3D shearing box approximation by Fleming & Stone (2003) showed that MRI-driven, non-axisymmetric density waves can propagate far into the dead zone. As a result, a purely HD turbulence is excited there, providing an effective viscosity.
Our original intention was to study how the evolution of the dead zone may be affected by the radiative heating from the inner active part of the disc. To that end, we performed extensive two-dimensional (axisymmetric) simulations of the evolution of layered discs, allowing for radiative energy transfer in both radial and vertical directions. Unexpectedly, we found that the dead zone is unstable in a way that has not been reported before; namely, it tends to decompose into rings. In the present paper we illustrate this ”ring” instability with numerical simulations, and we discuss it analytically, providing a physical explanation of the observed phenomena. The remaining results of our simulations will be reported in a forthcoming paper.
The outline of the present paper is as follows. In §2 we briefly describe the numerical code and we list the assumptions underlying our simulations. A layered-disc model whose dead zone decomposes into rings is presented in §3. §4 contains an analytical discussion of the ring instability. Finally, in §5 we summarize our results and discuss the effects of the ring instability in more sophisticated disc models.
## 2 Numerical methods
The simulations are performed with the help of the 3-D code TRAMP (Klahr et al., 1999). The equations of continuity
$$\frac{\rho }{t}+(\rho 𝐯)=0,$$
(1)
momentum conservation
$$\frac{\rho 𝐯}{t}+(𝐯)\rho 𝐯=P+\rho \mathrm{\Phi }+𝐟_{\mathrm{cen}}+𝐓,$$
(2)
and internal energy
$$c_v\rho \left[\frac{T}{t}+(𝐯)T\right]=P𝐯+𝐓:(𝐯)$$
(3)
are solved in spherical coordinates $`(r,\theta ,\varphi )`$ using an explicit operator-splitting method. $`\mathrm{\Phi }`$ is here the gravitational potential, $`𝐟_{\mathrm{cen}}`$ is the centrifugal force, $`𝐓`$ is the viscous stress tensor, $`𝐯`$ is the rate of strain tensor composed of derivatives of velocity components, $`T`$ is the temperature (assumed to be the same for gas, dust and radiation), $`c_v`$ is the specific heat, and the colon denotes a double scalar product of two tensors. The self-gravity of the disc is neglected, so that
$$\mathrm{\Phi }=\frac{GM_{}}{r},$$
where $`M_{}`$ is the mass of the central star. We use the equation of state of the ideal gas
$$P=\frac{k_\mathrm{B}T}{\mu m_\mathrm{H}}\rho ,$$
(4)
where $`k_\mathrm{B}`$ is Boltzmann constant, $`\mu `$ is the average molecular weight of the disc gas, and $`m_\mathrm{H}`$ is the proton mass.
The radiation transport is treated at the end of each hydrodynamic time-step by solving the equation
$$\frac{E_r}{t}=𝐅,$$
(5)
where $`E_r=aT^4`$ is the radiation energy density, and $`𝐅`$ is the radiative flux. We adopt
$$𝐅=\frac{\lambda c}{\rho \kappa }E_r,$$
(6)
where $`\kappa `$ is the Rosseland mean opacity for the gas-dust mixture (Bell & Lin, 1994), and $`\lambda `$ is the flux limiter used to interpolate between optically thin and optically thick cases (Levermore & Pomraning, 1981). Equations (5) and (6) are solved implicitly using an iterative SOR method. After the convergence has been achieved, the updated temperature is calculated from the updated radiation energy density. The advection routine for mass, momentum and energy is based on a second-order monotonic transport scheme originally introduced by van Leer (1977) and optimized by Kley & Hensler (1987). For more details about the numerical code we refer to Klahr et al. (1999).
The models are axially symmetric. We allow for the flow through the midplane of the disc, i.e. we do not impose a reflecting boundary condition there. Grid points are spaced logarithmically in $`r`$, resulting in a progressively smaller radial extent of grid cells near the centre of the disc, where their vertical extent is also progressively smaller. Thus, the shape of grid cells is more nearly regular throughout the grid.
In the inner active part of the disc, approximate initial distributions of density and temperature are obtained from analytical $`\alpha `$-disc models (Shakura & Sunayev, 1973) assuming a constant temperature profile in the direction perpendicular to the midplane. The outer, layered part of the disc is initiated with the analytical model of Gammie (1996). The two solutions merge at the radius $`r_{\mathrm{D}Z}`$, where the mid-plane temperature of the layered part reaches $`1000`$ K. The accretion rate of the layered part at $`r_{\mathrm{D}Z}`$ defines the accretion rate in the inner active part. Similarly, the surface density of the dead zone is calculated assuming continuity of the total surface density at $`r_{\mathrm{D}Z}`$. Initially, the surface density of the dead zone is constant in $`r`$. The initial rotational velocity is chosen so as to balance the gravity reduced due to the radial pressure gradient (in a thin disc the rotation is nearly Keplerian). Other components of velocity are set to zero.
At the outer edge of the disc mass is injected at every time-step at a fixed rate by the following procedure:
1. densities and temperatures are copied into the ghost zone from the adjacent active cells
2. densities in the ghost zone are normalized to the initial surface density
3. a uniform radial velocity is set across the ghost zone such that the accretion rate is the same as in the analytical model.
The angular rotational velocity at both inner and outer boundary is extrapolated using a $`r^{3/2}`$ power-law. At the inner edge of the disc, and at the upper and lower boundary of the grid, a free outflow boundary condition is imposed. A constant temperature of $`10`$ K is maintained at the upper and lower boundary.
At each time-step the location of the dead zone is found based on two conditions that have to be fulfilled simultaneously:
1. $`T<T_{\mathrm{lim}}`$, where $`T_{\mathrm{lim}}`$ is the minimum temperature at which the coupling between the gas and the magnetic field still operates
2. $`\mathrm{\Sigma }>\mathrm{\Sigma }_a`$, where $`\mathrm{\Sigma }`$ is the column density integrated from the surface of the disc along a line perpendicular to the midplane, and $`\mathrm{\Sigma }_a`$ is the column density ionized by the cosmic rays.
Beyond the dead zone the viscosity coefficient is defined according to the $`\alpha `$-prescription of Shakura & Sunayev (1973)
$$\nu =\alpha c_\mathrm{s}H_a,$$
(7)
where $`\alpha `$ is a dimensionless parameter, $`c_s`$ is the sound speed and $`H_a`$ is the thickness of the active layer which is determined from the vertical distribution of density for each radius at each time-step. We chose this particular form of Shakura-Sunyaev prescription because it is closer to the analytical, vertically averaged model. Within the dead zone we set $`\nu =0`$.
To avoid numerical problems, we introduce a density limiter: whenever the density falls below $`\rho _{min}`$ (which may be different for different models), it is doubled, and the temperature is adjusted so as to keep the pressure unchanged. This procedure leads to the formation of a low-density ”atmosphere” surrounding the disc. To prevent it from collapsing onto the disc, we artificially damp vertical and radial velocities in this region. However, this artificial readjustment only affects a negligible number of grid cells in the ”atmosphere”, and the overall evolution of the disc is not influenced.
## 3 Results of simulations
We obtained a broad sample of models with different parameters, which will be presented elsewhere. Here we describe a representative model with ”canonical” parameters $`\mu =2.353`$, $`T_{\mathrm{lim}}=1000`$ K, $`\mathrm{\Sigma }_a=100\mathrm{g}\mathrm{cm}^2`$ and $`\alpha =0.01`$. Note that at $`T=1000`$ K the ionization degree increases by five orders of magnitude due to ionization of potassium (Umebayashi, 1983), and the magnetic Reynolds number, which is directly proportional to the ionization degree, exceeds unity. $`\mathrm{\Sigma }=100\mathrm{g}\mathrm{cm}^2`$ is the standard stopping column density of cosmic rays (Umebayashi & Nakano, 1981), and 0.01 is a value of the viscosity parameter often adopted for protoplanetary discs (e.g. Hawley, Gammie & Balbus 1995). The simulations were performed on a grid of $`256\times 45`$ points in ($`r`$,$`\theta `$), with $`r`$ varying between 0.05 and 0.35 AU and $`\theta `$ varying between $`5\mathrm{°}`$ and $`+5\mathrm{°}`$. The model was integrated for $`124`$ yr, i.e. for 600 orbits at the outer edge of the disc.
The initial model is far from thermal equilibrium, however it quickly relaxes due radiative heat diffusion. The density and mass flux across the computational domain shortly after the relaxation, at the beginning of the instability and at the end of the simulation are shown in Fig. 1: one clearly sees how the dead zone successively breaks into more and more rings.
As we have not imposed any explicit perturbations, there is a good reason to believe that the rings result from a linear instability of layered discs. The mechanism of this instability is illustrated in Fig. 2. Assume a small axially symmetric enhancement of the surface density at a radius $`r_0`$ (which, of course, must be accompanied by an increase of the dead zone thickness $`H_{\mathrm{DZ}}`$). Since at higher distances from the mid-plane the vertical component of gravity is stronger, the active layer (whose column density is at all times fixed by cosmic rays) must get thinner at $`r_0`$.
In the standard prescription, the viscosity coefficient is proportional to the local thickness of the active zone $`H_a`$ (this is motivated by the idea that $`H_a`$ is a natural scale that limits the maximum size of the edies of the MHD turbulence). Therefore, smaller $`H_a`$ results in a lower viscosity in the elevated part of the active layer. The accretion rate, which depends on the derivative of viscosity (see Eq. 8 below), decreases at the inner edge of the ring. This causes a bottle-neck effect in the accretion flow, and the mass accumulates in the ring. On the other hand, the higher accretion rate at the outer edge of the ring makes the ring more compact. This positive feedback between the dead zone thickness and the mass accumulation rate leads to the formation of the rings.
The rings formed in the simulation are shown in the radial profile of the surface density in Fig. 3. The accretion rate profile exhibits the described minima at the inner edges of the rings and maxima at the outer ones. The plot also shows the deviations from the Keplerian angular velocity – it can be seen that the inner edges of the rings rotate at a super-Keplerian velocity. Thus, the rings may capture and concentrate the radially drifting boulders of meter-size, preventing them from accretion onto the central star (e.g. Klahr & Lin 2000). Such concentration of solids in peaks of gas density was observed by Haghighipour & Boss (2003) and Rice et al. (2004) in the case of spiral density waves formed by the gravitational instability.
## 4 Analytical description of the ring instability
### 4.1 Basic assumptions and definitions
Let us consider a layered disc consisting of a dead zone with thickness $`2H_{\mathrm{DZ}}`$ and two active surface layers with thickness $`H_a`$ each. The surface layer accretes at a rate
$$\dot{M}=6\pi r^{1/2}\frac{}{r}(2\mathrm{\Sigma }_a\nu r^{1/2}),$$
(8)
where $`\mathrm{\Sigma }_a`$ is the column density of the surface layer, $`\nu `$ is the kinematic viscosity, and $`r`$ is the cylindrical radial coordinate. Equation (8) is a direct consequence of angular momentum conservation (Gammie, 1996). Assuming that the accretion energy is radiated locally, we obtain the standard formula for the effective temperature $`T_e`$
$$\frac{9}{4}\mathrm{\Sigma }_a\nu \mathrm{\Omega }^2=\sigma T_e^4,$$
(9)
where $`\sigma `$ is Stefan-Boltzmann constant. In the optically thick approximation the temperature $`T_i`$ at the boundary between the dead zone and the active layer is given by the formula
$$T_i^4=\frac{3}{8}\tau T_e^4,$$
(10)
where $`\tau =\mathrm{\Sigma }_a\kappa `$ is the optical depth of the active layer, and $`\kappa `$ is the Rosseland mean opacity (Hubeny, 1990). In general, the opacity coefficient is a complex function of density and temperature. In a protoplanetary disc it can be approximated by a set of power laws
$$\kappa (\rho ,T)=\kappa _0\rho ^aT^b$$
(11)
where the constants $`\kappa _0`$, $`a`$ and $`b`$ have different values in different opacity regimes (Bell & Lin, 1994). At temperatures lower than $`2290\rho ^{2/49}`$ K (i.e. nearly everywhere in a protoplanetary disc) $`\kappa `$ practically does not depend on density, so we assume $`a=0`$. As before, we assume that the viscosity coefficient is given by equation (7).
Let $`\delta `$ be the ratio of the total disc half-thickness to the active layer thickness $`H_a`$:
$$\delta \frac{H_a+H_{\mathrm{DZ}}}{H_a}.$$
(12)
From the equation of hydrostatic equilibrium in the direction perpendicular do the midplane of the disc we get an approximate formula
$$\frac{P_i}{\rho _i}=\mathrm{\Omega }^2\delta H_a^2=c_{si}^2,$$
(13)
where $`P_i`$, $`\rho _i`$ and $`c_{\mathrm{si}}`$ is, respectively, pressure, density and sound speed at the boundary between the dead zone and the active layer. Combining (13) with equations (9) and (10), and using (7) with $`c_s=c_{\mathrm{si}}`$ we get
$$T_i=\left[\frac{3\kappa _0}{8}\frac{9}{4\sigma }\frac{k_B}{\mu m_\mathrm{H}}\right]^{\frac{1}{3b}}\mathrm{\Sigma }_a^{\frac{2}{3b}}\alpha ^{\frac{1}{3b}}\delta ^{\frac{1}{2(3b)}}\mathrm{\Omega }^{\frac{1}{3b}}.$$
(14)
Inserting equations (9) – (14) into equation (8), we arrive at the following formula for the accretion rate in the active layer:
$$\dot{M}=r^{1/2}\frac{}{r}\left(\mathrm{\Omega }^{\frac{2+b}{3b}}\delta ^{\frac{4+b}{2(3b)}}r^{1/2}\right),$$
(15)
where
$$=12\pi \left[\frac{3\kappa _0}{8}\frac{9}{4\sigma }\right]^{\frac{1}{3b}}\frac{k_B}{\mu m_\mathrm{H}}^{\frac{4b}{3b}}\mathrm{\Sigma }_a^{\frac{5b}{3b}}\alpha ^{\frac{4b}{3b}}.$$
(16)
In the following, we shall derive an equation which relates the angular rotational velocity $`\mathrm{\Omega }`$ to the thickness of the disc.
Let us assume that the disc is thin ($`\delta H_ar`$), and neglect the dependence of $`\mathrm{\Omega }`$ on $`z`$ (see e.g. Urpin 1983). We may write
$$\mathrm{\Omega }^2=\mathrm{\Omega }_0^2+\frac{1}{r\rho _m}\frac{P_m}{r},$$
(17)
where $`P_m=c_{\mathrm{si}}^2\rho _m`$, $`\rho _m`$ and $`c_{\mathrm{si}}`$ are, respectively, mid-plane values of pressure, density, and sound speed. Note that because there is no heat generation in the dead zone, the mid-plane sound speed is the same as at the boundary between the dead zone and the active layer (i.e. at each $`r`$ the dead zone is isothermal along a line perpendicular to the midplane).
The mid-plane density $`\rho _m`$ is given by the vertical hydrostatic equilibrium
$$\rho _m=\rho _i\mathrm{exp}\left(\frac{\mathrm{\Omega }^2H_{\mathrm{DZ}}^2}{2c_{\mathrm{si}}^2}\right)=\rho _i\mathrm{exp}\left(\frac{(\delta 1)^2}{2\delta }\right),$$
(18)
and for $`\delta 1`$ we have
$$\rho _m=\rho _i\mathrm{exp}\left(\frac{\delta }{2}1\right)$$
(19)
where it was assumed that the mass accumulated in dead zone is already large (i.e. $`\delta 1`$), which allows us to neglect terms of the second order in $`1/\delta `$. Inserting $`\rho _\mathrm{m}`$ and $`P_\mathrm{m}`$ into equation (17) we obtain
$$\mathrm{\Omega }^2=\mathrm{\Omega }_0^2+\frac{c_{\mathrm{si}}^2}{2r}\frac{\delta }{r}+\frac{1}{r}\frac{c_{\mathrm{si}}^2}{r}+\frac{c_{\mathrm{si}}^2}{\rho _ir}\frac{\rho _i}{r}.$$
(20)
The sound speed $`c_{\mathrm{si}}`$ only weakly depends on disc thickness ($`c_{\mathrm{si}}^2\delta ^{1/5}`$ for $`b=0.5`$). The dependence of $`\rho _i`$ on $`\delta `$ is even weaker ($`\rho _i\delta ^{1/10}`$ for $`b=0.5`$). Therefore, we assume that for a small perturbation of $`\delta `$ the last two rhs terms in equation (20), which are proportional to derivatives of $`c_{\mathrm{si}}^2`$ and $`\mathrm{log}(\rho _i)`$, are small compared to the second rhs term, which is directly proportional to the derivative of $`\delta `$.
The final relation between $`\mathrm{\Omega }`$ and $`\delta `$ is
$$\mathrm{\Omega }^2=\mathrm{\Omega }_0^2+\frac{c_{\mathrm{si}}^2}{2r}\frac{\delta }{r}.$$
(21)
Equations (15) and (21) will be used in the next subsections to derive the dispersion relation for the perturbations of the disc.
### 4.2 The linear analysis
Let us perturb $`\delta `$ at a radius $`r_0`$, assuming that in a small region $`(r_0R,r_0+R)`$, $`Rr`$ $`\delta `$ consists of the unperturbed part $`\delta _0`$ (which in that region may be regarded as independent of $`r`$) plus a cosine term with a wavenumber $`k`$
$$\delta =\delta _0+\delta _k\mathrm{cos}(kR),$$
(22)
where $`\delta _k`$ is the amplitude of the perturbation.
According to (21), the angular velocity consists of the Keplerian part $`\mathrm{\Omega }_0`$ and the pressure-correction $`\mathrm{\Omega }_1`$. Neglecting terms of the second order in $`\mathrm{\Omega }_1`$, and using equation (21) we obtain
$$\mathrm{\Omega }^2=\mathrm{\Omega }_0^2+2\mathrm{\Omega }_0\mathrm{\Omega }_1\mathrm{\Omega }_1=\frac{c_{\mathrm{si}}^2}{4r\mathrm{\Omega }_0}\frac{\delta }{r}.$$
(23)
The first and second derivatives of $`\delta `$ and $`\mathrm{\Omega }`$ are
$$\begin{array}{cc}\frac{\delta }{R}=\delta _kk\mathrm{sin}(kR)\hfill & \frac{\mathrm{\Omega }}{R}=\frac{3\mathrm{\Omega }_0}{2r}\frac{c_{\mathrm{si}}^2}{4r\mathrm{\Omega }_0}\delta _kk^2\mathrm{cos}(kR)\hfill \\ \frac{^2\delta }{R^2}=\delta _kk^2\mathrm{cos}(kR)\hfill & \frac{^2\mathrm{\Omega }}{R^2}=\frac{15\mathrm{\Omega }_0}{4r^2}+\frac{c_{\mathrm{si}}^2}{4r\mathrm{\Omega }_0}\delta _kk^3\mathrm{sin}(kR)\hfill \end{array}$$
(24)
Since we only want to obtain a rough idea about the growth rate of the perturbation with a specific $`k`$, we may approximate $`\mathrm{cos}(kR)`$ with 1 and $`\mathrm{sin}(kR)`$ with 0. Then the disc thickness $`\delta `$, the angular velocity $`\mathrm{\Omega }`$ and their derivatives are:
$$\delta =\delta _0+\delta _k,\frac{\delta }{R}=0,\frac{^2\delta }{R^2}=\delta _kk^2$$
(25)
$$\mathrm{\Omega }=\mathrm{\Omega }_0,\frac{\mathrm{\Omega }}{R}=\frac{3\mathrm{\Omega }_0}{2r_0}\frac{c_{\mathrm{si}}^2}{4r\mathrm{\Omega }_0}\delta _kk^2,\frac{^2\mathrm{\Omega }}{R^2}=\frac{15\mathrm{\Omega }_0}{4r_0}.$$
(26)
Due to mass accumulation, the surface density of the dead zone $`\dot{\mathrm{\Sigma }}_{\mathrm{DZ}}`$ grows at a rate
$$\dot{\mathrm{\Sigma }}_{\mathrm{DZ}}=\frac{1}{2\pi r}\frac{\dot{M}}{r}.$$
(27)
Inserting equation (15) into (27) and using approximations (25) and (26) we obtain
$`\dot{\mathrm{\Sigma }}_{\mathrm{DZ}}`$ $`=`$ $`{\displaystyle \frac{}{2\pi r_0}}[{\displaystyle \frac{3}{2}}{\displaystyle \frac{2+b}{3b}}\mathrm{\Omega }_0^{\frac{5+2b}{3b}}\delta ^{\frac{4+b}{2(3b)}}{\displaystyle \frac{\mathrm{\Omega }}{R}}`$ (28)
$`+`$ $`{\displaystyle \frac{4+b}{2(3b)}}r_0\mathrm{\Omega }_0^{\frac{2+b}{3b}}\delta ^{\frac{10+3b}{2(3b)}}{\displaystyle \frac{^2\delta }{R^2}}`$
$`+`$ $`{\displaystyle \frac{2+b}{3b}}{\displaystyle \frac{5+2b}{2(3b)}}r_0\mathrm{\Omega }_0^{\frac{8+3b}{3b}}\delta ^{\frac{4+b}{2(3b)}}\left({\displaystyle \frac{\mathrm{\Omega }}{R}}\right)^2`$
$`+`$ $`{\displaystyle \frac{2+b}{3b}}r_0\mathrm{\Omega }_0^{\frac{5+2b}{3b}}\delta ^{\frac{4+b}{2(3b)}}{\displaystyle \frac{^2\mathrm{\Omega }}{R^2}}]`$
The linearized unperturbed part of the last equation is
$$\dot{\mathrm{\Sigma }}_{\mathrm{DZ},0}=\frac{3}{8\pi r_0^2}\frac{2+b}{3b}\frac{9+4b}{3b}\mathrm{\Omega }_0^{\frac{2+b}{3b}}\delta _0^{\frac{4+b}{2(3b)}},$$
(29)
and the linearized equation which describes the evolution of surface density perturbation with wavenumber $`k`$ is
$`\dot{\mathrm{\Sigma }}_{\mathrm{DZ},k}`$ $`=`$ $`{\displaystyle \frac{}{2\pi r_0^2}}\delta _k`$ (30)
$`\times `$ $`[{\displaystyle \frac{3}{4}}{\displaystyle \frac{2+b}{3b}}{\displaystyle \frac{4+b}{2(3b)}}{\displaystyle \frac{9+4b}{3b}}\mathrm{\Omega }_0^{\frac{2+b}{3b}}\delta _0^{\frac{10+3b}{2(3b)}}`$
$`+`$ $`{\displaystyle \frac{3c_{\mathrm{si}}^2}{4}}{\displaystyle \frac{2+b}{3b}}{\displaystyle \frac{73b}{2(3b)}}\mathrm{\Omega }_0^{\frac{8+3b}{3b}}\delta _0^{\frac{4+b}{2(3b)}}k^2`$
$``$ $`{\displaystyle \frac{4+b}{2(3b)}}r_0^2\mathrm{\Omega }_0^{\frac{2+b}{3b}}\delta _0^{\frac{10+3b}{2(3b)}}k^2].`$
The surface density of the dead zone $`\mathrm{\Sigma }_{\mathrm{DZ}}`$ is related to the disc thickness $`\delta `$ through
$`\mathrm{\Sigma }_{\mathrm{DZ}}`$ $`=`$ $`2{\displaystyle _0^{H_{\mathrm{DZ}}}}\rho (z)𝑑z=2{\displaystyle _0^{H_{\mathrm{DZ}}}}\rho _\mathrm{m}\mathrm{exp}\left({\displaystyle \frac{\mathrm{\Omega }^2z^2}{2c_{\mathrm{si}}^2}}\right)𝑑z`$ (31)
$`=`$ $`2{\displaystyle _0^{H_{\mathrm{DZ}}}}\rho _i\mathrm{exp}\left({\displaystyle \frac{\mathrm{\Omega }^2H_{\mathrm{DZ}}^2}{2c_{\mathrm{si}}^2}}\right)\mathrm{exp}\left({\displaystyle \frac{\mathrm{\Omega }^2z^2}{2c_{\mathrm{si}}^2}}\right)𝑑z`$
where we used the mid-plane density given by (19). The integral can be evaluated using an error function to yield
$$\mathrm{\Sigma }_{\mathrm{DZ}}=\frac{\sqrt{2\pi }c_{\mathrm{si}}}{\mathrm{\Omega }}\rho _i\mathrm{exp}\left(\frac{\delta }{2}1\right)\mathrm{erf}\left(\frac{\delta 1}{\sqrt{2\delta }}\right).$$
(32)
Remembering that $`\delta 1`$ we get
$$\mathrm{\Sigma }_{\mathrm{DZ}}=\mathrm{\Sigma }_a\sqrt{\delta }\mathrm{exp}\left(\frac{\delta }{2}1\right)$$
(33)
Differentiating the last equation with respect to time leads to a formula connecting $`\dot{\mathrm{\Sigma }}_{\mathrm{DZ}}`$ with $`\dot{\delta }`$
$$\dot{\mathrm{\Sigma }}_{\mathrm{DZ}}=\mathrm{\Sigma }_a\frac{\dot{\delta }\sqrt{\delta }}{2}\left(\frac{1}{\delta }+1\right)\mathrm{exp}\left(\frac{\delta }{2}1\right)$$
(34)
whose unperturbed part and equation (29) can be combined into an equation
$`\dot{\delta }_0`$ $`=`$ $`{\displaystyle \frac{3}{4\pi r_0^2\mathrm{\Sigma }_a}}{\displaystyle \frac{2+b}{3b}}{\displaystyle \frac{9+4b}{3b}}\mathrm{\Omega }_0^{\frac{2+b}{3b}}\delta _0^{\frac{7+2b}{2(3b)}}`$ (35)
$`\mathrm{exp}\left(1{\displaystyle \frac{\delta _0}{2}}\right)\left(1+{\displaystyle \frac{1}{\delta _0}}\right)^1,`$
which describes the evolution of the unperturbed disc thickness due to the accumulation of mass. Next, from equation (30), and equation (34) written for a specific wavenumber $`k`$ we get the dispersion relation which describes the growth rate of the perturbation with the wavenumber $`k`$.
$`\dot{\delta }_k`$ $`=`$ $`{\displaystyle \frac{}{\pi r_0^2\mathrm{\Sigma }_a}}\mathrm{exp}\left(1{\displaystyle \frac{\delta _0}{2}}\right)\left(1+{\displaystyle \frac{1}{\delta _0}}\right)^1\delta _k`$ (36)
$`\times `$ $`[{\displaystyle \frac{3}{4}}{\displaystyle \frac{2+b}{3b}}{\displaystyle \frac{9+4b}{3b}}\mathrm{\Omega }_0^{\frac{2+b}{3b}}\delta _0^{\frac{7+2b}{2(3b)}}`$
$`\times `$ $`\left({\displaystyle \frac{\delta _0+2}{2(\delta _0+1)}}+{\displaystyle \frac{4+b}{2(3b)}}{\displaystyle \frac{1}{\delta _0}}\right)`$
$`+`$ $`{\displaystyle \frac{3c_{\mathrm{si}}^2}{4}}{\displaystyle \frac{2+b}{3b}}{\displaystyle \frac{73b}{2(3b)}}\mathrm{\Omega }_0^{\frac{8+3b}{3b}}\delta _0^{\frac{7+2b}{2(3b)}}k^2`$
$``$ $`{\displaystyle \frac{4+b}{2(3b)}}r_0^2\mathrm{\Omega }_0^{\frac{2+b}{3b}}\delta _0^{\frac{13+4b}{2(3b)}}k^2]`$
The perturbation growth rate diverges at short wavelengths. This is because our simple analysis does not include any damping mechanisms. In reality, however, thin rings would diffuse due to the thermal motion of particles.
Note also that that the radial pressure scale height cannot become smaller than the vertical one. Therefore, rings with circular r-z profiles should form, which is in agreement with our numerical model (see Fig. 1).
Fig. 4 shows growth rates of the disc thickness $`\delta `$ in the four inner-most rings. They are compared to the growth rates given by the dispersion relation (36). The growth rates in the analytical model are more than one order of magnitude higher. The possible explanation of this discrepancy can be that the analytical model does not allow for heat transfer from the ring into the active layer above it or for convective flows that mix the mass inside the rings. The radiative transfer in the radial direction may also be important. All those processes make the dip in the average viscosity shallower and effectively decrease the growth rate of the instability.
### 4.3 The influence of stellar irradiation
In this section we estimate how strongly the ring instability is affected by the radiation flux from the star. In the preceding section the stellar flux was not included self-consistently since the linear analysis presented there serves as the comparison model for the numerical simulation which does not include irradiation. On the other hand, with the irradiation included explicitly the calculations would become too complex (or just impossible) to perform.
Firstly, we estimate how important the irradiation is for the unperturbed disc. We do it by computing the change of the temperature at the boundary between the dead zone and the surface layer (from which the thickness of the surface layer and the viscosity are determined).
Stellar heating can be included in Eq. (10) in the following way
$$T_i^4=\frac{3}{8}\tau T_e^4+WT_{}^4,$$
(37)
where $`T_e`$ is the effective temperature resulting solely from the viscous dissipation, $`T_{}`$ is the temperature of the stellar photosphere, and $`W`$ is the dilution factor which depends on stellar radius, distance from the star, and geometry of the disc (Hubeny, 1990). For a conical disc (in which the aspect ratio $`H/r`$ does not depend on $`r`$) and $`rR_{}`$ we have
$$W=\frac{2}{3\pi }\left(\frac{R_{}}{r}\right)^3$$
(38)
(see e.g. Ruden & Pollack 1991).
Evaluating the rhs terms in Eq. (37) for our disc model and typical parameters of a T Tauri star ($`R_{}=3`$ R, $`T_{}=4400`$ K) we find that the first (viscous) term always dominates. More susceptible to effects of irradiation are flaring discs, in which $`H/rr^\gamma `$. The flaring index $`\gamma `$ depends on the model, e.g. a vertically isothermal model in which the intercepted stellar flux is equal to the blackbody emission from the disc yields $`\gamma =2/7`$ (Chiang & Goldreich, 1997). However, even in such model the irradiation term dominates only for ($`r10`$ AU), where the grazing angle becomes sufficiently large.
Effects of irradiation may be more important for the ring instability itself. The inner edge of a growing ring is more strongly heated by stellar radiation, because the grazing angle $`\alpha _{\mathrm{gr}}`$ is larger there. As a result, the temperature and the thickness of the active layer increase locally, leading to an increased viscosity. Then, the dip in the accretion rate, which is responsible for the ring growth, becomes shallower or it may even disappear entirely.
The importance of this effect may be roughly estimated by comparing the amplitude of viscosity increase due to stellar irradiation to the amplitude of viscosity decrease due to the stronger vertical component of the gravitational force (the latter effect is explained in Fig. 2).
Let us assume a ring-like perturbation of the disc thickness with an amplitude $`\delta _k`$ and a wavelength $`\lambda =2\pi /k`$. Since the most unstable wavelength is $`\lambda _{\mathrm{max}}\delta _0+\delta _k`$, let us parametrize the perturbation wavelength by the dimensionless value $`l=\lambda /\lambda _{\mathrm{max}}`$. The grazing angle at which the stellar radiation strikes the inner edge of the ring is
$$\alpha _{\mathrm{gr}}\frac{\delta _k}{l(\delta _0+\delta _k)}+\frac{2}{3\pi }\frac{R_{}}{r}.$$
(39)
Evaluating the dilution factor $`W=\alpha _{\mathrm{gr}}(R_{}/r)^2`$, inserting it into Eq. (37), and combining it with Eqs. (9), (13) and (7) we obtain the viscosity in a form
$`\nu (\delta _k,\delta _0,r,l)`$ $`=`$ $`N_1(r)(\delta _0+\delta _k)^{\frac{4+b}{2(3b)}}+N_2(r)`$ (40)
$`\times `$ $`(\delta _0+\delta _k)^{1/2}\left[{\displaystyle \frac{\delta _k}{l(\delta _0+\delta _k)}}+N_3(r)\right]^{1/4}`$
where
$$N_1(r)=\left(\frac{3\kappa _0}{8}\frac{9}{4\sigma }\right)^{\frac{1}{3b}}\alpha ^{\frac{4b}{3b}}\left(\frac{k_B}{\mu m_\mathrm{H}}\right)^{\frac{4b}{3b}}\mathrm{\Omega }^{\frac{2+b}{3b}},$$
(41)
$$N_2(r)=\frac{\alpha }{\mathrm{\Omega }}\frac{k_B}{\mu m_\mathrm{H}}\left(\frac{R_{}}{r}\right)^{1/2}T_{},$$
(42)
and
$$N_3(r)=\frac{2}{3\pi }\frac{R_{}}{r}.$$
(43)
Initially, the function $`\nu (\delta _k)`$ always grows as $`\delta _k`$ increases. However, depending on the remaining parameters ($`\delta _0`$, $`r`$ and $`l`$), it may start to decrease and quickly drop below the initial value $`\nu (\delta _k=0)`$. Therefore, a very small perturbation of the disc thickness is always stabilized by the stellar irradiation, but if the amplitude of the perturbation reaches some value, the ring instability may start to work. This critical value $`\delta _{k,\mathrm{crit}}`$ ($`\nu (\delta _{k,\mathrm{crit}})=\nu (0)`$) strongly depends on $`l`$: it is smaller for higher $`l`$, where the grazing angle is smaller. The dependence on the remaining parameters ($`r`$ and $`\delta _0`$) for $`l=3`$ is shown by Fig. 5.
We see that the stabilizing effect of the irradiation may become unimportant for broad rings (large $`l`$), small $`\delta _0`$ (less mass accumulated in the dead zone) and/or at small radii.
## 5 Discussion
We described a new instability in layered accretion discs. The accretion rate in such discs is in general a function of radius. As a result, mass may accumulate in the non-viscous dead zone near the mid-plane of the disc. However, a small ring-like perturbation of the dead zone thickness leads to the deviation of the rotational velocity which results in non-uniform accumulation rate of the mass supporting growing of that ring perturbation. This ring instability may eventually lead to a decomposition of the dead zone into rings.
We observed the formation of such rings in the 2D axially symmetric radiation-hydrodynamic simulation of the layered disc. To illustrate how the instability works, we give its approximate analytical description. According to the analytical results, the narrowest rings grow most rapidly. Therefore, we may expect the formation of the radially thin rings whose radial extent will be given just by the thermal pressure of the gas. The comparison of ring sizes shows a reasonable agreement between the numerical simulation and the analytical model.
The irradiation by the central star may substantially decelerate or even stop the growth of the instability in some regions of the disc. However, broad rings are less affected, and the effects of irradiation become less important in thin discs and/or at small distances from the star. On the other hand, once the innermost ring has developed, and the flaring index $`\gamma `$ is not too high, the disc at larger radii is shadowed and more rings may develop there. If the disc is truncated (e.g. by the magnetosphere of the star) the shadowing effect may also be caused by its inner edge, allowing for an efficient growth of the instability.
The rings created by this instability may be important in terms of planet formation, because they can be the places where the solid particles (gravel and boulders) accumulate. It is because of the higher rotational velocity at the inner edge of the ring and the lower one at the outer edge. The drag force which makes the grains to move inwards is smaller at the inner edge of the ring and higher at the outer edge. The solid particles may merge into larger bodies (planetesimals) necessary for the formation of planets.
Massive rings are subject to a hydrodynamical instability in 3D, e.g. with respect to non-axisymmetric perturbations (Papaloizou & Pringle, 1984, 1985). This instability occurs if the rotational angular velocity decreases with radius steeper than $`r^\sqrt{3}`$. In our simulation, this occurs for the innermost ring at a time of $`280`$ yr. In the non-viscous case such ring would most likely decay into so-called ”planets”, i.e. large scale vortices, as found in numerical simulations by Hawley (1987), and further investigated by Goodman, Naryan & Goldreich (1987). The fate of such a ring in 3D viscous hydro-simulations including the effects of layered accretion still have to be performed.
Our model assumes zero viscosity in the dead zone. However, there are some indications (Fleming & Stone, 2003) that even in the dead one there might be some turbulence present, induced by waves originating in the active parts of the disc. In such a case, the excess of surface density could result in a higher accretion rate in the rings, and the growth of the rings would be stopped or even they might not form at all. This issue will be addressed in a forthcoming paper.
## 6 Acknowledgments
This research was supported in part by the European Community’s Human Potential Programme under contract HPRN-CT-2002-00308,PLANETS. RW acknowledges financial support by the Grant Agency of the Academy of Sciences of the Czech Republic under the grants AVOZ 10030501 and B3003106. MR acknowledges the support from the grant No. 1 P03D 026 26 from the Polish Ministry of Science. |
warning/0506/gr-qc0506071.html | ar5iv | text | # Discrete quantum gravity in the framework of Regge calculus formalism
## Abstract
An approach to the discrete quantum gravity based on the Regge calculus is discussed which was developed in a number of our papers. Regge calculus is general relativity for the subclass of general Riemannian manifolds called piecewise flat ones. Regge calculus deals with the discrete set of variables, triangulation lengths, and contains continuous general relativity as a particular limiting case when the lengths tend to zero. In our approach the quantum length expectations are nonzero and of the order of Plank scale $`10^{33}cm`$. This means the discrete spacetime structure on these scales.
PACS numbers: 04.60.-m Quantum gravity
1.Introduction
An interest to formulation of general relativity (GR) in a discrete form is provided, not in the last place, by complexity of the theory. In classical aspect, rewriting the essentially nonlinear equations of the theory, Einstein equations, in terms of a discrete set of physical quantities, that is, discretising them simplifies using numerical methods for solving these equations. In quantum aspect, discretisation might be introduced, as in any other field theory, in order to regularise the originally divergent expressions. However, in the case of GR the following two distinctive features arise. First, according to the standard classification, GR is nonrenormalisable theory, therefore dependence of the result on the specific way of the regularisation cannot be removed by the renormalisation. Consequently, here discretisation must not only be a mathematical approximation like finite-difference one of the originally continuum theory but present some realizable physics specifying the form of the theory at small distances. Second, the covariance with respect to the arbitrary coordinate transformatioms is specific for GR, and this property is badly consistent with quantum theory in which the time plays the special role. To avoid this difficulty, one might try to formulate GR in the explicitly coordinateless form.
In the framework of Regge calculus suggested in 1961 the exact GR is considered on the particular case of general Riemannian spacetime, the so-called piecewise flat manifolds which are flat everywhere with exception of the subset of points of zero measure. Any such spacetime can be represented as collection of a number of the flat 4-dimensional simplices(tetrahedrons). In the $`n`$-dimensional case the $`n`$-dimensional simplices $`\sigma ^n`$ are considered. The $`n`$-dimensional simplex $`\sigma ^n`$ contains the $`n+1`$ vertices each of which being connected by the edges with the other $`n`$ vertices. All the geometrical characteristics of the $`n`$-simplex are uniquely defined by the (freely chosen) lengths of the $`n\frac{n+1}{2}`$ edges of it. Geometry of the Regge spacetime is defined by the freely chosen lengths of all edges (or 1-simplices). The linklengths of the two $`n`$-simplices sharing some $`(n1)`$-simplex as their common face should coincide on this face. As for all the $`n`$-simplices containing some $`(n2)`$-simplex as $`(n2)`$-dimensional face this manifold cannot be embedded into the flat $`n`$-dimensional spacetime if linklengths are freely chosen since the sum of hyperdihedral angles at $`(n2)`$-face in all the $`n`$-simplices meeting at this face is $`2\pi `$ \- $`\alpha `$ where the so-called angle defect $`\alpha `$ does not generally vanish. As a result of the parallel transport of a vector along a closed contour contained in the given $`n`$-simplices and enclosing the given $`(n2)`$-simplex, the vector is rotated by the angle $`\alpha `$. This corresponds to a $`\delta `$-function-like curvature distribution with support on $`(n2)`$-simplices proportional to the angle defects on these simplices. The action for the 4-dimensional Regge gravity is proportional to
$$\underset{\sigma ^2}{}\alpha _{\sigma ^2}|\sigma ^2|,$$
(1)
where $`|\sigma ^2|`$ is the area of a triangle (the 2-simplex) $`\sigma ^2`$, $`\alpha _{\sigma ^2}`$ is the angle defect on this triangle, and summation run over all the 2-simplices $`\sigma ^2`$. It is shown in the work that the action (1) can be obtained from the expression
$$\frac{1}{2}R\sqrt{g}\mathrm{d}^4x,$$
(2)
to which the Einstein action is proportional when passing to the $`\delta `$-function limit of the curvature $`R`$ distribution. Thus Regge calculus is GR in which some degrees of freedom are ”frozen”, that is, the so-called minisuperspace theory for GR. Thereby the first of above mentioned requirements is satisfied, namely, Regge manifold is a particular (although somewhat singular) case of general Riemannian manifold. Besides that, the mutual location of vertices (the 0-dimensional simplices) and thereby the geometry is uniquely fixed by the values of invariant lengths of the edges (the 1-simplices $`\sigma ^1`$) which therefore play the role of the field variables. Thus the second requirement of the coordinateless description is also fullfilled.
Although Regge calculus is only some subset in the configuration superspace of GR, this subset is dense in this superspace. That is, each nonsingular Riemannian manifold can be approximated with arbitrarily high accuracy by an appropriately chosen Regge manifold. To construct such the Regge manifold one divides, for example, the Riemannian manifold into the sufficiently small regions topologically equivalent to the simplices $`\sigma ^4`$, the edges of these being geodesics. As the piecewise-flat manifold of interest the manifold of such type can be taken possessing the same scheme (topology) of connection of the different vertices by the edges and the linklengths being geodesic lengths. It is shown in the paper that the Einstein action (2) follows as limiting case of the Regge action (1) for such the approximating spaces if a typical edge length for triangulation length) tends to zero. A more general statement is that the so-called Lipshitz-Killing curvatures converge to their continuum counterparts in the sense of measures if the decomposition into the 4-simplices becomes finer and finer, that is, it is integrals of the considered values over the spacetime regions which converge (to the integrals of the continuum counterparts). The volume of a spacetime region, contribution of the region to the Einstein action and to the Gauss-Bonnet topological term are examples of such the integrals.
Regge calculus possesses exact discrete analogs of many quantities which can be defined in the continuum GR. The first example are the Einstein equations whose discrete analog was derived by Regge upon varying the action (1) over the linklengths. It turns out that the variation of $`\alpha _{\sigma ^2}`$ does not contribute into (1), and the equation obtained by variation of a particular edge $`\sigma ^1`$ takes the form
$$\underset{\sigma ^2\sigma ^1}{}\alpha _{\sigma ^2}\mathrm{cot}\vartheta (\sigma ^1,\sigma ^2)=0.$$
(3)
Here $`\vartheta (\sigma ^1,\sigma ^2)`$ is the angle in the triangle $`\sigma ^2`$ opposite to the edge $`\sigma ^1`$ while summation is over all the triangles sharing $`\sigma ^1`$. Evidently, the discrete coordinateless formulation in terms of physical quantities (lengths) is an ideal means for numerical simulations, and originally Regge calculus was just used for numerical analysis of the Einstein equations .
However Regge calculus is of the maximal interest if applied to quantum gravity. In this aspect the main problem is in constructing the Hamiltonian formalism analogous to the Arnowitt-Deser-Misner construction in the continuum GR . In accordance with their result, the GR Lagrangian can be reduced to the form
$$L=\underset{A}{}p_A\dot{q}_A\underset{\alpha }{}\lambda _\alpha \mathrm{\Phi }_\alpha (p,q)$$
(4)
with the canonical variables $`p_A`$, $`q_A`$ and yet another variables $`\lambda _\alpha `$ playing the role of the Lagrange multipliers whose dynamics is not fixed by the equations of motion. Thus, GR is a theory described by the set of constraints $`\mathrm{\Phi }_\alpha (p,q)`$ = 0 and zero Hamiltonian. In the case of coordinateless Regge calculus theory we need to particularly return to the coordinate description but with respect to only time coordinate $`t`$, the discrete field distribution (here lengths) having been reduced to that one smooth in $`t`$. Passing to the resulting so-called $`(3`$ $`+`$ $`1)`$ Regge calculus had been undertaken in a number of works \- . Usually, the authors of the cited papers tried to define in some way the discrete analogs of the variables $`p_A`$, $`q_A`$ and of the constraints $`\mathrm{\Phi }_\alpha (p,q)`$, the largest attention having been paid to providing the algebra of Poisson brackets being maximally close to that in the continuum GR. If we stick to the strategy requiring to deal with a realizable Riemannian manifold at each stage, the $`(3`$ $`+`$ $`1)`$ Regge calculus follows as the limiting case of the 4-dimensional Regge calculus while sizes of the 4-simplices tend to zero in some direction taken as direction of time. This limiting procedure had been studied in the papers , although not all the degrees of freedom could be taken into account in these works. The reason is in the singular nature of description of Regge manifold with the help of the linklengths when the scale of sizes along any direction tends to zero. As an example, one can imagine a triangle one of the edges of which is infinitesimal; then infinitesimal variations of the two other (finite) edge lengths lead to finite variations of the angles. As a result, description in terms of only the lengths had left some (singular) degrees of freedom missing, and therefore not all the discrete counterparts of the constraints $`\mathrm{\Phi }_\alpha (p,q)`$ could be found.
2.The problem of constructing the quantum measure in Regge calculus
The singular nature of passing to the continuous time is thus connected with using only the lengths as fundamental set of variables in Regge calculus. While we are studying the quantum measure on the completely discrete Regge manifold, this circumstance is not of importance for us. However the issue concept of quantum theory which might be used to construct quantum measure is canonical quantization. The latter is defined just in the continuous time. Therefore the quantum measure of interest should be defined from the requirement that it would tend in some sense to the canonical quantization measure (Feynman path integral) whenever continuum limit is taken along any of the coordinates, the coordinate chosen playing the role of time. In other words, the continuous time limit serves as probing tool for defining the quantum measure in the completely discrete Regge calculus.
The way to avoid singularities in the continuous time limit is in extending the set of variables via adding the new ones having the sense of angles and considered as independent variables. Such the variables are the finite rotation matrices which are the discrete analogs of the connections in the continuum GR.
3.Representation of the Regge calculus in terms of finite rotation matrices as independent variables
The situation considered is analogous to that one occurred when recasting the Einstein action (2) in the Hilbert-Palatini form,
$$\frac{1}{2}R\sqrt{g}\mathrm{d}^4x\frac{1}{8}ϵ_{abcd}ϵ^{\lambda \mu \nu \rho }e_\lambda ^ae_\mu ^b[_\nu +\omega _\nu ,_\rho +\omega _\rho ]^{cd}\mathrm{d}^4x,$$
(5)
where the tetrad $`e_\lambda ^a`$ and connection $`\omega _\lambda ^{ab}`$ = $`\omega _\lambda ^{ba}`$ are independent variables, the equation (5) being reduced to (2) in terms of $`g_{\lambda \mu }`$ = $`e_\lambda ^ae_{a\mu }`$ if we substitute for $`\omega _\lambda ^{ab}`$ solution of the equations of motion for these variables in terms of $`e_\lambda ^a`$. The Latin indices $`a`$, $`b`$, $`c`$, … are the vector ones with respect to the local Euclidean frames which are introduced at the each point $`x`$. The Regge calculus analog of the representation (5) follows if the local Euclidean frames are introduced in all the 4-simplices. Then the analogs of the connection are defined on the 3-simplices $`\sigma ^3`$ and are the matrices $`\mathrm{\Omega }_{\sigma ^3}`$ connecting the frames of the pairs of the 4-simplices $`\sigma ^4`$ sharing the 3-faces $`\sigma ^3`$. These matrices are the finite SO(4) rotations in the Euclidean case (or SO(3,1) rotations in the Lorentzian case) in contrast with the continuum connections $`\omega _\lambda ^{ab}`$ which are the elements of the Lee algebra so(4)(so(3,1)) of this group. This definition includes pointing out the direction in which the connection $`\mathrm{\Omega }_{\sigma ^3}`$ acts (and, correspondingly, the opposite direction, in which the $`\mathrm{\Omega }_{\sigma ^3}^1`$ = $`\overline{\mathrm{\Omega }}_{\sigma ^3}`$ acts), that is, the connections $`\mathrm{\Omega }`$ are defined on the oriented 3-simplices $`\sigma ^3`$.
Also we can define the curvature matrix $`R_{\sigma ^2}`$ on each the 2-simplex $`\sigma ^2`$ as the product of the connections $`\mathrm{\Omega }_{\sigma ^3}^{\pm 1}`$ on the 3-simplices $`\sigma ^3`$ sharing $`\sigma ^2`$ which act in certain direction along the contour enclosing $`\sigma ^2`$ once and contained in these 3-simplices. The matrix $`R_{\sigma ^2}`$ should be rotation around $`\sigma ^2`$ by an angle $`\alpha _{\sigma ^2}`$. Besides the direction along the contour it is necessary to specify the 4-simplex $`\sigma ^4`$ $``$ $`\sigma ^2`$ in which the contour begins and is ended, that is, the simplex in the local Euclidean frame of which we define
$$R_{\sigma ^2}=\underset{\sigma ^3\sigma ^2}{}\mathrm{\Omega }_{\sigma ^3}^{\pm 1}.$$
(6)
The discrete analogs of the connection and curvature had been considered by Bander as functions of the lengths. Our approach is based on treatment of the connections as independent variables and using a representation of the Regge calculus action (1) analogous to the Hilbert-Palatini form of the Einstein action (5). To write out this representation let us define the dual bivector of the triangle $`\sigma ^2`$ in terms of the vectors of its edges $`l_1^a`$, $`l_2^a`$ given in some 4-simplex containing $`\sigma ^2`$,
$$v_{\sigma ^2ab}=\frac{1}{2}ϵ_{abcd}l_1^cl_2^d.$$
(7)
Then the discrete analog of the expression (5) as suggested in our work reads
$$S(v,\mathrm{\Omega })=\underset{\sigma ^2}{}|v_{\sigma ^2}|\mathrm{arcsin}\frac{v_{\sigma ^2}R_{\sigma ^2}(\mathrm{\Omega })}{|v_{\sigma ^2}|}$$
(8)
where we have defined $`AB`$ = $`\frac{1}{2}A^{ab}B_{ab}`$, $`|A|`$ = $`(AA)^{1/2}`$ for the two tensors $`A`$, $`B`$, in particular, $`|v_{\sigma ^2}|`$ = $`|\sigma ^2|`$ is the area of the triangle. It is important that $`v_{\sigma ^2}`$ and $`R_{\sigma ^2}`$ in (8) be defined in the same 4-simplex containing $`\sigma ^2`$. As we can show, when substituting as $`\mathrm{\Omega }_{\sigma ^3}`$ the genuine rotations connecting the neighbouring local frames as functions of the genuine Regge lengths into the equations of motion for $`\mathrm{\Omega }_{\sigma ^3}`$ with the action (8) we get the closure condition for the surface of the 3-simplex $`\sigma ^3`$ (vanishing the sum of the bivectors of its 2-faces) written in the frame of one of the 4-simplices containing $`\sigma ^3`$, that is, the identity. This means that (8) is the exact representation for (1). At the same time the paper has appeared which suggests a representation which reminds (8) but differs from it by replacing $`\mathrm{arcsin}()`$ by $`()`$. This can be considered as representation for an approximate (at small angle defects) Regge calculus with action $`\underset{\sigma ^2}{}|\sigma ^2|\mathrm{sin}\alpha _{\sigma ^2}`$. Here we should point out remarkable feature namely of the exact Regge calculus since it is namely for the action (8) the equations of motion for $`\mathrm{\Omega }_{\sigma ^3}`$ are shown by us to hold exactly as closure condition for the surface of $`\sigma ^3`$. In other words, representation for an approximate Regge action inevitably proves to be, in turn, approximate one.
4.Naturalness of arising area tensor Regge calculus
In the representation in terms of rotation matrices it is possible to pass to the continuous time and develop the canonical formalism in Regge calculus which turns out to possess the second class constraints (that is, noncommuting ones). As a result, Feynman path integral contains the determinant of the Poisson brackets of the second class constraints as a factor which is singular when approaching the flat geometry. The matter is that the Regge geometry usually changes at arbitrary variations of the edge lengths with exception of the flat case in which these variations are the symmetry transformations. In other words, division of the constraints into those ones of the first and the second class is being changed in the flat case. An exception is the 3-dimensional case. Due to the local triviality of the 3-dimensional gravity all the dynamical constraints are of the first class, and therefore the path integral takes a simple form. In this case the problem of constructing the discrete quantum measure as formulated in the subsection 2 can be solved yielding a simple form for this measure .
The conditions of the subsection 2 on the discrete measure are rather restrictive ones, and existence of the solution is not evident. Singularity of the path integral in 4 dimensions in the vicinity of the flat background is by itself not an obstacle for existence of such the solution; crucial is the occurrence of the above determinant factor in the path integral. This factor depends on the variables which are lattice artefacts connected with specific coordinate along which the continuum limit is taken, and it can not be obtained from some universal expression by taking the continuous time limit.
Let us try to modify the 4-dimensional Regge calculus so that it would remind the 3-dimensional case as far as canonical structure is concerned. The 3-dimensional Regge calculus in the representation analogous to (8) has the edge vectors $`𝒍_{\sigma ^1}`$ instead of the area tensors $`v_{\sigma ^2}`$. The edge vectors are the independent variables thus ensuring the local triviality of the 3-dimensional gravity. Contrary to this, area tensors are not independent. For example, tensors of the two triangles $`\sigma _1^2`$, $`\sigma _2^2`$ sharing an edge satisfy the relation
$$ϵ_{abcd}v_{\sigma _1^2}^{ab}v_{\sigma _2^2}^{cd}=0.$$
(9)
The idea is to construct quantum measure first for the system with formally independent area tensors. That is, originally we concentrate on quantization of the dynamics while kinematical relations of the type (9) are taken into account at the second stage.
Area tensor Regge calculus admits solution to the problem of constructing the discrete quantum measure, the latter taking a simple form . Consider the Euclidean case. The Einstein action is not bounded from below, therefore the Euclidean path integral itself requires careful definition. In particular, the result of for vacuum expectations of the functions of our field variables $`v`$, $`\mathrm{\Omega }`$ can be written with the help of the integration over imaginary areas with the help of the formal replacement of the tensors of a certain subset of areas $`\pi `$ over which integration in the path integral is to be performed,
$$\pi i\pi ,$$
in the form
$`<\mathrm{\Psi }(\{\pi \},\{\mathrm{\Omega }\})>`$ $`=`$ $`{\displaystyle \mathrm{\Psi }(i\{\pi \},\{\mathrm{\Omega }\})\mathrm{exp}\left(\underset{\stackrel{t\mathrm{like}}{\sigma ^2}}{}\tau _{\sigma ^2}R_{\sigma ^2}(\mathrm{\Omega })\right)}`$ (10)
$`\mathrm{exp}\left(i{\displaystyle \underset{\stackrel{\stackrel{\mathrm{not}}{t\mathrm{like}}}{\sigma ^2}}{}}\pi _{\sigma ^2}R_{\sigma ^2}(\mathrm{\Omega })\right){\displaystyle \underset{\stackrel{\stackrel{\mathrm{not}}{t\mathrm{like}}}{\sigma ^2}}{}}\mathrm{d}^6\pi _{\sigma ^2}{\displaystyle \underset{\sigma ^3}{}}𝒟\mathrm{\Omega }_{\sigma ^3}`$
$``$ $`{\displaystyle \mathrm{\Psi }(i\{\pi \},\{\mathrm{\Omega }\})d\mu _{\mathrm{area}}(i\{\pi \},\{\mathrm{\Omega }\})},`$
where we define $`AB`$ = $`\frac{1}{2}A^{ab}B_{ab}`$ for the two tensors $`A`$, $`B`$. The equation implies attributing a certain structure to our Regge lattice which suggests constructing it of leaves being themselves the 3-dimensional Regge geometries of the same structure. The leaves are labelled by the values of some coordinate $`t`$. The corresponding vertices in the neighbouring leaves are connected by the $`t`$-like edges, and, in addition, there are the diagonal edges connecting a vertex with the neighbours of the corresponding vertex in the neighbouring leaf. Then it is natural to define the $`t`$-like simplices and the leaf simplices as the simplices either containing a $`t`$-like edge or completely contained in the leaf, respectively, and also the diagonal simplices as all others. Then $`\tau _{\sigma ^2}`$ is $`v_{\sigma ^2}`$ when $`\sigma ^2`$ is $`t`$-like while $`\pi _{\sigma ^2}`$ is $`v_{\sigma ^2}`$ when $`\sigma ^2`$ is not $`t`$-like, that is, is the leaf or diagonal simplex. As far as considered here area tensor Regge calculus is concerned (the area tensors are independent), the $`\pi _{\sigma ^2}`$ can be chosen as dynamical variables, then the $`\tau _{\sigma ^2}`$ should be considered as parameters.
The equation (10) is in many aspects similar to the intuitively expected one for the quantum measure. In particular, the expected from symmetry considerations invariant (Haar) measure on SO(4) $`𝒟\mathrm{\Omega }`$ arises in the formal path integral expression corresponding in the continuous time limit to the canonical quantization with the kinetic term $`\pi _{\sigma ^2}\overline{\mathrm{\Omega }}_{\sigma ^2}\dot{\mathrm{\Omega }}_{\sigma ^2}`$ in the Lagrangian (the connection variables $`\mathrm{\Omega }`$ in the continuous time limit naturally correspond not to the tetrahedra $`\sigma ^3`$ but to the triangles $`\sigma ^2`$).
Specific features of the quantum measure are, first, the absence of the inverse trigonometric functions $`\mathrm{arcsin}`$ in the exponential, whereas the Regge action (8) contains such functions. This is connected with using the canonical quantization at the intermediate stage of derivation: in gravity this quantization is completely defined by the constraints, the latter being equivalent to those ones without $`\mathrm{arcsin}`$ (in some sense on-shell).
Second, there are no integrations over some subset of area tensors, $`\tau _{\sigma ^2}`$, thereby the symmetry between the different triangles turns out to be incomplete. However, this violation of symmetry can be considered as spontaneous one when some a’priori arbitrary direction coordinatised in (10) by $`t`$ turns out to be singled out. The curvature matrices $`R(\mathrm{\Omega })`$ on all but $`t`$-like triangles can be chosen as independent variables, then such matrices on the $`t`$-like triangles are by means of the Bianchi identities the functions of these variables. Integrations over all the area tensors would lead to singularities of the type of $`[\delta (R`$ $``$ $`\overline{R})]^2`$.
This specific feature of the discrete quantum measure, incomplete symmetry with respect to the different coordinate directions, complies with the above conditions of subsection 2 in the following way. In the continuous limit along some coordinate $`x`$ (which may not coincide with $`t`$) the absence of integrations over the tensors of the $`t`$-like triangles will mean some of the simplest kinds of gauge fixing in the limiting measure, namely, fixing the tensors of some subset of the triangles .
With taking into account the properties of the invariant Haar measure and with the negligibly small values of $`\tau _{\sigma ^2}`$ we get factorisation of the quantum measure obtained into the ”elementary” measures on the separate areas (this just corresponds to the local triviality of the theory) of the form
$$\mathrm{exp}(i\pi R)\mathrm{d}^6\pi 𝒟R.$$
(11)
In turn, use the group property SO(4) = SU(2) $`\times `$ SU(2) to split the variables ($`\pi `$ and generator of $`R`$) into the self- and antiselfdual parts, in particular, $`\pi `$ is mapped into the two 3-vectors $`{}_{}{}^{+}𝝅`$, $`{}_{}{}^{}𝝅`$ in the adjoint representation SO(3). As a result, the measure (11) is the product of the two measures each of which act in 3-dimensional configuration space of area vectors,
$$\mathrm{exp}(i^+\pi ^+R)\mathrm{d}^3{}_{}{}^{+}𝝅𝒟^+R\mathrm{exp}(i^{}\pi ^{}R)\mathrm{d}^3{}_{}{}^{}𝝅𝒟^{}R.$$
(12)
As a result, expression for the expectation of any function on the triangle reads
$`<f(\pi )>`$ $`=`$ $`{\displaystyle f(i\pi )\mathrm{d}^6\pi e^{i\pi R}𝒟R}`$
$`=`$ $`{\displaystyle f(\pi )\frac{\nu _2(|^+𝝅|)}{|^+𝝅|^2}\frac{\nu (|^{}𝝅|)}{|^{}𝝅|^2}\frac{\mathrm{d}^3{}_{}{}^{+}𝝅}{4\pi }\frac{\mathrm{d}^3{}_{}{}^{}𝝅}{4\pi }},`$
$`\nu (l)={\displaystyle \frac{l}{\pi }}{\displaystyle \underset{0}{\overset{\pi }{}}}{\displaystyle \frac{\mathrm{d}\phi }{\mathrm{sin}^2\phi }}e^{l/\mathrm{sin}\phi }.`$
In particular, expectations of powers of area squared,
$$|\pi |^2=\pi \pi =\frac{1}{2}(^+𝝅)^2+\frac{1}{2}(^{}𝝅)^2,$$
(14)
and of the dual product,
$$\pi \pi =\frac{1}{2}(^+𝝅)^2\frac{1}{2}(^{}𝝅)^2,$$
(15)
easily follow by means of averaging the powers of $`{}_{}{}^{\pm }𝝅`$,
$$<(^\pm 𝝅)^{2k}>=\frac{4^k(2k+1)!(2k)!}{k!^2}.$$
(16)
5.Reducing to genuine Regge calculus case
Thus, we have come to the finite nonzero area expectation values in area tensor Regge calculus. However, we need the length expectations in the genuine ordinary Regge calculus. The ordinary Regge calculus follows upon imposing unambiguity conditions on the lengths computed in the different 4-simplices. These conditions are equivalent to the continuity conditions for the induced on 3-faces metric. In the configuration space of area tensor Regge calculus these conditions single out some hypersurface $`\mathrm{\Gamma }_{\mathrm{Regge}}`$. The quantum measure can be considered as a linear functional $`\mu _{\mathrm{area}}(\mathrm{\Psi })`$ on the space of functionals $`\mathrm{\Psi }(\{v\})`$ on the configuration space (for our purposes here it is sufficient to restrict ourselves to the functional dependence on the area tensors $`\{v\}`$; the dependence on the connections is unimportant). The physical assumption is that we can consider ordinary Regge calculus as a kind of the state of the more general system with independent area tensors. This state is described by the following functional,
$$\mathrm{\Psi }(\{v\})=\psi (\{v\})\delta _{\mathrm{Regge}}(\{v\}),$$
(17)
where $`\delta _{\mathrm{Regge}}(\{v\})`$ is the (many-dimensional) $`\delta `$-function with support on $`\mathrm{\Gamma }_{\mathrm{Regge}}`$. The derivatives of $`\delta _{\mathrm{Regge}}`$ have the same support, but these violate positivity in our subsequent construction. To be more precise, delta-function is distribution, not function, but can be treated as function being regularised. If the measure on such the functionals exists in the limit when regularisation is removed, this allows to define the quantum measure on $`\mathrm{\Gamma }_{\mathrm{Regge}}`$,
$$\mu _{\mathrm{Regge}}()=\mu _{\mathrm{area}}(\delta _{\mathrm{Regge}}(\{v\})).$$
(18)
Uniqueness of the construction of $`\delta _{\mathrm{Regge}}`$ follows under quite natural assumption of the minimum of lattice artefacts. Let the system be described by the metric $`g_{\lambda \mu }`$ constant in each of the two 4-simplices $`\sigma _1^4`$, $`\sigma _2^4`$ separated by the 3-face $`\sigma ^3`$ = $`\sigma _1^4\sigma _2^4`$ formed by three 4-vectors $`\iota _a^\lambda `$. These vectors also define the metric induced on the 3-face, $`g_{ab}^{}`$ = $`\iota _a^\lambda \iota _b^\mu g_{\lambda \mu }`$. The continuity condition for the induced metric is taken into account by the $`\delta `$-function of the induced metric variation,
$$\mathrm{\Delta }_{\sigma ^3}g_{ab}^{}\stackrel{\mathrm{def}}{=}g_{ab}^{}(\sigma _1^4)g_{ab}^{}(\sigma _2^4).$$
(19)
As for the $`\delta _{\mathrm{Regge}}`$, it is of course defined up to a factor which is arbitrary function nonvanishing at nondegenerate field configurations. In the spirit of just mentioned principle of minimizing the lattice artefacts it is natural to choose this factor in such the way that the resulting $`\delta `$-function factor would depend only on hyperplane defined by the 3-face but not on the form of this face, that is, would be invariant with respect to arbitrary nondegenerate transformations $`\iota _a^\lambda `$ $``$ $`m_a^b\iota _b^\lambda `$. To provide this, the $`\delta `$-function should be multiplied by the determinant $`g_{ab}^{}`$ squared. This gives
$$[\mathrm{det}(\iota _a^\lambda \iota _b^\mu g_{\lambda \mu })]^2\delta ^6(\iota _a^\lambda \iota _b^\mu \mathrm{\Delta }_{\sigma ^3}g_{\lambda \mu })=V_{\sigma ^3}^4\delta ^6(\mathrm{\Delta }_{\sigma ^3}S_{\sigma ^3}).$$
(20)
Here $`S_{\sigma ^3}`$ is the set of the 6 edge lengths squared of the 3-face $`\sigma ^3`$, $`V_{\sigma ^3}`$ is the volume of the face.
Further, the product of the factors (20) over all the 3-faces should be taken. As a result, we have for each edge the products of the $`\delta `$-functions of the variations of its length between the 4-simplices taken along closed contours, $`\delta (s_1s_2)\delta (s_2s_3)\mathrm{}\delta (s_Ns_1)`$ containing singularity of the type of the $`\delta `$-function squared. In other words, the conditions equating (19) to zero on the different 3-faces are not independent. The more detailed consideration allows to cancel this singularity in a way symmetrical with respect to the different 4-simplices (thus extracting irreducible conditions), the resulting $`\delta `$-function factor remaining invariant with respect to arbitrary deformations of the faces of different dimensions keeping each face in the fixed plane spanned by it .
Qualitatively, it is important that our $`\delta `$-function factor (see, for example, the simplest version (20)) automatically turns out to be invariant with respect to overall length scaling. Remind, however, that dynamical variables to be averaged are $`\pi _{\sigma ^2}`$ but not $`\tau _{\sigma ^2}`$. The more detailed analysis shows that upon fixing the scale of the tensors $`\tau _{\sigma ^2}`$ at the level $`\epsilon `$ $``$ 1 the $`\delta `$-function factor is invariant also with respect to overall scaling of only the dynamical variables $`\pi _{\sigma ^2}`$. This leads to the finite nonzero length expectation values in the ordinary Regge calculus as far as the area expectation values in area tensor Regge calculus are finite and nonzero .
Strictly speaking, when passing from area tensor Regge calculus to the ordinary Regge calculus we need first to impose the conditions ensuring that tensors of the 2-faces in the given 4-simplex define a metric in this simplex. These conditions of the type of (9) can be easily written in general form. Let a vertex of the given 4-simplex be the coordinate origin and the edges emitted from it be the coordinate lines $`\lambda `$, $`\mu `$, $`\nu `$, $`\rho `$, …= 1, 2, 3, 4. Then the (ordered) pair $`\lambda \mu `$ means the (oriented) triangle formed by the edges $`\lambda `$, $`\mu `$. The conditions of interest take the form
$$ϵ_{abcd}v_{\lambda \mu }^{ab}v_{\nu \rho }^{cd}ϵ_{\lambda \mu \nu \rho }.$$
(21)
The 20 equations (21) define the 16-dimensional surface $`\gamma (\sigma ^4)`$ in the 36-dimensional configuration space of the six antisymmetric tensors $`v_{\lambda \mu }^{ab}`$<sup>1</sup><sup>1</sup>1There are also the linear constraints of the type $`\pm v`$ = 0 providing closing surfaces of the 3-faces of our 4-simplex. These are assumed to be already resolved.. The factor of interest in quantum measure is the product of the $`\delta `$-functions with support on $`\gamma (\sigma ^4)`$ over all the 4-simplices $`\sigma ^4`$. The covariant form of the constraints (21) with respect to the world index means that these $`\delta `$-functions are the scalar densities of a certain weight with respect to the world index, that is, the scalars up to powers of the volume of the 4-simplex $`V_{\sigma ^4}`$. Therefore introducing the factors of the type $`V_{\sigma ^4}^\eta `$ we get the scalar at some parameter $`\eta `$. Namely, the product of the factors
$$\underset{\sigma ^4}{}V_{\sigma ^4}^\eta \delta ^{21}(ϵ_{abcd}v_{\lambda \mu |\sigma ^4}^{ab}v_{\nu \rho |\sigma ^4}^{cd}V_{\sigma ^4}ϵ_{\lambda \mu \nu \rho })dV_{\sigma ^4}$$
(22)
at $`\eta `$ = 20 is the scale invariant value as required by the principle of the minimum of the lattice artefacts (that is, the factor of interest should not depend on the size of the 4-simplex). As a result, the conclusion made in the previous paragraph that the length expectation values in ordinary Regge calculus are finite and nonzero as soon as the area expectation values in area tensor Regge calculus are finite and nonzero remains valid .
6.Conclusion
Thus, our approach to quantization of Regge calculus ”from the first principles” includes the following steps and conditions.
1. Constructing the quantum measure which is reduced to the Feynman path integral corresponding to the canonical quantization in the continuous time limit irrespectively of what coordinate is taken as time.
2. Using the exact representation of the Regge action in terms of the rotation matrices as independent variables.
3. Extending the configuration space of the theory by considering area tensors as independent variables (considering the so-called area tensor Regge calculus).
4. Reducing the quantum measure from area tensor Regge calculus to the hypersurface corresponding to the ordinary Regge calculus under assumption of the minimal lattice artefacts, that is, minimal dependence on the form and size of the simplices.
As a result, the length expectations are found to be of the order of the Plank scale $`10^{33}cm`$. Were these values vanishing, this would mean that the quantum measure is saturated by the arbitrarily small lengths, that is, by the smooth Riemannian manifolds, and we would return to the continuum GR. Here the remarkable property of the Regge calculus is displayed: it is minisuperspace GR theory, that is, it is exact GR for a class of certain (piecewise flat) spacetimes. Therefore Regge calculus in quantum theory does not mean complete exclusion of the continuum GR (rather it contains the continuum GR as the limiting point), but rather presents an alternative way of description of the system with the help of the triangulation lengths. Our result, the nonzero length expectations, means adequacy of namely such the description. The GR becomes discrete at the Plank scale dynamically, that is, as a result of competition between the different contributions into the functional integral including also that one from the smooth manifolds.
I am grateful to I.B. Khriplovich for attention to the work and discussion. The present work was supported in part by the Russian Foundation for Basic Research through Grant No. 05-02-16627-a. |
warning/0506/math0506306.html | ar5iv | text | # Three amalgams with remarkable normal subgroup structures
## 1. Introduction
Motivated by expected analogies between cocompact lattices in products of automorphism groups of regular trees and cocompact lattices in higher rank semisimple Lie groups, Burger and Mozes discovered in their study of groups acting on products of trees the first examples of finitely presented torsion-free simple groups . These groups are moreover amalgamated products of finitely generated non-abelian free groups, thus answering Neumann’s question on the existence of simple amalgams of free groups. One crucial step in the construction of Burger-Mozes is a deep theorem, which states that certain cocompact lattices in the product of automorphism groups of locally finite trees $`\mathrm{Aut}(T_1)\times \mathrm{Aut}(T_2)`$ cannot have non-trivial normal subgroups of infinite index. Applying this theorem to a cocompact lattice which contains as a subgroup a non-residually finite group constructed by Wise in , we give an example of a finitely presented torsion-free simple group $`\mathrm{\Lambda }_1`$ of the form $`F_9_{F_{81}}F_9`$, where $`F_k`$ denotes the free group of rank $`k`$. See for a list of $`32`$ other finitely presented torsion-free simple groups emerging from the same method. Note that the simple groups of Burger-Mozes are also explicitly given in principle, but not very manageable in practice, because of their extremely long finite presentations. In addition to the simple group $`\mathrm{\Lambda }_1`$, we construct two other groups $`\mathrm{\Lambda }_2`$ and $`\mathrm{\Lambda }_3`$, also having amalgam decompositions $`F_9_{F_{81}}F_9`$. They are not simple, but $`\mathrm{\Lambda }_2`$ is virtually simple and $`\mathrm{\Lambda }_3`$ has no non-trivial finite quotients. An amalgam $`F_3_{F_{13}}F_3`$ without proper subgroups of finite index has already been constructed by Bhattacharjee in , using different techniques. Our search for groups with the desired properties was made possible by several computer programs written in GAP . See \[11, Appendix B\] for the program code used to construct the examples. We refer to , , and for detailed background on automorphism groups of trees, lattices in products of trees, and square complexes.
## 2. Definition of the groups $`\mathrm{\Gamma }_i`$ and $`\mathrm{\Lambda }_i`$
Let always $`i\{1,2,3\}`$. Our groups $`\mathrm{\Lambda }_i`$ will be normal subgroups of index $`4`$ of groups $`\mathrm{\Gamma }_i`$ defined by their finite presentations
$$\mathrm{\Gamma }_i=a_1,a_2,a_3,a_4,a_5,b_1,b_2,b_3,b_4,b_5r_1,\mathrm{},r_{25},$$
where the relators $`r_1,\mathrm{},r_{25}`$ (depending on $`i`$) are given in Table 1. Capital letters in this table indicate inverses, for example $`r_1=a_1b_1A_2B_2=a_1b_1a_2^1b_2^1`$.
Observe that the twelve relators $`r_1,\mathrm{},r_{12}`$ are the same for each group $`\mathrm{\Gamma }_i`$. The reason for this will become clear in the proof of Theorem 1 in Section 3. To describe the geometric nature of $`\mathrm{\Gamma }_i`$, we recall the following general construction which associates to a finite presentation of a group $`G`$ its *standard $`2`$-complex* $`X`$ with fundamental group $`G`$: by definition, the one-skeleton of $`X`$ has a single vertex $`x`$ and an oriented loop for each generator of the given presentation of $`G`$. Furthermore, for each relator $`r`$, a $`2`$-cell with boundary labelled by $`r`$ is glued into this one-skeleton to get $`X`$. Then $`G=\pi _1(X,x)`$. By construction of the $`25`$ relators of $`\mathrm{\Gamma }_i`$, its associated standard $`2`$-complex $`X_i`$ is a finite square complex (all relators have length four, hence all $`2`$-cells are squares) having the additional property that its universal cover $`\stackrel{~}{X_i}`$ is the affine building $`𝒯_{10}\times 𝒯_{10}`$, the product of two $`10`$-regular trees. Equivalently, this property requires that to each pair $`(a,b)A\times B`$, there is a uniquely determined pair $`(\stackrel{~}{a},\stackrel{~}{b})A\times B`$ such that $`ab=\stackrel{~}{b}\stackrel{~}{a}`$ in $`\mathrm{\Gamma }_i`$, where $`A:=\{a_1,\mathrm{},a_5\}^{\pm 1}`$ and $`B:=\{b_1,\mathrm{},b_5\}^{\pm 1}`$. This can be easily verified for our three given examples. In the terminology of , $`X_i`$ is a finite $`1`$-vertex VH-T-square complex, and in the terminology of , $`\mathrm{\Gamma }_i=\pi _1(X_i)`$ is a $`(10,10)`$–group. The group of automorphisms $`\mathrm{Aut}(𝒯_{10})`$, equipped with the usual topology of simple convergence, is a locally compact group. Taking the product topology, $`\mathrm{\Gamma }_i`$ can be seen as a discrete subgroup of $`\mathrm{Aut}(𝒯_{10})\times \mathrm{Aut}(𝒯_{10})`$ with compact quotient, in other words as a cocompact lattice. A crucial role in deducing interesting results on the normal subgroup structure of $`\mathrm{\Gamma }_i`$ play the so-called local groups of $`\mathrm{\Gamma }_i`$. The idea to define them is the following: take the projection of $`\mathrm{\Gamma }_i`$ to one factor of $`\mathrm{Aut}(𝒯_{10})\times \mathrm{Aut}(𝒯_{10})`$ (say the projection $`\mathrm{pr}_1`$ to the first factor) and fix any vertex $`x_h`$ of $`𝒯_{10}`$. Then the elements in the closure $`\overline{\mathrm{pr}_1(\mathrm{\Gamma }_i)}<\mathrm{Aut}(𝒯_{10})`$ stabilizing $`x_h`$, induce a finite permutation group $`P_h^{(1)}(\mathrm{\Gamma }_i)<S_{10}`$ on the $`10`$ neighbouring vertices of $`x_h`$ in $`𝒯_{10}`$ (or more generally, for $`k`$, subgroups $`P_h^{(k)}(\mathrm{\Gamma }_i)`$ of the symmetric group $`S_{109^{k1}}`$, taking the induced action on the $`k`$-sphere in $`𝒯_{10}`$ around $`x_h`$). The same procedure can be done with the second projection $`\mathrm{pr}_2`$ to get local groups $`P_v^{(k)}(\mathrm{\Gamma }_i)<S_{109^{k1}}`$. It is important to note that these local groups (more precisely, their generators in $`S_{109^{k1}}`$) can be directly computed, given the relators $`r_1,\mathrm{},r_{25}`$ of Table 1, see \[7, Chapter 1\] or \[11, Section 1.4\] for details. Here, we get for $`k=1`$ the groups
$`P_h^{(1)}(\mathrm{\Gamma }_1)=`$ $`(7,8)(9,10),(1,2)(3,4),(1,2)(3,4)(7,8)(9,10),`$
$`(1,8,4,5)(2,7,3,10),(1,9,4,8)(3,10,6,7)=A_{10},`$
$`P_h^{(1)}(\mathrm{\Gamma }_2)=`$ $`(7,8)(9,10),(1,2)(3,4),(1,2)(3,4)(7,8)(9,10),`$
$`(1,8,4,9)(2,10,7,3),(1,9,8,6,4)(2,7,5,3,10)=A_{10},`$
$`P_h^{(1)}(\mathrm{\Gamma }_3)=`$ $`(5,6)(7,8)(9,10),(1,2)(3,4),(1,2)(3,4)(7,8)(9,10),`$
$`(1,4,8,9,2,3,7,10)(5,6),(1,9,2,10)(3,5,7)(4,6,8),`$
$`P_v^{(1)}(\mathrm{\Gamma }_1)=`$ $`(1,2)(4,6,7,5)(8,10,9),(1,2,3)(4,5,7,6)(9,10),(1,2)(4,5,7,6)(8,10,9),`$
$`(1,2,3)(4,6,7,5)(9,10),(1,3,10,8)(2,4,6,9,7,5)=A_{10},`$
$`P_v^{(1)}(\mathrm{\Gamma }_2)=`$ $`(1,2)(4,6)(8,10,9),(1,2,3)(5,7)(9,10),(1,2)(4,6,5,7)(8,10,9),`$
$`(1,2,3)(4,6,5,7)(9,10),(1,2,4,3,10,9,7,8)(5,6)=A_{10},`$
$`P_v^{(1)}(\mathrm{\Gamma }_3)=`$ $`(1,2)(4,7,5,6)(8,10,9),(1,2,3)(4,7,5,6)(9,10),(1,2)(4,5,6,7)(8,10,9),`$
$`(1,2,3)(4,5,6,7)(9,10),(1,7)(2,8)(3,9)(4,10)(5,6)=S_{10}.`$
The transitivity of the permutation groups given above will be important in the proof of Theorem 1. Recall that a group $`G<S_{10}`$ is *transitive* if for any pair $`m,n\{1,\mathrm{},10\}`$ there exists a $`gG`$ such that $`g(m)=n`$. Moreover, $`G`$ is called *$`2`$-transitive* if for any $`m_1,m_2,n_1,n_2\{1,\mathrm{},10\}`$ with $`m_1m_2`$ and $`n_1n_2`$ there is an element $`gG`$ such that $`g(m_1)=n_1`$ and $`g(m_2)=n_2`$. Note that the group $`P_h^{(1)}(\mathrm{\Gamma }_3)`$ is a transitive (but not $`2`$-transitive) subgroup of $`S_{10}`$ of order $`3840`$, whereas the alternating group $`A_{10}`$ and the symmetric group $`S_{10}`$ are obviously $`2`$-transitive.
We define now $`\mathrm{\Lambda }_i`$ to be the kernel of the surjective homomorphism
$`\mathrm{\Gamma }_i`$ $`/2\times /2`$
$`a_1,\mathrm{},a_5`$ $`(1+2,0+2)`$
$`b_1,\mathrm{},b_5`$ $`(0+2,1+2),`$
where $`\mathrm{\Gamma }_i`$ is given by its finite presentation described above. Each group $`\mathrm{\Lambda }_i`$ can be decomposed in two ways as amalgamated products $`F_9_{F_{81}}F_9`$, such that $`F_{81}`$ has index $`10`$ in both factors $`F_9`$. More precisely, this means that for any $`i\{1,2,3\}`$ there exist injective homomorphisms $`j_1,j_3:F_{81}F_9s_1,\mathrm{},s_9`$ and $`j_2,j_4:F_{81}F_9t_1,\mathrm{},t_9`$ such that
$$[F_9:j_1(F_{81})]=[F_9:j_2(F_{81})]=[F_9:j_3(F_{81})]=[F_9:j_4(F_{81})]=10$$
and
$`\mathrm{\Lambda }_i`$ $`s_1,\mathrm{},s_9,t_1,\mathrm{},t_9j_1(u_1)=j_2(u_1),\mathrm{},j_1(u_{81})=j_2(u_{81})`$
$`s_1,\mathrm{},s_9,t_1,\mathrm{},t_9j_3(u_1)=j_4(u_1),\mathrm{},j_3(u_{81})=j_4(u_{81}),`$
where $`\{u_1,\mathrm{},u_{81}\}`$ are the free generators of $`F_{81}`$. This is a direct consequence of a result of Wise (see \[14, Theorem I.1.18\]), describing each of the two decompositions of certain square complex groups $`\mathrm{\Gamma }`$ as a fundamental group of a finite graph of finitely generated free groups (in the language of the Bass-Serre theory). If the local groups of $`\mathrm{\Gamma }`$ are “sufficiently transitive” (which always happens in our examples), the two finite graphs corresponding to $`\mathrm{\Lambda }_i`$ in Wise’s construction each consist of two vertices and a single edge. Therefore we get amalgams of finitely generated free groups. It is well-known that amalgams of free groups are always torsion-free, since every element of finite order in an amalgam is conjugate to an element of finite order in one of the two factors (see for example \[9, Theorem IV.2.7\]). Note that following Wise’s proof of \[14, Theorem I.1.18\], it is not difficult (but quite laborious by hand) to give explicit descriptions of the injective homomorphisms $`F_{81}F_9`$ in the amalgam decompositions of $`\mathrm{\Lambda }_i`$.
## 3. Results and Proofs
In the following theorem, we discuss the normal subgroups of $`\mathrm{\Lambda }_i`$.
###### Theorem 1.
Let $`\mathrm{\Lambda }_1`$, $`\mathrm{\Lambda }_2`$, $`\mathrm{\Lambda }_3`$ be the groups defined in Section 2. Then
* $`\mathrm{\Lambda }_1`$ is simple.
* Every non-trivial normal subgroup of $`\mathrm{\Lambda }_2`$ has finite index, but $`\mathrm{\Lambda }_2`$ is not simple.
* $`\mathrm{\Lambda }_3`$ has no proper subgroups of finite index, but is not simple.
###### Proof.
Let $`W`$ be the group with finite presentation
$$a_1,a_2,a_3,a_4,b_1,b_2,b_3r_1,\mathrm{},r_{12},$$
where the relators $`r_1,\mathrm{},r_{12}`$ are again taken from Table 1. Wise showed in \[14, Main Theorem 5.5\], that the non-trivial element $`w:=a_2a_1^1a_3a_4^1W`$ is contained in each finite index subgroup of $`W`$. In particular, $`W`$ is non-residually finite. Moreover, $`W<\mathrm{Aut}(𝒯_8)\times \mathrm{Aut}(𝒯_6)`$ is the fundamental group of a $`1`$-vertex VH-T square complex which embeds into the square complex $`X_i`$ associated to $`\mathrm{\Gamma }_i`$ ($`i=1,2,3`$), inducing an injection on the level of fundamental groups, i.e. $`W<\mathrm{\Gamma }_i=\pi _1(X_i)`$ (the fact that we get an injection can be deduced from the non-positive curvature of the product of trees $`𝒯_{10}\times 𝒯_{10}`$, see \[4, Proposition II.4.14(1)\]). Hence we have
$$1w\underset{N\stackrel{\text{f.i.}}{}W}{}N<\underset{N\stackrel{\text{f.i.}}{}\mathrm{\Gamma }_i}{}N=\underset{N\stackrel{\text{f.i.}}{<}\mathrm{\Gamma }_i}{}N\mathrm{\Gamma }_i,$$
where “f.i.” stands for “finite index”. In particular, $`\mathrm{\Gamma }_i`$ (and hence its finite index subgroup $`\mathrm{\Lambda }_i`$) is non-residually finite. Observe that $`w\mathrm{\Lambda }_i\mathrm{\Gamma }_i`$. One important point in the construction of $`\mathrm{\Gamma }_i`$ is to guarantee that the normal closure of $`w`$ in $`\mathrm{\Gamma }_i`$, denoted by $`w_{\mathrm{\Gamma }_i}`$, has finite index in $`\mathrm{\Lambda }_i`$. (Note that however $`[W:w_W]=\mathrm{}`$.) This already implies that $`w_{\mathrm{\Gamma }_i}`$ has no proper subgroups of finite index. Indeed, assume that $`M<w_{\mathrm{\Gamma }_i}`$ is a subgroup of finite index. Then
$$\underset{N\stackrel{\text{f.i.}}{<}\mathrm{\Gamma }_i}{}N<M\stackrel{\text{f.i.}}{<}w_{\mathrm{\Gamma }_i}\stackrel{\text{f.i.}}{<}\mathrm{\Lambda }_i\stackrel{\text{f.i.}}{<}\mathrm{\Gamma }_i.$$
Using
$$w_{\mathrm{\Gamma }_i}<\underset{N\stackrel{\text{f.i.}}{}\mathrm{\Gamma }_i}{}N=\underset{N\stackrel{\text{f.i.}}{<}\mathrm{\Gamma }_i}{}N,$$
we get
$$M=w_{\mathrm{\Gamma }_i}=\underset{N\stackrel{\text{f.i.}}{<}\mathrm{\Gamma }_i}{}N.$$
We proceed now separately for the three groups $`\mathrm{\Lambda }_1`$, $`\mathrm{\Lambda }_2`$ and $`\mathrm{\Lambda }_3`$.
* We have $`w_{\mathrm{\Gamma }_1}=\mathrm{\Lambda }_1`$. This can be checked by hand, or more easily, using a computer algebra system like GAP , which shows that adding the relator $`w`$ to the presentation of $`\mathrm{\Gamma }_1`$ gives the group $`/2\times /2`$ of order $`4`$. It remains to prove that $`\mathrm{\Lambda }_1`$ has no non-trivial normal subgroups of infinite index. But this follows directly from the normal subgroup theorem of Burger-Mozes \[7, Theorem 4.1, Corollary 5.4\] applied to the “irreducible” cocompact lattice $`\mathrm{\Gamma }_1<\mathrm{Aut}(𝒯_{10})\times \mathrm{Aut}(𝒯_{10})`$ with local groups $`P_h^{(1)}(\mathrm{\Gamma }_1)P_v^{(1)}(\mathrm{\Gamma }_1)=A_{10}`$, and applied to its finite index subgroup $`\mathrm{\Lambda }_1<\mathrm{\Gamma }_1`$.
* For the second group, we compute $`[\mathrm{\Lambda }_2:w_{\mathrm{\Gamma }_2}]=2`$, thus $`\mathrm{\Lambda }_2`$ is not simple. By exactly the same argument as in part (1), every non-trivial normal subgroup of $`\mathrm{\Gamma }_2`$ (and of $`\mathrm{\Lambda }_2`$, respectively) has finite index. Observe that $`w_{\mathrm{\Gamma }_2}`$ is a simple group with amalgam decomposition $`F_{17}_{F_{161}}F_{17}`$. In particular, $`\mathrm{\Gamma }_2`$ and $`\mathrm{\Lambda }_2`$ are virtually simple groups.
* As in part (1), $`w_{\mathrm{\Gamma }_3}=\mathrm{\Lambda }_3`$ proves that $`\mathrm{\Lambda }_3`$ has no proper subgroup of finite index. However, in contrast to what happens in part (1) and (2), the local group $`P_h^{(1)}(\mathrm{\Gamma }_3)`$ is transitive, but not $`2`$-transitive. Therefore, the normal subgroup theorem of Burger-Mozes cannot be applied here. Indeed, $`\mathrm{\Lambda }_3`$ is not simple, since $`1a_5^4_{\mathrm{\Lambda }_3}\mathrm{\Lambda }_3`$. This comes from the fact that $`a_5^4`$ acts trivially on the second factor of $`𝒯_{10}\times 𝒯_{10}`$. In other words, $`a_5^4\mathrm{ker}(\mathrm{pr}_2)\mathrm{\Gamma }_3`$. To see this, let
$$A^{}:=\{(a_1a_2^1)^2,(a_2^1a_1)^2,(a_3a_4^1)^2,(a_4^1a_3)^2,a_5^4\}^{\pm 1}$$
and check that for all $`a^{}A^{}`$ and $`bB=\{b_1,\mathrm{},b_5\}^{\pm 1}`$, we have $`b^1a^{}bA^{}`$. This in fact implies that $`A^{}\mathrm{ker}(\mathrm{pr}_2)`$. Note that no element of $`\mathrm{\Gamma }_3`$ acts trivially on the *first* factor of $`𝒯_{10}\times 𝒯_{10}`$ (by \[6, Proposition 3.1.2, 1)\] and \[6, Proposition 3.3.2\]). As a consequence, $`\mathrm{\Lambda }_3`$ has two decompositions $`F_9_{F_{81}}F_9`$, where one amalgam is effective and the other one is not effective.
We conclude by giving two remarks:
###### Remark 2.
Recall that a group $`G`$ is called *SQ-universal* if every countable group can be embedded in a quotient of $`G`$. It is mentioned in \[1, Chapter 9.15\] that Ilya Rips can prove any amalgamated product $`A_CB`$ to be SQ-universal, provided that $`BC`$ and the number of double cosets $`|C\backslash A/C|`$ is at least $`3`$ (if $`C`$ is seen as usual as a subgroup of $`A`$ and $`B`$ via the two injections $`j_1:CA`$ and $`j_2:CB`$ in the amalgam), but there is no published proof as far as we know. If Rips’ statement is true, we could apply it to exactly one decomposition $`F_9_{F_{81}}F_9`$ of $`\mathrm{\Lambda }_3`$ (to the effective one), where $`|F_{81}\backslash F_9/F_{81}|=3`$. Note however that in the second decomposition of $`\mathrm{\Lambda }_3`$ (where the corresponding local group $`P_v^{(1)}(\mathrm{\Gamma }_3)`$ is $`S_{10}`$) and in both decompositions of $`\mathrm{\Lambda }_1`$ and $`\mathrm{\Lambda }_2`$, we always have $`|F_{81}\backslash F_9/F_{81}|=2`$, since their local actions on $`𝒯_{10}`$ are $`2`$-transitive.
###### Remark 3.
By construction, the three groups $`\mathrm{\Lambda }_1`$, $`\mathrm{\Lambda }_2`$, $`\mathrm{\Lambda }_3`$ are non-residually finite. As a contrast, if one takes a *double* $`F_9_{F_{81}}F_9`$ (i.e. an amalgam where the two injections $`j_1,j_2:F_{81}F_9`$ are identical), such that $`F_{81}`$ has finite index in both factors $`F_9`$ (consequently index $`10=(811)/(91)`$), then one directly gets a surjective homomorphism $`F_9_{F_{81}}F_9F_9`$ (the obvious folding map), and moreover $`F_9_{F_{81}}F_9`$ contains by \[2, Theorem 1.4\] a subgroup of finite index which is a direct product of two non-abelian free groups of finite rank. In particular, such a double $`F_9_{F_{81}}F_9`$ is SQ-universal and residually finite. The residual finiteness also follows from . |
warning/0506/hep-ph0506210.html | ar5iv | text | # Lorentz invariant ensembles of vector backgrounds
## I Introduction
Vector backgrounds appear in numerous sectors of physics. For example they can be used to include the influence of mass dimension two condensates into quantum chromodynamics (QCD) by shifting the gauge field and subsequently restoring Lorentz invariance by averaging over a Lorentz invariant ensemble of backgrounds hp ; ls . There this construction removes quarks and gluons from the spectrum of freely propagating particles. Lately, those condensates have attracted much attention in various respects (see e.g. a2 ).
Analogous vector terms also occur in other domains. Introduced without the subsequent restoration of Lorentz invariance they are used in extensions of the standard model explicitely breaking Lorentz invariance lv . The inclusion of such terms is motivated from string theory string and non-commutative field theory noncom . Lorentz invariance is one of the best-tested postulates of field theory but the current experimental status still leaves room for small deviations. However the approach under investigation in this article is Lorentz invariant.
We also identify the background vector with a shift of the gauge field. First however, in section II, we discuss the general framework for the inclusion of a dependence on an additional vector into a previously Lorentz invariant theory which a priori breaks Lorentz invariance explicitely but where the correlators are defined as the ensemble average over a Lorentz invariant set. Thereby a modified theory is obtained which preserves Lorentz invariance. A Lorentz invariant ensemble is a set of vectors which is mapped onto itself under any Lorentz transformation, while, of course, almost every single element changes. As a next step, in section IIA, we carry out a classification of the weight functions which characterise those ensembles. At variance with euclidean space some subtleties arise in Minkowski space. Commonly the used sets contain configurations leading to vanishing and non-vanishing field tensors hp ; ls . Here, we limit the ensembles further by constraining them to pure gauge configurations. In this framework we analyse the objects central to the modified theory, i.e., the generating functional for the Green functions in section IIB and the fermionic two-point function by solving its equation of motion explicitly ddd in section IIC.
Finally, in section III, we summarise the paper, the main result being that the propagation of fermions over arbitrarily long distances is already stopped in ensembles of pure gauge configurations of the background for euclidean and Minkowski spaces characterised by their respective metrics. This result is not only derived to all orders in the background and otherwise at tree level but for the exact propagator.
Other observations are non-gaussianity of the resulting theory, structural similarities between the present approach and Lorentz invariant generalisations of chemical potentials as well as a technical relationship to stochastic field theory. Last but not least, in Minkowski space the modification of the fermionic two-point Green function amounts to a contribution of a scalar to the fermion’s self energy but without external legs. This again indicates in a diagrammatic way that the ultraviolet degrees of freedom are removed from the asymptotic spectrum.
## II Breaking and restoring Lorentz invariance
Regard a gauge field theory which is modified by including a dependence on a vector $`\mathrm{\Phi }`$. The translational invariance of the system remains intact, because the vector is constant. The Lorentz invariance of the theory is to be restored by taking the average over an ensemble of vectors $`\mathrm{\Phi }`$ characterised by a Lorentz invariant weight $`W(\mathrm{\Phi })`$ <sup>1</sup><sup>1</sup>1Integrations over the $`^4`$ are denoted by a subscript, e.g.: $`d^4\mathrm{\Phi }=:_\mathrm{\Phi }`$.:
$`𝒪_W={\displaystyle _\mathrm{\Phi }}W(\mathrm{\Phi })𝒪,`$ (1)
where $`𝒪`$ stands for a generic operator, here and in the following, and with the normalisation condition:
$`{\displaystyle _\mathrm{\Phi }}W(\mathrm{\Phi })=1.`$ (2)
Apart from the case where $`\mathrm{\Phi }=0`$, which corresponds to the original theory, functions of $`\mathrm{\Phi }^2`$ are the only Lorentz invariant quantities that can be constructed from the vector $`\mathrm{\Phi }`$. The most general Lorentz invariant weight $`W(\mathrm{\Phi })`$ is given by the sum of an arbitrary normalisable function $`w=w(\mathrm{\Phi }^2)`$ and a delta distribution $`\delta ^{(4)}(\mathrm{\Phi })`$:
$`W(\mathrm{\Phi })=c\delta ^{(4)}(\mathrm{\Phi })+w(\mathrm{\Phi }^2)`$ (3)
In hp the vector $`\mathrm{\Phi }`$ represent a vector condensate translating the gauge boson field $`AA+\mathrm{\Phi }`$. That system is investigated with a euclidean metric for quantum chromodynamics (QCD). In the sense $`c=0`$ the weight chosen there ($`\mathrm{exp}\{\mathrm{\Phi }^2/\mathrm{\Lambda }^2\}`$) does not contain the unmodified theory. This manifests itself in the one-particle pole being removed from the quark and gluon propagators determined to all orders in $`\mathrm{\Phi }`$, meaning that there the partons do not propagate over arbitrarily long distances. They are no longer part of the asymptotic spectrum. The lowest order in an expansion for momenta large compared to the scale $`\mathrm{\Lambda }^2`$ reproduces the standard free propagators.
In QCD the vector $`\mathrm{\Phi }`$ also carries colour indices: $`\mathrm{\Phi }^2=\mathrm{\Phi }_\mu ^a\mathrm{\Phi }^{a\mu }`$. The vector $`\mathrm{\Phi }`$ is to transform homogeneously under gauge transformations whence any function of $`\mathrm{\Phi }^2`$ is gauge covariant. This also establishes the connection of the present approach with mass dimension two condensates because now $`\mathrm{\Phi }`$ acts as a contribution to the gauge field hp ; ls ; a2 . Due to the non-abelian nature of the gauge theory, the ensemble of constant vectors characterised by a function of its square contains members which are pure gauge configurations and such leading to a non-vanishing field tensor. In our investigation, we will distinguish these cases by limiting the ensemble to vanishing field tensors. This in itself is a gauge invariant criterion. For the fermionic sector this leads to an analogue of quantum electrodynamics (QED) which we will study in the following. Further, we will compare the results for the different metrics.
### II.1 Weight classification
Let us begin with a classification of the weight functions. In principle, in euclidean space the case $`\mathrm{\Phi }=0`$ is already included in $`w(\mathrm{\Phi }^2)`$ as there $`\mathrm{\Phi }^2=0`$ implies $`\mathrm{\Phi }=0`$. Nevertheless, in order to mark the potential contribution from the unmodified theory clearly, i.e., from $`\mathrm{\Phi }=0`$, let us split it off in form of a delta distribution in accordance with Eq. (3). The normalisation condition (2) then implies
$`\pi ^2{\displaystyle _0^+\mathrm{}}v𝑑vw_\mathrm{E}(v)=1c`$ (4)
with $`v:=\mathrm{\Phi }^2`$ and where the subscript <sub>E</sub> marks the Euclidean case.
Every possible Lorentz invariant weight function $`w_\mathrm{E}(\mathrm{\Phi }^2)`$ can be reconstructed by a convolution with a delta weight
$`w_\mathrm{E}(\mathrm{\Phi }^2)={\displaystyle 𝑑\lambda \delta (\mathrm{\Phi }^2\lambda )w_\mathrm{E}(\lambda )}.`$ (5)
In this sense the delta weight:
$`w_\mathrm{E}^{\{\lambda \}}(\mathrm{\Phi }):=(4\pi \lambda )^1\delta (\mathrm{\Phi }^2\lambda )`$ (6)
can be seen as fundamental.
However, if, in the presence of a space with Minkowski metric, one wants to work in a time ordered formalism also in the theory with the background a different choice for the basis is better adapted. Noticing that:
$`2\pi \mathrm{i}\delta (\mathrm{\Phi }^2\lambda )=S_\lambda ^{}(\mathrm{\Phi })S_\lambda ^+(\mathrm{\Phi }),`$ (7)
where
$`S_\lambda ^\pm (\mathrm{\Phi })=(\mathrm{\Phi }^2\lambda \pm iϵ)^1`$ (8)
with $`+`$ ($``$) is the time ordered (anti time-orderded) propagator of a scalar with the squared mass equal to $`\lambda `$. In the framework of a time-ordered formalism $`S_\lambda ^+(\mathrm{\Phi })`$ could be seen as the elementary weight.
However, with a Minkowski metric—apart from the fact that the case $`\mathrm{\Phi }=0`$ is not included in the function $`w_\mathrm{M}(\mathrm{\Phi }^2)`$ and has to be added separately—the hyperboloid pair characterised by $`\mathrm{\Phi }^2=\mathrm{const}.`$ has infinite content. Thus with one single elementary weight the normalisation condition (2) cannot be satisfied. Further, even the difference in content between two hyperboloid pairs is in general infinite whereby a superposition of two weights does not suffice to satisfy the normalisation condition in a non-trivial way. For these reasons the minimal construction has to be:
$`w_\mathrm{M}(\mathrm{\Phi }^2)={\displaystyle \underset{j=1}{\overset{3}{}}}a_jS_{\lambda _j}^+(\mathrm{\Phi }),`$ (9)
with
$`{\displaystyle \underset{j=1}{\overset{3}{}}}a_j=0,`$ (10)
and
$`{\displaystyle \underset{j=1}{\overset{3}{}}}a_j\lambda _j=0.`$ (11)
Then the normalisation condition (2) becomes <sup>2</sup><sup>2</sup>2If the first two conditions are taken into account, the normalisation condition can be expressed in terms of logarithms of ratios of $`\lambda _j`$.:
$`{\displaystyle \frac{4\pi ^2}{4}}{\displaystyle \underset{j=0}{\overset{3}{}}}a_j\lambda _j\mathrm{ln}\lambda _j=1c.`$ (12)
The conditions (10) to (12) can be derived by putting a Fourier phase into the normalisation integral (2) and letting the variable conjugate to $`\mathrm{\Phi }`$ go to zero afterwards. The conditions follow from requiring that the limit exist. Then, in general, it will also be non-zero \[see Eq. (12)\].
As a consequence of these conditions, $`w_\mathrm{M}`$, as opposed to $`w_\mathrm{E}`$, cannot be positive definite. Condition (10) resembles the one used in Pauli-Villars regularisation.
Any discrete or continuous superposition of delta weights or (time ordered) scalar propagators (8), respectively, fulfilling the normalisation condition (2) is an allowed weight function, but in what follows we will concentrate on the minimal forms given in the previous equations.
### II.2 Generating functional
The generating functional for the time ordered Green function of QED is given by:
$`Z=Z_{\mathrm{i}nt}Z_AZ_\psi ,`$ (13)
with the interaction:
$`Z_{\mathrm{i}nt}=\mathrm{exp}\left\{i{\displaystyle _x}\delta _\eta \mathit{\delta ̸}_J\delta _\eta \right\},`$ (14)
the bosonic:
$`Z_A=\mathrm{exp}\left\{{\displaystyle \frac{i}{2}}{\displaystyle _{x,y}}J(x)\mathrm{\Gamma }_0(xy)J(y)\right\},`$ (15)
and the fermionic part:
$`Z_\psi =\mathrm{exp}\left\{i{\displaystyle _{x,y}}\overline{\eta }(x)G_0(xy)\eta (y)\right\},`$ (16)
with functional derivatives $`\delta `$ with respect to the currents $`J,\eta ,\overline{\eta },`$ and the free time-ordered propagators for the bosons $`\mathrm{\Gamma }_0`$ and the fermions $`G_0`$, respectively.
The modified theory’s generating functional $`𝒵`$ is obtained by shifting the gauge field $`A`$ by $`\mathrm{\Phi }`$ and subsequent averaging with the weight $`W`$. Here this amounts to a modification of the fermionic part leading to:
$`𝒵=Z_{\mathrm{i}nt}Z_A𝒵_\psi ,`$ (17)
where
$`𝒵_\psi =\mathrm{exp}\left\{i{\displaystyle _{x,y}}\overline{\eta }(x)G_\mathrm{\Phi }(xy)\eta (y)\right\}_W.`$ (18)
The other two factors of the generating functional $`Z_{\mathrm{i}nt}`$ and $`Z_A`$ can always be taken inside the averaging integral. $`G_\mathrm{\Phi }`$ is the time ordered fermion propagator in the field $`\mathrm{\Phi }`$. Under the usually made assumption hp ; ls that all other condensates are absent it obeys the equation of motion:
$`[i\partial ̸(x)+\mathrm{\Phi ̸}m]G_\mathrm{\Phi }(xy)=\delta ^{(4)}(xy)`$ (19)
which is solved by:
$`G_\mathrm{\Phi }(z)=e^{i\mathrm{\Phi }z}G_0(z).`$ (20)
Remember that the fermionic propagator in the presence of a medium resulting in a chemical potential $`\mu `$ reads $`e^{i\mu z_0}G_0(z)`$, i.e., technically the chemical potential corresponds to the temporal component of a vector and physically to a conserved charge. Carrying out the $`\mathrm{\Phi }`$ integral corresponds to a Fourier transformation of the weight function:
$`𝒵_\psi `$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(i)^n}{n!}}{\displaystyle _{\{x_m\},\{y_m\}}}\stackrel{~}{W}(z_n)\times `$ (21)
$`\times {\displaystyle \underset{m=0}{\overset{n}{}}}[\overline{\eta }(x_m)G_0(x_my_m)\eta (y_m)],`$
where $`\stackrel{~}{W}(z_n)`$ is the Fourier transformation of the weight function evaluated at $`z_n:=_{m=0}^n(x_my_m)`$. Thence and due to the Fourier integral the prefactor can be seen as a Lorentz invariant superposition of chemical potential-like factors.
Eq. (20) has also similarities with the expressions occuring in the context of twisted boundary conditions on compact spaces which are used in lattice calculations b . There only the spatial components of the vector are non-zero.
The special form of the generating functional leads to the following relation for the (higher) correlators:
$`0|T{\displaystyle \underset{m=1}{\overset{n}{}}}\psi (x_m)\overline{\psi }(y_m)|0_W=`$
$`=\stackrel{~}{W}(z_n)0|T{\displaystyle \underset{m=1}{\overset{n}{}}}\psi (x_m)\overline{\psi }(y_m)|0,`$ (22)
which remains essentially the same if bosonic operators are added. Eq. (22) evidences why $`\stackrel{~}{W}`$ has to be time ordered, if the functional is to generate time ordered Green functions.
Even if the second factor on the right-hand side of the previous equation should show gaussianity—on a given level—, i.e., factorise into two point correlators, the first factor is a genuine $`2n`$-point function. Therefore the new theory is not gaussian.
Through the limitation to pure gauge configurations of the background, i.e. such with vanishing field tensor, the modified theory can be interpreted as one with a non-trivial vacuum structure without background energy density. One could write:
$`0|T{\displaystyle \underset{m=1}{\overset{n}{}}}\psi (x_m)\overline{\psi }(y_m)|0_W=:`$
$`=:\mathrm{\Omega }|T{\displaystyle \underset{m=1}{\overset{n}{}}}\psi (x_m)\overline{\psi }(y_m)|\mathrm{\Omega },`$ (23)
where $`|\mathrm{\Omega }`$ stands for the new vacuum. The vacuum expectation values of the new theory $`\mathrm{\Omega }|𝒪|\mathrm{\Omega }`$ are the averaged vacuum expectation values $`0|𝒪|0_W`$ of the old theory.
For two-point functions Eq. (22) together with the elementary weight (9) makes the modification of the propagator look like a contribution of a scalar to the self energy of the fermion without external fermion legs. The superposition of multiple scalars with different ”mass squares” $`\lambda `$ leads to the summation of the related self-energy bubbles. If, for example, a three-point function was constructed by including a gauge boson in the previous correlator, the modification would look like a vertex correction but still without the outer fermion legs. For correlators with more than two external fermions the correspondence to standard scalar loops does no longer persist, whence the new theory is not identical to one, where the terms for a scalar degree of freedom are added to the lagrangian density.
### II.3 The time ordered fermion propagator
Now we study the Green function central to the modified theory, i.e., the fermionic two-point function to all orders in $`\mathrm{\Phi }`$.
#### II.3.1 Euclidean metric
With a euclidean metric the details of the $`ϵ`$-prescription are not important. Therefore the elementary weight of choice is the normalised delta weight $`w_E^{\{\lambda \}}(\mathrm{\Phi }^2)`$ and $`c=0`$. As mentioned before $`\mathrm{\Phi }^2=0`$ in $`w_\mathrm{E}`$ also means $`\mathrm{\Phi }=0`$. However, assuming that $`w_\mathrm{E}`$ is not divergent at this point, this contribution is negligible. Denoting this special averaging procedure by $`𝒪_\lambda ^\mathrm{E}`$ we get:
$`G_\mathrm{\Phi }(z)_\lambda ^\mathrm{E}={\displaystyle \frac{\mathrm{sin}\sqrt{\lambda z^2}}{\sqrt{\lambda z^2}}}G_0(z).`$ (24)
This shows that, apart from an oscillatory behaviour, the propagator is suppressed over large distances $`\sqrt{z^2}`$. In the limit of short distances $`\sqrt{z^2}`$ the free propagator is recovered.
Interestingly, Eq. (24) holds not only for the free propagator $`G_0(z)`$ but for the full fermionic propagator, i.e., to all orders in perturbation theory. That is so, because Eq. (20) is also satisfied by the full propagator in the presence of the background $`\mathrm{\Phi }`$. Therefore, the result that the standard propagator is recovered at small distances and that the background causes a suppression at long distances.
Back to tree-level, in momentum space we get:
$`G_\mathrm{\Phi }(k)_\lambda ^\mathrm{E}={\displaystyle \frac{\mathit{}+m}{4\sqrt{k^2\lambda }}}\mathrm{ln}\left|{\displaystyle \frac{(\sqrt{k^2}+\sqrt{\lambda })^2m^2}{(\sqrt{k^2}\sqrt{\lambda })^2m^2}}\right|+`$
$`+{\displaystyle \frac{\mathit{}}{4k^2}}\left[2{\displaystyle \frac{k^2+\lambda m^2}{2\sqrt{k^2\lambda }}}\mathrm{ln}\left|{\displaystyle \frac{(\sqrt{k^2}+\sqrt{\lambda })^2m^2}{(\sqrt{k^2}\sqrt{\lambda })^2m^2}}\right|\right]`$ (25)
One can see that the on-shell pole has been removed from the propagator. It has been replaced by one proportional to $`1/\sqrt{k^2}`$. Consequently the elementary fermions have been removed from the spectrum of asymptotic states.
This result is similar to the one in hp ; ls , but where also background configurations with non-vanishing field tensors were admitted in addition to the pure gauge configurations used here exclusively. In coordinate space the weight chosen in hp constrained to pure gauge backgrounds yields :
$`G_\mathrm{\Phi }(z)_{\mathrm{H}P}^\mathrm{E}=\mathrm{exp}(z^2\mathrm{\Lambda }^2/4)G_0(z)`$ (26)
Here as well the free propagator is recovered at small $`z^2`$ and damped at large $`z^2`$. Even taking the average with this special weight for $`2n`$-point fermion correlators does not lead to a factorisation into two-point correlators (gaussianity).
#### II.3.2 Minkowski metric
In Minkowski space, if one wants to stick to a time ordered treatment the adapted weight function has to be chosen for the additional contribution. In coordinate space, taking the weight given by Eq. (3) with Eq. (9) and $`c=0`$ leads to:
$`G_\mathrm{\Phi }(z)_\mathrm{M}^{\{\lambda _j\}}=4\pi ^2{\displaystyle \underset{j=1}{\overset{3}{}}}a_j\sqrt{\lambda _j}{\displaystyle \frac{\mathrm{K}_1(\sqrt{\lambda _j}\sqrt{z^2+iϵ})}{\sqrt{z^2+iϵ}}}G_0(z)`$
(27)
At $`z^2=0`$ the prefactor of the free propagator $`G_0(z)`$ in the previous equation goes to 1 due to Eqs. (10) to (12) with $`c=0`$. Thence, the free propagator is recovered in the limit of small $`z^2`$. If $`c`$ is chosen different from zero, once the free propagator is still reproduced taken together with the explicitely free contribution. The reproduction of the free propgator at $`z^2=0`$ is an intrinsic consequence of the need to normalise.
Like in euclidean space the previous relation also holds for the full propagator, for the same reason as there.
If $`\lambda _j>0j\{1;2;3\}`$, for $`z^2+\mathrm{}`$ the envelope of this function decays proportionally to $`(z^2)^{(3/4)}`$; for $`z^2\mathrm{}`$ proportionally to $`(z^2)^{(3/4)}\mathrm{exp}[\sqrt{\mathrm{min}(\{|\lambda _j|\})z^2}]`$. If $`\lambda _j<0j\{1,2,3\}`$ the two cases are exchanged. Thus, for large absolute values of $`z^2`$, $`G_\mathrm{\Phi }(z)_\mathrm{M}^{\{\lambda _j\}}`$ is suppressed relative to a free propagator, which shows that the fermions cannot propagate over arbitrarily large distances.
In momentum space the form of $`G_\mathrm{\Phi }(k)_\mathrm{M}^{\{\lambda _j\}}`$ can be determined best by making use of its correspondence to the one-loop contribution of a scalar to the self-energy of the fermion. One obtains:
$`G_\mathrm{\Phi }(k)_\mathrm{M}^{\{\lambda _j\}}=`$
$`=i\pi ^2{\displaystyle \underset{j=0}{\overset{3}{}}}a_j{\displaystyle _0^1}𝑑x(x\mathit{}+m)\mathrm{ln}|(xx_j^+)(xx_j^{})|`$ (28)
with
$`x_j^\pm ={\displaystyle \frac{\lambda _j+k^2m^2}{2k^2}}\pm \sqrt{\left({\displaystyle \frac{\lambda _j+k^2m^2}{2k^2}}\right)^2+{\displaystyle \frac{\lambda _j}{k^2}}},`$ (29)
$`j\{1;2;3\}`$. For $`\lambda _j>0j\{1;2;3\}`$ the $`x`$-integration can be carried out yielding:
$`{\displaystyle _0^1}𝑑x\mathrm{ln}|xx_j^\pm |=`$
$`=(1x_j^\pm )\mathrm{ln}|1x_j^\pm |+x_j^\pm \mathrm{ln}|x_j^\pm |1`$ (30)
and
$`{\displaystyle _0^1}x𝑑x\mathrm{ln}|xx_j^\pm |=`$
$`={\displaystyle \frac{1(x_j^\pm )^2}{2}}\mathrm{ln}|1x_j^\pm |+{\displaystyle \frac{(x_j^\pm )^2}{2}}\mathrm{ln}|x_j^\pm |{\displaystyle \frac{x_j^\pm }{2}}{\displaystyle \frac{1}{4}}.`$ (31)
$`G_\mathrm{\Phi }(k)_\mathrm{M}^{\{\lambda _j\}}`$ is free of poles. For small $`k^2`$ and small mass $`m^2`$, the propagator becomes:
$`G_\mathrm{\Phi }(k)_\mathrm{M}^{\{\lambda _j\}}i\pi ^2(\mathit{}/2+m){\displaystyle \underset{j=1}{\overset{3}{}}}a_j\mathrm{ln}|\lambda _j|.`$ (32)
For large $`k^2`$, $`G_\mathrm{\Phi }(k)_\mathrm{M}^{\{\lambda _j\}}`$ becomes proportional to $`k^2`$ reproducing the behaviour of the free propagator.
Independent of the details, for $`c=0`$ no freely propagating particles are described by the propagator in an ensemble of pure gauge backgrounds. This also explains its aformentioned correspondence to self-energy contributions from scalars without external fermion legs.
## III Summary
We have studied gauge field theories with restored Lorentz invariance. Starting out with a manifestly Lorentz invariant field theory, this symmetry is broken through the inclusion of a (non-trivial) dependence on a four-vector $`\mathrm{\Phi }`$. In the explicitely investigated examples said vector plays the rôle of a contribution to the gauge field. The symmetry is restored by defining correlators as average over a Lorentz invariant ensemble of vectors. Apart from the original contribution with $`\mathrm{\Phi }=0`$ the additional term is characterised by a weight function of the only Lorentz invariant $`\mathrm{\Phi }^2`$. Therefore these theories are connected to mass dimension two vector condensates. The modifications can also be interpreted as a means to include the effect of a non-trivial Lorentz invariant vacuum structure into the original theory. Some of the structures also appear in stochastic field theories. The resulting theories do not show gaussianity.
The presence of the background can bar the asymptotically free propagation of the matter fields used to write down the lagrangian density, in euclidean and Minkowski space. The bare propagator without the background is still reproduced at short distances and high momenta. The dressed propagator has no on-shell pole and is suppressed at large distances relative to the undressed one. In non-abelian gauge theories also a constant gauge field can contribute to the field tensor. Commonly, within this setting, ensembles of backgrounds are used which contain pure gauge configurations and field configurations leading to a non-zero field tensor. Therefore it was not clear a priori to which contributions the observed effect is connected. We find that when constraining the ensemble to pure gauge configurations, the effect is still present. That the pure gauge configurations do lead to non-trivial phenomena can be understood from the mentioned analogy to a chemical potential. Remarkably we have been able to show the suppression of the long-range propagation not only based on the fermionic propagator to all orders in the background and otherwise at tree level but for the full fermionic propagator, i.e., to all loops.
For correlators involving two fermion fields in general and thus especially for the fermion propagator the modifications due to the background resemble scalar loops bridging the fermion lines. The corresponding diagrams do not carry external fermion legs although they are direct contributions to the propagator indicating in this way that they do not involve freely propagating particles.
Some extensions of the standard model violating Lorentz invariance are based on the concept of non-commutative field theories. As here Lorentz invariance has been restored, it would be interesting to study the relationship between the present approach and non-commutative field theories not breaking Lorentz invariance nonbreak .
## Acknowledgments
The authors are grateful for inspiring and informative discussions with Carsten Greiner, Andrew D. Jackson, Kerstin Paech, Olivier P$`\stackrel{`}{\mathrm{e}}`$ne, Dirk Rischke, Francesco Sannino, and Kim Splittorff as well as for the hospitality of the Institute for Theoretical Physics of the Johann Wolfgang Goethe-University in Frankfurt/Main where part of the work for this project has been carried out. Thanks are again due to Kim Splittorff for the careful reading of and useful comments on the manuscript. |
warning/0506/nucl-ex0506023.html | ar5iv | text | # Tomographic Studies of the sQGP at RHIC
## 1 Introduction
Reactions between Au ions at the Relativistic Heavy Ion Collider (RHIC), indicate the creation of a fireball of nuclear matter having energy density well above that required for a de-confined phase of quarks and gluons (QGP) . The decay of this matter results in large azimuthal anisotropies in the particle emission patterns, suggesting early thermalization and the development of substantial pressure gradients which drive the dynamical evolution of the system . Hydrodynamic evolution of the fireball is further corroborated by the observation of strong radial flow and first hints of a long-range emitting source from a recent Imaging analysis . Strong indications for hydrodynamic evolution of the emitting system implies the production of strongly interacting high energy density matter in energetic RHIC collisions . Indeed, this matter has been observed to strongly suppress the yield of hadrons with large transverse momenta and to suppress the away-side jet in central Au+Au collisions . It is believed that this suppression results from energy loss of hard-scattered partons traversing the high energy density matter prior to the formation of hadrons .
An important open question of great current interest is the influence of the parton-medium interactions on jet properties. Such an influence is of paramount importance if one wants to use jets as a probe of the properties of strongly interacting high energy density matter. Several recent works have outlined a possible influence of the coupling between jets and a strongly interacting medium . A particularly important proposal is the conjecture that the energy deposited in the medium could lead to the creation of a shock wave around the propagating parton, thereby creating “conical flow” or “bow waves” analogous to a sonic boom in a fluid. The experimental observation of such conical flow could serve to pin down the sound speed in the nuclear matter created at RHIC.
In order to probe the influence of possible parton-medium interactions on jet properties, we use azimuthal angular correlation functions to investigate jet topologies and yields in d+Au and and Au+Au collisions. Here, the operational strategy is that d+Au measurements provide a good baseline for comparison to the Au+Au measurements which are expected to show much stronger modifications to jet properties.
## 2 Data Analysis
The analysis presented in this paper uses Au+Au and d+Au data ($`\sqrt{s_{NN}}`$=200 GeV) provided by RHIC in the second and third running periods (2001,2003), respectively. The full PHENIX detector setup is described elsewhere . Charged tracks relevant to this analysis were reconstructed in the central arms of PHENIX, each of which covers 90 degrees in azimuth. Tracking was performed via the drift chamber and two layers of multi-wire proportional chambers with pad readout (PC1,PC3). A combinatorial Hough transform in the track bend plane was used for pattern recognition. Most conversions, albedo and decays were rejected by requiring a confirmation hit within a 2 $`\sigma `$ matching window in the PC3. Collision centrality was determined via cuts in the space of BBC versus ZDC analog response .
The correlation function in relative azimuthal angle between particle pairs, $`\mathrm{\Delta }\varphi =(\varphi _1\varphi _2)`$, is defined as the ratio of two distributions
$$C\left(\mathrm{\Delta }\varphi \right)\frac{N_{cor}\left(\mathrm{\Delta }\varphi \right)}{N_{mix}\left(\mathrm{\Delta }\varphi \right)}.$$
(1)
The foreground distribution $`N_{cor}`$, measures coincident particle pairs from the same event by pairing particles from a high-$`p_T`$ “trigger” bin ($`2.5<p_T<4.0`$ GeV/c, hereafter labeled A) with associated particles from a lower $`p_T`$ selection ($`1.0<p_T<2.5`$ GeV/c, hereafter labeled B). The background distribution $`N_{mix}`$, is generated in an analogous way by mixing particle pairs from different events within the same multiplicity and vertex class. The Azimuthal acceptance and detector efficiency effects cancel in the ratio of foreground to background distributions and the correlation function yields the probability distribution for detecting correlated particle pairs per event within the PHENIX pseudorapidity acceptance ($`|\eta |<0.35`$).
Fig. 1 shows d+Au correlation functions for two different trigger and associated $`p_T`$ selections as indicated. The centrality selection is 0-80%. The dashed line represents a double Gaussian fit to the data. These d+Au correlation functions exhibit a shape reminiscent of what one would expect from di-jet fragmentation. That is, a relatively narrow near-side peak centered at $`\mathrm{\Delta }\varphi =0`$ and a somewhat wider away-side peak centered at $`\mathrm{\Delta }\varphi =\pi `$. Fig. 1 shows that both near-side and away-side peaks narrow and are more pronounced for higher $`p_T`$-selections of trigger and associated particles, respectively. Such a pattern is expected if (di-)jet fragmentation is the dominant particle production mechanism at high transverse momenta. It is noteworthy that the widths (of near- and away-side jets) and the yields obtained from d+Au correlation functions are rather similar to those obtained from p+p collisions. Consequently, we ascribe all correlation in d+Au collisions to jets.
The d+Au correlation functions (cf. Fig. 1) are to be compared to the correlation functions for Au+Au data shown in Fig. 2. The filled squares in the figure show correlation functions for AB charged hadron pairs for several indicated centralities. It is difficult to overlook the striking similarity between the correlation function obtained for the most peripheral collisions (cf. Fig. 2f) and that obtained for d+Au collisions (cf. Fig. 1). Both correlation functions show the two narrow peaks (located at $`\mathrm{\Delta }\varphi =0`$ and $`\mathrm{\Delta }\varphi =\pi `$) characteristic of (di-)jet fragmentation. By contrast, the correlation functions for more central Au+Au collisions show strong indications for a harmonic component and hints for a broad away-side jet ie. the correlation functions show a minimum below $`\mathrm{\Delta }\varphi =\pi /2`$. Such a shift away from the minimum expected for harmonic contributions, can only come about if the away-side jet is significantly broader than that observed in d+Au collisions. An important finding that should be stressed here is that the observed characteristics of all of the Au+Au correlation functions can be fully accounted for via two contributions to the correlation function: (i) a (di-)jet and (ii) a harmonic contribution .
## 3 Decomposition of jet and harmonic contributions
Careful investigation of possible modifications to the di-jet distributions in Au+Au collisions require access to the jet-pair distribution without the blurring effects of the harmonic contributions. Consequently, a reliable procedure for decomposing the measured correlation functions into their (di-)jet and harmonic (or flow) contributions are required. Detailed descriptions of such a procedure are given in Refs. . We give here only an outline of the main points.
It can be shown that the pair correlations from the combination of flow and jet sources is given by
$$C^{AB}(\mathrm{\Delta }\varphi )=a_o[C_H^{AB}(\mathrm{\Delta }\varphi )]+J(\mathrm{\Delta }\varphi ),$$
(2)
where $`C_H^{AB}(\mathrm{\Delta }\varphi )`$ is a harmonic function of effective amplitude v<sub>2</sub>,
$$C_H^{AB}(\mathrm{\Delta }\varphi )=[1+2v_2cos2(\mathrm{\Delta }\varphi )];v_2=(v_2^A\times v_2^B).$$
(3)
and J($`\mathrm{\Delta }\varphi )`$ is the (di-)jet function. No explicit or implicit assumption is required for the functional form of J($`\mathrm{\Delta }\varphi `$). Rearrangement of Eq. 2 gives
$$J(\mathrm{\Delta }\varphi )=C^{AB}(\mathrm{\Delta }\varphi )a_oC_H^{AB}(\mathrm{\Delta }\varphi ).$$
(4)
Thus, one only requires knowledge about a<sub>o</sub> and v<sub>2</sub> to evaluate J($`\mathrm{\Delta }\varphi `$). To constrain, a<sub>o</sub> we assume that the (di-)jet function has zero yield at the minimum (ZYAM) $`\mathrm{\Delta }\varphi _{min}`$, in the jet function, i.e. $`a_0C_H^{AB}(\mathrm{\Delta }\varphi _{min})=C^{AB}(\mathrm{\Delta }\varphi _{min})`$. This fixes the value of $`a_o`$.
The $`v_2`$ value reflects the average anisotropy of the particles from both sources, and can be obtained from the single particle distributions relative to the reaction plane $`\psi _R`$. However, this step requires that the reaction plane is itself determined by a procedure essentially free of non-flow effects. This is accomplished in the present analysis by demanding a large (pseudo)rapidity gap ($`\mathrm{\Delta }\eta 33.9`$) between the reaction plane and the particles correlated with it . It is expected that the $`v_2`$ values so obtained are much less affected by jet contributions . In addition, the reaction plane (at each centrality) and its dispersion correction, and $`v_2^A`$ and $`v_2^B`$ were obtained from the same data set used for jet function extraction in order to avoid potential biases. Correction for reaction plane dispersion followed the procedures outlined in Ref. .
### 3.1 Simulation tests
Prior to applying the decomposition method to PHENIX data, its reliability was thoroughly tested via extensive Monte Carlo simulations . These investigations included simulation tests which took account of the $`\varphi `$ and the $`\eta `$ acceptance of PHENIX.
Representative results from these Monte Carlo investigations are summarized in Fig. 3. Panel (a) depicts the azimuthal distribution of the simulated reaction products with respect to the event plane. The smooth line is a harmonic fit to the data to extract $`v_2`$. This $`v_2`$ determines the amplitude of the harmonic component (dashed line in panel (b)) to be subtracted from the correlation function (filled circles) shown in panel (b). The solid squares in panel (b) shows the ZYAM subtracted jet-pair distribution referenced to $`a_o`$. This distribution is to be compared to the input jet pair distribution (solid line) obtained via tagging of the jet-particles in the simulation. The rather good agreement shown between input and output jet-pair distributions in Fig. 3b serves to confirm the reliability of the method. The method is easily generalized to the case in which trigger particle detection is constrained within a cut angle parallel or perpendicular to the reaction plane . Results from simulations in which such constraints have been applied are summarized in panels (c) and (d) of Fig.3. In this case, the harmonic function is determined following the techniques outlined in . Here again, panels (c) and (d) clearly indicate that the input jet function is reproduced in detail. It is noteworthy that a wide range of tests for a variety of input jet-pair distributions including those that might be expected from conical flow, were made with equally good recovery of the input jet-pair distributions .
### 3.2 Decomposition of the measured correlation functions
The solid bands in each panel of Fig. 2 illustrate the application of the ZYAM condition to PHENIX data with the measured values of $`v_2`$ ($`v_2=(v_2^A\times v_2^B`$)). The dashed lines show the $`a_o`$ value obtained for each centrality. Following Eq. 4, the jet-pair distribution is obtained at each centrality via subtraction of the harmonic contribution from the correlation function. It is straightforward to show that the integral of this distribution is related to the average fraction of jet-correlated particle pairs per event and hence the conditional per trigger yield . The ratio of the sum of $`J(\mathrm{\Delta }\varphi )`$ and the sum of $`C(\mathrm{\Delta }\varphi )`$ (over all bins in $`\mathrm{\Delta }\varphi `$) gives the fraction of jet-correlated particle pairs per event $`PF`$,
$$PF=\frac{_iJ(\mathrm{\Delta }\varphi _i)}{_iC(\mathrm{\Delta }\varphi _i)}$$
(5)
Subsequent multiplication of this fraction by the average number of detected particle pairs per event $`N_d^{AB}`$, followed by a division by the product of the detected singles rates $`N_d^A`$, $`N_d^B`$, gives the event averaged jet-pair production in excess of the combinatoric pair production. A final product with the efficiency corrected singles rate $`N_{eff}^B`$, for bin B, gives the efficiency corrected pairs per trigger or conditional yield $`CY`$ ,
$$CY=PF\times \frac{N_d^{AB}}{N_d^A\times N_d^B}\times N_{eff}^B.$$
(6)
## 4 Results
The ZYAM-subtracted conditional yield distributions, normalized to the number of jet-pairs per trigger particle, are shown in Fig. 4. The distribution for the 60-90% peripheral event sample (cf. Fig. 4f) shows the typical (di-)jet shape that is familiar from p+p and d+Au collisions at RHIC. It consists of two distinct peaks, one narrow near-side peak centered at $`\mathrm{\Delta }\varphi =0`$ and a broader away-side peak at $`\mathrm{\Delta }\varphi =\pi `$. It is interesting to trace the evolution of these peaks with collision centrality, as one expects to make increasing amounts of hot and dense matter in the more central collisions. Scanning the mid-central and central jet-pair distributions (Fig. 4a-e), one observes that the near-side peak topology remains essentially unmodified. We attribute this to a possible trigger bias of the near-side jet. However, it should be noted that the near-side yield does show a mild rise with centrality as discussed below. In contrast to the near-side peak, the away-side peaks show indication for strong modifications (both in magnitude and in shape) for all centralities other than the most peripheral event selection. More precisely, in the 0-5%, 5-10% and 40-60% samples, the away-side peak evidences a plateau like shape which is decidedly non-gaussian and much broader in width than that for the 60-90% sample. For the 10-20% and 20-40% samples (Fig. 4c-d), the away-side peak remains broad but also indicates an apparent local minimum at $`\mathrm{\Delta }\varphi =\pi `$ and a maximum at $`\mathrm{\Delta }\varphi =2\pi /3`$. This latter pattern is similar to recent predictions of jet-induced ”conical flow” . It should be pointed out however that these results do not preclude an alternative scenario which conjectures the combined influence of energy loss and the inclination angle of the jet with the flow field . Nonetheless, both approaches require relatively strong coupling between jets and the high energy density medium.
The solid (dashed) lines in Fig.4 represent the conditional yield distributions that would result from subtracting out a $`v_2`$ product lowered (raised) by one systematic error, respectively. The systematic error on $`v_2`$ is dominated by the uncertainty on the reaction plane dispersion. The dotted line indicates the jet-pair distribution resulting from a subtraction with $`v_2`$ product lowered by twice the systematic uncertainty. From these curves, one can see that a decrease in $`v_2`$ by two intervals of the systematic error can recover the local minimum at $`\mathrm{\Delta }\varphi =\pi `$. However, the result of a broadened away-side jet remains robust and the away-side jet-shapes remain non-Gaussian. Several studies are currently underway to firm up the mechanism/s responsible for the atypical away-side jet topologies observed in Au+Au collisions.
To better quantify these jet properties, we split the jet-pair distributions into a near-side range, 0 - $`\mathrm{\Delta }\varphi _{min}`$ and an away-side range $`\mathrm{\Delta }\varphi _{min}`$ \- $`\pi `$. The near- and away-side parts of the distribution are then further characterized by their RMS (taken around 0 and $`\pi `$) and their yield of associated pairs per trigger. These results are summarized in Fig.5 as a function of centrality. For comparison, similar results are included for the 0-20% most central d+Au collisions (open circles) obtained at $`\sqrt{s_{NN}}=200GeV`$. The systematic uncertainty in $`v_2`$ has been propagated into the systematic errors for yields and RMS values. The systematic error on the yields also accounts for the systematic uncertainty on the single particle reconstruction efficiency.
Fig.5b shows that both the near- and away-side widths for peripheral (60-90%) Au+Au collisions compare well with those obtained for d+Au collisions. This is consistent with the expectation of very little, if any, medium induced modifications to jet-topologies in peripheral Au+Au collisions. An inspection of the centrality dependence of the near-side RMS shows no apparent change with centrality. On the other hand, the away-side width is significantly broadened for all but the most peripheral event sample, possibly indicating strong modifications to the fragmentation process by the hot nuclear medium. Although the near-side widths are centrality independent, Fig.5a points to a mild increase in the near-side conditional yield from peripheral to central Au+Au collisions, possibly indicating that even the near-side fragmentation process might still be influenced by the medium. The apparent differences in the evolution of near- and away-side jet characteristics could be signaling the contribution of several different mechanisms to jet-modification.
## 5 Conclusions
In summary, we have used a novel correlation function technique to decompose jet correlations from collective long range harmonic correlations (elliptic flow). We find, that the extracted jet-pair distributions show a strong centrality dependent change in shape and associated per-trigger yield, especially for the away-side jet. The jet-pair distributions obtained for peripheral Au+Au collisions are very similar to those obtained for d+Au collisions but the distributions for more central Au+Au collisions are markedly different and are qualitatively consistent with several recent theoretical predictions of possible modification to jet fragmentation by a strongly interacting medium . Further experimental and theoretical studies are clearly required to establish the detailed mechanism/s responsible for the observed jet modification/s and to pin down the properties of the high energy density strongly interacting matter produced in energetic Au+Au collisions at RHIC. Several such studies are currently being pursued with vigor.
## Acknowledgments
I would like to thank the organizers for giving me the opportunity to present these results in such a stimulating and friendly environment. I am deeply grateful to the Nuclear Chemistry Gang at Stony Brook for much more than I have space left to write. Finally, these acknowledgments would not be complete without a big Thank You to the PHENIX collaboration and the RHIC team. |
warning/0506/cond-mat0506130.html | ar5iv | text | # Introduction to the Keldysh Formalism and Applications to Time-Dependent Density-Functional theory
## I Introduction
We will in this paper give an introduction to the Keldysh formalism, which is an extremely useful tool for first-principles studies of nonequilibrium many-particle systems. Of particular interest for TDDFT is the relation to non-equilibrium Green functions (NEGF), which allows us to construct exchange-correlation potentials with memory by using diagrammatic techniques. For many problems, such as, e.g., quantum transport or atoms in intense laser pulses, one needs exchange-correlation functionals with memory, and Green function techniques offer a systematic method for developing these. The Keldysh formalism is also necessary for defining response functions in TDDFT and for defining an action functional needed for deriving TDDFT from a variational principle. We will in this section give an introduction to the nonequilibrium Green function formalism, intended to illustrate the usefulness of the theory. The formalism does not differ much from ordinary equilibrium theory, the main difference being that all time-dependent functions are definied for time-arguments on a contour, known as the Keldysh contour.
The Green function, $`G(\stackrel{}{r},t;\stackrel{}{r}^{},t^{})`$ is a function of two space- and time-coordinates, and is obviously more complicated than the one-particle density $`n(\stackrel{}{r},t)`$, which is the main ingredient in TDDFT. However, the advantage of the NEGF methods is that we can systematically improve the approximations by taking into account particular physical processes (represented in the form of Feynman diagrams) that we believe to be important. The Green function provides us directly with all expectation values of one-body operators (such as the density and the current), and also the total energy, ionization potentials, response functions, spectral functions, etc.. In relation to TDDFT, this is useful not only for developing orbital functionals and exchange-correlation functionals with memory, but also for providing insight in the exact properties of the non-interacting Kohn-Sham system.
In the following, we shall focus on systems that are initially in thermal equilibrium. We will start by introducing the Keldysh contour and the nonequilbrium Green functions, and then explain how to combine and manipulate functions with time variables on the contour. While we in TDDFT take exchange- and correlation-effects into account through $`v_{\mathrm{xc}}[n]`$, the corresponding quantity in Green function theory is the self-energy $`\mathrm{\Sigma }[G]`$. Just like $`v_{\mathrm{xc}}`$, the self-energy functional must be approximated. For a given functional $`\mathrm{\Sigma }[G]`$, it is important that the resulting observables obey the macroscopic conservation laws, such as, e.g., the continuity equation. These approximations are known as conserving, and will be discussed briefly. In the last part of this section we will discuss the applications of the Keldysh formalism in TDDFT, including the relation between $`\mathrm{\Sigma }`$ and $`v_{\mathrm{xc}}`$, the derivation of the Kohn-Sham equations from an action functional, and the derivation of an $`f_{\mathrm{xc}}`$ functional. As an illustrative example, we will discuss the time-dependent exchange-only optimized effective potential approximation.
## II The Keldysh Contour
In quantum mechanics we associate with any observable quantity $`O`$ a hermitean operator $`\widehat{O}`$. The expectation value $`\mathrm{Tr}\{\widehat{\rho }_0\widehat{O}\}`$ gives the value of $`O`$ when the system is described by the density operator $`\widehat{\rho }_0`$ and the trace denotes a sum over a complete set of states in Hilbert space. For an isolated system the Hamiltonian $`\widehat{H}_0`$ does not depend on time, and the expectation value of any observable quantity is constant, provided $`[\widehat{\rho }_0,\widehat{H}_0]=0`$. In these notes we want to discuss how to describe systems that are isolated for times $`t<0`$, such that $`\widehat{H}(t<0)=\widehat{H}_0`$, but disturbed by an external time-dependent field at $`t>0`$. The expectation value of $`\widehat{O}`$ at $`t>0`$ is then given by the average on the initial density operator $`\widehat{\rho }_0`$ of the operator $`\widehat{O}`$ in the Heisenberg representation,
$$O(t)=\widehat{O}_H(t)\mathrm{Tr}\{\widehat{\rho }_0\widehat{O}_H(t)\}=\mathrm{Tr}\{\widehat{\rho }_0\widehat{S}(0;t)\widehat{O}\widehat{S}(t;0)\},$$
(1)
where the operator in the Heisenberg picture has a time-dependence according to $`\widehat{O}_H(t)=\widehat{S}(0;t)\widehat{O}\widehat{S}(t;0)`$. The evolution operator $`\widehat{S}(t;t^{})`$ is the solution of the equations
$$i\frac{\mathrm{d}}{\mathrm{d}t}\widehat{S}(t;t^{})=\widehat{H}(t)\widehat{S}(t;t^{})\text{and}i\frac{\mathrm{d}}{\mathrm{d}t^{}}\widehat{S}(t;t^{})=\widehat{S}(t;t^{})\widehat{H}(t^{}),$$
(2)
with the boundary condition $`\widehat{S}(t;t)=1`$. It can be formally written as
$$\widehat{S}(t;t^{})=\{\begin{array}{cc}Te^{i_t^{}^t𝑑\overline{t}\widehat{H}(\overline{t})}\hfill & t>t^{}\hfill \\ \overline{T}e^{i_t^{}^t𝑑\overline{t}\widehat{H}(\overline{t})}\hfill & t<t^{}\hfill \end{array}.$$
(3)
In Eq. (3), $`T`$ is the time-ordering operator that rearranges the operators in chronological order with later times to the left; $`\overline{T}`$ is the anti-chronological time-ordering operator. The evolution operator satisfies the group property $`\widehat{S}(t;t_1)\widehat{S}(t_1;t^{})=\widehat{S}(t;t^{})`$ for any $`t_1`$. Notice that if the Hamiltonian is time-independent in the interval between $`t`$ and $`t^{}`$, then the evolution operator becomes $`\widehat{S}(t;t^{})=e^{i\widehat{H}(tt^{})}`$. If we now let the system be initially in thermal equilibrium, with an inverse temperature $`\beta 1/k_BT`$ and chemical potential $`\mu `$, the initial density matrix is $`\widehat{\rho }_0=e^{\beta (\widehat{H}_0\mu \widehat{N})}/\mathrm{Tr}\{e^{\beta (\widehat{H}_0\mu \widehat{N})}\}`$. Assuming that $`\widehat{H}_0`$ and $`\widehat{N}`$ commute, $`\widehat{\rho }_0`$ can be rewritten using the evolution operator $`\widehat{S}`$ with a complex time-argument, $`t=i\beta `$, according to $`\widehat{\rho }_0=e^{\beta \mu \widehat{N}}\widehat{S}(i\beta ;0)/\mathrm{Tr}\{e^{\beta \mu \widehat{N}}\widehat{S}(i\beta ;0)\}`$. Inserting this expression in (1), we find
$$O(t)=\frac{\mathrm{Tr}\left[e^{\beta \mu \widehat{N}}\widehat{S}(i\beta ;0)\widehat{S}(0;t)\widehat{O}\widehat{S}(t;0)\right]}{\mathrm{Tr}\left[e^{\beta \mu \widehat{N}}\widehat{S}(i\beta ;0)\right]}.$$
(4)
Reading the arguments in the numerator from the right to the left, we see that we can design a time-contour $`\gamma `$ with a forward branch going from $`0`$ to $`t`$, a backward branch coming back from $`t`$ and ending in $`0`$, and a branch along the imaginary time-axis from $`0`$ to $`i\beta `$. This contour is illustrated in Fig. 1. Note that the group property of $`\widehat{S}`$ means that we are free to extend this contour up to infinity. We can now generalize (4), and let $`z`$ be a time-contour variable on $`\gamma `$. Letting the variable $`\overline{z}`$ run along this same contour, (4) can be formally recast as
$$O(z)=\frac{\mathrm{Tr}\left[e^{\beta \mu \widehat{N}}T_\mathrm{c}\left\{e^{i_\gamma d\overline{z}\widehat{H}(\overline{z})}\widehat{O}(z)\right\}\right]}{\mathrm{Tr}\left[e^{\beta \mu \widehat{N}}T_\mathrm{c}\left\{e^{i_\gamma d\overline{z}\widehat{H}(\overline{z})}\right\}\right]}.$$
(5)
The contour ordering operator $`T_\mathrm{c}`$ moves the operators with “later” contour variable to the left. In (5), $`\widehat{O}(z)`$ is not the operator in the Heisenberg representation \[the latter is denoted with $`\widehat{O}_H(t)`$\]. The contour-time argument in $`\widehat{O}`$ is there only to specify the position of the operator $`\widehat{O}`$ on $`\gamma `$. A point on the real axis can be either on the forward (we denote these points $`t_{}`$), or on the backward branch (denoted $`t_+`$), and a point which is earlier in real time, can therefore be later on the contour, as illustrated in Fig. 1.
If $`z`$ lies on the vertical track, then there is no need to extend the contour along the real axis. Instead, we have
$$O(z)=\frac{\mathrm{Tr}\left[e^{\beta \mu \widehat{N}}e^{i_z^{i\beta }d\overline{z}\widehat{H}_0}\widehat{O}e^{i_0^zd\overline{z}\widehat{H}_0}\right]}{\mathrm{Tr}\left[e^{\beta (\widehat{H}_0\mu \widehat{N})}\right]}=\frac{\mathrm{Tr}\left[e^{\beta (\widehat{H}_0\mu \widehat{N})}\widehat{O}\right]}{\mathrm{Tr}\left[e^{\beta (\widehat{H}_0\mu \widehat{N})}\right]},$$
where the cyclic property of the trace has been used. The right hand side is independent of $`z`$ and coincides with the thermal average $`\mathrm{Tr}\{\widehat{\rho }_0\widehat{O}\}`$. It is easy to verify that (5) would give exactly the same result for $`O(t)`$, where $`t`$ is real, if the Hamiltonian was time-independent, i.e. $`\widehat{H}(t)=\widehat{H}_0`$ also for $`t>0`$.
To summarize, in (5) the variable $`z`$ lies on the contour of Fig. 1; the r.h.s. gives the time-dependent statistical average of the observable $`O`$ when $`z`$ lies on the forward or backward branch, and the statistical average before the system is disturbed when $`z`$ lies on the vertical track.
## III Nonequilibrium Green Functions
We now introduce the nonequilibrium Green function (NEGF), which is a function of two contour time-variables. In order to keep the notation as light as possible, we here discard the spin degree of freedom; the spin index may be restored later as needed. The field operators $`\widehat{\psi }(\stackrel{}{r})`$, $`\widehat{\psi }^{}(\stackrel{}{r})`$ destroy and create an electron in $`\stackrel{}{r}`$ and obey the anticommutation relations $`\{\widehat{\psi }(\stackrel{}{r}),\widehat{\psi }^{}(\stackrel{}{r}^{})\}=\delta (\stackrel{}{r}\stackrel{}{r}^{})`$. We write the Hamiltonian $`\widehat{H}(t)`$ as the sum of a quadratic term
$$\widehat{h}(t)=d\stackrel{}{r}d\stackrel{}{r}^{}\widehat{\psi }^{}(\stackrel{}{r})\stackrel{}{r}|𝐡(t)|\stackrel{}{r}^{}\widehat{\psi }(\stackrel{}{r}^{})$$
(6)
and the interaction operator
$$\widehat{H}_w=\frac{1}{2}d\stackrel{}{r}d\stackrel{}{r}^{}\widehat{\psi }^{}(\stackrel{}{r})\widehat{\psi }^{}(\stackrel{}{r}^{})w(\stackrel{}{r},\stackrel{}{r}^{})\widehat{\psi }(\stackrel{}{r}^{})\widehat{\psi }(\stackrel{}{r}).$$
(7)
We use boldface to indicate matrices in one-electron labels, e.g., $`𝐡`$ is a matrix and $`\stackrel{}{r}|𝐡|\stackrel{}{r}^{}`$ is the $`(\stackrel{}{r},\stackrel{}{r}^{})`$ matrix element of $`𝐡`$. When describing electrons in an electro-magnetic field, the quadratic term is given by $`\stackrel{}{r}|𝐡(t)|\stackrel{}{r}^{}=\delta (\stackrel{}{r}\stackrel{}{r}^{})\left\{[/i+\stackrel{}{A}(\stackrel{}{r},t)]^2/2+v(\stackrel{}{r},t)\right\}`$.
The definition of an expectation value in (1) can be generalized to the expectation value of two operators. The Green function is defined as
$$G(\stackrel{}{r},z;\stackrel{}{r}^{},z^{})=\stackrel{}{r}|𝐆(z;z^{})|\stackrel{}{r}^{}iT_\mathrm{c}[\widehat{\psi }_H(\stackrel{}{r},z)\widehat{\psi }_H^{}(\stackrel{}{r}^{},z^{})],$$
(8)
where the contour variable in the field operators specifies the position in the contour ordering. The operators have a time-dependence according to the definition of the Heisenberg picture, e.g. $`\widehat{\psi }_H^{}(\stackrel{}{r},z)=\widehat{S}(0;z)\widehat{\psi }^{}(\stackrel{}{r})\widehat{S}(z;0)`$. Notice that if the time-argument $`z`$ is located on the real axis, then $`\widehat{\psi }_H(\stackrel{}{r},t_+)=\widehat{\psi }_H(\stackrel{}{r},t_{})`$. If the time-argument is on the imaginary axis, then $`\widehat{\psi }(\stackrel{}{r},i\tau )`$ is not the adjoint of $`\widehat{\psi }(\stackrel{}{r},i\tau )`$ since $`\widehat{S}^{}(i\tau ;0)\widehat{S}(0;i\tau )`$. The Green function can be written
$$𝐆(z;z^{})=\theta (z,z^{})𝐆^>(z;z^{})+\theta (z^{},z)𝐆^<(z;z^{}).$$
(9)
The function $`\theta (z,z^{})`$ is defined to be 1 if $`z`$ is later on the contour than $`z^{}`$, and 0 otherwise. From the definition of the time-dependent expectation value in Eq. (4), it follows that the greater Green function $`𝐆^>(z;z^{})`$, where $`z`$ is later on the contour than $`z^{}`$, is
$$G^>(\stackrel{}{r},z;\stackrel{}{r}^{},z^{})=\frac{1}{i}\frac{\mathrm{Tr}\left[e^{\beta \mu \widehat{N}}\widehat{S}(i\beta ;0)\widehat{\psi }_H(\stackrel{}{r},z)\widehat{\psi }_H^{}(\stackrel{}{r}^{},z^{})\right]}{\mathrm{Tr}\left[e^{\beta \mu \widehat{N}}\widehat{S}(i\beta ;0)\right]}.$$
(10)
If $`z^{}`$ is later on the contour than $`z`$, then the Green function equals
$$G^<(\stackrel{}{r},z;\stackrel{}{r}^{},z^{})=\frac{1}{i}\frac{\mathrm{Tr}\left[e^{\beta \mu \widehat{N}}\widehat{S}(i\beta ;0)\widehat{\psi }_H^{}(\stackrel{}{r}^{},z^{})\widehat{\psi }_H(\stackrel{}{r},z)\right]}{\mathrm{Tr}\left[e^{\beta \mu \widehat{N}}\widehat{S}(i\beta ;0)\right]}.$$
(11)
The extra minus sign on the right hand side comes from the contour ordering. More generally, rearranging the field operators $`\widehat{\psi }`$ and $`\widehat{\psi }^{}`$ (later arguments to the left), we also have to multiply by $`(1)^P`$, where $`P`$ is the parity of the permutation. From the definition of the Green function, it is easily seen that the electron density, $`n(\stackrel{}{r},z)=\widehat{\psi }_H^{}(\stackrel{}{r},z)\widehat{\psi }_H(\stackrel{}{r},z)`$ and current is obtained according to
$$n(\stackrel{}{r},z)=iG(\stackrel{}{r},z;\stackrel{}{r},z^+)$$
(12)
$$\stackrel{}{j}(\stackrel{}{r},z)=i\left\{\left[\frac{}{2i}\frac{^{}}{2i}+\stackrel{}{A}(\stackrel{}{r},z)\right]G(\stackrel{}{r},z;\stackrel{}{r}^{},z^{})\right\}_{z^{}=z^+}.$$
(13)
where $`z^+`$ indicates that this time-argument is infinitesimally later on the contour.
The Green function $`𝐆(z;z^{})`$ obeys an important cyclic relation on the Keldysh contour. Choosing $`z=0_{}`$, which is the earliest time on the contour, we find $`𝐆(0_{};z^{})=𝐆^<(0;z^{})`$, given by (11) with $`\widehat{\psi }_H(\stackrel{}{r},0)=\widehat{\psi }(\stackrel{}{r})`$. Inside the trace we can move $`\widehat{\psi }(\stackrel{}{r})`$ to the left. Furthermore, we can exchange the position of $`\widehat{\psi }(\stackrel{}{r})`$ and $`e^{\beta \mu \widehat{N}}`$ by noting that $`\widehat{\psi }(\stackrel{}{r})e^{\beta \mu \widehat{N}}=e^{\beta \mu (\widehat{N}+1)}\widehat{\psi }(\stackrel{}{r})`$. Using the group identity $`\widehat{S}(i\beta ;0)\widehat{S}(0;i\beta )=1`$, we obtain
$`G(\stackrel{}{r},0_{};\stackrel{}{r}^{},z^{})`$ $`=`$ $`{\displaystyle \frac{1}{i}}{\displaystyle \frac{\mathrm{Tr}\left[\widehat{\psi }_H(\stackrel{}{r})e^{\beta \mu \widehat{N}}\widehat{S}(i\beta ;0)\widehat{\psi }_H^{}(\stackrel{}{r}^{},z^{})\right]}{\mathrm{Tr}\left[e^{\beta \mu \widehat{N}}\widehat{S}(i\beta ;0)\right]}}.`$ (14)
$`=`$ $`{\displaystyle \frac{e^{\beta \mu }}{i}}{\displaystyle \frac{\mathrm{Tr}\left[e^{\beta \mu \widehat{N}}\widehat{S}(i\beta ;0)\widehat{\psi }_H(\stackrel{}{r},i\beta )\widehat{\psi }_H^{}(\stackrel{}{r}^{},z^{})\right]}{\mathrm{Tr}\left[e^{\beta \mu \widehat{N}}\widehat{S}(i\beta ;0)\right]}}.`$
The r.h.s. equals $`e^{\beta \mu }\stackrel{}{r}|𝐆(i\beta ;z^{})|\stackrel{}{r}^{}`$. Together with a similar analysis for $`𝐆(z;0_{})`$, we conclude that
$$𝐆(0_{};z^{})=e^{\beta \mu }𝐆(i\beta ;z^{})\text{and}𝐆(z;0_{})=e^{\beta \mu }𝐆(z;i\beta ).$$
(15)
These equations constitute the so called Kubo-Martin-Schwinger (KMS) boundary conditions kubo ; martin . From the definition of the Green function in (8), it is easily seen that the $`𝐆(z;z)`$ has a discontinuity in $`z=z^{}`$,
$$𝐆^>(z;z)=𝐆^<(z;z)i\mathrm{𝟏}.$$
(16)
Furthermore, for both time-arguments on the real axis we have the important symmetry $`\left[𝐆^{}(t^{};t)\right]^{}=𝐆^{}(t;t^{})`$. As we shall see, these relations play a crucial role in solving the equation of motion.
## IV The Keldysh Book-Keeping
The Green function belongs to a larger class of functions of two time-contour variables that we will refer to as Keldysh space. These functions can be written on the form
$$k(z;z^{})=\delta (z,z^{})k^\delta (z)+\theta (z,z^{})k^>(z;z^{})+\theta (z^{},z)k^<(z;z^{}),$$
(17)
where the $`\delta `$-function on the contour is defined as $`\delta (z,z^{})=\mathrm{d}\theta (z,z^{})/\mathrm{d}z`$ <sup>1</sup><sup>1</sup>1In general, functions containing singularity of the form $`\mathrm{d}^n\delta (z,z^{})/\mathrm{d}z^n`$ belongs to the Keldysh space, see daniele . These functions are somewhat complicated due to the fact that each of the time-arguments can be located on three different brances of the contour, as illustrated in Fig. 1. Below we systematically derive a set of identities that are commonly used for dealing with such functions and will be used extensively in the following sections. Most of the relations are well known langreth , while others, equally important wagner , are not. Our aim is to provide a self-contained derivation of all of them. A table at the end of the Section summarizes the main results. For those who are not familiar with the Keldysh contour, we strongly recommend to scan what follows with pencil and paper.
It is straightforward to show that if $`a(z;z^{})`$ and $`b(z;z^{})`$ belong to the Keldysh space, then
$$c(z;z^{})=_\gamma d\overline{z}a(z;\overline{z})b(\overline{z};z^{})$$
(18)
also belongs to the Keldysh space. For any $`k(z;z^{})`$ in the Keldysh space we define the greater and lesser functions on the physical time axis
$$k^>(t;t^{})k(t_+;t_{}^{}),k^<(t,t^{})k(t_{};t_+^{}).$$
We also define the following two-point functions with one argument $`t`$ on the physical time axis and the other $`\tau `$ on the vertical track
$$k^{}(t;\tau )k(t_\pm ;\tau ),k^{}(\tau ,t)k(\tau ;t_\pm ).$$
(19)
In the definition of $`k^{}`$ and $`k^{}`$ we can arbitrarily choose $`t_+`$ or $`t_{}`$ since $`\tau `$ is later than both of them. The symbols “$``$” and “$``$” have been chosen in order to help the visualization of the time arguments. For instance, “$``$” has a horizontal segment followed by a vertical one; correspondingly, $`k^{}`$ has a first argument which is real (and thus lies on the horizontal axis) and a second argument which is imaginary (and thus lies on the vertical axis). We will also use the convention of denoting the real time with latin letters and the imaginary time with greek letters.
If we write out the contour integral in (18) in detail, we see with the help of Fig. 1 that the integral consists of four main parts. First, we must integrate along the real axis from $`\overline{z}=0_{}`$ to $`\overline{z}=t_{}^{}`$, for which $`a=a^>`$ and $`b=b^<`$. Then, the integral goes from $`\overline{z}=t_{}^{}`$ to $`\overline{z}=t_+`$, where $`a=a^>`$ and $`b=b^>`$. The third part of the integral goes along the real axis from $`\overline{z}=t_+`$ to $`\overline{z}=0_+`$, with $`a=a^<`$ and $`b=b^>`$. The last integral is along the imaginary track, from $`0_+`$ to $`i\beta `$, where $`a=a^{}`$ and $`b=b^{}`$. In addition, we have the contribution from the singular parts, $`a^\delta `$ and $`b^\delta `$, which is trivial since these integrals involve a $`\delta `$-function. With these specifications, we can drop the $`\pm `$-subscripts on the time-arguments and write
$`c^>(t;t^{})`$ $`=`$ $`a^>(t,t^{})b^\delta (t^{})+a^\delta (t)b^>(t,t^{})+{\displaystyle _0^t^{}}d\overline{t}a^>(t;\overline{t})b^<(\overline{t};t^{})`$
$`+{\displaystyle _t^{}^t}d\overline{t}a^>(t;\overline{t})b^>(\overline{t};t^{})+{\displaystyle _t^0}d\overline{t}a^<(t;\overline{t})b^>(\overline{t};t^{})+{\displaystyle _0^{i\beta }}d\overline{\tau }a^{}(t;\overline{\tau })b^{}(\overline{\tau };t^{}).`$
The second integral on the r.h.s. is an ordinary integral on the real axis of two well defined functions and may be rewritten as
$$_t^{}^td\overline{t}a^>(t;\overline{t})b^>(\overline{t};t^{})=_t^{}^0d\overline{t}a^>(t;\overline{t})b^>(\overline{t};t^{})+_0^td\overline{t}a^>(t;\overline{t})b^>(\overline{t};t^{}).$$
Using this relation, the expression for $`c^>`$ becomes
$`c^>(t;t^{})`$ $`=`$ $`a^>(t,t^{})b^\delta (t^{})+a^\delta (t)b^>(t,t^{}){\displaystyle _0^t^{}}d\overline{t}a^>(t;\overline{t})[b^>(\overline{t};t^{})b^<(\overline{t};t^{})]`$ (20)
$`+`$ $`{\displaystyle _0^t}d\overline{t}[a^>(t;\overline{t})a^<(t;\overline{t})]b^>(\overline{t};t^{})+{\displaystyle _0^{i\beta }}d\overline{\tau }a^{}(t;\overline{\tau })b^{}(\overline{\tau };t^{}).`$
Next, we introduce two other functions on the physical time axis
$`k^\mathrm{R}(t;t^{})`$ $``$ $`\delta (t,t^{})k^\delta +\theta (tt^{})[k^>(t;t^{})k^<(t;t^{})],`$ (21)
$`k^\mathrm{A}(t;t^{})`$ $``$ $`\delta (t,t^{})k^\delta \theta (t^{}t)[k^>(t;t^{})k^<(t;t^{})].`$ (22)
The retarded function $`k^\mathrm{R}(t;t^{})`$ vanishes for $`t<t^{}`$, while the advanced function $`k^\mathrm{A}(t;t^{})`$ vanishes for $`t>t^{}`$. The retarded and advanced functions can be used to rewrite (20) in a more compact form
$$c^>(t;t^{})=_0^{\mathrm{}}d\overline{t}[a^>(t;\overline{t})b^\mathrm{A}(\overline{t};t^{})+a^\mathrm{R}(t;\overline{t})b^>(\overline{t};t^{})]+_0^{i\beta }d\overline{\tau }a^{}(t;\overline{\tau })b^{}(\overline{\tau };t^{}).$$
It is convenient to introduce a short hand notation for integrals along the physical time axis and for those between 0 and $`i\beta `$. The symbol “$``$” will be used to write $`_0^{\mathrm{}}d\overline{t}f(\overline{t})g(\overline{t})`$ as $`fg`$, while the symbol “$``$” will be used to write $`_0^{i\beta }d\overline{\tau }f(\overline{\tau })g(\overline{\tau })`$ as $`fg`$. Then
$$c^>=a^>b^\mathrm{A}+a^\mathrm{R}b^>+a^{}b^{}.$$
(23)
Similarly, one can prove that
$$c^<=a^<b^\mathrm{A}+a^\mathrm{R}b^<+a^{}b^{}.$$
(24)
Equations (23-24) can be used to extract the retarded and advanced component of $`c`$. By definition
$`c^\mathrm{R}(t;t^{})`$ $`=`$ $`\delta (tt^{})c^\delta (t)+\theta (tt^{})[c^>(t;t^{})c^<(t;t^{})]`$
$`=`$ $`a^\delta (t)b^\delta (t^{})\delta (tt^{})+\theta (tt^{}){\displaystyle _0^{\mathrm{}}}d\overline{t}a^\mathrm{R}(t;\overline{t})[b^>(\overline{t};t^{})b^<(\overline{t};t^{})]`$
$`+\theta (tt^{}){\displaystyle _0^{\mathrm{}}}d\overline{t}[a^>(t;\overline{t})a^<(t;\overline{t})]b^\mathrm{A}(\overline{t};t^{}).`$
Using the definitions (21) and (22) to expand the integrals on the r.h.s. of this equation, it is straightforward to show that
$$c^\mathrm{R}=a^\mathrm{R}b^\mathrm{R}.$$
(25)
Proceeding along the same lines, one can show that the advanced component is given by $`c^\mathrm{A}=a^\mathrm{A}b^\mathrm{A}`$. It is worth noting that in the expressions for $`c^\mathrm{R}`$ and $`c^\mathrm{A}`$ no integration along the imaginary track is required.
Next, we show how to extract the components $`c^{}`$ and $`c^{}`$. We first define the Matzubara function $`k^\mathrm{M}(\tau ;\tau ^{})`$ with both the arguments in the interval $`(0,i\beta )`$:
$$k^\mathrm{M}(\tau ;\tau ^{})k(z=\tau ;z^{}=\tau ^{}).$$
Let us focus on $`k^{}`$. Without any restrictions we may take $`t_{}`$ as the first argument in (19). In this case, we find
$$c^{}(t;\tau )=a^\delta (t)b^{}(t;\tau )+_0_{}^t_{}d\overline{z}a^>(t_{};\overline{z})b^<(\overline{z};\tau )+_{t_+}^{0_+}d\overline{z}a^<(t_{};\overline{z})b^<(\overline{z};\tau )+_{0_+}^{i\beta }d\overline{z}a^<(t_{};\overline{z})b(\overline{z};\tau ).$$
(26)
Converting the contour integrals in integrals along the real time axis and along the imaginary track, and taking into account the definition in (21)
$$c^{}=a^\mathrm{R}b^{}+a^{}b^\mathrm{M}.$$
(27)
The relation for $`c^{}`$ can be obtained in a similar way and reads $`c^{}=a^{}b^\mathrm{A}+a^\mathrm{M}b^{}`$. Finally, it is straightforward to prove that the Matzubara component of $`c`$ is simply given by $`c^\mathrm{M}=a^\mathrm{M}b^\mathrm{M}`$.
There is another class of identities we want to discuss for completeness. We have seen that the convolution (18) of two functions belonging to the Keldysh space also belongs to the Keldysh space. The same holds true for the product
$$c(z;z^{})=a(z;z^{})b(z^{};z).$$
Omitting the arguments of the functions, one readily finds (for $`zz^{}`$)
$$c^>=a^>b^<,c^<=a^<b^>,c^{}=a^{}b^{},c^{}=a^{}b^{},c^\mathrm{M}=a^\mathrm{M}b^\mathrm{M}.$$
(28)
The retarded function is then obtained exploiting the identities in (28). We have (for $`tt^{}`$)
$$c^\mathrm{R}(t;t^{})=\theta (tt^{})[a^>(t;t^{})b^<(t^{};t)a^<(t;t^{})b^>(t^{};t)].$$
We may get rid of the $`\theta `$-function by adding and subtracting $`a^<b^<`$ or $`a^>b^>`$ to the above relation and rearranging the terms. The final result is
$$c^\mathrm{R}=a^\mathrm{R}b^<+a^<b^\mathrm{A}=a^\mathrm{R}b^>+a^>b^\mathrm{A}.$$
Similarly one finds $`c^\mathrm{A}=a^\mathrm{A}b^<+a^<b^\mathrm{R}=a^\mathrm{A}b^>+a^>b^\mathrm{R}`$. The time-ordered and anti-time-ordered functions can be obtained in a similar way and the reader can look at Table 1 for the complete list of definitions and identities.
For later purposes, we also consider the case of a Keldysh function $`k(z;z^{})`$ multiplied on the left by a scalar function $`l(z)`$. The scalar function is equivalent to the singular part of a function belonging to Keldysh space, $`\stackrel{~}{l}(z;z^{})=l(z)\delta (z,z^{})`$, meaning that $`\stackrel{~}{l}^{\mathrm{R}/\mathrm{A}}=\stackrel{~}{l}^\mathrm{M}=\stackrel{~}{l}`$ and $`\stackrel{~}{l}^{}=\stackrel{~}{l}^{}=\stackrel{~}{l}^{}=0`$. Using Table 1, one immediately realizes that the function $`l`$ is simply a prefactor: $`_\gamma 𝑑\overline{z}\stackrel{~}{l}(z;\overline{z})k^\mathrm{x}(\overline{z};z^{})=l(z)k^\mathrm{x}(z;z^{})`$, where $`\mathrm{x}`$ is one of the Keldysh components ($``$, $`\mathrm{R},\mathrm{A}`$, $`,`$, $`\mathrm{M}`$). The same is true for $`_\gamma k^\mathrm{x}(z;\overline{z})\stackrel{~}{r}(\overline{z};z^{})=k^\mathrm{x}(z;z^{})r(z^{})`$, where $`\stackrel{~}{r}(z;z^{})=r(z)\delta (z,z^{})`$ and $`r(z)`$ is a scalar function.
## V The Kadanoff-Baym equations
The Green function, as defined in (9), satisfies the equation of motion
$$i\frac{\mathrm{d}}{\mathrm{d}z}𝐆(z;z^{})=\mathrm{𝟏}\delta (z,z^{})+𝐡(z)𝐆(z;z^{})+_\gamma 𝑑\overline{z}𝚺(z,\overline{z})𝐆(\overline{z};z^{}),$$
(29)
as well as the adjoint equation
$$i\frac{\mathrm{d}}{\mathrm{d}z^{}}𝐆(z;z^{})=\mathrm{𝟏}\delta (z,z^{})+𝐆(z;z^{})𝐡(z^{})+_\gamma 𝑑\overline{z}𝐆(z;\overline{z})𝚺(\overline{z},z^{}).$$
(30)
The external potential is included in $`𝐡`$, while the self-energy $`𝚺`$ is a functional of the Green function, and describes the effects of the electron interaction. The self-energy belongs to Keldysh space and can therefore be written on the form $`𝚺(z,z^{})=\delta (z,z^{})𝚺^\delta (z)+\theta (z,z^{})𝚺^>(z,z^{})+\theta (z^{},z)𝚺^<(z,z^{})`$. The singular part of the self-energy can be identified as the Hartree–Fock potential, $`𝚺^\delta (z)=𝐯_\mathrm{H}(z)+𝚺_\mathrm{x}(z)`$. The self-energy obeys the same anti-periodic boundary conditions at $`z=0_{}`$ and $`z=i\beta `$ as $`𝐆`$. We will discuss self-energy approximations in more detail below.
Calculating the Green function on the time-contour now consists of two steps: 1) First one has to find the Green function for imaginary times, which is equivalent to finding the equilibrium Matzubara Green function $`𝐆^\mathrm{M}(\tau ,\tau ^{})`$. This Green function depends only on the difference between the time-coordinates, and satisfies the KMS boundary conditions according to $`𝐆^\mathrm{M}(\tau +i\beta ,\tau ^{})=e^{\beta \mu N}𝐆^\mathrm{M}(\tau ,\tau ^{})`$. Since the self-energy depends on the Green function, this amounts to solving the finite-temperature Dyson equation to self-consistency. 2) The Green function with one or two time-variables on the real axis can now be found by propagating according to (29) and (30). Starting from $`t=0`$, this procedure corresponds to extending the time-contour along the real time-axis. The process is illustrated in Fig. 2.
Writing out the equations for the components of $`𝐆`$ using Table 1, we obtain the equations known as the Kadanoff-Baym equations kb-book ,
$$i\frac{\mathrm{d}}{\mathrm{d}t}𝐆^{}(t;t^{})=𝐡(t)𝐆^{}(t;t^{})+\left[𝚺^\mathrm{R}𝐆^{}\right](t;t^{})+\left[𝚺^{}𝐆^\mathrm{A}\right](t;t^{})+\left[𝚺^{}𝐆^{}\right](t,t^{}),$$
(31)
$$i\frac{\mathrm{d}}{\mathrm{d}t^{}}𝐆^{}(t;t^{})=𝐆^{}(t;t^{})𝐡(t^{})+\left[𝐆^{}𝚺^\mathrm{A}\right](t;t^{})+\left[𝐆^\mathrm{R}𝚺^{}\right](t;t^{})+\left[𝐆^{}𝚺^{}\right](t,t^{}),$$
(32)
$$i\frac{\mathrm{d}}{\mathrm{d}t}𝐆^{}(t;\tau )=𝐡(t)𝐆^{}(t;\tau )+\left[𝚺^\mathrm{R}𝐆^{}\right](t;\tau )+\left[𝚺^{}𝐆^\mathrm{M}\right](t,\tau ),$$
(33)
and
$$i\frac{\mathrm{d}}{\mathrm{d}t}𝐆^{}(\tau ;t)=𝐆^{}(\tau ;t)𝐡(t)+\left[𝚺^{}𝐆^\mathrm{A}\right](\tau ,t)+\left[𝚺^\mathrm{M}𝐆^{}\right](\tau ;t).$$
(34)
It is easily seen that if we denote by $`T`$ the largest of the two time-arguments $`t`$ and $`t^{}`$, then the right hand sides of (31) and (32) depend on $`𝐆^{}(t_1;t_2)`$, $`𝐆^{}(\tau _1,t_2)`$ and $`𝐆^{}(t_1,\tau _2)`$ for $`t_1,t_2T`$. When propagating the Kadanoff-Baym equations one therefore starts at $`t=t^{}=0`$, with the initial conditions given by $`𝐆^<(0;0)=lim_{\eta 0}𝐆^\mathrm{M}(0;i\eta )`$, $`𝐆^>(0;0)=lim_{\eta 0}𝐆^\mathrm{M}(i\eta ;0)`$, $`𝐆^{}(\tau ,0)=𝐆^M(\tau ,0)`$ and $`𝐆^{}(0,\tau )=𝐆^M(0,\tau )`$. One then calculates $`𝐆^{}(t,t^{})`$ for time-arguments within the expanding square given by $`t,t^{}T`$. Simultaneously, one calculates $`𝐆^{}(t,\tau )`$ and $`𝐆^{}(\tau ,t)`$ for $`tT`$. The resulting $`𝐆`$ then automatically satisfies the KMS boundary conditions. The Kadanoff-Baym equations (31) and (33) can both be written in the form
$$i\frac{\mathrm{d}}{\mathrm{d}t}𝐆^\mathrm{x}(t;z^{})=𝐡^{\mathrm{HF}}(t)𝐆^\mathrm{x}(t;z^{})+𝐈^\mathrm{x}(t;z^{}),$$
(35)
while (32) and (34) can be written as the adjoint equations. The term proportional to $`𝐡^{\mathrm{HF}}𝐡+𝚺^\delta `$ describes a free-particle propagation, while $`𝐈^\mathrm{x}`$ is a collision term, which introduces memory effects and dissipation. As can be seen from (3134), the only contribution to $`𝐈^\mathrm{x}(0;0)`$ comes from terms containing time-arguments on the imaginary axis. These terms therefore contain the effect of initial correlations, since the time-derivative of $`𝐆`$ would otherwise correspond to that of a non-interacting system, i.e., $`𝐈^\mathrm{x}(0;0)=0`$.
An example of a time-propagation is given in Fig. 3, which shows the Green function for an H<sub>2</sub> molecule. In this example, the Green function is represented in a basis of Hartree–Fock molecular orbitals, $`\stackrel{}{r}|𝐆(z,z^{})|\stackrel{}{r}^{}=_{ij}\varphi _i(\stackrel{}{r})G_{ij}(z,z^{})\varphi _j^{}(\stackrel{}{r}^{})`$, where $`\varphi _i(\stackrel{}{r})=\stackrel{}{r}|\varphi _i`$.
We have here propagated the Kadanoff-Baym equations using the second Born approximation, illustrated in Fig. 4b. The plots show the imaginary part of the matrix element $`G_{\sigma _g,\sigma _g}^<(t,t^{})`$ calculated for time-variables within the square $`t_1,t_2T=20.0`$ atomic units. In the plot to the left, there is no added external potential and the molecule remains in equilibrium. This means that the Green function only depends on the difference $`t_2t_1`$, and the oscillations in this time-coordinate has a frequency given by the ionization potential of the molecule, in agreement with equilibrium Green function theory fetterwalecka . The figure on the right shows the same matrix element, but now in the presence of an additional electric field which is switched on at $`t=0`$. The oscillations along the ridge $`t_1=t_2`$ can be interpreted as oscillations in the occupation number.
## VI Conserving approximations
In the Dyson-Schwinger equations (29) and (30), we introduced the electronic self-energy functional $`𝚺`$, which accounts for the effects of the electron interaction. The self-energy is a functional of the Green function, and will have to be approximated in practical calculations. Diagrammatic techniques provide a natural scheme for generating approximate self-energies and for systematically improving the approximations. There are no general prescriptions for how to select the relevant diagrams, which means that this selection must be guided by physical intuition. There are, however, important conservation laws, like the number conservation law or the energy conservation law, that should always be obeyed. We will in the following discuss an exact framework for generating such conserving approximations.
Let us first discuss the conservation laws obeyed by a system of interacting electrons, in an external field given by the electrostatic potential $`v(\stackrel{}{r},t)`$ and vector potential $`\stackrel{}{A}(\stackrel{}{r},t)`$. An important relation between these two quantities is provided by the continuity equation
$$\frac{\mathrm{d}}{\mathrm{d}t}n(\stackrel{}{r},t)+\stackrel{}{j}(\stackrel{}{r},t)=0.$$
(36)
The density and the current density can be calculated from the Green function using (12) and (13). Whether these quantities will agree with the continuity equation will depend on whether the Green function is obtained from a conserving self-energy approximation. If we know the current density we can also calculate the total momentum and angular momentum expectation values in the system from the equations
$$\stackrel{}{P}(t)=𝑑\stackrel{}{r}\stackrel{}{j}(\stackrel{}{r},t)\text{and}\stackrel{}{L}(t)=𝑑\stackrel{}{r}\stackrel{}{r}\times \stackrel{}{j}(\stackrel{}{r},t).$$
(37)
For these two quantities the following relations should be satisfied
$`{\displaystyle \frac{\mathrm{d}}{\mathrm{d}t}}\stackrel{}{P}(t)`$ $`=`$ $`{\displaystyle 𝑑\stackrel{}{r}\left[n(\stackrel{}{r},t)\stackrel{}{E}(\stackrel{}{r},t)+\stackrel{}{j}(\stackrel{}{r},t)\times \stackrel{}{B}(\stackrel{}{r},t)\right]}`$ (38)
$`{\displaystyle \frac{\mathrm{d}}{\mathrm{d}t}}\stackrel{}{L}(t)`$ $`=`$ $`{\displaystyle 𝑑\stackrel{}{r}\left[n(\stackrel{}{r},t)\stackrel{}{r}\times \stackrel{}{E}(\stackrel{}{r},t)+\stackrel{}{r}\times (\stackrel{}{j}(\stackrel{}{r},t)\times \stackrel{}{B}(\stackrel{}{r},t))\right]}.`$ (39)
where $`\stackrel{}{E}`$ and $`\stackrel{}{B}`$ are the electric and magnetic fields calculated from
$$\stackrel{}{E}(\stackrel{}{r},t)=v(\stackrel{}{r},t)_t\stackrel{}{A}(\stackrel{}{r},t)\text{and}\stackrel{}{B}(\stackrel{}{r},t)=\times \stackrel{}{A}(\stackrel{}{r},t).$$
(40)
The equations (38) and (39) tell us that the change in momentum and angular momentum is equal to the total force and total torque on the system. In the absence of external fields these equations express momentum and angular momentum conservation. Since the right hand sides of (38) and (39) can also directly be calculated from the density and the current and therefore from the Green function, we may wonder whether they are satisfied for a given approximation to the Green function.
Finally we will consider the case of energy conservation. Let $`E(t)=\widehat{H}(t)`$ be the energy expectation value of the system, then we have
$$\frac{\mathrm{d}}{\mathrm{d}t}E(t)=𝑑\stackrel{}{r}\stackrel{}{j}(\stackrel{}{r},t)\stackrel{}{E}(\stackrel{}{r},t)$$
(41)
This equation tells us that the energy change of the system is equal to the work done on the system. The total energy is calculated from the Green function using the expression
$$E(t)=\frac{i}{2}𝑑\stackrel{}{r}\stackrel{}{r}|\left[i\frac{\mathrm{d}}{\mathrm{d}t}+𝐡(t)\right]𝐆^<(t,t^{})|\stackrel{}{r}|_{t^{}=t}.$$
(42)
The question is now whether the energy and the current calculated from an approximate Green function satisfy the relation in (41).
Baym and Kadanoff baym1 ; baym2 showed that conserving approximations follow immediately if the self-energy is obtained as the functional derivative,
$$\mathrm{\Sigma }(1;2)=\frac{\delta \mathrm{\Phi }}{\delta G(2;1)}.$$
(43)
Here, and in the following discussion, we use numbers to denote the contour coordinates, such that $`1=(\stackrel{}{r}_1,z_1)`$. A functional $`\mathrm{\Phi }`$ can be constructed, as first shown in a seminal paper by Luttinger and Ward lw , by summing over irreducible self-energy diagrams closed with an additional Green function line and multiplied by appropriate numerical prefactors,
$$\mathrm{\Phi }[G]=\underset{n,k}{}\frac{1}{2n}𝑑\overline{1}𝑑\overline{2}\mathrm{\Sigma }_n^{(k)}(\overline{1};\overline{2})G(\overline{2};\overline{1}).$$
(44)
In this summation, $`\mathrm{\Sigma }_n^{(k)}`$ denotes a self-energy diagram of $`n`$-th order, i.e. containing $`n`$ interaction lines. The time-integrals go along the contour, but the rules for constructing Feynman diagrams are otherwise exactly the same as those in the ground-state formalism fetterwalecka . Notice that the functional derivative in (43) may generate other self-energy diagrams in addition to those used in the construction of $`\mathrm{\Phi }`$ in (44).
In Fig. 4 we show some examples of typical $`\mathrm{\Phi }`$ diagrams. Examples of $`\mathrm{\Phi }`$-derivable approximations include Hartree–Fock, the second Born approximation, the $`GW`$ approximation and the $`T`$-matrix approximation.
When the Green function is calculated from a conserving approximation, the resulting observables agree with the conservation laws of the underlying Hamiltonian, as given in (36), (38), (39), and (41). This guarantees the conservation of particles, energy, momentum, and angular momentum. All these conservation laws follow from the invariance of $`\mathrm{\Phi }`$ under specific changes in $`G`$. We will here only outline the principles of the proofs, without going into the details, which can be found in baym1 ; baym2 . 1) Number conservation follows from the gauge invariance of $`\mathrm{\Phi }`$. A gauge transformation $`\stackrel{}{A}(1)\stackrel{}{A}(1)+\mathrm{\Lambda }(1)`$, where $`\mathrm{\Lambda }(\stackrel{}{r},0_{})=\mathrm{\Lambda }(\stackrel{}{r},i\beta )`$ leaves $`\mathrm{\Phi }`$ unchanged. A consequence of the gauge invariance is that a pure gauge cannot induce a change in the density or current. The invariance is therefore closely related to the Ward-identities and to the $`f`$-sum rule for the density response function vanleeuwen04 . 2) Momentum conservation follows from the invariance of $`\mathrm{\Phi }`$ under spatial translations, $`\stackrel{}{r}\stackrel{}{r}+\stackrel{}{R}(z)`$. The invariance is a consequence of the electron interaction $`v(1,2)=\delta (z_1,z_2)/|\stackrel{}{r}_1\stackrel{}{r}_2|`$ being instantaneous and only depending on the difference between the spatial coordinates. 3) Angular momentum conservation follows from the invariance of $`\mathrm{\Phi }`$ under a rotation of the spatial coordinates. 4) Energy conservation follows from the invariance of $`\mathrm{\Phi }`$ when described by an observer using a ”rubbery clock”, measuring time according to the function $`s(z)`$. The invariance relies on the electron interaction being instantaneous.
## VII Non-interacting Electrons
In this Section we focus on non-interacting electrons. This is particularly relevant for TDDFT, where the electrons are described by the non-interacting Kohn-Sham system. While the Kohn-Sham Green function differs from the true Green function, they both produce the same time-dependent density. This is important since the density is not only an important observable in, e.g., quantum transport, but also since the density is the central ingredient in TDDFT. The use of NEGFs in TDDFT is therefore important due to the relation between $`v_{\mathrm{xc}}`$ and the self-energy.
For a system of non-interacting electrons $`\widehat{H}_v=0`$ and it is straightforward to show that the Green function obeys the equations of motion (29) and (30), with $`𝚺=0`$. For any $`zz^{}`$, the equations of motion can be solved by using the evolution operator on the contour,
$$𝐒(z,z^{})=T_\mathrm{c}\left\{e^{i_z^{}^zd\overline{z}𝐡(\overline{z})}\right\},$$
which solves $`i\frac{\mathrm{d}}{\mathrm{d}z}𝐒(z,z^{})=𝐡(z)𝐒(z,z^{})`$ and $`i\frac{\mathrm{d}}{\mathrm{d}z^{}}𝐒(z,z^{})=𝐒(z,z^{})𝐡(z^{})`$. Therefore, any Green function
$$𝐆(z;z^{})=\theta (z,z^{})𝐒(z,0_{})\stackrel{}{f}^>𝐒(0_{},z^{})+\theta (z^{},z)𝐒(z,0_{})\stackrel{}{f}^<𝐒(0_{},z^{}),$$
(45)
satisfying the constraint (16) on the form
$$\stackrel{}{f}^>\stackrel{}{f}^<=i\mathrm{𝟏},$$
(46)
is a solution of the (29-30). In order to fix the matrix $`\stackrel{}{f}^>`$ or $`\stackrel{}{f}^<`$ we impose the KMS boundary conditions. The matrix $`𝐡(z)=𝐡_0`$ for any $`z`$ on the vertical track, meaning that $`𝐒(i\beta ,0_{})=e^{\beta 𝐡_0}`$. Equation (15) then implies $`\stackrel{}{f}^<=e^{\beta (𝐡_0\mu )}\stackrel{}{f}^>`$, and taking into account the constraint (46) we conclude that
$$\stackrel{}{f}^<=\frac{i}{e^{\beta (𝐡_0\mu )}+1}=if(𝐡_0),$$
where $`f(\omega )=1/[e^{\beta (\omega \mu )}+1]`$ is the Fermi distribution function. The matrix $`\stackrel{}{f}^>`$ takes the form $`\stackrel{}{f}^>=i[f(𝐡_0)\mathrm{𝟏}]`$.
The Green function $`𝐆(z;z^{})`$ for a system of non-interacting electrons is now completely fixed. Both $`𝐆^>`$ and $`𝐆^<`$ depend on the initial distribution function $`f(𝐡_0)`$, as it should according to the discussion of Section III. Another way of writing $`i𝐆^<`$ is in terms of the eigenstates $`|\phi _n|\phi _n(0)`$ of $`𝐡_0`$ with eigenvalues $`\epsilon _n`$. From the time-evolved eigenstate $`|\phi _n(t)=𝐒(t,0)|\phi _n`$ we can calculate the time-dependent wavefunction $`\phi _n(\stackrel{}{r},t)=\stackrel{}{r}|\phi _n(t)`$. Inserting $`_n|\phi _n(0)\phi _n(0)|`$ in the expression for $`𝐆^<`$ we find
$$iG^<(\stackrel{}{r},t;\stackrel{}{r}^{},t^{})=i\underset{m,n}{}\stackrel{}{r}|𝐒(t;0)|\phi _m\phi _m|\stackrel{}{f}^<|\phi _n\phi _n|𝐒(0;t)\stackrel{}{r}^{}=\underset{n}{}f(\epsilon _n)\phi _n(\stackrel{}{r},t)\phi _n^{}(\stackrel{}{r}^{},t^{}),$$
(47)
which for $`t=t^{}`$ reduces to the time-dependent density matrix. The Green function $`𝐆^>`$ becomes
$$iG^>(\stackrel{}{r},t;\stackrel{}{r}^{},t^{})=\underset{n}{}[1f(\epsilon _n)]\phi _n(\stackrel{}{r},t)\phi _n^{}(\stackrel{}{r}^{},t^{}),$$
(48)
Knowing the greater and lesser Green functions we can also calculate $`𝐆^{\mathrm{R},\mathrm{A}}`$. By definition we have
$$𝐆^\mathrm{R}(t;t^{})=\theta (tt^{})[𝐆^>(t;t^{})𝐆^<(t;t^{})]=i\theta (tt^{})𝐒(t,t^{}),$$
and similarly
$$𝐆^\mathrm{A}(t;t^{})=i\theta (t^{}t)𝐒(t,t^{})=[𝐆^\mathrm{R}(t^{};t)]^{}.$$
(49)
In the above expressions the Fermi distribution function has disappeared. The information carried by $`𝐆^{\mathrm{R},\mathrm{A}}`$ is the same contained in the one-particle evolution operator. There is no information on how the system is prepared (how many particles, how they are distributed, etc). We use this observation to rewrite $`𝐆^{}`$ in terms of $`𝐆^{\mathrm{R},\mathrm{A}}`$
$$𝐆^{}(t;t^{})=𝐆^\mathrm{R}(t;0)𝐆^{}(0;0)𝐆^\mathrm{A}(0;t^{}).$$
(50)
Thus, $`𝐆^{}`$ is completely known once we know how to propagate the one-electron orbitals in time and how they are populated before the system is perturbed blandin ; cini ; stefanucci . We also observe that an analogous relation holds for $`𝐆^,`$
$$𝐆^{}(t;\tau )=i𝐆^\mathrm{R}(t;0)𝐆^{}(0;\tau ),𝐆^{}(\tau ;t)=i𝐆^{}(\tau ;0)𝐆^\mathrm{A}(0;\tau ).$$
Let us now focus on a special kind of disturbance, namely $`𝐡(t)=𝐡_0+\theta (t)𝐡_1`$. In this case
$$𝐆^\mathrm{R}(t;t^{})=i\theta (tt^{})e^{i(𝐡_0+𝐡_1)(tt^{})}$$
(51)
depends only on the difference between the time arguments. Let us define the Fourier transform of $`𝐆^{\mathrm{R},\mathrm{A}}`$ from
$$𝐆^{\mathrm{R},\mathrm{A}}(t;t^{})=\frac{\mathrm{d}\omega }{2\pi }e^{i\omega (tt^{})}𝐆^{\mathrm{R},\mathrm{A}}(\omega ).$$
The step function can be written as $`\theta (tt^{})=\frac{\mathrm{d}\omega }{2\pi i}\frac{e^{i\omega (tt^{})}}{\omega i\eta }`$, with $`\eta `$ an infinitesimally small positive constant. Substituting this representation of the $`\theta `$-function into (51) and shifting the $`\omega `$ variable one readily finds
$$𝐆^\mathrm{R}(\omega )=\frac{1}{\omega 𝐡_0𝐡_1+i\eta },$$
and therefore $`𝐆^\mathrm{R}(\omega )`$ is analytic in the upper half plane. On the other hand, from (49) it follows that $`𝐆^\mathrm{A}(\omega )=[𝐆^\mathrm{R}(\omega )]^{}`$ is analytic in the lower half plane. What can we say about the greater and lesser component? Do they also depend only on the difference $`tt^{}`$? The answer to the latter question is negative. Indeed, we recall that they contain information on how the system is prepared before $`𝐡_1`$ is switched on. In particular the original eigenstates are eigenstates of $`𝐡_0`$ and in general are not eigenstates of the Hamiltonian $`𝐡_0+𝐡_1`$ at positive times. From (50) one can see that $`𝐆^{}(t;t^{})`$ cannot be expressed only in terms of the time difference $`tt^{}`$. For instance
$$𝐆^<(t;t^{})=e^{i(𝐡_0+𝐡_1)t}if(𝐡_0)e^{i(𝐡_0+𝐡_1)t^{}},$$
and unless $`𝐡_0`$ and $`𝐡_1`$ commute, it is a function of $`t`$ and $`t^{}`$ separately.
It is sometimes useful to split $`𝐡(t)`$ in two parts and treat one of them perturbatively. Let us think, for instance, of a system composed of two connected subsystems $`A+B`$. In case we know how to calculate the Green function of the isolated subsystems $`A`$ and $`B`$, it is convenient to treat the connecting part as a perturbation. Thus, we write $`𝐡(t)=𝓔(t)+𝓥(t)`$, and we define $`𝐠`$ as the Green function when $`𝓥=0`$. The $`𝐠`$ is a solution of
$$\left\{i\frac{\mathrm{d}}{\mathrm{d}z}𝓔(z)\right\}𝐠(z;z^{})=\mathrm{𝟏}\delta (z,z^{}),$$
and of the corresponding adjoint equation of motion. Furthermore, the Green function $`𝐠`$ obeys the KMS boundary conditions. With these we can use $`𝐠`$ to convert the equations of motion for $`𝐆`$ into integral equations
$$𝐆(z;z^{})=𝐠(z;z^{})+_\gamma d\overline{z}𝐠(z;\overline{z})𝓥(\overline{z})𝐆(\overline{z};z^{})=𝐠(z;z^{})+_\gamma d\overline{z}𝐆(z;\overline{z})𝓥(\overline{z})𝐠(\overline{z};z^{});$$
(52)
the integral on $`\overline{z}`$ is along the generalized Keldysh contour of Fig. 1. One can easily check that this $`𝐆`$ satisfies both (29) and (30). $`𝐆`$ also obeys the KMS boundary conditions since the integral equation is defined on the contour of Fig. 1.
In order to get some familiarity with the above perturbation scheme, we consider explicitly the system $`A+B`$ already mentioned. We partition the one-electron Hilbert space in states of the subsystem $`A`$ and states of the subsystem $`B`$. The “unperturbed” system is described by $`𝓔`$, while the connecting part by $`𝓥`$ and
$$𝓔=\left[\begin{array}{cc}𝓔_{AA}& 0\\ 0& 𝓔_{BB}\end{array}\right],𝓥=\left[\begin{array}{cc}0& 𝓥_{AB}\\ 𝓥_{BA}& 0\end{array}\right].$$
Taking into account that $`𝐠`$ has no off-diagonal matrix elements, the Green function projected on one of the two subsystems, e.g., $`𝐆_{BB}`$, is
$$𝐆_{BB}(z;z^{})=𝐠_{BB}(z;z^{})+_\gamma d\overline{z}𝐠_{BB}(z;\overline{z})𝓥_{BA}(\overline{z})𝐆_{AB}(\overline{z};z^{})$$
and
$$𝐆_{AB}(z;z^{})=_\gamma d\overline{z}𝐠_{AA}(z;\overline{z})𝓥_{AB}(\overline{z})𝐆_{BB}(\overline{z};z^{}).$$
Substituting this latter equation into the first one, we obtain a closed equation for $`𝐆_{BB}`$:
$$𝐆_{BB}(z;z^{})=𝐠_{BB}(z;z^{})+_\gamma d\overline{z}d\overline{z}^{}𝐠_{BB}(z;\overline{z})𝚺_{BB}(\overline{z};\overline{z}^{})𝐆_{BB}(\overline{z}^{};z^{}),$$
(53)
with
$$𝚺_{BB}(\overline{z};\overline{z}^{})=𝓥_{BA}(\overline{z})𝐠_{AA}(\overline{z};\overline{z}^{})𝓥_{AB}(\overline{z}^{})$$
the embedding self-energy. The retarded and advanced component can now be easily computed. With the help of Table 1 one finds
$$𝐆_{BB}^{\mathrm{R},\mathrm{A}}=𝐠_{BB}^{\mathrm{R},\mathrm{A}}+𝐠_{BB}^{\mathrm{R},\mathrm{A}}𝚺_{BB}^{\mathrm{R},\mathrm{A}}𝐆_{BB}^{\mathrm{R},\mathrm{A}}.$$
Next, we have to compute the lesser or greater component. As for the retarded and advanced components, this can be done starting from (53). The reader can soon realize that the calculation is rather complicated, due to the mixing of pure real-time functions with function having one real time argument and one imaginary time argument, see Table 1. Below, we use (50) as a feasible short-cut. A closed equation for the retarded and advanced component has been already obtained. Thus, we simply need an equation for $`𝐆^{}(0;0)`$. Let us focus on the lesser component $`𝐆^<(0;0)=i\stackrel{}{f}^<`$. Assuming that the Hamiltonian $`𝐡_0`$ is hermitian, the matrix $`(\omega 𝐡_0)^1`$ has poles at frequencies equal to the eigenvalues of $`𝐡_0`$. These poles are all on the real frequency axis, and we can therefore write
$$𝐆^<(0;0)=if(𝐡_0)=_\gamma \frac{d\zeta }{2\pi }f(\zeta )\frac{1}{\zeta 𝐡_0},$$
(54)
where the contour $`\gamma `$ encloses the real frequency axis.
## VIII Action functional and TDDFT
We define the action functional
$$\stackrel{~}{A}=i\mathrm{ln}\mathrm{Tr}\left\{e^{\beta \mu \widehat{N}}\widehat{S}(i\beta ;0)\right\},$$
(55)
where the evolution operator $`\widehat{S}`$ is the same as defined in (3). The action functional is a tool for generating equations of motion, and is not interesting per se. Nevertheless, one should notice that the action, as defined in (55) has a numerical value equal to $`i\mathrm{ln}Z`$, where $`Z`$ is the thermodynamic partition function.
It is easy to show that if we make a perturbation $`\delta \widehat{V}(z)`$ in the Hamiltonian, the change in the evolution operator is given by
$$i\frac{\mathrm{d}}{\mathrm{d}z}\delta \widehat{S}(z;z^{})=\delta \widehat{V}(z)\widehat{S}(z;z^{})+\widehat{H}(z)\delta S(z;z^{}).$$
(56)
A similar equation for the dependence on $`z^{}`$, and the boundary condition $`\delta \widehat{S}(z;z)=0`$ gives
$$\delta \widehat{S}(z;z^{})=i_z^{}^z𝑑\overline{z}\widehat{S}(z;\overline{z})\delta \widehat{V}(\overline{z})\widehat{S}(\overline{z},z^{}).$$
(57)
We stress that the time-coordinates are on a contour going from $`0`$ to $`i\beta `$. The variation in, e.g., $`V(t_+)`$ is therefore independent of the variation in $`V(t_{})`$. If we let $`\delta \widehat{V}(z)=𝑑\stackrel{}{r}\delta v(\stackrel{}{r},z)\widehat{n}(\stackrel{}{r})`$, a combination of (55) and (57) yields \[compare to (4)\] the expectation values of the density,
$`{\displaystyle \frac{\delta \stackrel{~}{A}}{\delta v(\stackrel{}{r},z)}}`$ $`=`$ $`{\displaystyle \frac{i}{\mathrm{Tr}\left\{e^{\beta \mu \widehat{N}}\widehat{S}(i\beta ;0)\right\}}}{\displaystyle \frac{\delta }{\delta v(\stackrel{}{r},z)}}\mathrm{Tr}\left\{e^{\beta \mu N}\widehat{S}(i\beta ;0)\right\}`$ (58)
$`=`$ $`{\displaystyle \frac{\mathrm{Tr}\left\{e^{\beta \mu \widehat{N}}\widehat{S}(i\beta ;0)\widehat{S}(0;z)\widehat{n}(\stackrel{}{r})\widehat{S}(z,0)\right\}}{\mathrm{Tr}\left\{e^{\beta \mu \widehat{N}}\widehat{S}(i\beta ;0)\right\}}}=n(\stackrel{}{r},z).`$
A physical potential is the same on the positive and on the negative branch of the contour, and the same is true for the corresponding time-dependent density, $`n(\stackrel{}{r},t)=n(\stackrel{}{r},t_\pm )`$. A density response function defined for time-arguments on the contour is found by taking the functional derivative of the density with respect to the external potential. Using the compact notation $`1=(\stackrel{}{r}_1,z_1)`$, the response function is written
$$\chi (1;2)=\frac{\delta n(1)}{\delta v(2)}=\frac{\delta ^2\stackrel{~}{A}}{\delta v(1)\delta v(2)}=\chi (2;1).$$
(59)
This response function is symmetric in the space and time-contour coordinates. We again stress that the variations in the potentials at $`t_+`$ and $`t_{}`$ are independent. If, however, one uses this response function to calculate the density response to an actual physical perturbing electric field, we obtain
$$\delta n(\stackrel{}{r},t)=\delta n(\stackrel{}{r},t_\pm )=_\gamma 𝑑z^{}𝑑\stackrel{}{r}^{}\chi (\stackrel{}{r},t_\pm ;\stackrel{}{r}^{}z^{})\delta v(\stackrel{}{r}^{},z^{}),$$
(60)
where $`\gamma `$ indicates an integral along the contour. In this expression, the perturbing potential (as well as the induced density response) is independent of whether it is located on the positive or negative branch, i.e. $`\delta v(\stackrel{}{r}^{},t_\pm ^{})=\delta v(\stackrel{}{r}^{},t^{})`$. We consider a perturbation of a system initially in equilibrium, which means that $`\delta v(\stackrel{}{r}^{},t^{})0`$ only for $`t^{}>0`$, and we can therefore ignore the integral along the imaginary track of the time-contour. The contour integral then consists of two parts: 1) First an integral from $`t^{}=0`$ to $`t^{}=t`$, in which $`\chi =\chi ^>`$, and 2) an integral from $`t^{}=t`$ to $`t^{}=0`$, where $`\chi =\chi ^<`$. Writing out the contour integral in (60) explicitly then gives
$$\delta n(\stackrel{}{r},t)=_0^t𝑑t^{}𝑑\stackrel{}{r}^{}\left[\chi ^>(\stackrel{}{r}t;\stackrel{}{r}^{}t^{})\chi ^<(\stackrel{}{r}t;\stackrel{}{r}^{}t^{})\right]\delta v(\stackrel{}{r}^{},t^{})=_0^{\mathrm{}}𝑑t^{}𝑑\stackrel{}{r}^{}\chi ^\mathrm{R}(\stackrel{}{r}t;\stackrel{}{r}^{}t^{})\delta v(\stackrel{}{r}^{},t^{}).$$
(61)
The response to a perturbing field is therefore given by the retarded response function, while $`\chi (1,2)`$ defined on the contour is symmetric in $`(12)`$.
If we now consider a system of non-interacting electrons in some external potential $`v_\mathrm{s}`$, we can similarly define a non-interacting action-functional $`\stackrel{~}{A}_\mathrm{s}`$. The steps above can be repeated to calculate the non-interacting response function. The derivation is straightforward, and gives
$$\chi _\mathrm{s}(1;2)=\frac{\delta ^2\stackrel{~}{A}_\mathrm{s}}{\delta v_\mathrm{s}(1)\delta v_\mathrm{s}(2)}=iG_\mathrm{s}(1;2)G_\mathrm{s}(2;1).$$
(62)
The non-interacting Green function $`G_\mathrm{s}`$ has the form given in (45), (47) and (48). The retarded response-function is
$`\chi _\mathrm{s}^\mathrm{R}(\stackrel{}{r}_1,t_1;\stackrel{}{r}_2,t_2)`$ $`=`$ $`i\theta (t_1t_2)\left[G_\mathrm{s}^>(\stackrel{}{r}_1,t_1;\stackrel{}{r}_2,t_2)G_\mathrm{s}^<(\stackrel{}{r}_2,t_2;\stackrel{}{r}_1,t_1)G_\mathrm{s}^<(\stackrel{}{r}_1,t_1;\stackrel{}{r}_2,t_2)G_\mathrm{s}^>(\stackrel{}{r}_2,t_2;\stackrel{}{r}_1,t_1)\right]`$ (63)
$`=`$ $`i{\displaystyle \underset{n,m}{}}[f(\epsilon _m)f(\epsilon _n)]\phi _n(\stackrel{}{r}_1,t_1)\phi _m^{}(\stackrel{}{r}_1,t_1)\phi _m(\stackrel{}{r}_2,t_2)\phi _n^{}(\stackrel{}{r}_2,t_2),`$
where we have used (47) and (48) in the last step.
Having defined the action functional for both the interacting and the non-interacting systems, we now make a Legendre transform, and define
$$A[n]=\stackrel{~}{A}[v]+d(1)n(1)v(1),$$
(64)
which has the property that $`\delta A[n]/\delta n(1)=v(1)`$. Similarly, we define the action functional
$$A_\mathrm{s}[n]=\stackrel{~}{A}_\mathrm{s}[v_\mathrm{s}]+d(1)n(1)v_\mathrm{s}(1).$$
(65)
with the property $`\delta A_\mathrm{s}[n]/\delta n(1)=v_\mathrm{s}(1)`$. The Legendre transforms assume the existence of a one-to-one correspondence between the density and the potential. From these action functionals, we now define the exchange-correlation part to be
$$A_{\mathrm{xc}}[n]=A_\mathrm{s}[n]A[n]\frac{1}{2}d(12)\delta (z_1,z_2)\frac{n(1)n(2)}{|\stackrel{}{r}_1\stackrel{}{r}_2|}.$$
(66)
Taking the functional derivative with respect to the density gives
$$v_\mathrm{s}[n](1)=v(1)+v_\mathrm{H}(1)+v_{\mathrm{xc}}[n](1)$$
(67)
where $`v_H(1)`$ is the Hartree potential and $`v_{\mathrm{xc}}(1)=\delta A_{\mathrm{xc}}/\delta n(1)`$. Again, for time-arguments on the real axis, these potentials are independent of whether the time is on the positive or the negative branch. If we, however, want to calculate the response function from the action functional, then it is indeed important which part of the contour the time-arguments are located on.
We already described how to define response function on the contour, both in the interacting (59) and the non-interacting (62) case. Given the exact Kohn-Sham potential, the TDDFT response function should give exactly the same density change as the exact response function,
$$\delta n(1)=d(2)\chi (1;2)\delta v(2)=d(2)\chi _\mathrm{s}(1;2)\delta v_\mathrm{s}(2).$$
(68)
The change in the Kohn-Sham potential is given by
$`\delta v_\mathrm{s}(1)`$ $`=`$ $`\delta v(1)+{\displaystyle d(2)\frac{\delta v_\mathrm{H}(1)}{\delta n(2)}\delta n(2)}+{\displaystyle d(2)\frac{\delta v_{\mathrm{xc}}(1)}{\delta n(2)}\delta n(2)}`$ (69)
$`=`$ $`\delta v(1)+{\displaystyle d(2)f_{\mathrm{Hxc}}(1;2)\delta n(2)}=\delta v(1)+{\displaystyle d(23)f_{\mathrm{Hxc}}(1;2)\chi (2;3)\delta v(3)}`$
where $`f_{\mathrm{Hxc}}(1;2)=\delta (z_1,z_2)/|\stackrel{}{r}_1\stackrel{}{r}_2|+\delta v_{\mathrm{xc}}(1)/\delta n(2)`$. Inserted in (68), we obtain
$$\chi (1;2)=\chi _\mathrm{s}(1;2)+d(34)\chi _\mathrm{s}(1;3)f_{\mathrm{Hxc}}(3;4)\chi (4;2).$$
(70)
This is the response function defined for time-arguments on the contour. If we want to calculate the response induced by a perturbing potential, the density change will be given by the retarded response function. Using Table 1, we can just write down
$$\chi ^\mathrm{R}(\stackrel{}{r}_1,t_1;\stackrel{}{r}_2,t_2)=\chi _\mathrm{s}^\mathrm{R}(\stackrel{}{r}_1,t_1;\stackrel{}{r}_2,t_2)+𝑑t_3𝑑t_4𝑑\stackrel{}{r}_3𝑑\stackrel{}{r}_4\chi _\mathrm{s}^\mathrm{R}(\stackrel{}{r}_1,t_1;\stackrel{}{r}_3,t_3)f_{\mathrm{Hxc}}^\mathrm{R}(\stackrel{}{r}_3,t_3;\stackrel{}{r}_4,t_4)\chi ^\mathrm{R}(\stackrel{}{r}_4,t_4;\stackrel{}{r}_2,t_2).$$
(71)
The time-integrals in the last expression go from $`0`$ to $`\mathrm{}`$. As expected, only the retarded functions are involved in this expression. We stress the important result that while the function $`f_{\mathrm{Hxc}}(1,2)`$ is symmetric under the coordinate-permutation ($`12`$), it is the retarded function
$`f_{\mathrm{Hxc}}^\mathrm{R}(\stackrel{}{r}_1,t_1;\stackrel{}{r}_2,t_2)={\displaystyle \frac{\delta (t_1,t_2)}{|\stackrel{}{r}_1\stackrel{}{r}_2|}}+f_{\mathrm{xc}}^\mathrm{R}(\stackrel{}{r}_1,t_1;\stackrel{}{r}_2,t_2)`$ (72)
which is used when calculating the response to a perturbing potential.
## IX Example: Time-dependent OEP
We will close this section by discussing the time-dependent optimized effective potential (TDOEP) method in the exchange-only approximation. This is a useful example of how to use functions on the Keldysh contour. While the TDOEP equations can be derived from an action functional, we will here use the time-dependent Sham-Schlüter equations as starting point vanleeuwen96 . This equation is derived by employing a Kohn-Sham Green function, $`G_\mathrm{s}(1,2)`$ which satisfies the equation of motion
$$\left\{i\frac{\mathrm{d}}{\mathrm{d}z_1}+\frac{_1^2}{2}v_\mathrm{s}(\stackrel{}{r}_1,z_1)\right\}G_\mathrm{s}(\stackrel{}{r}_1,z_1;\stackrel{}{r}_2,z_2)=\delta (z_1,z_2)\delta (\stackrel{}{r}_1\stackrel{}{r}_2),$$
(73)
as well as the adjoint equation. The Kohn-Sham Green function is given by (47) and (48) in terms of the time-dependent Kohn-Sham orbitals. Comparing (73) to the Dyson-Schwinger equation (29), we see that we can write an integral equation for the interacting Green function in terms of the Kohn-Sham quantities,
$$G(1;2)=G_\mathrm{s}(1;2)+d(\overline{1}\overline{2})G_\mathrm{s}(1;\overline{1})\left[\mathrm{\Sigma }(\overline{1};\overline{2})+\delta (\overline{1},\overline{2})[v(\overline{1})v_\mathrm{s}(\overline{1})]\right]G(\overline{2};2).$$
(74)
It is important to keep in mind that this integral equation for $`G(1;2)`$ differs from the differential equations (29) and (30) in the sense that we have imposed the boundary conditions of $`G_\mathrm{s}`$ on $`G`$ in (74). This means that if $`G_\mathrm{s}(1;2)`$ satisfies the KMS boundary conditions (15), then so will $`G(1;2)`$.
If we now assume that for any density $`n(1)=iG(1;1^+)`$ there is a potential $`v_\mathrm{s}(1)`$ such that $`n(1)=iG_\mathrm{s}(1;1^+)`$, we obtain the time-dependent Sham-Schlüter equation,
$$d(\overline{1}\overline{2})G_\mathrm{s}(1;\overline{1})\mathrm{\Sigma }(\overline{1};\overline{2})G(\overline{2};1)=d(\overline{1})G_\mathrm{s}(1;\overline{1})[v_\mathrm{s}(\overline{1})v(\overline{1})]G(\overline{1};1).$$
(75)
This equation is formally correct, but not useful in practice since solving it would involve first calculating the nonequilibrium Green function. Instead, one sets $`G=G_\mathrm{s}`$ and $`\mathrm{\Sigma }[G]=\mathrm{\Sigma }[G_\mathrm{s}]`$. For a given self-energy functional, we then have an integral equation for the Kohn-Sham equation. This equation is known as the time-dependent OEP equation. Defining $`\mathrm{\Sigma }=v_H+\mathrm{\Sigma }_{\mathrm{xc}}`$ and $`v_\mathrm{s}=v+v_\mathrm{H}+v_{\mathrm{xc}}`$, the TDOEP equation can be written
$$d(\overline{1}\overline{2})G_\mathrm{s}(1;\overline{1})\mathrm{\Sigma }_{\mathrm{xc}}[G_\mathrm{s}](\overline{1};\overline{2})G_\mathrm{s}(\overline{2};1)=d(\overline{1})G_\mathrm{s}(1;\overline{1})v_{\mathrm{xc}}(\overline{1})G_\mathrm{s}(\overline{1};1).$$
(76)
In the simplest approximation, $`\mathrm{\Sigma }_{\mathrm{xc}}`$ is given by the exchange-only self-energy of Fig. 4a,
$$\mathrm{\Sigma }_\mathrm{x}(1;2)=iG_\mathrm{s}^<(1;2)v(1,2)=\underset{j}{}n_j\varphi _j(1)\varphi _j^{}(2)w(1,2)$$
(77)
where $`n_j`$ is the occupation number. This approximation leads to what is known as the exchange-only TDOEP equations Ullrichetal:PRL95 ; Ullrichetal:BBG95 ; Gorling:PRA97 . Since the exchange self-energy $`\mathrm{\Sigma }_\mathrm{x}`$ is local in time, there is only one time-integration in (76). The x-only solution for the potential will be denoted $`v_\mathrm{x}`$. With the notation $`\stackrel{~}{\mathrm{\Sigma }}(3;4)=\mathrm{\Sigma }_\mathrm{x}(\stackrel{}{r}_3t_3;\stackrel{}{r}_4t_3)\delta (\stackrel{}{r}_3\stackrel{}{r}_4)v_\mathrm{x}(\stackrel{}{r}_3t_3)`$ we obtain from (76)
$`0`$ $`=`$ $`{\displaystyle _0^{t_1}}𝑑t_3{\displaystyle 𝑑\stackrel{}{r}_3𝑑\stackrel{}{r}_4\left[G_\mathrm{s}^<(1;3)\stackrel{~}{\mathrm{\Sigma }}(3;4)G_\mathrm{s}^>(4;1)G_\mathrm{s}^>(1;3)\stackrel{~}{\mathrm{\Sigma }}(3;4)G_\mathrm{s}^<(4;1)\right]}`$ (78)
$`+{\displaystyle _0^{i\beta }}𝑑t_3{\displaystyle 𝑑\stackrel{}{r}_3𝑑\stackrel{}{r}_4G_\mathrm{s}^{}(1;3)\stackrel{~}{\mathrm{\Sigma }}(3;4)G_\mathrm{s}^{}(4;1)}.`$
Let us first work out the last term which describes a time-integral from $`0`$ to $`i\beta `$. On this part of the contour, the Kohn-Sham Hamiltonian is time-independent, with $`v_\mathrm{x}(\stackrel{}{r},0)v_\mathrm{x}(\stackrel{}{r})`$, and $`\phi _i(\stackrel{}{r},t)=\phi _i(\stackrel{}{r})\mathrm{exp}(i\epsilon _it)`$. Since $`\mathrm{\Sigma }_\mathrm{x}`$ is time-independent on this part of the contour, we can integrate
$$_0^{i\beta }𝑑t_3G_\mathrm{s}^{}(1;\stackrel{}{r}_3,t_3)G_\mathrm{s}^{}(\stackrel{}{r}_4,t_3;1)=i\underset{i,k}{}n_i(1n_k)\phi _i(1)\phi _i^{}(\stackrel{}{r}_3)\phi _k(\stackrel{}{r}_4)\phi _k^{}(1)\frac{e^{\beta (\epsilon _i\epsilon _k)}1}{\epsilon _i\epsilon _k}$$
(79)
If we then use $`n_i(1n_k)(e^{\beta (\epsilon _i\epsilon _k)}1)=n_kn_i`$ and define the function $`u_{\mathrm{x},j}`$ by
$$u_{\mathrm{x},j}(1)=\frac{1}{\phi _j^{}(1)}\underset{k}{}n_kd2\phi _j^{}(2)\phi _k(2)\phi _k^{}(1)w(1,2)$$
(80)
we obtain from (79) and (77)
$$_0^{i\beta }𝑑t_3𝑑\stackrel{}{r}_3𝑑\stackrel{}{r}_4G_\mathrm{s}^{}(1,3)\stackrel{~}{\mathrm{\Sigma }}(3,4)G_\mathrm{s}^{}(4,1)=𝑑\stackrel{}{r}_2\underset{j}{}n_j\underset{k_j}{}\frac{\phi _j^{}(\stackrel{}{r}_2)\phi _k(\stackrel{}{r}_2)}{\epsilon _j\epsilon _k}\phi _j(1)\phi _k^{}(1)\left[u_{\mathrm{x},j}(\stackrel{}{r}_2)v_x(\stackrel{}{r}_2)\right]+c.c.$$
(81)
The integral along the real axis on the lhs of (78) can similarly be evaluated. Collecting our results we obtain the OEP equations on the same form as in Gorling:PRA97 ,
$`0`$ $`=`$ $`i{\displaystyle \underset{j}{}}{\displaystyle \underset{kj}{}}n_j{\displaystyle _0^{t_1}}𝑑t_2{\displaystyle 𝑑\stackrel{}{r}_2\left[v_\mathrm{x}(2)u_{\mathrm{x},j}(2)\right]\phi _j(1)\phi _j^{}(2)\phi _k^{}(1)\phi _k(2)}+c.c.`$ (82)
$`+{\displaystyle \underset{j}{}}{\displaystyle \underset{kj}{}}n_j{\displaystyle \frac{\phi _j(1)\phi _k^{}(1)}{\epsilon _j\epsilon _k}}{\displaystyle 𝑑\stackrel{}{r}_2\phi _j^{}(\stackrel{}{r}_2)\left[v_\mathrm{x}(\stackrel{}{r}_2)u_{\mathrm{x},j}(\stackrel{}{r}_2)\right]\phi _k(\stackrel{}{r}_2)}.`$
The last term represents the initial conditions, expressing that the system is in thermal equilibrium at $`t=0`$. The equations have exactly the same form if the initial condition is specified at some other initial time $`t_0`$. In the case that we let $`t_0\mathrm{}`$, the term due to the initial conditions vanish and the remaining expression equals the one given in vanleeuwen96 ; Ullrichetal:PRL95 ; Ullrichetal:BBG95 . The OEP-equations (82) in the so-called KLI-approximation have been successfully used by Ullrich et al.Ullrichetal:BBG95 to calculate properties of atoms in strong laser fields. |
warning/0506/cond-mat0506477.html | ar5iv | text | # Quantum Transport from the Perspective of Quantum Open Systems
## Abstract
By viewing the non-equilibrium transport setup as a quantum open system, we propose a reduced-density-matrix based quantum transport formalism. At the level of self-consistent Born approximation, it can precisely account for the correlation between tunneling and the system internal many-body interaction, leading to certain novel behavior such as the non-equilibrium Kondo effect. It also opens a new way to construct time-dependent density functional theory for transport through large-scale complex systems.
Conventionally, quantum transport and quantum dissipation are treated with different approaches. For instance, the former (in mesoscopic context) is usually described by the Landauer-Büttiker theory and the non-equilibrium Green’s function (nGF) approach Dat95 ; Hau96 , whereas the latter by the reduced density matrix equation Gar00 . Nevertheless, both are quantum open systems, with either the non-equilibrium electron reservoirs (electrodes) or the dissipative thermal bath as their environments, as schematically shown in Fig. 1. It is thus of great interest to develop a unified language to bridge them.
This motivation can be historically dated back to the phenomenological rate equation and quantum Bloch equation approaches to transport Gla88 ; Gur96 . There, either implicitly or explicitly, the electrodes are viewed as dissipative reservoirs. Along this line and based on our work in quantum measurement of solid-state qubit Li04a ; Li04b , we developed recently a quantum master equation approach to quantum transport Li04c . The reduced dynamics involved there was originally constructed under the cumulant second-order expansion (Born approximation). In this letter, we re-formalize it in the spirit of self-consistent Born approximation (SCBA), in order to make the formalism not only convenient but also accurate enough in practice. Moreover, by reducing the many-particle density matrix formalism to single-particle one, we will also construct an efficient approach for large-scale (e.g. molecular) system applications.
The typical transport setup, see Fig. 1(A), can be described by
$`H`$ $`=`$ $`H_S(a_\mu ^{},a_\mu )+{\displaystyle \underset{\alpha =L,R}{}}{\displaystyle \underset{\mu k}{}}ϵ_{\alpha \mu k}d_{\alpha \mu k}^{}d_{\alpha \mu k}`$ (1)
$`+{\displaystyle \underset{\alpha =L,R}{}}{\displaystyle \underset{\mu k}{}}(t_{\alpha \mu k}a_\mu ^{}d_{\alpha \mu k}+\mathrm{H}.\mathrm{c}.).`$
$`H_S`$ is the system Hamiltonian, which can be rather general (e.g. including many-body interaction). $`a_\mu ^{}`$ ($`a_\mu `$) is the creation (annihilation) operator of electron in state labelled by “$`\mu `$”, which labels both the multi-orbital and distinct spin states of the transport system. The second and third terms describe, respectively, the two (left and right) electrodes and the tunneling between the electrodes and the system.
To contact with the quantum dissipation theory for quantum open systems, let us introduce the reservoir operators $`F_\mu =_{\alpha k}t_{\alpha \mu k}d_{\alpha \mu k}f_{L\mu }+f_{R\mu }`$. Accordingly, the tunneling Hamiltonian $`H^{}`$ reads $`H^{}=_\mu (a_\mu ^{}F_\mu +\mathrm{H}.\mathrm{c}.)`$. Treating $`H^{}`$ perturbatively up to the cumulant second-order expansion, it yields Yan98
$$\dot{\rho }(t)=i\rho (t)_0^t𝑑\tau ^{}(t)𝒢(t,\tau )^{}(\tau )𝒢^{}(t,\tau )\rho (t).$$
(2)
The reduced density matrix is defined as $`\rho (t)=\mathrm{Tr}_\mathrm{B}[\rho _\mathrm{T}(\mathrm{t})]`$, by tracing out the reservoir degrees of freedom from the total system-plus-reservoirs density matrix. The Liouvillian superoperators are defined as $`(\mathrm{})[H_S,(\mathrm{})]`$, $`^{}(\mathrm{})[H^{},(\mathrm{})]`$, and $`𝒢(t,\tau )(\mathrm{})G(t,\tau )(\mathrm{})G^{}(t,\tau )`$ with $`G(t,\tau )`$ the usual propagator (Green’s function) associated with the system Hamiltonian $`H_S`$.
The integral kernel in Eq. (2), which is in the so-called partial ordering prescription (POP) (or time-local) form Yan98 , describes the second-order tunneling self-energy. At the second-order level, one may replace $`\rho (t)`$ in the last term of Eq. (2) with $`𝒢(t,\tau )\rho (\tau )`$, leading the tunneling integral kernel to $`^{}(t)𝒢(t,\tau )^{}(\tau )\rho (\tau )`$, being in the chronological ordering prescription (COP) (or memory) form Yan98 . The corresponding four terms in the conventional Hilbert space Li04b ; Li04c , depicted on the real-time Keldysh contour in Fig. 2, provide a clear diagrammatic interpretation for the second-order tunneling self-energy process.
Explicitly tracing out the states of electrodes, Eq. (2) gives rise to
$`\dot{\rho }`$ $`=`$ $`i\rho {\displaystyle \underset{\mu }{}}\{[a_\mu ^{},A_\mu ^{()}\rho \rho A_\mu ^{(+)}]+\mathrm{H}.\mathrm{c}.\}.`$ (3)
For time-independent system Hamiltonian, $`A_\mu ^{(\pm )}=_{\alpha =L,R}A_{\alpha \mu }^{(\pm )}=_{\alpha \nu }_{\mathrm{}}^{\mathrm{}}𝑑tC_{\alpha \mu \nu }^{(\pm )}(t)[i\stackrel{~}{a}_\nu (t)]`$, with $`C_{\alpha \mu \nu }^{(+)}(t)f_{\alpha \mu }^{}(t)f_{\alpha \nu }(0)`$, $`C_{\alpha \mu \nu }^{()}(t)f_{\alpha \mu }(t)f_{\alpha \nu }^{}(0)`$, and $`\stackrel{~}{a}_\nu (t)=i\mathrm{\Theta }(t)𝒢(t,0)a_\nu \mathrm{\Pi }^{(0)}(t,0)a_\nu `$. Note that the step function $`\mathrm{\Theta }(t)`$ extends the lower bound of the time integral from $`0`$ to $`\mathrm{}`$, whereas the extension of the upper bound from $`t`$ to $`\mathrm{}`$ results from the consideration of Markovian approximation.
For time-dependent system Hamiltonian, the time-translational invariance breaks down, we thus define $`\stackrel{~}{a}_\nu (t,t^{})=i\mathrm{\Theta }(tt^{})𝒢(t,t^{})a_\nu `$. The backward evolution of $`\stackrel{~}{a}_\nu (t,t^{})`$ with respect to $`t^{}`$, starting from $`t^{}=t`$, can be carried out via $`_t^{}a_\nu (t,t^{})=i\delta (tt^{})a_\nu +i[H_S(t^{}),a_\nu (t,t^{})]`$. Thus, the time integral in $`A_\mu ^{(\pm )}`$, which becomes now the type of $`_0^t𝑑t^{}C_{\alpha \mu \nu }^{(\pm )}(tt^{})\stackrel{~}{a}_\nu (t,t^{})`$, can be calculated accordingly. Inserting the obtained $`A_\mu ^{(\pm )}`$ into Eq. (3), the time-dependent phenomena associated with either quantum dissipative dynamics or transport current can be easily treated. For clarity, we hereafter assume the system Hamiltonian to be time-independent, unless further specification.
Now we consider the possibility to go beyond the second-order self-energy process diagrammatically shown in Fig. 1. An efficient scheme follows the idea of the well-known self-consistent Born approximation (SCBA), i.e., the free propagator defined above, $`\mathrm{\Pi }^{(0)}(t)i\mathrm{\Theta }(t)𝒢(t,0)`$, is replaced by an effective propagator $`\mathrm{\Pi }(t)`$. The latter is obtained formally via the Dyson equation note-2
$`\dot{\mathrm{\Pi }}(t)=i\delta (t)i\mathrm{\Pi }(t)i{\displaystyle _{\mathrm{}}^t}𝑑t^{}\mathrm{\Sigma }(tt^{})\mathrm{\Pi }(t^{}),`$ (4)
or $`\mathrm{\Pi }(\omega )=[\omega \mathrm{\Sigma }(\omega )]^1`$ in frequency domain, with $`\mathrm{\Sigma }`$ being the irreducible self-energy defined again by Fig. 2. Accordingly, $`\stackrel{~}{a}_\nu (t)=\mathrm{\Pi }(t)a_\nu `$, and
$`A_\mu ^{(\pm )}={\displaystyle \underset{\alpha \nu }{}}{\displaystyle \frac{d\omega }{2\pi }C_{\alpha \mu \nu }^{(\pm )}(\pm \omega )[i\mathrm{\Pi }(\omega )a_\nu ]}.`$ (5)
Equations (3)–(5) constitute the basic ingredients of the proposed SCBA scheme. This type of self-consistently partial summation correction has included an infinite number of higher order tunneling processes into the reduced system dynamics. The resulting non-trivial effect on quantum transport will be demonstrated soon.
So far, the trace is performed over all the electrode states, and the resulting Eq. (3) is a unconditional master equation for the system reduced dynamics. To characterize the transport problem, we should keep track of the record of electron numbers entering the right reservoir (electrode). Following Refs. Li04b, and Li04c, , one can obtain a conditional master equation for the reduced system density matrix, $`\rho ^{(n)}(t)`$, under the condition that $`n`$ electrons have arrived at the right electrode until time $`t`$. On the basis of $`\rho ^{(n)}(t)`$, one is readily able to compute various transport properties, such as the transport current, the probability distribution function $`P(n,t)\mathrm{Tr}[\rho ^{(\mathrm{n})}(\mathrm{t})]`$, and the noise spectrum Li04b . For transport current, it can be carried out via $`I(t)=e_nn\mathrm{Tr}[\dot{\rho }^{(\mathrm{n})}(\mathrm{t})`$, giving rise to
$`I(t)=e{\displaystyle \underset{\mu }{}}\mathrm{Tr}[(\mathrm{a}_\mu ^{}\mathrm{A}_{\mathrm{R}\mu }^{()}\mathrm{A}_{\mathrm{R}\mu }^{(+)}\mathrm{a}_\mu ^{})\rho (\mathrm{t})+\mathrm{H}.\mathrm{c}.].`$ (6)
Compared to other transport formalisms, Eqs. (3)-(6) provide a convenient framework for quantum transport.
As an illustrative application, we consider the non-trivial problem of quantum transport through strongly interacting quantum dot, under the well-known Anderson impurity model Hamiltonian: $`H_S=_\mu (ϵ_0a_\mu ^{}a_\mu +\frac{U}{2}n_\mu n_{\overline{\mu }})`$. Here the index $`\mu `$ labels the spin up (“$``$”) and spin down (“$``$”) states, and $`\overline{\mu }`$ stands for the opposite spin orientation. The electron number operator $`n_\mu =a_\mu ^{}a_\mu `$, and the Hubbard term $`Un_{}n_{}`$ describe the charging effect. Apparently, the reservoir correlation function is diagonal with respect to the spin indices, i.e., $`C_{\alpha \mu \nu }^{(\pm )}(t)=\delta _{\mu \nu }C_{\alpha \mu \mu }^{(\pm )}(t)`$, and $`C_{\alpha \mu \mu }^{(\pm )}(t)=_k|t_{\alpha \mu k}|^2e^{\pm iϵ_k(t)}n_\alpha ^{(\pm )}(ϵ_k)`$. Here $`n_\alpha ^{(+)}(ϵ_k)=n_\alpha (ϵ_k)`$ is the Fermi distribution function, and $`n_\alpha ^{()}(ϵ_k)=1n_\alpha (ϵ_k)`$. Accordingly, we have $`A_{\alpha \mu }^{(\pm )}=\mathrm{\Gamma }_{\alpha \mu }\frac{dϵ}{2\pi }n_\alpha ^{(\pm )}(ϵ)[i\mathrm{\Pi }(ϵ)]a_\mu `$, where, under the wide-band approximation, we have introduced $`\mathrm{\Gamma }_{\alpha \mu }=2\pi g_\alpha |t_{\alpha \mu k}|^2`$, and assumed it energy independent. From Eqs. (6) and (3), the stationary current is obtained as
$`I`$ $`=`$ $`{\displaystyle \frac{e\mathrm{\Gamma }_L\mathrm{\Gamma }_R}{\mathrm{\Gamma }_L+\mathrm{\Gamma }_R}}{\displaystyle \frac{dϵ}{2\pi }\mathrm{Im}[\mathrm{\Pi }(ϵ)][n_L(ϵ)n_R(ϵ)]}.`$ (7)
For the single level system under study, the propagator in energy space simply reads $`\mathrm{\Pi }(ϵ)=[ϵϵ_0\mathrm{\Sigma }(ϵ)]^1`$.
Within the SCBA scheme, the self-energy $`\mathrm{\Sigma }`$ can be explicitly carried out via Fig. 2. However, in the case of strong Coulomb repulsion, the dot can be occupied at most by one electron. As a result, it can be easily proven that only Fig. 2(C) and (D) contribute to the self-energy. Physically, replacing the bare system propagator with the effective propagator corresponds to including the infinitely multiple forward and backward tunnelings between the system and the same electrode. This is in fact a tunneling-induced quantum fluctuation, which would lead to the level broadening and the non-trivial interference between tunneling and system internal interaction. Explicitly, in large-$`U`$ limit, the real and imaginary parts of the self-energy read $`\mathrm{Re}\mathrm{\Sigma }(ϵ)=(m1)_{\alpha =L,R}\frac{\mathrm{\Gamma }_{\alpha \mu }}{2\pi }\left[\mathrm{ln}\left(\frac{\beta U}{2\pi }\right)\mathrm{Re}\psi \left(\frac{1}{2}+i\frac{\beta }{2\pi }(ϵ\mu _\alpha )\right)\right]`$ and $`\mathrm{Im}\mathrm{\Sigma }(ϵ)=_{\alpha =L,R}\frac{\mathrm{\Gamma }_{\alpha \mu }}{2}\left[1+(m1)n_\alpha (ϵ)\right]`$, respectively Sch94 ; Kon96 . Here $`\beta 1/(k_BT)`$ is the inverse temperature, $`\mu _\alpha `$ the chemical potential of the electrode, $`\psi `$ the digamma function, and $`m`$ denotes the spin degeneracy. (i) For $`m=1`$, i.e. neglecting the spin degree of freedom, $`\mathrm{Im}\mathrm{\Pi }(ϵ)`$ gives the well-known Breit-Wigner formula, which appropriately includes the level broadening effect. (ii) For $`m2`$ (e.g., $`m=2`$ for spin $`1/2`$), the above self-energy correction would result in rich behaviors, depending on the relative values of the parameters such as the temperature and the position of $`ϵ_0`$ with respect to the Fermi levels. Detailed discussions, in particular the non-equilibrium Kondo effect, are referred to literature, e.g. Refs. Kon96, -Mei93, .
Application to Large Scale Systems.— By far, the transport-related density matrix formalism has been constructed in many-particle Hilbert space, which may restrict its direct application only in small systems. For large-scale systems in the absence of many-electron interaction, we first recast the formalism to a very simple version in terms of the reduced single-particle density matrix (RSPDM), $`\sigma _{\mu \nu }(t)\mathrm{Tr}[\mathrm{a}_\nu ^{}\mathrm{a}_\mu \rho (\mathrm{t})]`$, which greatly reduces the dimension of Hilbert space, thus saves computing expense. To account for the electron-electron interaction, we then propose an efficient time-dependent density functional theory (TDDFT) scheme. Note that it is quite natural to combine the TDDFT technique with the present RSPDM formalism, since the former self-consistently amounts to the many-body interaction but still keeps the single-particle picture note-5 .
(i) Time Independent System Hamiltonian: For simplicity, we first proceed our derivation in the single-particle eigenstate basis, which is denoted as $`\{|\mu ,|\nu ,\mathrm{}\}`$. In this representation, $`A_{\alpha \mu }^{(\pm )}=_\nu C_{\alpha \mu \nu }^{(\pm )}(ϵ_\nu )a_\nu `$, and the equation of motion for the RSPDM can be readily obtained by applying Eq. (3) directly for $`\sigma _{\mu \nu }(t)=\mathrm{Tr}[\mathrm{a}_\nu ^{}\mathrm{a}_\mu \rho (\mathrm{t})]`$. We have note-4 ; Yok04
$`\dot{\sigma }=i[h,\sigma ]{\displaystyle \frac{1}{2}}\{[C^{()}\sigma C^{(+)}\overline{\sigma }]+\mathrm{H}.\mathrm{c}.\}.`$ (8)
Here, $`h`$ is the single-particle Hamiltonian or the Fock matrix within the TDDFT framework which will be identified soon. $`\overline{\sigma }1\sigma `$ denotes the “hole” density matrix. The involving matrix products are defined as usual; e.g., $`[C^{()}\sigma ]_{\mu \nu }_{\alpha =L,R}_\nu ^{}C_{\alpha \mu \nu ^{}}^{()}(ϵ_\nu ^{})\sigma _{\nu ^{}\nu }`$. Straightforwardly, the current can be expressed in terms of the RSPDM as
$`I(t)=e\mathrm{Re}\left\{\mathrm{Tr}\left[\mathrm{C}_\mathrm{R}^{()}\sigma (\mathrm{t})\mathrm{C}_\mathrm{R}^{(+)}\overline{\sigma }(\mathrm{t})\right]\right\}.`$ (9)
In arbitrary state basis, derivation is the same as above. The difference lies only at the expression of $`A_{\alpha \mu }^{(\pm )}`$, which in a non-eigenstate representation is given formally as $`A_{\alpha \mu }^{(\pm )}_\nu \stackrel{~}{C}_{\alpha \mu \nu }^{(\pm )}a_\nu =_{\nu \nu ^{},m}C_{\alpha \mu \nu ^{}}^{(\pm )}(ϵ_m)D_{\nu ^{}m}^1D_{m\nu }a_\nu `$. Here $`ϵ_m`$ is the eigen-energy of eigenstate $`|m`$, and $`D`$ is the transformation matrix from the non-eigestate representation to the eigenstate one. Obviously, with this identification, the resultant master equation and current formula are the same as Eqs. (8) and (9), only replacing the matrices $`C^{(\pm )}`$ by $`\stackrel{~}{C}^{(\pm )}`$.
As an illustrative application of Eqs. (8) and (9), we consider the simple non-interacting multi-level model studied in Ref. Li04c, . In the non-equilibrium stationary state, $`\sigma (t\mathrm{})`$ is diagonal in the eigenstate basis, thus $`[h,\sigma ]=0`$. As a consequence, the stationary state solution is determined by $`C_{\mu \mu }^{()}(ϵ_\mu )\sigma _{\mu \mu }=C_{\mu \mu }^{(+)}(ϵ_\mu )(1\sigma _{\mu \mu })`$, leading to the well-known result Li04c , $`\sigma _{\mu \mu }=[\mathrm{\Gamma }_L(ϵ_\mu )n_L(ϵ_\mu )+n_R(ϵ_\mu )\mathrm{\Gamma }_R(ϵ_\mu )]/[\mathrm{\Gamma }_L(ϵ_\mu )+\mathrm{\Gamma }_R(ϵ_\mu )]`$. In particular, in the special case of equilibrium, $`\sigma _{\mu \mu }`$ reduces to the Fermi-Dirac function. Substituting $`\sigma _{\mu \mu }`$ into Eq. (9), the well-known resonant tunnel current is obtained.
(ii) Time Dependent System Hamiltonian: In this case, the RSPDM can be introduced in a similar manner. Consider, for example, $`\mathrm{Tr}[a_\mu ^{}A_\nu ^{()}\rho (t)]=_{\alpha \nu ^{}}_0^t𝑑t^{}C_{\alpha \nu \nu ^{}}^{()}(t,t^{})\sigma _{\nu ^{}\mu }(t^{},t)[C^{()}\sigma ]_{\nu \mu }`$. Here, $`\sigma _{\nu ^{}\mu }(t^{},t)\mathrm{Tr}\{a_\mu ^{}[𝒢(t,t^{})a_\nu ^{}]\rho (t)\}`$, which can be solved via $`_t^{}\sigma _{\nu ^{}\mu }(t^{},t)=i[h(t^{})\sigma (t^{},t)]_{\nu ^{}\mu }`$, with the initial condition $`\sigma _{\nu ^{}\mu }(t,t)=\mathrm{Tr}[a_\mu ^{}a_\nu ^{}\rho (t)]`$. Similarly, we have $`\mathrm{Tr}[A_\nu ^{(+)}a_\mu ^{}\rho (t)]=_{\alpha \nu ^{}}_0^t𝑑t^{}C_{\alpha \nu \nu ^{}}^{(+)}(t,t^{})\overline{\sigma }_{\nu ^{}\mu }(t^{},t)[C^{(+)}\overline{\sigma }]_{\nu \mu }`$. Here, $`\overline{\sigma }_{\nu ^{}\mu }(t^{},t)\mathrm{Tr}\{[𝒢(t,t^{})a_\nu ^{}]a_\mu ^{}\rho (t)\}`$, satisfying an equation of the same form as $`\sigma _{\nu ^{}\mu }(t^{},t)`$, but with initial condition $`\overline{\sigma }_{\nu ^{}\mu }(t,t)=\delta _{\nu ^{}\mu }\sigma _{\nu ^{}\mu }(t)`$. As a result, in the time-dependent case, the resultant master equation and transport current can also be expressed as Eqs. (8) and (9), only keeping in mind that the matrices product needs not only the inner-state summation, but also the “inner-time” integration.
Now we extend the above RSPDM formalism, i.e., Eqs. (8) and (9), to interacting systems. Within the TDDFT framework RG84 , this can be straightforwardly done by replacing the single particle Hamiltonian by the Fock matrix
$`h_{mn}(t)=h_{mn}^0(t)+v_{mn}^{\mathrm{xc}}(t)+{\displaystyle \underset{ij}{}}\sigma _{ij}(t)V_{mnij}.`$ (10)
In first-principles calculation the state basis is usually chosen as the local atomic orbitals, $`\{\varphi _m(𝐫),m=1,2,\mathrm{}\}`$. Here $`h^0(t)`$ is the non-interacting Hamiltonian which can be in general time-dependent; $`V_{mnij}`$ is the two-electron Coulomb integral, $`V_{mnij}=𝑑𝐫𝑑𝐫^{}\varphi _m^{}(𝐫)\varphi _n(𝐫)\frac{1}{|𝐫𝐫^{}|}\varphi _i^{}(𝐫^{})\varphi _j(𝐫^{})`$; and $`v_{mn}^{\mathrm{xc}}(t)=𝑑𝐫\varphi _m^{}(𝐫)v^{\mathrm{xc}}[n](𝐫,t)\varphi _n(𝐫)`$, with $`v^{\mathrm{xc}}[n](𝐫,t)`$ the exchange-correlation potential, which is defined by the functional derivative of the the exchange-correlation functional $`A^{\mathrm{xc}}`$. In practice, especially in the time-dependent case, the unknown functional $`A^{\mathrm{xc}}`$ can be approximated by the energy functional $`E^{\mathrm{xc}}`$, obtained in the Kohn-Sham theory and further with the local density approximation (LDA). Notice that the density function $`n(𝐫,t)`$ appeared in the Fock operator is related to the RSPDM via $`n(𝐫,t)=_{mn}\varphi _m(𝐫)\sigma _{mn}(t)\varphi _n^{}(𝐫)`$. Thus, Eqs. (8)-(10) constitute a closed form of TDDFT approach for the first-principles study of quantum transport, which is currently an intensive research subject Bur05 .
To summarize, we have proposed a compact transport formalism from the perspective of quantum open systems. The new formulation is constructed in terms of an improved reduced density matrix approach at the SCBA level, which is shown to be accurate enough in practice. Based on the established density matrix formalism, we also developed a new TDDFT scheme for first-principles study of transport through complex large-scale systems. Systematic applications and numerical implementations are in progress and will be published elsewhere.
Acknowledgments. Support from the National Natural Science Foundation of China and the Research Grants Council of the Hong Kong Government is gratefully acknowledged. |
warning/0506/math0506127.html | ar5iv | text | # A Note on the Ruin Problem with Risky Investments
## 1 Introduction
The classical Cramér-Lundberg risk model may be written as
$$X_t=u+ct\underset{k=1}{\overset{N_t}{}}Y_t$$
(1.1)
where $`X_t`$ is the capital of the insurance company, $`u`$ is the initial capital, $`c`$ is the premium rate, $`N_t`$ is a Poisson process of rate $`\lambda `$, and $`\{Y_t\}_{k=1}^{N_t}`$ is a sequence of positive i.i.d. random variables, modelling the claim sizes (independently of $`N_t`$).
The probability of ruin is written as
$$\mathrm{\Psi }(u)=P(X_t=u+ct\underset{k=1}{\overset{N_t}{}}Y_t0,\text{for some}t>0)$$
(1.2)
However, this model assumes no return on investments. For many insurers, the extremely high competitiveness of today’s financial market means that they actually have a zero or negative operating profit, thus relying on investing to make up the shortfall and make a profit. The classical ruin model with risky investments has been considered by several authors, , , , , - (and the references contained therein) being but a few.
In this paper, we shall firstly consider loss distributions with infinite support, reproving a result originating in , and generalised in and , that if the volatility of the investments is of a certain magnitude, then ruin is inevitable. Our method of proof of this result is somewhat simpler and intuitively easier to understand. This method is also quite flexible, and we are able to prove generalisations of this result for more general risk processes, as well as when the investment is modelled by a certain Lévy process.
Additionally, we provide the transition density when the risk process is regarded as a diffusion, and conclude with some remarks on Lévy processes.
## 2 The main idea
Consider the following: Take the classical Cramér-Lundberg risk model, and then act on the capital position by a geometric Brownian motion that models the investment:
$$dX_t^{}=e^{\sigma B_t+\alpha t}dX_t$$
That is,
$$X_t^{}=e^{\sigma B_t+\alpha t}u+_0^te^{\sigma B_s+\alpha s}𝑑X_s$$
(2.1)
This is a version of risk process modified for investment considered in , and . The continuous interaction of the combined (standard) risk process and investment can be described in the manner of a semi-direct product:
$$(X_t,e^{\sigma B_t+\alpha t})(X_s,e^{\sigma B_s+\alpha s})=(X_t+X_se^{\sigma B_t+\alpha t},e^{\sigma B_{t+s}+\alpha (t+s)})$$
The group with this operation is better known as real hyperbolic space. Some knowledge of hyperbolic space will be needed for Brownian motion, which we will be considering in section 3. It is not a requisite for Theorem 1, although the idea of the action of investments on the risk process will quickly prove this result.
### 2.1 Certain ruin with risky investments
Consider the typical model for a share at time $`t`$, $`S_t`$, expressed as geometric Brownian motion with drift and diffusion parameters of $`a`$ and $`\sigma ^2`$, respectively. That is,
$$dS_t=aS_tdt+\sigma S_tdB_t$$
(2.2)
The solution of this S.D.E. is
$$S_t=S_0\mathrm{exp}\left((a\frac{\sigma ^2}{2})t+\sigma B_t\right)=S_0\mathrm{exp}\left(\alpha t+\sigma B_t\right)$$
setting $`\alpha =a\frac{\sigma ^2}{2}`$. The classical Cramér-Lundberg ruin model (1.1) is typically modified for investments by putting
$$X_t^{}=u+a_0^tX_s^{}𝑑s+\sigma _0^tX_s^{}𝑑B_s+ct\underset{k=1}{\overset{N_t}{}}Y_t$$
(2.3)
so that $`X_t^{}`$ describes the evolution of the capital of an insurer which is continuously invested in an asset which follows a geometric Brownian motion (independent of $`Y_t`$ and $`N_t`$) with parameters $`a`$ and $`\sigma `$. However, by using the model from (2.1):
$$X_t^{}=_0^te^{\sigma B_s+\alpha s}𝑑X_s$$
(2.4)
we view the geometric Brownian motion (share price) as acting as a dilation on the classical risk model. From this fact, our first result follows rather easily, proving the following:
###### Theorem 1.
Consider the classical Cramér-Lundberg risk model with investments as in (2.4). Assume further that the distribution of $`Y_1`$ does not have finite support, that is, $`P(Y_1>y)>0`$ for all $`y>0`$.
If $`\alpha <0`$ (or equivalently, $`\frac{2a}{\sigma ^2}<1`$), then ruin is certain.
###### Proof.
It is enough to prove that the capital position $`X_t^{}`$ is bounded, since the claim size can be large enough to ruin the company. Now consider the classical risk model. We have
$$X_t=u+ct\underset{k=1}{\overset{N_t}{}}Y_tu+ct$$
(2.5)
for all $`t0`$, since the $`Y_t`$ are positive distributions. Thus from our model in (2.4),
$`X_t^{}`$ $`=e^{\sigma B_t+\alpha t}u+{\displaystyle _0^t}e^{\sigma B_s+\alpha s}𝑑X_s`$
$`e^{\sigma B_t+\alpha t}u+{\displaystyle _0^t}e^{\sigma B_s+\alpha s}d(cs)`$
$`=e^{\sigma B_t+\alpha t}u+c{\displaystyle _0^t}e^{\sigma B_s+\alpha s}𝑑s`$
We are interested in what happens for $`t`$ large. We have
$$\mathrm{exp}\left(\sigma B_t+\alpha t\right)=\mathrm{exp}\left(t(\sigma B_t/t+\alpha )\right)$$
(2.6)
By the strong law of large numbers, $`B_t/t0`$ as $`t\mathrm{}`$, so our dilation term acts as $`e^{t\alpha }0`$ as $`t\mathrm{}`$. From this we deduce that $`e^{\sigma B_t+\alpha t}u0`$ as $`t0`$, and $`_0^te^{\sigma B_s+\alpha s}𝑑s`$ is bounded for all $`t>0`$. Therefore, the capital position is bounded, and the theorem follows.
Theorem 1 also holds for many generalisations of the Cramér-Lundberg ruin model. The first concerns varying premium rates (c.f. )
###### Corollary 1.
Suppose the premium rate in the risk model (1.1) is a bounded function, $`c_t`$. Then with the assumptions of Theorem 1, ruin is certain.
###### Proof.
This follows by putting $`c=sup_{t>0}c_t`$ in (2.5) above. ∎
Another concerns when $`N_t`$ is a counting process other than the Poisson process (see, for example, ).
###### Corollary 2.
Suppose $`N_t`$ is an arbitrary counting process in the risk model (1.1). Then with the assumptions of Theorem 1, ruin is certain.
We may also consider variations in the investment model where the Brownian motion is replaced by a more general Lévy process:
###### Corollary 3.
Suppose the investments in the risk model (2.4) are modelled by the exponential functional
$$e^{\sigma L_t+\alpha t}$$
where $`L_t`$ is Lévy process with mean 0. Then with the assumptions of Theorem 1, ruin is certain.
###### Proof.
This follows again by the strong law of large numbers since $`lim_t\mathrm{}L_t/t=0`$. ∎
Remark: These results also sheds some light on dividend constraints, since paying a dividend may be regarded as subtracting from the value of $`a`$ above, thus contributing to the overall probability of ruin. Additionally, setting $`\sigma =0`$ describes the risk model with a deterministic force of interest, $`a`$.
## 3 The diffusion limit of the probability of ruin
### 3.1 The diffusion model
The following characterisation of the probability of ruin was first introduced by Grandell in , who constructed a sequence of risk processes that converged weakly in the Skorohood topology to a diffusion process.
Put $`C_t=_{k=1}^{N_t}Y_t`$, which has mean $`\lambda \mu t`$ and variance $`\lambda mt`$. If the premium rate is set equal to $`\lambda \mu `$ then ruin will be certain. To avoid this, a safety loading $`\rho `$ is added as follows:
$$c=(1+\rho )\lambda \mu \rho =\frac{c\lambda \mu }{\lambda \mu }$$
(3.1)
By regarding the risk process as a diffusion, we may use the classical result of the hitting times of Brownian motion to give the well known “back of the envelope” calculation of the probability of ruin by considering the diffusion limit of the probability of ruin:
$$\mathrm{\Psi }_D=\mathrm{exp}\{2.\frac{\rho \mu u}{m}\}$$
(3.2)
We now considner the classical Cramér-Lundberg ruin model (considered as a diffusion), and then act on the capital position by a geometric Brownian motion that models the investment. These models follow shifted and dilated Brownian motion on the groups $`(,+)`$ and $`(^+,\times )`$, respectively, with their interaction again described by the semi-direct product from (2.1). That is, we will consider $`_{k=1}^{N_t}Y_t`$ to be a Brownian motion with a drift of $`\lambda \mu `$ and diffusion co-efficient $`\lambda m`$. This is somewhat at odds with , who argue that no financial asset may be correctly modelled on a continuous martingale. However, some rigour for using the pure diffusion model has been provided in , who give conditions for a weak convergence to a diffusion (ie, in the Skorohood topology as in ), so that we may consider the “unscaled version” of the diffusion limit of the probability of ruin with investments.
### 3.2 Real hyperbolic space
Real hyperbolic space may be defined in several ways. It is usually recognised as the Poincaré upper half-plane (c.f. , and , Ch III). For our purposes, we identify real hyperbolic space as the semi-direct product $`^+`$, where $``$ is the group of reals under addition, and $`^+`$ is the group of positive reals under multiplication. The semi-direct product $`^+`$ is the topological space $`\times ^+`$ with group multiplication given by
$$(x,y)(x^{},y^{})=(x+x^{}y,yy^{})$$
There are some slight technical details when performing analysis on this group. Firstly, it is a prime example of a non-unimodular group, ie, the left Haar measure is not equal to right Haar measure. More importantly, it is the only simple, simply connected Lie group whose Laplacian cannot be written as a sum of squares of its vector fields (c.f. , Thm. 4.1). Since the generator of Brownian motion is the Laplacian, this fact has implications for Brownian motion (c.f. , Cor. 4.4). However, this is overcome in and by considering a distinguished sub-Laplacian to generate a Brownian motion.
### 3.3 Brownian motion on hyperbolic space
Brownian motion on $`^+`$ was studied explicitly in and . The random variable considered in was
$$(_0^te^{B_s}𝑑W_s,e^{B_t})$$
In this was shown to be equivalent under the Dambis, Dubins-Schwarz theorem to the process
$$(W_{A_t},e^{B_t})$$
where $`(W_t)_{t0}`$ and $`(B_t)_{t0}`$ are standard Brownian motions, and
$$A_t=_0^te^{2B_s}𝑑s$$
Its characteristic function is
$`E(\mathrm{exp}(i\xi W_{A_t}+i\zeta e^{B_t}))`$ $`=E(\mathrm{exp}(i\zeta e^{B_t})(\mathrm{exp}(i\xi W_{A_t})|_t^B))`$
$`=E(\mathrm{exp}(i\zeta e^{B_t})\mathrm{exp}(\frac{1}{2}\xi ^2A_t))`$
This is then determined, and then inverted to obtain the density:
$$p_t(z,e^x)=\frac{e^{x^2/2t}}{\sqrt{2\pi t}}_0^{\mathrm{}}\frac{1}{\sqrt{2\pi y^2}}\mathrm{exp}\left(\frac{z^2}{y^2}\right)a_t(x,y)𝑑y$$
(3.3)
However, we wish to consider $`X_t`$ as a Brownian motion on $``$ starting from $`u`$, with mean $`\rho \lambda \mu t`$ and variance $`\lambda mt`$, being acted upon by a geometric Brownian motion on $`^+`$ with the parameters $`\alpha `$ and $`\sigma `$. That is, we are considering
$$(X_{A_{\sigma ^2t}^{(\alpha )}},e^{\sigma B_t+\alpha t})$$
(3.4)
where $`(X_t)_{t0}`$ is as in (1.2), $`(B_t)_{t0}`$ is a standard Brownian motion, and
$$A_{\sigma ^2t}^{(\alpha )}=_0^te^{2(\sigma B_s+\alpha s)}𝑑s$$
(3.5)
Equivalently from we have
$$(_0^te^{\sigma B_s+\alpha s}𝑑X_s,e^{\sigma B_t+\alpha t})$$
which are the random variables - considered separately - in . Although as remarked in that the density of $`_0^te^{\sigma B_s+\alpha s}𝑑X_s`$ cannot be put in a closed form, in the next section we shall derive the combined density of these processes.
The conditional distribution for $`A_{\sigma ^2t}^{(\alpha )}`$ in (3.5) comes from the functional studied in
$$A_t^{(\alpha )}=_0^te^{2(B_s+\alpha s)}𝑑s$$
and write $`A_t`$ when $`\alpha =0`$. This functional originates in Yor’s work with Bessel functions and its application to pricing Asian options. For a given $`B_t`$, if the density of $`A_t`$ is given by
$$P(A_tdu|B_t=x)=a_t(x,u)du$$
then
$$\frac{1}{\sqrt{2\pi t}}e^{\frac{x^2}{2t}}a_t(x,u)=\frac{1}{u}\mathrm{exp}((1+e^{2x})/2u)\mathrm{\Theta }_{\frac{1}{u}e^x}(t)$$
where
$$\mathrm{\Theta }_{\frac{1}{u}e^x}(t)=\frac{x}{u\sqrt{2\pi ^3t}}_0^{\mathrm{}}e^{\frac{y^2}{2t}}\mathrm{exp}(e^x\mathrm{cosh}(y)/u)\mathrm{sinh}(y)\mathrm{sin}(\pi y/t)𝑑y$$
### 3.4 Transition density of the diffusion model
Using the method described above, we are now in a position to obtain the transition density of the classical Cramér-Lundberg ruin model with investments, where both are considered to be diffusions.
###### Theorem 2.
The transition density of the risk process (1.1) with investments
$$(_0^te^{\sigma B_s+\alpha s}𝑑X_s,e^{\sigma B_t+\alpha t})$$
(3.6)
is given by
$$p_t(z,e^x)=\frac{e^{x^2/2\sigma ^2t}}{\sqrt{2\pi \sigma ^2t}}_0^{\mathrm{}}\frac{1}{\sqrt{2\pi \lambda \mu y^2}}\mathrm{exp}\left(\frac{(z\rho \lambda \mu tu)^2}{\lambda \mu y^2}\right)a_{\sigma ^2t}(x,y)𝑑y$$
(3.7)
###### Proof.
By the scalar invariance of $`B_t`$, and applying the Girsanov theorem, it is readily seen that the density of
$$_0^te^{2(\sigma B_s+\alpha s)}𝑑s=A_{\sigma ^2t}^{(\alpha )}$$
is given by
$$P(A_{\sigma ^2t}^{(\alpha )}du|\sigma B_t+\alpha t=x)=a_{\sigma ^2t}(x,u)du$$
The characteristic function for the random variable in (3.6) is
$`E[\mathrm{exp}(i\xi X_{A_{\sigma ^2t}^{(\alpha )}}`$ $`+i\zeta e^{(\sigma B_t+\alpha t)})]=E[\mathrm{exp}(i\zeta e^{(\sigma B_t+\alpha t)})(\mathrm{exp}(i\xi X_{A_{\sigma ^2t}^{(\alpha )}})|_t^B)]`$
$`=E\left[\mathrm{exp}(i\zeta e^{(\sigma B_t+\alpha t)})\mathrm{exp}((i\rho \lambda \mu t\xi +iu\xi \frac{1}{2}\lambda m\xi ^2)A_{\sigma ^2t}^{(\alpha )})\right]`$
$`=e^{i\rho \lambda \mu t\xi +iu\xi }{\displaystyle _{\mathrm{}}^{\mathrm{}}}\mathrm{exp}(i\zeta e^{(\sigma x+\alpha t)}){\displaystyle _0^{\mathrm{}}}\mathrm{exp}(\frac{1}{2}(\sqrt{\lambda m}\xi )^2)`$
$`\times a_{\sigma ^2t}(x,u)e^{x^2/(2\sigma ^2t)}dudx`$
This is merely a shifted and dilated version of the characteristic function of the Brownian motion in . So, mutatis mutandis, we invert the Fourier transform to obtain the density:
$$p_t(z,e^x)=\frac{e^{x^2/2\sigma ^2t}}{\sqrt{2\pi \sigma ^2t}}_0^{\mathrm{}}\frac{1}{\sqrt{2\pi \lambda \mu y^2}}\mathrm{exp}\left(\frac{(z\rho \lambda \mu tu)^2}{\lambda \mu y^2}\right)a_{\sigma ^2t}(x,y)𝑑y$$
as required.
Although the expression for $`p_t(z,e^x)`$ is far from tractable, we may use it to say the following:
###### Corollary 4.
Suppose $`p_t(z,e^x)`$ is the transition density for the risk process, then the probability that the company is ruined at time $`t`$ is given by
$$_{\mathrm{}}^0_{\mathrm{}}^{\mathrm{}}p_t(z,e^x)𝑑x𝑑z$$
Evaluating this integral will (more than likely) need to be done using numerical techniques. This expression is particularly important in actuarial practice, since it may only be necessary for the insurer to be solvent at certain times, that is, review dates.
Remark: The above calculations would appear to be valid when the generator is that of the classical risk process. That is, the process on $``$ is a (discontinuous) semimartingale, rather than just a Brownian motion. This allow us to compute the transition density in the case of a compound Poisson process, rather than just the diffusion model, where the combined density would be an expression similar to (3.7). More generally, we could consider the risk process as a Lévy process. These models were discussed in detail in , which we refer the reader to for many explicit examples. The critical step is a “Dambis, Dubins-Schwarz”-type theorem for discontinuous semimartingales, which is outside the scope of this paper.
Remark: Similarly, the transition density when the investment model is other than geometric Brownian motion requires significant modifications to Yor’s work in generalising the density of the functional $`A_t`$ when the investment model is a Lévy process rather than geometric Brownian motion. That said, a functional of the form $`A_t^{}=_0^te^{2L_s}𝑑s`$ where $`L_t`$ is a Lévy process does provide us with a cádlág modification to the risk model in (3.6). |
warning/0506/astro-ph0506458.html | ar5iv | text | # The age of the Galactic thin disk from Th/Eu nucleocosmochronology
## 1 Introduction
The age of the Galactic thin disk<sup>1</sup><sup>1</sup>1All references to the *Galactic disk* must be regarded, in this work, as references to the *thin* disk, unless otherwise specified. is an important constraint for Galactic formation models. It is usually estimated by dating the oldest open clusters with isochrones or white dwarfs with cooling sequences. These methods are strongly dependent on stellar evolution models and on numerous physical parameters known at different levels of uncertainty. Nucleocosmochronology is only weakly dependent on main sequence stellar evolution models, allowing for a nearly independent crosscheck of other techniques.
We were the first to determine an age for the Galactic disk from Th/Eu nucleocosmochronology – del Peloso et al. 2005b (Paper I) and del Peloso et al. 2005a (Paper II), based on the PhD thesis of one of us (del Peloso 2003). This work, the last part in a series of three articles, aims at reducing the uncertainty in this determination by expanding the stellar sample from Papers I and II with objects observed only after the publication of del Peloso (2003) and deriving a new age from this extended sample. With this intent, we determined \[Th/Eu\] abundance ratios for a sample of Galactic disk stars and employed Galactic disk chemical evolution models developed by us in the chronological analysis.<sup>2</sup><sup>2</sup>2In this paper we obey the customary spectroscopic notation: abundance ratio $`\text{[A/B]}\mathrm{log}_{10}(N_\mathrm{A}/N_\mathrm{B})_{\mathrm{star}}\mathrm{log}_{10}(N_\mathrm{A}/N_\mathrm{B})_{\text{}}`$, where $`N_\mathrm{A}`$ and $`N_\mathrm{B}`$ are the abundances of elements A and B, respectively, in atoms cm<sup>-3</sup>.
## 2 Sample selection, observations and data reduction
Sample selection, observations and data reduction were carried out following exactly the same procedures described in detail in Paper I. In what follows we provide a brief overview of these topics.
The stellar sample of this work is composed of seven F8–G5 dwarfs and subgiants (Table 1). As the objective of this work is the determination of the age of the Galactic *disk*, we performed a kinematic analysis of the objects to help ensure that they do not belong to the Galactic halo. We calculated the objects’ $`U`$, $`V`$, and $`W`$ spatial velocity components (Table 2) and constructed a $`V`$ vs. \[Fe/H\] diagram (Figure 1) using metallicities from the literature (Table 3). According to Schuster et al. (1993), objects located above the displayed cut-off line belong to the Galactic halo, whereas those located below the line belong to the halo. It can be seen that all sample stars are located far above the cut-off line. For a star to cross the line, the literature metallicities would have to have been *overestimated* by at least 1.2 dex, which is very unlikely. After having determined our own metallicities (Table 4), we confirmed that results from the literature agree with our values to 0.1 dex.
High resolution, high signal-to-noise ratio spectra were obtained for all objects with the Fiber-fed Extended Range Optical Spectrograph (FEROS; Kaufer et al. 1999) fed by the 1.52 m European Southern Observatory (ESO) telescope, in the ESO-Observatório Nacional, Brazil, agreement (March and August 2001). Spectra were also obtained with a coudé spectrograph fed by the 1.60 m telescope of the Observatório do Pico dos Dias (OPD), LNA/MCT, Brazil (May and October 2000; May, August, and October 2002), and with the Coudé Échelle Spectrometer (CES) fiber-fed by ESO’s 3.60 m telescope (August 2003). FEROS spectra are reduced automatically during observation by a script executed using the European Southern Observatory Munich Image Data Analysis System (ESO-MIDAS) immediately after the CCD read out. OPD and CES spectra were reduced by us using the Image Reduction and Analysis Facility (IRAF<sup>3</sup><sup>3</sup>3IRAF is distributed by the National Optical Astronomy Observatories, which are operated by the Association of Universities for Research in Astronomy, Inc., under cooperative agreement with the National Science Foundation.), following the usual steps of bias, scattered light and flat field corrections, and extraction.
## 3 Atmospheric parameters
A set of homogeneous, self-consistent atmospheric parameters was determined for the sample stars, following the procedure described in detail in Paper I. Effective temperatures were derived from photometric calibrations (Table 3) and H$`\alpha `$ profile fitting; surface gravities were obtained from $`T_{\mathrm{eff}}`$, stellar masses and luminosities; microturbulence velocities and metallicities were obtained from detailed, differential spectroscopic analysis, relative to the Sun, using EWs of Fe i and Fe ii lines. The final, adopted values of all atmospheric parameters are presented in Table 4.
## 4 Abundances of contaminating elements
Eu and Th abundances were determined by spectral synthesis of the Eu ii line at 4129.72 Å and of the Th ii line at 4019.13 Å, respectively. In the synthesis calculations, abundances of the elements that contaminate these spectral regions (Ti, V, Cr, Mn, Co, Ni, Ce, Nd, and Sm) were kept fixed in the values determined using EWs and the atmospheric parameters obtained by us; see Paper I for a full description of the employed method. Table 5 presents a sample of the EW data. Its complete content, composed of the EWs of all measured lines, for the Sun and all sample stars, is only available in electronic form at the CDS. Column 1 lists the central wavelength (in angstroms), Col. 2 gives the element symbol and degree of ionization, Col. 3 gives the excitation potential of the lower level of the electronic transition (in eV), Col. 4 presents the solar $`\mathrm{log}gf`$ derived by us, and the subsequent columns present the EWs, in mÅ, for the Sun and the other stars, from HD 1461 to HD 210 918.
Abundance results are presented in Table 6. No detailed uncertainty assessment was carried out. We have rather adopted an average of the uncertainties presented in Paper I: 0.08 dex for Mn, 0.09 dex for Fe, Ti, and Co, 0.10 dex for V, Cr, Ce, and Nd, and 0.11 dex for Ni and Sm. These values are used as error bars in Fig. 2, which shows the abundance patterns for all elements. Note that the abundances for the sample of this work (filled squares) agree very well with those from Paper I (open squares).
## 5 Eu and Th abundances
Our determination of Eu abundances by spectral synthesis used FEROS spectra. As the abundances used for age determination in Paper I were obtained with the CES coupled to the Coudé Échelle Spectrograph (CAT), we converted our results to the CAT system, using the equation derived in Sect. 5.1.1 of Paper I (see their Fig. 15): $`[\mathrm{Eu}/\mathrm{H}]_{\mathrm{CAT}+\mathrm{CES}}=0.00462+0.96480[\mathrm{Eu}/\mathrm{H}]_{\mathrm{FEROS}}`$. Table 7 presents the \[Eu/H\], \[Th/H\], and \[Th/Eu\] abundance ratios for our sample. As uncertainty, we adopted an average of the values presented in Paper I: 0.04 dex for \[Eu/H\], 0.11 dex for \[Th/H\], and 0.08 dex for \[Th/Eu\].
In Fig. 3, we compare our \[Eu/H\] abundance ratios to those from Woolf et al. (1995, WTL95), Koch & Edvardsson (2002, KE02), and Paper I. These are the best works available with Eu abundances for Galactic disk stars, in terms of care of analysis and sample size. Although two of our objects (HD 1461 and HD 157 089) exhibit considerable discrepancy with Paper I, our results present behavior and dispersion similar to WTL95/KE02.
Th spectral synthesis employed CES spectra fed by the 3.60 m ESO telescope, whereas the abundances in Paper I were obtained with the CES fed by the CAT. Our results were converted to the CAT system, using the equation derived in Sects. 5.2.1 of Paper I (see their Fig. 20): $`[\mathrm{Th}/\mathrm{H}]_{\mathrm{CAT}+\mathrm{CES}}=0.01500+0.86000[\mathrm{Th}/\mathrm{H}]_{\mathrm{CAT}+3.60\mathrm{m}}`$.
In Fig. 4, we compare our \[Eu/H\] abundance ratios to those from Morell et al. (1992, MKB92) and Paper I. These are the best works available with Th abundances for Galactic disk stars, in terms of care of analysis and sample size. Our results are in good accord with them.
## 6 Chronological analysis
In order to estimate the age of the Galactic disk, we compared the stellar \[Th/Eu\] abundance ratios with curves calculated using a Galactic disk chemical evolution (GDCE) model. In this model, developed by us based on Pagel & Tautvaišienė (1995), it was assumed that the so-called “universality of the r-process abundances” is valid for second and third r-process peaks and can be extended to the actinides. Such extension may not be legitimate, as two ultra-metal-poor stars – CS 31 082-001 (Cayrel et al. 2001 and Hill et al. 2002) and CS 30 306-132 (Honda et al. 2004) – have been recently shown to have Th/Eu abundance ratios much higher than expected for their age. This could indicate that they have been formed from matter enriched in actinides, relatively to second r-process peak elements. However, it is not yet clear if CS 31 082-001 and CS 30 306-132 are merely chemically peculiar objects or if their discrepancies could be present in other yet unobserved stars. A detailed description of the GDCE model was presented in Paper II.
The abundances for our sample were determined in exactly the same way as those from Paper I, as the objective of this work is to expand the sample used for chronological analysis. Accordingly, we merged our abundance data with those from Paper I, resulting in a set of 28 objects. The theoretical model curves and the observed abundance data are presented in Fig. 5.
The age of the curve that best fits the observed data was computed by minimizing the total deviation between the curves and the data. 26 of the 28 objects in the merged sample were actually used in the determination; two objects from Paper I were disregarded because they are too metal-rich, falling out of the interval where the curves are defined. Considering that in Paper II 19 objects were actually employed in the analysis, this work accomplishes a 37% increase in sample size. An uncertainty for the age, related solely to the uncertainties in the abundances, was computed through a Monte Carlo simulation. It must be noted that this is only an assessment of the *internal* uncertainty of the analysis, and does not take into consideration the uncertainties of the GDCE model itself, which are very difficult to estimate. The uncertainty related to the model could very well be the main source of age uncertainty. These procedures are fully described in Paper II. The final value obtained using the merged abundance data set is $`(8.8\pm 1.8)\mathrm{Gyr}`$. This result, 0.6 Gyr larger and with an uncertainty 0.1 Gyr lower than that of Paper II, matches very well the estimate obtained from literature data and the GDCE model ($`8.7_{4.1}^{+5.8}\mathrm{Gyr}`$). These estimates were combined using the maximum likelihood method, assuming that each one follows a Gaussian probability distribution, which results in a weighted average using the reciprocal of the square uncertainties as weights. The final, adopted Galactic disk age is
$$T_\mathrm{G}=(8.8\pm 1.7)\text{Gyr.}$$
## 7 Conclusions
We determined \[Th/Eu\] abundance ratios for a sample of seven Galactic disk F8–G5 dwarfs and subgiants. The analysis was carried out in exactly the same way as that of Paper I, so that we could merge the data, resulting in a totally homogeneous extended sample of 28 objects; 26 of these were actually used in the nucleocosmochronological analysis. A GDCE model, developed in Paper II, was used in conjunction with the stellar abundance data to compute an age for the Galactic disk: $`(8.8\pm 1.8)\mathrm{Gyr}`$. In Paper II, Th/Eu production and solar abundance ratio data taken from the literature were analyzed with our GDCE model, yielding $`8.7_{4.1}^{+5.8}\mathrm{Gyr}`$. These two results were combined using the maximum likelihood method, resulting in $`\text{FINAL }T_\mathrm{G}=(8.8\pm 1.7)\text{Gyr}`$.
The inclusion of seven more stars in the abundance data base had two main consequences: it increased the age obtained from our stellar data by 0.6 Gyr, rendering it more compatible with the age determined from literature data, and it decreased the uncertainty of the final, adopted age by 0.1 Gyr, i.e., a 6% reduction. Our result remains compatible with the most up-to-date white dwarf ages derived from cooling sequence calculations, which indicate a low age ($`10\mathrm{Gyr}`$) for the disk (Oswalt et al. 1995; Bergeron et al. 1997; Leggett et al. 1998; Knox et al. 1999; Hansen et al. 2002). Considering that the age of the oldest halo globular clusters are currently estimated at $`(13.5\pm 0.7)\mathrm{Gyr}`$ (Pont et al. 1998; Jimenez 1999; Gratton et al. 2003; Krauss & Chaboyer 2003), an hiatus of $`(4.7\pm 1.8)\mathrm{Gyr}`$ between the formation of halo and disk must be taken into consideration in future Galactic formation models.
###### Acknowledgements.
The authors wish to thank the staff of the Observatório do Pico dos Dias, LNA/MCT, Brazil and of the European Southern Observatory, La Silla, Chile. We thank R. de la Reza for his contributions to this work. The suggestions of Dr. N. Christlieb, the referee, were greatly appreciated. EFP acknowledges financial support from CAPES/PROAP, FAPERJ/FP (grant E-26/150.567/2003), and CNPq/DTI (grant 382814/2004-5). LS thanks the CNPq, Brazilian Agency, for the financial support 453529.0.1 and for the grants 301376/86-7 and 304134-2003.1. GFPM acknowledges financial support from CNPq/Conteúdos Digitais, CNPq/Institutos do Milênio/MEGALIT, FINEP/PRONEX (grant 41.96.0908.00), FAPESP/Temáticos (grant 00/06769-4), and FAPERJ/APQ1 (grant E-26/170.687/2004). |
warning/0506/gr-qc0506001.html | ar5iv | text | # Stability of phantom wormholes
## I Introduction
It is now generally accepted that the Universe is undergoing an accelerated phase of expansion Riess2 ; Perlmutter ; Bennet ; Hinshaw , where the scale factor obeys $`\ddot{a}>0`$. This cosmic acceleration is one of the most challenging current problems in cosmology. Several candidates, responsible for this expansion, have been proposed in the literature, namely, dark energy models, generalizations of the Chaplygin gas, modified gravity and scalar-tensor theories, tachyon scalar fields and braneworld models, amongst others. The dark energy models are parametrized by an equation of state given by $`\omega =p/\rho `$, where $`p`$ is the spatially homogeneous pressure and $`\rho `$ is the dark energy density. For the cosmic expansion, a value of $`\omega <1/3`$ is required, as dictated by the Friedman equation $`\ddot{a}/a=4\pi (p+\rho /3)`$. A specific exotic form of dark energy, denoted phantom energy, has also been proposed, possessing the peculiar property of $`\omega <1`$. This parameter range is not excluded by observation, and possesses peculiar properties, such as the violation of the null energy condition and an infinitely increasing energy density, resulting in a Big Rip, at which point the Universe blows up in a finite time Weinberg . However, recent fits to supernovae, CMB and weak gravitational lensing data indicate that an evolving equation of state $`\omega `$ crossing the phantom divide $`1`$, is mildly favored, and several models have been proposed in the literature Zhang2 ; Zhang3 ; Periv1 ; Wei-Cai ; Li-Feng ; Stef ; Periv2 ; Feng ; Vikman ; Tsujikawa ; Sami . In particular, models considering a redshift dependent equation of state, $`\omega (z)`$, provide significantly ameliorated fits to the most recent and reliable SN Ia supernovae Gold dataset Riess3 .
As the phantom energy equation of state, $`p=\omega \rho `$ with $`\omega <1`$, violates the null energy condition, $`p+\rho <0`$, the fundamental ingredient to sustain traversable wormhole Morris ; mty ; Visser , one now has at hand a possible source for these exotic spacetimes. In fact, this possibility has recently been explored Sushkov ; Lobo-phantom , and it was shown that traversable wormholes can be theoretically supported by phantom energy. However, a subtlety needs to be pointed out, as emphasized in Refs. Sushkov ; Lobo-phantom . The notion of phantom energy is that of a homogeneously distributed fluid. When extended to inhomogeneous spherically symmetric spacetimes, the pressure appearing in the equation of state is now a radial pressure, and the transverse pressure is then determined via the field equations. In this context, it is interesting to note that wormhole solutions with an isotropic pressure were found in Lobo-phantom , although these geometries are not asymptotically flat. Sushkov, in Ref. Sushkov , found wormhole geometries by considering specific choices for the distribution of the energy density, and in Lobo-phantom , a complementary approach was traced out, by imposing appropriate choices for the form function and/or the redshift function, and the stress-energy tensor components were consequently determined. In Ref. Lobo-phantom it was also shown, using the “volume integral quantifier” VKD1 ; VKD2 , that these geometries can be theoretically constructed with infinitesimal amounts of averaged null energy condition violating phantom energy.
It is also of a fundamental importance to investigate the stability of these phantom wormhole geometries (It is also interesting to note that a stability analysis of a specific class of traversable wormholes was carried out in Ref. Armen , in a rather different context). As in Ref. Lobo-phantom , we shall model these spacetimes by matching an interior traversable wormhole geometry with an exterior Schwarzschild vacuum solution at a junction interface Lobo-CQG ; LLQ ; Lobo ; LL-PRD . In this work, we analyze the stability of these phantom wormholes to linearized perturbations around static solutions. Work along these lines was done by considering thin-shell Schwarzschild wormholes, using the cut-and-paste technique Poisson . It was later shown that the inclusion of a charge Eiroa and of a cosmological constant LC-CQG significantly increases the stable equilibrium configurations found in Ref. Poisson . The advantage of this analysis resides in using a parametrization of the stability of equilibrium, so that there is no need to specify a surface equation of state. Note that the stability analysis of these thin-shell wormholes to linearized spherically symmetric perturbations about static equilibrium solutions was carried out by assuming that the shells remain transparent under perturbation Ishak . This amounts to considering specific spacetimes that do not contribute with the momentum flux term in the conservation identity, which provides the conservation law for the surface stress-energy tensor. The inclusion of this term, corresponding to the discontinuity of the momentum impinging on the shell, severely complicates the analysis. However, we shall follow the approach of Ishak and Lake Ishak , with the respective inclusion of the momentum flux term, and deduce a master equation responsible for dictating the stability equilibrium configurations for the specific phantom wormhole geometries found in Ref. Lobo-phantom . We shall separate the cases of a positive and a negative surface energy density, and find that the stability may be significantly increased by varying the wormhole throat.
This paper is outlined in the following manner. In Section II, we present solutions of a phantom energy traversable wormhole. In Section III, we outline a general linearized stability analysis procedure, and deduce a master equation dictating stable equilibrium configurations. We then apply this analysis to phantom wormhole geometries and determine their respective stability regions. Finally in Section IV, we conclude.
## II Phantom energy traversable wormholes
### II.1 Field equations
The interior wormhole spacetime is given by the following metric Morris
$$ds^2=e^{2\mathrm{\Phi }(r)}dt^2+\frac{dr^2}{1b(r)/r}+r^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2),$$
(1)
where $`\mathrm{\Phi }(r)`$ and $`b(r)`$ are arbitrary functions of the radial coordinate, $`r`$, denoted as the redshift function and the form function, respectively Morris . The wormhole throat is located at $`b(r_0)=r=r_0`$. For the wormhole to be traversable, one must demand that there are no horizons present, which are identified as the surfaces with $`e^{2\mathrm{\Phi }}0`$, so that $`\mathrm{\Phi }(r)`$ must be finite everywhere. The condition $`1b/r>0`$ is also imposed. The stress-energy tensor components are given by (with $`c=G=1`$)
$`\rho (r)`$ $`=`$ $`{\displaystyle \frac{1}{8\pi }}{\displaystyle \frac{b^{}}{r^2}},`$ (2)
$`p_r(r)`$ $`=`$ $`{\displaystyle \frac{1}{8\pi }}\left[{\displaystyle \frac{b}{r^3}}+2\left(1{\displaystyle \frac{b}{r}}\right){\displaystyle \frac{\mathrm{\Phi }^{}}{r}}\right],`$ (3)
$`p_t(r)`$ $`=`$ $`{\displaystyle \frac{1}{8\pi }}(1{\displaystyle \frac{b}{r}})[\mathrm{\Phi }^{\prime \prime }+(\mathrm{\Phi }^{})^2{\displaystyle \frac{b^{}rb}{2r(rb)}}\mathrm{\Phi }^{}`$ (4)
$`{\displaystyle \frac{b^{}rb}{2r^2(rb)}}+{\displaystyle \frac{\mathrm{\Phi }^{}}{r}}],`$
where $`\rho (r)`$ is the energy density; $`p_r(r)`$ the radial pressure; and $`p_t(r)`$ the transverse pressure. The conservation of the stress-energy tensor, $`T^{\mu \nu }{}_{;\nu }{}^{}=0`$, provides us with the following relationship
$$p_r^{}=\frac{2}{r}(p_tp_r)(\rho +p_r)\mathrm{\Phi }^{}.$$
(5)
A fundamental ingredient of traversable wormholes and phantom energy is the violation of the null energy condition (NEC), which is defined as $`T_{\mu \nu }k^\mu k^\nu 0`$, where $`k^\mu `$ is any null vector and the $`T_{\mu \nu }`$ the stress-energy tensor. Note that for phantom energy, governed by the equation of state $`\omega =p/\rho `$ with $`\omega <1`$, one readily verifies that the NEC is violated, i.e., $`p+\rho <0`$. For wormhole spacetimes, consider an orthonormal reference frame with $`k^{\widehat{\mu }}=(1,\pm 1,0,0)`$, so that we have
$$T_{\widehat{\mu }\widehat{\nu }}k^{\widehat{\mu }}k^{\widehat{\nu }}=\frac{1}{8\pi }\left[\frac{b^{}rb}{r^3}+2\left(1\frac{b}{r}\right)\frac{\mathrm{\Phi }^{}}{r}\right].$$
(6)
Thus, using the flaring out condition of the throat, $`(bb^{}r)/2b^2>0`$ Morris ; Visser , and considering the finite character of $`\mathrm{\Phi }(r)`$, we verify that evaluated at the throat the NEC is violated, i.e., $`T_{\widehat{\mu }\widehat{\nu }}k^{\widehat{\mu }}k^{\widehat{\nu }}<0`$. Matter that violates the NEC is denoted as exotic matter.
Note that the notion of phantom energy is that of a homogeneously distributed cosmic fluid. However, as emphasized in Sushkov ; Lobo-phantom , it may be extended to inhomogeneous spherically symmetric spacetimes by regarding that the pressure in the equation of state $`p=\omega \rho `$ is now a radial pressure $`p_r`$. The transverse pressure $`p_t`$ may then be determined from the field equation, in particular, from Eq. (4). Thus, to find phantom energy traversable wormhole spacetimes, we use the equation of state $`p_r=\omega \rho `$ with $`\omega <1`$, representing phantom energy, and thus deduce the following relationship
$$\mathrm{\Phi }^{}(r)=\frac{b+\omega rb^{}}{2r^2\left(1b/r\right)},$$
(7)
by taking into account Eq. (2) and Eq. (3).
To model a traversable wormhole, one now considers appropriate choices for $`b(r)`$ and/or $`\mathrm{\Phi }(r)`$. Note that this is necessary as we only have four equations, namely, Eqs. (2)-(4), and Eq. (7), with five unknown functions of $`r`$, i.e., $`\rho (r)`$, $`p_r(r)`$, $`p_t(r)`$, $`b(r)`$ and $`\mathrm{\Phi }(r)`$. We shall only consider form functions of the type $`b^{}(r)>0`$, as in cosmology the phantom energy density is considered positive. Now, using the flaring out condition evaluated at the throat Morris ; Visser , we also have the condition $`b^{}(r_0)<1`$.
One may construct asymptotically flat spacetimes, in which $`b(r)/r0`$ and $`\mathrm{\Phi }0`$ as $`r\mathrm{}`$. However, one may also consider solutions with a cut-off of the stress-energy, by matching the interior solution to an exterior vacuum spacetime, at a junction interface, $`a`$. For simplicity, in this paper, we shall consider that the exterior spacetime is the Schwarzschild solution, so that the matching occurs at a junction interface, $`r=a`$, situated outside the event horizon, i.e., $`a>r_b=2M`$, in order to avoid a black hole solution.
### II.2 Specific phantom wormhole models
The physical properties and characteristics of specific phantom energy traversable wormhole models were analyzed in Ref. Lobo-phantom , by considering asymptotically flat spacetimes and by imposing an isotropic pressure. Using the “volume integral quantifier” it was found that it is theoretically possible to construct these geometries with vanishing amounts of averaged null energy condition violating phantom energy. Specific wormhole dimensions and the traversal velocity and time were also deduced from the traversability conditions for a particular wormhole geometry. We shall briefly summarize two specific phantom wormhole models, found in Ref. Lobo-phantom , and for which we shall further analyze the respective stable equilibrium configurations.
#### Asymptotically flat spacetimes
To construct an asymptotically flat wormhole solution Lobo-phantom , consider $`\mathrm{\Phi }(r)=\mathrm{const}`$. Thus, from Eq. (7) one obtains
$$b(r)=r_0(r/r_0)^{1/\omega },$$
(8)
so that $`b(r)/r=(r_0/r)^{(1+\omega )/\omega }0`$ for $`r\mathrm{}`$. We also verify that $`b^{}(r)=(1/\omega )(r/r_0)^{(1+\omega )/\omega }`$, so that at the throat the condition $`b^{}(r_0)=1/|\omega |<1`$ is satisfied.
The stress-energy tensor components are given by
$`p_r(r)`$ $`=`$ $`\omega \rho (r)={\displaystyle \frac{1}{8\pi r_0^2}}\left({\displaystyle \frac{r_0}{r}}\right)^{3+\frac{1}{\omega }},`$ (9)
$`p_t(r)`$ $`=`$ $`{\displaystyle \frac{1}{16\pi r_0^2}}\left({\displaystyle \frac{1+\omega }{\omega }}\right)\left({\displaystyle \frac{r_0}{r}}\right)^{3+\frac{1}{\omega }}.`$ (10)
Thus, determining the parameter $`\omega `$ from observational cosmology, assuming the existence of phantom energy, one may theoretically construct traversable phantom wormholes by considering the above-mentioned form function and a constant redshift function.
#### Isotropic pressure, $`p_r=p_t=p`$
It was found that considering an isotropic pressure, $`p_r=p_t=p`$, for $`\mathrm{\Phi }(r)`$ to be finite one cannot construct asymptotically flat traversable wormholes Lobo-phantom . By taking into account the form function given by $`b(r)=r_0(r/r_0)^\alpha `$, with $`0<\alpha <1`$, and using Eq. (5) and Eq. (2), one finds that the redshift function is given by
$$\mathrm{\Phi }(r)=\left(\frac{3\omega +1}{1+\omega }\right)\mathrm{ln}\left(\frac{r}{r_0}\right),$$
(11)
where the relationship $`\alpha =1/\omega `$ is imposed (see Ref. Lobo-phantom for details). The stress-energy tensor components are provided by
$`p(r)`$ $`=`$ $`\omega \rho (r)={\displaystyle \frac{1}{8\pi r_0^2}}\left({\displaystyle \frac{r_0}{r}}\right)^{3+\frac{1}{\omega }}.`$ (12)
As noted above, the spacetime is not asymptotically flat. Nevertheless, one may match the interior wormhole solution to an exterior vacuum spacetime at a finite junction surface.
## III Stability analysis
### III.1 Junction conditions
We shall model specific phantom wormholes by matching an interior traversal wormhole geometry, satisfying the equation of state $`p_r=\omega \rho `$ with $`\omega <1`$, with an exterior Schwarzschild solution at a junction interface $`\mathrm{\Sigma }`$, situated outside the event horizon, $`a>r_b=2M`$.
Using the Darmois-Israel formalism Darmois ; Israel , the surface stress-energy tensor, $`S^i_j`$, at the junction interface $`\mathrm{\Sigma }`$ is provided by the Lanczos equations
$$S_j^i=\frac{1}{8\pi }(\kappa _j^i\delta _j^i\kappa _k^k),$$
(13)
where $`\kappa _{ij}`$ is the discontinuity of the extrinsic curvatures across the surface $`\mathrm{\Sigma }`$, i.e., $`\kappa _{ij}=K_{ij}^+K_{ij}^{}`$. The extrinsic curvature is defined as $`K_{ij}=n_{\mu ;\nu }e_{(i)}^\mu e_{(j)}^\nu `$, where $`n^\mu `$ is the unit normal $`4`$vector to $`\mathrm{\Sigma }`$, and $`e_{(i)}^\mu `$ are the components of the holonomic basis vectors tangent to $`\mathrm{\Sigma }`$.
Taking into account the wormhole spacetime metric (1) and the Schwarzschild solution, the non-trivial components of the extrinsic curvature are given by
$`K_\tau ^{\tau +}`$ $`=`$ $`{\displaystyle \frac{\frac{M}{a^2}+\ddot{a}}{\sqrt{1\frac{2M}{a}+\dot{a}^2}}},`$ (14)
$`K_\tau ^\tau `$ $`=`$ $`{\displaystyle \frac{\mathrm{\Phi }^{}\left(1\frac{b}{a}+\dot{a}^2\right)+\ddot{a}\frac{\dot{a}^2(bb^{}a)}{2a(ab)}}{\sqrt{1\frac{b(a)}{a}+\dot{a}^2}}},`$ (15)
and
$`K_\theta ^{\theta +}`$ $`=`$ $`{\displaystyle \frac{1}{a}}\sqrt{1{\displaystyle \frac{2M}{a}}+\dot{a}^2},`$ (16)
$`K_\theta ^\theta `$ $`=`$ $`{\displaystyle \frac{1}{a}}\sqrt{1{\displaystyle \frac{b(a)}{a}}+\dot{a}^2}.`$ (17)
The Lanczos equation, Eq. (13), then provide us with the following expressions for the surface stresses
$`\sigma `$ $`=`$ $`{\displaystyle \frac{1}{4\pi a}}\left(\sqrt{1{\displaystyle \frac{2M}{a}}+\dot{a}^2}\sqrt{1{\displaystyle \frac{b(a)}{a}}+\dot{a}^2}\right),`$ (18)
$`𝒫`$ $`=`$ $`{\displaystyle \frac{1}{8\pi a}}[{\displaystyle \frac{1\frac{M}{a}+\dot{a}^2+a\ddot{a}}{\sqrt{1\frac{2M}{a}+\dot{a}^2}}}`$ (19)
$`{\displaystyle \frac{(1+a\mathrm{\Phi }^{})\left(1\frac{b}{a}+\dot{a}^2\right)+a\ddot{a}\frac{\dot{a}^2(bb^{}a)}{2(ab)}}{\sqrt{1\frac{b(a)}{a}+\dot{a}^2}}}],`$
where $`\sigma `$ and $`𝒫`$ are the surface energy density and the tangential surface pressure, respectively.
We shall make use of the conservation identity, which is obtained from the second contracted Gauss-Kodazzi equation or the “ADM” constraint $`G_{\mu \nu }e_{(i)}^\mu n^\nu =K_{i|j}^jK,_i`$ with the Lanczos equations, and is given by
$`S_{j|i}^i=\left[T_{\mu \nu }e_{(j)}^\mu n^\nu \right]_{}^+.`$ (20)
The momentum flux term in the right hand side corresponds to the net discontinuity in the momentum which impinges on the shell.
Using $`S_{\tau |i}^i=\left[\dot{\sigma }+2\dot{a}(\sigma +𝒫)/a\right]`$, Eq. (20) provides us with
$$\sigma ^{}=\frac{2}{a}(\sigma +𝒫)+\mathrm{\Xi },$$
(21)
where $`\mathrm{\Xi }`$, defined for notational convenience, is given by
$`\mathrm{\Xi }={\displaystyle \frac{1}{4\pi a^2}}\left[{\displaystyle \frac{b^{}ab}{2a\left(1\frac{b}{a}\right)}}+a\mathrm{\Phi }^{}\right]\sqrt{1{\displaystyle \frac{b}{a}}+\dot{a}^2}.`$ (22)
For self-completeness, we shall also include the $`\sigma +𝒫`$ term, which is given by
$`\sigma +𝒫`$ $`=`$ $`{\displaystyle \frac{1}{8\pi a}}[{\displaystyle \frac{(1a\mathrm{\Phi }^{})\left(1\frac{b}{a}+\dot{a}^2\right)a\ddot{a}+\frac{\dot{a}^2(bb^{}a)}{2(ab)}}{\sqrt{1\frac{b(a)}{a}+\dot{a}^2}}}`$ (23)
$`{\displaystyle \frac{1\frac{3M}{a}+\dot{a}^2a\ddot{a}}{\sqrt{1\frac{2M}{a}+\dot{a}^2}}}].`$
Thus, taking into account Eq. (23), and the definition of $`\mathrm{\Xi }`$, we verify that Eq. (21) finally takes the form
$`\sigma ^{}`$ $`=`$ $`{\displaystyle \frac{1}{4\pi a^2}}({\displaystyle \frac{1\frac{3M}{a}+\dot{a}^2a\ddot{a}}{\sqrt{1\frac{2M}{a}+\dot{a}^2}}}`$ (24)
$`{\displaystyle \frac{1\frac{3b}{2a}+\frac{b^{}}{2}+\dot{a}^2a\ddot{a}}{\sqrt{1\frac{b}{a}+\dot{a}^2}}}),`$
which, evaluated at a static solution $`a_0`$, shall play a fundamental role in determining the stability regions. Note that Eq. (24) can also be deduced by taking the radial derivative of the surface energy density, Eq. (18).
### III.2 Equation of motion
Rearranging Eq. (18) into the form
$$\sqrt{1\frac{2M}{a}+\dot{a}^2}=\sqrt{1\frac{b(a)}{a}+\dot{a}^2}4\pi \sigma a,$$
(25)
we deduce the thin shell’s equation of motion, i.e.,
$$\dot{a}^2+V(a)=0,$$
(26)
with the potential given by
$$V(a)=1+\frac{2Mb(a)}{m_s^2}\left[\frac{m_s}{2a}+\frac{\left(M+\frac{b(a)}{2}\right)}{m_s}\right]^2,$$
(27)
where $`m_s=4\pi \sigma a^2`$ is the surface mass of the thin shell. However, for computational purposes and notational convenience, we define the following factors
$`F(a)`$ $`=`$ $`1{\displaystyle \frac{b(a)/2+M}{a}},`$ (28)
$`G(a)`$ $`=`$ $`{\displaystyle \frac{Mb(a)/2}{a}},`$ (29)
so that the potential $`V(a)`$ takes the form
$$V(a)=F(a)\left(\frac{m_s}{2a}\right)^2\left(\frac{aG}{m_s}\right)^2.$$
(30)
Linearizing around a stable solution situated at $`a_0`$, we consider a Taylor expansion of $`V(a)`$ around $`a_0`$ to second order, given by
$`V(a)`$ $`=`$ $`V(a_0)+V^{}(a_0)(aa_0)`$ (31)
$`+{\displaystyle \frac{1}{2}}V^{\prime \prime }(a_0)(aa_0)^2+O[(aa_0)^3].`$
The first and second derivatives of $`V(a)`$ are given by
$`V^{}(a)`$ $`=`$ $`F^{}2\left({\displaystyle \frac{m_s}{2a}}\right)\left({\displaystyle \frac{m_s}{2a}}\right)^{}2\left({\displaystyle \frac{aG}{m_s}}\right)\left({\displaystyle \frac{aG}{m_s}}\right)^{}`$ (32)
$`V^{\prime \prime }(a)`$ $`=`$ $`F^{\prime \prime }2\left[\left({\displaystyle \frac{m_s}{2a}}\right)^{}\right]^22\left({\displaystyle \frac{m_s}{2a}}\right)\left({\displaystyle \frac{m_s}{2a}}\right)^{\prime \prime }`$ (33)
$`2\left[\left({\displaystyle \frac{aG}{m_s}}\right)^{}\right]^22\left({\displaystyle \frac{aG}{m_s}}\right)\left({\displaystyle \frac{aG}{m_s}}\right)^{\prime \prime },`$
respectively. Evaluated at the static solution, at $`a=a_0`$, we verify that $`V(a_0)=0`$ and $`V^{}(a_0)=0`$. From the condition $`V^{}(a_0)=0`$, one extracts the following useful equilibrium relationship
$`\mathrm{\Gamma }\left({\displaystyle \frac{m_s}{2a_0}}\right)^{}=\left({\displaystyle \frac{a_0}{m_s}}\right)\left[F^{}2\left({\displaystyle \frac{a_0G}{m_s}}\right)\left({\displaystyle \frac{a_0G}{m_s}}\right)^{}\right],`$ (34)
which will be used in determining the master equation, responsible for dictating the stable equilibrium configurations.
The solution is stable if and only if $`V(a)`$ has a local minimum at $`a_0`$ and $`V^{\prime \prime }(a_0)>0`$ is verified. The latter condition takes the form
$`\left({\displaystyle \frac{m_s}{2a}}\right)\left({\displaystyle \frac{m_s}{2a}}\right)^{\prime \prime }<\mathrm{\Psi }\mathrm{\Gamma }^2,`$ (35)
where $`\mathrm{\Psi }`$ is defined as
$`\mathrm{\Psi }={\displaystyle \frac{F^{\prime \prime }}{2}}\left[\left({\displaystyle \frac{aG}{m_s}}\right)^{}\right]^2\left({\displaystyle \frac{aG}{m_s}}\right)\left({\displaystyle \frac{aG}{m_s}}\right)^{\prime \prime }.`$ (36)
### III.3 The master equation
Using $`m_s=4\pi a^2\sigma `$, and taking into account the radial derivative of $`\sigma ^{}`$, Eq. (21) can be rearranged to provide the following relationship
$$\left(\frac{m_s}{2a}\right)^{\prime \prime }=\mathrm{{\rm Y}}4\pi \sigma ^{}\eta ,$$
(37)
with the parameter $`\eta `$ defined as $`\eta =𝒫^{}/\sigma ^{}`$, and $`\mathrm{{\rm Y}}`$ given by
$$\mathrm{{\rm Y}}\frac{4\pi }{a}(\sigma +𝒫)+2\pi a\mathrm{\Xi }^{}.$$
(38)
Equation (37) will play a fundamental role in determining the stability regions of the respective solutions. Note that the parameter $`\sqrt{\eta }`$ is normally interpreted as the speed of sound, so that one would expect that $`0<\eta 1`$, based on the requirement that the speed of sound should not exceed the speed of light. However, in the presence of exotic matter this cannot naively de done so. Therefore, in this work the above range will be relaxed. We refer the reader to Ref. Poisson for an extensive discussion on the respective physical interpretation of $`\eta `$ in the presence of exotic matter.
We shall use $`\eta `$ as a parametrization of the stable equilibrium, so that there is no need to specify a surface equation of state. Thus, substituting Eq. (37) into Eq. (35), one deduces the master equation given by
$$\sigma ^{}m_s\eta _0>\mathrm{\Theta },$$
(39)
where $`\eta _0=\eta (a_0)`$ and $`\mathrm{\Theta }`$ is defined as
$$\mathrm{\Theta }\frac{a_0}{2\pi }\left(\mathrm{\Gamma }^2\mathrm{\Psi }\right)+\frac{1}{4\pi }m_s\mathrm{{\rm Y}}.$$
(40)
Now, from the master equation we find that the stable equilibrium regions are dictated by the following inequalities
$`\eta _0`$ $`>`$ $`\overline{\mathrm{\Theta }},\mathrm{if}\sigma ^{}m_s>0,`$ (41)
$`\eta _0`$ $`<`$ $`\overline{\mathrm{\Theta }},\mathrm{if}\sigma ^{}m_s<0,`$ (42)
with the definition
$`\overline{\mathrm{\Theta }}{\displaystyle \frac{\mathrm{\Theta }}{\sigma ^{}m_s}}.`$ (43)
We shall now model the phantom wormhole geometries by choosing the specific form and redshift functions considered in Ref. Lobo-phantom , and consequently determine the stability regions dictated by the inequalities (41)-(42). In the specific cases that follow, the explicit form of $`\overline{\mathrm{\Theta }}`$ is extremely messy, so that as in Ishak , we find it more instructive to show the stability regions graphically.
### III.4 Stability regions
#### Asymptotically flat spacetimes
Consider the specific choices for the redshift and form functions given by
$`\mathrm{\Phi }(r)`$ $`=`$ $`\mathrm{const},`$ (44)
$`b(r)`$ $`=`$ $`r_0(r/r_0)^{1/\omega },`$ (45)
respectively. These are solutions to Eq. (7), for an asymptotically flat spacetime.
The factor related to the net discontinuity of the momentum flux impinging on the shell, $`\mathrm{\Xi }`$, is provided by
$$\mathrm{\Xi }=\frac{1}{8\pi a_0^2}\frac{\left(\frac{1+\omega }{\omega }\right)\left(\frac{r_0}{a_0}\right)^{\frac{1+\omega }{\omega }}}{\sqrt{1\left(\frac{r_0}{a_0}\right)^{\frac{1+\omega }{\omega }}}}.$$
(46)
The factor deduced from the equilibrium condition, $`\mathrm{\Gamma }`$, is given by
$$\mathrm{\Gamma }=\frac{1}{2a_0}\left[\frac{1}{2}\frac{\left(\frac{1+\omega }{\omega }\right)\left(\frac{r_0}{a_0}\right)^{\frac{1+\omega }{\omega }}}{\sqrt{1\left(\frac{r_0}{a_0}\right)^{\frac{1+\omega }{\omega }}}}\frac{\frac{M}{a_0}}{\sqrt{1\frac{2M}{a_0}}}\right].$$
(47)
The radial derivative of the surface energy density, $`\sigma ^{}`$, evaluated at the static solution, which will be fundamental in determining the stability regions, takes the following form
$`\sigma ^{}={\displaystyle \frac{1}{4\pi a_0^2}}\left({\displaystyle \frac{1\frac{3M}{a_0}}{\sqrt{1\frac{2M}{a_0}}}}{\displaystyle \frac{1\left(\frac{1+\omega }{2\omega }\right)\left(\frac{r_0}{a_0}\right)^{\frac{1+\omega }{\omega }}}{\sqrt{1\left(\frac{r_0}{a_0}\right)^{\frac{1+\omega }{\omega }}}}}\right).`$ (48)
We shall not write down the explicit forms of the remaining functions, i.e., $`\mathrm{{\rm Y}}`$, $`\mathrm{\Psi }`$ and $`\mathrm{\Theta }`$, as they are extremely lengthy. However, the stability regions shall be shown graphically.
To determine the stability regions of this solution, we shall separate the cases of $`b(a_0)<2M`$ and $`b(a_0)>2M`$. From Eq. (18) and the definition of $`m_s=4\pi a_0^2\sigma `$, this corresponds to $`m_s>0`$ and $`m_s<0`$, respectively. Here, we shall relax the condition that the surface energy density be positive, as in considering traversable wormhole geometries, one is already dealing with exotic matter. Note that for $`\sigma <0`$, the weak energy condition is readily violated.
For $`b(a_0)<2M`$, i.e., for a positive surface energy density, and using the form function, Eq. (45), we need to impose the condition $`r_0<2M`$, so that the junction radius lies outside the event horizon, $`a_0>2M`$. Thus, the junction radius lies in the following range
$$2M<a_0<2M\left(\frac{2M}{r_0}\right)^{(1+\omega )}.$$
(49)
For a fixed value of $`\omega `$, we verify that as $`r_00`$, then $`a_0\mathrm{}`$. The range decreases, i.e., $`a_02M`$, as $`r_02M`$. Note that by fixing $`r_0`$ and decreasing $`\omega `$, the range of $`a_0`$ is also significantly increased.
For a fixed value of the parameter, for instance $`\omega =2`$, we shall consider the following cases: $`r_0/M=1.0`$, so that $`2<a_0/M<4`$; and for $`r_0/M=0.25`$, we have $`2<a_0/M<16`$. The respective stability regions are depicted in Fig. 1. From Eq. (48) we find that $`\sigma ^{}<0`$, and as we are considering a positive surface energy density, this implies $`m_s\sigma ^{}<0`$. Thus, the stability regions, dictated by the inequality (42), lie beneath the solid lines in the plots of Fig. 1. Note that for decreasing values of $`r_0/M`$, despite the fact that the range of $`a_0`$ increases, the values of $`\eta _0`$ are further restricted. Thus, adopting a conservative point of view, using positive surface energy densities, we note that stable phantom wormhole geometries may be found well within the bound of $`0<\eta _01`$, and the stability regions increase for increasing values of $`r_0/M`$.
For $`b(a_0)>2M`$, the surface mass of the thin shell is negative, $`m_s(a_0)<0`$. We shall separate the cases of $`r_0<2M`$ and $`r_0>2M`$.
If $`r_0<2M`$, the range of the junction radius is given by
$$a_0>2M\left(\frac{2M}{r_0}\right)^{(1+\omega )}.$$
(50)
For this specific case, $`\sigma ^{}`$ possesses one real positive root, $`R`$, in the range of Eq. (50), signalling the presence of an asymptote, $`\sigma ^{}|_R=0`$. We verify that $`\sigma ^{}<0`$ for $`2M(2M/r_0)^{(1+\omega )}<a_0<R`$, and $`\sigma ^{}>0`$ for $`a_0>R`$. Thus, the stability regions are given by
$`\eta _0`$ $`>`$ $`\overline{\mathrm{\Theta }},\mathrm{if}2M\left({\displaystyle \frac{2M}{r_0}}\right)^{(1+\omega )}<a_0<R,`$ (51)
$`\eta _0`$ $`<`$ $`\overline{\mathrm{\Theta }},\mathrm{if}a_0>R.`$ (52)
Consider for $`\omega =2`$, the particular cases of $`r_0/M=0.5`$, so that $`a_0/M>8`$, and $`r_0/M=1.5`$, so that $`a_0/M>2.667`$. The asymptotes, $`\sigma ^{}|_R=0`$, for these cases exist at $`R/M13.9`$ and $`R/M4.24`$, respectively. These cases are represented in Fig. 2. Note that for increasing values of $`r_0/M`$, the range of $`a_0`$ decreases, and the values of $`\eta _0`$ are less restricted. Thus, one may conclude that the stability regions increase, for increasing values of $`r_0/M`$.
If $`r_0>2M`$, then obviously $`a_0>r_0`$. We verify that $`\sigma ^{}>0`$, and consequently $`m_s\sigma ^{}<0`$, so that the stability region is given by inequality (42). We verify that the values of $`\eta _0`$ are always negative. However, by increasing $`r_0/M`$, the values of $`\eta _0`$ become less restricted, and the range of $`a_0`$ decreases.
#### Isotropic pressure, $`p_r=p_t=p`$
Consider the following functions
$`\mathrm{\Phi }(r)`$ $`=`$ $`\left({\displaystyle \frac{3\omega +1}{1+\omega }}\right)\mathrm{ln}\left(r/r_0\right),`$ (53)
$`b(r)`$ $`=`$ $`r_0(r/r_0)^{1/\omega },`$ (54)
which are solutions of a phantom wormhole possessing an isotropic pressure Lobo-phantom .
The factor related to the momentum flux term, $`\mathrm{\Xi }`$, is given by
$$\mathrm{\Xi }=\frac{1}{8\pi a_0^2}\sqrt{1\left(\frac{r_0}{a_0}\right)^{\frac{1+\omega }{\omega }}}\left[\frac{\left(\frac{1+\omega }{\omega }\right)\left(\frac{r_0}{a_0}\right)^{\frac{1+\omega }{\omega }}}{1\left(\frac{r_0}{a_0}\right)^{\frac{1+\omega }{\omega }}}\frac{6\omega }{1+\omega }\right].$$
(55)
The $`\mathrm{\Gamma }`$ and $`\sigma ^{}`$ are identical to the previous case of an asymptotically flat spacetime, and as before we shall not show the specific forms of the remaining functions, as they are extremely lengthly.
To determine the stability regions of this solution, as in the previous case, we shall separate the cases of $`b(a_0)<2M`$ and $`b(a_0)>2M`$.
For $`b(a_0)<2M`$, we have $`m_s>0`$, and the condition $`r_0<2M`$ is imposed. Therefore, the junction radius lies in the same range as the previous case, i.e., Eq. (49). We also verify that $`\sigma ^{}<0`$ in the respective range. Thus the stability region is given by
$$\eta _0<\overline{\mathrm{\Theta }},\mathrm{if}2M<a_0<2M\left(\frac{2M}{r_0}\right)^{(1+\omega )}.$$
(56)
Consider, for simplicity, $`\omega =2`$, and the cases for $`r_0/M=1`$ and $`r_0/M=1.5`$ are analyzed in Fig. 3. The ranges are given by $`2<a_0/M<4`$ and $`2<a_0/M<2.667`$, respectively. Note that as $`r_0/M`$ decreases, the range of $`a_0`$ increases. However, the values of the parameter $`\eta _0`$ become more restricted. Thus, one may conclude that the stability regions increase, as $`r_0/M`$ increases.
For $`b(a_0)>2M`$, then $`m_s(a_0)<0`$. As before, we shall separate the cases of $`r_0<2M`$ and $`r_0>2M`$. For $`r_0<2M`$, the range of $`a_0`$ is given by $`a_0>2M(2M/r_0)^{(1+\omega )}`$, as in the previous case of the asymptotically flat wormhole spacetime.
For this case $`\sigma ^{}`$ also possesses one real positive root, $`R`$, in the respective range. We have $`\sigma ^{}<0`$ for $`(2M/r_0)^{1/\omega }<a<R`$, and $`\sigma ^{}>0`$ for $`a_0>R`$. The stability regions are also given by the conditions (51)-(52). We have considered the specific cases of $`r_0/M=1`$ so that the respective range is $`a_0/M>4`$; and $`r_0/M=1.5`$, so that $`a_0/M>2.667`$. The asymptotes, $`\sigma ^{}|_R=0`$, for these cases exist at $`R/M6.72`$ and $`R/M4.24`$, respectively. This analysis is depicted in the plots of Fig. 4. Note that the plots given by $`\overline{\mathrm{\Theta }}`$ are inverted relatively to the asymptotically flat spacetime. For decreasing values of $`r_0/M`$, note that the value of the stability parameter $`\eta _0`$ becomes less restricted and the range of the junction radius increases. Thus, one may conclude that the stability regions increase for decreasing values of $`r_0/M`$.
If $`r_0>2M`$, then $`a_0>r_0`$. We find that $`\sigma ^{}>0`$, which implies $`m_s\sigma ^{}<0`$. Consider $`\omega =2`$, and the specific case of $`r_0/M=2.5`$, so that the stability region lies below the solid line in Fig. 5. We also verify that for increasing values of $`r_0/M`$, the values of the parameter $`\eta _0`$ become further restricted. Thus, one may conclude that the stability regions decrease for increasing values of $`r_0/M`$.
## IV Summary and Discussion
As the Universe is probably constituted of approximately 70% of null energy condition violating phantom energy, this cosmic fluid may be used as a possible source to theoretically construct traversable wormholes. In fact, it was found that infinitesimal amounts of phantom energy may support traversable wormholes Lobo-phantom . In this paper, we have modelled phantom wormholes by matching an interior traversable wormhole geometry, satisfying the equation of state $`p=\omega \rho `$ with $`\omega <1`$, to an exterior vacuum solution at a finite junction interface. We have analyzed the stability of these phantom wormholes, an issue of fundamental importance, to linearized perturbations around static solutions, by including the momentum flux term in the conservation identity. We have considered two particularly interesting cases, namely, that of an asymptotically flat spacetime, and that of an isotropic pressure wormhole geometry. The latter solution is of particular interest, as the notion of phantom energy is that of a spatially homogeneous cosmic fluid, although it may be extended to inhomogeneous spherically symmetric spacetimes. We have separated the cases of positive and negative surface energy densities and found that the stable equilibrium regions may be significantly increased by strategically varying the wormhole throat. As we are considering exotic matter, we have relaxed the condition $`0<\eta _01`$, and found stability regions for phantom wormholes well beyond this range. There are several known examples of exotic $`\eta _0<0`$ behavior, namely the Casimir effect and the false vacuum Poisson , so that one cannot a priori impose $`0<\eta _01`$ until a detailed microphysical model of exotic matter is devised.
As emphasized in Ref. Lobo-phantom , these stable phantom wormholes have far-reaching physical and cosmological implications. First, apart from being used for interstellar travel, they may be transformed into time-machines mty ; Visser , consequently violating causality with the associated time travel paradoxes. Relative to the cosmological consequences, the existence of phantom energy presents us with a natural scenario for traversal wormholes. It was shown by González-Díaz gonzalez2 , that due to the fact of the accelerated expansion of the Universe, macroscopic wormholes could naturally be grown from the quantum foam. It was shown that the wormhole’s size increases by a factor which is proportional to the scale factor of the Universe, and still increases significantly if the cosmic expansion is driven by phantom energy diaz-phantom3 . However, it was also found that using wormholes modelled by thin shells accreting phantom energy Var-Israel , the wormholes become asymptotically comoving with the cosmological background as the Big Rip is approached, so that the future of the universe is shown to be causal. |
warning/0506/astro-ph0506712.html | ar5iv | text | # The first miniquasar
## 1 Introduction
The first ‘seed’ black holes (BHs) that later grew to become the supermassive variety that power active galactic nuclei (AGNs) must have appeared at very early epochs, $`z>10`$: this is in order to have had sufficient time to build up via gas accretion a mass of several $`\times 10^9\mathrm{M}_{}`$ by $`z=6.4`$, the redshift of the most distant quasars discovered in the Sloan Digital Sky Survey (Fan et al. 2003). The origin and nature of this seed population remain uncertain. Numerical simulations performed in the context of hierarchical structure formation theories show that the first stars (the so-called ‘Population III’) in the Universe formed out of metal-free gas in dark matter ‘minihaloes’ of mass above a few $`\times 10^5h^1\mathrm{M}_{}`$ (Abel, Bryan, & Norman 2000; Fuller & Couchman 2000; Yoshida et al. 2003; Reed et al. 2005) condensing from the rare high-$`\sigma `$ peaks of the primordial density fluctuation field at $`z>20`$, and were likely very massive (e.g. Abel, Bryan, & Norman 2002; Bromm, Coppi, & Larson 2002; see Bromm & Larson 2004 for a recent review). Non-rotating very massive stars in the mass window $`150m_{}250\mathrm{M}_{}`$ are expected to disappear as pair-instability supernovae (Bond, Arnett, & Carr 1984) and leave no compact remnants. Stars with $`40m_{}150\mathrm{M}_{}`$ and $`m_{}250\mathrm{M}_{}`$ are predicted instead to collapse to BHs with masses comparable to those of their progenitors (Fryer, Woosley, & Heger 2001). Barring any fine tuning of the initial mass function of Population III stars, intermediate-mass BHs – with masses above the 5–20$`\mathrm{M}_{}`$ range of known ‘stellar-mass’ holes – may then be the inevitable endproduct of the first episodes of pregalactic star formation (Madau & Rees 2001). Another route for the creation of more massive black hole seeds may be the formation – and subsequent collapse as a result of the post-Newtonian instability – of supermassive stars with $`m_{}10^3\mathrm{M}_{}`$ (see, e.g., Shapiro 2004a for a recent review) out of the lowest angular momentum gas in rare haloes above the minimum mass for atomic cooling at high redshifts (Bromm & Loeb 2003; Koushiappas, Bullock, & Dekel 2004). Our simulations are related more closely to the former type of scenario.
Physical conditions in the central potential wells of gas-rich protogalaxies may have been propitious for BH accretion. In the absence of an H<sub>2</sub> photodissociating UV flux and of ionizing X-ray radiation, three-dimensional simulations of early structure formation show that the fraction of cold, dense gas available for accretion on to seed holes or star formation exceeds 20 per cent for haloes more massive than $`10^6\mathrm{M}_{}`$ (Machacek, Bryan, & Abel 2003). There is also the possibility that the collapse of a rotating very massive star may lead to the formation of a massive BH and a ‘ready-made’ accretion disc (Fryer et al. 2001; Shapiro 2004b) that may provide a convenient source of fuel to power a ‘miniquasar’.
The presence of accreting BHs powering Eddington-limited miniquasars at such a crucial formative stage in the evolution of the Universe presents a challenge to models of the epoch of first light and of the thermal and ionization early history of the intergalactic medium (IGM), and serves as the main motivation of this paper. Feedback processes from the first stars and their remnants likely played a key role in reheating and structuring the IGM and in regulating gas cooling and star formation in pregalactic objects. Energetic photons from miniquasars may make the low-density IGM warm and weakly ionized prior to the epoch of reionization breakthrough (Madau et al. 2004; Ricotti, Ostriker, & Gnedin 2005). X-ray radiation may boost the free-electron fraction and catalyze the formation of H<sub>2</sub> molecules in dense regions, counteracting their destruction by UV Lyman-Werner radiation (Haiman, Abel, & Rees 2000; Machacek et al. 2003). Or it may furnish an entropy floor to the entire IGM, preventing gas contraction and therefore impeding rather than enhancing H<sub>2</sub> formation (Oh & Haiman 2003). Photoionization heating may evaporate back into the IGM some of the gas already incorporated into haloes with virial temperatures below a few thousand kelvins (Barkana & Loeb 1999; Haiman, Abel, & Madau 2001; Shapiro, Iliev, & Raga 2004). The detailed consequences of all these effects is poorly understood.
In this paper, we describe the results of fully 3D Eulerian cosmological hydrodynamical simulations of the effect of the first miniquasars on the thermal properties of the high-redshift IGM. The focus of our investigation is not the modeling of the processes that lead to gas accretion on to BHs, but rather the radiative feedback of this population on their environment and on structure formation as a whole. Since their host haloes form from the collapse of rare density fluctuations, miniquasars powered by accreting BHs are expected to be strongly clustered and highly biased tracers of the underlying dark matter (DM) distribution. We use the adaptive mesh refinement (AMR) technique to home in, with progressively finer resolution, on the densest parts of the ‘cosmic web’. Rather than approximating radiation fields as isotropic, we study the impact of a point-source of X-ray radiation that starts shining in a rare high-$`\sigma `$ peak at $`z=21`$, before a universal ionizing background is actually established. The outline of this paper is as follows. In § 2 we describe our suite of numerical simulations. § 3 presents general results of the simulations, with and without X-rays, and on the influence on the thermal and ionization state of the IGM of a non-uniform radiation field in the optically thin limit. We discuss the distance-dependent radiative feedback effect of a miniquasar on the amount of primordial gas that can cool and collapse, thus becoming available for star formation, in § 4. Finally, we present our conclusions in §5.
## 2 Simulations
High-resolution hydrodynamics simulations of early structure formation in $`\mathrm{\Lambda }`$CDM cosmologies are a powerful tool to track in detail the thermal and ionization history of a clumpy IGM and guide studies of early reheating. We have used a modified version of enzo, an adaptive mesh refinement (AMR), grid-based hybrid (hydro$`+`$N-body) code developed by Bryan & Norman (see http://cosmos.ucsd.edu/enzo/) to solve the cosmological hydrodynamics equations, study the cooling and collapse of primordial gas in DM haloes, and simulate the large-scale effect of a miniquasar turning on at very high redshift. All the results shown below assume a $`\mathrm{\Lambda }`$CDM world model with parameters $`\mathrm{\Omega }_M=0.3`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$, $`h=0.7`$, $`\mathrm{\Omega }_b=0.05`$, $`\sigma _8=0.85`$, and $`n=1`$.
The primordial distribution of gas and DM is initialized in a comoving box of 1 Mpc on a side using the linear power spectrum of Eisenstein and Hu (1999). We first identify in a pure N-body simulation the Lagrangian volume of the most massive DM halo at $`z=25`$. At this time its total mass is $`7\times 10^5\mathrm{M}_{}`$, corresponding to a 3.5 sigma peak in the density fluctuation field. Because its mass is comparable to the mass threshold for gas cloud formation by molecular cooling, this halo will likely harbor the first collapsed baryonic object in our simulated volume. It is this halo that we flag as the host of the first massive BH. We then generate new initial conditions centered around this density peak, consisting of two nested $`128^3`$ static grids, of which the inner one covers the central 0.5 Mpc volume. The DM density field is also sampled with $`128^3`$ particles in the inner region, leading to a mass resolution of $`m_{\mathrm{DM}}=2000\mathrm{M}_{}`$. This ensures that haloes above the cosmological Jeans mass are well resolved at all redshifts $`z<21`$. At $`z=15`$, the five most massive haloes in the box have between $`5\times 10^3`$ and $`10^4`$ DM particles.
During the evolution from $`z=99`$ to $`z=15`$, refined (child) grids are introduced with twice the spatial resolution of the coarser (parent) grid. The AMR is restricted to the inner 0.5 Mpc and is triggered when a cell reaches a DM overdensity (baryonic overdensity) of 2.0 (4.0). This low overdensity refinement criterion has been shown to yield comparable results in equivalent simulations with enzo and the smoothed particle hydrodynamics (SPH) code GADGET over the entire range of the dark halo mass function (O’Shea et al. 2005a). The region of interest is allowed to dynamically refine further to a total of 8 levels on a $`128^3`$ top grid, resulting in a maximum dynamic range (ratio of box size to the smallest spatial scale that can be resolved) of 32,768 or a maximum resolution of 30 pc (comoving). At the end of a typical simulation, the code uses $`2.2\times 10^7`$ computational grid cells on $`9800`$ grids, with a median number of cells per grid of 768.
The spatial distribution of $`128^3\times 15/8`$ collisionless dark matter particles determines at every time-step the large-scale gravitational field in the box. We evolve the non-equilibrium rate equations for nine species (H, H<sup>+</sup>, H<sup>-</sup>, e, He, He<sup>+</sup>, He<sup>++</sup>, H<sub>2</sub>, and H<sub>2</sub><sup>+</sup>) in primordial self-gravitating gas, including radiative losses from atomic and molecular line cooling, Compton heating by X-rays and Compton cooling by the cosmic background radiation. The rate coefficients for the reaction network for these species are primarily those described in Abel et al. (1997): the H<sub>2</sub> cooling function is that of Lepp & Shull (1983). We impose the constraint that the minimum temperature attainable by radiative cooling is that of the cosmic microwave background.
A zero-metallicity progenitor star of mass $`m_{}=150\mathrm{M}_{}`$ turning on at $`z=25`$ in the host halo will emit $`10^{64}`$ photons above 13.6 eV over a lifetime of $`2.3`$ Myr (Schaerer 2002), about 10 such photons per halo baryon. The ensuing ionization front will overrun the host halo, photoevaporating most of the surrounding gas (Whalen, Abel, & Norman 2004). Gas accretion on to the BH remnants of the first stars that create H ii regions may have to wait for new cold material to be made available through the hierarchical merging of many gaseous subunits. The simulations show that, by $`z=21`$, our massive BH has been incorporated into a larger DM halo that contains 7 times more gas than the original host: it is at this time that we turn on the miniquasar radiation field. Our miniquasar is powered by a $`m_{\mathrm{BH}}=150\mathrm{M}_{}`$ hole accreting at the Eddington rate and shining for 1 (2 in one of our 7 simulations) Salpeter time-scale, $`t_S=450ϵ/(1ϵ)`$ Myr, with a radiative efficiency of $`ϵ=10`$ per cent. Its exponentially growing mass and luminosity are recomputed at every time-step, and its trajectory is followed by flagging the DM particle corresponding to the maximum density in the host halo as our BH. The location of this flagged particle is determined at every time-step and used as the origin of an isotropic $`1/r^2`$ radiation field: from $`z=21`$ to $`z=15`$ our BH typically moves by less than 6 kpc (comoving). Note that, because of the considerable mass difference between our DM particle resolution ($`2000\mathrm{M}_{}`$) and the initial mass of the black hole ($`150\mathrm{M}_{}`$), we are not able to accurately track the true trajectory of the hole.
The photon energy distribution of putative high-z miniquasars is uncertain. The spectra of ‘ultraluminous’ X-ray sources in nearby galaxies appear to require both a soft component (well fit by a cool multicolour disc blackbody with $`kT_{\mathrm{max}}0.15`$ keV, which may indicated intermediate-mass BHs; Miller & Colbert 2004) and a non-thermal power-law component, $`L_EE^\alpha `$, of comparable luminosity and slope $`\alpha 1`$. Here we have run simulations with three different spectra: one (‘PL’) assuming a simple power-law miniquasar spectrum with $`\alpha =1`$ for photons with energies in the range $`0.210`$ keV; one (‘PLhard’) with a similar spectrum, but with a low-energy cutoff of 0.4 keV; and one (‘MCD’) with luminosities equally divided between a multicolour disc component and a power law with $`\alpha =1.2`$, both with a 0.2 keV low-energy cutoff. In the multicolour disc component, each annulus of a thin accretion disc is assumed to radiate as a blackbody with a radius dependent temperature, $`T(r)r^{3/4}`$, and the temperature of the innermost portion of the disc is related to the mass of the BH as $`T_{\mathrm{in}}m_{\mathrm{BH}}^{1/4}`$ (e.g. Makishima et al. 2000). All spectral energy distributions are normalized to the Eddington luminosity, which is exponentially increasing during the active phase of the miniquasar. A plot of the different input spectra is shown in Fig. 1. For the MCD case only, we included in addition to the X-rays an H<sub>2</sub>-dissociating flux in the Lyman-Werner band (LW, $`11.213.6`$ eV), with an intensity determined by the Raleigh-Jeans tail of the multicolour disc. The implicit assumption for all of these spectra is that photons with energies between 13.6 eV and 0.2 keV are absorbed within the miniquasar host halo. For this reason our simulation is not strictly valid within the host halo, and we have excised it from all of the following analysis.
In the pre-reionization Universe, when the IGM is predominantly neutral, soft X-ray photons will be absorbed as they photoionize hydrogen or helium atoms. For a mixture of H and He with cosmic abundances, the effective bound-free absorption cross-section can be approximated to an accuracy of 30 per cent in the range $`50\mathrm{eV}<h\nu <2`$keV as $`\sigma _{\mathrm{bf}}4\times 10^{20}\mathrm{cm}^2(h\nu /0.1\mathrm{keV})^3`$. The mean free path of ionizing radiation in the neutral medium with overdensity $`\delta \rho /\rho _b`$ is then
$$\lambda =\frac{1}{n_\mathrm{H}\sigma _{\mathrm{bf}}}5\mathrm{kpc}\left(\frac{1+z}{20}\right)^3\left(\frac{h\nu }{0.1\mathrm{keV}}\right)^3\delta ^1.$$
(1)
Throughout the simulation box, gas at the mean density will be transparent to photons with energies $`>0.2`$ keV and we take advantage of this by working in the optically thin approximation. At $`z=21`$ even the most massive haloes in our simulation have total hydrogen columns below $`10^{22}`$ cm<sup>-2</sup>, i.e. are transparent to photons above 0.7 keV. Since a power-law spectrum with $`\nu I_\nu `$ const is characterized by equal power per logarithmic frequency interval, photoelectric absorption by intervening haloes will not significantly attenuate the ionizing energy flux: this is then only a function of distance from the miniquasar, and grows linearly with the mass of the hole in the PL and PLhard cases. Photo-ionizations and photodissociations couple the radiation field to the hydrodynamics through a heating term in the energy conservation equation and source and sink terms in the species abundance equations (Anninos et al. 1997). The radiative heating $`H_j`$ and photoionization and photodissociation $`\mathrm{\Gamma }_j`$ rate coefficients are given by
$`H_j`$ $`=`$ $`{\displaystyle _{\nu _{0,j}}^{\mathrm{}}}\sigma _j(\nu )I(\nu ){\displaystyle \frac{h\nu h\nu _{0,j}}{h\nu }}𝑑\nu `$ (2)
$`\mathrm{\Gamma }_j`$ $`=`$ $`{\displaystyle _{\nu _{0,j}}^{\mathrm{}}}\sigma _j(\nu ){\displaystyle \frac{I(\nu )}{h\nu }}𝑑\nu ,`$ (3)
where $`\sigma _j(\nu )`$ is the cross-section and $`\nu _{0,j}`$ is the frequency threshold for the $`j^{th}`$ reaction, and $`I(\nu )`$ is the intensity of the radial radiation field.
## 3 General results
### 3.1 ‘NoBH’ simulation
A reference simulation with no radiation source (‘NoBH’) was run for comparison. The clustered structure around the selected high-$`\sigma `$ peak of the density field is clearly seen in Fig. 2, a 3D volume rendering of the inner 0.5 Mpc box at redshifts 21 and 15.5. The figure shows gas at $`4<\delta <10`$, with the locations of dark matter minihaloes marked by spheres colored and sized according to their mass (the spheres are only markers, the actual shape of the haloes is typically non-spherical). Several interleaving filaments are visible, at the intersections of which minihaloes are typically found. We identify the locations of DM minihaloes in the simulation volume by employing the HOP halo finder developed by Eisenstein & Hut (1998). This algorithm identifies haloes by grouping together particles that are associated with the same local density maximum, without enforcing spherical symmetry around the densest particle. We consider only haloes containing more than 100 particles. At $`z=21`$, 55 haloes in our volume satisfy this criterion, by $`z=17.5`$ this number has grown to 149, and by $`z=15.5`$ to 262 haloes. At this epoch, only four haloes have reached the critical virial temperature for atomic cooling, $`T_{\mathrm{vir}}=10^4`$K, where
$$T_{\mathrm{vir}}=1.7\times 10^4\left(\frac{\mu }{1.2}\right)\left(\frac{M_{\mathrm{halo}}}{10^7h^1\mathrm{M}_{}}\right)^{2/3}$$
$$\times \left[\frac{\mathrm{\Omega }_M}{\mathrm{\Omega }_M(z)}\frac{\mathrm{\Delta }_{\mathrm{vir}}}{18\pi ^2}\right]^{1/3}\left(\frac{1+z}{20}\right)\mathrm{K},$$
(4)
$`\mu `$ is the mean molecular weight, and $`\mathrm{\Delta }_{\mathrm{vir}}`$ is the density contrast at virialization.
The primordial fractional abundance of H<sub>2</sub> in the IGM is small, $`x_{\mathrm{H}_2}2\times 10^6`$, as at $`z>100`$ H<sub>2</sub> formation is inhibited because the required intermediaries, either H<sub>2</sub><sup>+</sup> or H<sup>-</sup>, are destroyed by cosmic microwave background (CMB) photons (e.g. Galli & Palla 1998). Most of the gas in the simulation therefore cools by adiabatic expansion. Within collapsing minihaloes, however, gas densities and temperatures are large enough that H<sub>2</sub> formation is catalyzed by H<sup>-</sup> ions through the associative detachment reaction H+H$`{}_{}{}^{}`$ H<sub>2</sub>+$`e^{}`$, and the molecular fraction increases at the rate $`dx_{\mathrm{H}_2}/dtx_en_{\mathrm{HI}}T_{\mathrm{vir}}^{0.88}`$, where $`x_e`$ is the number of electrons per hydrogen atom. For $`T_{\mathrm{vir}}`$ a few thousand kelvins the virialization shock is not ionizing, the free electrons left over from recombination are depleted in the denser regions, and the production of H<sub>2</sub> stalls at a temperature-dependent asymptotic molecular fraction $`x_{\mathrm{H}_2}10^8T_{\mathrm{vir}}^{1.5}\mathrm{ln}(1+t/t_{\mathrm{rec}})`$, where $`t_{\mathrm{rec}}`$ is the hydrogen recombination time-scale (Tegmark et al. 1997). A typical H<sub>2</sub> fraction in excess of 200 times the primordial value is therefore produced after the collapse of structures with virial temperatures of order $`10^3`$K. This is large enough to efficiently cool the gas and allow it to collapse within a Hubble time unless significant heating occurs during this phase (Abel et al. 2000; Yoshida et al. 2003).
Fig. 3 (left panel) shows the fraction of cold gas within the virial radius as a function of halo mass for all the haloes identified at redshift 17.5. Note that the presence of halo masses below the HOP selection limit of $`2\times 10^5\mathrm{M}_{}`$ is caused by differences between the extent of a halo as defined by the HOP algorithm and the virial radius measured by the spherically-averaged radial profile analysis. Following Machacek et al. (2001), we define $`f_c`$ as the fraction of gas with temperature $`<0.5T_{\mathrm{vir}}`$ and density $`>1000`$ times the background (this is the halo gas that is able to cool below the virial temperature because of H<sub>2</sub>), and $`f_{\mathrm{cd}}`$ as the fraction of gas with temperature $`<0.5T_{\mathrm{vir}}`$ and (physical) density $`>10^{19}\mathrm{M}_{}`$ Mpc<sup>-3</sup> (this is the self-gravitating gas available for star formation). As in Machacek et al. (2001), we find that both $`f_c`$ and $`f_{\mathrm{cd}}`$ are correlated with the halo mass, with the fraction of cold$`+`$dense gas increasing less rapidly than $`f_c`$ with $`M_{\mathrm{halo}}`$. The mass threshold for significant baryonic condensation (non-zero $`f_{\mathrm{cd}}`$) is approximately $`5\times 10^5\mathrm{M}_{}`$ at these redshifts (Haiman, Thoul, & Loeb 1996). Also depicted in Fig. 3 (right panel) is the mass-weighted mean molecular fraction of all haloes with $`T_{\mathrm{vir}}>400`$K. Filled circles represent haloes with $`f_{\mathrm{cd}}>0.1`$, while open circles represent the others. The straight line marks the scaling of the asymptotic molecular fraction in the electron-depletion transition regime. Haloes with significant cold$`+`$dense gas fraction have larger $`x_{\mathrm{H}_2}`$ than those shown in fig. 3 of Yoshida et al. (2003). This is due to the high resolution attainable by an AMR code compared to SPH. The maximum gas density reached at redshift 15 in the most refined region of our simulation is $`4\times 10^5`$cm<sup>-3</sup> (corresponding to an overdensity of $`3\times 10^8`$): within this cold pocket the excited states of H<sub>2</sub> are in LTE and the cooling time is nearly independent of density (e.g. Lepp & Shull 1983).
At $`z=17.5`$, $`90`$ per cent of the simulated volume contains gas with $`T<10`$ K, as shown by the differential temperature distributions in Fig. 4. The distribution of the entropy parameter $`\kappa =kTn^{2/3}`$ is even more strongly peaked than the temperature distribution, as entropy is conserved during infall on to non-linear structures or Hubble expansion, which heat/cool the gas adiabatically. After decoupling from the CMB, the IGM temperatures drops as $`T(z)2.73(1+z_d)[(1+z)/(1+z_d)]^2`$K for $`z<z_d=150`$, giving rise to an entropy floor of $`0.04`$ eV cm<sup>2</sup>, independent of redshift. Only gas that is shock heated during accretion on to DM haloes or that cools radiatively via H<sub>2</sub> can depart from this peak. Shock heating produces the extended tail out to $`\kappa =10`$ eV cm<sup>2</sup> at a volume fraction in the range $`f_V=10^510^4`$. At these epochs, we find a collapsed gas mass fraction of $`f_M=2.5`$ per cent. The fraction of baryons that gets shock heated ($`\kappa >0.1`$ eV cm<sup>2</sup>) is $`f_M=6.7`$ per cent. Radiative cooling processes affects such a tiny fraction of the volume that it remains outside the plotted range.<sup>1</sup><sup>1</sup>1The effect of radiative cooling can be observed in the ‘phase diagram’ of Fig. 6 as the bluish tail at $`\delta >100`$, $`T<1000`$ K, and volume fraction $`f_V<10^8`$. Due to the lack of ionizing sources the medium is almost completely neutral, with a distribution of electron fractions sharply peaked at $`x_e=1.5\times 10^4`$. Note that the distribution of molecular hydrogen fraction, $`x_{\mathrm{H}_2}`$, is considerably broader than that of $`x_e`$.
Another quantity of interest is the gas clumping factor $`C=n^2/n^2`$, where the brackets denote volume-averaged quantities. The shorter radiative recombination time-scale, $`t_{\mathrm{rec}}=(\alpha n_eC)^1`$ of a clumpy medium relative to a homogeneous one increases the total number of photons per baryon required for reionization. We have calculated the average clumping factor within our simulation box of volume $`L^3`$ as
$$C=C_{\mathrm{IGM}}+C_h=\underset{\delta <70}{}\frac{n_i^2V_i}{n_H^2L^3}+\underset{\delta >70}{}\frac{n_i^2V_i}{n_H^2L^3},$$
(5)
where $`n_i`$ is the number density of hydrogen in the $`i^{th}`$ grid cell of volume $`V_i`$, $`n_H`$ is the background hydrogen density, and the two sums are taken over all grid cells with baryonic overdensity smaller and larger than $`70`$. We use this density threshold – corresponding to the mean baryonic overdensity at the virial radius of our haloes – to differentiate dense gas belonging to virialized haloes ($`C_h`$) and diffuse uncollapsed ‘intergalactic’ material ($`C_{\mathrm{IGM}}`$). Gas within haloes only contributes to recombinations when it is photoionized: self-shielding of a cold gaseous disc and photoevaporation of halo gas by an external radiation field all make the contribution of dense material with $`\delta 70`$ to the overall clumping very uncertain (e.g. Haiman et al. 2001). In Fig. 5 we have plotted the evolution with redshift of $`C_{\mathrm{IGM}}`$ and $`C_h`$ in the NoBH simulation. The clumping factor increases rapidly with time as structure formation progresses, with $`C_{\mathrm{IGM}}2.65\mathrm{exp}[0.2(21z)]`$ over the range in redshift we probe.
### 3.2 Heating and ionization by the miniquasar
X-ray radiation from the miniquasar partially ionizes most of the gas in the simulated volume both by direct photoionization of H i and He i and indirectly by collisional ionization of H i by the fast photolectrons. We have adopted the fitting formulae of Shull & van Steenberg (1985), cast as a function of the hydrogen ionized fraction, $`x=n_{\mathrm{H}^+}/n_\mathrm{H}`$, to determine the fraction of energy of a photoelectron deposited as heat versus further ionization. For photon energies above 0.2 keV, the ratio of the H i /He i photoionization cross-sections drops below 4 per cent: since the primordial ratio of helium to hydrogen is about 8 per cent, the photoionization and heating rates are dominated by helium absorption, which exceeds the hydrogen contribution by $`2:1`$. While He i is the main source of hot primary photoelectrons, however, it is H i that undergoes the bulk of secondary ionizations. A primary nonthermal photoelectron of energy $`E=1`$keV in a medium with residual ionization fraction (from the recombination epoch) $`x=2\times 10^4`$ will create over two dozens secondary electrons, depositing a fraction $`f_1=0.37`$ of its initial energy as secondary ionizations of H i , $`f_2=0.05`$ as secondary ionizations of He i , and $`f_3=0.13`$ as heat. The time-scale for electron-electron encounters resulting in a fractional energy loss $`f=\mathrm{\Delta }E/E`$,
$$t_{\mathrm{ee}}140\mathrm{yr}Ef\left(\frac{1+z}{20}\right)^3\left(\frac{\mathrm{ln}\mathrm{\Lambda }}{20}\right)^1x^1\delta ^1$$
(6)
(where $`E`$ is measured in keV), is typically much shorter than the electron Compton cooling time-scale: the primary photoelectron will therefore ionize and heat the surrounding medium before it is cooled by the CMB. The fraction of primary energy going into heat increases gradually with $`x`$, reaching $`f_3=0.5`$ at $`x=0.04`$. Secondary ionizations enhance the rates for ionization of H i and He i (but not for He ii , Shull & Van Steenberg 1985), and introduce an additional coupling between the abundance rate equations. The rate coefficients for H i and He i ionizations can be written as
$$\mathrm{\Gamma }_{\mathrm{HI}}+f_1(x)\left(H_{\mathrm{HI}}+\frac{n_{\mathrm{HeI}}}{n_{\mathrm{HI}}}H_{\mathrm{HeI}}\right)\frac{1}{13.6\mathrm{eV}}$$
(7)
and
$$\mathrm{\Gamma }_{\mathrm{HeI}}+f_2(x)\left(H_{\mathrm{HeI}}+\frac{n_{\mathrm{HI}}}{n_{\mathrm{HeI}}}H_{\mathrm{HI}}\right)\frac{1}{24.6\mathrm{eV}},$$
(8)
respectively. We self-consistently evolve the hydrogen ionization fraction $`x`$ and the hydrogen and helium abundances throughout the simulations (note that the Shull & van Steenberg fitting formulae were derived assuming equal ionization fractions of hydrogen and singly-ionized helium).
The two-dimensional distribution of gas overdensity and temperature at $`z=17.5`$ is shown in Fig. 6 for the NoBH and PL simulations. The color coding in this phase diagram indicates the fraction of the simulated volume at a given ($`\delta ,T)`$. Most of the gas in the NoBH run lies on the yellow line ($`\kappa =KTn^{2/3}T=const`$) representing the initial adiabat. At low overdensities the temperature either drops because of Hubble expansion or rises because of adiabatic compression until the gas is shock heated (red and green swath) to virial values, $`T=10^310^4`$K. At higher densities, the blue cooling branch follows the evolutionary tracks in the temperature-density plane for spherically collapsing clouds (Yoshida et al. 2003). H<sub>2</sub> line emission lowers the temperature down to $`100`$K, the minimum value attainable by molecular cooling. The onset of the gravitational instability can further compress the gas and cause a modest rise in temperature again (Abel et al. 2000).
After shining for a Salpeter time-scale, the miniquasar has heated up the box to a volume-averaged temperature of 2800 K. The mean electron fraction and entropy are now 0.03 and 17 eV cm<sup>2</sup> (see Fig. 4): farther than 20 kpc (comoving) from the source hydrogen is never ionized to more than 30 per cent. Gas near the miniquasar is heated above $`10^4`$K and quickly cools down via efficient atomic processes: the green finger at $`\mathrm{log}\delta =0,\mathrm{log}(T/\mathrm{K})>4`$ represents baryonic material in this phase. The increased electron fraction and gas temperature boost the gas-phase H<sub>2</sub> production, which occurs on a time-scale,
$$t_{\mathrm{H}_2}=30\mathrm{Myr}\left(\frac{x_{\mathrm{H}_2}}{10^5}\right)\left(\frac{0.01}{x_e}\right)x_{\mathrm{HI}}^1T_3^{0.88}\delta ^1\left(\frac{1+z}{20}\right)^3,$$
(9)
that is much shorter than the hydrogen recombination time,
$$t_{\mathrm{rec}}=9\times 10^3\mathrm{Myr}\left(\frac{0.01}{x_e}\right)T_3^{0.64}\delta ^1\left(\frac{1+z}{20}\right)^3$$
(10)
(here $`T_3T/10^3\mathrm{K}`$). The volume-averaged molecular fraction is now 20 times larger than the primordial value, with denser filaments in the IGM ($`\delta 1020`$) being traced out by a smaller electron fraction and by a molecular fraction in the range $`10^510^3`$ (see Fig. 7). This large molecular fraction will hinder the buildup of a uniform UV photodissociating background, as the maximum optical depth of the IGM in the H<sub>2</sub> LW bands can now exceed 10 (cf. Ricotti, Gnedin, & Shull 2001; Haiman et al. 2000). It will also promote H<sub>2</sub> radiative cooling in filaments on a time-scale,
$$t_{\mathrm{cool}}650\mathrm{Myr}\left(\frac{10^3}{x_{\mathrm{H}_2}}\right)x_{\mathrm{HI}}^1\delta ^1\left(\frac{1+z}{20}\right)^3e^{730\mathrm{K}/T},$$
(11)
(Machacek et al. 2001) that can be shorter than the Hubble expansion time, $`1/H285\mathrm{Myr}[(1+z)/20]^{3/2}`$, and the Compton cooling time for mostly neutral primordial gas, $`t_C=8\mathrm{Myr}x_e^1[(1+z)/20]^4`$. In Fig. 8 we plot the mean H<sub>2</sub> cooling time in units of the expansion time-scale $`H^1`$ versus overdensity, at $`z=17.5`$. In the NoBH simulation only gas with $`\delta >600`$ can cool in a Hubble time. This critical cooling density is lowered to $`\delta 40`$ in the PL simulation, due to the enhanced H<sub>2</sub> fraction. Molecular cooling can then affect the dynamics of baryonic material before it has fallen into the potential well of DM haloes and virialized. Gas in the simulation box is now able to collapse and cool more rapidly within subsequent star-forming minihaloes, without the need to produce much additional H<sub>2</sub>.
The global environmental impact of the miniquasar is illustrated in Fig. 9, where we plot in the top panel the volume ($`f_V`$) and mass fraction ($`f_M`$) as a function of overdensity in the NoBH and PL runs. The bottom panel shows the corresponding percentage change from the NoBH to the PL case ($`100\times (f_{\mathrm{PL}}f_{\mathrm{NoBH}})/f_{\mathrm{NoBH}}`$). Three prominent features are clearly seen in this figure: 1) over most of the simulation volume, the net effect of X-rays is to reduce gas clumping due to the smoothing of sheets and filaments in the ‘cosmic web’ by gas pressure. This ‘Jeans smoothing’ removes gas from intermediate $`\delta `$’s and puffs it up to low overdensities ($`\mathrm{\Delta }f>0`$ for $`1<\delta <10`$), thereby filling in underdense regions ($`\mathrm{\Delta }f<0`$ for $`\delta <1`$);<sup>2</sup><sup>2</sup>2In the miniquasar simulations the IGM clumping factor remains approximately constant, $`C_{\mathrm{IGM}}3`$, from $`z=21`$ to $`z=15`$ (see the left panel in Fig. 5). Because of the decreasing background density with time, the IGM recombination time-scale at a given temperature therefore increases as $`(1+z)^3`$ over this redshift interval. 2) the suppression of baryonic infall and the photoevaporation back into the IGM of some of the gas already incorporated into haloes both lower the gas mass fraction at $`20<\delta <2000`$; and 3) enhanced molecular cooling increases the amount of dense material at $`\delta >2000`$. While the vast majority of the baryons, both in volume and mass, are heated up, the densest gas at the centres of haloes is actually cooled, due to the X-ray catalysis of molecular hydrogen. The feedback effect of X-rays on the formation and cooling of pregalactic minihaloes will be discussed in more details in the next section.
In the PLlong simulation we let the miniquasar shine for two Salpeter times, from $`z=21`$ to $`z=15.5`$, leading to a final black hole mass of $`1100\mathrm{M}_{}`$. The increased intensity and duration of the radiation further enhances the effects described above (see Fig. 4). At $`z=15.5`$, the volume averages in the box are: $`T=6150`$K, $`\kappa =49\mathrm{eVcm}^2`$, $`x_e=0.07`$, and $`x_{\mathrm{H}_2}=9.3\times 10^5`$. The additional heating, however, does not further decrease the IGM clumping factor, and the halo clumping factor is still not affected (see Fig.5). The increase in $`x_{\mathrm{H}_2}`$ by more than a factor of two is the most significant difference between the PL to PLlong simulation. This will make it even harder for subsequent generation of sources to establish a uniform Lyman-Werner background. Note that if early black holes undergo a period of super-Eddington mass growth (Volonteri & Rees 2005; Haiman 2004), a possible subsequent luminous miniquasar phase, powered by a more massive hole, might have feedback effects more similar to, or even exceeding those of our PLlong simulation.
In order to probe the dependence of our results on the hardness of the radiation field, we ran the PLhard simulation, in which the lower energy cutoff was raised to $`\nu _{\mathrm{min}}`$=0.4 keV. Normalizing the total energy output to $`L_{\mathrm{Edd}}`$ leads to a specific radiation intensity which is 20 per cent higher (for $`\alpha =1`$) compared to the PL simulation. The heating rate, on the other hand, scales roughly as $`\nu _{\mathrm{min}}^3`$. A factor of two increase in $`\nu _{\mathrm{min}}`$ thus leads to a factor of eight decrease in the heating rate, and we expect the feedback from the miniquasar to be much weaker. This expectation is confirmed by the PLhard simulation, as shown by the dashed histogram in Fig. 4. At $`z=17.5`$ the volume-averaged temperature is only 380 K, with $`x_e=0.005`$ and $`x_{\mathrm{H}_2}=6.2\times 10^6`$. By the same arguments we would expect the miniquasar feedback to strengthen if the spectrum extended to lower energies, although the optically thin assumption would then break down.
## 4 X-ray radiative feedback: positive or negative?
The impact of the radiation produced by Population III stars and their remnants on H<sub>2</sub> chemistry and the properties of star-forming gas in low-mass haloes has been studied by many authors. Molecular hydrogen is fragile and easily photodissociated by soft UV photons below 13.6 eV (e.g. Haiman, Rees, & Loeb 1997; Ciardi, Ferrara, & Abel 2000; Machacek et al. 2001; Glover & Brand 2001). While this negative feedback may suppress gas cooling and collapse, the formation of H<sub>2</sub> molecules is catalyzed by free electrons, and any process that increased their abundance would boost the abundance of H<sub>2</sub> as well. Fossil H ii regions (Ricotti et al. 2001) and X-ray photons may provide such positive feedback (Haiman et al. 2000; Venkatesan, Giroux, & Shull 2001; Machacek et al. 2003; Glover & Brand 2003; Cen 2003). The issue of whether positive feedback processes dominate over negative feedback is still open. All we offer here is an assessment of the effect of X-rays in the absence of a strong photodissociating flux. This condition would naturally be fulfilled if accreting intermediate-mass BHs were the endproduct of the first episodes of Population III star formation, as miniquasars are much more efficient sources of radiation than their stellar-progenitors (Madau et al. 2004). As already mentioned in § 2, our MCD run includes in addition to the X-rays a LW H<sub>2</sub>-dissociating flux with an intensity determined by the Raleigh-Jeans tail of the multicolour disc, $`I_{\mathrm{LW}}(z=21)=10^{21}`$ergs cm<sup>-2</sup> s<sup>-1</sup> Hz<sup>-1</sup> sr<sup>-1</sup> at a comoving distance of 7.2 kpc, which is 10 (3) times lower than the 0.2 keV (1 keV) X-ray flux (cf. Machacek et al. 2003). We find the results of the MCD simulation to be very similar to those of the PL run, and it is the latter that we analyze in detail below.
Fig. 10 depicts the change in halo gas content between the PL and NoBH simulations as a function of halo mass at two different epochs, $`z=17.5`$ (when the miniquasar stops shining) and $`z=15.5`$: haloes in the two simulations have been matched by location and DM mass. Different symbols represent total, cold, and cold$`+`$dense gas mass both for haloes closer than 75 comoving kpc from the radiation source and for those farther away. X-ray heating lowers the total gas mass in nearly all virialized DM clumps, both by photoevaporating baryons that were already incorporated into haloes and by suppressing baryonic infall into newly forming structures. This decrease in gas mass is more apparent below $`M_{\mathrm{halo}}10^6\mathrm{M}_{}`$, and it is the strongest in pregalactic clouds close to the miniquasars that have been exposed to the highest level of X-ray flux. Only a small fraction of the gas accreted by haloes in the NoBH run is retained by the host gravitational potential in the PL run. One of the potential wells that experiences the largest decrease in gas content is fairly massive, $`M_{\mathrm{halo}}=10^6\mathrm{M}_{}`$, is located only 12 kpc (comoving) away from the miniquasar, and has a total gas fraction of 0.7 per cent, more than 20 times lower than in the NoBH simulation!
Population III stars, however, can only form from gas that is able to cool and condense. Fig. 11 shows the mass weighted mean H<sub>2</sub> fraction inside haloes in the PL simulation at $`z=17.5`$ (cf. right panel of Fig. 3). Overall $`x_{\mathrm{H}_2}`$ increases by slightly more than one order of magnitude, independently of $`T_{\mathrm{vir}}`$. Compared with the NoBH case, haloes with $`f_{\mathrm{cd}}>0.1`$ can be found at slightly lower $`T_{\mathrm{vir}}`$. None of the haloes with $`T_{\mathrm{vir}}<800`$K has any cold+dense gas, even though their molecular fraction is quite high, $`x_{\mathrm{H}_2}10^3`$. These haloes are the least massive haloes in our simulation and were formed only after the miniquasar began to shine. Continuous X-ray heating has prevented the gas from cooling and falling into the halo. As a consequence their mean gas mass fraction ($`M_{\mathrm{gas}}/M_{\mathrm{tot}}`$) is only 0.08 and their mean central temperature exceeds their virial temperature by a factor of 3. Haloes with larger $`T_{\mathrm{vir}}`$, in contrast, formed earlier and their gas has contracted to densities where H<sub>2</sub> cooling is able to overcome X-ray heating.
Evidence for X-ray enhanced cooling is seen in many haloes above $`2\times 10^5\mathrm{M}_{}`$. The largest relative increase in the amount of cold material occurs in the mass range $`2\times 10^5<M_{\mathrm{halo}}<10^6\mathrm{M}_{}`$, where the boosting effect can exceed 1-2 orders of magnitude! In this sense, X-ray preheating appears to somewhat decrease the threshold mass required for efficient gas cooling and star formation, albeit not by a large factor. In more massive peaks most baryons are already cold in the absence of X-rays (see Fig. 3), and positive feedback can only promote a little additional cooling. Yet in many haloes in the proximity of our miniquasar the H<sub>2</sub>boosting effect of X-rays is too weak to overcome heating, and the cold and dense gas mass actually decreases. Table 2 summarizes the miniquasar’s feedback effect on the total, cold, and cold+dense gas mass. We find evidence for a global positive feedback in minihaloes more than 75 comoving kpc away from the miniquasar, as the total mass in cold and cold$`+`$dense baryons increase by 30 and 44 per cent, respectively. Farther than 150 kpc this positive feedback is slightly larger yet, reaching $`+50`$ per cent for cold gas. Within 75 kpc of the miniquasar, however, the feedback is negative, as the total mass in cold$`+`$dense baryons decreases by more than 50 per cent. The total amount of gas mass in haloes is reduced at all distances within our simulated volume, with haloes closer than 75 kpc losing on average more than half their gas mass. These effects are long-lasting, and persist at $`z=15.5`$, 40 Myr after the miniquasar has stopped shining (see right panel of Fig. 10). Reducing the total amount of gas in the outer regions of haloes, while raising the amount of cold and dense gas available for star formation in the centres, could make it easier for these haloes to chemically enrich their surroundings.
Overall, however, positive and negative feedbacks nearly compensate each other, and the net impact of X-rays on the total amount of cold material available for star formation within the simulated volume is remarkably mild. By $`z=17.5`$, exposure to X-ray radiation has boosted the total mass fractions of cold and cold$`+`$dense gas by only 14 and 6 per cent, respectively. There is hardly an X-ray ‘sterilizing effect’: star formation can proceed in minihaloes above the threshold mass, including those that are near the radiation source. In the most massive peaks at $`z=15.5`$ the fraction of gas available for star formation in the NoBH control simulation is comparable to the value found when X-rays are present. Qualitatively similar conclusions were reached by Machacek et al. (2003), who showed that an early X-ray background cannot overcome the negative feedback from H<sub>2</sub> photodissociation by the soft UV radiation spectrum of the first stellar sources. Note that, at a distance of 50 comoving kpc, the X-ray flux from our miniquasars is comparable to the strongest X-ray background adopted by Machacek et al. (their $`ϵ_\mathrm{x}=10`$ case). Here we find that, even in the absence of a strong LW flux, the radiative feedback from X-rays is subtle. It enhances gas cooling in lower-$`\sigma `$ peaks that are far away from the initial site of star formation, thus decreasing the clustering bias of the early pregalactic population, but does not dramatically reverse or promote the collapse of pregalactic clouds as a whole.
It has been suggested recently by Oh & Haiman (2003) that an early X-ray background would establish an entropy floor over the entire IGM, thus preventing gas contraction, H<sub>2</sub> formation, cooling, and the build up of dense cores in minihaloes. The implication is a large reduction in the collapsed gas fraction and a pause in the cosmic star formation history, before more massive haloes with $`T_{\mathrm{vir}}>10^4`$K (which can undergo atomic cooling) start forming. We find little evidence for preheating suppressing H<sub>2</sub> formation. None of our minihaloes exhibits an entropy floor and, as we have shown above, the amount of cold and dense gas in the centres is actually enhanced in haloes sufficiently removed from the miniquasar. Oh & Haiman considered the evolution of a pre-heated, isolated, uniform density gas parcel with initially primordial H<sub>2</sub> fraction, and showed that the entropy floor prevented subsequent H<sub>2</sub> formation. In our simulations, however, the heating source is also a catalyst of H<sub>2</sub>. This causes the H<sub>2</sub> cooling time to become shorter than the Hubble expansion time at $`\delta >40`$ (see Fig. 8), and thus gas entropy is no longer a conserved quantity inside and in the vicinity of minihaloes.
Fig. 12 shows spherically averaged mass-weighted profiles of $`\delta `$, $`T`$, $`\kappa `$, $`x_e`$ and $`x_{\mathrm{H}_2}`$ at $`z=17.5`$ of a typical low mass minihalo in our simulations. It has a dark matter mass of $`3\times 10^5\mathrm{M}_{}`$ ($`T_{\mathrm{vir}}925`$ K), a virial radius of 127 pc (proper), and is located 200 comoving kpc from the miniquasar. In the NoBH simulation it has a total gas mass of $`5.7\times 10^4\mathrm{M}_{}`$, a cold gas mass of $`1.1\times 10^4\mathrm{M}_{}`$, and no cold+dense gas. The central regions of this halo have only triggered refinement down to level 7, one shy of the maximum refinement level. It has a central core, where the density profile levels off at $`\delta =7000`$. The gas temperature rises from about 100 K to $`T_{\mathrm{vir}}`$ at the virial radius, and then drops down to the central value of $`350`$ K. This drop is possible due to the slightly elevated central H<sub>2</sub> fraction of $`2.7\times 10^4`$, which, however, is not high enough to allow the gas to condense to the cold+dense threshold density of $`330\mathrm{cm}^3`$ in a Hubble time. The X-ray flux from the miniquasar changes this picture dramatically. The H<sub>2</sub> catalysis has led to an enhancement of almost 2 orders of magnitude in $`x_{\mathrm{H}_2}`$ outside the virial radius, from $`10^5`$ to $`8\times 10^4`$. Towards the centre the increase in temperature and density allows the H<sub>2</sub> abundance to grow further, peaking at a central value of $`x_{\mathrm{H}_2}=6.6\times 10^3`$. This increase in $`x_{\mathrm{H}_2}`$ causes the gas to cool down to 140 K and greatly increases the central overdensity ($`\delta =3\times 10^5`$), triggering refinement all the way to the maximum refinement level. In the PL simulation the external medium is at a temperature of $`3000`$ K, which is much higher than $`T_{\mathrm{vir}}`$. Thus no virial shock developes, and H<sub>2</sub> cooling causes the temperature profile to decreases monotonically towards the centre, dropping below $`T_{\mathrm{vir}}`$ at $`40`$ pc, a third of the virial radius. Gas that is hotter than $`T_{\mathrm{vir}}`$ cannot accrete on to the halo, and as a consequence the outer halo regions show a reduced gas density. The total gas mass decreases by 30 per cent down to $`4.0\times 10^4\mathrm{M}_{}`$. The increased cooling, however, has boosted the amount of cold gas by 60 per cent, up to $`1.8\times 10^4\mathrm{M}_{}`$, and even allowed $`6.5\times 10^3\mathrm{M}_{}`$ to reach the cold+dense threshold. The entropy profile reflects the additional cooling in the centre: it also decreases monotonically and reaches a minimum of $`2.2\times 10^4\mathrm{eV}\mathrm{cm}^2`$, a factor of 30 below the central value in the NoBH case ($`6.4\times 10^3\mathrm{eV}\mathrm{cm}^2`$).
The spatial resolution of our AMR simulations is determined by the maximum refinement level, which is usually set by computational costs. Mesh cells that have refined down to the maximum refinement level can pose a problem. The physics of cooling and collapsing gas is inherently unstable: Jeans unstable cells should refine further, but the simulation does not allow it. This leads to an error in the numerical solution of the differential equations governing the system. One way to prevent the gas from continuing to trigger refinment is to introduce Artificial Pressure Support (APS) to stabilize the density peaks that have reached maximal refinement against further collapse. With APS the internal gas energy of every mesh cell that reaches maximum refinement and is Jeans unstable is raised to 10 times the value that would stabilize it against gravitational collapse. This provides support against further refinement, but comes at the cost of introducing an artificially high central pressure and temperature inside haloes that have reached maximum refinement. Neither applying APS nor letting the simulation progress without it is correctly modelling the physics of the problem. The only solution is to continue to refine further (Abel et al. 2002) or to introduce sink particles (Krumholz, McKee, & Klein 2004). In our fiducial simulations we have 50,000-60,000 mesh cells at the maximum refinement level of 8. In order to determine whether the associated numerical error affects the conclusions of our study, we have re-run the PL simulation from $`z=21`$ to $`z=17.5`$, with a maximum refinement level of 10 (for a maximum dynamic range of 131,072) and APS. In the right panel of Fig. 12 we plot for the PL and new ’PL10’ runs a comparison of the five profiles described above. As expected the two profiles agree excellently at radii outside of the central core ($`>10`$ pc). Even inside the core, the temperature and $`x_{\mathrm{H}_2}`$ profiles agree very well. The gas in the PL10 halo, however, has been able to reach an overdensity a factor of $`3.5`$ higher than in the PL case, and has $`50`$ per cent more cold+dense gas. Although we show here only profiles for one halo, we see similar trends for all haloes that reach maximum refinement, and so we conclude that the amount of cold+dense gas in haloes in the PL simulation should be viewed as a lower limit. The numerical error due to evolving the simulation with cells at the maximum refinement level is unlikely to affect the global conclusions of this paper.
## 5 Discussion
Active galactic nuclei powered by supermassive holes keep the Universe ionized at $`z4`$, structure the IGM, and probably regulate star formation in their host galaxies. Their seeds were likely planted at very early epochs through the collapse of massive stars. Intermediate-mass holes accreting gas from the surrounding medium may shine as miniquasars at redshifts as high as $`z20`$. In this paper we have carried out AMR cosmological simulations using enzo to address the thermodynamic effect of miniquasars on the IGM at early times.
X-ray radiation from the miniquasar efficiently heats the gas in the simulated box to a volume-averaged temperature of 2800 K after one Salpeter time ($`z=17.5`$), and 6150 K after two ($`z=15.5`$). The main effect of this sharp increase in temperature is a reduction of gas clumping in the IGM by as much as a factor of 3, due to the Jeans smoothing of sheets and filaments in the ‘cosmic web’ by increased gas pressure. This smoothing will lower the number of hydrogen recombinations and thus the number of UV photons per baryon required to reionize the universe. Since X-rays are more efficient at heating than at ionizing, the free electron fraction is never raised to more than $`10`$ per cent, reaching a volume average of 3 per cent at $`z=17.5`$. Provided that either collisions or Ly$`\alpha `$ radiation can couple the spin temperature to the kinetic temperature of the largely neutral gas, it may be possible to observe the X-ray heated region in redshifted 21-cm line emission against the CMB (Madau, Meiksin, & Rees 1997) with future facilities like the LOw Frequency ARray (LOFAR).
The elevated free electron fraction leads to a strong enhancement of molecular hydrogen, both inside haloes and in the more tenuous filaments. The volume-averaged H<sub>2</sub> fraction is raised to $`x_{\mathrm{H}_2}=4\times 10^5`$, 20 times larger than the primordial value. This will delay the buildup of a uniform UV photodissociating background by subsequent sources of Lyman-Werner photons. We plan to carry out a detailed calculation of this suppression and its consequences in future work. The increased H<sub>2</sub> abundance allows gas with overdensities $`\delta >40`$ to cool within a Hubble time. As a consequence the evolution of gas within haloes and filaments is not adiabatic, and no entropy floor is established. While global heating suppresses baryonic infall and lowers the gas mass fraction at overdensities $`\delta `$ in the range 20-2000, enhanced molecular cooling increases the amount of dense material at $`\delta >2000`$. Thus, while the vast majority of the baryons, both in volume and mass, is heated up, the densest gas at the centres of haloes is actually cooled down. The largest relative increase in the amount of cold material occurs in the mass range $`2\times 10^5<M_{\mathrm{halo}}<10^6\mathrm{M}_{}`$, where the boosting effect can exceed 1-2 orders of magnitude. Yet for many haloes in the proximity of our miniquasar ($`<75`$ kpc) the H<sub>2</sub>-boosting effect is too weak to overcome heating, and the cold and dense gas mass actually decreases. Overall, the radiative feedback from X-rays enhances gas cooling in lower-$`\sigma `$ peaks that are far away from the initial site of star formation, thus decreasing the clustering bias of the early pregalactic population, but does not appear to dramatically reverse or promote the collapse of pregalactic clouds as a whole.
While we have included the essential physics necessary to model the radiative feedback from the first miniquasars, our study, like any, has its limitations. In the following we briefly address some of these:
1. We have not considered the environmental impact of the progenitor of the black hole. In the case of direct collapse of a primordial gas cloud to an intermediate-mass black hole the issue of a progenitor star would not arise. The more typical formation mechanism of the black hole, however, may be the collapse of a Population III star. In this case the copious UV radiation emitted during its main-sequence lifetime could have significantly altered the initial conditions for our simulation. Recent three-dimensional simulations (O’Shea et al. 2005b) seem to indicate that the heating and ionizing effect from the Population III progenitor are rather short lived. In their simulation the gas quickly recombines and cools, H<sub>2</sub> reforms, and gas is able to collapse to form a second generation star in a neighboring halo only 265 pc from the Population III host halo.
2. The strong increase in H<sub>2</sub> abundance catalyzed by the X-rays is dependent on the absence of a Lyman-Werner H<sub>2</sub>-dissociating background. As shown by Machacek et al. (2003) the positive feedback effects from an X-ray background in the presence of a LW background are rather mild, and they did not find a global increase in H<sub>2</sub> abundance by a factor of 20. For this reason our work should be considered as a simulation of the first miniquasar in its volume, i.e. before other sources establish a LW background. How soon after the miniquasar a LW background can be established, given the large increase in H<sub>2</sub>, remains to be seen.
3. Recently there has been a renewed interest in the importance of deuterium hydride (HD) to the cooling of primordial gas clouds (Johnson & Bromm 2005; Lipovka, Nunez-Lopez, & Avila-Reese 2005; Nagakura & Omukai 2005). HD becomes more efficient than H<sub>2</sub> cooling at $`T<200`$K, and, given a sufficient HD abundance, can cool gas down to below 100 K. In our simulations deuterium chemistry has been neglected. An increased free-electron fraction, as produced by our miniquasar, would lead to a catalysis of HD molecules, so it is quite possible that HD cooling may become important in the cores of our minihaloes.
4. As the emission spectrum of a miniquasar is uncertain, so are the photoheating and photoionization rates. We have mainly considered a pure power-law spectrum, ranging from 0.2 to 10 keV, with a slope of $`1`$. The power-law component in our MCD simulation is slightly softer ($`\alpha =1.2`$), but the differences between PL and MCD are negligible.
Note also that we have neglected UV radiation with $`13.6\mathrm{eV}<h\nu <200\mathrm{eV}`$. If a significant amount of the emitted Lyman-continuum flux escapes the host halo, it will be necessary to follow the radiative transfer in order to fully assess the environmental influence of the first miniquasars.
5. We have not taken into account the possibility of a feedback effect from the miniquasar’s radiation on the black hole’s mass accretion rate. Instead of assuming a constant exponential growth rate, it would be better to determine $`\dot{M}`$ as a function of time from the amount of cold gas available for accretion. Unfortunately our simulations don’t resolve the length scale of the accretion flow, and a self-consistent treatment of the mass accretion rate will have to wait for higher resolution studies.
## Acknowledgments
We thank T. Abel, G. Bryan, Z. Haiman, and B. O’Shea for many informative discussions, Ryan Montgomery for help with the phase diagrams, and the referee, Marie Machacek, for helpful comments that improved the flow of the paper. Fig. 2 was created with Nick Gnedin’s Ionization Front Interactive Tool (IFrIT). Support for this work was provided by NSF grants PHY99-07949 and AST02-05738, and by NASA grants NAG5-11513 and NNG04GK85G (P.M.). P.M. acknowledges support from the Alexander von Humboldt Foundation. M.K. thanks the Graduate Fellowship Program at Kavli Institute for Theoretical Physics, Santa Barbara, where part of this work was done. All computations were performed on NASA’s Project Columbia supercomputer system and on UpsAnd, a Beowulf cluster at UCSC. Movies of the simulations are available at http://www.ucolick.org/$``$mqk/miniqso. |
warning/0506/cond-mat0506752.html | ar5iv | text | # Nonequilibrium Perturbative Formalism and Spectral Function for the Anderson Model
## 1 Introduction
### 1.1 Nonequilibrium Perturbative Formalism
The basic idea on the nonequilibrium perturbation theory grounded on the time-contour which starts and ends at $`t=\mathrm{}`$ via $`t=\mathrm{}`$ has been proposed by Schwinger. After that, the frame of the nonequilibrium perturbation theory has been built up using the nonequilibrium Green’s functions given after the time-contour by Keldysh. The perturbative equation is expressed in matrix form:
$`𝐆`$ $`=`$ $`𝐠+𝐠𝚺𝐆,`$ (1)
where
$`𝐆=\left[\begin{array}{cc}G^{}\hfill & G^<\hfill \\ G^>\hfill & G^{++}\hfill \end{array}\right],𝚺=\left[\begin{array}{cc}\mathrm{\Sigma }^{}\hfill & \mathrm{\Sigma }^<\hfill \\ \mathrm{\Sigma }^>\hfill & \mathrm{\Sigma }^{++}\hfill \end{array}\right].`$ (6)
The nonequilibrium Green’s functions are given in the Heisenberg representation by
$`G^{}(t_1,t_2)i\mathrm{T}\widehat{d}(t_1)\widehat{d}^{}(t_2),`$ (7)
$`G^{++}(t_1,t_2)i\stackrel{~}{\mathrm{T}}\widehat{d}(t_1)\widehat{d}^{}(t_2),`$ (8)
$`G^>(t_1,t_2)i\widehat{d}(t_1)\widehat{d}^{}(t_2),`$ (9)
$`G^<(t_1,t_2)i\widehat{d}^{}(t_2)\widehat{d}(t_1).`$ (10)
Here, the time ordering operator $`\mathrm{T}`$ arranges in chronological order and $`\stackrel{~}{\mathrm{T}}`$ is the anti time ordering operator which arranges in the reverse of chronological order. The angular brackets denote thermal average in nonequilibrium.
The present work is undertaken on the basis of the nonequilibrium perturbative formalism, Eq. ( 1 ). The retarded and advanced self-energies up to the fourth-order are formulated. Then it is confirmed that the nonequilibrium perturbative expansion can be connected with the Matsubara imaginary-time perturbative expansion for equilibrium.
### 1.2 The Kondo effect
The Kondo effect was discovered forty years ago and after that, the Kondo physics has been clarified from Landau’s Fermi liquid theory , the renormalization group , scaling , etc.. Besides, generalized Kondo problem, that of more than one channel or one impurity has been investigated. Then, the Kondo effect in electron transport through a quantum dot has been predicted theoretically at the end of 1980s and after a decade, this phenomenon has been observed. The Kondo effect has been studied theoretically using the Anderson model and the predictions have been confirmed experimentally. In the Kondo regime, the conductance has been observed to reach the unitarity limit and the Kondo temperatures estimated from observation are in excellent agreement with the expression derived using the Anderson model. Furthermore, the Kondo effect in a quantum dot has been studied for nonequilibrium system where the bias voltage is applied. The Yamada-Yosida theory, perturbation theory for equilibrium based on the Fermi liquid theory has been extended to nonequilibrium system and it has been shown that for bias voltage higher than the Kondo temperatures, the Kondo resonance disappears in the spectral function with the second-order self-energy of the Anderson model. This is in good agreement with the experiments of two terminal systems that it has been observed that the Kondo effect is suppressed when source-drain bias voltage is comparable to or exceeds the Kondo temperatures.
In the present work, using self-energies derived up to the fourth-order, the behavior of Kondo resonance is investigated for nonequilibrium state caused by bias voltage. The numerical results are still that the Kondo resonance is broken, as supported by two-terminal experiments.
## 2 Nonequilibrium Perturbation Theory
### 2.1 Formalism
A thermal average can be obtained on the basis of the nonequilibrium perturbation theory. It is assumed that we can know only the state at $`t=\mathrm{}`$, that is, initially at $`t=\mathrm{}`$, the system is equilibrium and/or noninteracting state. The perturbation is turned on at $`t=\mathrm{}`$ and introduced adiabatically and then, brought wholly into the system at $`t=0`$; around $`t=0`$, the system is regarded as stationary nonequilibrium and/or interacting state. After that, the perturbation is taken away adiabatically and vanishes at $`t=\mathrm{}`$. When the time evolution of the state is irreversible, then, the state at $`t=\mathrm{}`$ cannot be well-defined. The time evolution is therefore, performed along the real-time contour which starts and ends at $`t=\mathrm{}`$, as illustrated in Fig. 1.
S matrix is defined by
$`𝒮(t,t_0)`$ $`=`$ $`1+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n!}}\left({\displaystyle \frac{i}{\mathrm{}}}\right)^n{\displaystyle _{t_0}^t}𝑑t_1\mathrm{}{\displaystyle _{t_0}^t}𝑑t_n\mathrm{T}\left[\stackrel{~}{}_\mathrm{I}(t_1)\mathrm{}\stackrel{~}{}_\mathrm{I}(t_n)\right]`$ (11)
$`=`$ $`\mathrm{T}\left[\mathrm{exp}\left\{{\displaystyle \frac{i}{\mathrm{}}}{\displaystyle _{t_0}^t}𝑑t^{^{}}\stackrel{~}{}_\mathrm{I}(t^{^{}})\right\}\right],`$
$`𝒮(t,t_0)^{}`$ $`=`$ $`𝒮(t_0,t)=\stackrel{~}{\mathrm{T}}\left[\mathrm{exp}\left\{{\displaystyle \frac{i}{\mathrm{}}}{\displaystyle _{t_0}^t}𝑑t^{^{}}\stackrel{~}{}_\mathrm{I}(t^{^{}})\right\}\right].`$ (12)
Here $`\stackrel{~}{}_\mathrm{I}`$ is perturbation term in interaction representation.
For thermal equilibrium, the statistical operator ( density matrix ) is written in Gibbs form for the grand canonical ensemble by
$`\varrho _G={\displaystyle \frac{e^{\beta (\mu N)}}{\mathrm{Tr}e^{\beta (\mu N)}}}=e^{\beta (\mathrm{\Omega }+\mu N)}.`$ (13)
Equation ( 8 ) is not valid exactly for nonequilibrium. We have no specific limitations upon the statistical operator. The statistical operator can generally be expressed in Schrödinger representation by
$`\varrho _S(t)={\displaystyle \underset{m}{}}|m_S(t)>P_m<m_S(t)|.`$ (14)
Here, $`P_m`$ is probability that the system is in state $`m`$ and $`|m_S(t)>`$ is the state in Schrödinger representation. $`\varrho _S`$ satisfies the Liouville equation by
$`i\mathrm{}{\displaystyle \frac{\varrho _S}{t}}=[,\varrho _S].`$ (15)
The statistical operator in the interaction representation is given by $`\stackrel{~}{\varrho }(t)=e^{i_0t/\mathrm{}}\varrho _S(t)e^{i_0t/\mathrm{}}`$ and satisfies
$`i\mathrm{}{\displaystyle \frac{\stackrel{~}{\varrho }}{t}}`$ $`=`$ $`[\stackrel{~}{}_\mathrm{I},\stackrel{~}{\varrho }].`$ (16)
As a matter of course, $`\varrho _S(0)=\varrho (0)=\stackrel{~}{\varrho }(0).`$ Here $`\varrho (t)`$ is in the Heisenberg representation. The time evolution is expressed using S matrix by
$`\stackrel{~}{\varrho }(t)=S(t,t_0)\stackrel{~}{\varrho }(t_0)S(t_0,t).`$ (17)
The thermal average in the Heisenberg representation at $`t=0`$ can be obtained, for example by
$`\mathrm{T}A(t)B(t^{^{}})`$
$``$ $`\mathrm{Tr}[\varrho (0)\mathrm{T}A(t)B(t^{^{}})]`$
$`=`$ $`\mathrm{Tr}[\stackrel{~}{\varrho }(\mathrm{})𝒮(\mathrm{},0)\mathrm{T}A(t)B(t^{^{}})𝒮(0,\mathrm{})]`$
$`=`$ $`\mathrm{Tr}[\stackrel{~}{\varrho }(\mathrm{})𝒮(\mathrm{},\mathrm{})\{\mathrm{T}𝒮(\mathrm{},\mathrm{})\stackrel{~}{A}(t^{})\stackrel{~}{B}(t^{{}_{}{}^{}})\}]`$
$`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n!}}{\displaystyle \frac{1}{m!}}\left({\displaystyle \frac{i}{\mathrm{}}}\right)^n\left({\displaystyle \frac{i}{\mathrm{}}}\right)^m{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_1\mathrm{}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_n{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_1^{^{}}\mathrm{}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_m^{^{}}`$
$`\times \left\{\stackrel{~}{\mathrm{T}}\stackrel{~}{}_\mathrm{I}(t_1^+)\mathrm{}\stackrel{~}{}_\mathrm{I}(t_n^+)\right\}\left\{\mathrm{T}\stackrel{~}{}_\mathrm{I}(t_1^{{}_{}{}^{}})\mathrm{}\stackrel{~}{}_\mathrm{I}(t_m^{{}_{}{}^{}})\stackrel{~}{A}(t^{})\stackrel{~}{B}(t^{{}_{}{}^{}})\right\}_{av},`$
where $`\mathrm{}_{av}`$$`=\mathrm{Tr}[\stackrel{~}{\varrho }(\mathrm{})\mathrm{}].`$ Then, the thermal average is derived by following the ordinary procedure via the Wick’s theorem.
### 2.2 Relation of Self-Energy
After the perturbative expansion is executed, the retarded and advanced self-energies are formulated. According to the definition, the retarded and advanced Green’s functions are given by
$`G^r(t_1,t_2)i\theta (t_1t_2)\{\widehat{d}(t_1),\widehat{d}^{}(t_2)\},`$ (18)
$`G^a(t_1,t_2)i\theta (t_2t_1)\{\widehat{d}(t_1),\widehat{d}^{}(t_2)\}.`$ (19)
Here, the curly brackets denote anticommutator. The Dyson’s equations for the retarded and advanced Green’s functions are given by
$`G^r`$ $`=`$ $`g^r+g^r\mathrm{\Sigma }^rG^r,`$ (20)
$`G^a`$ $`=`$ $`g^a+g^a\mathrm{\Sigma }^aG^a.`$ (21)
As the necessity to Eqs. ( 15 ) and ( 16 ), the self-energies $`\mathrm{\Sigma }^r`$ and $`\mathrm{\Sigma }^a`$ are also required to be retarded and advanced functions in time, respectively. In accordance with the ordinary procedure of nonequilibrium perturbative formalism, there
$`𝐋=[𝐋^{}]^1={\displaystyle \frac{1}{\sqrt{2}}}\left[\begin{array}{cc}1\hfill & 1\hfill \\ 1\hfill & 1\hfill \end{array}\right],`$ (24)
and using this, then,
$`𝚺=\left[\begin{array}{cc}\mathrm{\Sigma }^{}\hfill & \mathrm{\Sigma }^<\hfill \\ \mathrm{\Sigma }^>\hfill & \mathrm{\Sigma }^{++}\hfill \end{array}\right]𝐋𝚺𝐋^{}=\left[\begin{array}{cc}\mathrm{\Omega }\hfill & \mathrm{\Sigma }^r\hfill \\ \mathrm{\Sigma }^a\hfill & 0\hfill \end{array}\right].`$ (29)
The relationship for self-energies ought to be obtained here by comparison of Eq. ( 1 ) with Eqs. ( 15 ) and ( 16 ):
$`\mathrm{\Sigma }^r(t)`$ $`=`$ $`\mathrm{\Sigma }^{}(t)+\mathrm{\Sigma }^<(t)=\mathrm{\Sigma }^{++}(t)\mathrm{\Sigma }^>(t),`$ (30)
$`\mathrm{\Sigma }^a(t)`$ $`=`$ $`\mathrm{\Sigma }^{}(t)+\mathrm{\Sigma }^>(t)=\mathrm{\Sigma }^{++}(t)\mathrm{\Sigma }^<(t).`$ (31)
## 3 Expressions of Self-Energy for Anderson model
### 3.1 Anderson Model
We consider equilibrium and nonequilibrium stationary states. Nonequilibrium state is caused by finite bias voltage, that is, the difference of chemical potentials; after bias voltage was turned on, long time has passed enough to reach stationary states. Since the states are stationary, Hamiltonian has no time dependence. The system is described by the Anderson model connected to leads. The impurity with on-site energy $`E_0`$ and the Coulomb interaction $`U`$ is connected to the left and right leads by the mixing matrix elements, $`v_L`$ and $`v_R`$. The Anderson Hamiltonian is given by
$`=`$ $`E_0{\displaystyle \underset{\sigma }{}}\widehat{n}_{d\sigma }+\mu _L{\displaystyle \underset{\sigma }{}}\widehat{n}_{L\sigma }+\mu _R{\displaystyle \underset{\sigma }{}}\widehat{n}_{R\sigma }+U\left(\widehat{n}_d\widehat{n}_d\right)\left(\widehat{n}_d\widehat{n}_d\right)`$ (32)
$`{\displaystyle \underset{\sigma }{}}v_L(\widehat{d}_\sigma ^{}\widehat{c}_{L\sigma }+\mathrm{H}.\mathrm{c}.){\displaystyle \underset{\sigma }{}}v_R(\widehat{d}_\sigma ^{}\widehat{c}_{R\sigma }+\mathrm{H}.\mathrm{c}.).`$
$`\widehat{d}^{}`$ ($`\widehat{d}`$) is creation (annihilation) operator for electron on the impurity, and $`\widehat{c}_L^{}`$ and $`\widehat{c}_R^{}`$ ($`\widehat{c}_L`$ and $`\widehat{c}_R`$) are creation (annihilation) operators in the left and right leads, respectively. $`\sigma `$ is index for spin. The chemical potentials in the isolated left and right leads are $`\mu _L`$ and $`\mu _R`$, respectively. The applied voltage is, therefore defined by $`eV\mu _L\mu _R`$.
We consider that the band-width of left and right leads is large infinitely, so that the coupling functions, $`\mathrm{\Gamma }_L`$ and $`\mathrm{\Gamma }_R`$ can be taken to be independent of energy, $`E`$. On-site energy $`E_0`$ is set canceling with the Hartree term, i.e. the first-order contribution to self-energy for electron correlation: $`\mathrm{\Sigma }_\sigma ^{r(1)}(E)`$$`=\mathrm{\Sigma }_\sigma ^{a(1)}(E)`$ $`=Un_\sigma `$.
Accordingly, the Fourier components of the noninteracting ( unperturbed ) Green’s functions reduce to
$`g^r(E)={\displaystyle \frac{1}{E+i\mathrm{\Gamma },}}`$ (33)
$`g^a(E)={\displaystyle \frac{1}{Ei\mathrm{\Gamma },}}`$ (34)
where $`\mathrm{\Gamma }=(\mathrm{\Gamma }_L+\mathrm{\Gamma }_R)/2`$. Hence, the inverse Fourier components can be written by
$`g^r(t)`$ $`=`$ $`i\theta (t)e^{\mathrm{\Gamma }t},`$
$`g^a(t)`$ $`=`$ $`i\theta (t)e^{\mathrm{\Gamma }t}.`$
In addition, from Eq. ( 1 ), we have
$`g^<(E)=g^r(E)\left[if_L(E)\mathrm{\Gamma }_L+if_R(E)\mathrm{\Gamma }_R\right]g^a(E),`$ (35)
$`g^>(E)=g^r(E)\left[i(f_L(E)1)\mathrm{\Gamma }_L+i(f_R(E)1)\mathrm{\Gamma }_R\right]g^a(E).`$ (36)
$`f_L`$ and $`f_R`$ are the Fermi distribution functions in the isolated left and right leads, respectively. By Eqs. ( 22 ) and ( 23 ), the nonequilibrium state is introduced as the superposition of the left and right leads. Then, the effective Fermi distribution function can be expressed by
$`f_{\mathrm{eff}}(E)={\displaystyle \frac{f_L(E)\mathrm{\Gamma }_L+f_R(E)\mathrm{\Gamma }_R}{\mathrm{\Gamma }_L+\mathrm{\Gamma }_R}}.`$ (37)
### 3.2 Self-Energy
The retarded and advanced self-energies are derived up to the forth-order. Equations ( 17 ) and ( 18 ) are divided into retarded and advanced terms in time:
$`\mathrm{\Sigma }^r(t)`$ $`=`$ $`[\mathrm{\Sigma }^{}(t)+\mathrm{\Sigma }^<(t)]\theta (t)+[\mathrm{\Sigma }^{}(t)+\mathrm{\Sigma }^<(t)]\theta (t)`$
$`=`$ $`[\mathrm{\Sigma }^{++}(t)+\mathrm{\Sigma }^>(t)]\theta (t)[\mathrm{\Sigma }^{++}(t)+\mathrm{\Sigma }^>(t)]\theta (t),`$
$`\mathrm{\Sigma }^a(t)`$ $`=`$ $`[\mathrm{\Sigma }^{}(t)+\mathrm{\Sigma }^>(t)]\theta (t)+[\mathrm{\Sigma }^{}(t)+\mathrm{\Sigma }^>(t)]\theta (t)`$
$`=`$ $`[\mathrm{\Sigma }^{++}(t)+\mathrm{\Sigma }^<(t)]\theta (t)[\mathrm{\Sigma }^{++}(t)+\mathrm{\Sigma }^<(t)]\theta (t).`$
Then it is found that for the self-energies drawn using the Wick’s theorem,
$`\mathrm{\Sigma }^{}(t)\theta (t)=\mathrm{\Sigma }^>(t)\theta (t),`$ $`\mathrm{\Sigma }^{++}(t)\theta (t)=\mathrm{\Sigma }^<(t)\theta (t),`$
$`\mathrm{\Sigma }^{}(t)\theta (t)=\mathrm{\Sigma }^<(t)\theta (t),`$ $`\mathrm{\Sigma }^{++}(t)\theta (t)=\mathrm{\Sigma }^>(t)\theta (t).`$
It leads to
$`[\mathrm{\Sigma }^{}(t)+\mathrm{\Sigma }^<(t)]\theta (t)=[\mathrm{\Sigma }^{++}(t)+\mathrm{\Sigma }^>(t)]\theta (t)=0,`$
$`[\mathrm{\Sigma }^{}(t)+\mathrm{\Sigma }^>(t)]\theta (t)=[\mathrm{\Sigma }^{++}(t)+\mathrm{\Sigma }^<(t)]\theta (t)=0.`$
As a consequence, the retarded and advanced self-energies are obtained as retarded and advanced functions of time, respectively:
$`\mathrm{\Sigma }^r(t)`$ $`=`$ $`[\mathrm{\Sigma }^{}(t)+\mathrm{\Sigma }^<(t)]\theta (t)=[\mathrm{\Sigma }^{++}(t)+\mathrm{\Sigma }^>(t)]\theta (t),`$
$`\mathrm{\Sigma }^a(t)`$ $`=`$ $`[\mathrm{\Sigma }^{}(t)+\mathrm{\Sigma }^>(t)]\theta (t)=[\mathrm{\Sigma }^{++}(t)+\mathrm{\Sigma }^<(t)]\theta (t).`$
In addition, it proves
$`\mathrm{\Sigma }^r(t)`$ $`=`$ $`[\mathrm{\Sigma }^<(t)\mathrm{\Sigma }^>(t)]\theta (t),`$
$`\mathrm{\Sigma }^a(t)`$ $`=`$ $`[\mathrm{\Sigma }^>(t)\mathrm{\Sigma }^<(t)]\theta (t).`$
Hence the following relations still stand: $`\mathrm{\Sigma }^r\mathrm{\Sigma }^a=\mathrm{\Sigma }^<\mathrm{\Sigma }^>`$ and furthermore
$`G^<=(1+G^r\mathrm{\Sigma }^r)g^<(1+G^a\mathrm{\Sigma }^a)G^r\mathrm{\Sigma }^<G^a,`$
$`G^>=(1+G^r\mathrm{\Sigma }^r)g^>(1+G^a\mathrm{\Sigma }^a)G^r\mathrm{\Sigma }^>G^a.`$
As the results, the second-order self-energy is written by
$`\mathrm{\Sigma }^{r(2)}(E)=U^2{\displaystyle _0^{\mathrm{}}}𝑑t_1e^{iEt_1}\left[\begin{array}{cccc}g^>(t_1)g^>(t_1)g^<(t_1)\hfill & & & \\ g^<(t_1)g^<(t_1)g^>(t_1)\hfill & & & \end{array}\right]`$ (40)
$`=U^2{\displaystyle _0^{\mathrm{}}}𝑑t_1e^{iEt_1}\left[\begin{array}{cccc}g^\pm (t_1)g^>(t_1)g^<(t_1)\hfill & & & \\ +g^<(t_1)g^\pm (t_1)g^>(t_1)\hfill & & & \\ +g^<(t_1)g^>(t_1)g^\pm (t_1)\hfill & & & \end{array}\right],`$ (44)
$`\mathrm{\Sigma }^{a(2)}(E)=U^2{\displaystyle _{\mathrm{}}^0}𝑑t_1e^{iEt_1}\left[\begin{array}{cccc}g^<(t_1)g^<(t_1)g^>(t_1)\hfill & & & \\ g^>(t_1)g^>(t_1)g^<(t_1)\hfill & & & \end{array}\right]`$ (47)
$`=U^2{\displaystyle _{\mathrm{}}^0}𝑑t_1e^{iEt_1}\left[\begin{array}{cccc}g^\pm (t_1)g^>(t_1)g^<(t_1)\hfill & & & \\ +g^<(t_1)g^\pm (t_1)g^>(t_1)\hfill & & & \\ +g^<(t_1)g^>(t_1)g^\pm (t_1)\hfill & & & \end{array}\right].`$ (51)
Here $`g^\pm (t)=g^r(t)+g^a(t)`$, that is, $`g^+(t)=g^r(t)=i\theta (t)e^{\mathrm{\Gamma }t}`$ for $`t0`$ and $`g^{}(t)=g^a(t)=i\theta (t)e^{\mathrm{\Gamma }t}`$ for $`t<0`$. Additionally, $`g^<(t)`$ and $`g^>(t)`$ are the inverse Fourier components of Eqs. ( 22 ) and ( 23 ). Figure 2 shows the diagram for the second-order self-energy. As shown numerically later, the second-order contribution coincide with those derived by Hershfield et al. . In the symmetric equilibrium case, the asymptotic behavior at low energy is expressed by
$`\mathrm{\Sigma }^{r(2)}(E)\mathrm{\Gamma }\left(3{\displaystyle \frac{\pi ^2}{4}}\right)\left({\displaystyle \frac{U}{\pi \mathrm{\Gamma }}}\right)^2{\displaystyle \frac{E}{\mathrm{\Gamma }}}i{\displaystyle \frac{\mathrm{\Gamma }}{2}}\left({\displaystyle \frac{U}{\pi \mathrm{\Gamma }}}\right)^2\left({\displaystyle \frac{E}{\mathrm{\Gamma }}}\right)^2,`$ (52)
the exact results based on the Bethe ansatz method.
The third-order terms corresponding to the diagram in Fig. 3(a) are expressed by
$`\mathrm{\Sigma }_{pp}^{r(3)}(E)`$ $`=`$ $`U^3{\displaystyle _0^{\mathrm{}}}𝑑t_1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_2e^{iEt_1}\left[\begin{array}{cccc}g^<(t_1)g^>(t_1t_2)g^>(t_1t_2)\hfill & & & \\ g^>(t_1)g^<(t_1t_2)g^<(t_1t_2)\hfill & & & \end{array}\right]`$ (55)
$`\times \left[\begin{array}{cccc}g^\pm (t_2)g^>(t_2)+g^<(t_2)g^\pm (t_2)\hfill & & & \end{array}\right],`$
$`\mathrm{\Sigma }_{pp}^{a(3)}(E)`$ $`=`$ $`U^3{\displaystyle _{\mathrm{}}^0}𝑑t_1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_2e^{iEt_1}\left[\begin{array}{cccc}g^>(t_1)g^<(t_1t_2)g^<(t_1t_2)\hfill & & & \\ g^<(t_1)g^>(t_1t_2)g^>(t_1t_2)\hfill & & & \end{array}\right]`$ (58)
$`\times \left[\begin{array}{cccc}g^\pm (t_2)g^>(t_2)+g^<(t_2)g^\pm (t_2)\hfill & & & \end{array}\right].`$
Figure 3(b) illustrates the diagram for the following terms:
$`\mathrm{\Sigma }_{ph}^{r(3)}(E)`$ $`=`$ $`U^3{\displaystyle _0^{\mathrm{}}}𝑑t_1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_2e^{iEt_1}\left[\begin{array}{cccc}g^>(t_1)g^>(t_1t_2)g^<(t_2t_1)\hfill & & & \\ g^<(t_1)g^<(t_1t_2)g^>(t_2t_1)\hfill & & & \end{array}\right]`$ (61)
$`\times \left[\begin{array}{cccc}g^\pm (t_2)g^<(t_2)+g^<(t_2)g^\pm (t_2)\hfill & & & \end{array}\right],`$
$`\mathrm{\Sigma }_{ph}^{a(3)}(E)`$ $`=`$ $`U^3{\displaystyle _{\mathrm{}}^0}𝑑t_1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_2e^{iEt_1}\left[\begin{array}{cccc}g^<(t_1)g^<(t_1t_2)g^>(t_2t_1)\hfill & & & \\ g^>(t_1)g^>(t_1t_2)g^<(t_2t_1)\hfill & & & \end{array}\right]`$ (64)
$`\times \left[\begin{array}{cccc}g^\pm (t_2)g^<(t_2)+g^<(t_2)g^\pm (t_2)\hfill & & & \end{array}\right].`$
Equations ( 28 )-( 31 ) for equilibrium agree exactly with those derived from the Matsubara imaginary-time perturbative expansion for equilibrium and analytical continuity by Zlatić et al. As mentioned later, it is numerically confirmed that the third-order contribution vanishes for the symmetric Anderson model; this is in good agreement with both the results derived from the Yamada-Yosida theory and obtained on the basis of the Bethe ansatz method.
Furthermore, the fourth-order contribution to the self-energy is formulated. ( See Appendix. ) The twelve terms for the proper fourth-order self-energy can be divided into four groups, each of which comprises three terms. The four groups correspond to the diagrams denoted in Figs. 4 (a)-(c), Figs. 4 (d)-(f), Figs. 4 (g)-(i), and Figs. 4 (j)-(l), respectively. For symmetric Anderson model at equilibrium, the asymptotic behavior at low energy is approximately in agreement with those based on the Bethe ansatz method :
$`\mathrm{\Sigma }^{r(4)}(E)\mathrm{\Gamma }\left(105{\displaystyle \frac{45\pi ^2}{4}}+{\displaystyle \frac{\pi ^4}{16}}\right)\left({\displaystyle \frac{U}{\pi \mathrm{\Gamma }}}\right)^4{\displaystyle \frac{E}{\mathrm{\Gamma }}}i{\displaystyle \frac{\mathrm{\Gamma }}{2}}\left(303\pi ^2\right)\left({\displaystyle \frac{U}{\pi \mathrm{\Gamma }}}\right)^4\left({\displaystyle \frac{E}{\mathrm{\Gamma }}}\right)^2.`$
## 4 Numerical Results and Discussion
### 4.1 Self-Energy
The third-order terms, Eqs. ( 28 )-( 31 ) cancel under electron-hole symmetry not only at equilibrium but also at nonequilibrium: $`\mathrm{\Sigma }_{ph}^{r(3)}(E)=\mathrm{\Sigma }_{pp}^{r(3)}(E)`$ and $`\mathrm{\Sigma }_{ph}^{a(3)}(E)=\mathrm{\Sigma }_{pp}^{a(3)}(E)`$. As a consequence, the third-order contribution to self-energy vanishes in the symmetric case. It agrees with the results of Refs. based on the Yamada-Yosida theory that all odd-order contributions except the Hartree term vanish for equilibrium in the symmetric single-impurity Anderson model; probably, it is just the same with nonequilibrium state. On the other hand, the third-order terms contribute to the asymmetric system where electron-hole symmetry breaks and furthermore, the third-order terms for spin-up and for spin-down contribute respectively when the spin degeneracy is lifted for example, by magnetic field. For the fourth-order contribution, three terms which constitute each of four groups contribute equivalently under electron-hole symmetry. Moreover, to the asymmetric system, the terms brought by the diagrams of Figs. 4(a) and 4(b) contribute equivalently and the terms by the diagrams of Figs. 4(j) and 4(k) make equivalent contribution, and the rest, the eight terms contribute respectively. Further, the twenty-four terms for spin-up and spin-down take effect severally in the presence of magnetic field.
The second-order and the fourth-order contributions to self-energy for zero temperature symmetric Anderson model are shown in Figs. 5(a) and 5(b) and in Figs. 6(a) and 6(b), respectively. Equation ( 27 ) represents the curves around $`E=0`$ denoted by solid line in Figs. 5(a) and 5(b), respectively, and Equation ( 32 ) represents approximately those shown in Figs. 6(a) and 6(b), respectively. The curves of the second-order self-energy shown in Figs. 5(a) and 5(b) are identical with those of expressions derived by Hershfield et al. . In comparison of Figs. 6(a) and 6(b) with Figs. 5(a) and 5(b), it is found that the fourth-order contribution for equilibrium has the same but narrow curves at low energy with those of the second-order contribution. In addition, the broad curves are attached at high energy for the fourth-order self-energy. ( The higher-order contribution is, the more the curves must oscillate as a function of energy. ) When the voltage, $`eV/\mathrm{\Gamma }`$ exceeds $`2.0`$, the behavior of curves of self-energy changes distinctly and comes to present striking contrasts to that for the second-order contribution. Especially, the curve for the imaginary part of the fourth-order contribution rises up with maximum at $`E=0`$. On the other hand, for the second-order contribution, a valley appears with minimum at energy of zero$``$it is quite the contrary. Moreover, from these results, it is expected that the sixth-order contribution to imaginary part of self-energy has minimum at $`E=0`$. Because of these, the perturbative expansion is hard to converge for $`eV/\mathrm{\Gamma }>2.0`$, as mentioned later.
Besides, the current conservation is mentioned. In Ref. , it is shown that the continuity of current entering and leaving the impurity stands exactly at any strength of $`U`$ within the approximation up to the second-order for the symmetric single-impurity Anderson model. In comparison of Figs. 6(a) and 6(b) with Figs. 5(a) and 5(b), it is found that curves of fourth-order self-energy have the symmetry similar to those of the second-order. From this, it is anticipated that the current conservation are satisfied perfectly with approximation up to the fourth-order in the single-impurity system where electron-hole symmetry holds. The continuity of current can be maintained perfectly in single-impurity system as far as electron-hole symmetry stands. On the other hand, current comes to fail to be conserved with increasing $`U`$ in asymmetric single-impurity case and in two-impurity case.
### 4.2 Spectral Function
The spectral function with the second-order self-energy is generally known. It is plotted for $`U/\mathrm{\Gamma }=10.0`$ and zero temperature in Fig. 7. For equilibrium, the Kondo peak at energy of zero is very sharp and the two-side broad peaks appear at $`E`$ $`\pm U/2`$. The curve for $`eV=0`$ is identical with that shown in Ref. . As $`eV`$ becomes higher than the Kondo temperatures, $`k_BT_K`$ , the Kondo peak becomes lower and finally vanishes, while the two-side broad peaks rise at $`E`$ $`\pm U/2`$
Figure 8 shows the spectral function with the self-energy up to the fourth-order for equilibrium and zero temperature. With strengthening $`U`$, two-side narrow peaks come to occur in the vicinity of $`E=`$ $`\pm U/2`$ in addition to the Kondo peak. At $`U`$ large enough, the Kondo peak becomes very acute and two-side narrow peaks rise higher and sharpen; the energy levels for the atomic limit are produced distinctly. The fourth-order self-energy has the same but narrow curves as functions of energy with those of the second-order and those curves make the peaks at $`E=`$ $`\pm U/2`$.
For the present approximation up to the fourth-order, the Kondo peak at $`E=0`$ reaches the unitarity limit and the charge, $`n`$ corresponds to $`1/2`$, that is, the Friedel sum rule is correctly satisfied:
$`\rho (E_f)=\mathrm{sin}^2(\pi n)/\pi \mathrm{\Gamma },`$ (66)
where $`\rho (E_f)`$ is the local density of states at the Fermi energy. Here, the discussions should be made on the ranges of $`U`$ in which the present approximation up to the fourth-order stands. From the results, it is found that the approximation within the fourth-order holds up to $`U/\mathrm{\Gamma }`$ $`5.0`$ and is beyond the validity for $`U/\mathrm{\Gamma }`$$`>6.0`$. In addition, the curve for imaginary part of the fourth-order contribution is positive partly, as shown in Fig. 6(b) and as a consequence, the curve of the spectral function becomes negative partly for too large $`U`$. In such a case, the present approximation is out of validity and the higher-order terms are required.
Next, the results for nonequilibrium and zero temperature are shown. The expression for the Friedel sum rule, Eq. ( 33 ) does not stand for nonequilibrium, since the charge cannot be expressed with respect to the local density of states. All the same, the Kondo peak reaches the unitarity limit and $`n=1/2`$ in the symmetric and noninteracting case. The spectral functions with the self-energy up to the fourth-order are plotted for $`eV/\mathrm{\Gamma }=0.5`$ and $`eV/\mathrm{\Gamma }=1.0`$ in Figs. 9, respectively. When $`U`$ is strengthened and $`eV`$ exceeds $`k_BT_K`$ ( approximately, $`k_BT_K/`$$`\mathrm{\Gamma }`$ $``$$`0.5`$ for $`U/`$$`\mathrm{\Gamma }`$ $`=3.5`$ and $`k_BT_K/`$$`\mathrm{\Gamma }`$ $``$$`0.3`$ for $`U/`$$`\mathrm{\Gamma }`$$`=5.0`$ ), the Kondo peak for $`eV/\mathrm{\Gamma }=0.5`$ falls in and instead, the two-side narrow peaks remain to sharpen in the vicinity of $`E=`$ $`\pm U/2`$. For $`eV/\mathrm{\Gamma }=1.0`$, the Kondo peak becomes broad and disappears for $`U`$ large enough. The two-side peaks is generated small in the vicinity of $`E=`$ $`\pm U/2`$. The Kondo resonance is quite broken for bias voltage exceeding the Kondo temperatures; this accords with the experimental results of two terminal systems that the Kondo effect is suppressed when source-drain bias voltage is comparable to or exceeds the Kondo temperatures, $`eVk_BT_K`$ For $`eV/\mathrm{\Gamma }>2.0`$, the Kondo peak does not lower even when $`eV`$ is much larger than $`k_BT_K`$. The perturbative expansion is hard to converge on account of the imaginary part of the self-energy for $`eV/\mathrm{\Gamma }>2.0`$, as described before; thereby, the higher-order contribution to self-energy is probably needed for high voltage.
In the present work, nonequilibrium state is represented as the superposition of the two leads and the effective Fermi distribution function, Eq. ( 24 ) is qualitatively similar to that for finite temperatures. From the analogy in the Fermi distribution function, it is inferred that there are nonequilibrium fluctuations similar to thermal fluctuations. Because of the effective Fermi distribution function, not only for the second-order but also for the fourth-order, the Kondo resonance is destroyed, qualitatively the same as for finite temperatures. In contrast, if the finite voltage state is expressed as two localized states, the numerical results of the Kondo peak splitting can be obtained. All the same, for two terminal systems, the Kondo resonance splitting may not take place for finite bias voltage.
Summary: The present work is based on the nonequilibrium perturbative formalism. Here the self-energies are derived and then it is indicated that the nonequilibrium ( real-time ) perturbative expansion can be related to the Matsubara imaginary-time perturbative expansion for equilibrium. As the numerical results, the Kondo peak disappears as bias voltage exceeding the Kondo temperatures. Because of the analogy of the effective Fermi distribution function for nonequilibrium with that for finite temperatures, the present result is qualitatively similar to that for finite temperatures. This characteristic appears in the experiments of two terminal systems.
## Acknowledgements
The numerical calculations were executed at the Yukawa Institute Computer Facility. The multiple integrals were performed using the computer subroutine, MQFSRD of NUMPAC.
## Appendix
The twelve terms for the fourth-order contribution can be divided into four groups, each of which is composed of three terms. The four groups are brought from diagrams denoted in Figs. 4 (a)-(c), Figs. 4 (d)-(f), Figs. 4 (g)-(i), and Figs. 4 (j)-(l), respectively. The terms for the diagrams illustrated in Figs. 4(a) and 4(b) are equivalent except for the spin indices and expressed by
$`\mathrm{\Sigma }_{a,b}^{r(4)}(E)`$ $`=`$ $`U^4{\displaystyle _0^{\mathrm{}}}𝑑t_1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_2{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_3e^{iEt_1}`$
$`\times \left[\begin{array}{cccc}g^<(t_1)g^<(t_1t_2t_3)g^>(t_1+t_2+t_3)\hfill & & & \\ g^>(t_1)g^>(t_1t_2t_3)g^<(t_1+t_2+t_3)\hfill & & & \end{array}\right]`$
$`\times \left[\begin{array}{cccc}g^\pm (t_2)g^<(t_2)+g^<(t_2)g^\pm (t_2)\hfill & & & \end{array}\right]`$ (69)
$`\times \left[\begin{array}{cccc}g^\pm (t_3)g^<(t_3)+g^<(t_3)g^\pm (t_3)\hfill & & & \end{array}\right],`$ (71)
$`\mathrm{\Sigma }_{a,b}^{a(4)}(E)`$ $`=`$ $`U^4{\displaystyle _{\mathrm{}}^0}𝑑t_1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_2{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_3e^{iEt_1}`$
$`\times \left[\begin{array}{cccc}g^>(t_1)g^>(t_1t_2t_3)g^<(t_1+t_2+t_3)\hfill & & & \\ g^<(t_1)g^<(t_1t_2t_3)g^>(t_1+t_2+t_3)\hfill & & & \end{array}\right]`$
$`\times \left[\begin{array}{cccc}g^\pm (t_2)g^<(t_2)+g^<(t_2)g^\pm (t_2)\hfill & & & \end{array}\right]`$ (74)
$`\times \left[\begin{array}{cccc}g^\pm (t_3)g^<(t_3)+g^<(t_3)g^\pm (t_3)\hfill & & & \end{array}\right].`$ (76)
Additionally, Figure 4(c) shows the diagram for the following terms:
$`\mathrm{\Sigma }_c^{r(4)}(E)`$ $`=`$ $`U^4{\displaystyle _0^{\mathrm{}}}𝑑t_1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_2{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_3e^{iEt_1}`$
$`\times \left[\begin{array}{cccc}g^>(t_1)g^<(t_1t_2t_3)g^<(t_1t_2t_3)\hfill & & & \\ g^<(t_1)g^>(t_1t_2t_3)g^>(t_1t_2t_3)\hfill & & & \end{array}\right]`$
$`\times \left[\begin{array}{cccc}g^\pm (t_2)g^>(t_2)+g^<(t_2)g^\pm (t_2)\hfill & & & \end{array}\right]`$ (79)
$`\times \left[\begin{array}{cccc}g^\pm (t_3)g^>(t_3)+g^<(t_3)g^\pm (t_3)\hfill & & & \end{array}\right],`$ (81)
$`\mathrm{\Sigma }_c^{a(4)}(E)`$ $`=`$ $`U^4{\displaystyle _{\mathrm{}}^0}𝑑t_1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_2{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_3e^{iEt_1}`$
$`\times \left[\begin{array}{cccc}g^<(t_1)g^>(t_1t_2t_3)g^>(t_1t_2t_3)\hfill & & & \\ g^>(t_1)g^<(t_1t_2t_3)g^<(t_1t_2t_3)\hfill & & & \end{array}\right]`$
$`\times \left[\begin{array}{cccc}g^\pm (t_2)g^>(t_2)+g^<(t_2)g^\pm (t_2)\hfill & & & \end{array}\right]`$ (84)
$`\times \left[\begin{array}{cccc}g^\pm (t_3)g^>(t_3)+g^<(t_3)g^\pm (t_3)\hfill & & & \end{array}\right].`$ (86)
Next, the terms brought from diagram in Fig. 4(d) are expressed by
$`\mathrm{\Sigma }_d^{r(4)}(E)`$ $`=`$ $`U^4{\displaystyle _0^{\mathrm{}}}𝑑t_1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_2{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_3e^{iEt_1}`$ (90)
$`\times \left[\begin{array}{cccc}g^>(t_1t_3)g^>(t_1t_2)g^<(t_2t_1)\hfill & & & \\ g^<(t_1t_3)g^<(t_1t_2)g^>(t_2t_1)\hfill & & & \end{array}\right]`$
$`\times g^\pm (t_2)\mathrm{sgn}(t_3)\left[\begin{array}{cccc}g^>(t_2+t_3)g^>(t_3)g^<(t_3)\hfill & & & \\ g^<(t_2+t_3)g^<(t_3)g^>(t_3)\hfill & & & \end{array}\right],`$
$`\mathrm{\Sigma }_d^{a(4)}(E)`$ $`=`$ $`U^4{\displaystyle _{\mathrm{}}^0}𝑑t_1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_2{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_3e^{iEt_1}`$ (94)
$`\times \left[\begin{array}{cccc}g^<(t_1t_3)g^<(t_1t_2)g^>(t_2t_1)\hfill & & & \\ g^>(t_1t_3)g^>(t_1t_2)g^<(t_2t_1)\hfill & & & \end{array}\right]`$
$`\times g^\pm (t_2)\mathrm{sgn}(t_3)\left[\begin{array}{cccc}g^>(t_2+t_3)g^>(t_3)g^<(t_3)\hfill & & & \\ g^<(t_2+t_3)g^<(t_3)g^>(t_3)\hfill & & & \end{array}\right].`$
The terms for diagram in Fig. 4(e) are written by
$`\mathrm{\Sigma }_e^{r(4)}(E)`$ $`=`$ $`U^4{\displaystyle _0^{\mathrm{}}}𝑑t_1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_2{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_3e^{iEt_1}`$ (98)
$`\times \left[\begin{array}{cccc}g^>(t_1t_2)g^>(t_1t_2)g^<(t_3t_1)\hfill & & & \\ g^<(t_1t_2)g^<(t_1t_2)g^>(t_3t_1)\hfill & & & \end{array}\right]`$
$`\times g^\pm (t_2)\mathrm{sgn}(t_3)\left[\begin{array}{cccc}g^>(t_2t_3)g^>(t_3)g^<(t_3)\hfill & & & \\ g^<(t_2t_3)g^<(t_3)g^>(t_3)\hfill & & & \end{array}\right],`$
$`\mathrm{\Sigma }_e^{a(4)}(E)`$ $`=`$ $`U^4{\displaystyle _{\mathrm{}}^0}𝑑t_1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_2{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_3e^{iEt_1}`$ (102)
$`\times \left[\begin{array}{cccc}g^<(t_1t_2)g^<(t_1t_2)g^>(t_3t_1)\hfill & & & \\ g^>(t_1t_2)g^>(t_1t_2)g^<(t_3t_1)\hfill & & & \end{array}\right]`$
$`\times g^\pm (t_2)\mathrm{sgn}(t_3)\left[\begin{array}{cccc}g^>(t_2t_3)g^>(t_3)g^<(t_3)\hfill & & & \\ g^<(t_2t_3)g^<(t_3)g^>(t_3)\hfill & & & \end{array}\right].`$
In addition, Figure 4(f) denotes the diagram for the following terms:
$`\mathrm{\Sigma }_f^{r(4)}(E)`$ $`=`$ $`U^4{\displaystyle _0^{\mathrm{}}}𝑑t_1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_2{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_3e^{iEt_1}`$ (106)
$`\times \left[\begin{array}{cccc}g^>(t_1t_3)g^>(t_1t_2)g^<(t_2t_1)\hfill & & & \\ g^<(t_1t_3)g^<(t_1t_2)g^>(t_2t_1)\hfill & & & \end{array}\right]`$
$`\times g^\pm (t_2)\mathrm{sgn}(t_3)\left[\begin{array}{cccc}g^<(t_3)g^<(t_3)g^>(t_2t_3)\hfill & & & \\ g^>(t_3)g^>(t_3)g^<(t_2t_3)\hfill & & & \end{array}\right],`$
$`\mathrm{\Sigma }_f^{a(4)}(E)`$ $`=`$ $`U^4{\displaystyle _{\mathrm{}}^0}𝑑t_1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_2{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_3e^{iEt_1}`$ (110)
$`\times \left[\begin{array}{cccc}g^<(t_1t_3)g^<(t_1t_2)g^>(t_2t_1)\hfill & & & \\ g^>(t_1t_3)g^>(t_1t_2)g^<(t_2t_1)\hfill & & & \end{array}\right]`$
$`\times g^\pm (t_2)\mathrm{sgn}(t_3)\left[\begin{array}{cccc}g^<(t_3)g^<(t_3)g^>(t_2t_3)\hfill & & & \\ g^>(t_3)g^>(t_3)g^<(t_2t_3)\hfill & & & \end{array}\right].`$
Next, the terms formulated from diagram illustrated in Fig. 4(g) are expressed by
$`\mathrm{\Sigma }_g^{r(4)}(E)`$ $`=`$ $`U^4{\displaystyle _0^{\mathrm{}}}𝑑t_1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_2{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_3e^{iEt_1}`$
$`\times \left[\begin{array}{cccc}g^>(t_1)g^>(t_1t_2t_3)g^<(t_2t_1)\hfill & & & \\ g^<(t_1)g^<(t_1t_2t_3)g^>(t_2t_1)\hfill & & & \end{array}\right]`$
$`\times g^\pm (t_2)\mathrm{sgn}(t_3)\left[\begin{array}{cccc}g^>(t_2+t_3)g^>(t_3)g^<(t_3)\hfill & & & \\ g^<(t_2+t_3)g^<(t_3)g^>(t_3)\hfill & & & \end{array}\right],`$ (114)
$`\mathrm{\Sigma }_g^{a(4)}(E)`$ $`=`$ $`U^4{\displaystyle _{\mathrm{}}^0}𝑑t_1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_2{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_3e^{iEt_1}`$
$`\times \left[\begin{array}{cccc}g^<(t_1)g^<(t_1t_2t_3)g^>(t_2t_1)\hfill & & & \\ g^>(t_1)g^>(t_1t_2t_3)g^<(t_2t_1)\hfill & & & \end{array}\right]`$
$`\times g^\pm (t_2)\mathrm{sgn}(t_3)\left[\begin{array}{cccc}g^>(t_2+t_3)g^>(t_3)g^<(t_3)\hfill & & & \\ g^<(t_2+t_3)g^<(t_3)g^>(t_3)\hfill & & & \end{array}\right].`$ (118)
Figure 4(h) illustrates the diagram for the following terms:
$`\mathrm{\Sigma }_h^{r(4)}(E)`$ $`=`$ $`U^4{\displaystyle _0^{\mathrm{}}}𝑑t_1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_2{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_3e^{iEt_1}`$
$`\times \left[\begin{array}{cccc}g^<(t_1)g^<(t_1t_2t_3)g^>(t_2t_1)\hfill & & & \\ g^>(t_1)g^>(t_1t_2t_3)g^<(t_2t_1)\hfill & & & \end{array}\right]`$
$`\times g^\pm (t_2)\mathrm{sgn}(t_3)\left[\begin{array}{cccc}g^>(t_3)g^>(t_3)g^<(t_2t_3)\hfill & & & \\ g^<(t_3)g^<(t_3)g^>(t_2t_3)\hfill & & & \end{array}\right],`$ (122)
$`\mathrm{\Sigma }_h^{a(4)}(E)`$ $`=`$ $`U^4{\displaystyle _{\mathrm{}}^0}𝑑t_1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_2{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_3e^{iEt_1}`$
$`\times \left[\begin{array}{cccc}g^>(t_1)g^>(t_1t_2t_3)g^<(t_2t_1)\hfill & & & \\ g^<(t_1)g^<(t_1t_2t_3)g^>(t_2t_1)\hfill & & & \end{array}\right]`$
$`\times g^\pm (t_2)\mathrm{sgn}(t_3)\left[\begin{array}{cccc}g^>(t_3)g^>(t_3)g^<(t_2t_3)\hfill & & & \\ g^<(t_3)g^<(t_3)g^>(t_2t_3)\hfill & & & \end{array}\right].`$ (126)
Besides, the terms formulated from the diagram in Fig. 4(i) are written by
$`\mathrm{\Sigma }_i^{r(4)}(E)`$ $`=`$ $`U^4{\displaystyle _0^{\mathrm{}}}𝑑t_1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_2{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_3e^{iEt_1}`$
$`\times \left[\begin{array}{cccc}g^>(t_1)g^<(t_1t_2t_3)g^<(t_1t_2)\hfill & & & \\ g^<(t_1)g^>(t_1t_2t_3)g^>(t_1t_2)\hfill & & & \end{array}\right]`$
$`\times g^\pm (t_2)\mathrm{sgn}(t_3)\left[\begin{array}{cccc}g^>(t_2+t_3)g^>(t_3)g^<(t_3)\hfill & & & \\ g^<(t_2+t_3)g^<(t_3)g^>(t_3)\hfill & & & \end{array}\right],`$ (130)
$`\mathrm{\Sigma }_i^{a(4)}(E)`$ $`=`$ $`U^4{\displaystyle _{\mathrm{}}^0}𝑑t_1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_2{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_3e^{iEt_1}`$
$`\times \left[\begin{array}{cccc}g^<(t_1)g^>(t_1t_2t_3)g^>(t_1t_2)\hfill & & & \\ g^>(t_1)g^<(t_1t_2t_3)g^<(t_1t_2)\hfill & & & \end{array}\right]`$
$`\times g^\pm (t_2)\mathrm{sgn}(t_3)\left[\begin{array}{cccc}g^>(t_2+t_3)g^>(t_3)g^<(t_3)\hfill & & & \\ g^<(t_2+t_3)g^<(t_3)g^>(t_3)\hfill & & & \end{array}\right].`$ (134)
Next, the terms for diagrams denoted in Figs. 4 (j) and 4(k) are equivalent except for the spin indices and written by
$`\mathrm{\Sigma }_{j,k}^{r(4)}(E)`$ $`=`$ $`U^4{\displaystyle _0^{\mathrm{}}}𝑑t_1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_2{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_3e^{iEt_1}`$
$`\times \left[\begin{array}{cccc}g^>(t_1)g^<(t_1)g^>(t_1t_2t_3)\hfill & & & \\ g^<(t_1)g^>(t_1)g^<(t_1t_2t_3)\hfill & & & \end{array}\right]`$
$`\times g^\pm (t_2)\left[\begin{array}{cccc}g^\pm (t_3)g^>(t_3)g^<(t_3)\hfill & & & \\ +g^<(t_3)g^\pm (t_3)g^>(t_3)\hfill & & & \\ +g^<(t_3)g^>(t_3)g^\pm (t_3)\hfill & & & \end{array}\right],`$ (139)
$`\mathrm{\Sigma }_{j,k}^{a(4)}(E)`$ $`=`$ $`U^4{\displaystyle _{\mathrm{}}^0}𝑑t_1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_2{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_3e^{iEt_1}`$
$`\times \left[\begin{array}{cccc}g^<(t_1)g^>(t_1)g^<(t_1t_2t_3)\hfill & & & \\ g^>(t_1)g^<(t_1)g^>(t_1t_2t_3)\hfill & & & \end{array}\right]`$
$`\times g^\pm (t_2)\left[\begin{array}{cccc}g^\pm (t_3)g^>(t_3)g^<(t_3)\hfill & & & \\ +g^<(t_3)g^\pm (t_3)g^>(t_3)\hfill & & & \\ +g^<(t_3)g^>(t_3)g^\pm (t_3)\hfill & & & \end{array}\right].`$ (144)
In addition, the terms for diagram illustrated in Fig. 4(l) are expressed by
$`\mathrm{\Sigma }_l^{r(4)}(E)`$ $`=`$ $`U^4{\displaystyle _0^{\mathrm{}}}𝑑t_1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_2{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_3e^{iEt_1}`$
$`\times \left[\begin{array}{cccc}g^>(t_1)g^>(t_1)g^<(t_1+t_2+t_3)\hfill & & & \\ g^<(t_1)g^<(t_1)g^>(t_1+t_2+t_3)\hfill & & & \end{array}\right]`$
$`\times g^\pm (t_2)\left[\begin{array}{cccc}g^\pm (t_3)g^>(t_3)g^<(t_3)\hfill & & & \\ +g^<(t_3)g^\pm (t_3)g^>(t_3)\hfill & & & \\ +g^<(t_3)g^>(t_3)g^\pm (t_3)\hfill & & & \end{array}\right],`$ (149)
$`\mathrm{\Sigma }_l^{a(4)}(E)`$ $`=`$ $`U^4{\displaystyle _{\mathrm{}}^0}𝑑t_1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_2{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_3e^{iEt_1}`$
$`\times \left[\begin{array}{cccc}g^<(t_1)g^<(t_1)g^>(t_1+t_2+t_3)\hfill & & & \\ g^>(t_1)g^>(t_1)g^<(t_1+t_2+t_3)\hfill & & & \end{array}\right]`$
$`\times g^\pm (t_2)\left[\begin{array}{cccc}g^\pm (t_3)g^>(t_3)g^<(t_3)\hfill & & & \\ +g^<(t_3)g^\pm (t_3)g^>(t_3)\hfill & & & \\ +g^<(t_3)g^>(t_3)g^\pm (t_3)\hfill & & & \end{array}\right].`$ (154) |
warning/0506/quant-ph0506066.html | ar5iv | text | # Bell’s Jump Process in Discrete Time
(September 27, 2005)
## Abstract
The jump process introduced by J. S. Bell in 1986, for defining a quantum field theory without observers, presupposes that space is discrete whereas time is continuous. In this letter, our interest is to find an analogous process in discrete time. We argue that a genuine analog does not exist, but provide examples of processes in discrete time that could be used as a replacement.
MSC (2000): 81T25, 60J10. PACS: 02.50.Ga; 03.65.Ta. Key words: Bell’s jump process; Markov chain; quantum theory on a lattice.
One of the central challenges for “hidden variable” approaches to quantum mechanics, such as the de Broglie-Bohm pilot wave theory, is to provide an adequate account of relativistic quantum field theory. To address this, Bell introduced a jump process on a discrete lattice , intended to reproduce the quantum mechanical predictions for fermion number density in space. The same method can be used to generate stochastic trajectories for any discrete observable, both in field theory and in nonrelativistic quantum mechanics. For a discretized position observable in nonrelativistic quantum mechanics, Bell’s process reduces to the de Broglie-Bohm pilot wave theory in the continuum limit , so it is a natural analog of this theory for discrete “beables” .
Although the “beables” in Bell’s process are discrete, it still contains a continuous time parameter. However, there are several reasons for developing a discrete-time version of the process. Firstly, some approaches to quantum gravity are based on fundamentally discrete space-time structures, so a realist account of these theories along Bohmian lines would have to be fully discrete. Secondly, “hidden variable” theories, no matter whether they are realized in nature or not, can be useful for numerical simulations , visualizations , bookkeeping , and obtaining better intuitions about quantum phenomena. Numerical simulations are discrete by nature, and a fully discrete theory may also be useful when dealing with quantum phenomena usually described in a discrete setting, such as those considered in quantum information and computation. Thirdly, Valentini has recently proposed that matter in quantum nonequilibrium, i.e. beables with distributions other than $`|\mathrm{\Psi }|^2`$, if existant, may provide astonishing computational resources, enabling us to solve NP-complete problems in polynomial time. However, since classical analog computers can also outperform Turing machines if the continuous variables can be manipulated with perfect accuracy, this claim would be simpler to verify in a fully discrete model.
In this letter, we highlight the difficulties inherent in discretizing Bell’s jump process, and propose two concrete discretized processes that circumvent them and converge to Bell’s process as the time step $`\tau `$ goes to zero. Other possibilities exist, along the lines of recent proposals by Aaronson , and these will be developed in future work.
Bell’s process is a Markovian pure jump process $`(Q_t)_t`$ on a lattice $`𝒬`$ with rate for the jump $`q^{}q`$ given by
$$\sigma _t(q|q^{})=\frac{\left[\frac{2}{\mathrm{}}\mathrm{Im}\mathrm{\Psi }_t|P(q)HP(q^{})|\mathrm{\Psi }_t\right]^+}{\mathrm{\Psi }_t|P(q^{})|\mathrm{\Psi }_t},$$
(1)
where $`x^+=\mathrm{max}(x,0)`$ denotes the positive part of $`x`$, $`\mathrm{\Psi }_t`$ is the state vector of a quantum (field) theory, evolving in some Hilbert space $``$ according to
$$i\mathrm{}\frac{d\mathrm{\Psi }_t}{dt}=H\mathrm{\Psi }_t,$$
(2)
$`H`$ is the Hamiltonian, and $`P(q)`$ is the projection to the subspace $`_q`$, where the $`_q`$ form an orthogonal decomposition, $`=_{q𝒬}_q`$. Relevant properties of Bell’s process are that at every time $`t`$, the distribution of $`Q_t`$ is the quantum distribution
$$\mathrm{\Psi }_t|P(q)|\mathrm{\Psi }_t,$$
(3)
and that its net probability current between $`q^{}`$ and $`q`$, $`\sigma _t(q|q^{})(Q_t=q^{})\sigma _t(q^{}|q)(Q_t=q)`$ where $``$ denotes “probability,” agrees with the quantum expression for the probability current,
$$\frac{2}{\mathrm{}}\mathrm{Im}\mathrm{\Psi }|P(q)HP(q^{})|\mathrm{\Psi }.$$
(4)
Since many constructions are easier in discrete time than in continuous time, one might have expected that there is an analogous Markov chain $`(\stackrel{~}{Q}_t)_{t\tau }`$ on $`𝒬`$ with discrete time step $`\tau `$ such that the probability $`_t(q^{}q)`$ for the transition $`q^{}q`$, i.e., the conditional probability $`(\stackrel{~}{Q}_{t+\tau }=q|\stackrel{~}{Q}_t=q^{})`$, is given by a formula similar to (1), with $`H`$ replaced by a simple function of the unitary $`U`$ defining the time evolution
$$\mathrm{\Psi }_{t+\tau }=U\mathrm{\Psi }_t,$$
(5)
and that one could arrive at this formula by a reasoning similar to the one leading to (1) from (3) and (4), as given in \[9, Sec. 2.5\].
However, this is not possible in any obvious way. The obstacle is that in the time-discrete case there is no obvious formula for the net probability current $`J(q,q^{})`$ between $`q^{}`$ and $`q`$, replacing (4) of the continuous case. Given an expression for $`J(q,q^{})`$ in terms of $`\mathrm{\Psi },P`$, and $`U`$, we could set
$$_t(q^{}q)=\frac{J_t(q,q^{})^+}{\mathrm{\Psi }_t|P(q^{})|\mathrm{\Psi }_t}\text{ for }qq^{},$$
(6)
which would define a Markov chain $`(\stackrel{~}{Q}_t)_{t\tau }`$ whose probability current
$$_t(q^{}q)(\stackrel{~}{Q}_t=q^{})_t(qq^{})(\stackrel{~}{Q}_t=q)$$
(7)
coincides with $`J(q,q^{})`$ and whose distribution at any time $`t`$ coincides with the quantum distribution (3), provided $`J(q,q^{})`$ has the following properties:
$`J(q,q^{})`$ $``$ (8a)
$`J(q^{},q)`$ $`=J(q,q^{})`$ (8b)
$`{\displaystyle \underset{q𝒬}{}}J(q,q^{})^+`$ $`\mathrm{\Psi }|P(q^{})|\mathrm{\Psi }`$ (8c)
$`{\displaystyle \underset{q^{}𝒬}{}}J(q,q^{})`$ $`=\mathrm{\Psi }|U^{}P(q)U|\mathrm{\Psi }\mathrm{\Psi }|P(q)|\mathrm{\Psi }.`$ (8d)
Currents of the form (7), with transition probabilities (6) and distribution (3), have these properties by construction. Property (8c) expresses that no greater amount of probability can get transported away from $`q^{}`$ than present at $`q^{}`$, and (8d) guarantees the quantum distribution (3) at the next time step. The obvious way of guessing a formula for $`J(q,q^{})`$ is to start from one of the expressions
$`\mathrm{\Psi }|U^{}P(q)UP(q^{})|\mathrm{\Psi }`$ (9a)
$`\mathrm{\Psi }|P(q)UP(q^{})|\mathrm{\Psi },`$ (9b)
to multiply it by any numerical constant, to take the real or imaginary parts to ensure (8a), and to anti-symmetrize in $`q`$ and $`q^{}`$ to ensure (8b). However, all expressions thus obtained generically violate (8d), except for the anti-symmetrization of $`2\mathrm{R}\mathrm{e}`$(9a),
$$\begin{array}{c}J(q,q^{})=\frac{1}{2}\mathrm{\Psi }|(U^{}P(q)UP(q^{})+P(q^{})U^{}P(q)U\hfill \\ \hfill U^{}P(q^{})UP(q)P(q)U^{}P(q^{})U)|\mathrm{\Psi },\end{array}$$
(10)
which can violate (8c) (numerically we found 46 examples of such violations among one thousand randomly chosen $`U`$ and $`\psi `$ in $`=^3`$ with fixed one-dimensional projections $`P(q)`$ and $`P(q^{})`$).
However, a different reasoning leads to a process in discrete time that has some features in common with Bell’s process. Choose $`H`$ such that
$$U=e^{i\tau H},$$
(11)
so that the evolution (2) generated by $`H`$ is a continuation of the evolution (5) generated by $`U`$. (The degree of non-uniqueness of this choice is discussed later.) Then, consider Bell’s process $`(Q_t)_t`$ in continuous time for this $`H`$. By restriction to just the integer times, we obtain a Markov process $`\stackrel{~}{Q}_t:=Q_t`$ for $`t\tau `$.
The process $`(\stackrel{~}{Q}_t)_{t\tau }`$ has the quantum distribution (3) at every time. It is important for this that the two evolution laws (2) and (5) for $`\mathrm{\Psi }`$ lead to the same $`\mathrm{\Psi }_t`$ at every $`t`$ that is an integer multiple of $`\tau `$. It makes no sense to ask whether the probability current of this process, $`(\stackrel{~}{Q}_{t+\tau }=q,\stackrel{~}{Q}_t=q^{})(\stackrel{~}{Q}_{t+\tau }=q^{},\stackrel{~}{Q}_t=q)`$, agrees with the one prescribed by quantum theory, since, as discussed above, quantum theory does not prescribe a unique current in the discrete-time case. Note that in the limit $`\tau 0`$ the process approaches Bell’s process. This fact and the simple and straightforward construction of $`(\stackrel{~}{Q}_t)_{t\tau }`$ suggest that this may be the closest one can get to an analog of Bell’s process in the time-discrete case.
The transition probability $`_t(q^{}q)=(\stackrel{~}{Q}_{t+\tau }=q|\stackrel{~}{Q}_t=q^{})`$ does not, however, possess a simple formula in terms of $`\mathrm{\Psi }_t`$, $`U`$, and $`P()`$ analogous to (1), only the following one:
$$\begin{array}{cc}& _{t_0}(q^{}q)=\underset{n=0}{\overset{\mathrm{}}{}}\underset{q_0,\mathrm{},q_n𝒬}{}\delta _{q^{},q_0}\delta _{q,q_n}\underset{t_0}{\overset{t_0+\tau }{}}dt_1\underset{t_1}{\overset{t_0+\tau }{}}dt_2\mathrm{}\underset{t_{n1}}{\overset{t_0+\tau }{}}dt_n\times \hfill \\ & \times \mathrm{exp}\left(\underset{t_0}{\overset{t_0+\tau }{}}\sigma _s(𝒬|q_{\mathrm{max}\{k:t_k<s\}})𝑑s\right)\underset{k=1}{\overset{n}{}}\sigma _{t_k}(q_k|q_{k1}),\hfill \end{array}$$
(12)
with $`\sigma _s(q|r)`$ given by (1) and $`\sigma _s(𝒬|r):=_{q𝒬}\sigma _s(q|r)`$. Eq. (12) is a fact about any jump process in continuous time with jump rates $`\sigma `$ (applied here to Bell’s process $`Q_t`$).<sup>4</sup><sup>4</sup>4To get a grasp of (12), begin with noting that $`\sigma _s(𝒬|r)`$ is the total jump rate at time $`s`$ in the configuration $`r`$. The probability that no jump takes place before time $`t`$, if the process starts at $`t_0`$ in $`q_0`$, is $`\mathrm{exp}\left(_{t_0}^t\sigma _s(𝒬|q_0)𝑑s\right)`$. Thus, the probability that the first jump takes place between time $`t`$ and $`t+dt`$ is $`\mathrm{exp}\left(_{t_0}^t\sigma _s(𝒬|q_0)𝑑s\right)\sigma _t(𝒬|q_0)dt`$. The probability that the destination of the first jump is $`q_1`$, given that the jump takes place at time $`t`$, is $`\sigma _t(q_1|q_0)/\sigma _t(𝒬|q_0)`$. Conditional on that the first jump occurs at $`t`$ and leads to $`q_1`$, the distribution of the times and destinations of the further jumps is the same as for a process starting at time $`t`$ in $`q_1`$. Thus, the probability of a path $`q_0,\mathrm{},q_n`$ with the $`k`$-th jump between $`t_k`$ and $`t_k+dt_k`$ and no further jump before $`t_0+\tau `$ is the integrand of (12) times $`dt_1\mathrm{}dt_n`$. Now add (respectively integrate) the probabilities of all ways the process can move from $`q^{}`$ to $`q`$ in the time interval $`[t_0,t_0+\tau ]`$, namely by means of $`n`$ jumps at times $`t_1,\mathrm{},t_n`$ with destinations $`q_1,\mathrm{},q_n`$. For a more detailed discussion of such probability formulas, see and .
The process $`\stackrel{~}{Q}`$ is not completely determined by $`\mathrm{\Psi }_0`$, $`U`$, and $`P()`$ since $`H`$ is not completely determined by (11), even though in many cases there may be a natural choice of $`H`$. For example, if $`U`$ has an eigenvalue $`e^{i\theta }`$, then $`H`$ may have as the corresponding eigenvalue any of the numbers $`\frac{\theta }{\tau }+\frac{2\pi }{\tau }k`$ with $`k`$. More generally, for any self-adjoint operator $`S`$ with spectrum contained in $`\frac{2\pi }{\tau }`$ and commuting with $`H`$ (in the sense of commuting spectral projections), $`H+S`$ is another solution of (11) for given $`U`$. A unique $`H`$ could be selected by the additional condition that the spectrum of $`H`$ be contained in $`(\frac{\pi }{\tau },\frac{\pi }{\tau }]`$.
In the particularly simple situation $`|𝒬|=2`$, there does exist a time-discrete analog $`(\widehat{Q}_t)_{t\tau }`$ to Bell’s process. In this case, the expression (10) satisfies (8) and thus defines a process; in fact, the net probability current between the two configurations $`q^{}`$ and $`q`$ is already determined by the distribution (3) and must be
$$\mathrm{\Psi }_t|U^{}P(q)U|\mathrm{\Psi }_t\mathrm{\Psi }_t|P(q)|\mathrm{\Psi }_t,$$
(13)
since any increase or decrease can occur only by transitions from or to the other configuration. Just as Bell’s process has the smallest jump rates compatible with the current (4) , we may choose now the smallest transition probabilities compatible with the current (13), which are
$$_t\left(\widehat{Q}_{t+\tau }q|\widehat{Q}_t=q\right)=\frac{\mathrm{\Psi }_t|(P(q)U^{}P(q)U)|\mathrm{\Psi }_t^+}{\mathrm{\Psi }_t|P(q)|\mathrm{\Psi }_t}.$$
(14)
This need not coincide with the transition probability (12) of $`(\stackrel{~}{Q}_t)`$, even though in the limit $`\tau 0`$ also $`(\widehat{Q}_t)`$ converges to Bell’s process. The same construction can be applied to the case $`|𝒬|>2`$ if $`U`$ involves only pairs of configurations, i.e., if there is a partition of $`𝒬`$ into subsets, all of which are either pairs or singlets, such that $`P(q)UP(q^{})=0`$ whenever $`q`$ and $`q^{}`$ do not belong to the same subset. Then (10) still satisfies (8) and thus defines a process. An example of this is a quantum computing circuit, realized through a time sequence of single qubit unitaries and CNOT gates. (Here, a configuration $`q`$ corresponds to a definite value for the computational basis observable for each qubit.)
To contrast the previous processes with an example of a process that does not converge to Bell’s process in the limit $`\tau 0`$ but has the quantum distribution (3) at every time, we define the process $`(Q_t^{})_{t\tau }`$ by the transition probability
$$(Q_{t+\tau }^{}=q|Q_t^{}=q^{})=\mathrm{\Psi }_{t+\tau }|P(q)|\mathrm{\Psi }_{t+\tau }.$$
(15)
This means that for every $`t`$, $`Q_t^{}`$ is independent of the past and has the quantum distribution. Its limit as $`\tau 0`$, in a suitable sense, is simply the process $`(\stackrel{~}{Q}_t^{})_t`$ for which every $`\stackrel{~}{Q}_t^{}`$ is independent of the past and has the quantum distribution, a process reminiscent of Bell’s description of a precise version of the “many worlds” interpretation of quantum mechanics: “\[I\]nstantaneous classical configurations \[$`Q`$\] are supposed to exist, and to be distributed \[…\] with probability $`|\psi |^2`$. But no pairing of configurations at different times, as would be effected by the existence of trajectories, is supposed.”
Acknowledgments. We thank Shelly Goldstein of Rutgers University (USA) for helpful discussions and two anonymous referees for their comments. R.T. thanks the Institut des Hautes Études Scientifiques at Bures-sur-Yvette, France, for hospitality. This work has been partially supported by the European Commission through its 6th Framework Programme “Structuring the European Research Area” and the contract Nr. RITA-CT-2004-505493 for the provision of Transnational Access implemented as Specific Support Action. |
warning/0506/gr-qc0506077.html | ar5iv | text | # Numerical validation of the Kerr metric in Bondi-Sachs form
## I Introduction
It is widely believed that the end-state of most massive objects is a black hole described by the Kerr geometry. The most powerful source of gravitational radiation is the merger of a black hole with another black hole (or other compact object such as a neutron star). Such processes need to be calculated numerically, and the field of numerical relativity has been developed to tackle such problems – see, for example the review lehner . The most popular approach to numerical relativity is based on the ADM formalism, and variations thereof, as discussed in the review cited. In these formalisms, it is understood how to set initial data for a Kerr black hole cook .
An alternative approach to numerical relativity is based on an evolution of the metric variables in the Bondi-Sachs metric bondi ; sachs . The approach has been used by a number of different groups – for example newnews ; particle ; roberto ; rdi2 ; siebel1 ; bartnik , and for a review see winLR . Applications involving black holes have been restricted to the Schwarzschild case, because of the difficulty in setting initial and boundary data for a rotating black hole. Also, the problem of obtaining just the horizon data is much simpler win , and from such data one could then use a numerical code to construct a complete Bondi-Sachs representation of the Kerr geometry.
The representation of the Kerr geometry obtained by Pretorius and Israel PI-1 (and hereafter this reference will be denoted as PI) is not in Bondi-Sachs form. However, the PI coordinates are based on outgoing light cones, and as such should provide a useful starting-point for the construction of the Bondi-Sachs form of the Kerr metric. We have constructed a coordinate transformation on the PI coordinates that brings the transformed metric into Bondi-Sachs form. In order to validate our metric (and also that of PI) we need to show that the Ricci tensor vanishes. The metric is not explicit with respect to the new coordinates, so an analytic approach would seem to be hopeless. Instead we wrote code that numerically computes the Ricci tensor for the metric, and confirmed that the metric is Ricci flat.
In order to describe the Kerr geometry, a metric must not only be Ricci flat but must also be regular on the axis of symmetry. Thus, we also find the leading order terms of the metric near the axis, and thereby confirm regularity.
Fletcher and Lun lun have proposed a metric in Bondi-Sachs form that represents the Kerr geometry. We will call this the FL metric. The metric has the convenient property that its coefficients are explicit functions of its coordinates, so a numerical evaluation of the Ricci tensor is straightforward. This we did, confirming that the Ricci tensor vanishes. However, we also found that it does not have a normal behaviour near the pole, thus rendering it inappropriate for use in numerical relativity.
Sec. II summarizes results obtained in PI that will be used later. Then in Sec. III we obtain in explicit form the equations that relate the Boyer-Lindquist and PI coordinates, and also find explicit expressions for the PI metric quantities. Sec. IV derives the coordinate transformation needed to change the PI metric into Bondi-Sachs form. Sec. V discusses the numerical implementation of the computation of the Ricci tensor, and presents the results. The properties of the FL metric are discussed in Sec. VI. We end with a Conclusion, Sec. VII.
## II Background: the PI metric
The Kerr metric in standard Boyer-Lindquist coordinates $`(t,r,\theta ,\varphi )`$ is
$$ds^2=\frac{\mathrm{\Sigma }}{\mathrm{\Delta }}dr^2+\mathrm{\Sigma }d\theta ^2+R^2\mathrm{sin}^2\theta d\varphi ^2\frac{4mar\mathrm{sin}^2\theta }{\mathrm{\Sigma }}d\varphi dt\left(1\frac{2mr}{\mathrm{\Sigma }}\right)dt^2$$
(1)
where
$$\mathrm{\Sigma }=r^2+a^2\mathrm{cos}^2\theta ,\mathrm{\Delta }=r^2+a^22mr,R^2=r^2+a^2+\frac{2ma^2r\mathrm{sin}^2\theta }{\mathrm{\Sigma }}.$$
(2)
The metric obtained by PI (Eq. (PI-31)) uses coordinates $`(t,r_{},\lambda ,\varphi )`$, with $`r_{}=r_{}(r,\theta ),\lambda =\lambda (r,\theta )`$. The metric is
$$ds^2=\frac{\mathrm{\Delta }}{R^2}(dr_{}^2dt^2)+\frac{L^2}{R^2}d\lambda ^2+R^2\mathrm{sin}^2\theta (d\varphi \omega _Bdt)^2,$$
(3)
where (see Eq. (PI-30))
$$\omega _B=\frac{2mar}{\mathrm{\Sigma }R^2},$$
(4)
and $`L`$ is defined in Eq. (9) below. PI often use $`\theta _{}`$ instead of $`\lambda `$ as a coordinate, the two quantities being related by
$$\lambda =\mathrm{sin}^2(\theta _{})$$
(5)
(see Eqs. (PI-36) and (PI-37)). Thus an alternative form of the metric (3) is
$$ds^2=\frac{\mathrm{\Delta }}{R^2}(dr_{}^2dt^2)+\frac{L^2\mathrm{sin}^22\theta _{}}{R^2}d\theta _{}^2+R^2\mathrm{sin}^2\theta (d\varphi \omega _Bdt)^2.$$
(6)
PI define the quantities (see Eqs. (PI-7))
$$P^2(\theta ,\lambda )=a^2(\lambda \mathrm{sin}^2\theta ),Q^2(r,\lambda )=(r^2+a^2)^2a^2\lambda \mathrm{\Delta },$$
(7)
and (see Eqs. (PI-36) and (PI-14))
$$F(r,\theta ,\lambda )=_r^{\mathrm{}}\frac{dr^{}}{Q(r^{},\lambda )}_\theta ^\theta _{}\frac{d\theta ^{}}{P(\theta ^{},\lambda )}\text{ subject to the constraint }F=0,$$
(8)
as well as (see Eqs. (PI-28) and (PI-25))
$$L=\mu PQ\text{ with }\mu =\frac{F}{\lambda }.$$
(9)
Eq. (PI-15), and using Eqs. (PI-9) and (PI-39), defines $`r_{}`$
$$r_{}=\frac{r^2+a^2}{\mathrm{\Delta }(r)}𝑑r+_r^{\mathrm{}}\frac{r^2+a^2Q(r^{},\lambda )}{\mathrm{\Delta }(r^{})}𝑑r^{}+_\theta _{}^\theta P(\theta ^{},\lambda )𝑑\theta ^{}$$
(10)
There is one minor difference in the notation used in PI and here. We regard the mass of the black hole $`m`$ and the normalized angular momentum $`a`$ as parameters: thus we do not show $`a`$ or $`m`$ as variables in any functional dependence list.
In order to evaluate the metric, the implicit equations need to be reduced to a numerically solvable form. This is now done.
## III Properties of the coordinate transformation and metric variables
### III.1 Relation between $`(r,\theta )`$ and $`(r_{},\theta _{})`$
Eq. (8) implies
$$_r^{\mathrm{}}\frac{dr^{}}{Q(r^{},\lambda )}_\theta ^\theta _{}\frac{d\theta ^{}}{P(\theta ^{},\lambda )}=0,$$
(11)
which constitutes a relationship of the form function$`(r,\theta ,\theta _{})=0`$. We found an exact solution for this equation but it was more efficient to solve the $`r`$-integral in Eq.(11) numerically and the $`\theta `$-integral exactly as,
$`{\displaystyle _\theta ^\theta _{}}{\displaystyle \frac{d\theta ^{}}{P(\theta ^{},\lambda )}}`$ $`=`$ $`{\displaystyle \frac{1}{a}}\left[\mathrm{\Phi }(\theta _{},\theta _{})\mathrm{\Phi }(\theta ,\theta _{})\right]`$ (12)
$`\mathrm{\Phi }(\alpha ,\theta _{})`$ $`=`$ $`ϵ(\mathrm{cos}\alpha )\mathrm{\Xi }({\displaystyle \frac{\mathrm{sin}\alpha }{\mathrm{sin}\theta _{}}},\mathrm{sin}\theta _{})`$ (13)
$`\mathrm{\Xi }(z,k)`$ $`=`$ $`{\displaystyle _0^z}{\displaystyle \frac{dt}{\sqrt{1t^2}\sqrt{1k^2t^2}}}`$ (14)
with the condition that $`a0`$ and where $`ϵ(x)`$ is the $`signum`$ of $`x`$.
Eq. (10) can be written as
$$r_{}=I_1(r)+I_2(r,\lambda )+I_3(\theta ,\lambda )$$
(15)
and we now evaluate the integrals $`I_1`$, $`I_2`$ and $`I_3`$.
$$I_1=r+m\mathrm{ln}(r^2+a^22mr)+\frac{2m^2}{\sqrt{\nu }}Q_0\left(\frac{\sqrt{\nu }}{mr}\right)+C,$$
(16)
where $`\nu =m^2a^2`$; $`Q_0(x)`$ is the zero-order special Legendre function of the second kind; $`C`$ is an arbitrary constant of integration, and we set $`C=0`$. $`I_2`$ is expressed as
$$I_2=_r^{\mathrm{}}\frac{\zeta ^2+a^2\sqrt{(\zeta ^2+a^2)^2a^2\lambda \mathrm{\Delta }(\zeta )}}{\mathrm{\Delta }(\zeta )}𝑑\zeta $$
(17)
and will be evaluated numerically. $`I_3`$ reduces to
$$I_3=a\left[[E(\gamma ,t)E(\frac{\pi }{2},t)](1t^2)\left[F_E(\gamma ,t)F_E(\frac{\pi }{2},t)\right]\right]$$
(18)
$$\gamma =\mathrm{sin}^1\frac{\mathrm{sin}\theta }{\mathrm{sin}\theta _{}},t=\mathrm{sin}\theta _{}$$
(19)
and where $`E`$ and $`F_E`$ are Legendre elliptic integrals of the first and second kinds (we use the notation $`F_E`$, rather than the usual $`F`$, to avoid any confusion with the $`F`$ introduced by PI and defined in Eq. (8)).
### III.2 The metric variable $`L`$
The function $`L`$ in the metric needs to be written in a form that can be evaluated explicitly. Now, $`L=\mu PQ`$ and the problem is to find an explict expression for $`\mu `$. Using Eqs. (8) and (9)
$$\mu =T_1+T_2=_r^{\mathrm{}}\frac{\frac{}{\lambda }Q(r^{},\lambda )}{Q(r^{},\lambda )^2}𝑑r^{}+\frac{}{\lambda }_\theta ^{\mathrm{sin}^1\sqrt{\lambda }}\frac{d\theta }{P(\theta ,\lambda )}.$$
(20)
The first term $`T_1`$ reduces to
$$T_1=\frac{1}{2}_r^{\mathrm{}}\frac{a^2(r^2+a^22mr^{})}{\left[(r^2+a^2)^2a^2\lambda (r^2+a^22mr^{})\right]^{3/2}}𝑑r^{}$$
(21)
and this term will need to be evaluated numerically. The second term $`T_2`$ can be simplified to
$$T_2=\frac{1}{2}\frac{a\mathrm{sin}\theta \mathrm{cos}\theta }{a^2\lambda (\lambda 1)\sqrt{\lambda \mathrm{sin}^2\theta }}+\frac{1}{2}\frac{\mathrm{\Omega }(\theta )ϵ(\mathrm{cos}\theta )\mathrm{\Omega }(\theta _{})}{a^2\lambda (\lambda 1)\sqrt{\lambda \mathrm{sin}^2\theta }}$$
(22)
where
$$\mathrm{\Omega }(\alpha )=a\left[E(\eta ,t)(1t^2)F_E(\eta ,t)\right],\eta =\mathrm{sin}^1\frac{\mathrm{sin}\alpha }{\mathrm{sin}\theta _{}},t=\mathrm{sin}\theta _{}.$$
(23)
### III.3 Regularity on the axis $`\theta _{}=0`$
The Kerr geometry is axially symmetric, and therefore the metric needs to satisfy a regularity condition on the axis $`\theta _{}=0`$. In order to check the condition we need to know the values on the axis of the non-zero metric components, as well as the leading term in a series expansion in $`\theta _{}`$ for those components that are zero on the axis. Clearly the following metric functions have non-zero limits at $`\theta _{}=\theta =0`$: $`\mathrm{\Delta },\mathrm{\Sigma },R^2,\omega _B`$; but we will need to work out $`Q`$ and $`P`$ near the axis so that we can obtain expressions for $`\theta `$ and $`L`$ near the axis.
We assume that $`\theta _{}`$ and $`\theta `$ are small so that
$$P=a\sqrt{\theta _{}^2\theta ^2}$$
(24)
and
$$Q=\left(r^2+a^2\right)\left(1\frac{a^2\theta _{}^2\mathrm{\Delta }}{2r^2+2a^2}\right),Q^1=\left(r^2+a^2\right)^1\left(1+\frac{a^2\theta _{}^2\mathrm{\Delta }}{2r^2+2a^2}\right).$$
(25)
Then Eq. (11) leads to
$$\theta =\frac{r\theta _{}}{\sqrt{\left(r^2+a^2\right)}}.$$
(26)
Next, we evaluate $`\mu `$ as
$$\mu =\frac{F}{\lambda }=\frac{1}{2\theta _{}}\frac{F}{\theta _{}}=\frac{r}{2a^2\theta _{}^2},$$
(27)
and then $`L=\mu PQ`$ is
$$L=\frac{r\sqrt{r^2+a^2}}{2\theta _{}}.$$
(28)
Thus the non-trivial metric components at, or near, the $`\theta _{}=0`$ axis can be evaluated. We find that the PI metric satisfies
$$g_{\theta _{}\theta _{}}=r^2,g_{\varphi \varphi }=r^2\theta _{}^2,g_{t\varphi }=\frac{2r^3\theta _{}^2ma}{(r^2+a^2)^2}.$$
(29)
Since $`\theta _{}^2g_{\theta _{}\theta _{}}/g_{\varphi \varphi }1`$ as $`\theta _{}0`$, the metric is regular on the axis.
## IV Coordinate transformation to Bondi-Sachs form
The transformation
$$tu=tr_{}$$
(30)
brings the metric (3) into Bondi-Sachs form except for the last term
$$R^2sin^2\theta (d\varphi \omega _B(du+dr_{}))^2;$$
(31)
but this term can be brought into the required form with the transformation
$$\varphi \varphi _{}=\varphi +H(r_{},\theta _{})$$
(32)
if the arbitrary function $`H`$ satisfies the condition
$$\frac{H(r_{},\theta _{})}{r_{}}=\omega _B.$$
(33)
Substitution of the coordinate transforms (30) and (32) into Eq. (3) results in
$$ds^2=\frac{(du^22dudr_{})\mathrm{\Delta }}{R^2}+4\frac{L^2sin^2\theta _{}cos^2\theta _{}d\theta _{}^2}{R^2}+(\omega _Bdu+\frac{H(r_{},\theta _{})}{\theta _{}}d\theta _{}d\varphi _{})^2R^2sin^2\theta .$$
(34)
We now need to obtain an expression for $`\frac{H(r_{},\theta _{})}{\theta _{}}`$. This is non-trivial as $`\omega _B`$ is given in terms of $`r`$ and $`\theta `$ rather than in terms of $`r_{}`$ and $`\theta _{}`$. Since $`H`$ in the coordinate transform (32) is arbitrary up to the condition (33), the remaining $`\frac{H}{\theta _{}}`$ component is still undefined. We proceed by defining
$$H(r_{},\theta _{})=_r^{\mathrm{}}\frac{2mas}{\mathrm{\Delta }(s)Q(s,\theta _{})}𝑑s.$$
(35)
It then follows that Eq. (33) is met, since
$$\frac{H(r_{},\theta _{})}{r_{}}=\frac{2mar}{\mathrm{\Sigma }R^2}=\omega _B.$$
(36)
After integration of (35), $`\frac{H(r_{},\theta _{})}{\theta _{}}`$ is obtained as
$$\frac{H(r_{},\theta _{})}{\theta _{}}=\omega _B\beta \mathrm{sin}(2\theta _{})$$
(37)
where
$$\beta =\mu P^2+\frac{a^3m}{\omega _B}_r^{\mathrm{}}\frac{s}{Q^3}𝑑s.$$
(38)
Note that the integral in Eq. (38) is proportional to $`T_1`$ defined in Eq. (21) and, as remarked earlier, will need to be evaluated numerically.
The metric now becomes
$$ds^2=\frac{(du^22dudr_{})\mathrm{\Delta }}{R^2}+4\frac{L^2sin^2\theta _{}cos^2\theta _{}d\theta _{}^2}{R^2}+(\omega _B\left[du+\beta \mathrm{sin}(2\theta _{})d\theta _{}\right]d\varphi _{})^2R^2sin^2\theta $$
(39)
which is the Kerr metric in Bondi-Sachs form.
### IV.1 Regularity on the axis $`\theta _{}=0`$
The situation is very similar to that discussed earlier for the PI metric in Sec. III.3. The only new metric function that has been introduced is $`\beta `$, which evaluates on the axis to
$$\beta =\frac{a^2(5r^2+a^2)}{8r(r^2+a^2)}.$$
(40)
The extra terms introduced into the metric components by way of $`\beta `$ are multipled by $`\theta _{}^3`$, and are therefore too small to affect the issue of regularity.
## V Numerical method and results
### V.1 Coding issues
We mention only matters that are not routine. The expressions containing elliptic integrals $`E`$ and $`F_E`$ were evaluated using Bulirsch’s elliptical numerical solvers bulirsch-1 . The quantity $`\mathrm{\Xi }(z,k)`$ defined in Eq. (14) can be expressed as a linear combination of elliptic integrals of various kinds, also evaluated using Bulirsch elliptical numerical solvers bulirsch-1 . The quantity $`I_2`$ is defined in Eq. (17) by means of an integral whose integrand decreases monotonically as $`\zeta \mathrm{}`$, therefore standard solvers for infinite boundaries can be used. In similar fashion the first integral in Eq. (11) and the integral in Eq. (21) were evaluated using the same Romberg integration methods.
The quantity $`\mu `$ is singular as $`\theta _{}0`$, and this proves to be problematic for accurate numerical evaluation of metric quantities near the axis. It is for this reason that, in the Tables below, the smallest value of $`\theta _{}`$ that is used is 0.3.
### V.2 Construction of data on a regular grid
In order to obtain numerically the Ricci tensor of the metric to $`O(h^2)`$, a 13-point regular grid was used. First we find first metric derivatives using the central difference formula, then we calculate the Christoffel symbols, then we calculate first derivatives of the Christoffel symbols using the central difference formula, and finally we calculate the various components of the Ricci tensor.
For given $`(r_{},\theta _{})`$ we obtain the Boyer-Lyndquist coordinates $`(r,\theta )`$ as follows. We make an estimate for $`r`$ and use Eq. (11) to find $`\theta `$. Once $`\theta `$ has been obtained in this manner, $`(r,\theta ,\theta _{})`$ is substituted into Eq. (15) establishing a value for $`r_{}`$. Of course, in general this will not be the desired value of $`r_{}`$, so we repeat the calculation with a different estimate of $`r`$, and then use the bisection method to find, to an accuracy of $`10^{12}`$ the value of $`r`$ which leads to the correct value of $`r_{}`$. Now, having obtained $`(r,\theta )`$ for given $`(r_{},\theta _{})`$, we find all the metric coefficients, including $`L`$, at the point $`(r_{},\theta _{})`$.
### V.3 Results
First, the numerical code was validated by using it to find the Ricci tensor of the Kerr metric in standard Boyer-Linquist coordinates. Second order convergence to zero was observed.
The Ricci tensor of the metric (39) was evaluated. Tables 1, 2, 3 and 4 show values for $`r`$, $`\theta `$, $`L`$ and $`\beta `$ on a regular $`(r_{},\theta _{})`$ grid centered at $`(0.4,0.3)`$ with $`m=1`$, $`a=0.1`$, and $`h=\mathrm{\Delta }r_{}=\mathrm{\Delta }\theta _{}=0.01`$. From this data it is straightforward to find numerical values for all components of the metric $`g_{\alpha \beta }`$ on the regular $`(r_{},\theta _{})`$ grid, and thus to compute the Ricci tensor $`R_{\alpha \beta }`$ at the centre of the grid. The results are shown in the second column of Table 5. Also in that Table we show $`R_{\alpha \beta }`$ evaluated using coarser grids, $`h=0.02`$ and $`h=0.04`$, centered at the same point $`(0.4,0.3)`$ (but in these cases we do not give the intermediate data $`r`$, $`\theta `$, $`L`$ or $`\beta `$). The last columns of the Table show the result of convergence testing, and it is clear that all components are second order convergent to zero.
The values of $`(r,\theta )`$ and of $`g_{\alpha \beta }`$ were also found for the cases $`m=1`$, $`a\{0.1,0.2\}`$ on regular grids $`(r_{},\theta _{})`$ centered at $`r_{}\{0.4,0.5\}5`$, $`\theta _{}\{0.3,0.6,0.9,1.2\}`$, and for three different discretization parameters namely $`h=0.01`$, $`h=0.02`$ and $`h=0.04`$. Rather than present the convergence rates of all the components of the 16 sets of Ricci tensors, we instead calculated the $`L_2`$-norm defined by
$$R_{\alpha \beta }=\sqrt{\underset{\alpha =0}{\overset{3}{}}\underset{\beta =0}{\overset{3}{}}\frac{(R_{\alpha \beta })^2}{16}}.$$
(41)
These results are given in Table 6. In all cases, second-order convergence to zero was observed.
Because the metric (39) is derived from the PI metric, the above results imply that the PI metric is also Ricci flat; we have performed computations similar to those described above and explicitly confirmed this – see Table 7.
## VI The metric of Fletcher and Lun
Using coordinates $`(u,r,\theta ,\varphi )`$ Fletcher and Lun lun have obtained a metric in Bondi-Sachs form, denoted here by FL. The metric coefficients do not contain implicit functions, so a numerical evaluation of the Ricci tensor is straightforward, and convergence to zero was confirmed (see Table 8). However, the metric suffers from an important restriction in that it is not regular at $`\theta =0`$.
The FL metric coefficients at $`\theta =0`$ were found to be
$`l_{tt}`$ $`=`$ $`{\displaystyle \frac{r^2a^2+a^2\mathrm{tanh}^2\alpha +2mr}{r^2a^2+a^2\mathrm{tanh}^2\alpha }}`$
$`l_{tr}`$ $`=`$ $`{\displaystyle \frac{r^2a^2+a^2\mathrm{tanh}^2\alpha }{\sqrt{r\left(r^3+ra^2+2a^2m\right)}}}`$
$`l_{t\theta }`$ $`=`$ $`{\displaystyle \frac{\left(r^2\mathrm{cosh}^2\alpha 2mr\mathrm{cosh}^2\alpha +a^2\right)a}{\left(r^2\mathrm{cosh}^2\alpha +a^2\right)\mathrm{cosh}^2\alpha }}`$
$`l_{t\varphi }`$ $`=`$ $`2{\displaystyle \frac{amr\mathrm{tanh}^2\alpha }{r^2a^2+a^2\mathrm{tanh}^2\alpha }}`$
$`l_{\theta \theta }`$ $`=`$ $`{\displaystyle \frac{\left(r^3\mathrm{cosh}^2\alpha +2a^2m+ra^2\right)r}{\left(r^2\mathrm{cosh}^2\alpha +a^2\right)\mathrm{cosh}^2\alpha }}`$
$`l_{\theta \varphi }`$ $`=`$ $`2{\displaystyle \frac{mra^2\mathrm{sinh}^2\alpha }{\left(r^2\mathrm{cosh}^2\alpha +a^2\right)\mathrm{cosh}^2\alpha }}`$
$`l_{\varphi \varphi }`$ $`=`$ $`{\displaystyle \frac{\mathrm{sinh}^2\alpha \left(r^4\mathrm{cosh}^2\alpha +r^2a^2\mathrm{cosh}^2\alpha +2mr\mathrm{cosh}^2\alpha a^2+r^2a^22mra^2+a^4\right)}{\left(r^2\mathrm{cosh}^2\alpha +a^2\right)\mathrm{cosh}^2\alpha }}`$
where $`\alpha =\alpha (r,a,m)`$ is defined by
$$\alpha =a_r^{\mathrm{}}\frac{d\nu }{\sqrt{\nu ^4+a^2\nu ^2+2a^2m\nu }}.$$
(43)
In general, all the angular components of the FL metric are non-zero at $`\theta =0`$, so the metric cannot be regular there. Now, lun suggests that the axis of symmetry is not at $`\theta =0`$ but its location is a function of $`r`$, explicitly $`\theta =\mathrm{arcsin}(\mathrm{tanh}(\alpha ))`$. We do not comment on the analytic validity of such a statement. Our aim is to find a form for the Kerr metric that can be used in numerical codes based on the Bondi-Sachs metric, and whatever the case – irregularity at $`\theta =0`$, or an axis of symmetry that is not a straight coordinate line – the metric of lun is unsuitable for this purpose.
## VII Conclusion
We have computed numerically the Ricci tensors of the metrics (39), PI and FL. In all cases, second order convergence to zero was observed. We also investigated regularity near the axis $`\theta _{}=0`$ (or $`\theta =0`$ in the FL case), finding that the metrics (39) and PI were regular, but that FL was not.
The aim of this work is to find a metric in Bondi-Sachs form that represents the Kerr geometry and that can be used in a numerical computation. The metric (39) partially meets that goal. In order to meet it completely, the radial coordinate $`r_{}`$ will need to be transformed to an area coordinate say $`r_{}`$ (i.e., so that the area of the 2-surface $`u=r_{}=\text{constant}`$ is $`4\pi r_{}^2`$); the procedure for doing so has been described in vishbook . Also, there will need to be a transformation to the angular coordinates (stereographic) actually used in the codes. Finally, the numerical difficulty of finding the metric coefficients near the axis will need to be overcome. While not trivial, this is all in principle straightforward, and is deferred to further work.
## Acknowledgments
We wish to thank Jeffrey Winicour for discussion; and Frans Pretorius for looking at an earlier draft of the paper, and informing us about an inconsistency in the calculation of a norm. NTB thanks Max-Planck-Institut für GravitationsPhysik, Albert-Einstein-Institüt, for hospitality. LRV thanks Cornelius Du Toit, Michael van Canneyt for assistance with the Lazarus and ppc386 compilers. The work was supported in part by the National Research Foundation, South Africa, under Grant number 2053724. |
warning/0506/math0506498.html | ar5iv | text | # On the algebra of quasi-shuffles
## Introduction
The shuffle algebra is a commutative algebra structure on the tensor module $`T(V)`$ of the vector space $`V`$. For instance $`uv=uv+vu`$. If $`R`$ is a commutative algebra, then the shuffle product on $`T(R)`$ can de deformed into a new product $``$ called the quasi-shuffle product. For instance
$$ab=ab+ab+ba.$$
Observe that this product is not a graded product, but only a filtered product, since $`ab+baR^2`$ and $`abR`$.
The quasi-shuffle algebra appears in different places in mathematics, for instance in the Multi-Zeta Values topic, cf. for instance , where it plays an important role.
The aim of this paper is to show, first, that the quasi-shuffle algebra over the polynomial algebra has a universal property: it is the free “Commutative TriDendriform algebra”. This notion of CTD-algebra is a commutative version of the notion of dendriform algebras, which plays a role in several themes (cf. ). A CTD-algebra is determined by two operations, denoted $``$ (left) and $``$ (dot), related by the axioms:
$$xy=yx$$
and
$`(xy)z`$ $`=`$ $`x(yz+zy+yz),`$
$`(xy)z`$ $`=`$ $`x(yz),`$
$`(xy)z`$ $`=`$ $`x(yz).`$
We show that the functor $`T^q:Com\text{-alg}CTD\text{-alg}`$ is left adjoint to the forgetful functor (which forgets $``$). A quasi-shuffle algebra is not only a CTD-algebra, but it is even a CTD-bialgebra. In the second part we state and prove a structure theorem for connected CTD-bialgebras, cf. 4.1. It means that the triple of operads
$$(As,CTD,Com)$$
is good in the sense that it satisfies the analogues of the Poincaré-Birkhoff-Witt theorem and the Cartier-Milnor-Moore theorem satisfied by the triple of operads $`(Com,As,Lie)`$.
In the last part we state and prove a non-commutative analogue of the universal property: the quasi-shuffle algebra over the tensor algebra is free in the category of involutive tridendriform algebras.
Notation. In this paper $`K`$ is a field and all vector spaces are over $`K`$. Its unit is denoted by 1. The tensor product of vector spaces over $`K`$ is denoted by $``$. The tensor product of $`n`$ copies of the space $`V`$ is denoted $`V^n`$. For $`v_iV`$ the element $`v_1\mathrm{}v_n`$ of $`V^n`$ is denoted by the concatenation of the elements: $`v_1\mathrm{}v_n`$. The *tensor module* over $`V`$ is the direct sum
$$T(V):=K1VV^2\mathrm{}V^n\mathrm{}$$
and the *reduced tensor module* is $`\overline{T}(V):=T(V)/K1`$. The *symmetric module* over $`V`$ is the direct sum
$$S(V):=K1VS^2(V)\mathrm{}S^n(V)\mathrm{}$$
and the *reduced symmetric module* is $`\overline{S}(V):=S(V)/K1`$, where $`S^n(V)`$ is the quotient of $`V^n`$ by the action of the symmetric group. We still denote by $`v_1\mathrm{}v_n`$ the image in $`S^n(V)`$ of $`v_1\mathrm{}v_nV^n`$. If $`V`$ is generated by $`x_1,\mathrm{},x_n`$, then $`S(V)`$ (resp. $`T(V)`$) can be identified with the polynomials (resp. noncommutative polynomials) in $`n`$ variables.
Acknowledgement I thank Muriel Livernet for comments on the first version of this paper. This work has been partially supported by the “Agence Nationale pour la Recherche”.
## 1. Quasi-shuffle algebra
### 1.1. Definition, notation
For any associative (not necessarily unital) algebra $`(R,)`$ the *quasi-shuffle algebra*, or *stuffle algebra*, on $`R`$ is the tensor module $`T(R)`$ equipped with the associative unital product $``$ uniquely determined by the relation:
$$axby=(ab)(xy)+a(xby)+b(axy)$$
for any $`a,bR`$, any $`x,yT(R)`$, 1 being a unit for $``$.
The product $``$ is associative and unital, we denote by $`T^q(R):=(T(R),)`$ the quasi-shuffle algebra on $`R`$. It is immediate to check that if $`(R,)`$ is commutative, then so is $`T^q(R)`$.
Explicitly the product $``$ is the following. Let $`\gamma `$ be a piece-wise linear path in the plane from $`(0,0)`$ to $`(p,q)`$ using only the steps $`(0,1),(1,0)`$ and $`(1,1)`$. Define $`\gamma (a_1\mathrm{}a_p,a_{p+1}\mathrm{}a_{p+q})`$ to be the element of $`R^r,rp+q`$, obtained by writing $`a_1`$ if the first step is $`(0,1)`$, $`a_{p+1}`$ if the first step is $`(1,0)`$, and $`a_1a_{p+1}`$ if the first step is $`(1,1)`$, etc. until the last step. Observe that, if there is no $`(1,1)`$ step, then we obtain a shuffle. If the path is the diagonal ($`p=q`$), then we get $`a_1a_{p+1}\mathrm{}a_pa_{2p}R`$. With this notation we have
$$a_1\mathrm{}a_pa_{p+1}\mathrm{}a_{p+q}=\underset{\gamma }{}\gamma (a_1\mathrm{}a_p,a_{p+1}\mathrm{}a_{p+q}).$$
The elements in the sum are sometimes called “mixable shuffles”, cf. .
### 1.2. Examples
1. If the product $``$ is zero (recall that $`R`$ need not be unital), then $`T^q(R)`$ is the shuffle algebra on the vector space $`R`$.
2. Let $`y_i`$, $`i1`$, be a set of variables and let $`R_0:=Ky_1,y_2,\mathrm{}`$ be the noncommutative polynomial algebra on a countable number of variables. Define a product $``$ on $`R_0`$ by
$$y_k\omega y_k^{}\omega ^{}=y_{k+k^{}}(\omega \omega ^{})+y_k(\omega y_k^{}\omega ^{})+y_k^{}(y_k\omega \omega ^{}),$$
see formula (94) of . Then it is easy to check that $`R_0`$ is the quasi-shuffle algebra on the algebra of non-constant polynomials in one variable.
Let us denote by $`T^c(V)`$ the cofree coalgebra on the vector space $`V`$. As a vector space it is the tensor module $`T(V)`$. The comultiplication is the *deconcatenation* given by
$$\mathrm{\Delta }(v_1\mathrm{}v_n):=\underset{i=0}{\overset{n}{}}v_1\mathrm{}v_iv_{i+1}\mathrm{}v_n.$$
###### 1.3 Proposition.
The unique coalgebra morphism
$$T^c(R)T^c(R)T^c(R)$$
induced by the algebra morphism $`:RRR`$ is the quasi-shuffle product $``$.
* *Proof.* Since $`T^c(R)`$ is cofree in the category of connected coassociative coalgebras, there is a unique coalgebra morphism $`\mu :T^c(R)T^c(R)T^c(R)`$ which makes the following diagram commutative:
Since $``$ is a coalgebra morphism, in order to show that $`\mu =`$, it is sufficient to show that the two composites $`T^c(R)T^c(R)RRR`$ and $`T^c(R)T^c(R)T^c(R)R`$ are equal. By definition $`ab=ab+ba+ab`$ for any $`a,bR`$, hence the projection of this element in $`R`$ is $`ab`$ and we are done. $`\mathrm{}`$
### 1.4. Remark
If, instead of starting with an algebra $`R`$, we start with a coalgebra $`C`$, then the dual construction gives rise to a bialgebra structure on the tensor algebra $`T(C)`$. This construction is used in to construct renormalization Hopf algebras in various settings. See also where such objects were studied.
### 1.5. Operations on the quasi-shuffle algebra
Let us denote by $``$ (called left), by $``$ (called right) and by $``$ (called dot) the three binary operations on $`T(R)`$ uniquely determined by the following requirements (inductive definition):
$$\begin{array}{cccccc}\hfill axby& =& (ab)(xy)\hfill & and& 1x=0& x1=0,\\ \hfill axby& =& a(xby)\hfill & and& 1x=0& x1=x,\\ \hfill axby& =& b(axy)\hfill & and& 1x=x& x1=0,\end{array}$$
where, for any $`x,yT(R)`$, we have defined
$$xy:=xy+xy+xy.$$
It is immediate to check that this operation is the same as the quasi-shuffle introduced in 1.1.
###### 1.6 Proposition.
The three operations $``$, $``$ and $``$ on $`T(R)`$ satisfy the following relations:
(1) $`(xy)z`$ $`=`$ $`x(yz)`$
(2) $`(xy)z`$ $`=`$ $`x(yz)`$
(3) $`(xy)z`$ $`=`$ $`x(yz)`$
(4) $`(xy)z`$ $`=`$ $`x(yz)`$
(5) $`(xy)z`$ $`=`$ $`x(yz)`$
(6) $`(xy)z`$ $`=`$ $`x(yz)`$
(7) $`(xy)z`$ $`=`$ $`x(yz).`$
* *Proof.* The proof is by induction on the degree of the elements (an element of $`V^n`$ is of degree $`n`$). First, we prove the seven relations when the elements are in $`R`$.
1. $`(ab)c=(ab)c=a(bc)=a(bc)`$,
2. $`(ab)c=(ba)c=b(ac)=a(bc)=a(bc)`$,
3. $`(ab)c=c(ab)=a(cb)=a(bc)`$,
4. $`(ab)c=(ab)c=a(bc)=a(bc)`$,
5. $`(ab)c=(ab)c=(ac)b=a(cb)=a(bc)`$,
6. $`(ab)c=(ba)c=(bc)a=a(bc)`$,
7. $`(ab)c=a(bc)`$.
Second, we show by induction the first relation for elements in $`T(R)`$. On the one hand we have
$`(axy)z`$ $`=`$ $`a(xy)z`$
$`=`$ $`a((xy)z)`$
On the other hand we have
$`ax(yz)`$ $`=`$ $`a(x(yz)).`$
Since $``$ is associative by the inductive assumption, we get the first relation $`(xy)z=x(yz)`$. The other relations are proved similarly. $`\mathrm{}`$
###### 1.7 Proposition.
If $`(R,)`$ is commutative, then $`xy=yx`$, $`xy=yx`$ and $`xy=yx`$ for any $`x,yT(R)`$. Under this hypothesis the seven relations above come down to the following three:
$`(xy)z`$ $`=`$ $`x(yz),`$
$`(xy)z`$ $`=`$ $`x(yz),`$
$`(xy)z`$ $`=`$ $`x(yz).`$
* *Proof.* The first part of the statement is immediate. For the second part: by the symmetry hypothesis the 7 relations can be written in terms of $``$ and $``$. Then relation (3) is equivalent to relation (1), and relation (2) is implied by relation (1). Relation (6) is equivalent to relation (4) and relation (5) is implied by relation (4). So we are left with relations (1), (4) and (7) which can be written as in the statement because of the symmetry hypothesis. $`\mathrm{}`$
### 1.8. Remark
A $`𝐁_{\mathrm{}}`$-algebra $`R`$ is a vector space equipped with $`(p+q)`$-ary operations $`M_{pq}`$ verifying the relations which enable us to put a bialgebra structure on $`T^c(R)`$. If $`M_{pq}=0`$ for $`p2`$, then it is called a *brace algebra*. If $`M_{pq}=0`$ for $`p2`$ and $`q2`$, then only $`M_{11}`$ survives and $`R`$ is an associative algebra. The general procedure to construct the bialgebra out of the $`𝐁_{\mathrm{}}`$-algebra (cf. for instance ), gives precisely the quasi-shuffle algebra in the latter case.
## 2. Commutative TriDendriform algebra
The properties of the quasi-shuffle algebra shown in the previous section suggest the notion of commutative tridendriform algebra that we now introduce. We show that the free CTD-algebra is the quasi-shuffle algebra over the polynomials.
The notion of dendriform algebra was first introduce in in connection with some problems in algebraic $`K`$-theory, and the notion of tridendriform algebra was introduced in .
### 2.1. Definition
A *commutative tridendriform algebra* or $`ComTriDend`$-algebra (CTD-algebra for short) is a vector space $`A`$ equipped with two operations $``$ and $``$ satisfying:
$$xy=yx$$
and
(8) $`(xy)z`$ $`=`$ $`x(yz+zy+yz),`$
(9) $`(xy)z`$ $`=`$ $`x(yz),`$
(10) $`(xy)z`$ $`=`$ $`x(yz).`$
It is immediate to check from these relations that the product $``$, defined as $`xy:=xy+yx+xy`$, is associative and commutative.
It will prove helpful to introduce the notion of *unital CTD-algebra* which is a vector space of the form $`A=K1\overline{A}`$ where $`\overline{A}`$ is CTD-algebra and the operations $``$ and $``$ are extended to $`1`$ by the following formulas:
$$1x=0=x1,1x=0,x1=x\text{ for any }x\overline{A}.$$
Observe that $`11`$ and $`11`$ are not defined but $`1`$ is a unit for $``$. Therefore $`\overline{A}`$ is the augmentation ideal of $`A`$.
### 2.2. Examples
(a) From Proposition 1.7 it is clear that the quasi-shuffle algebra over a commutative ring is a CTD-algebra. In particular the algebra of shuffles and $`T^q(\overline{S}(V))`$ are CTD-algebras.
(b) Let $`(R,)`$ be a commutative algebra and let $`P:RR`$ be a *Rota-Baxter operator*, i.e. a linear map which satisfies the identity:
$$P(a)P(b)=P(aP(b)+P(a)b+ab).$$
If we define $`ab:=aP(b)`$, then one can check easily that $`(R,,)`$ is a $`CTD`$-algebra, cf. . Moreover the operator $`P`$ becomes an associative algebra morphism: $`P:(R,)(R,)`$.
### 2.3. Free CTD-algebra
Let $`V`$ be a vector space. By definition the *free CTD-algebra* over $`V`$ is a CTD-algebra $`CTD(V)`$ equipped with a linear map $`\iota :VCTD(V)`$ which satisfies the following universal condition:
any linear map $`\varphi :VA`$, where $`A`$ is a CTD-algebra, admits a unique extension $`\mathrm{\Phi }:CTD(V)A`$ as morphism of CTD-algebras:
In the unital case it is assumed that the image of $`\varphi `$ lies in the augmentation ideal of $`A`$.
### 2.4. The operad CTD
Let us work over a field of characteristic zero field. Since the relations of a CTD-algebra are multilinear, the $`n`$-multilinear part $`CTD(n)`$ of $`CTD(Kx_1\mathrm{}Kx_n)`$ inherits a structure of $`S_n`$-module from the action of the symmetric group on the set of variables $`\{x_1,\mathrm{},x_n\}`$, and we have
$$CTD(V)=\underset{n}{}CTD(n)_{S_n}V^n.$$
So the functor $`CTD`$ is a Schur functor. It is called the operad of the CTD-algebras (cf. the first Chapter of for instance). As a representation of $`S_n`$, $`CTD(n)`$ is spanned by the ordered-unordered partitions of $`\{1,\mathrm{},n\}`$. The dimension $`d_n`$ of $`CTD(n)`$ is sometimes called the *Fubini number* (sequence number A000670 in the On-Line Encyclopedia of Integer Sequences):
$$\begin{array}{cccccccc}n& 1& 2& 3& 4& 5& 6& \mathrm{}\\ d_n& 1& 3& 13& 75& 541& 4683& \mathrm{}\end{array}$$
It can be shown combinatorially that the generating series $`f^{CTD}(x):=_{n1}\frac{d_n}{n!}x^n`$ is given by:
$$f^{CTD}(x)=\frac{\mathrm{exp}(x)1}{2\mathrm{exp}(x)}.$$
It will be a consequence of our main result, cf. 4.2.
###### 2.5 Theorem.
The free unital CTD-algebra on the vector space $`V`$ is isomorphic to the quasi-shuffle algebra over the reduced symmetric algebra on $`V`$ (non-constant polynomials):
$$\mathrm{\Phi }:CTD(V)T^q(\overline{S}(V)).$$
* *Proof.* Since by Proposition 1.7 the quasi-shuffle algebra $`T^q(\overline{S}(V))`$ is a CTD-algebra and since it contains $`V`$, the inclusion map $`\iota `$ induces a CTD-morphism $`\mathrm{\Phi }(V):CTD(V)T^q(\overline{S}(V)).`$ It is clear from the definition of the operations $``$ and $``$ in $`T^q(\overline{S}(V))`$ that $`\mathrm{\Phi }`$ is surjective. So we need only to check that it is injective. From the naturality of $`\mathrm{\Phi }`$ it suffices to check injectivity on the multilinear parts (since both functors are Schur functors). We denote them by $`CTD(n)`$ and $`T^q\overline{S}(n)`$ respectively.
Let $`V`$ be an $`n`$-dimensional vector space, spanned by $`v_1,\mathrm{},v_n`$. In $`T^q(\overline{S}(V))`$ the multilinear part of degree $`n`$ admits as a basis the vectors
$$(w_1\mathrm{}w_{i_1})(w_{i_1+1}\mathrm{}w_{i_1+i_2})\mathrm{}(\mathrm{}\mathrm{}w_{i_1+\mathrm{}+i_k})$$
where the $`w_i`$’s are a permutation of the $`v_i`$’s and the subwords in parenthesis are in increasing order. We call such an element an *ordered-unordered partition*. Example: $`(v_2v_5)(v_1v_4)(v_3)`$. It suffices to show that $`dimCTD(n)dimT^q\overline{S}(n)=:d_n`$.
First, we show that any monomial in $`CTD(V)`$ can be written as
$$x_1(x_2(\mathrm{}(x_{k1}x_k)\mathrm{}))$$
where $`x_i`$ is a monomial of the form $`v_{i_1}\mathrm{}v_{i_l}`$ where the indices $`i_j`$ are in increasing order.
Any monomial of degree larger than 1 is of the form $`XY`$ or $`XY`$ for elements of strictly smaller degree. We freely use the fact that the operation $``$ is associative and commutative. We work by induction on the degree of the components.
(a) Case $`XY`$.
(a1) If $`X=v_{i_1}\mathrm{}v_{i_l}`$, we are done since the degree of $`Y`$ is strictly smaller than the degree of $`XY`$.
(a2) Otherwise, either $`X=X_1X_2`$ where the symbol $``$ appears in $`X_1`$ or $`X_2`$, or $`X=X_1X_2`$. In the second case we apply relation (8) and we are back to case (a) with the degree of the second component being strictly smaller. In the first case, on can suppose that $`X_1=X_1^{}X_1^{\prime \prime }`$. We apply relation (9) to obtain $`((X_1^{}X_2)X_1^{\prime \prime })Y`$ to which we apply relation (8). Again, we are back to case (a) with the degree of the first component strictly smaller.
(b) Case $`XY`$. If the product $``$ does not appear in $`X`$ nor $`Y`$ we are done. Otherwise, assume that $`X=X_1X_2`$. We use relation (9) to get back to case (a) as before and we are done.
We have shown that $`CTD(n)`$ is spanned by $`d_n`$ elements, where $`d_n`$ is the number of ordered-unordered partitions. Since $`CTD(n)T^q\overline{S}(n)`$ is surjective we are done. $`\mathrm{}`$
### 2.6. Universal enveloping CTD-algebra
Let $`A`$ be a CTD-algebra. Ignoring the left operation yields a forgetful functor from the category of CTD-algebras to the category of commutative algebras (not necessarily unital). We denote by $`U_{CTD}`$ the left adjoint of this forgetful functor:
$$U_{CTD}:Com\text{-alg}CTD\text{-alg}.$$
Explicitly, for any commutative algebra $`R`$, the CTD-algebra $`U_{CTD}(R)`$ is given by
$$U_{CTD}(R)=CTD(R)/$$
where the equivalence relation $``$ consists in identifying, for any $`a,bR`$, the product in $`R`$, $`abRCTD(R)`$, with the dot-product in $`CTD(R)`$, $`abCTD(R)`$. Observe that only the vector space structure of $`R`$ is used to construct the free CTD-algebra $`CTD(R)`$. The commutative structure of $`R`$ is used when making the quotient.
###### 2.7 Theorem.
For any (not necessarily unital) commutative algebra $`R`$ there is an isomorphism of CTD-algebras:
$$U_{CTD}(R)T^q(R).$$
* *Proof.* The inclusion $`(R,)(T^q(R),)`$ determines an adjoint map $`U_{CTD}(R)T^q(R)`$ which is a CTD-morphism. In order to prove that it is an isomorphism it suffices, by a classical argument, cf. Appendix B, to prove it for the free commutative algebra $`R=\overline{S}(V)`$. Since $`U_{CTD}`$ and $`\overline{S}`$ are both left adjoint functors, their composite is left adjoint to the forgetful functor $`CTD\text{-alg}Vect`$ and therefore we have $`U_{CTD}(\overline{S}(V))=CTD(V)`$. The restriction of the adjoint map $`CTD(V)=U_{CTD}(\overline{S}(V))T^q(\overline{S}(V))`$ to $`V`$ is simply the inclusion $`V\overline{S}(V)T^q(\overline{S}(V))`$, therefore this adjoint map is the isomorphism $`\mathrm{\Phi }`$ of Theorem 2.5. $`\mathrm{}`$
## 3. CTD-bialgebras
We use the technique developed in to show that the free CTD-algebra is in fact a Hopf algebra. We introduce the notion of CTD-bialgebra, which is a CTD-algebra equipped with a compatible coproduct. Quasi-shuflle algebras are CTD-bialgebras. In this section we always assume unitality.
### 3.1. Tensor product of CTD-algebras
Let $`A`$ and $`B`$ be two CTD-algebras. One can put a CTD-algebra structure on $`AB`$ as follows:
$`(ab)(a^{}b^{})`$ $`=`$ $`(aa^{})(bb^{}),`$
$`(ab)(a^{}b^{})`$ $`=`$ $`(aa^{})(bb^{}).`$
Since the elements $`11`$ and $`11`$ do not make sense, we take the following convention:
$`11bb^{}`$ $`:=`$ $`1(bb^{}),`$
$`11bb^{}`$ $`:=`$ $`1(bb^{}).`$
This trick is due to M. Ronco. Observe that if $`A,B`$ and $`C`$ are three CTD-algebras, then the two CTD-algebra structures that one can put on $`ABC`$ are the same.
### 3.2. Bialgebra structure on the free CTD-algebra
Let CTD(V) be the free (unital) CTD-algebra on $`V`$. Since $`CTD(V)CTD(V)`$ is a CTD-algebra there is a unique CTD-morphism
$$\mathrm{\Delta }:CTD(V)CTD(V)CTD(V)$$
which extends the linear map $`VCTD(V)CTD(V)`$ given by $`vv1+1v`$.
As in we see that $`\mathrm{\Delta }`$ is a coassociative counital morphism.
This example justifies the following definition.
### 3.3. $`As^c\text{-}CTD`$-bialgebra
By definition an *$`As^c\text{-}CTD`$-bialgebra*, or CTD-bialgebra for short, is a vector space $`=K1\overline{}`$ equipped with a structure of unital CTD-algebra, a structure of coassociative counital coalgebra, and these two structures are related by the following *compatibility relation*:
$`\mathrm{\Delta }(xy)`$ $`=`$ $`x_{(1)}y_{(1)}x_{(2)}y_{(2)},`$
$`\mathrm{\Delta }(xy)`$ $`=`$ $`x_{(1)}y_{(1)}x_{(2)}y_{(2)},`$
where we use the notation $`\mathrm{\Delta }(x)=x_{(1)}x_{(2)}`$. Here we adopt the convention stated in 3.1 and $`\mathrm{\Delta }(1)=11`$.
Following Quillen (cf. , p. 282) we say that a coaugmented coalgebra $``$ is connected if $`=_{r0}F_r`$ where $`F_r`$ is the coradical filtration of $``$ defined recursively by the formulas
$`F_0`$ $`:=`$ $`K1,`$
$`F_r`$ $`:=`$ $`\{x\mathrm{\Delta }(x)x11xF_{r1}F_{r1}\}.`$
From the fact that $`CTD(V)`$ is an $`As^c\text{-}CTD`$-bialgebra, cf. 3.2, it follows that $`U_{CTD}(R)`$ is also an $`As^c\text{-}CTD`$-bialgebra.
###### 3.4 Proposition.
For any CTD-algebra $``$ the dot operation is stable on the primitive part $`\mathrm{Prim}:=\{x|\mathrm{\Delta }(x)x11x=0\}`$. So $`\mathrm{Prim}`$ is a commutative algebra.
* *Proof.* It suffices to compute for $`x`$ and $`y`$ primitive:
$`\mathrm{\Delta }(xy)`$ $`=`$ $`(x1+1x)(y1+1y)`$
$`=`$ $`xy11+x11y+1yx1+11xy`$
$`=`$ $`xy1+0+0+1xy`$
$`\mathrm{}`$
###### 3.5 Proposition.
The isomorphism $`\mathrm{\Phi }:U_{CTD}(R)T^q(R)`$, cf. Theorem 2.5, is compatible with the coproduct, and so is an isomorphism of CTD-bialgebras.
* *Proof.* Since the coproduct on $`U_{CTD}(R)`$ is induced by the coproduct on $`CTD(R)`$, it suffices to show that $`\mathrm{\Phi }:CTD(R)T^q(R)`$ is compatible with the coproduct. From the definition of the coproduct on $`CTD(R)`$, cf. 3.2, it suffices to show that the coproduct on $`T^q(R)`$, that is deconcatenation, satisfies the compatibility relations of 3.3. We adopt the notation $`a(xy):=(ax)y`$ for $`aR`$ and $`xyT(R)T(R)`$ and we remark that $`\mathrm{\Delta }(az)=a\mathrm{\Delta }(z)+1az`$. For the dot operation we already verified the compatibility for elements of $`R`$ in Proposition 3.4. For elements of $`T^q(R)`$ we have on the one hand:
$`\mathrm{\Delta }(axby)`$ $`=`$ $`\mathrm{\Delta }((ab)(xy))`$
$`=`$ $`(ab)\mathrm{\Delta }(xy)+1(ab)(xy)`$
On the other hand we get:
$`\mathrm{\Delta }(ax)\mathrm{\Delta }(by)`$ $`=`$ $`(a\mathrm{\Delta }(x)+1ax)(b\mathrm{\Delta }(y)1+1by)`$
$`=`$ $`(ab)\mathrm{\Delta }(x)\mathrm{\Delta }(y)+1axby`$
$`=`$ $`(ab)\mathrm{\Delta }(xy)+1(ab)(xy)`$
Therefore we have proved that $`\mathrm{\Delta }(axby)=\mathrm{\Delta }(ax)\mathrm{\Delta }(by)`$ as expected.
For the left operation, we first verify it on the elements of $`R`$.
One hand hand we get:
$`\mathrm{\Delta }(ab)`$ $`=`$ $`\mathrm{\Delta }(ab)`$
$`=`$ $`ab1+ab+1ab.`$
On the other hand we get:
$`\mathrm{\Delta }(a)\mathrm{\Delta }(b)`$ $`=`$ $`(a1+1a)(b1+1b)`$
$`=`$ $`ab1+a11b+1ba1+11ab`$
$`=`$ $`ab1+a1b+0+1ab`$
$`=`$ $`ab1+ab+1ab.`$
We now prove this relation for any elements of $`T^q(R)`$ by induction. On the one hand we get:
$`\mathrm{\Delta }(axby)`$ $`=`$ $`\mathrm{\Delta }(a(xby))`$
$`=`$ $`a\mathrm{\Delta }(xby)+1a(xby)`$
$`=`$ $`a(\mathrm{\Delta }(x)\mathrm{\Delta }(by))+1a(xby)`$
$`=`$ $`a(\mathrm{\Delta }(x)b\mathrm{\Delta }(y))+a(\mathrm{\Delta }(x)(1by))+1a(xby)`$
On the other hand we get:
$`\mathrm{\Delta }(ax)\mathrm{\Delta }(by)`$ $`=`$ $`(a\mathrm{\Delta }(x)+1ax)(b\mathrm{\Delta }(y)+1by)`$
$`=`$ $`a\mathrm{\Delta }(x)b\mathrm{\Delta }(y)+a\mathrm{\Delta }(x)(1by)+1axby`$
$`=`$ $`a(\mathrm{\Delta }(x)b\mathrm{\Delta }(y))+a(\mathrm{\Delta }(x)(1by))+1a(xby).`$
Therefore we have proved that $`\mathrm{\Delta }(axby)=\mathrm{\Delta }(ax)\mathrm{\Delta }(by)`$ as expected. $`\mathrm{}`$
## 4. Structure theorem for CTD-bialgebras
We prove a structure theorem for $`As^c\text{-}CTD`$-bialgebras analogous to the structure theorem for cocommutative Hopf algebras (PBW+CMM theorem), cf. . Here the structure of the primitive part is the commutative structure (instead of the Lie structure in the classical case). Therefore $`(As,CTD,Com)`$ is a good triple of operads in the sense of .
###### 4.1 Theorem.
We suppose that the ground field $`K`$ is of characteristic zero. If $``$ is a CTD-bialgebra over $`K`$, then the following are equivalent:
(a) $``$ is a connected CTD-bialgebra,
(b) $``$ is isomorphic to $`U_{CTD}(\mathrm{Prim})`$ as a CTD-bialgebra,
(c) $``$ is cofree among the connected coalgebras: $`T^c(\mathrm{Prim})`$.
* *Proof.* We apply the main theorem of to the $`As^c\text{-}CTD`$-bialgebra type. In order to get the equivalence of the three assertions, it is sufficient to verify three hypotheses, which, in our case, read as follows:
(H0) the compatibility relation is distributive,
(H1) the free CTD-algebra is naturally a $`As^c\text{-}CTD`$-bialgebra,
(H2epi) the natural coalgebra map $`\phi (V):CTD(V)T^c(V)`$ admits a natural splitting.
Hypothesis (H0) is immediate by direct inspection, cf. 3.3. Hypothesis (H1) is immediate by construction, cf. 3.2. Before proving hypothesis (H2epi), let us recall the construction of $`\phi (V)`$. Since by hypothesis (H1), $`CTD(V)`$ is a coassociative coalgebra, the natural projection $`CTD(V)V`$ induces a unique coalgebra map $`\phi (V):CTD(V)T^c(V)`$. By Theorem 2.5 there is an isomorphism $`CTD(V)T(\overline{S}(V))`$. The projection $`CTD(V)V`$ is the composite of the projections $`T(\overline{S}(V))\overline{S}(V)V`$. The splitting $`s:T^c(V)T(\overline{S}(V))`$ is simply induced by the inclusion $`V\overline{S}(V)`$ to which we apply the functor $`T^c`$. $`\mathrm{}`$
###### 4.2 Corollary.
There is an isomorphism of Schur functors $`CTD=AsCom`$ and therefore:
$$f^{CTD}(x)=\frac{\mathrm{exp}(x)1}{2\mathrm{exp}(x)}.$$
* *Proof.* Since the composition of two left adjoint functors is still a left adjoint functor, there is an isomorphism $`CTD(V)U_{CTD}(Com(V))`$ and therefore an isomorphism $`CTD(V)T^c(Com(V))`$. As a consequence we get $`f^{CTD}=f^{As}f^{Com}`$, since the generating series $`f^{As}(x)`$ of $`\overline{T}`$ is the geometric series $`\frac{x}{1x}`$ and the generating series $`f^{Com}(x)`$ of $`Com`$, that is $`\overline{S}`$, is the exponential series $`\mathrm{exp}(x)1`$. $`\mathrm{}`$
### 4.3. Remark
If we are interested only in the Hopf algebra structure, then $`T^q(R)`$ is isomorphic to the shuffle algebra $`T^{sh}(R)`$ of the vector space $`R`$. Indeed the dual Hopf algebra of $`T^q(R)`$ is connected and cocommutative. So, by the Cartier-Milnor-Moore theorem, it is isomorphic to the universal enveloping algebra of its primitive part. Since this primitive part is the free Lie algebra $`Lie(R)`$, we get $`U(Lie(R))=T(R)`$ where only the vector space structure of $`R`$ is involved. Since the dual of the tensor Hopf algebra is the shuffle Hopf algebra, we are done.
### 4.4. Triple of operads
Theorem 4.1 provides a new example of “good triple of operads” as defined in :
$$(As,CTD,Com).$$
Its quotient triple is $`(As,Zinb,Vect)`$, where $`Zinb`$ is the operad of Zinbiel algebras. Recall that a Zinbiel algebra (cf. ) is a vector space equipped with a binary operation $``$ verifying
$$(xy)z=x(yz+zy).$$
Since $`(As,CTD,Com)`$ is a good triple, so is $`(As,Zinb,Vect)`$ by a Theorem of . It is also a consequence of a theorem due to M. Ronco which asserts that the triple $`(As,Dend,Brace)`$ is good, where $`Brace`$ is the operad of brace algebras, cf. . For a self-contained proof concerning $`(As,Zinb,Vect)`$, see .
### 4.5. Graded version
Until now we have worked in the symmetric category of vector spaces over $`K`$. We could as well work in the symmetric category of *graded* vector spaces over $`K`$, taking the sign into account in the symmetry operator: $`xy(1)^{|x||y|}yx`$. In this context the first relation of a graded CTD-algebra reads
$$(xy)z=x(yz+(1)^{|y||z|}zy+yz)$$
on homogeneous elements. The other two relations are unchanged and the dot operation is supposed to be graded-symmetric.
All the preceding results can be written in the graded framework, provided that the functors are replaced by their graded analogues. For instance, the symmetric functor $`S`$ has to be replaced by the graded-symmetric functor $`\mathrm{\Lambda }`$ given by
$$\mathrm{\Lambda }(V^{ev}V^{odd})=S(V^{ev})E(V^{odd})$$
where $`E`$ is the exterior algebra functor.
## 5. Structure of non-commutative quasi-shuffle algebras
### 5.1. Definition
A *tridendriform algebra* (also called dendriform trialgebra) is a vector space $`A`$ over $`K`$ equipped with three binary operations $`,,`$ satisfying the following relations:
$`(xy)z`$ $`=`$ $`x(yz)`$
$`(xy)z`$ $`=`$ $`x(yz)`$
$`(xy)z`$ $`=`$ $`x(yz)`$
$`(xy)z`$ $`=`$ $`x(yz)`$
$`(xy)z`$ $`=`$ $`x(yz)`$
$`(xy)z`$ $`=`$ $`x(yz)`$
$`(xy)z`$ $`=`$ $`x(yz),`$
where $`xy:=xy+xy+xy`$. Observe that the product $``$ is associative. As in the commutative case (cf. ) we can introduce a unit by requiring:
$`1x=0=x1,`$
$`1x=0,x1=x,`$
$`1x=x,x1=0.`$
### 5.2. Examples
(a) The quasi-shuffle algebra $`T^q(R)`$ over the (nonunital) associative algebra $`R`$ is a tridendriform algebra by Proposition 1.6.
(b) The free tridendriform algebra on one generator can be described explicitly in terms of planar trees, cf. .
(c) Let $`(R,)`$ be a (nonunital) associative algebra and let $`P`$ be a Baxter operator on $`R`$, cf. 2.2. If we put $`ab:=aP(b)`$ and $`ab:=P(a)b`$, then $`(R,,,)`$ is a tridendriform algebra.
(d) A CDT-algebra is a particular case of tridendriform algebra as proved in 1.7.
### 5.3. Tridendriform algebra with involution
An *involution* on a tridendriform algebra $`A`$ is a linear map $`a\iota (a)`$ verifying the following conditions:
$`\iota (xy)`$ $`=`$ $`\iota (y)\iota (x)`$
$`\iota (xy)`$ $`=`$ $`\iota (y)\iota (x)`$
$`\iota (xy)`$ $`=`$ $`\iota (y)\iota (x)`$
Observe that on a CTD-algebra the identity is an involution.
###### 5.4 Proposition.
Let $`R`$ be an associative algebra equipped with an involution $`\iota `$, i.e. $`\iota (ab)=\iota (b)\iota (a)`$. Extending $`\iota `$ to $`T(R)`$ by
$$\iota (a_1\mathrm{}a_n)=\iota (a_1)\mathrm{}\iota (a_n)$$
(the order of the $`a_i`$’s is unchanged) gives a structure of tridendriform algebra with involution to $`T^q(R)`$.
* *Proof.* We work by induction. We compute on one hand:
$`\iota (axby)`$ $`=`$ $`\iota (a(xby))`$
$`=`$ $`\iota (a)(\iota (xby))`$
$`=`$ $`\iota (a)(\iota (by)\iota (x))`$
$`=`$ $`\iota (a)(\iota (b)\iota (y)\iota (x))`$
and on the other hand:
$`\iota (by)\iota (ax)`$ $`=`$ $`\iota (b)\iota (y)\iota (a)\iota (x)`$
$`=`$ $`\iota (a)(\iota (b)\iota (y)\iota (x)),`$
Hence we have proved that $`\iota (axby)=\iota (by)\iota (ax)`$. The other cases are similar. $`\mathrm{}`$
### 5.5. Example
Let $`\overline{T}(V)`$ be the tensor algebra over $`V`$ without constants. Equip $`\overline{T}(V)`$ with the involution
$$\iota (v_1\mathrm{}v_n)=v_nv_{n1}\mathrm{}v_1$$
(the order of the $`v_i`$’s is reversed). So $`T^q(\overline{T}(V))`$ is an involutive tridendriform algebra by Proposition 5.4.
###### 5.6 Theorem.
The quasi-shuffle algebra $`T^q(\overline{T}(V))`$ is free over $`V`$ (equipped with the trivial involution) in the category of involutive tridendriform algebras with involution.
* *Proof.* Let $`ITD(V)`$ be the the free involutive tridendriform algebra on the space $`V`$ equipped with the trivial involution. Since $`T^q(\overline{T}(V))`$ is an involutive tridendriform algebra, the inclusion $`V\overline{T}(V)T^q(\overline{T}(V))`$ induces a map $`ITD(V)T^q(\overline{T}(V))`$. Then the proof is similar to the proof of Theorem 2.5. In this variation the multilinear space $`T^q\overline{T}(n)`$ is spanned by the ordered-ordered partitions of $`\{1,\mathrm{},n\}`$ and so its dimension is $`2^{n1}n!`$. $`\mathrm{}`$ |
warning/0506/astro-ph0506134.html | ar5iv | text | # Bright X-ray flares in Orion young stars from COUP: evidence for star-disk magnetic fields?
## 1. Introduction
Beginning in their protostellar (Class I) phase, young stellar objects (YSOs) have long been known as copious sources of X-ray emission (Feigelson & Montmerle, 1999; Favata & Micela, 2003). The presence of frequent flares, observed both in accreting and non-accreting sources, shows that the emitting plasma is magnetically confined, and in most cases also magnetically heated. In more evolved (Class III, non-accreting) YSOs the X-ray emission mechanism is likely to be very similar to the one operating in similarly active older stars, involving a scaled up version of the solar corona, with the emitting plasma being entirely confined in magnetic loop structures and heated by similar mechanisms as in less active stars (most likely magnetic energy dissipation caused by the shear at the loop’s footpoints, e.g. Peter et al., 2004). In accreting sources (Class I and Class II), on the other hand, recent evidence points to the X-ray emission mechanism being, at least in part, different: statistically, in the ONC (Flaccomio et al., 2003; Preibisch et al., 2005) accreting YSOs are observed to be less X-ray luminous than non-accreting ones, although the actual mechanism causing this difference is not understood. Two accreting YSOs have been observed to date in X-rays at high spectral resolution: TW Hya, (Kastner et al., 2002; Stelzer & Schmitt, 2004), and BP Tau (Schmitt et al., 2005). In both cases the O vii He-like triplet shows very low $`f/i`$ ratios, which can be caused by very high densities and/or by a very high UV ambient flux. Such low $`f/i`$ is not observed in any “normal” (ZAMS, MS or more evolved) coronal source. As e.g. discussed by Drake (2005) both effects are most likely associated with the accretion stream (obviously not present in non-accreting stars). While the X-ray emission from TW Hya is very soft ($`T3`$ MK), and thus could be explained as (largely) driven by the accretion shock, both cool and hot plasma (up to 30 MK) is present in BP Tau; the hot plasma cannot originate in the accretion shock (not enough gravitational energy is available), and thus in these cases accretion-driven X-ray emission likely coexists with a more “canonical” corona.
An obvious question to ask, also in the light of the recent results concerning TW Hya and BP Tau, is therefore whether the different types of YSOs have different types of X-ray emitting structures (and thus of confining magnetic fields) in terms of size, location, structuring and density. One of the few diagnostic tools which can provide an insight in this question is the analysis of intense X-ray flares. In this context, a flare is defined as a sudden impulsive rise in the plasma temperature, immediately followed by a rise in the X-ray luminosity of the source (typically by a factor of at least a few times the quiescent luminosity), and later followed by a slower, roughly exponential decay of both the X-ray luminosity and temperature. These types of events are often observed both in the Sun (see e.g. the review by Priest & Forbes, 2002) and in most active stars, and their study has been the subject of a copious literature (see e.g. the relevant chapter reviewing the topic in Favata & Micela, 2003).
The *Chandra* Orion Ultradeep Project (COUP) is an unique long (13 days span with 9.7 days of effective exposure) X-ray observation of the Orion Nebula Cluster (ONC). The details of the COUP observation are described by Getman et al. (2005). With its long, nearly uninterrupted observation, the COUP sample offers an unique opportunity to study flares in YSOs, and thus to study the type of magnetic structures present in young stars, and in particular in stars which are either actively accreting or which are surrounded by inactive disks. The magnetic field structure in YSOs, is expected to be different from the one in older, isolated stars. Magnetic field lines (and thus possibly magnetic loops) connecting the stellar photosphere and the disk have been postulated in the context of the magnetospheric (Hartmann et al., 1994) model of accretion, in which the accreting plasma originates from the disk and is channeled onto the stellar photosphere by the magnetic field. Magnetically funneled accretion both explains the observed broad optical emission line profiles (Hartmann et al., 1994) and the ejection of high-velocity bipolar outflows (Shu et al., 2000).
However, no direct observational evidence for such extended magnetic structures has been available up to now. Strong magnetic fields have been measured in the photosphere of YSOs using Zeeman splitting, and fields strengths up to 6 kG have been determined (see the review of Johns-Krull & Valenti, 2005; also Johns-Krull et al., 2004). Such measurements cannot however determine the field’s structure, and whether it is dominated by a large-scale, dipole structure, or by higher-order multipoles imposing a smaller scale to the field. In particular, the measurement of photospheric fields does not allow determination of how the field extends into the circumstellar region, and in particular into the disk. Magnetic flux tubes are required by the magnetospheric model of accretion to channel accreting material from the inner edge of the disk to the photosphere. Analysis of the flares’ decay can help to determine if the same flux tubes are also the seat of more energetic processes, i.e. if the plasma is heated to temperatures resulting in X-ray emission and in the attendant flaring activity.
A number of intense X-ray flares on PMS stars have been previously observed and studied. Prior to the XMM-*Newton* and *Chandra* observations some 6 flares on PMS stars had been analyzed in detail (see review by Favata & Micela, 2003), although a large number of fainter events (for which therefore physical parameters could not be derived) had also been observed – see e.g. Stelzer et al. (2000) for a survey of flares in the Taurus-Auriga-Perseus complex. Favata et al. (2001) analyzed four of the bright events on YSOs known at the time, finding, in all cases, that the flaring loop structures were of modest size ($`LR_{}`$), similar to what observed in active main sequence and more evolved stars.
A number of flaring events on YSOs have been observed with *Chandra* and XMM-*Newton* (see e.g. Imanishi et al., 2003 for a *Chandra* study of flares in the $`\rho `$ Oph region), but very few have been analyzed in detail. Grosso et al. (2004) analyzed a flare observed with *Chandra* on the YSO LkH$`\alpha `$312 in M78, also finding modest-sized loop structures ($`L=0.2`$$`0.5R_{}`$).
Therefore, the analysis of the flares observed up to now on YSOs has resulted in magnetic structures which are in all respects similar to the ones found in older stars, supporting a view in which the coronae of YSOs are similar to the ones of their more evolved counterparts, and supplying no evidence for the long magnetic structures expected to connect the star with its disk. As it will be shown in this paper, the analysis of the large flares occurring in the COUP sample provides a different picture: in addition to flares occurring in star-sized or smaller structures, similar to what found up to now in both PMS and older stars, a number of events are found taking place in very long magnetic structures, with lengths of several stellar radii. As discussed in Sect. 6, loops of this size around fast rotating stars are unlikely to be stable if anchored on the star alone, and therefore such structures are likely to be extending from the star to the disk, providing observational evidence for the type of structures postulated by magnetospheric accretion models.
While COUP is the best sample available to date for the study of the size of coronal structures on YSOs, its limitations should also be stressed. The main one is the relatively limited photon statistics of *Chandra* observations of YSOs in the ONC: only the most intense and longer lasting flares will have sufficient statistics to allow for the time-resolved spectral analysis needed by the approach used here. Therefore, the physical parameters derived in the present paper are not necessarily representative of the “average” flaring structure in the coronae of YSOs (in the ONC or elsewhere); rather they must be interpreted as the structures associated with the most intense flares present in these coronae. Structures associated with less intense flares (most likely smaller magnetic loops than the ones associated with very intense flares) are certainly present, and given the large large number of shorter and less intense flares visible in the light curve of these stars (which cannot be analyzed in detail due to the low statistics) these are common. Indeed, even in our biased sample a number of compact flaring events is present, and, like what observed in older stars, most likely confined in loops anchored on the stellar photosphere, with no interaction with the disk.
The present paper is structured as follows: the selection of our sample is described in Sect. 2, the analysis procedures adopted in Sect. 3, the results from the analysis are presented in Sect. 4, which includes detailed modeling of a sample flare (on COUP 1343, Sect. 4.1). An example of a complex flaring event (COUP 450, which cannot be analyzed with the method adopted here) is discussed in Sect. 5, and the implications from the present work are discussed in Sect. 6. The conclusions are summarized in Sect. 7, while Appendix Appendix A: Notes on selected individual events discusses in detail a number of notable individual events.
## 2. Sample selection
We have studied all flares in the COUP sample (described in detail elsewhere, Getman et al., 2005) which have sufficient statistics for the analysis performed here. From a set of simulations we initially determined the minimum number of photons in a spectrum necessary to derive reliable spectral parameters given the typical X-ray spectra of COUP stars. A minimum of 750 photons was found to represent a reasonable compromise between the need for statistics and being able to analyze as large as possible a number of events.
To subdivide the light curve in different interval, a Maximum Likelihood (ML) algorithm was used, the same as also used and discussed by Wolk et al. (2005). Briefly, the rationale behind the ML approach is to recursively subdivide the light curve in a number of segments, with the property that the source’s count rate is compatible with being constant during each segment. The minimum number of photons which are comprised in a given ML block is a configurable parameter. While for many applications allowing blocks with as little as 1 photon is useful (e.g. to allow to track fast variability), for our analysis a minimum of 750 photons per ML blocks has been imposed (to allow for meaningful spectral analysis). In all the following whenever the term “ML block” is used, it implies ML blocks with a minimum of 750 photons each.
The ML block approach identifies (as discussed by Wolk et al., 2005) a ’characteristic’ level, essentially a quiescent level during the observation, as well as blocks which are significantly above this level and during which the X-ray luminosity is constant. Note that our constraint of $`750`$ photons per block results, in a number of cases, in significant “under-binning” of the light curve, i.e. the light curve is clearly variable within a given block. This is for example visible in the case of COUP 649 (see Sect. A.5).
Candidate flares for analysis have been initially hand-selected by inspection of the atlas of light curves produced by the ML blocking for the complete COUP sample, and choosing all sources which had at least three contiguous ML blocks above the characteristic values. These are candidate flaring events with sufficient information to grant a more detailed analysis. For example the light curve of COUP 1343 is plotted in Fig. 1, where (as in all other light curves) the ML blocks at the characteristic level are plotted in light blue and the ML blocks with a higher level are plotted in orange. The latter ML blocks form the candidate flaring events. For COUP 1343, in addition to the flare being analyzed (at the beginning of the observation) there are additional ML blocks above the characteristic level, which have been ignored as they are not part of the main flare. This is also common in a number of other sources.
Following the spectral analysis described below, the evolution of each candidate flaring event in the $`\mathrm{log}T`$$`\mathrm{log}\sqrt{EM}`$ plane (a proxy for the $`\mathrm{log}T`$$`\mathrm{log}\sqrt{n_e}`$ plane, as described in Sect. 3) has been examined. Events evolving according to the pattern discussed by Reale et al. (1997) show, after the flare peak, a decay of both the temperature and the emission measure, more or less regular depending on the event characteristics. The event shown in Fig. 1 shows this ’canonical’ behavior: as visible in the the $`\mathrm{log}T`$ vs. $`\mathrm{log}\sqrt{EM}`$ diagram (where the event evolves temporally in the clockwise direction), the temperature peaks while the emission measure is still rising (block 2), and has already started to decay when the emission measure peaks (block 3). Then (blocks 4–9) both the temperature and the emission measure decay regularly. Two smaller flares are also visible in the light curve, one on day 13 and the other on day 21. As not enough statistics is available in each case, they have not been further considered.
We discarded for the purpose of the present analysis all events for which the evolution did not proceed in a sufficiently regular way (e.g. events for which the temperature did not decay). These can either be flares with different underlying mechanisms, which cannot therefore be analyzed within the framework used here (such as e.g. the two-ribbon flares seen in the Sun) or different types of variability (e.g. induced by rotational modulation, or by the emergence of new active regions, etc.). Only events which could be analyzed in an useful way (i.e. producing well constrained physical parameters within the assumed framework) have been retained in the final sample. One example of a particularly complex and intense flare, which cannot be analyzed with the approach used here, is shown and discussed in Sect. 5.
Clearly, as mentioned in Sect. 1, such procedure does not produce an unbiased sample. The most obvious bias is toward events which have sufficient X-ray flux and duration to yield enough time-resolved spectra to allow the analysis. The resulting physical parameters are therefore unlikely to be representative of the ’typical’ flaring structures present in YSOs in the ONC, but will rather represent the more extreme cases (and thus the largest magnetic structures present). Quantifying the biases introduced by our selection procedure is not obvious. An indication can however come from comparison with the study of Wolk et al. (2005), who have performed an analysis of flares in young solar analogs in COUP using an automatic procedure. The young solar analog sample consists of 28 out of 1616 COUP sources. A total of 41 distinct flares were identified on 26 of them, for an average of $`1.5`$ flares per stars during the COUP observation. If the statistics on solar analogs were valid for the complete COUP sample (an assumption which needs to be verified) this would imply some 2400 flares in the observation. The 32 events analyzed here represent only 1% of the COUP population of flares expected extrapolating the statistics for solar-mass stars. We are therefore studying (by necessity, imposed by photon statistics) the tip of an much larger iceberg.
## 3. Analysis method
The analysis of the decay of flares is a classic tool to derive the size of the flaring structure, and thus by inference other quantities such as the plasma density and the confining magnetic field (as e.g. discussed for the Sun by Culhane et al., 1970). Stellar flares have been observed since the beginning of stellar X-ray astronomy, and they have been analyzed in a variety of ways, most of which proceeded by analogy to the characteristics of (spatially resolved) solar flares. In the case of the Sun, flares occur in localized regions in the corona, often in single magnetic loop structures, which can remain mostly unchanged during the flare evolution.
While the heating mechanism causing the flare still remains not understood in its details, it is generally assumed to be unlocked by the release of energy stored in the magnetic field. The impulsive energy release generates either a strong thermal front or particle beams, which, channeled by the magnetic field, hit the photosphere and cause plasma evaporation; the plasma fills the coronal loop(s), and produces the observed increase in X-ray luminosity. The decay of the temperature and the X-ray luminosity is caused jointly by the radiative losses and by heat conduction back to the photosphere.
Both cooling mechanisms depend on the geometry of the flaring structure, and thus the decay time scale of flares is related to the size of the coronal structure (or structures). Stellar flares have been used to infer coronal structure sizes since the early stages of coronal X-ray astronomy (e.g. Pallavicini et al., 1990). However, up and until the late ’90s, it was largely assumed that the heating event was impulsive, and concentrated at the beginning of the flare, coincident with the observed impulsive rise in temperature, so that long decay times always led to infer large coronal structures, in the form of very long loops. In the ’90s studies of solar flares showed that the heating mechanisms often extends well into the decay phase, so that sustained heating and thermodynamic cooling processes are competing, in some cases resulting in long decay times also for flares confined in compact structures. Thus, analyses of stellar flares based on the assumption of impulsive heating would likely result in longer coronal structures than actually present on the star (see for a detailed review of past results Favata & Micela, 2003).
Sylwester et al. (1993) showed that, in solar flares, the slope of the flare decay in the temperature-density plane is a sensitive diagnostic of the presence of sustained heating. Reale et al. (1997) later extended the approach showing that the slope $`\zeta `$ of the flare decay in a $`\mathrm{log}T`$ vs. $`\mathrm{log}n_e`$ diagram gives a quantitative measure of the time-scale of sustained heating, which can be used to correct the observed decay time scale and to derive the intrinsic (thermodynamic) decay time, thus allowing us to obtain the actual size of the flaring structure. Reale et al. (1997) developed a methodology for the analysis of stellar flares, based on grids of hydrodynamic models of flaring loops, which was verified on solar flares (for which the size of individual flaring structures can be obtained from imaging observations), and which has been applied to a significant number of stellar events on a variety of stars, including YSOs. In the vast majority of cases, the analysis of stellar events has shown that sustained heating is present, so that flaring structures in stellar coronae are generally smaller than previously thought.
### 3.1. Uniform cooling loop modeling
The rise time in both solar and stellar flares is almost invariably observed to be significantly shorter than the decay time, and therefore it has been often assumed, in deriving the flare’s physical parameters, that the heating is impulsive, concentrated at the beginning of the event, and that the decay takes place in an “undisturbed” fashion from a loop in near-equilibrium at the flare’s peak, according to the thermodynamic cooling time of the plasma. The two cooling processes which determine the decay time of the flare are thermal conduction downwards to the chromosphere and radiation, each with its characteristic $`1/e`$ decay time,
$$\tau _{\mathrm{cond}}\frac{3nkT}{\kappa T^{7/2}/L^2}\tau _{\mathrm{rad}}\frac{3nkT}{n^2P(T)}$$
(1)
where $`n`$ is the plasma density, $`\kappa `$ is the thermal conductivity and $`P(T)`$ is the plasma emissivity per unit emission measure. Both decay times depend on the loop’s length, $`\tau _{\mathrm{cond}}`$ explicitly through the $`L^2`$ term, $`\tau _{\mathrm{rad}}`$ implicitly through the density’s dependence. The effective cooling time of the loop is the combination of the two,
$$\frac{1}{\tau _{\mathrm{th}}}\frac{1}{\tau _{\mathrm{cond}}}+\frac{1}{\tau _{\mathrm{rad}}}$$
(2)
More generally it has been shown by Serio et al. (1991) that the decay time of a flaring loop starting from equilibrium and decaying freely is linearly related to its half length, through
$$L=\frac{\tau _{\mathrm{th}}\sqrt{T_{\mathrm{pk}}}}{3.7\times 10^4}$$
(3)
where $`L`$ is in units of $`10^9`$ cm, $`\tau _{\mathrm{th}}`$ in sec and $`T_{\mathrm{pk}}`$ is the peak temperature of the plasma in the flaring loop in units of $`10^7`$ K. A fast decay thus implies a short loop, while a slow decay implies a long loop. In all the following, unless otherwise specified, the term “length” will be used to indicate the loop’s half length, from one of the footpoints to the loop apex.
In the 80’s and 90’s a number of approaches to the analysis of stellar flare decays have been used (as discussed in detail by Reale, 2002; see also Favata & Micela, 2003 for a review of previous literature). While differing in their details, they were mostly equivalent to the use of Eq. 3 above, and their key assumption was that the flaring loop always decays, after a short heating episode, in an undisturbed fashion.
If heating does not switch off abruptly but rather decays slowly, the plasma is subject to prolonged heating which extends into the decay phase of the flare. As a result, the actual flare decay time determined from the flare’s light curve, $`\tau _{\mathrm{lc}}`$, will be longer than the intrinsic decay time $`\tau _{\mathrm{th}}`$, and simple application of Eq. 3 will result in an over-estimate of the size of the flaring loop. If sustained heating is present the observed flare decay time must be “corrected” to obtain a reliable estimate of the flaring region’s size.
In the case of solar flares Sylwester et al. (1993) showed that the slope $`\zeta `$ of the flare decay in a $`\mathrm{log}T`$ vs. $`\mathrm{log}n_e`$ diagram provides a diagnostic of the presence of sustained heating, with a shallower slope (slower temperature decay) indicating an event with strongly sustained heating. This approach was extended to the analysis of stellar flares by Reale et al. (1997), who showed that $`\zeta `$ provides a quantitative diagnostic of the ratio between the intrinsic and observed decay times, i.e.
$$\frac{\tau _{\mathrm{lc}}}{\tau _{\mathrm{th}}}=F(\zeta )$$
(4)
The approach of Reale et al. (1997) was based on numerical simulations of flaring loops, under the assumption of an exponentially decaying heating function (with a time scale $`\tau _\mathrm{H}`$), and it was validated by comparing its results with the measured size of a number of (spatially resolved) solar flares.
The actual functional form and parameters of $`F(\zeta )`$ depends on the bandpass and spectral response of the instrument used to observe the X-ray emission from the flare, and therefore needs to be separately determined for each instrument. Note that for stellar observations no density determination is normally available; in this case, the quantity $`\sqrt{EM}`$ is used as a proxy to the density, under the assumption that the geometry of the flaring loop does not vary during the decay.
For ACIS observations the relationship between $`\zeta `$ and the ratio between the observed and intrinsic flare decay time has been calibrated for the present analysis, resulting in
$$\frac{\tau _{\mathrm{lc}}}{\tau _{\mathrm{th}}}=F(\zeta )=\frac{0.63}{\zeta 0.32}+1.41$$
(5)
which is valid for $`0.32<\zeta 1.5`$. The limits of applicability correspond on one side ($`\zeta 1.5`$) to a freely decaying loop, with no heating ($`\tau _\mathrm{H}=0`$), to the other ($`\zeta =0.32`$) to a sequence of quasi-static states for the loop ($`\tau _\mathrm{H}\tau _{\mathrm{th}}`$), in which the heating time scale is so long as to mask the loop’s intrinsic decay. Note that while it is not possible to have loops with $`\zeta >1.5`$, flares with $`\zeta <0.32`$ are possible if the heating is not simply exponentially decaying but if it increases again after an initial decay.
The temperature in Eq. 3 is the peak temperature in the flaring loop, which will contain plasma with a distribution of temperatures. The relationship between the maximum temperature present in the loop and the temperature determined by performing a simple 1-temperature fit to the integrated X-ray spectrum emitted by the loop has also been calibrated here. In the case of observations performed with ACIS is
$$T_{\mathrm{pk}}=0.068\times T_{\mathrm{obs}}^{1.20}$$
(6)
where both temperatures are in K.
Putting together the above elements, the length of the flaring structure is therefore given by
$$L=\frac{\tau _{\mathrm{lc}}\sqrt{T_{\mathrm{pk}}}}{3.7\times 10^4F(\zeta )}$$
(7)
Once the length of the flaring loop has been derived, it is possible to infer additional physical parameters of the flaring regions, with some additional assumptions. Given (from the analysis of the peak spectrum) a peak emission measure for the flare, a density can be derived if a volume is available. The above analysis only gives the loop length; if an assumption is made about the ratio between the loop’s radius and its length ($`\beta =r/L`$) one can derive a volume,
$$V=2\pi \beta ^2L^3$$
(8)
and the resulting plasma density at flare peak
$$n\sqrt{\frac{EM}{V}}\sqrt{\frac{EM}{2\pi \beta ^2L^3}}$$
(9)
The minimum magnetic field necessary to confine the flaring plasma can be simply estimated as
$$B\sqrt{8\pi p}=\sqrt{8\pi knT}$$
(10)
where $`p`$ is the plasma pressure derived from the density and the temperature.
For the solar corona (the only one accessible to imaging observations) typically $`\beta 0.1`$, and this value has been used in most previous analyses of stellar flares. We have also assumed, in all of our estimates, $`\beta =0.1`$, unless otherwise specified. The consequences of differences in $`\beta `$, depending on the loops characteristics, are discussed in Sects. 3.2 and 4.
### 3.2. Uncertainties and error analysis
Both statistical and systematic uncertainties are of course present in the length estimate for the flaring loop obtained through Eq. 7. The main statistical uncertainty comes from the uncertainty in the determination of $`\zeta `$ and therefore of $`F(\zeta )`$. Uncertainties in the determination of the peak temperature and of the decay time are usually significantly smaller than the uncertainty deriving from the $`F(\zeta )`$, which, given the hyperbolic form of the function, becomes potentially very significant at the lower end of the range of validity of $`\zeta `$. In our analysis, the error bar reported for the length estimate is the one obtained by propagation of the statistical errors on $`\zeta `$, $`T_{\mathrm{pk}}`$ and $`\tau _{\mathrm{lc}}`$, and it is dominated by the error on $`\zeta `$.
One possible additional source of uncertainty would be the lack of accurate knowledge about the star’s gravity. Gravity may in principle bring systematic effects to the diagnostics, because it makes the plasma drain faster during the decay phase, flattening the flare’s decay in the $`\mathrm{log}T`$$`\mathrm{log}\sqrt{EM}`$ plane, mimicking the presence of sustained heating (Reale & Micela, 1998). In a number of cases, no estimates of either the radius or the mass of the parent star are known, so that no determination of the surface gravity is possible. However, given the low values of surface gravity in YSOs (with their large radii, typically $`g0.1g_{}`$), the reduction of the gravity at distances comparable to the stellar radius from the surface, and the very high temperatures observed at flare peaks, the pressure scale height of the flaring plasma is likely to be always larger than the length of the flaring structures, so that gravity is unlikely to be an additional source of uncertainty.
More subtle is the issue of systematics, i.e. of uncertainties deriving from assumptions in the models which are not reflected by reality. The modeling of the flaring plasma structure is based on relatively simple physics, robust in its assumptions. The modeling is simplified by the one-dimensional nature of the problem, with the magnetic field confining the plasma except along the field lines, so that the precise shape of the flaring structure is not important. Another assumption is that the sustained heating is an exponentially decaying function; of course this is a simple parameterization which has no underlying physical reasoning behind it. However, as long as the temperature indeed decreases, and the slope in the $`\mathrm{log}T`$$`\mathrm{log}\sqrt{EM}`$ plane is steeper than the one of the quasi-static sequence locus ($`\zeta >0.32`$) the heating must be a decreasing function of time. Its detailed functional form in this case has a limited impact on the loop length, so that the resulting values are robust against the details of the heating function. Also, as shown by Reale et al. (2002) strong, long-lasting flares must be strongly confined by the magnetic and would decay much faster than observed if the plasma were to break confinement.
The derived parameters (density, magnetic fields) are subject to an additional uncertainty relative to the loop’s aspect ratio $`\beta `$. Our simple assumption about $`\beta `$ clearly cannot hold for a loop with constant cylindrical cross section when this has a length of several stellar radii (as is the case in some of the events analyzed here). In this case the footpoints would cover the whole stars or be larger than it, a clearly unrealistic configuration. Thus, such loops must, at the footpoints, have smaller aspect ratios. At the same time, though, a loop of such length is likely to follow a dipole-like magnetic field configuration at large distances from the star (with the higher-degree multipole terms becoming less important with distance). This implies that it will likely have an expanded cross-section at the apex (where most of the flaring plasma will be) than at the footpoints. These two effects compensate each other (although to an unknown degree), and we have thus decided to still maintain $`\beta =0.1`$ in the analysis, keeping in mind that for the longer loops the “effective $`\beta `$” may be somewhat smaller.
### 3.3. Application of the analysis to the COUP data
To determine the evolution of COUP flares in the $`\mathrm{log}T`$ vs. $`\mathrm{log}\sqrt{EM}`$ diagram, the temperatures and emission measure of the flaring component have been separately determined for each ML block. All spectral fits were performed within the XSPEC package, using the mekal spectral emissivity model for coronal equilibrium plasma, for consistency with the rest of the COUP analysis as described in Getman et al. (2005). An example of the analysis is shown in Fig. 1: the top panel shows the raw ACIS light curve (at a resolution of 2 ks), with superimposed the light curve rebinned in ML blocks. The segments in light blue are the ones representative of the “characteristic” level for the source’s emission, while the ones in orange are the ones significantly above the characteristic level, and which form the basis for the flare analysis. The count rate for the light curve is given in cts/ks on the left vertical axis, and in (approximate) counts per 2 ks bin on the right vertical axis. The time scale is given in ks from the observation start on the bottom axis and in calendar days (January 2003) in the top axis. Also indicated are the radius and mass of the star (if known, in solar units) and the total number of photons in the light curve.
In a number of cases the flaring source showed significant pile-up in the ACIS detector during the flare. In this case the pile-up free light-curve was examined (obtained by extracting the photons from a ring around the peak of the PSF, as described in Sect. 6 of Getman et al., 2005). However, in most cases the statistics of the pile-up free light curve are so much lower that insufficient photons were left for a detailed analysis of the flare. Only in the case of COUP 891 (Fig. 13) could the flare be analyzed, in spite of the presence of pile-up. The flares in COUP 107, 245,290, 394, 430 and 881, which would have had sufficient statistics and sampling, were lost from our sample because of pile-up.
The source’s “characteristic spectrum” has been determined by performing an absorbed multi-temperature fit to the integrated spectrum of all ML blocks compatible with the characteristic level of the source. The purpose of this fit was only to provide a noise-free, phenomenological description of the source’s characteristic emission, which is effectively a form of background to be subtracted from the flaring emission. Our fit provides a model of the background to be subtracted, which is more robust and less sensitive to statistical fluctuations than the subtraction of the characteristics spectrum itself, which in a number of cases has rather low statistics. Therefore enough temperature components have been included until the fit has yielded a satisfactory $`\chi ^2`$, and individual metal abundances have been left free to vary, independent of whether their best-fit value was physically meaningful or not.
For each ML block with emission above the characteristic level, a spectral fit to the emission in excess of the characteristic level has been performed. In addition to the characteristic component model with frozen fit parameters, we include a single additional absorbed thermal component, with global metal abundance set to $`Z=0.3Z_{}`$, thus determining the emission measure and temperature of the flaring component. The absorbing column density $`N(\mathrm{H})`$ was left free to vary. Previous works on CCD resolution X-ray spectroscopy of very active stars, including PMS stars, shows that they typically have low coronal abundance, around $`Z=0.3Z_{}`$, which we have therefore chosen here. While coronal abundance variations have been observed during intense stellar flares, at the temperatures of interest here the plasma is fully dominated by continuum emission, and therefore the exact value of the metallicity – which often cannot be well constrained – has little influence on the best-fit temperature and emission measures. In all cases this simple one-temperature fit resulted in a satisfactory description of the spectrum, although a number of events have spectral peculiarities (regarding variations of abundance and absorbing column density) which will be discussed in detail in a future paper. These do not affect any of our conclusions.
The bottom four-quadrant panel in Fig. 1 shows the results of our analysis in the case of COUP 1343. The bottom right quadrant shows the ML blocked light curve with the time axis expanded for clarity (only the flaring event is shown, and none of the characteristic level blocks), while the upper right and lower left quadrants show, respectively, the time evolution of the emission measure and of the temperature of the flaring component, as determined from the analysis of each of the ML flaring blocks, in a semi-log scale. The best fit exponential to the decay is also plotted for both quantities as a continuous line. The top left quadrant shows the evolution of the flare in the the $`\mathrm{log}T`$$`\mathrm{log}\sqrt{EM}`$ plane, with the points connected by a dashed line to allow to follow the event’s evolution. Points are numbered accordingly to the sequential number of the ML block. The blocks which have been used to determine the flare’s evolution are marked by a dot, and the best-fitting decay is plotted as continuous line. Its slope $`\zeta `$ is the one used in Eq. 5. For all such plots we have chosen to adopt the same scale in both axes, to allow an immediate comparison of different events.
In Fig. 2 we plot four representative ACIS spectra for COUP 1343, illustrating the type and quality of the spectra used in the present paper. The top left panel of Fig. 2 shows the characteristic spectrum of COUP 1343, with the 3-temperature fit used to describe it. This spectrum has been integrated for $`540`$ ks, and has 3821 source file counts. The other panels show (clockwise) the spectrum at the beginning of the flare’s rise (block 2, when the plasma is hottest), at the beginning of the flare’s decay (block 4) and at the end of the decay (block 9, when the plasma has already cooled to a temperature very similar to the quiescent value). The flare spectra contain 738 photons over 4.7 ks, 1058 photons over 5.6 ks and 1214 photons over 36 ks, respectively.
For a large number of sources (14 out of 32), the peak temperature derived by the above procedure converges to a value higher than 10 keV, although invariably with large error bars. Detailed analysis of the fits shows that the fit procedure cannot determine with good accuracy the temperature of thermal components hotter than $`100`$ MK on ACIS spectra; when hotter plasma is present, statistically equivalent fits can be obtained by arbitrarily fixing the temperature for this hot component to any value between 100 and 600 MK. We have thus decided to adopt a conservative approach, and for all cases for which the maximum best-fit temperature $`T_{\mathrm{obs}}>100`$ MK the value 100 MK has been used in Eq. 6 (yielding $`T_{\mathrm{pk}}=270`$ MK). This will result in smaller loop sizes derived from Eq. 7 than it would be the case if the actual maximum peak temperature were used. However the dependence of $`L`$ on $`T_{\mathrm{pk}}`$ is relatively weak (going as $`\sqrt{T_{\mathrm{pk}}}`$); thus, even if the peak temperature were to be underestimated, in a few cases perhaps by up to a factor of two, this would result in an additional uncertainty in the size of the flaring loop of $`40\%`$.
## 4. Results
The physical parameters derived for each flare are listed in Table 1. For each source we report the maximum observed temperature $`T_{\mathrm{obs}}`$, the peak X-ray luminosity $`L_\mathrm{X}`$, the peak temperature at the loop apex $`T_{\mathrm{pk}}`$ determined from Eq. 6, the slope $`\zeta `$ of the flare decay in the $`\mathrm{log}T`$ vs. $`\mathrm{log}\sqrt{EM}`$ diagram, the ratio of the observed to the intrinsic decay times $`F(\zeta )`$ (determined from Eq. 5), the observed decay time scale (in ks) $`\tau _{\mathrm{lc}}`$ determined by fitting the flare decay binned in ML blocks. The length of the flaring loop $`L`$ (from Eq. 7) is given in units of $`10^{10}`$ cm, as well as, for stars which have a radius estimate published, in units of the stellar radius. The density $`n_e`$ computed from the peak emission measure and from the volume of the flaring loop (Eq. 9) is given in units of $`10^{10}`$ cm<sup>-3</sup>, and the minimum magnetic field necessary to confine the flaring plasma at the loop apex $`B`$ (Eq. 10) is given in Gauss. Also given (when known) are the equivalent width of the Ca ii IR triplet, a diagnostic of active accretion processes, and whether the source has $`\mathrm{\Delta }(KL)`$ and $`\mathrm{\Delta }(IK)`$ excesses. Note that the peak X-ray luminosity is computed after subtraction of the characteristic emission from the flaring one.
The derived quantities ($`n_e,B`$) are computed assuming $`\beta =0.1`$. For long loops (of several stellar radii) this assumption may be incorrect, as such large loops would have footpoints covering much of the stellar surface (even though loop expansion may partially offset this, see Sect. 3.2). However, the dependence of $`n_e`$ and $`B`$ on $`\beta `$ is moderate, with $`n_e\beta ^1`$ and $`B\beta ^{0.5}`$. Thus, the values in Table 1 can be easily scaled to any value of $`\beta `$.
The flares analyzed here are all X-ray bright, but not exceptionally so, when compared to other intense flares already observed in YSOs. The peak X-ray luminosities reported in Table 1, determined by fitting the spectrum of the peak ML block for the flare, range from $`7\times 10^{30}`$ erg s<sup>-1</sup> up to $`8\times 10^{32}`$ erg s<sup>-1</sup> for the most X-ray bright event. The latter one takes place on COUP 1568 (discussed in detail in Appendix Appendix A: Notes on selected individual events), at $`M=2.6M_{}`$ one of the most massive stars in our sample. Thanks to the large offset angle, the source does not suffer from pile-up even if the peak ACIS count rate approaches 1 cts/s. While certainly very high, the peak luminosity of $`8\times 10^{32}`$ erg s<sup>-1</sup> is comparable to the values observed in other PMS flares, e.g. the $`10^{33}`$ erg s<sup>-1</sup> determined at peak in a similar band by Tsuboi et al. (1998) for the flare observed with the ASCA satellite on the weak line T Tau V773 Tau.
More remarkable are the high temperatures present in COUP flares. Up to recently, flaring events with temperatures reaching up to 100 MK where considered exceptional. In the COUP sample, on the other hand, very high temperatures are common. While the ACIS spectral response does not allow to properly constrain temperatures in excess of 100 MK, the best fit values for the hottest ML block in the flare exceed 100 MK in about half of the events studied here. The fact that the temperature decay shows a clear regular pattern even when the best-fit values are above 100 MK (as for example is the case for COUP 1568), gives confidence in the fact that the peak observed temperatures are, in a number of cases, very high, hundreds of MK, even though their values cannot be accurately determined using ACIS spectra (see Sect. 3.3).
For the ONC YSO sources studied here, the nominal sizes of the flaring structures vary from a fraction of the stellar radius (e.g. COUP 7, 90, 141), with absolute sizes of order $`10^{11}`$ cm, to very large structures, up to a few times $`10^{12}`$ cm, i.e., 10–20 stellar radii. As discussed below, the largest events with a reliable size estimate (with small uncertainties) typically have $`L5R_{}`$.
A number of intense flares on active coronal sources have been analyzed to date with the approach used here. These include a number of YSOs, both accreting and non-accreting (YLW 15, HD 283572, LkH$`\alpha `$,92, V773 Tau, Favata et al., 2001), the nearby ZAMS star AB Dor (Maggio et al., 2000), a number of dMe flare stars (Favata et al., 2000b; Favata et al., 2000a, Reale et al., 2004), and a number of active binary systems as well as Algol (Favata & Schmitt, 1999). In all cases, without exception, the size of the flaring structure derived from the analysis is at most comparable to the stellar radius, and in most cases it’s smaller. On the other hand the present analysis of intense flares on YSOs in the ONC results, in a number of cases, in very large loop sizes, extending to several stellar radii, showing them to be events of a different nature from the flares analyzed to date.
The strongest evidence for large loop sizes in the COUP sample comes from well resolved long-lasting events with little sustained heating. In these cases, one is largely observing the undisturbed thermodynamic decay of the flaring structure, which, as discussed in Sect. 3, is directly related to the size of the flaring structure. One such example is COUP 1343, for which our analysis gives $`\zeta =1.95\pm 0.51`$, and $`F(\zeta )=1.8`$, resulting in a loop size of $`1\times 10^{12}`$ cm. The presence of a small amount of residual heating is derived from the tail of the light curve, which decays more slowly, and indeed the detailed simulation of Sect. 4.1 shows that for the first 50 or so ks the decay of both the emission measure and the temperature is well reproduced by a freely decaying loop of $`L=1\times 10^{12}`$ cm. In fact, if we were to relax our conservative assumption on the peak temperature of the loop (with the best-fit temperature clipped to 100 MK) and use the nominal best-fit temperature, the resulting loop size would be even larger.
Other similar cases of well determined flaring structure sizes include COUP 669, 891, 971 and 1246. In all these cases $`F(\zeta )2.5`$, with a well constrained value of $`\zeta `$, leading to a reliable determination of $`L`$ from Eq. 7. Both COUP 971 and 1246 have $`L4\times 10^{11}`$ cm, while COUP 669 and 1343 have $`L1\times 10^{12}`$ cm. With $`L=1.7\times 10^{12}`$ cm, COUP 891 is the largest of the well constrained loop structures. In all these cases, the uncertainties for $`L`$ resulting from the uncertainty in the determination of $`F(\zeta )`$ are of order 20–30%. Apart from COUP 971 (which, at $`L1.6R_{}`$, is likely to be a large “normal” coronal flare, with both loop footpoints anchored on the stellar photosphere), the well determined loops have $`L4`$$`5R_{}`$, which, as discussed in Sect. 6, is the typical corotation radius for a low-mass YSO (also likely to be the disk truncation radius), supporting the hypothesis that these loops may indeed be structures linking the star and the accretion disk.
For the large events with a well constrained loop size, the peak densities are all in the range $`1`$$`8\times 10^{10}`$ cm<sup>-3</sup>, and the equipartition magnetic fields (determined from Eq. 10, at the top of the loop) range from 40 to 250 G. In the assumption of a simple dipole geometry (very likely to dominate at large distance from the stellar surface) the photospheric field would be
$$B_{\mathrm{ph}}(LR_{})^3\times B_{\mathrm{eq}}$$
(11)
Assuming an average equipartition field of 100 G and a typical loop length of $`5R_{}`$, the corresponding field strength at the stellar surface would therefore be of order 6 kG, at the top of the range of the field strengths determined in YSOs by Zeeman splitting (Johns-Krull & Valenti, 2005).
In a few flares of Table 1, we obtain $`\zeta 1.5`$ (e.g. COUP 332, 454, 848, 976), outside the range of validity of Eq. 5. Probably, in these events the analyzed part of the decay is too short to cover a significant part of the decay path, and the slope is not yet well defined. Cases which result in long flaring structures in the presence of strong sustained heating are to be regarded with caution, as the strong heating dominates the decay, hiding the intrinsic decay of the cooling plasma and potentially introducing a larger uncertainty in the results. Also, the hyperbolic form of Eq. 5 will tend to include lengths much smaller than the ’nominal’ results of Eq. 7 within the final error bar for $`L`$. One such example is COUP 1410, for which $`\zeta =0.45\pm 0.12`$ and $`F(\zeta )=6.1`$. While the nominal loop length is $`L=1.1\times 10^{12}`$ cm, the uncertainty range resulting from the uncertainty in $`\zeta `$ is $`L=10^{11}`$$`2\times 10^{12}`$ cm, with a factor of ten error bar at the lower end.
A number of individual flaring events of particular interest are discussed in detail in Appendix Appendix A: Notes on selected individual events, where the individual light curves and the time evolution of the spectral parameters are also shown.
### 4.1. Detailed simulation of the flare on COUP 1343
Although the method used here has been widely applied in the literature, and its reliability tested for example in the case of the eclipsed flare observed by SAX on Algol (Schmitt & Favata, 1999; Favata & Schmitt, 1999), as discussed in the above Section no previous analysis based on this approach has resulted in such large loop sizes (of order $`0.1`$ AU). It is thus legitimate to ask: 1) whether such gigantic loops can be still described with standard loop models, closed magnetic structures where plasma is confined and moves and transports energy along the magnetic field, and 2) whether the diagnostics we are using (in particular the slope $`\zeta `$ in the $`\mathrm{log}T`$$`\mathrm{log}\sqrt{EM}`$ plane) are still valid. To check on the applicability of the method to such extreme regimes, and to validate its use on the present sample, we have modeled one of the well-resolved large flares from the COUP sample in detail. We have used a time-dependent hydrodynamic model of plasma confined inside a coronal loop, as already done for other more regular stellar flares (Reale et al., 1988; Reale et al., 2004).
We have chosen to model the large event on COUP 1343, in which the peak count rate is approximately $`30\times `$ larger than the characteristic rate and which appears, from the analysis performed using the approach described in Sect. 3.3, to have evidence for little sustained heating. The event lasts some 130 ks, and both the rise phase and most of the decay are not interrupted by gaps in the observations. Additionally, the decay in the $`\mathrm{log}T`$$`\mathrm{log}\sqrt{EM}`$ plane is well defined. The event and its analysis with the analysis method used here (as described in Sect. 3.3) is shown in Fig. 1.
To model the event on COUP 1343 in detail, we have assumed the plasma to be confined in a loop with constant cross-section and half-length $`L=10^{12}`$ cm, symmetric around the loop apex. The flare simulation is triggered by injecting a heating pulse in the loop which is initially at a temperature of $`20`$ MK (which is the dominant temperature in the characteristic spectrum as well as the temperature observed at the end of flare decay). We have chosen to deposit two heat pulses with a Gaussian spatial distribution of intensity 10 erg cm<sup>-3</sup> s<sup>-1</sup> and width $`10^{10}`$ cm (1/100 of the loop half-length) at a distance of $`2\times 10^{10}`$ cm from the footpoints, i.e. very close to the footpoints themselves (Reale et al., 2004). After 20 ks the heat pulses are switched off completely, with no residual heating in the decay, as indicated by the density-temperature diagnostics for this specific flare. We have computed the evolution of the loop plasma by solving the time-dependent hydrodynamic equations of mass, momentum and energy conservation for a compressible plasma confined in the loop (Peres et al., 1982; Betta et al., 1997), including the relevant physical effects such as the plasma thermal conduction and radiative losses (computed assuming solar metal abundances – at these high temperatures however the line losses are a minor contributor, so that the results are robust against differences in the plasma metal abundance). The gravity component along the loop is computed assuming a radius $`R_{}=3R_{}`$ and a surface gravity $`g_{}=0.1g_{}`$, typical of low-mass pre-main sequence objects. The results of the simulations are shown in Fig. 3, together with the observational results for COUP 1343.
The computed evolution of the flare’s parameters largely resembles the evolution computed for other stellar flares (e.g. Reale et al., 2004), although on larger scales. The heat pulses make the loop plasma heat rapidly ($`1`$ hour) to temperatures above 200 MK and the initially denser chromospheric plasma near the loop’s footpoints expands dynamically upwards at speeds above 2000 km/s, to reach the loop apex on similar timescales (also about 1 hour). After this first impulsive phase, the temperature does not change much and the evaporation continues substantially but less dynamically. After the end of the heat pulse, the plasma begins to cool while still filling the loop. The density begins to decrease only two hours later.
From the density and temperature distribution of the plasma along the model loop, we can synthesize the expected plasma X-ray spectrum filtered through the ACIS spectral response, and the relative spectral parameters as they would be determined through a spectral fit to the ACIS data. We have computed the emission in the coronal part of the loop, i.e. above the loop transition region, assuming a hydrogen column density $`N(\mathrm{H})=10^{22}`$ cm<sup>-2</sup> (i.e. the value obtained from the fit to the source’s characteristic level spectrum). The dashed line in Fig. 3 (top left panel) shows the ACIS light curve resulting from the simulation, integrated above 1 keV, assuming a loop cross-section radius of $`2\times 10^{10}`$ cm, i.e. $`\beta =0.02`$. The smaller $`\beta `$ than customary ($`\beta =0.1`$ being assumed elsewhere) is necessary to provide a good fit to the flare’s decay; the small $`\beta `$ also results in footpoints still small relative to the photosphere. The spectra obtained from the loop modeling have been fit with 1-$`T`$ model spectra, and the resulting evolution of the emission measure and of the temperature, as well as the $`\mathrm{log}T`$$`\mathrm{log}\sqrt{EM}`$ diagram are also shown in the figure. A visual comparison of the model results with the actual light curve and with the evolution of the measured spectral parameters indicates that the loop simulation results are in good qualitative agreement with the data and that, therefore, the data are compatible with the flare occurring in a single giant flaring structure. The modeling could be further improved by changing some of the simulation’s parameters; for example the light curve’s agreement for $`t50`$ ks could be improved by including a low residual heating active during the late decay (as done e.g. by Reale et al., 2004 in their simulation of a large flare on Proxima Cen). However, such changes would not alter the large-scale features of the flaring loop, in particular its large size, and thus for the purpose of the present work they are not necessary.
One additional diagnostic provided by detailed simulations such as the one discussed here is the rise time of the event. This is not available with the more general analysis employed for the complete sample of flares discussed in the present paper, which only uses the decay phase. In the case of COUP 1343 the rise phase is very steep, with a very peaked light curve. Not all impulsive heating functions will produce such steep rise; we find that to produce the steep rise it is necessary to have strongly localized heating deposited very close to the loop footpoints. More diffuse heating, or heating localized near the loop apex would produce a smoother light curve.
## 5. Complex flaring events
While the present paper deals with all the flaring events in the COUP sample which can be modeled within the framework described by Reale et al. (1997), a number of events are present in the sample which do not show such regular behavior, and which cannot therefore be analyzed with the same approach. Perhaps the most dramatic one is the flare on COUP 450, which, coincidentally, was also observed at mm wavelengths. COUP 450 is associated with a well known radio source (GMR-A, Garay et al., 1987) and is identified with star HC 573 from Hillenbrand & Carpenter (2000).
The X-ray event has been briefly reported by Getman et al. (2003), and the COUP light curve for the event is shown in Fig. 4. As discussed in detail by Bower et al. (2003), the source also flared in the mm regime, although the mm flare started significantly later (about two days after the start of the X-ray flare). The mm flare lasted some 20 days, with erratic variability continuing for tens of days thereafter. The COUP observation only covers the first day or so of the mm event (which has a much spottier time coverage), and no obvious correlation between the details of the two light curves is present. The counterpart to COUP 450 has been studied by Bower et al. (2003) who (using IR Keck spectroscopy) identify it as a heavily absorbed K5 T Tauri star with average photospheric magnetic fields $`B23`$ kG determined from Zeeman-splitted infrared lines.
The COUP 450 event is not exceptional either in terms of its peak luminosity (with a rise of $`20\times `$ in X-ray count rate at peak) or in terms of its duration (other events with durations of a few days are present in COUP). It is however one of the most complex events in terms of its light curve shape, and it lies at the extreme of complex flares, which lack a clear, well defined decay phase. Such events cannot therefore be analyzed with the method used here. In this specific case, the analysis is also hampered by the very strong degree of pile-up present in the *Chandra* observation, so that while the raw count rate ensures a ’high-resolution’ view of the light curve, most events are lost to pile-up, and the residual pile-up free events suffer from very limited statistics. Nevertheless it is possible to still resolve the light curve in a number of intervals, allowing us to monitor the evolution of the spectral parameters.
The time evolution of the flare can be seen in the bottom panels of Fig. 4. The very significant amount of pile-up in the ACIS observation results in a very small pile-up free count rate. In the light curve shown in the top panel of Fig. 4 the purple line shows the ML blocks of the pile-up free events: although the raw count rate approaches 1 cts/sec, the pile-up free ML blocks have durations of tens of ks, which clearly subsample the event’s evolution. Even so, however, the irregular evolution of the event is readily visible. No regular decay is present; the emission measure does not decrease, and the temperature shows an irregular behavior, with a strong increase in the next to last bin. Whether this is evidence of strong reheating within the same flaring structure or whether this is a separate structure undergoing flaring (perhaps induced by the first event) cannot be inferred from our data, specially because of the coarse time sampling caused by pile-up.
The lack of a regular evolution for this event does not allow us to make (using the approach adopted in the present paper) any inference about the size and characteristics of the flaring region, and thus about the flaring mechanism and the underlying physics. The clear lack of a decay does point in this case to a complex heating function, with a highly time-variable release of energy in the flaring plasma, and – perhaps – to the involvement of different structures in the flaring, rather than to the single, confined plasma structure inferred from the analysis of the simpler events discussed in Sects. 3 and 4.
## 6. Discussion
In a significant number of cases the size of the flaring structures derived here for intense flares in COUP sources is large (several stellar radii) with absolute sizes up to $`10^{12}`$ cm. In a few of them, the evidence for large loop sizes is strong, because we could constrain negligible heating to be present during the decay. Significant heating would have made the diagnosed lengths shorter and at the same time more uncertain. Given the high plasma temperatures at the peak of the flare ($`T>100`$ MK) the magnetic field necessary to confine the plasma typically is of a few hundred G. This scenario is very different from the one observed in more evolved stars, ZAMS or MS, in which the same type of analysis of flare decay as performed here has invariably yielded relatively compact magnetic structures $`LR_{}`$, compatible with scaled-up version of the geometry observed on the Sun. Also the limited number of flares previously studied with the same technique on YSOs (Favata et al., 2001) have shown evidence for compact structures. This is not in contradiction with the evidence from the COUP sample, as the YSOs studied by Favata et al. (2001) all are in the nearest SFRs (thus allowing studies of much less intense flares) and were the result of shorter observations (and therefore biased against the long lasting events well represented in COUP).
Such long loops are unlikely to be anchored on the star alone, as the centrifugal force would rip them open and eject the plasma. The corotation radius is defined as the distance at which a test gas particle in keplerian orbit around the central star has the same angular velocity as the star,
$$R_\mathrm{c}=(GM_{}/\mathrm{\Omega }_{}^2)^{1/3}$$
(12)
The effective gravitational potential (gravity plus centrifugal) takes the form of a vanishing gradient at the corotation radius, which is a saddle point, allowing only for magnetically dominated plasma to be present above and below the equatorial plane. The plasma favors moving inward while inside of the corotation radius, but flowing outward while outside of the corotation radius. Unless the magnetic field is strong enough to provide sufficient tension therefore enforcing corotation, the field itself would easily be opened up by plasma flowing downhill in the effective gravitational potential.
Following Mestel & Spruit (1987) – hereafter MS87 – and assuming the simplest stellar dipole field, a scaled-up version of the Sun for the stellar activity (with the magnetic field strength proportional to the rotational velocity of the star, i.e. the $`p=1`$ case of MS87, so that $`L_\mathrm{X}B^2`$), Table 1 in MS87 gives a list of sizes of the largest closed loops that can possibly exist for a rotating solar-mass star under similar condition. By simply re-scaling with our stellar parameters, for loops of temperatures on the order of 100 MK or higher, the largest equatorial extent for closed magnetic loops (in the notation of MS87, $`C`$-cusp point, see their Fig. 1), will likely be in the range of 3–4 $`R_{}`$, for a slow rotator (6–8 days, or $`\alpha /\alpha _{}4`$$`5`$ in the notation of MS87, where $`\alpha `$ is the rotation rate of the star), and high plasma temperature ($`\zeta _d\alpha `$ and $`(\zeta _d)_{}=4`$, where $`\zeta _d`$ is the ratio of magnetic to thermal energy in the corona). Similar considerations have recently been applied by Jardine (2004) in the context of ’coronal stripping’ in fast rotating, supersaturated stars.
Our flaring sources have typical rotational periods (when they are known) of a few days, up to a week. For typical CTTSs periods of 6–8 days, the corresponding corotation radius is about $`9\times 10^{11}`$ to $`1.2\times 10^{12}`$ cm, equivalent to $`5R_{}`$, if $`R_{}=3R_{}`$. Loop structures larger than this would have difficulty anchoring on the star alone without very strong magnetic field tension to counter balance the centrifugal force. What is then the magnetic geometry which can support coherent fields of hundreds of G over $`10^{12}`$ cm? In the case of YSOs one obvious possibility is to have magnetic flux tubes connecting the stellar photosphere with the disk, at the co-rotation radius or somewhat outside it. As e.g. discussed by Shu et al. (1997), the twisting of the magnetic field lines induced by the differential (Keplerian) rotation of the inner disk rim and the photosphere will twist the flux tubes, resulting in longer loops and presumably in the stressed field configurations which drive the flaring. A good candidate among COUP flaring sources for star-disk flaring structures is COUP 1083, a moderate mass star with $`P_{\mathrm{rot}}=5.9`$ d and with a flare implying $`L=2.4\times 10^{12}\mathrm{cm}=34R_{}`$, in the presence of moderate sustained heating, with an uncertainty range $`1.9`$$`3.0\times 10^{12}\mathrm{cm}`$. Even assuming a very young YSO, with a large radius $`R=3R_{}`$,the loop length is $`L10R_{}`$. Given the above consideration on stability of large magnetic loops in rotating stars, it is unlikely that this type of loops can be anchored on the star alone.
Magnetic flux tubes extending from the inner rim of the disk to the stellar photosphere provide a natural location for the intense flares found here on a number of sources. An additional open question is what heating mechanism can heat the plasma in these long loop to such high temperatures. While the general coronal heating mechanism is still not fully established, the most modern magnetohydrodynamic simulations (e.g. Peter et al., 2004 and Schrijver et al., 2004) point to the heating being due to shuffling of the loop’s footpoint, a natural result of convective motions in the photosphere. The photospheric convection is certainly present in accreting YSOs, and shearing of the disk would provide an equivalent, natural source of loop footpoint shuffling on the disk side.
It should be stressed that the COUP flaring sources analyzed in the present paper present a large variety of flaring structures. While the long flaring loops which likely connect the star with the accretion disk are well represented in the sample, so are also the more compact structures which have very similar characteristics with the ones observed on more evolved stars, ZAMS or MS. Thus, it appears from the present analysis that two types of (flaring) coronal structures can coexist on YSOs: one which appears similar to the one present in older, more evolved stars (and thus likely to be a scaled-up version of the solar corona) and one made of very extended magnetic structures, which is instead peculiar to YSOs, being dependent on the presence of disks on which to anchor one footpoint of the very long loops. While unfortunately none of the COUP sources has more than one flare sufficiently bright to be analyzed with the present technique, some sources (notably COUP 597 and COUP 891, both discussed in Appendix Appendix A: Notes on selected individual events) show, in addition to the long events linked with long loops, some short events, too short to be analyzed, but likely originating in compact loops (given the small $`\tau _{\mathrm{lc}}`$), so that it would appear that both types of coronae could exist, at the same time, in a given YSO.
Wolk et al. (2005) have presented an analysis of flare characteristics and frequency in “young suns” in the COUP database. Only two of our objects are in common with the sample of Wolk et al. (2005), namely COUP 223 and 262. While no statistical conclusions can be drawn from only two objects, the flare on COUP 262 results in a long coronal structure ($`L3.6R_{}`$, see discussion in Sect. A.3), showing that these extended magnetic structures were indeed present in the young Sun.
Very recently, Loinard et al. (2005) have obtained evidence for the presence of long magnetic structures, in T Tau S: their radio VLBI observations show that, in addition to a compact source which they identify as the star’s magnetosphere, a fainter streak of radio emission is present, extending to a distance of about $`30R_{}`$ from the star, which they state “may result from reconnection flares at the star-disk interface”. The size of the radio structure ($`10R_{}`$, T Tau S having a radius $`R3R_{}`$) is very similar to the loop lengths we have derived here from the analysis of the large X-ray flares, providing independent support for the presence of such long, coherent magnetic structures in YSOs.
### 6.1. The role of disks and of accretion
To investigate the role of disks in the presence of long magnetic structures one would ideally like to have a clear separation of the sample in stars with disk and stars without. Unfortunately, the available diagnostics do not allow us to do this in an unambiguous way. The most complete investigation on the presence of disks in the ONC is the work of Hillenbrand et al. (1998), who have compared the results of a number of indicators for the presence of disks in the ONC, including IR excesses in a number of color bands and Ca ii emission. Emission in Ca ii should indicate the presence of active accretion, and thus will miss “quiescent” disks. NIR excesses are often used as an indicator of the presence of disks, and Hillenbrand et al. (1998) have quantitatively modeled the expected excesses in a number of bands, using models very similar to the ones of Meyer et al. (1997) and converging on similar conclusions. As discussed in detail in Hillenbrand et al. (1998) the IR excess is expected to be produced by the disk material being heated both by the stellar radiation and by the friction induced by accretion. In practice, its magnitude depends on a number of factors, such as the relative importance of accretion vs. the size of the disk’s inner hole, the relative contrast between the disk emission and the photosphere and the system inclination. Disks with low accretion rate and with a large inner hole, generating only a far-IR excess could be difficult to see in the NIR (see also the discussion in Meyer et al., 1997).
Hillenbrand et al. (1998) conclude that the excess $`\mathrm{\Delta }(IK)`$ (determined by taking into account the star’s spectral type and thus effective temperature and intrinsic colors and its reddening) is the most effective single indicator of the presence of disks; measurements in in the $`L`$ bands would have be even better but were not available to Hillenbrand et al. (1998) except for few stars. However, while $`\mathrm{\Delta }(IK)`$ clearly is an effective indicator on a statistical basis, its value for determining whether individual stars have a disk is much less clear. The distributions shown e.g. in Fig. 10 of Hillenbrand et al. (1998) have strong tails to negative values (as far as $`\mathrm{\Delta }(IK)1`$) and a significant number of stars with Ca ii emission (thus, accreting systems) do not have $`\mathrm{\Delta }(IK)`$ excesses. Even more significantly, as discussed in Sect. 7.3 of Hillenbrand et al. (1998), a number of sources with disks visible in the HST images (“proplyds”, from the O’Dell & Wong, 1996 sample) do not have IR excesses, clearly showing that some sources with disks escape detection through IR excesses.
In the present paper, we have tried to analyze whether a relationship exists between the detection of long flaring structures and the presence of disks from the available indicators, bearing however in mind the above caveat. In Table 1 we report, for all the COUP sources analyzed in our sample, the available data from the compilation of Getman et al. (2005) – the original source for most of the IR data is Hillenbrand et al. (1998) – regarding Ca ii emission and excess in $`IK`$ colors. We also report whether the sources have an excess in the NIR $`JHKL`$ bands from inspection of Fig. 5, although (quoting Hillenbrand et al., 1998 literally) “calculating disk fractions simply by counting the relative number of stars outside of and inside of reddening vectors in observed color-color diagrams is clearly naive”.
Fig. 5 plots the IR color-color diagrams for all COUP sources for which $`JHKL`$ NIR photometry is available (from Getman et al., 2005), with the subsample studied in the present paper singled out. The slanted blue line in Fig. 5 is the observational mean locus of CTTSs in Taurus derived by Meyer et al. (1997). All sources lying to the right of the reddened photospheric locus in Fig. 5 (bounded by the green dashed lines) have NIR colors compatible (bearing in mind the above caveat) with the presence of disks similar to the ones present in Taurus YSOs. For objects whose NIR colors place them in the reddened photospheric region the situation is however much less clear as they could still have disks, perhaps with low accretion and large interior holes, generating negligible IR excesses. Indeed, as discussed by Meyer et al. (1997) – see their Fig. 5 – the most important parameter driving the NIR excess is the accretion rate, and for a given color excess there is a minimum accretion rate, thus disk heating rate, needed to generate a NIR excess, independent of central hole size and of the disk inclination. Also, large holes ($`R_{\mathrm{hole}}R_{}`$) will produce no NIR excess in the $`JHKL`$ bands, with any excess shifted redward of $`3\mu `$m.
From Fig. 5 it is evident that none of the COUP sources with strong flares has a significant $`HK`$ excess (with the exception of COUP 1608 and, marginally, COUP 597), while 6 sources (out of 13 for which $`L`$ photometry is available) show a $`KL`$ band excess. Of these 6 (COUPs 332, 454, 597, 848, 976, 1343), 5 have no $`HK`$ band excess, while COUP 597 is the only one possibly displaying excesses in both $`HK`$ and $`KL`$. In the framework of the Meyer et al. (1997) models, COUP 597 is the only one whose $`JHKL`$ colors would be compatible with the presence of an accreting disk with a moderate-size central hole, while the $`JHKL`$ colors of COUPs 332, 454, 848, 976 and 1343 (with $`KL`$ excess but no $`HK`$ excess) would be compatible with the presence of a central disk with a larger size central hole (few stellar radii) and low accretion rate. The remaining sources in our flaring sample have no measurable excess in these bands, which would be compatible either with their being diskless sources or their having disks characterized by a very low accretion rate and with large inner holes.
An additional complexity comes from the fact that star-disk magnetic coupling requires a gaseous disk at the corotation radius, while NIR excesses are only sensitive to the dust content of the disk. But, given sufficient heating, the dust might be sublimated relatively near the star leaving behind the needed gas disk without an NIR excess. As discussed by Muzerolle et al. (2003), for sufficiently large accretion rates, the inner rim lies beyond the corotation radius, so that pure gaseous disks must extend inside the dust rim.
The situation is however clearly more complex, as evident from the fact that some sources which fall, in the color-color diagrams of Fig. 5, in the region of normal reddenened photospheres with no disks have significant $`\mathrm{\Delta }(IK)`$ excesses. Unfortunately the overlap between the two samples is not very large (13 out of 32 sources have $`JHKL`$ colors, 19 out of 32 have $`\mathrm{\Delta }(IK)`$ determinations, and only 6 have both), but all 6 sources for which both indicators are available have a $`\mathrm{\Delta }(IK)`$ excess, while 3 of them fall in the region of reddened normal photospheres in Fig. 5. In the whole sample of 1616 COUP sources for which optical data are available, 472 have a $`\mathrm{\Delta }(IK)`$ determination; of these, 388 have $`\mathrm{\Delta }(IK)>0.02`$ (82%), compatible with an excess due to circumstellar material, while 71 have $`\mathrm{\Delta }(IK)<0.02`$. In the present sample, $`\mathrm{\Delta }(IK)>0.02`$ for 30 out of 32 sources (94%), a comparable number as the complete sample given the low statistics.
Further evidence that disks may be present also when no excess in $`JHKL`$ NIR colors is observed comes from the COUP data directly. Tsujimoto et al. (2005) have recently searched the COUP databases for sources showing significant Fe 6.4 keV fluorescent emission, of which they found 7 cases. In all cases, fluorescent emission was detected during intense flares, a fact easily explained by the need for a sufficient number of hard ($`E>7.11`$ keV) photons to excite the fluorescence. Two of the Tsujimoto et al. (2005) sources are also in our large flare sample, namely COUP 649 and 1040. The remaining 5 fluorescing sources have too low statistics (either because of the short duration of the flare or because of the low peak count rate) for our analysis. Flares are however not a necessary condition for fluorescent emission, as shown by the detection of 6.4 keV fluorescence in the $`\rho `$ Oph YSO Elias 29 in quiescence by Favata et al. (2005). At the same time, the large equivalent widths of the 6.4 keV line observed in the COUP sample (as well as in Elias 29) are compatible only (as discussed by Favata et al., 2005 and Tsujimoto et al., 2005) with the fluorescence originating in a centrally illuminated disk observed face on. A disk illuminated from above, or other topologies (such as fluorescence from the photosphere) would result in a much lower equivalent width, likely too low to be detected on COUP sources. Yet most of the Tsujimoto et al. (2005) sources (and in particular COUP 649 and 1040), for which a disk is certainly present, given the observed X-ray fluorescence, show no $`JHKL`$ NIR excess, although some of them (e.g. COUP 649) do show a significant excess in $`\mathrm{\Delta }(IK)`$. Again, $`JHKL`$ NIR excesses appear to be, in the ONC, a sufficient condition for a disk to be present, but not a necessary one, with $`\mathrm{\Delta }(IK)`$ providing a more sensitive indicator.
Given the biases affecting the samples and the problems affecting disk diagnostics, we refrain from formal correlation analysis on e.g. flaring loop size and the presence of disk, limiting the analysis to qualitative considerations. Using $`\mathrm{\Delta }(IK)`$ as an indicator, the only two sources with $`\mathrm{\Delta }(IK)<0`$ in the present sample (which we take to imply that disks are most likely not present) both have flares confined to compact magnetic structures (COUP 7 with $`L=0.1R_{}`$ and COUP 960 with $`L=0.3R_{}`$), while both compact and large magnetic structures are present in the sources with $`\mathrm{\Delta }(IK)>0`$.
The 6 sources with significant $`KL`$ excess comprise flares spanning both large loops and more compact structures (comparable to the stellar size): large structures appear present on COUP 332 (although the determination of loop size suffers from a large uncertainty), COUP 454, with $`L3\times 10^{12}`$ cm, COUP 848, with $`L2\times 10^{12}`$ cm, and COUP 1343, with $`L10^{12}`$ cm. On the other hand COUP 597 has $`L<6\times 10^{11}`$ cm and COUP 976 has $`L<8\times 10^{11}`$ cm. Large loops are also present in sources with no $`KL`$ excess, such as COUP 669 ($`L10^{12}`$) and COUP 1083 ($`L2\times 10^{12}`$ cm), so that no clear correlation appears present between the presence of large loops and $`KL`$ excess. With one single exception (COUP 1608) none of the sources in the present sample show significant $`HK`$ excesses.
Only one strongly accreting source (COUP 141, with Ca ii equivalent width $`17.8`$ Å) is present in our sample of large flares; two sources (COUP 1114 and COUP 1608) show very modest emission (and presumably accretion rates), and all other sources show (when available) Ca ii in absorption or filled in and thus little or no ongoing accretion. The statistical significance of this is difficult to assess given the biases of our sample (discussed in Sect. 2), and the fact (Preibisch et al., 2005) that accreting sources in the COUP sample are statistically intrinsically fainter than the non-accreting sample (as well as more absorbed), introducing a small bias against strong flares from accreting sources being present in our sample. Interestingly, COUP 141 is one of the cases of compact loops, with $`LR_{}`$, so that in this case no disk-photosphere magnetic structures need to be postulated, and the flare is likely to be a normal coronal event, similar to the ones observed in more evolved stars. COUP 1114 and COUP 1608 both only result in upper limits to the size of the flaring structure.
In the complete COUP sample, out of 1616 sources for which optical data are available, Ca ii data are available for 537 sources. Of these, 198 have Ca ii in absorption, 189 have Ca ii in emission and 150 have filled lines (equivalent width of the Ca ii lines reported as “0.0” in Getman et al., 2005). Among the sources with Ca ii determination the fraction of sources in emission is then 35%, while in the flaring sample discussed in the present paper only 3 sources out of 20 for which the data are available have Ca ii in emission, i.e. 15%. Given the small numbers, and the biases present in the various samples, it’s difficult to assess the significance of the difference. If the lack of long flaring structures in accreting YSOs is however real, given that long magnetic structures connecting the disk to the star are certainly present, it implies that active accretion inhibits the formation of strong flares in the extended magnetic structures connecting the disk to the star. Long and intense flares (and thus very hot plasma) appear to be present in non-accreting structures connecting the disk to the star.
## 7. Conclusions
In the complete COUP sample of YSOs in the ONC, 32 flares have sufficient statistics (in terms of flare duration and photon count rate) and regular evolution to grant a detailed analysis of their decay. The type of flare decay analysis employed here allow us to derive the size of the flaring loops in the presence of heating extending into the decay phase, and, under simple assumptions, estimates of the plasma density and confining magnetic field.
The most notable result from our analysis is the strong evidence for very large flaring structures in these stars. The magnetic structures confining the plasma in a number of cases are much larger than the stars themselves. Among the events analyzed, a large fraction have very high peak temperatures ($`T>100`$ MK), and some are very long-lasting, with the longest flares extending to up to a week in duration. The peak X-ray luminosity for the most intense flare reaches $`8\times 10^{32}`$ erg s<sup>-1</sup>, comparable to other very bright flares observed in YSOs.
Our sample is limited to the brightest 1% of COUP flares, and the results may not be representative of the average magnetic reconnection event in Orion stars. Nevertheless, the present results show that very extended magnetic structures confining hot flaring plasma exist in YSOs. Structures of comparable sizes have never been seen in more evolved stars, and, given the short rotational periods of many of the flaring COUP stars, would not be stable if anchored onto the photosphere with both footpoints. These long magnetic structures are in the present paper interpreted as linking the stellar photosphere with the inner rim of the circumstellar disk.
The available indicators of the presence of disks do not allow for a formal analysis of whether large flaring loops are linked to the presence of disks; in particular stars without NIR excesses are not necessarily diskless, as they may e.g. have disks with little dust. However in a limited number of sources the NIR data clearly point to disks being present, either as $`KL`$ excess in color-color diagrams or as strong excesses above the de-reddened photospheric colors ($`\mathrm{\Delta }(IK)>0.4`$). Both small flaring structures (likely anchored on the photosphere only) and long flaring structures, with sizes which would extend well into the disk (assuming the disk to be truncated at the co-rotation radius) are present in these stars. Both types of structures (loops anchored on the star only and loops connecting the star to the disk) are then perhaps likely to coexist in YSOs with disks.
For the two sources in our sample without any $`\mathrm{\Delta }(IK)`$ excess, and for the star with a strong Ca ii accretion signature, the flaring structures are relatively compact. This suggests that inner disks are needed for reconnection in star-disk magnetic loops and that heating to X-ray temperatures is inhibited by mass loading associated with YSO accretion. More statistics are clearly needed to corroborate the present conclusions, and hopefully future, already approved long X-ray observations of other star-forming regions will be able to supply the needed additional evidence.
The data presented in the present paper thus constitute the first direct evidence for the presence of very long magnetic structures in YSOs, and, by inference, of magnetic structures linking the stellar photosphere with the circumstellar disk. Such structures are postulated by magnetospheric models of YSO accretion, but had not been directly detected to date.
We would like to thank Frank Shu (National Tsing Hua Univ.) and Ronald Taam (Northwestern Univ.) for the thoughtful and stimulating discussions, and an anonymous referee for the careful reading of the original manuscript. COUP is supported by *Chandra* guest observer grant SAO GO3-4009A (E. D. Feigelson, PI). EDF is also supported by NASA contract NAS8-38252. E. F., F. R., G. M. and S. S. acknowledge financial support from the Ministero dell’Istruzione dell’Università e della Ricerca.
## Appendix A: Notes on selected individual events
The present Appendix presents a detailed discussion (together with the relevant light curves and diagrams from the analysis) for a number of individually selected COUP sources. Its purpose is both to discuss in detail some peculiar or particularly representative event and to show the variety of flare shapes and characteristics present among the ONC YSOs.
### A.1. COUP 28
COUP 28 is a $`0.53M_{}`$, $`2.3R_{}`$ star of spectral type M0, showing a NIR excess ($`\mathrm{\Delta }(IK)=0.30`$) and no evidence for active accretion. Its rotational period is 4.4 d. It displays one of the longest flares (Fig. 6) of the COUP sample, with the total duration of the event extending over more than a week. Unfortunately an observation gap has blocked much of the flare rise, which however appears to be slow and extending over more than one day. The peak seems to be observed, and this allows a reliable analysis of the decay phase. The sources is at a large off-axis angle, so that no pile-up is present even if the flare at peak reaches a high count rate. The flare’s peak count rate is almost $`100\times `$ the source characteristic level. The peak temperature, at $`80`$ MK, is a moderate one for COUP flaring sources (although it is very high by the standard of flares observed in older stars), and is still well determined on ACIS spectra. The $`\mathrm{log}T`$ vs. $`\mathrm{log}\sqrt{EM}`$ diagram shows, as typical in both solar and stellar flares, the temperature peaking well before the $`EM`$. The excellent stastistics allows us to follow the flare’s evolution in detail, and the temperature decay shows evidence for reheating (in block 9, Fig. 6), as also reflected in the change in slope of the light curve decay, which is not a simple exponential but rather shows evidence for two different time scales. In this case, fitting an average decay to both the light curve and the slope in the $`\mathrm{log}T`$ vs. $`\mathrm{log}\sqrt{EM}`$ diagram still provides a good estimate of the parameters of the flaring region (Reale et al., 2004). At $`L=5.5\times 10^{11}\mathrm{cm}=1.9R_{}`$ the flare is a relatively compact one among the ones observed in COUP.
### A.2. COUP 43
COUP 43 is a known SB2 binary; the mass and radius computed by Hillenbrand (1997) for the unresolved source are $`0.40M_{}`$ and $`2.9R_{}`$. Its spectral type is M1, and it shows no evidence of active accretion; an $`IK`$ excess ($`\mathrm{\Delta }(IK)=0.50`$) is present. The flare (Fig. 7) is a double event, with the second flare beginning while the first one is still in its decay. The rise phase of the first event is lost in an observation gap, and therefore only the second event, with a duration of about two days, and which shows an impulsive rise, has been analyzed. While the limited statistics only allow us to subdivide the decay in three intervals, the steep $`\zeta `$ implies very limited sustained heating, and thus a long structure ($`L=1.1\times 10^{12}\mathrm{cm}=5.5R_{}`$) driven by the slow decay ($`\tau _{\mathrm{lc}}=62`$ ks). The two flares are sufficiently well separated, so that the first (incomplete) event does not significantly affect the spectral analysis of the second event. However, the lack of coverage of the latest phases of the decay imply caution in the interpretation of the analysis results, as a flattening of the temperature decay could imply a final shallower $`\zeta `$ and thus a smaller loop size.
### A.3. COUP 262
COUP 262, with $`1.1M_{}`$ and $`1.6R_{}`$ is a “young Sun”, and thus of particular interest for flare studies. In the context of COUP, its flaring activity is also discussed by Wolk et al. (2005). With spectral type K5, it shows no evidence for active accretion and has a rather strong $`IK`$ excess ($`\mathrm{\Delta }(IK)=2.2`$), pointing to a relatively massive but inactive disk being present (although no $`HK`$ excess is present). Although the flare (Fig. 8) has a very sharp peak, it decays slowly ($`\tau _{\mathrm{lc}}=120`$ ks), lasting for a few days. The event represents a relatively modest increase over the characteristic level of a factor of 4, but the shallow $`\zeta `$ implies a high level of sustained heating during the decay, so that (also given the high temperature, $`T>100`$ MK) the resulting flaring loop is very large, $`L=2.0\times 10^{12}\mathrm{cm}=18R_{}`$, although the significant uncertainty in $`\zeta `$ implies a lower confidence range of $`L=3.6R_{}`$. Nevertheless, even at the lower end of the confidence range, the flaring region is larger than normally found in active stars, providing evidence that large loops, likely connecting the star to the disk, must have existed for the young Sun.
### A.4. COUP 597
COUP 597 is a somewhat more massive star than COUP 262, with $`M=1.5M_{}`$ and $`R=2.0R_{}`$. The strong Ca ii absorption shows it not to be actively accreting. The flare (Fig. 9) appears to be the superposition of two events, a less intense and more long-lasting one and a more intense one of which only the very fast decay is visible, as the rest falls into an observation gap. As the analysis relies on the delay of the second, slower flare, the first, fast event should have no effect on the results. The peak temperature for the longer event (determined during the rise phase) may be affected at some level from the presence of the first flare. However, given the very moderate peak temperature, and the weak dependance (as $`\sqrt{T}`$) of the loop size on peak temperature, this will not have a strong effect on the results.
Of interest here is the slow rise of the longer flare, which appears to take place over a day or so. The loop size resulting from the analysis is $`L=1.6R_{}`$, thus a moderate size loop. The size is largely driven by the slow temperature decay and resulting shallow $`\zeta `$ (with $`F(\zeta )=9.7`$, implying strong sustained heating). The slow rise can therefore imply either a very slowly rising heating function, or perhaps (using the results of the modeling of COUP 1343, Sect. 4.1) a heating source distributed along the loop (or even concentrated near the loop top) rather than localized at the loop footpoints.
### A.5. COUP 649
COUP 649, is a typical low-mass member of the ONC, with $`0.4M_{}`$ and $`2.2R_{}`$. It shows no evidence for active accretion and has an $`IK`$ excess ($`\mathrm{\Delta }(IK)=0.4`$).
The flaring event has a very peculiar evolution (Fig. 10), with an initial very slow rise, lasting for about a day, followed by a plateau at about $`10\times `$ the characteristic rate; after another half day, an impulsive peak increases the count rate by another factor of two, and then the decay begins. Notwithstanding the oddly shaped light curve, the event shows some of the characteristics of a magnetically confined flaring structure, in particular the temperature which peaks while the count rate is still increasing, and has already started to decrease on the plateau. Somewhat peculiarly, the secondary sharp peak preceding the decay does not result in an increase in the temperature.
The limited statistics does not allow a detailed analysis, with only two points on the flare’s decay. The peak temperature ($`T=80`$ MK) is moderate by the standard of the other large flares analyzed here, and the nominal length of the flaring structure is $`L=6.4\times 10^{11}\mathrm{cm}=4.2R_{}`$, although with a large uncertainty driven by the uncertainty on $`\zeta `$.
### A.6. COUP 669
COUP 669 is a relatively massive object, with $`1.5M_{}`$ and $`2.6R_{}`$, and an $`IK`$ excess ($`\mathrm{\Delta }(IK)=0.4`$). The flare (Fig. 11) is a well defined impulsive event, for which the rise phase has however been cut by an observing gap. The peak temperature, at $`T=79`$ MK is relatively modest, but the fast decay of the temperature implies limited sustained heating, resulting in a flaring loop with a reasonably well determined size, $`L=9.2\times 10^{11}\mathrm{cm}=5.1R_{}`$, typical of the large flaring loops in the COUP sample.
### A.7. COUP 752
Another typical low-mass object, COUP 752, has $`0.5M_{}`$ and $`1.7R_{}`$, no evidence for ongoing accretion, and an $`IK`$ excess ($`\mathrm{\Delta }(IK)=0.6`$). COUP 752 undergoes a very intense flare (Fig. 12), with a peak count rate more than $`100\times `$ the characteristic one. The event lasts over 5 days, and the rise phase is long (half a day) and well observed. The excellent statistics allow for a detailed analysis of the event, although the initial part of the decay unfortunately falls in an observing gap. The flare’s peak temperature slightly exceeds 100 MK, still falling in the domain in which the ACIS detector gives reliable temperature determinations.
Notwithstanding the very large increase in count rate, the event follows well the evolution in the $`\mathrm{log}T`$$`\mathrm{log}\sqrt{EM}`$ plane for a flare confined in a single loop. The temperature peaks (block 1) when the emission measure has just started to rise (and still is only about one tenth of the peak value), and then decreases to $`50`$ MK when the emission measure has peaked. The decay does not follow a simple exponential law; rather, there is clear evidence for sustained heating, with well defined reheating events. In particular, in block 11 (immediately after the first observing gap) the temperature increases again briefly, and the decay of both the light curve and the emission measure changes slope, clear evidence for reheating of the plasma.
Our analysis shows a very shallow $`\zeta `$, indicative of the presence of strong ongoing heating during the flare decay, as it is also evident from the increase in plasma temperature. As a consequence, only an upper limit to the flaring structure’s size can be derived, at $`L7.4\times 10^{10}\mathrm{cm}`$, with a best-fit value $`L=6.5\times 10^{10}\mathrm{cm}=5.6R_{}`$. The resulting minimum confining magnetic field, at $`B250`$ G is relatively strong.
### A.8. COUP 891
COUP 891 is the most massive object in the COUP flaring sample, with $`2.4M_{}`$, $`4.9R_{}`$, a strong $`IK`$ excess ($`\mathrm{\Delta }(IK)=1.10`$), no evidence for accretion and significant obscuration ($`A_V=8.0`$ mag). The X-ray light curve suffers from significant pile-up, and therefore we have used, for our analysis, the list of photons extracted from an annular region free from pile-up, as described by Getman et al. (2005). This strongly reduces the number of photons available for a spectral analysis, so that only moderate statistics are available for this event.
Two flares are present in the COUP light curve (Fig. 13); however the first one is too short and cannot be analyzed with the approach adopted here. The second event is a well defined impulsive event, with a fast rise and a regular, exponential decay, well traced for more than four days. At the peak, the flare count rate is $`50\times `$ larger than in quiescence. The peak temperature is slighly above 100 MK, but still within the range of the ACIS instrument. The fast decay of temperature results in a fairly steep $`\zeta `$, and thus in little sustained heating. The long decay time ($`\tau _{\mathrm{lc}}=73`$ ks) therefore results in a long flaring structure, $`L=1.7\times 10^{12}\mathrm{cm}=5.1R_{}`$.
### A.9. COUP 971
COUP 971 is a low mass ($`0.7M_{}`$) star with $`R=3.3R_{}`$. The Ca ii in absorption points to a non-accreting system. The flaring event (Fig. 14) is a typical example of a short impulsive flare. The decay in the $`\mathrm{log}T`$$`\mathrm{log}\sqrt{EM}`$ plane is steep and regular, pointing to little if any residual heating. The result is a moderate size loop ($`L=1.6R_{}`$), similar to the ones resulting from the analysis of intense flares in more evolved stars; this is likely an example of an event confined in a structure anchored to the stellar photosphere only.
### A.10. COUP 1083
Little optical information is available for COUP 1083, and in particular no mass and radius estimate is known. A rotational period has however been derived ($`P=5.9`$ d) for this star. Its spectral type has been estimated at M0, showing that the object must have a moderate mass.
The flare present in the COUP 1083 lightcurve has a peculiar shape, well visible in Fig. 15, with a very slow rise, lasting more than one day, and a decay lasting about 3 days. The peak of the flare unfortunately falls in an observing gap, and it thus not visible. We have analyzed the event assuming that the observed maximum is the actual peak, but in practice the peak might be higher, although this would not substantially affect our conclusions.
Notwithstanding the peculiar light curve, the evolution of the event follows the one theoretically expected for a confined flare, with the temperature peaking prior to the $`EM`$ and a regular decay in the $`\mathrm{log}T`$$`\mathrm{log}\sqrt{EM}`$ plane. A second event may be starting in the final phase of the decay of the flare, on day 12, as shown both by the increasing count rate and by the increasing temperature. The second flare is cut by the observation gap and thus nothing can be derived for it. We have not included the last ML block in the analysis.
Although the statistics are limited (and thus the error bars large), the derived steep $`\zeta `$ is compatible with limited sustained heating, so that a long structure, $`L=2.4\times 10^{12}`$ cm, is derived from the analysis due to the slow flare decay ($`\tau _{\mathrm{lc}}=125`$ ks). The large size is also compatible with the very slow rise time, as discussed for the detailed model of COUP 1343 (Sect. 4.1).
### A.11. COUP 1246
COUP 1246 is one of the lowest-mass stars in the present sample, with $`0.2M_{}`$ and $`1.6R_{}`$. It shows no evidence for active accretion, but it has a sizable $`K`$-band IR excess ($`\mathrm{\Delta }(IK)=0.8`$) and is one of the sources showing an $`KL`$ excess but not a $`HK`$ band one; its rotational period has been determined at $`P=5.2`$ d. The flare (Fig. 16) is a very well defined impulsive event, with a fast rise and a peak count rate $`100\times `$ the characteristic rate. Also in this case the peak observed temperature is somewhat above 100 MK, but still within the range of reliability of ACIS. The decay in the $`\mathrm{log}T`$$`\mathrm{log}\sqrt{EM}`$ plane is well defined, and closely follows the one expected from the models of confined flares, allowing to analyze the event in some detail. The light curve plotted in the top panel of Fig. 16 shows clear evidence for a second, smaller flare developing on the tail of the first one, resulting in an increase in temperature at the flare end. Statistics are insufficient to analyze the second flare at any level of detail, and therefore the last block has not been included in the analysis. The resulting size of the flaring region is $`L=4.0\times 10^{10}\mathrm{cm}=3.8R_{}`$, i.e. a typical event for the COUP sample.
### A.12. COUP 1343
Very little is known about the optical counterpart to COUP 1343 (the analysis for the flare is shown in Fig. 1), with no mass or radius estimates known, although a rotational period measurement ($`P=8.7`$ d) is available. The flare is a very well defined event, with a clean impulsive rise and an undisturbed decay. The peak flare luminosity is $`>10\times `$ larger than the characteristic source luminosity, allowing a clean analysis of the flare spectral characteristics. In the analysis we have ignored the last segment of the flare (the one extending into day 13), for which the temperature has too large an error (also because of the data gap). The peak temperature of the flare is very high, with a nominal best-fit temperature of $`400`$ MK. As discussed above such high temperatures cannot be reliably determined with ACIS, and therefore the resulting flare size (based on the assumption $`T_{\mathrm{obs}}=100`$ MK) is to be considered as a lower limit.
The rapid temperature decrease implies that little sustained heating is present, so that the observed decay is close to the thermodynamic decay of the loop. The smooth decay also argues for the lack of significant reheating. Therefore, with $`\tau _{\mathrm{lc}}=39`$ ks and $`T_{\mathrm{obs}}100`$ MK, the analysis results in $`L1\times 10^{12}`$ cm, or $`L14R_{}`$. Even if COUP 1343 were a relatively massive star, with a radius of a few $`R_{}`$, the flaring structure would still be several times the size of the star. The resulting magnetic field, while moderate in itself (150 G), must be coherently organized over some 0.1 AU.
As described in Sect. 4.1 the evolution of the same event has also been simulated in detail, confirming that the temperature and emission measure evolution are the result of a flare confined in a large loop of the same size of the one found by the decay analysis discussed here, with impulsive heating.
### A.13. COUP 1410
COUP 1410 appears more evolved than the rest of the sample, with an estimated age $`t=3.6\times 10^7`$ yr, and a consequently small radius ($`0.51R_{}`$) for its mass ($`0.36M_{}`$). It shows no evidence for accretion from Ca ii emission, but it has a significant $`K`$-band excess ($`\mathrm{\Delta }(IK)=2.3`$). The rotational period has been measured at $`P=6.76`$ d.
The flare (Fig. 17) has a well defined impulsive rise, followed by a very long decay extending for over a week, with $`\tau _{\mathrm{lc}}=150`$ ks. Even though significant heating is present (as shown by the shallow $`\zeta `$), the resulting size of the flaring structure is $`L=1.1\times 10^{12}\mathrm{cm}=55R_{}`$, very long with respect to the (small) stellar radius. However, the uncertainty on $`\zeta `$ results in a significant uncertainty on $`L`$, which has a lower range of $`L=1\times 10^{11}\mathrm{cm}=5R_{}`$, the same large size as found for the best determined large events in the COUP sample, making this event a good candidate for a star-disk connecting structure.
### A.14. COUP 1568
COUP 1568 is another relatively massive object in the flaring sample, with $`2.6M_{}`$ and $`4.0R_{}`$. It has a small IR excess ($`\mathrm{\Delta }(IK)=0.2`$).
The flare (Fig. 18) is relatively short, lasting just one day, but it is intense, peaking at over $`100\times `$ the characteristic rate, with excellent statistics due to the source’s off-axis angle which prevents pile-up. The impulsive phase is very sharp, and the decay almost a perfect exponential. The peak temperature is very high: the best-fit values for the ACIS spectra at peak are $`T500`$ MK, however with a very large uncertainty which always include values as low as 100 MK, due to the ACIS spectral response.
Once in the decay phase, the temperature decays very slowly, resulting in a very shallow $`\zeta `$. Therefore only an upper limit to size of the flaring structure can be obtained, $`L1.5\times 10^{11}`$ cm, with a best fit value $`L=3.7\times 10^{10}\mathrm{cm}=0.4R_{}`$. Unlikely many of the large COUP flares studied in this paper, this event is very likely to be compact, with size smaller than the star itself. A small size is also implied by the very fast rise phase (which constrains the time scale with which the chromospheric evaporation will fill the flaring loop). Given the compact size, together with the large peak emission measure and high temperature, the peak flare density will likely be high ($`n_e=1.2\times 10^{12}`$ cm<sup>-3</sup>), as must the confining magnetic field ($`B=3.5`$ kG) |
warning/0506/hep-ph0506194.html | ar5iv | text | # Dirac-Schrödinger equation for quark-antiquark bound states and derivation of its interaction kernel
## I Introduction
It is the common recognition that the Bethe-Salpeter (B-S) equation which was proposed early in Refs. is a rigorous formalism for relativistic bound states. The prominent features of the equation are: (1) The equation is derived from the quantum field theory and hence set up on the firm dynamical basis; (2) The interaction kernel in the equation contains all the interactions taking place in the bound states and therefore the equation provides a possibility of exactly solving the problem of relativistic bound states; (3) The equation is elegantly formulated in a manifestly Lorentz-covariant form in the Minkowski space which allows us to discuss the equation in any coordinate frame. However, there are tremendous difficulties in practical applications of the equation, particularly, for solving the nuclear force in the nuclear physics and the quark confinement in hadron physics. One of the difficulties arises from the fact that the kernel in the equation was not given a closed form in the past. The kernel usually is defined as a sum of B-S (two-particle) irreducible Feynman diagrams each of which can only be individually determined by a perturbative calculation. This definition is, certainly, not suitable to investigate the subjects such as the nuclear force and the quark confinement which must necessarily be solved by a nonperturbative method. This is why, as said in Ref. ,” The Bethe-Salpeter equation has not led to a real breakthrough in our understanding of the quark-quark force”. Opposite to the conventional concept as commented in Ref. that ”The kernel $`K`$ can not be given in closed form expression”, we have derived a closed expression of the B-S kernel for quark-antiquark bound states in a recent publication . The expression derived contains only a few types of Green’s functions which not only are easily calculated by the perturbation method, but also suitable to be investigated by a certain nonperturbation approach. Another difficulty of solving the B-S equation was attributed to the four-dimensional nature of the equation because the relative time (or the relative energy) would lead to unphysical solutions. Therefore, many efforts were made in the past to recast the four-dimensional equation in three-dimensional ones in either approximate manners or exact versions such as the instantaneous approximation , the quasipotential approach \[7-12\] and the equal-time formalism \[13-16\].
As one knows, the four-dimensionally covariant B-S equation for a two-fermion system is ordinarily formulated in a second-order differential equation with respect to the space-time variables in the position space. This kind of equation has been shown to have unphysical solutions with the negative norm. It was pointed out in Ref. that” The appearance of the negative-norm B-S amplitude is a quite common phenomenon in the B-S equation”. A similar phenomenon was encountered in the Klein-Gordon (K-G) equation
$$(\mathrm{}_x+m^2)\psi (x)=0$$
(1)
which was originally viewed as the wave equation satisfied by the single free fermion states. It is well-known that the K-G equation, as a second-order differential equation, has a solution with negative probability. This is because the wave function of the equation is determined not only by its initial value $`\psi (0)`$, but also by the initial value of the time-differential$`\frac{\psi }{t}_{t=0}`$ which would possibly cause the solution to have negative probability . Nevertheless, the Dirac equation
$$(i_xm)\psi (x)=0$$
(2)
where $`_x=\gamma ^\mu _\mu ^x`$ has not the negative norm solution because it is a first-order differential equation. As widely recognized, the Dirac equation gives a correct description of the single free fermion states.
Analogous to the case of single fermion, the relativistic states for a two-fermion system may also be formulated by a set of first-order differential equations just as the Hamilton equation in Mechanics and the Maxwell equation in Electrodynamics which are equivalent to the second-order differential equations, i.e. the Lagrange equation and the D’Alembert equation respectively. Motivated by this idea, it was proposed in the literature \[20-28\] that the quark-antiquark bound system may be described by two coupled Dirac equations which are constructed in accord with the Dirac’s Hamiltonian constraint formalism such that \[23-25\]
$$[i_{x_1}m_1V_1(x_1,x_2)]\psi (x_1,x_2)=0$$
(3)
$$(i_{x_2}m_2V_2(x_1,x_2)]\psi (x_1,x_2)=0$$
(4)
where $`\psi (x_1,x_2)`$ denotes the two-fermion wave function, $`m_1`$ and $`m_2`$ are the masses of quark and antiquark, $`V_1`$ and $`V_2`$ stand for the effective potentials which are determined by the requirement of satisfying the Lorentz-invariance, the charge conjugation symmetry and a certain constraint (or say, compatibility) conditions. With a constraint imposed on the relative time, the above equations will be reduced to a three-dimensional eigenvalue equation.
As emphasized in the previous literature \[23-25\], Eqs. (1.3) and (1.4) are built up within the framework of relativistic quantum mechanics and the interaction potentials are given in a phenomenological way although they are inspired by the quantum field theory and, as demonstrated in Ref. , are linked with the corresponding B-S equation. Obviously, in order to understand the Dirac-type equations for the two-fermion system more precisely, it is necessary to give such equations an extensive investigation and an exact formulation from the viewpoint of quantum field theory. This just is the purpose of this paper. In this paper, we limit ourself to discuss the quark and antiquark ($`q\overline{q})`$ bound states. The results certainly suit to other two-fermion bound systems. First we derive two first-order differential equations for the quark-antiquark bound states from Quantum Chromodynamics (QCD) which describe the evolution of the bound state with the total (center of mass) time and the relative time respectively. These equations will be called Dirac-Schrödinger (D-S) equation because the Dirac equation is, in essence, the Schrödinger equation in the relativistic case which is identified with itself as the uniquely correct equation of describing the evolution of a quantum state with time in the quantum theory. Next, we concentrate our main attention on the interaction kernel appearing in the D-S equation. We are devoted to deriving a closed and explicit expression of the interaction kernel. The kernel will be derived by two different methods: one is to utilize equations of motion satisfied by the $`q\overline{q}`$ four-point Green’s function and some other four-point Green’s functions in which the gluon field is involved; another is to employ the technique of irreducible decomposition of the Green’s functions involved in the D-S equation. The first method is similar to that proposed previously in Ref. . The kernel derived by this method has a compact expression which contains only a few types of Green’s functions. The kernel derived by the second method is expressed in terms of the quark, antiquark and gluon propagators and some kinds of three, four and five-line proper vertices and therefore exhibits a more specific structure of the kernel. Especially, the kernel derived can not only be easily calculated by the perturbation method, but also provides a suitable basis for nonperturbative investigations. The D-S equation and its interaction kernel mentioned above are Lorentz-covariant. We will show how this equation and its kernel are reduced to the exact three-dimensional forms given previously in Ref. in the equal-time Lorentz frame. It is well-known that the D-S equation is represented in the Dirac spinor space. This equation actually is a coupled set of sixteen scalar equations. In practical applications, sometimes it is more convenient to reduce the D-S equation to the Pauli spinor space following the procedure proposed in Ref. . By this procedure, we will obtain an equivalent Pauli-Schrödinger (P-S) equation represented in the Pauli spinor space from the D-S equation. In the P-S equation, the interaction Hamiltonian is explicitly given in a series expression which has an one-to-one correspondence with the perturbative expansion of the S-matrix. To illustrate the applicability of the kernels derived and the equivalence between the aforementioned two differerent expressions of the kernel, we will show how the one-gluon exchange kernels in the D-S equation and the corresponding interaction Hamiltonian in the P-S equation can be derived from the closed expressions. Finally, we will discuss the relation between the D-S equation and the corresponding B-S equation.
The remainder of this paper is arranged as follows. In section 2, we will first derive two Dirac-type equations satisfied by the $`q\overline{q}`$ four-point Green’s function. From these equations, the D-S equations obeyed by the B-S amplitudes will be derived by making use of the Lehmann representation of the Green’s function . Then, we will show how the four-dimensional D-S equation is reduced to the three-dimensional one. In section 3, the first explicit expression of the interaction kernel in the D-S equation will be derived by virtue of the equations of motion satisfied by the Green’s functions involved in the Dirac-like equations. And, it will be shown how the closed expression of the exact three-dimensional kernel can be written out from the four-dimensional one. In section 4, the second expression of the interaction kernel will be derived by means of the technique of irreducible decomposition of the Green’s functions. In section 5, the D-S equation will be reduced to the corresponding P-S equation. Section 6 will be used to give a brief derivation and description of the one-gluon exchange kernels. The last section serve to discuss the relation between the D-S equation and the corresponding B-S equation and to make some remarks. In Appendix A, we will describe the derivation of the equations of motion satisfied by the Green’s functions which are necessary to be used in the derivation of the D-S equation and its interaction kernel. In Appendix B, the irreducible decomposition of the relevant Green’s functions will be performed for the purpose of deriving the second expression of the kernel.
## II Derivation of the Dirac-Schrödinger equation
The Dirac-Schrödinger (D-S) equation satisfied by the $`q\overline{q}`$ bound states may be derived from the corresponding equation for the $`q\overline{q}`$ four-point Green’s function which is defined in the Heisenberg picture as follows
$$𝒢(x_{1,}x_2;y_1,y_2)_{\alpha \beta \rho \sigma }=0^+\left|T\{N[\psi _\alpha (x_1)\psi _\beta ^c(x_2)]N[\overline{\psi }_\rho (y_1)\overline{\psi }_\sigma ^c(y_2)]\}\right|0^{}$$
(5)
where $`\psi (x)`$ and $`\psi ^c(x)`$ are the quark and antiquark field operators respectively, $`\overline{\psi }(x)`$ and $`\overline{\text{ }\psi }^c(x)`$ are their corresponding Dirac conjugates
$$\psi ^c(x)=C\overline{\psi }^T(x),\overline{\psi }^c(x)=\psi ^T(x)C^1$$
(6)
here $`C=i\gamma ^2\gamma ^0`$ is the charge conjugation operator, $`0^\pm `$ denote the physical vacuum states, $`T`$ symbolizes the time-ordering product and $`N`$ designates the normal product which is defined by
$$N[\psi _\alpha (x_1)\psi _\beta ^c(x_2)]=T[\psi _\alpha (x_1)\psi _\beta ^c(x_2)]0^+\left|T[\psi _\alpha (x_1)\psi _\beta ^c(x_2)]\right|0^{}$$
(7)
It is emphasized here that the above normal product can only be viewed as a definition in the Heisenberg picture. With the definition shown in Eq. (2.3), the Green’s function in Eq. (2.1) may be represented as
$$𝒢(x_{1,}x_2;y_1,y_2)_{\alpha \beta \rho \sigma }=G(x_{1,}x_2;y_1,y_2)_{\alpha \beta \rho \sigma }+S_F^{}(x_1x_2)_{\alpha \beta }\overline{S}_F^{}(y_1y_2)_{\rho \sigma }$$
(8)
where
$$G(x_{1,}x_2;y_1,y_2)_{\alpha \beta \rho \sigma }=0^+\left|T\{\psi _\alpha (x_1)\psi _\beta ^c(x_2)\overline{\psi }_\rho (y_1)\overline{\psi }_\sigma ^c(y_2)\}\right|0^{}$$
(9)
is the ordinary $`q\overline{q}`$ four-point Green’s function ,
| $`S_F^{}(x_1x_2)_{\alpha \beta }=\frac{1}{i}0^+\left|T\{\psi _\alpha (x_1)\psi _\beta ^c(x_2)\}\right|0^{}`$ |
| --- |
| $`=S_F(x_1x_2)_{\alpha \gamma }(C^1)_{\gamma \beta }=S_F^c(x_2x_1)_{\beta \lambda }C_{\lambda \alpha }`$ |
(10)
and
| $`\overline{S}_F^{}(y_1y_2)_{\rho \sigma }=\frac{1}{i}0^+\left|T\{\overline{\psi }_\rho (y_1)\overline{\psi }_\sigma ^c(y_2)\}\right|0^{}`$ |
| --- |
| $`=C_{\sigma \tau }S_F(y_2y_1)_{\tau \rho }=(C^1)_{\rho \delta }S_F^c(y_1y_2)_{\delta \sigma }`$ |
(11)
in which
$$S_F(x_1x_2)_{\alpha \gamma }=0^+\left|T\{\psi _\alpha (x_1)\overline{\psi }_\gamma (x_2)\}\right|0^{}$$
(12)
and
$$S_F^c(y_1y_2)_{\delta \sigma }=0^+\left|T\{\psi _\delta ^c(y_1)\overline{\psi }_\sigma ^c(y_2)\}\right|0^{}$$
(13)
are the ordinary quark and antiquark propagators respectively . It is clear that the propagators defined in Eqs. (2.6) and (2.7) are nonzero only for the quark and the antiquark which are of the same flavor. For the quark and antiquark of different flavors, the Green’s function defined in Eq. (2.1) is reduced to the ordinary form shown in Eq. (2.5) since the second term on the right hand side (RHS) of Eq. (2.4) vanishes. In the case of the quark and antiquark of the same flavor, the normal product in Eq. (2.1) plays a role of excluding the contraction between the quark field and the antiquark one from the Green’s function. Physically, this avoids the $`q\overline{q}`$ annihilation to break stability of a bound state. It would be pointed out that use of $`\psi ^c(x)`$ other than $`\overline{\psi }(x)`$ to represent the antiquark field in this paper has an advantage that the antiquark field would behave as a quark one in the D-S equation so that the quark-antiquark equation formally is the same as the corresponding two-quark equation in the case that the quark and antiquark have different flavors.
The equations of motion which describe the variation of the $`q\overline{q}`$ four-point Green’s function $`G(x_{1,}x_2;y_1,y_2)`$ with the coordinates $`x_1`$ and $`x_2`$ may easily be derived from the QCD generating functional as described in Appendix A. The results are
| $`(i_{x_1}m_1)_{\alpha \gamma }G(x_{1,}x_2;y_1,y_2)_{\gamma \beta \rho \sigma }=\delta _{\alpha \rho }\delta ^4(x_1y_1)S_F^c(x_2y_2)_{\beta \sigma }`$ |
| --- |
| $`+C_{\alpha \beta }\delta ^4(x_1x_2)\overline{S}_F^{}(y_1y_2)_{\rho \sigma }(\mathrm{\Gamma }^{a\mu })_{\alpha \gamma }G_\mu ^a(x_1x_1,x_2;y_1,y_2)_{\gamma \beta \rho \sigma }`$ |
(14)
| $`(i_{x_2}m_2)_{\beta \lambda }G(x_{1,}x_2;y_1,y_2)_{\alpha \lambda \rho \sigma }=\delta _{\beta \sigma }\delta ^4(x_2y_2)S_F(x_1y_1)_{\alpha \rho }`$ |
| --- |
| $`+C_{\alpha \beta }\delta ^4(x_1x_2)\overline{S}_F^{}(y_1y_2)_{\rho \sigma }(\overline{\mathrm{\Gamma }}^{b\nu })_{\beta \lambda }G_\nu ^b(x_2x_1,x_2;y_1,y_2)_{\alpha \lambda \rho \sigma }`$ |
(15)
in which
$$(\mathrm{\Gamma }^{a\mu })_{\alpha \gamma }=g(\gamma ^\mu T^a)_{\alpha \gamma },\text{ }(\overline{\mathrm{\Gamma }}^{b\nu })_{\beta \lambda }=g(\gamma ^\nu \overline{T}^b)_{\beta \lambda }$$
(16)
where $`g`$ is the coupling constant, $`T^a=\frac{\lambda ^a}{2}`$ and $`\overline{T}^a=\lambda ^a/2`$ are the quark and antiquark color matrices respectively,
| $`G_\mu ^a(x_ix_1,x_2;y_1,y_2)_{\alpha \beta \rho \sigma }`$ |
| --- |
| $`=0^+\left|T\{𝐀_\mu ^a(x_i)\psi _\alpha (x_1)\psi _\beta ^c(x_2)\overline{\psi }_\rho (y_1)\overline{\psi }_\sigma ^c(y_2)\}\right|0^{}`$ |
(17)
with $`i=1,2`$ are the new four-point Green’s function including a gluon field in it and the propagators were defined in Eqs. (2.6)-(2.9). It would be noted here that the terms related to $`\overline{S}_F^{}(y_1y_2)`$ in Eqs. (2.10) and (2.11) are absent when the quark and the antiquark have different flavors. The equations of motion satisfied by the Green’s function defined in Eq. (2.1) may be found by substituting the relation in Eq. (2.4) into Eqs. (2.10) and (2.11) and by making use of the following equations as mentioned in Appendix A
$$\begin{array}{c}[(i_{x_1}m_1)_{\alpha \gamma }S_F^{}(x_1x_2)_{\gamma \beta }=C_{\alpha \beta }\delta ^4(x_1x_2)(\mathrm{\Gamma }^{a\mu })_{\alpha \gamma }\mathrm{\Lambda }_\mu ^a(x_1x_1,x_2)_{\gamma \beta }\\ [(i_{x_2}m_2)_{\beta \lambda }S_F^{}(x_1x_2)_{\alpha \lambda }=C_{\alpha \beta }\delta ^4(x_1x_2)(\overline{\mathrm{\Gamma }}^{b\nu })_{\beta \lambda }\mathrm{\Lambda }_\nu ^b(x_2x_1,x_2)_{\alpha \lambda }\end{array}$$
(18)
where
$$\mathrm{\Lambda }_\mu ^a(x_ix_1,x_2)_{\alpha \beta }=\frac{1}{i}0^+\left|T\{𝐀_\mu ^a(x_i)\psi _\alpha (x_1)\psi _\beta ^c(x_2)\}\right|0^{}$$
(19)
with $`i=1,2`$ are a kind of quark-antiquark-gluon Green’s function. The results are
| $`(i_{x_1}m_1)_{\alpha \gamma }𝒢_{\gamma \beta \rho \sigma }(x_{1,}x_2;y_1,y_2)=\delta _{\alpha \rho }\delta ^4(x_1y_1)S_F^c(x_2y_2)_{\beta \sigma }`$ |
| --- |
| $`(\mathrm{\Gamma }^{a\mu })_{\alpha \gamma }𝒢_\mu ^a(x_1x_1,x_2;y_1,y_2)_{\gamma \beta \rho \sigma }`$ |
(20)
and
| $`(i_{x_2}m_2)_{\beta \lambda }𝒢_{\alpha \lambda \rho \sigma }(x_{1,}x_2;y_1,y_2)=\delta _{\beta \sigma }\delta ^4(x_2y_2)S_F(x_1y_1)_{\alpha \rho }`$ |
| --- |
| $`(\overline{\mathrm{\Gamma }}^{b\nu })_{\beta \lambda }𝒢_\nu ^b(x_2x_1,x_2;y_1,y_2)_{\alpha \lambda \rho \sigma }`$ |
(21)
where
| $`𝒢_\mu ^a(x_ix_1,x_2;y_1,y_2)_{\alpha \beta \rho \sigma }=0^+|T[\{N[𝐀_\mu ^a(x_i)\psi _\alpha (x_1)\psi _\beta ^c(x_2)]N[\overline{\psi }_\rho (y_1)\overline{\psi }_\sigma ^c(y_2)]\}|0^{}`$ |
| --- |
| $`=G_\mu ^a(x_ix_1,x_2;y_1,y_2)_{\alpha \beta \rho \sigma }+\mathrm{\Lambda }_\mu ^a(x_ix_1,x_2)_{\alpha \beta }\overline{S}_F^{}(y_1y_2)_{\rho \sigma }`$ |
(22)
here $`i=1,2`$. In the above, the normal products defined in the same way as that in Eq. (2.3). In comparison of Eqs. (2.16) and (2.17) with Eqs. (2.10) and (2.11), we see, the terms related to $`\overline{S}_F^{}(y_1y_2)`$ in Eqs. (2.10) and (2.11) disappear in Eqs. (2.16) and (2.17). Therefore, the equations (2.16) and (2.17) formally are the same as given in the case that the quark and the antiquark are of different flavors.
Multiplying Eqs. (2.16) and (2.17) with the matrices $`\gamma _1^0`$ and $`\gamma _2^0`$ respectively, we have
| $`[i\frac{}{t_1}h^{(1)}(\stackrel{}{x_1})]_{\alpha \gamma }𝒢_{\gamma \beta \rho \sigma }(x_1,x_2;y_1,y_2)`$ |
| --- |
| $`=(\gamma _1^0)_{\alpha \rho }\delta ^4(x_1y_1)S_F^c(x_2y_2)_{\beta \sigma }+𝒢_{\alpha \beta \rho \sigma }^{(1)}(x_1,x_2;y_1,y_2)`$ |
(23)
and
| $`[i\frac{}{t_2}h^{(2)}(\stackrel{}{x_2})]_{\beta \lambda }𝒢_{\alpha \lambda \rho \sigma }(x_1,x_2;y_1,y_2)`$ |
| --- |
| $`=(\gamma _2^0)_{\beta \sigma }\delta ^4(x_2y_2)S_F(x_1y_1)_{\alpha \rho }+𝒢_{\alpha \beta \rho \sigma }^{(2)}(x_1,x_2;y_1,y_2)`$ |
(24)
where
$$h^{(i)}(\stackrel{}{x_i})=i\stackrel{}{\alpha _i}_{\stackrel{}{x_i}}+m_i\gamma _i^0$$
(25)
is the i-th free fermion Hamiltonian,
$$\begin{array}{c}𝒢_{\alpha \beta \rho \sigma }^{(1)}(x_1x_1,x_2;y_1,y_2)=(\mathrm{\Omega }_1^{a\mu })_{\alpha \gamma }𝒢_\mu ^a(x_1x_1,x_2;y_1,y_2)_{\gamma \beta \rho \sigma }\\ 𝒢_{\alpha \beta \rho \sigma }^{(2)}(x_2x_1,x_2;y_1,y_2)=(\mathrm{\Omega }_2^{b\nu })_{\beta \lambda }𝒢_\nu ^b(x_2x_1,x_2;y_1,y_2)_{\alpha \lambda \rho \sigma }\end{array}$$
(26)
here
$$\mathrm{\Omega }_1^{a\mu }=g\gamma _1^0\gamma _1^\mu T_1^a,\text{ }\mathrm{\Omega }_2^{b\nu }=g\gamma _2^0\gamma _2^\nu \overline{T}_2^b$$
(27)
As will be proved in section 4, the Green’s functions $`𝒢_\mu ^a(x_ix_1,x_2;y_1,y_2)`$ are B-S reducible, therefore, we can write
$$𝒢^{(i)}(x_ix_1,x_2;y_1,y_2)=d^4z_1d^4z_2K^{(i)}(x_1,x_2;z_1,z_2)𝒢(z_1,z_2;y_1,y_2)$$
(28)
where $`K^{(i)}(x_1,x_2;z_1,z_2)`$ $`(i=1,2)`$ are just the interaction kernels. With the expression given in the above, Eqs. (2.19) and (2.20) can be represented as
| $`[i\frac{}{t_1}h^{(1)}(\stackrel{}{x_1})]_{\alpha \gamma }𝒢_{\gamma \beta \rho \sigma }(x_1,x_2;y_1,y_2)=(\gamma _1^0)_{\alpha \rho }\delta ^4(x_1y_1)S_F^c(x_2y_2)_{\beta \sigma }`$ |
| --- |
| $`+d^4z_1d^4z_2K^{(1)}(x_1,x_2;z_1,z_2)_{\alpha \beta \lambda \tau }𝒢_{\lambda \tau \rho \sigma }(z_1,z_2;y_1,y_2)`$ |
(29)
and
| $`[i\frac{}{t_2}h^{(2)}(\stackrel{}{x_2})]_{\beta \lambda }𝒢_{\alpha \lambda \rho \sigma }(x_1,x_2;y_1,y_2)=(\gamma _2^0)_{\beta \sigma }\delta ^4(x_2y_2)S_F(x_1y_1)_{\alpha \rho }`$ |
| --- |
| $`+d^4z_1d^4z_2K^{(2)}(x_1,x_2;z_1,z_2)_{\alpha \beta \lambda \tau }𝒢_{\lambda \tau \rho \sigma }(z_1,z_2;y_1,y_2)`$ |
(30)
From the above two equations, we may obtain two equivalent equations: one describes the evolution of the Green’s function $`𝒢(z_1,z_2;y_1,y_2)`$ with the center of mass time; another describes the evolution of the Green’s function with the relative time. The first equation is given by summing up the two equations in Eqs. (2.25) and (2.26). Introducing the cluster coordinates
| $`X=\eta _1x_1+\eta _2x_2,`$ $`x=x_1x_2;`$ |
| --- |
| $`Y=\eta _1y_1+\eta _2y_2,`$ $`y=y_1y_2,`$ |
| $`\eta _i=\frac{m_i}{m_1+m_2},i=1,2`$ |
(31)
the equation may be represented in the matrix notation as
| $`[i\frac{}{t}H_0(\stackrel{}{X},\stackrel{}{x})]𝒢(XY,x,y)`$ |
| --- |
| $`=S(XY,x,y)+d^4Zd^4zK(XZ,x,z)𝒢(ZY,z,y)`$ |
(32)
where
$$S(XY,x,y)=\delta ^4(x_1y_1)\gamma _1^0S_F^c(x_2y_2)+\delta ^4(x_2y_2)\gamma _2^0S_F(x_1y_1)$$
(33)
$$H_0(\stackrel{}{X},\stackrel{}{x})=h^{(1)}(\stackrel{}{x}_1)+h^{(2)}(\stackrel{}{x}_2)$$
(34)
is the total free Hamiltonian,
$$K(XZ,x,z)=\underset{i=1}{\overset{2}{}}K^{(i)}(XZ,x,z)$$
(35)
is the total interaction kernel and $`t=X_0`$ is the center of mass time. In the above, the translation-invariance of the Green’s function and the interaction kernel has been considered.
The second equation mentioned above is given by subtracting the equation in Eq. (2.26) with weight $`\eta _1`$ from the equation in Eq. (2.25) with weight $`\eta _2`$
| $`[i\frac{}{\tau }\overline{H}_0(\stackrel{}{X},\stackrel{}{x})]𝒢(XY,x,y)`$ |
| --- |
| $`=\overline{S}(XY,x,y)+d^4Zd^4z\overline{K}(XZ,x,z)𝒢(ZY,z,y)`$ |
(36)
where
$$\overline{S}(XY,x,y)=\eta _2\delta ^4(x_1y_1)\gamma _1^0S_F^c(x_2y_2)\eta _1\delta ^4(x_2y_2)\gamma _2^0S_F(x_1y_1)$$
(37)
$$\overline{H}_0(\stackrel{}{X},\stackrel{}{x})=\eta _2h^{(1)}(\stackrel{}{x}_1)\eta _1h^{(2)}(\stackrel{}{x}_2)$$
(38)
which is the relative Hamiltonian,
$$\overline{K}(XZ,x,z)=\eta _2K^{(1)}(XZ,x,z)\eta _1K^{(2)}(XZ,x,z)$$
(39)
which is the relative kernel and $`\tau =x_0`$ is the relative time.
By virtue of the well-known Lehmann representation of the Green’s function $`𝒢(z_1,z_2;y_1,y_2)`$ , one may derived the equations satisfied by the B-S amplitude from the above equations. The Lehmann representation as shown below can easily be written out by the procedure of inserting the complete set of the $`q\overline{q}`$ bound states into the Green’s function $`𝒢(z_1,z_2;y_1,y_2)`$ denoted in Eq. (2.1), then considering the translation-invariance property of the Green’s function and finally employing the integral representation of the step function,
| $`𝒢(XY,x,y)=\underset{n}{}\frac{i}{(2\pi )^4}d^4Q_ne^{iQ_n(XY)}`$ |
| --- |
| $`\times \frac{1}{2\omega _n}\{\frac{\chi _{Q_n}(x)\overline{\chi }_{Q_n}(y)}{Q_n^0\omega _n+iϵ}\frac{\chi _{Q_n}^+(y)\overline{\chi }_{Q_n}^+(x)}{Q_n^0+\omega _niϵ}\}`$ |
(40)
where
$$\begin{array}{c}\chi _{Q_n}(X,x)=e^{iQ_nX}\chi _{Q_n}(x)=0^+\left|N\{\psi (x_1)\psi ^c(x_2)\}\right|n\\ \overline{\chi }_{Q_n}(Y,y)=e^{iQ_nY}\overline{\chi }_{Q_n}(y)=n\left|N\{\overline{\psi }(y_1)\overline{\psi }^c(y_2)\}\right|0^{}\end{array}$$
(41)
are the B-S amplitudes describing the bound state and $`\omega _n`$ is the energy of the state $`n`$. Upon substituting Eq. (2.36) into Eqs. (2.28) and (2.32), then taking the limit: $`lim_{Q_n^0\omega _n}(Q_n^0\omega _n)`$ and finally performing the integration: $`d^4Ye^{iPY}`$, one can find that
$$[i\frac{}{t}H_0(\stackrel{}{X},\stackrel{}{x})]\chi _{P\varsigma }(X,x)=d^4Yd^4yK(XY,x,y)\chi _{P\varsigma }(Y,y)$$
(42)
and
$$[i\frac{}{\tau }\overline{H}_0(\stackrel{}{X},\stackrel{}{x})]\chi _{P\varsigma }(X,x)=d^4Yd^4y\overline{K}(XY,x,y)\chi _{P\varsigma }(Y,y)$$
(43)
where the subscript $`\varsigma `$ in the B-S amplitude designates the other quantum numbers of a bound state. In the above derivation, the fact that the functions $`S(XY,x,y)`$ and $`\overline{S}(XY,x,y)`$ have no the bound state poles has been considered. Eqs. (2.38) and (2.39) are just the wanted D-S equations satisfied by the B-S amplitudes. Eq. (2.38) describes the evolution of the $`q\overline{q}`$ bound state with the center of mass time $`t`$, while, Eq. (2.39) describes the evolution of the $`q\overline{q}`$ bound state with the relative time $`\tau .`$ Clearly, both of the equations are all the first-order differential equations whose solutions are only determined by the initial amplitudes at the origin of times.
Considering the translation-invariance of the B-S amplitude and the kernels as shown in Eq. (2.37) and in the following
$$\begin{array}{c}K(XY,x,y)=\frac{d^4Q}{(2\pi )^4}e^{iQ(XY)}K(Q,x,y)\\ \overline{K}(XY,x,y)=\frac{d^4Q}{(2\pi )^4}e^{iQ(XY)}\overline{K}(Q,x,y)\end{array}$$
(44)
one can obtain from Eqs. (2.38) and (2.39) the equations satisfied by the amplitude which describes the internal motion of the $`q\overline{q}`$ bound system
$$[EH_0(\stackrel{}{P},\stackrel{}{x})]\chi _{P\varsigma }(x)=d^4yK(P,x,y)\chi _{P\varsigma }(y)$$
(45)
$$[i\frac{}{\tau }\overline{H}_0(\stackrel{}{P},\stackrel{}{x})]\chi _{P\varsigma }(x)=d^4y\overline{K}(P,x,y)\chi _{P\varsigma }(y)$$
(46)
Furthermore, in view of the Fourier transformation
$$\begin{array}{c}\chi _{P\varsigma }(x)=\frac{d^4q}{(2\pi )^4}e^{iqx}\chi _{P\varsigma }(q)\\ K(P,x,y)=\frac{d^4q}{(2\pi )^4}\frac{d^4k}{(2\pi )^4}e^{iqx+iky}K(P,q,k),\end{array}$$
(47)
Eqs. (2.41) and (2.42) will immediately be transformed into the momentum space
$$[EH_0(\stackrel{}{P},\stackrel{}{q})]\chi _{P\varsigma }(q)=\frac{d^4k}{(2\pi )^4}K(P,q,k)\chi _{P\varsigma }(k)$$
(48)
$$[q_0\overline{H}_0(\stackrel{}{P},\stackrel{}{q})]\chi _{P\varsigma }(q)=\frac{d^4k}{(2\pi )^4}\overline{K}(P,q,k)\chi _{P\varsigma }(k)$$
(49)
where
$$\begin{array}{c}H_0(\stackrel{}{P},\stackrel{}{q})=h^{(1)}(\stackrel{}{p_1})+h^{(2)}(\stackrel{}{p_2})\\ \overline{H}_0(\stackrel{}{P},\stackrel{}{q})=\eta _2h^{(1)}(\stackrel{}{p}_1)\eta _1h^{(2)}(\stackrel{}{p}_2)\end{array}$$
(50)
in which
$$\begin{array}{c}h^{(i)}(\stackrel{}{p_i})=\stackrel{}{p_i}\stackrel{}{\alpha _i}+m_i\gamma _i^0,\text{ }i=1,2,\\ P=p_1+p_2,\text{ }q=\eta _2p_1\eta _1P_{2,}\text{ }k=\eta _2q_1\eta _1q_2\end{array}$$
(51)
here $`q_i`$ and $`p_{i\text{ }}`$ are the initial state and the final state momenta for i-th single particle. Both of the above equations are identified themselves with the eigenvalue equations for the $`q\overline{q}`$ bound system.
The D-S equations given in Eqs. (2.41) and (2.42) or (2.44) and (2.45) are Lorentz-covariant. They allow us to investigate the $`q\overline{q}`$ bound states in any coordinate frame. Based on the space-like property of a bound state, it is admissible to discuss the bound state in a special equal-time Lorentz frame. In this frame, the $`q\overline{q}`$ four-point Green’s function becomes
$$𝒢(\stackrel{}{x_1}_,\stackrel{}{x_2};\stackrel{}{y_1},\stackrel{}{y_2};t_1t_2)=0^+\left|T\{N[\psi (\stackrel{}{x_1},t_1)\psi ^c(\stackrel{}{x_2},t_1)]N[\overline{\psi }(\stackrel{}{y_1},t_2)\overline{\psi }^c(\stackrel{}{y_2},t_2)]\}\right|0^{}$$
(52)
The equation in Eq. (2.28) is now reduced to
| $`[i\frac{}{t_1}H_0(\stackrel{}{x_1}_,\stackrel{}{x_2})]𝒢(\stackrel{}{x_1}_,\stackrel{}{x_2};\stackrel{}{y_1},\stackrel{}{y_2};t_1t_2)`$ |
| --- |
| $`=\delta (t_1t_2)S(\stackrel{}{x_1}_,\stackrel{}{x_2};\stackrel{}{y_1},\stackrel{}{y_2})+d^3z_1d^3z_2𝑑t_zK(\stackrel{}{x_1}_,\stackrel{}{x_2};\stackrel{}{z_1},\stackrel{}{z_2};t_1t_z)𝒢(\stackrel{}{z_1}_,\stackrel{}{z_2};\stackrel{}{y_1},\stackrel{}{y_2};t_zt_2)`$ |
(53)
where
$$S(\stackrel{}{x_1}_,\stackrel{}{x_2};\stackrel{}{y_1},\stackrel{}{y_2})=\delta ^3(\stackrel{}{x_1}\stackrel{}{y_1})\gamma _1^0S_F^c(\stackrel{}{x_2}\stackrel{}{y_2})+\delta ^3(\stackrel{}{x_2}\stackrel{}{y_2})\gamma _2^0S_F(\stackrel{}{x_1}\stackrel{}{y_1})$$
(54)
in which
$$\begin{array}{c}S_F(\stackrel{}{x_1}\stackrel{}{y_1})=\frac{1}{i}0^+|T[\psi (\stackrel{}{x_1},t_1)[\overline{\psi }(\stackrel{}{y_1},t_1)]|0^{}\\ S_F^c(\stackrel{}{x_2}\stackrel{}{y_2})=0^+\left|T[\psi ^c(\stackrel{}{x_2},t_2)\overline{\psi }^c(\stackrel{}{y_2},t_2)]\right|0^{}\end{array}$$
(55)
are the equal-time quark and antiquark propagators respectively which are actually independent of time due to the translation-invariance property. By the Fourier transformations
$$\begin{array}{c}𝒢(\stackrel{}{x_1}_,\stackrel{}{x_2};\stackrel{}{y_1},\stackrel{}{y_2};t_1t_2)=_{\mathrm{}}^+\mathrm{}\frac{dE}{2\pi }e^{iE(t_1t_2)}𝒢(\stackrel{}{x_1}_,\stackrel{}{x_2};\stackrel{}{y_1},\stackrel{}{y_2};E)\\ K(\stackrel{}{x_1}_,\stackrel{}{x_2};\stackrel{}{z_1},\stackrel{}{z_2};t_1t_2)=_{\mathrm{}}^+\mathrm{}\frac{dE}{2\pi }e^{iE(t_1t_2)}K(\stackrel{}{x_1}_,\stackrel{}{x_2};\stackrel{}{z_1},\stackrel{}{z_2};E)\end{array}$$
(56)
Eq. (2.49) will be represented as
| $`[EH_0(\stackrel{}{x_1}_,\stackrel{}{x_2})]𝒢(\stackrel{}{x_1}_,\stackrel{}{x_2};\stackrel{}{y_1},\stackrel{}{y_2};E)`$ |
| --- |
| $`=S(\stackrel{}{x_1}_,\stackrel{}{x_2};\stackrel{}{y_1},\stackrel{}{y_2})+d^3z_1d^3z_2K(\stackrel{}{x_1}_,\stackrel{}{x_2};\stackrel{}{z_1},\stackrel{}{z_2};E)𝒢(\stackrel{}{z_1}_,\stackrel{}{z_2};\stackrel{}{y_1},\stackrel{}{y_2};E)`$ |
(57)
This just is the three-dimensional equation satisfied by the $`q\overline{q}`$ four-point Green’s function defined in the equal-time Lorentz frame.
In the equal-time frame, the relative time of the $`q\overline{q}`$ system is zero. Therefore, the equation in Eq. (2.32) becomes meaningless. We are left with only the equation given in Eq. (2.53). The Lehmann representation of the Green’s function $`𝒢(\stackrel{}{x_1}_,\stackrel{}{x_2};\stackrel{}{y_1},\stackrel{}{y_2};t)`$ is still represented in Eq. (2.36) except that the four-dimensional relative coordinates $`x`$ and $`y`$ in the B-S amplitudes are now replaced by the three-dimensional ones $`\stackrel{}{x}`$ and $`\stackrel{}{y}`$. Substituting such a Lehmann representation into Eq. (2.53), by the same procedure as stated in Eqs. (2.36)-(2.38), one may obtain a D-S equation represented in the three-dimensional coordinate space such that
$$[EH_0(\stackrel{}{X},\stackrel{}{x})]\chi _{P\varsigma }(\stackrel{}{X},\stackrel{}{x})=d^3Yd^3yK(\stackrel{}{X}\stackrel{}{Y},\stackrel{}{x},\stackrel{}{y})\chi _{P\varsigma }(\stackrel{}{Y},\stackrel{}{y})$$
(58)
where
$$\chi _{P\varsigma }(\stackrel{}{X},\stackrel{}{x})=e^{i\stackrel{}{P}\stackrel{}{X}}\chi _{P\varsigma }(\stackrel{}{x})$$
(59)
Apparently, in the momentum space, Eq. (2.54) becomes
$$[EH_0(\stackrel{}{P},\stackrel{}{q})]\chi _{P\varsigma }(\stackrel{}{q})=\frac{d^3k}{(2\pi )^3}K(\stackrel{}{P},\stackrel{}{q},\stackrel{}{k};E)\chi _{P\varsigma }(\stackrel{}{k})$$
(60)
This is precisely the three-dimensional D-S equation satisfied by the amplitude $`\chi _{P\varsigma }(\stackrel{}{k})`$ which describes the internal motion of the $`q\overline{q}`$ bound system and may be used to solve the eigenvalue problem for the system. This equation is rigorous because the retardation effect is completely included in the kernel of the equation.
## III Derivation of the interaction kernel
In this section, the interaction kernel in the D-S equation will be derived by making use of equations of motion which describe the variation of the Green’s functions with coordinates $`y_1`$ and $`y_2`$. The equations satisfied by the Green’s function $`G(x_{1,}x_2;y_1,y_2)`$ are derived in Appendix A and can directly be written out from Eqs. (A.22) and (A.31) by setting the source $`J`$ to be zero. The result is
| $`G(x_{1,}x_2;y_1,y_2)_{\alpha \beta \tau \sigma }(i\stackrel{}{}_{y_1}+m_1)_{\tau \rho }=\delta _{\alpha \rho }\delta ^4(x_1y_1)S_F^c(x_2y_2)_{\beta \sigma }`$ |
| --- |
| $`C_{\rho \sigma }\delta ^4(y_1y_2)S_F^{}(x_1x_2)_{\alpha \beta }+G_\mu ^a(y_1x_1,x_2;y_1,y_2)_{\alpha \beta \tau \sigma }(\mathrm{\Gamma }^{a\mu })_{\tau \rho }`$ |
(61)
| $`G(x_{1,}x_2;y_1,y_2)_{\alpha \beta \rho \delta }(i\stackrel{}{}_{y_2}+m_2)_{\delta \sigma }=\delta _{\beta \sigma }\delta ^4(x_2y_2)S_F(x_1y_1)_{\alpha \rho }`$ |
| --- |
| $`C_{\rho \sigma }\delta ^4(y_1y_2)S_F^{}(x_1x_2)_{\alpha \beta }+G_\nu ^b(y_2x_1,x_2;y_1,y_2)_{\alpha \beta \rho \delta }(\overline{\mathrm{\Gamma }}^{b\nu })_{\delta \sigma }`$ |
(62)
where $`G_\mu ^a(y_ix_1,x_2;y_1,y_2)`$ $`(`$ $`i=1,2)`$ were defined in Eq. (2.13) with the replacement of $`x_i`$ by $`y_i`$.
Substituting the relation in Eq. (2.4) into the above two equations and employing the following equations obeyed by the propagator $`\overline{S}_F^{}(y_1y_2)`$ whose derivation was mentioned in Appendix A
$$\begin{array}{c}\overline{S}_F^{}(y_1y_2)_{\tau \sigma }(i\stackrel{}{}_{y_1}+m_1)_{\tau \rho }=C_{\rho \sigma }\delta ^4(y_1y_2)+\overline{\mathrm{\Lambda }}_\nu ^b(y_1y_1,y_2)_{\tau \sigma }(\mathrm{\Gamma }^{b\nu })_{\tau \rho }\\ \overline{S}_F^{}(y_1y_2)_{\rho \delta }(i\stackrel{}{}_{y_2}+m_2)_{\delta \sigma }=C_{\rho \sigma }\delta ^4(y_1y_2)+\overline{\mathrm{\Lambda }}_\nu ^b(y_2y_1,y_2)_{\rho \delta }(\overline{\mathrm{\Gamma }}^{b\nu })_{\delta \sigma }\end{array}$$
(63)
where
$$\overline{\mathrm{\Lambda }}_\nu ^b(x_iy_1,y_2)_{\tau \sigma }=\frac{1}{i}0^+\left|T\{𝐀_\nu ^b(x_i)\overline{\psi }_\tau (y_1)\overline{\psi }_\sigma ^c(y_2)\}\right|0^{},$$
(64)
it is easy to find
| $`𝒢(x_{1,}x_2;y_1,y_2)_{\alpha \beta \tau \sigma }(i\stackrel{}{}_{y_1}+m_1)_{\tau \rho }=\delta _{\alpha \rho }\delta ^4(x_1y_1)S_F^c(x_2y_2)_{\beta \sigma }`$ |
| --- |
| $`+𝒢_\mu ^a(y_1x_1,x_2;y_1,y_2)_{\alpha \beta \tau \sigma }(\mathrm{\Gamma }^{a\mu })_{\tau \rho }`$ |
(65)
and
| $`𝒢(x_{1,}x_2;y_1,y_2)_{\alpha \beta \rho \delta }(i\stackrel{}{}_{y_2}+m_2)_{\delta \sigma }=\delta _{\beta \sigma }\delta ^4(x_2y_2)S_F(x_1y_1)_{\alpha \rho }`$ |
| --- |
| $`+𝒢_\nu ^b(y_2x_1,x_2;y_1,y_2)_{\alpha \beta \rho \delta }(\overline{\mathrm{\Gamma }}^{b\nu })_{\delta \sigma }`$ |
(66)
where
| $`𝒢_\mu ^a(y_ix_1,x_2;y_1,y_2)_{\alpha \beta \rho \sigma }=0^+\left|T\{N[\psi _\alpha (x_1)\psi _\beta ^c(x_2)]N[𝐀_\mu ^a(y_i)\overline{\psi }_\tau (y_1)\overline{\psi }_\sigma ^c(y_2)]\}\right|0^{}`$ |
| --- |
| $`=G_\mu ^a(y_ix_1,x_2;y_1,y_2)_{\alpha \beta \rho \sigma }+S_F^{}(x_1x_2)_{\alpha \beta }\overline{\mathrm{\Lambda }}_\mu ^a(y_iy_1,y_2)_{\rho \sigma }`$ |
(67)
here $`i=1,2`$.
To derive the interaction kernel, we also need equations obeyed by the Green’s functions $`𝒢_\mu ^a(x_ix_1,x_2;y_1,y_2)`$. According to the description given in Appendix A, the equations for the Green’s functions $`G_\mu ^a(x_ix_1,x_2;y_1,y_2)`$ can be derived by differentiating Eqs. (A.22) and (A.31) with respect to the source $`J^{a\mu }(x_i)`$ and then setting $`J=0`$. The result is
| $`G_\mu ^a(x_ix_{1,}x_2;y_1,y_2)_{\alpha \beta \tau \sigma }(i\stackrel{}{}_{y_1}+m_1)_{\tau \rho }=\delta _{\alpha \rho }\delta ^4(x_1y_1)\mathrm{\Lambda }_\mu ^{𝐜a}(x_ix_2,y_2)_{\beta \sigma }`$ |
| --- |
| $`C_{\rho \sigma }\delta ^4(y_1y_2)\mathrm{\Lambda }_\mu ^a(x_ix_1,x_2)_{\alpha \beta }+G_{\mu \nu }^{ab}(x_i,y_1x_1,x_2;y_1,y_2)_{\alpha \beta \tau \sigma }(\mathrm{\Gamma }^{a\mu })_{\tau \rho }`$ |
(68)
and
| $`G_\mu ^a(x_ix_{1,}x_2;y_1,y_2)_{\alpha \beta \rho \delta }(i\stackrel{}{}_{y_2}+m_2)_{\delta \sigma }=\delta _{\beta \sigma }\delta ^4(x_2y_2)\mathrm{\Lambda }_\mu ^a(x_ix_1,y_1)_{\alpha \rho }`$ |
| --- |
| $`C_{\rho \sigma }\delta ^4(y_1y_2)\mathrm{\Lambda }_\mu ^a(x_ix_1,x_2)_{\alpha \beta }+G_{\mu \nu }^{ab}(x_i,y_2x_1,x_2;y_1,y_2)_{\alpha \beta \rho \delta }(\overline{\mathrm{\Gamma }}^{b\nu })_{\delta \sigma }`$ |
(69)
where
$$\begin{array}{c}\mathrm{\Lambda }_\mu ^a(x_ix_1,y_1)_{\gamma \rho }=\frac{1}{i}0^+\left|T[𝐀_\mu ^a(x_i)\psi _\gamma (x_1)\overline{\psi }_\rho (y_1)]\right|0^{},\\ \mathrm{\Lambda }_\mu ^{𝐜a}(x_ix_2,y_2)_{\lambda \sigma }=\frac{1}{i}0^+\left|T[𝐀_\mu ^a(x_i)\psi _\lambda ^c(x_2)\overline{\psi }_\sigma ^c(y_2)]\right|0^{},\\ \mathrm{\Lambda }_\mu ^a(x_iy_1,y_2)_{\tau \sigma }=\frac{1}{i}0^+\left|T[𝐀_\mu ^a(x_i)\psi _\tau (y_1)\psi _\sigma ^c(y_2)]\right|0^{}\end{array}$$
(70)
and
$$G_{\mu \nu }^{ab}(x_i,y_jx_1,x_2;y_1,y_2)_{\alpha \beta \rho \sigma }=0^+\left|T\{𝐀_\mu ^a(x_i)𝐀_\nu ^b(y_j)\psi _\alpha (x_1)\psi _\beta ^c(x_2)\overline{\psi }_\rho (y_1)\overline{\psi }_\sigma ^c(y_2)\}\right|0^{}$$
(71)
here $`i,j=1,2.`$ On inserting the relation in Eq. (2.18) into Eqs. (3.8) and (3.9) and utilizing the equations in Eq. (3.3), we are led to
| $`𝒢_\mu ^a(x_ix_{1,}x_2;y_1,y_2)_{\alpha \beta \tau \sigma }(i\stackrel{}{}_{y_1}+m_1)_{\tau \rho }=\delta _{\alpha \rho }\delta ^4(x_1y_1)\mathrm{\Lambda }_\mu ^{𝐜a}(x_ix_2,y_2)_{\beta \sigma }`$ |
| --- |
| $`+𝒢_{\mu \nu }^{ab}(x_i,y_1x_1,x_2;y_1,y_2)_{\alpha \beta \tau \sigma }(\mathrm{\Gamma }^{a\mu })_{\tau \rho }`$ |
(72)
and
| $`𝒢_\mu ^a(x_ix_{1,}x_2;y_1,y_2)_{\alpha \beta \rho \delta }(i\stackrel{}{}_{y_2}+m_2)_{\delta \sigma }=\delta _{\beta \sigma }\delta ^4(x_2y_2)\mathrm{\Lambda }_\mu ^a(x_ix_1,y_1)_{\alpha \rho }`$ |
| --- |
| $`+𝒢_{\mu \nu }^{ab}(x_i,y_2x_1,x_2;y_1,y_2)_{\alpha \beta \rho \delta }(\overline{\mathrm{\Gamma }}^{b\nu })_{\delta \sigma }`$ |
(73)
where
| $`𝒢_{\mu \nu }^{ab}(x_i,y_jx_1,x_2;y_1,y_2)_{\alpha \beta \rho \sigma }`$ |
| --- |
| =$`0^+\left|T\{N[𝐀_\mu ^a(x_i)\psi _\alpha (x_1)\psi _\beta ^c(x_2)]N[𝐀_\nu ^b(y_j)\overline{\psi }_\rho (y_1)\overline{\psi }_\sigma ^c(y_2)]\}\right|0^{}`$ |
| $`=G_{\mu \nu }^{ab}(x_i,y_jx_1,x_2;y_1,y_2)_{\alpha \beta \rho \sigma }+\mathrm{\Lambda }_\mu ^a(x_ix_1,x_2)_{\alpha \beta }\overline{\mathrm{\Lambda }}_\nu ^b(y_jy_1,y_2)_{\rho \sigma }`$ |
(74)
Now, let us multiply Eqs. (3.5) and (3.6) respectively with $`\gamma _1^0`$ and $`\gamma _2^0`$ from the right and sum up the both equations thus obtained. In this way, writing in the matrix from, we obtain the following equation
| $`𝒢(x_1,x_2;y_1,y_2)[i\stackrel{}{\frac{}{t^{}}}+\stackrel{}{H_0}(y_1,y_2)]`$ |
| --- |
| $`=S(x_1,x_2;y_1,y_2)\underset{i=1}{\overset{2}{}}𝒢^{(i)}(y_ix_1,x_2;y_1,y_2)`$ |
(75)
where
$$\frac{}{t^{}}=\frac{}{y_1^0}+\frac{}{y_2^0}=\frac{}{Y^0},$$
(76)
$$H_0(y_1,y_2)=h^{(1)}(y_1)+h^{(2)}(y_2)$$
(77)
here $`h^{(1)}(y_1)`$ and $`h^{(2)}(y_2)`$ were represented in Eq. (2.21),
$$\begin{array}{c}𝒢^{(1)}(y_1x_1,x_2;y_1,y_2)_{\alpha \beta \rho \sigma }=𝒢_\mu ^a(y_1x_1,x_2;y_1,y_2)_{\alpha \beta \tau \sigma }(\overline{\mathrm{\Omega }}_1^{a\mu })_{\tau \rho }\\ 𝒢^{(2)}(y_2x_1,x_2;y_1,y_2)_{\alpha \beta \rho \sigma }=𝒢_\mu ^a(y_2x_1,x_2;y_1,y_2)_{\alpha \beta \rho \delta }(\overline{\mathrm{\Omega }}_2^{a\mu })_{\delta \sigma }\end{array}$$
(78)
in which
$$\overline{\mathrm{\Omega }}_1^{a\mu }=g\gamma _1^\mu \gamma _1^0T^a=g\gamma _1^0\gamma _1^{\mu +}T^a,\text{ }\overline{\mathrm{\Omega }}_2^{a\mu }=g\gamma _2^\mu \gamma _2^0\overline{T}^a=g\gamma _2^0\gamma _2^{\mu +}\overline{T}^a$$
(79)
and $`S(x_1,x_2;y_1,y_2)`$ was defined in Eq. (2.29).
Similarly, on multiplying Eqs. (3.12) and (3.13) respectively with $`\gamma _1^0`$ and $`\gamma _2^0`$ from the right and summing up the both equations thus obtained, one can get
| $`𝒢_\mu ^a(x_ix_1,x_2;y_1,y_2)_{\alpha \beta \lambda \tau }[i\stackrel{}{\frac{}{t^{}}}+\stackrel{}{H_0}(y_1,y_2)]_{\lambda \tau \rho \sigma }`$ |
| --- |
| $`=_\mu ^{(i)a}(x_1,x_2;y_1,y_2)_{\alpha \beta \rho \sigma }𝒬_\mu ^{(i)a}(x_1,x_2;y_1,y_2)_{\alpha \beta \rho \sigma }`$ |
(80)
where
| $`_\mu ^{(i)a}(x_1,x_2;y_1,y_2)_{\alpha \beta \rho \sigma }=\delta ^4(x_1y_1)(\gamma _1^0)_{\alpha \rho }\mathrm{\Lambda }_\mu ^{𝐜a}(x_ix_2,y_2)_{\beta \sigma }`$ |
| --- |
| $`+\delta ^4(x_2y_2)(\gamma _2^0)_{\beta \sigma }\mathrm{\Lambda }_\mu ^a(x_ix_1,y_1)_{\alpha \rho }`$ |
(81)
and
| $`𝒬_\mu ^{(i)a}(x_1,x_2;y_1,y_2)_{\alpha \beta \rho \sigma }=𝒢_{\mu \nu }^{ab}(x_i,y_1x_1,x_2;y_1,y_2)_{\alpha \beta \tau \sigma }(\overline{\mathrm{\Omega }}_1^{b\nu })_{\tau \rho }`$ |
| --- |
| $`+𝒢_{\mu \nu }^{ab}(x_i,y_2x_1,x_2;y_1,y_2)_{\alpha \beta \rho \delta }(\overline{\mathrm{\Omega }}_2^{b\nu })_{\delta \sigma }`$ |
(82)
With the definitions given in Eq. (2.22) and in the following
$$^{(i)}(x_1,x_2;y_1,y_2)=\mathrm{\Omega }_i^{a\mu }_\mu ^{(i)a}(x_1,x_2;y_1,y_2)$$
(83)
$$𝒬^{(i)}(x_1,x_2;y_1,y_2)=\mathrm{\Omega }_i^{a\mu }𝒬_\mu ^{(i)a}(x_1,x_2;y_1,y_2)$$
(84)
we can write from Eq. (3.20) the equation satisfied by the function $`𝒢_{\alpha \beta \rho \sigma }^{(i)}(x_ix_1,x_2;y_1,y_2)`$. In the matrix form, it reads
| $`𝒢^{(i)}(x_ix_1,x_2;y_1,y_2)[i\stackrel{}{\frac{}{t^{}}}+\stackrel{}{H_0}(y_1,y_2)]`$ |
| --- |
| $`=^{(i)}(x_1,x_2;y_1,y_2)𝒬^{(i)}(x_1,x_2;y_1,y_2)`$ |
(85)
Up to the present, we are ready to derive the interaction kernel. Acting on the both sides of Eq. (2.24) with the operator $`i\stackrel{}{\frac{}{t^{}}}+\stackrel{}{H_0}(y_1,y_2)`$ from the right and employing Eqs. (3.15) and (3.25), we have
| $`d^4z_1d^4z_2K^{(i)}(x_1,x_2;z_1,z_2)S(z_1,z_2;y_1,y_2)=^{(i)}(x_1,x_2;y_1,y_2)+𝒬^{(i)}(x_1,x_2;y_1,y_2)`$ |
| --- |
| $`d^4z_1d^4z_2K^{(i)}(x_1,x_2;z_1,z_2)\underset{j=1}{\overset{2}{}}𝒢^{(j)}(y_jz_1,z_2;y_1,y_2)`$ |
(86)
In order to obtain the kernel, we may operate on the above equation with the inverse of $`S(x_1,x_2;y_1,y_2)`$ and the kernel on the RHS of Eq. (3.26) may be eliminated by the following relation given by acting on Eq. (2.24) with the inverse of the Green’s function $`𝒢(x_1,x_2;y_1,y_2)`$
$$K^{(i)}(x_1,x_2;y_1,y_2)=d^4u_1d^4u_2𝒢^{(i)}(x_ix_1,x_2;u_1,u_2)𝒢^1(u_1,u_2;y_1,y_2)$$
(87)
With these operations, one may derive from Eq.(3.26) a closed expression of the kernel $`K^{(i)}(x_1,x_2;y_1,y_2)`$ such that
| $`K^{(i)}(x_1,x_2;y_1,y_2)=d^4z_1d^4z_2\{^{(i)}(x_1,x_2;z_1,z_2)+𝒬^{(i)}(x_1,x_2;z_1,z_2)`$ |
| --- |
| $`𝒟^{(i)}(x_1,x_2;z_1,z_2)\}S^1(z_1,z_2;y_1,y_2)`$ |
(88)
where
| $`𝒟^{(i)}(x_1,x_2;z_1,z_2)`$ |
| --- |
| $`=\underset{k=1}{\overset{2}{}}d^4u_kd^4v_k𝒢^{(i)}(x_ix_1,x_2;u_1,u_2)𝒢^1(u_1,u_2;v_1,v_2)\underset{j=1}{\overset{2}{}}𝒢^{(j)}(y_jv_1,v_2;z_1,z_2)`$ |
(89)
In the above derivation, existence of the inverses $`𝒢^1(u_1,u_2;z_1,z_2)`$ and $`S^1(z_1,z_2;y_1,y_2)`$ has been assumed. The rationality of the assumption will be illustrated later.
Clearly, the total interaction kernel appearing in Eq. (2.38) is given by the sum
| $`K(x_1,x_2;y_1,y_2)=\underset{j=1}{\overset{2}{}}K^{(i)}(x_1,x_2;y_1,y_2)=d^4z_1d^4z_2\{(x_1,x_2;z_1,z_2)`$ |
| --- |
| $`+𝒬(x_1,x_2;z_1,z_2)𝒟(x_1,x_2;z_1,z_2)\}S^1(z_1,z_2;y_1,y_2)`$ |
(90)
where
$$(x_1,x_2;z_1,z_2)=\underset{i=1}{\overset{2}{}}^{(i)}(x_1,x_2;z_1,z_2)$$
(91)
$$𝒬(x_1,x_2;z_1,z_2)=\underset{i=1}{\overset{2}{}}𝒬^{(i)}(x_1,x_2;z_1,z_2)=\underset{i,j=1}{\overset{2}{}}\mathrm{\Omega }_i^{a\mu }𝒢_{\mu \nu }^{ab}(x_i,z_jx_1,x_2;z_1,z_2)\overline{\mathrm{\Omega }}_j^{b\nu }$$
(92)
and
| $`𝒟(x_1,x_2;z_1,z_2)`$ |
| --- |
| $`=\underset{k=1}{\overset{2}{}}d^4u_kd^4v_k\underset{i,j=1}{\overset{2}{}}𝒢^{(i)}(x_ix_1,x_2;u_1,u_2)𝒢^1(u_1,u_2;v_1,v_2)𝒢^{(j)}(y_jv_1,v_2;z_1,z_2)`$ |
(93)
We would like here to discuss the role played by the last term in Eq. (3.30). In view of the relation in Eq. (2.24) and the following relation
$$𝒢^{(j)}(y_jx_1,x_2;y_1,y_2)=d^4z_1d^4z_2𝒢(x_1,x_2;z_1,z_2)K^{(j)}(z_1,z_2;y_1,y_2)$$
(94)
which also follows from the B-S reducibility of the Green’s function $`𝒢^{(j)}(y_jx_1,x_2;y_1,y_2)`$ and considering
$$d^4z_1d^4z_2𝒢(x_1,x_2;z_1,z_2)𝒢^1(z_1,z_2;y_1,y_2)=\delta ^4(x_1y_1)\delta ^4(x_2y_2)$$
(95)
the function $`𝒟(x_1,x_2;z_1,z_2)`$ in Eq. (3.33) may be represented as
$$𝒟(x_1,x_2;z_1,z_2)=\underset{k=1}{\overset{2}{}}d^4u_kd^4v_kK(x_1,x_2;u_1,u_2)𝒢(u_1,u_2;v_1,v_2)K(v_1,v_2;z_1,z_2)$$
(96)
which manifests itself the typical structure of the B-S reducibility of the kernel. Therefore, the last term in Eq. (3.30) just plays the role of cancelling the B-S reducible part of the remaining terms in Eq. (3.30). As a result of the cancellation, the interaction kernel given in Eq. (3.30) is really B-S irreducible, consistent with the ordinary concept for the kernel in a two-body relativistic equation. Inserting Eq. (3.36) into Eq. (3.30), writing in the operator form, we have
$$KS=+𝒬K𝒢K$$
(97)
This can be regarded as the integral equation satisfied by the kernel $`K.`$
Analogously, In accord with the definition in Eq. (2.35) and the expression in Eq. (3.28), one may write out an explicit expression of the kernel occurring in Eq. (2.39)
| $`\overline{K}(x_1,x_2;y_1,y_2)=\eta _2K^{(1)}(x_1,x_2;y_1,y_2)\eta _1K^{(2)}(x_1,x_2;y_1,y_2)`$ |
| --- |
| $`=d^4z_1d^4z_2\{\overline{}(x_1,x_2;z_1,z_2)+\overline{𝒬}(x_1,x_2;z_1,z_2)`$ |
| $`\overline{𝒟}(x_1,x_2;z_1,z_2)\}S^1(z_1,z_2;y_1,y_2)`$ |
(98)
in which
$$\overline{𝒜}(x_1,x_2;z_1,z_2)=\eta _2𝒜^{(1)}(x_1,x_2;z_1,z_2)\eta _1𝒜^{(2)}(x_1,x_2;z_1,z_2)$$
(99)
where $`𝒜`$ stands for $`,𝒬`$ or $`𝒟`$.
Now, we are in a position to write out the interaction kernel appearing in the three-dimensional D-S equation. As demonstrated in Ref. , the kernel in Eq. (2.54) can be derived by the same procedure as for the four-dimensional kernel. The expression of the three-dimensional kernel shown below formally is the same as the four-dimensional one described in Eq. (3.30) except that it is now represented in the three-dimensional space.
| $`K(\stackrel{}{x_1}_,\stackrel{}{x_2};\stackrel{}{y_1},\stackrel{}{y_2};E)=d^3z_1d^3z_2\{(\stackrel{}{x_1}_,\stackrel{}{x_2};\stackrel{}{z_1},\stackrel{}{z_2})`$ |
| --- |
| $`+𝒬(\stackrel{}{x_1}_,\stackrel{}{x_2};\stackrel{}{z_1},\stackrel{}{z_2};E)𝒟(\stackrel{}{x_1}_,\stackrel{}{x_2};\stackrel{}{z_1},\stackrel{}{z_2};E)\}S^1(\stackrel{}{z_1},\stackrel{}{z_2};\stackrel{}{y}_1,\stackrel{}{y}_2)`$ |
(100)
where $`(\stackrel{}{x_1}_,\stackrel{}{x_2};\stackrel{}{z_1},\stackrel{}{z_2}),𝒬(\stackrel{}{x_1}_,\stackrel{}{x_2};\stackrel{}{z_1},\stackrel{}{z_2};E)`$ and $`𝒟(\stackrel{}{x_1}_,\stackrel{}{x_2};\stackrel{}{z_1},\stackrel{}{z_2};E)`$ can be written out from Eqs. (3.31)-(3.33) as follows:
$$(\stackrel{}{x_1}_,\stackrel{}{x_2};\stackrel{}{z_1},\stackrel{}{z_2})=\underset{i=1}{\overset{2}{}}\mathrm{\Omega }_i^{a\mu }_\mu ^{(i)a}(\stackrel{}{x_1}_,\stackrel{}{x_2};\stackrel{}{z_1},\stackrel{}{z_2})$$
(101)
in which
$$_\mu ^{(i)a}(\stackrel{}{x_1}_,\stackrel{}{x_2};\stackrel{}{z_1},\stackrel{}{z_2})=\delta ^3(\stackrel{}{x_1}\stackrel{}{z_1})\gamma _1^0\mathrm{\Lambda }_\mu ^{𝐜a}(\stackrel{}{x_i}\stackrel{}{x_2},\stackrel{}{z_2})+\delta ^3(\stackrel{}{x_2}\stackrel{}{z_2})\gamma _2^0\mathrm{\Lambda }_\mu ^a(\stackrel{}{x_i}\stackrel{}{x_1},\stackrel{}{z_1})$$
(102)
here $`\mathrm{\Lambda }_\mu ^a(\stackrel{}{x_i}\stackrel{}{x_1},\stackrel{}{y_1})`$ and $`\mathrm{\Lambda }_\mu ^{𝐜a}(\stackrel{}{x_i}\stackrel{}{x_2},\stackrel{}{y_2})`$ are defined as in Eq. (3.10) except that the time variables in all the field operators are now taken to be the same and therefore they are time-independent due to the translation-invariance property of the Green’s functions,
$$𝒬(\stackrel{}{x_1}_,\stackrel{}{x_2};\stackrel{}{z_1},\stackrel{}{z_2};E)=\underset{i,j=1}{\overset{2}{}}\mathrm{\Omega }_i^{a\mu }𝒢_{\mu \nu }^{ab}(\stackrel{}{x_i},\stackrel{}{z_j}\stackrel{}{x_1}_,\stackrel{}{x_2};\stackrel{}{z_1},\stackrel{}{z_2};E)\overline{\mathrm{\Omega }}_j^{b\nu }$$
(103)
in which $`𝒢_{\mu \nu }^{ab}(\stackrel{}{x_i},\stackrel{}{z_j}\stackrel{}{x_1}_,\stackrel{}{x_2};\stackrel{}{z}_1,\stackrel{}{z}_2;E)`$ is the Fourier transform of the Green’s function defined by
| $`𝒢_{\mu \nu }^{ab}(\stackrel{}{x_i},\stackrel{}{z_j}\stackrel{}{x_1}_,\stackrel{}{x_2};\stackrel{}{z}_1,\stackrel{}{z}_2;t_1t_2)`$ |
| --- |
| $`=0^+\left|T\{N[𝐀_\mu ^a(\stackrel{}{x_i},t_1)\psi (\stackrel{}{x_1},t_1)\psi ^c(\stackrel{}{x_2},t_1)]N[𝐀_\nu ^b(\stackrel{}{z_j},t_2)\overline{\psi }(\stackrel{}{z_1},t_2)\overline{\psi }^c(\stackrel{}{z_2},t_2)]\}\right|0^{}`$ |
(104)
and
| $`𝒟(\stackrel{}{x_1}_,\stackrel{}{x_2};\stackrel{}{z_1},\stackrel{}{z_2};E)=\underset{k=1}{\overset{2}{}}d^3u_kd^3v_k\underset{i,j=1}{\overset{2}{}}\mathrm{\Omega }_i^{a\mu }𝒢_\mu ^{(i)a}(\stackrel{}{x_i}\stackrel{}{x_1}_,\stackrel{}{x_2};\stackrel{}{u_1},\stackrel{}{u_2};E)`$ |
| --- |
| $`\times 𝒢^1(\stackrel{}{u_1}_,\stackrel{}{u_2};\stackrel{}{v}_1,\stackrel{}{v}_2;E)𝒢_\nu ^{(j)b}(\stackrel{}{z_j}\stackrel{}{v}_1,\stackrel{}{v}_2;\stackrel{}{z_1},\stackrel{}{z_2};E)\overline{\mathrm{\Omega }}_j^{b\nu }`$ |
(105)
in which $`𝒢_\mu ^{(i)a}(\stackrel{}{x_i}\stackrel{}{x_1}_,\stackrel{}{x_2};\stackrel{}{u_1},\stackrel{}{u_2};E)`$ and $`𝒢_\nu ^{(j)b}(\stackrel{}{z_j}\stackrel{}{v}_1,\stackrel{}{v}_2;\stackrel{}{z_1},\stackrel{}{z_2};E)`$ are the Fourier transforms of the following Green’s functions
| $`𝒢_\mu ^{(i)a}(\stackrel{}{x_i}\stackrel{}{x_1}_,\stackrel{}{x_2};\stackrel{}{u_1},\stackrel{}{u_2};t_1t_2)`$ |
| --- |
| $`=0^+\left|T\{N[𝐀_\mu ^a(\stackrel{}{x_i},t_1)\psi (\stackrel{}{x_1},t_1)\psi ^c(\stackrel{}{x_2},t_1)]N[\overline{\psi }(\stackrel{}{u_1},t_2)\overline{\psi }^c(\stackrel{}{u_2},t_2)]\}\right|0^{}`$ |
(106)
and
| $`𝒢_\nu ^{(j)b}(\stackrel{}{z_j}\stackrel{}{v}_1,\stackrel{}{v}_2;\stackrel{}{z_1},\stackrel{}{z_2};t_1t_2)`$ |
| --- |
| $`=0^+\left|T\{N[\psi (\stackrel{}{v_1},t_1)\psi ^c(\stackrel{}{v_2},t_1)]N[𝐀_\nu ^b(\stackrel{}{z_j},t_2)\overline{\psi }(\stackrel{}{z_1},t_2)\overline{\psi }^c(\stackrel{}{z_2},t_2)]\}\right|0^{}`$ |
(107)
and $`S^1(\stackrel{}{z_1}_,\stackrel{}{z_2};\stackrel{}{y_1},\stackrel{}{y_2})`$ is the inverse of the function in Eq. (2.50).
## IV Another derivation of the interaction kernel
The aim of this section is to give a different expression of the interaction kernel in Eqs. (2.38) and (2.39) which will be derived by means of the irreducible decomposition of the Green’s functions $`𝒢_\mu ^a(x_ix_1,x_2;y_1,y_2)`$ denoted in Eq. (2.18). First, we start from the relation between the full $`q\overline{q}`$ four-point Green’s function $`G(x_{1,}x_2;y_1,y_2)`$ and its connected one $`G_c(x_{1,}x_2;y_1,y_2)`$ which is derived in the beginning of Appendix B
| $`G(x_{1,}x_2;y_1,y_2)=G_c(x_{1,}x_2;y_1,y_2)+S_F(x_1y_1)S_F^c(x_2y_2)`$ |
| --- |
| $`S_F^{}(x_1x_2)\overline{S}_F^{}(y_1y_2)`$ |
(108)
where all the fermion propagators were defined in Eqs. (2.6)-(2.9). Substituting Eq. (4.1) into Eq. (2.4), we get
$$𝒢(x_{1,}x_2;y_1,y_2)=G_c(x_{1,}x_2;y_1,y_2)+S_F(x_1y_1)S_F^c(x_2y_2)$$
(109)
where the last term is the unconnected part of the function $`𝒢(x_{1,}x_2;y_1,y_2).`$
From Eq. (B.5) given in Appendix B, we obtain by setting the source $`J=0`$ that
| $`G_\mu ^a(x_ix_{1,}x_2;y_1,y_2)=G_{c\mu }^a(x_ix_{1,}x_2;y_1,y_2)+\mathrm{\Lambda }_\mu ^a(x_ix_1;y_1)S_F^c(x_1y_1)`$ |
| --- |
| $`+S_F(x_1y_1)\mathrm{\Lambda }_\mu ^{𝐜a}(x_ix_2;y_2)\mathrm{\Lambda }_\mu ^a(x_ix_1,x_2)\overline{S}_F^{}(y_1y_2)S_F^{}(x_1x_2)\overline{\mathrm{\Lambda }}_\mu ^a(x_iy_1,y_2)`$ |
(110)
where $`i=1,2`$, $`G_{c\mu }^a(x_ix_{1,}x_2;y_1,y_2)`$ is the connected part of the Green’s function $`G_\mu ^a(x_ix_{1,}x_2;y_1,y_2)`$ and $`\mathrm{\Lambda }_\mu ^a(x_ix_1,x_2)`$, $`\mathrm{\Lambda }_\mu ^a(x_ix_1;y_1)`$, $`\mathrm{\Lambda }_\mu ^{𝐜a}(x_ix_2;y_2)`$ and $`\overline{\mathrm{\Lambda }}_\mu ^a(x_iy_1,y_2)`$ are the three-point Green’s functions which are given in Eqs. (2.15) and (3.10). On inserting the decomposition in Eq. (4.3) into Eq. (2.18), we see that the last unconnected term in Eq. (2.18) is cancelled out. Thus, we have
| $`𝒢_\mu ^a(x_ix_1,x_2;y_1,y_2)=G_{c\mu }^a(x_ix_{1,}x_2;y_1,y_2)+\mathrm{\Lambda }_\mu ^a(x_ix_1;y_1)S_F^c(x_2y_2)`$ |
| --- |
| $`+S_F(x_1y_1)\mathrm{\Lambda }_\mu ^{𝐜a}(x_ix_2;y_2)S_F^{}(x_1x_2)\overline{\mathrm{\Lambda }}_\mu ^a(x_iy_1,y_2)`$ |
(111)
where $`i=1,2`$. Substituting the above expression in Eq. (2.22), we will obtain the decomposition of the function $`𝒢^{(i)}(x_ix_1,x_2;y_1,y_2).`$
In the following, we are devoted to analyzing the terms on the RHS of Eq. (4.4) through the one-particle-irreducible decompositions of the connected Green’s functions included in those terms. The decompositions have been carried out in Appendix B.
### IV.1 The t-channel one-gluon exchange kernel
First we focus our attention on the second and third terms in Eq. (4.4). According to the decomposition in Eq. (B.15), the three-point gluon-quark Green’s functions $`\mathrm{\Lambda }_\mu ^a(x_ix_j;y_k)`$ which is fully connected can be represented in the form
$$\mathrm{\Lambda }_\mu ^a(x_ix_j;y_k)=d^4z_1\mathrm{\Sigma }_\mu ^a(x_ix_j;z_1)S_F(z_1y_k)$$
(112)
where
$$\mathrm{\Sigma }_\mu ^a(x_ix_j;z_1)=d^4u_1d^4u_2\mathrm{\Delta }_{\mu \nu }^{ab}(x_iu_1)S_F(x_ju_2)\mathrm{\Gamma }^{b\nu }(u_1u_2,z_1)$$
(113)
in which
$$\mathrm{\Delta }_{\mu \nu }^{ab}(x_iu_j)=\frac{1}{i}0^+\left|T[𝐀_\mu ^a(x_i)𝐀_\nu ^b(u_j)]\right|0^{}$$
(114)
is the exact gluon propagator and $`\mathrm{\Gamma }^{b\nu }(u_1u_2,z_1)`$ is the gluon-quark three-line proper vertex as defined in Eq. (B.17). The one-particle-irreducible decompositions of the three-point gluon-antiquark Green’s functions $`\mathrm{\Lambda }_\mu ^{𝐜a}(x_ix_j;y_k)`$ can be obtained from Eqs. (4.5) and (4.6) by the charge conjugation transformation. The result is
$$\mathrm{\Lambda }_\mu ^{𝐜a}(x_ix_j;y_k)=d^4z_2\mathrm{\Sigma }_\mu ^{𝐜a}(x_ix_j;z_2)S_F^c(z_2y_k)$$
(115)
where
$$\mathrm{\Sigma }_\mu ^{ca}(x_ix_j;z_2)=d^4u_1d^4u_2\mathrm{\Delta }_{\mu \nu }^{ab}(x_iu_1)S_F^c(x_ju_2)\mathrm{\Gamma }_c^{b\nu }(u_1u_2,z_2)$$
(116)
in which $`\mathrm{\Gamma }_c^{b\nu }(u_1u_2,z_2)`$ is the gluon-antiquark three-line proper vertex as defined in Eq. (B.18).
When we set $`i=j`$ in Eqs. (4.5) and (4.8), it is found that the $`\mathrm{\Sigma }_\mu ^a(x_1x_1;z_1)`$ in Eq. (4.6) and the $`\mathrm{\Sigma }_\mu ^{𝐜a}(x_2x_2;z_2)`$ in Eq. (4.9) are respectively related to the quark and antiquark self-energies in such a way
| $`\mathrm{\Omega }_1^{a\mu }\mathrm{\Sigma }_\mu ^a(x_1x_1;z_1)=\gamma _1^0\mathrm{\Sigma }(x_1,z_1)=\widehat{\mathrm{\Sigma }}(x_1,z_1),`$ |
| --- |
| $`\mathrm{\Omega }_2^{a\mu }\mathrm{\Sigma }_\mu ^{𝐜a}(x_2x_2;z_2)=\gamma _2^0\mathrm{\Sigma }^c(x_2,z_2)=\widehat{\mathrm{\Sigma }}^c(x_2,z_2)`$ |
(117)
Thus, from Eqs. (4.5), (4.8) and (4.10), we have
| $`\mathrm{\Omega }_1^{a\mu }\mathrm{\Lambda }_\mu ^a(x_1x_1;y_1)S_F^c(x_2y_2)+\mathrm{\Omega }_2^{a\mu }\mathrm{\Lambda }_\mu ^{𝐜a}(x_2x_2;z)S_F(x_1y_1)`$ |
| --- |
| $`=d^4z_1d^4z_2\mathrm{\Sigma }(x_1,x_2;z_1,z_2)S_F(z_1y_1)S_F^c(z_2y_2)`$ |
(118)
where
$$\mathrm{\Sigma }(x_1,x_2;z_1,z_2)=\widehat{\mathrm{\Sigma }}(x_1,z_1)\delta ^4(x_2z_2)+\widehat{\mathrm{\Sigma }}^c(x_2,z_2)\delta ^4(x_1z_1)$$
(119)
which is the total self-energy of the $`q\overline{q}`$ system. According to the definitions given in Eqs. (4.4), (2.22)-(2.24), (2.31) and (4.2), we see, $`\mathrm{\Sigma }(x_1,x_2;z_1,z_2)`$ as a self energy term to appear in the interaction kernel.
In the case of $`ij`$, from Eqs. (4.5) and (4.8), it can be written
| $`\mathrm{\Omega }_2^{a\mu }\mathrm{\Lambda }_\mu ^a(x_2x_1;y_1)S_F^c(x_2y_2)+\mathrm{\Omega }_1^{a\mu }\mathrm{\Lambda }_\mu ^{𝐜a}(x_1x_2;y_2)S_F(x_1y_1)`$ |
| --- |
| $`=d^4z_1d^4z_2K_t(x_1,x_2;z_1,z_2)S_F(z_1y_1)S_F^c(z_2y_2)`$ |
(120)
where
$$K_t(x_1,x_2;z_1,z_2)=\mathrm{\Omega }_2^{a\mu }\mathrm{\Sigma }_\mu ^a(x_2x_1;z_1)\delta ^4(x_2z_2)+\mathrm{\Omega }_1^{a\mu }\mathrm{\Sigma }_\mu ^{𝐜a}(x_1x_2;z_2)\delta ^4(x_1z_1)$$
(121)
Based on Eqs. (4.4), (2.22)-(2.24), (2.31) and (4.2), it is clear that the $`K_t(x_1,x_2;z_1,z_2)`$ acts as a part of the interaction kernel to appear in the D-S equation. As will be seen in section 6, this part is precisely the kernel of t-channel one-gluon exchange interaction..
### IV.2 The s-channel one-gluon exchange kernel
Next, we turn to the last term in Eq. (4.4). The one-particle irreducible decomposition of the three-point Green’s function in this term may also be found from Eqs. (4.5) and (4.6) by the charge conjugation transformation. The result is
$$\overline{\mathrm{\Lambda }}_\mu ^a(x_iy_1,y_2)=d^4z_1d^4z_2d^4z\mathrm{\Delta }_{\mu \nu }^{ab}(x_iz)\overline{\mathrm{\Gamma }}^{b\nu }(zz_1,z_2)S_F(z_1y_1)S_F^c(z_2y_2)$$
(122)
where $`\overline{\mathrm{\Gamma }}^{b\nu }(zz_1,z_2)`$ is the gluon-quark-antiquark proper vertex defined in Eq. (B.19). With the decomposition given above, according to Eqs. (2.22)-(2.24) and (2.31), the contribution of the last term in Eq. (4.4) to the kernel $`K(x_1,x_2;y_1,y_2)`$ can be found from the sum
$$\underset{i=1}{\overset{2}{}}\mathrm{\Omega }_i^{a\mu }S_F^{}(x_1x_2)\overline{\mathrm{\Lambda }}_\mu ^a(x_iy_1,y_2)=d^4z_1d^4z_2K_s(x_1,x_2;z_1,z_2)S_F(z_1y_1)S_F^c(z_2y_2)$$
(123)
where
$$K_s(x_1,x_2;z_1,z_2)=S_F^{}(x_1x_2)d^4z\underset{i=1}{\overset{2}{}}\mathrm{\Omega }_i^{a\mu }\mathrm{\Delta }_{\mu \nu }^{ab}(x_iz)\overline{\mathrm{\Gamma }}^{b\nu }(zz_1,z_2)$$
(124)
is just the s-channel one-gluon exchange kernel occurring in the D-S equation which will be discussed in section 6. It is noted here that the s-channel one-gluon exchange describes the $`q\overline{q}`$ annihilation and creation process which takes place between the quark (antiquark) in the initial state and the antiquark (quark) in the final state as indicated by the gluon propagator in Eq. (4.17) (not between the quark and the antiquark both of which are simultaneously related to the initial state or the final state B-S amplitude for a bound state).
### IV.3 The kernel from the Green’s function $`G_{c\mu }^a(x_ix_1,x_2;y_1,y_2)`$
Now let us concentrate our attention on the irreducible decomposition of the first term in Eq. (4.4). As stated in Appendix B, this decomposition may be derived from the functional differential of the Green’s function $`G_c(x_{1,}x_2;y_1,y_2)`$ with respect to the source $`J^{a\mu }(x_i)`$ by using the one-particle irreducible decomposition of the function $`G_c(x_{1,}x_2;y_1,y_2)`$. The latter decomposition whose derivation is sketched in Appendix B is well-known and can be represented in the form
| $`G_c(x_{1,}x_2;y_1,y_2)=\underset{i=1}{\overset{2}{}}d^4u_id^4v_iS_F(x_1u_1)S_F^c(x_2u_2)`$ |
| --- |
| $`\times \mathrm{\Gamma }(u_1,u_2;v_1,v_2)S_F(v_1y_1)S_F^c(v_2y_2)`$ |
(125)
where
$$\mathrm{\Gamma }(u_1,u_2;v_1,v_2)=\underset{i=1}{\overset{3}{}}\mathrm{\Gamma }_i(u_1,u_2;v_1,v_2)$$
(126)
in which
$$\mathrm{\Gamma }_1(u_1,u_2;v_1,v_2)=d^4z_1d^4z_2\mathrm{\Gamma }^{b\nu }(z_1u_1,v_1)D_{\nu \nu ^{}}^{bb^{}}(z_1z_2)\mathrm{\Gamma }_c^{b^{}\nu ^{}}(z_2u_2,v_2)$$
(127)
$$\mathrm{\Gamma }_2(u_1,u_2;v_1,v_2)=d^4z_1d^4z_2\mathrm{\Gamma }^{b\nu }(z_1u_1,u_2)D_{\nu \nu ^{}}^{bb^{}}(z_1z_2)\overline{\mathrm{\Gamma }}^{b^{}\nu ^{}}(z_2v_1,v_2)$$
(128)
here the three-line vertices are defined in Eqs. (B.17)-(B.19) and (B.21), $`D_{\nu \nu ^{}}^{bb^{}}(z_1z_2)=i\mathrm{\Delta }_{\nu \nu ^{}}^{bb^{}}(z_1z_2)`$ with $`\mathrm{\Delta }_{\nu \nu ^{}}^{bb^{}}(z_1z_2)`$ defined in Eq. (4.7) and $`\mathrm{\Gamma }_3(u_1,u_2;v_1,v_2)`$ defined in Eq. (B.22) is the quark-antiquark four-line proper vertex. After substituting the expressions in Eqs. (4.18)-(4.21), which are now given in the presence of source $`J`$, into Eq. (B.6) and completing the differentiation, the one-particle irreducible decomposition of the Green’s function $`G_{c\mu }^a(x_ix_{1,}x_2;y_1,y_2)`$ will be found and, thereby, we can write
| $`\underset{i=1}{\overset{2}{}}G_c^{(i)}(x_ix_{1,}x_2;y_1,y_2)\mathrm{\Omega }_1^{a\mu }G_{c\mu }^a(x_1x_{1,}x_2;y_1,y_2)+\mathrm{\Omega }_2^{a\mu }G_{c\mu }^a(x_2x_{1,}x_2;y_1,y_2)`$ |
| --- |
| $`=\underset{j=1}{\overset{3}{}}G_j(x_{1,}x_2;y_1,y_2)`$ |
(129)
where
| $`G_1(x_{1,}x_2;y_1,y_2)`$ |
| --- |
| $`=\underset{j=1}{\overset{2}{}}d^4u_jd^4v_j\underset{i=1}{\overset{2}{}}\mathrm{\Omega }_i^{a\mu }[\mathrm{\Lambda }_\mu ^a(x_ix_1;u_1)S_F^c(x_2u_2)+S_F(x_1u_1)`$ |
| $`\times \mathrm{\Lambda }_\mu ^{𝐜a}(x_ix_2;u_2)]\mathrm{\Gamma }(u_1,u_2;v_1,v_2)S_F(v_1y_1)S_F^c(v_2y_2)`$ |
(130)
| $`G_2(x_{1,}x_2;y_1,y_2)`$ |
| --- |
| $`=\underset{j=1}{\overset{2}{}}d^4u_jd^4v_j\underset{i=1}{\overset{2}{}}\mathrm{\Omega }_i^{a\mu }S_F(x_1u_1)S_F^c(x_2u_2)\mathrm{\Gamma }(u_1,u_2;v_1,v_2)`$ |
| $`\times [\mathrm{\Lambda }_\mu ^a(x_iv_1;y_1)S_F^c(v_2y_2)+\mathrm{\Lambda }_\mu ^{𝐜a}(x_iv_2;y_2)S_F(v_1y_1)]`$ |
(131)
and
| $`G_3(x_{1,}x_2;y_1,y_2)`$ |
| --- |
| $`=\underset{j=1}{\overset{2}{}}d^4u_jd^4v_j\underset{i=1}{\overset{2}{}}\mathrm{\Omega }_i^{a\mu }S_F(x_1u_1)S_F^c(x_2u_2)`$ |
| $`\times \mathrm{\Gamma }^{a\mu }(x_iu_1,u_2;v_1,v_2)S_F(v_1y_1)S_F^c(v_2y_2)`$ |
(132)
where $`\mathrm{\Gamma }^{a\mu }(x_iu_1,u_2;v_1,v_2)`$ is a kind of five-line vertex which is defined in Eq. (B.27) and will be specified soon later.
Before specifying the function $`\mathrm{\Gamma }^{a\mu }(x_iu_1,u_2;v_1,v_2)`$, we first analyze the expressions in Eqs. (4.23) and (4.24). According to the expressions in Eqs. (4.11), (4.13) and (4.18), Eq. (4.23) may be represented as
$$G_1(x_{1,}x_2;y_1,y_2)=d^4z_1d^4z_2[\mathrm{\Sigma }(x_1,x_2;z_1,z_2)+K_t(x_1,x_2;z_1,z_2)]G_c(z_{1,}z_2;y_1,y_2)$$
(133)
where $`\mathrm{\Sigma }(x_1,x_2;z_1,z_2)`$ and $`K_t(x_1,x_2;z_1,z_2)`$ were respectively described in Eqs. (4.12) and (4.14). In view of the decompositions in Eqs. (4.5) and (4.8), Eq. (4.24) may be written in the form
$$G_2(x_{1,}x_2;y_1,y_2)=d^4z_1d^4z_2K_1(x_1,x_2;z_1,z_2)S_F(z_1y_1)S_F^c(z_2y_2)$$
(134)
where
| $`K_1(x_1,x_2;z_1,z_2)=\underset{j=1}{\overset{2}{}}d^4u_jd^4v\underset{i=1}{\overset{2}{}}\mathrm{\Omega }_i^{a\mu }S_F(x_1u_1)S_F^c(x_2u_2)`$ |
| --- |
| $`\times [\mathrm{\Gamma }(u_1,u_2;v,z_2)\mathrm{\Sigma }_\mu ^a(x_iv;z_1)+\mathrm{\Gamma }(u_1,u_2;z_1,v)\mathrm{\Sigma }_\mu ^{𝐜a}(x_iv;z_2)]`$ |
(135)
in which $`\mathrm{\Sigma }_\mu ^a(x_iv;z_1)`$ and $`\mathrm{\Sigma }_\mu ^{𝐜a}(x_iv;z_2)`$ were represented in Eqs (4.6) and (4.9) respectively.
Let us turn to the five-line vertex function contained in Eq. (4.25). This vertex is two-particle reducible (or say, B-S reducible) although it is one-particle-irreducible. Therefore, it can be decomposed into a B-S irreducible part $`\mathrm{\Gamma }_{IR}^{a\mu }`$ and a B-S reducible part $`\mathrm{\Gamma }_{RE}^{a\mu }`$
$$\mathrm{\Gamma }^{a\mu }(x_iu_1,u_2;v_1,v_2)=\mathrm{\Gamma }_{IR}^{a\mu }(x_iu_1,u_2;v_1,v_2)+\mathrm{\Gamma }_{RE}^{a\mu }(x_iu_1,u_2;v_1,v_2)$$
(136)
In order to exhibit the above decomposition specifically, as mentioned in Appendix B, we may insert Eqs. (4.19)- (4.21), which are now given in the case of presence of the source $`J,`$ into Eq. (B.27) and complete the differentiation with respect to the source $`J^{a\mu }(x_i)`$. After completing the differentiations shown in Eqs. (B.23) and (B.25), we can write
$$\mathrm{\Gamma }^{a\mu }(x_iu_1,u_2;v_1,v_2)=\underset{j=1}{\overset{3}{}}\mathrm{\Gamma }_j^{a\mu }(x_iu_1,u_2;v_1,v_2)$$
(137)
where $`\mathrm{\Gamma }_j^{a\mu }(x_iu_1,u_2;v_1,v_2)`$ are given by the differential of the functions $`\mathrm{\Gamma }_j(u_1,u_2;v_1,v_2)`$ in Eq. (4.19) and separately shown below
| $`\mathrm{\Gamma }_1^{a\mu }(x_iu_1,u_2;v_1,v_2)`$ |
| --- |
| $`=d^4zd^4z_1d^4z_2D_{\mu \nu }^{ab}(x_iz)\{\mathrm{\Gamma }_\nu ^{bc,\lambda }(z,z_1u_1,v_1)D_{\lambda \lambda ^{}}^{cc^{}}(z_1z_2)\mathrm{\Gamma }_𝐜^{c^{}\lambda ^{}}(z_2u_2,v_2)`$ |
| $`+\mathrm{\Gamma }^{c\lambda }(z_1u_1,v_1)[D_{\lambda \lambda ^{}}^{cc^{}}(z_1z_2)\mathrm{\Gamma }_{𝐜\nu }^{bc^{},\lambda }(z,z_2u_2,v_2)+\mathrm{\Pi }_{\nu \lambda \lambda ^{}}^{bcc^{}}(z,z_1,z_2)\mathrm{\Gamma }^{c^{}\lambda ^{}}(z_2u_2,v_2)]\}`$ |
(138)
in which besides the gluon-quark and gluon-antiquark three-line vertices mentioned before, there occur the gluon-quark four-line proper vertex $`\mathrm{\Gamma }_{\nu \lambda }^{bc}(z,z_1u_1,v_1)`$ and the corresponding gluon-antiquark one $`\mathrm{\Gamma }_{𝐜\nu \lambda }^{bc}(z,z_1u_1,v_1)`$ which are defined respectively in Eqs. (B.28) and (B.29) and
$$\mathrm{\Pi }_{\nu \rho \sigma }^{bcd}(z,z_1,z_2)=d^4u_1d^4u_2D_{\rho \rho ^{}}^{cc^{}}(z_1u_1)\mathrm{\Gamma }_{bc^{}d^{}}^{\nu \rho ^{}\sigma ^{}}(z,u_1,u_2)D_{\sigma ^{}\sigma }^{d^{}d}(u_2z_2)$$
(139)
in which $`\mathrm{\Gamma }_{bc^{}d^{}}^{\nu \rho ^{}\sigma ^{}}(z,u_1,u_2)`$ is the gluon three-line proper vertex defined in Eq. (B.24),
| $`\mathrm{\Gamma }_2^{a\mu }(x_iu_1,u_2;v_1,v_2)`$ |
| --- |
| $`=d^4zd^4z_1d^4z_2D_{\mu \nu }^{ab}(x_iz)\{\mathrm{\Gamma }_\nu ^{bc,\lambda }(z,z_1u_1,u_2)D_{\lambda \lambda ^{}}^{cc^{}}(z_1z_2)\overline{\mathrm{\Gamma }}_𝐜^{c^{}\lambda ^{}}(z_2v_1,v_2)`$ |
| $`+\mathrm{\Gamma }^{c\lambda }(z_1u_1,u_2)[D_{\lambda \lambda ^{}}^{cc^{}}(z_1z_2)\overline{\mathrm{\Gamma }}_{𝐜\nu }^{bc^{},\lambda }(z,z_2v_1,v_2)+\mathrm{\Pi }_{\nu \lambda \lambda ^{}}^{bcc^{}}(z,z_1,z_2)\overline{\mathrm{\Gamma }}^{c^{}\lambda ^{}}(z_2v_1,v_2)]\}`$ |
(140)
in which $`\mathrm{\Gamma }_{\nu \lambda }^{bc}(z,z_1u_1,u_2)`$ and $`\overline{\mathrm{\Gamma }}_{𝐜\nu \lambda }^{bc^{}}(z,z_1v_1,v_2)`$ are the gluon-quark-antiquark four-line proper vertices defined in Eqs. (B.30) and (B.31) and especially
$$\mathrm{\Gamma }_3^{a\mu }(x_iu_1,u_2;v_1,v_2)=d^4zD_{\mu \nu }^{ab}(x_iz)\widehat{\mathrm{\Gamma }}^{b\nu }(zu_1,u_2;v_1,v_2)$$
(141)
in which $`\widehat{\mathrm{\Gamma }}^{b\nu }(zu_1,u_2;v_1,v_2)`$ stands for the gluon-quark-antiquark five-line proper vertex defined in Eq. (B.33).
It is pointed out that the vertices formulated in Eqs. (4.31) and (4.33) are not only one-particle-irreducible, but also B-S irreducible. This point can be seen more clearly when the vertices are represented by their Feynman diagrams. From the diagrams, one can find that it is impossible to divide each of the diagrams into two unconnected parts by cutting two fermion lines. However, the five-line proper vertex $`\widehat{\mathrm{\Gamma }}^{b\nu }(zu_1,u_2;v_1,v_2)`$ in Eq. (4.34) is B-S reducible. It can be decomposed into a B-S reducible part $`\widehat{\mathrm{\Gamma }}_{RE}^{b\nu }`$ and a B-S irreducible part $`\widehat{\mathrm{\Gamma }}_{IR}^{b\nu },`$
$$\widehat{\mathrm{\Gamma }}^{b\nu }(zu_1,u_2;v_1,v_2)=\widehat{\mathrm{\Gamma }}_{IR}^{b\nu }(zu_1,u_2;v_1,v_2)+\widehat{\mathrm{\Gamma }}_{RE}^{b\nu }(zu_1,u_2;v_1,v_2)$$
(142)
This enables us to write Eq. (4.34) as follows:
$$\mathrm{\Gamma }_3^{a\mu }(x_iu_1,u_2;v_1,v_2)=\mathrm{\Gamma }_{31}^{a\mu }(x_iu_1,u_2;v_1,v_2)+\mathrm{\Gamma }_{32}^{a\mu }(x_iu_1,u_2;v_1,v_2)$$
(143)
where
$$\mathrm{\Gamma }_{31}^{a\mu }(x_iu_1,u_2;v_1,v_2)=d^4zD_{\mu \nu }^{ab}(x_iz)\widehat{\mathrm{\Gamma }}_{IR}^{b\nu }(zu_1,u_2;v_1,v_2)$$
(144)
and
$$\mathrm{\Gamma }_{32}^{a\mu }(x_iu_1,u_2;v_1,v_2)=d^4zD_{\mu \nu }^{ab}(x_iz)\widehat{\mathrm{\Gamma }}_{RE}^{b\nu }(zu_1,u_2;v_1,v_2)$$
(145)
From the above statement, it is clearly seen that the B-S irreducible part of the vertex in Eq. (4.29) is given by the sum
| $`\mathrm{\Gamma }_{IR}^{a\mu }(x_iu_1,u_2;v_1,v_2)=\mathrm{\Gamma }_1^{a\mu }(x_iu_1,u_2;v_1,v_2)`$ |
| --- |
| $`+\mathrm{\Gamma }_2^{a\mu }(x_iu_1,u_2;v_1,v_2)+\mathrm{\Gamma }_{31}^{a\mu }(x_iu_1,u_2;v_1,v_2)`$ |
(146)
where the three terms on the RHS of Eq. (4.39) were given in Eqs. (4.31), (4.33) and (4.37) respectively. While, the B-S reducible part in Eq. (4.29) is given by Eq. (4.38)
$$\mathrm{\Gamma }_{RE}^{a\mu }(x_iu_1,u_2;v_1,v_2)=\mathrm{\Gamma }_{32}^{a\mu }(x_iu_1,u_2;v_1,v_2)$$
(147)
Based on the decomposition in Eq. (4.29), the function in Eq. (4.25) will be decomposed into
$$G_3(x_{1,}x_2;y_1,y_2)=G_3^{(1)}(x_{1,}x_2;y_1,y_2)+G_3^{(2)}(x_{1,}x_2;y_1,y_2)$$
(148)
where
$$G_3^{(1)}(x_{1,}x_2;y_1,y_2)=d^4z_1d^4z_2K_2(x_1,x_2;z_1,z_2)S_F(z_1y_1)S_F^c(z_2y_2)$$
(149)
in which
$$K_2(x_1,x_2;z_1,z_2)=\underset{j=1}{\overset{2}{}}d^4u_j\underset{i=1}{\overset{2}{}}\mathrm{\Omega }_i^{a\mu }S_F(x_1u_1)S_F^c(x_2u_2)\mathrm{\Gamma }_{IR}^{a\mu }(x_iu_1,u_2;z_1,z_2)$$
(150)
and
| $`G_3^{(2)}(x_{1,}x_2;y_1,y_2)=\underset{j=1}{\overset{2}{}}d^4u_jd^4v_j\underset{i=1}{\overset{2}{}}\mathrm{\Omega }_i^{a\mu }S_F(x_1u_1)S_F^c(x_2u_2)`$ |
| --- |
| $`\times \mathrm{\Gamma }_{RE}^{a\mu }(x_iu_1,u_2;v_1,v_2)S_F(v_1y_1)S_F^c(v_2y_2)`$ |
(151)
It is emphasized that due to the B-S irreducibility of the vertex $`\mathrm{\Gamma }_{IR}^{a\mu }(x_iu_1,u_2;v_1,v_2)`$, the function $`G_3^{(1)}(x_{1,}x_2;y_1,y_2)`$ can only be written in the form as shown in Eq. (4.42). While, since the vertex $`\mathrm{\Gamma }_{RE}^{a\mu }(x_iu_1,u_2;v_1,v_2)`$ is B-S reducible, as was similarly done for the Green’s function $`𝒢^{(i)}(x_1,x_2;y_1,y_2)`$, the function $`G_3^{(2)}(x_{1,}x_2;y_1,y_2),`$ as a part of the connected Green’s function, must be represented in the B-S reducible form
$$G_3^{(2)}(x_{1,}x_2;y_1,y_2)=d^4z_1d^4z_2\stackrel{~}{K}(x_1,x_2;z_1,z_2)G_c(z_{1,}z_2;y_1,y_2)$$
(152)
where $`\stackrel{~}{K}(x_1,x_2;z_1,z_2)`$ is a kind of kernel needed to be determined later.
### IV.4 Complete expression of the interaction kernel
Up to the present, the irreducible decompositions of the functions $`𝒢^{(i)\text{ }}(x_1,x_2;y_1,y_2)`$ $`(i=1,2)`$ appearing in Eq. (2.24) have been completed. Now, let us first sum up the expressions given in Eqs. (4.11) and (4.13) which correspond to the second and third terms in Eq. (4.4) and the expression in Eq. (4.26) for the function $`G_1(x_{1,}x_2;y_1,y_2)`$ which is contained in the connected Green’s functions $`G_{c\mu }^a(x_ix_{1,}x_2;y_1,y_2)`$. The summation yields
$$d^4z_1d^4z_2[\mathrm{\Sigma }(x_1,x_2;z_1,z_2)+K_t(x_1,x_2;z_1,z_2)]𝒢(z_{1,}z_2;y_1,y_2)$$
(153)
where the relation in Eq. (4.2) has been considered. Then, we combine the expression in Eq. (4.16) which corresponds to the last term in Eq. (4.4) and the expressions in Eqs. (4.27) and (4.42) which come from the B-S irreducible part of connected Green’s functions $`G_{c\mu }^a(x_ix_{1,}x_2;y_1,y_2)`$ and obtain
$$d^4z_1d^4z_2[K_s(x_1,x_2;z_1,z_2)+K_1(x_1,x_2;z_1,z_2)+K_2(x_1,x_2;z_1,z_2)]S_F(z_1y_1)S_F^c(z_2y_2)$$
(154)
The final decomposition of the sum of the functions $`𝒢^{(1)}(x_1x_1,x_2;y_1,y_2)`$ and $`𝒢^{(2)}(x_2x_1,x_2;y_1,y_2)`$ will be given by the sum of Eqs. (4.45)-(4.47). Obviously, in order to make the D-S equation to be closed, the kernel In Eq. (4.45) must be
$$\stackrel{~}{K}(x_1,x_2;z_1,z_2)=K_s(x_1,x_2;z_1,z_2)+K_1(x_1,x_2;z_1,z_2)+K_2(x_1,x_2;z_1,z_2)$$
(155)
Thus, the summation of Eqs. (4.45) and (4.47) gives
$$d^4z_1d^4z_2[K_s(x_1,x_2;z_1,z_2)+K_1(x_1,x_2;z_1,z_2)+K_2(x_1,x_2;z_1,z_2)]𝒢(z_{1,}z_2;y_1,y_2)$$
(156)
Combining Eqs. (4.46) and (4.49), we eventually arrive at
$$\underset{i=1}{\overset{2}{}}𝒢^{(i)}(x_ix_1,x_2;y_1,y_2)=d^4z_1d^4z_2K(x_1,x_2;z_1,z_2)𝒢(z_{1,}z_2;y_1,y_2)$$
(157)
where
| $`K(x_1,x_2;z_1,z_2)=\mathrm{\Sigma }(x_1,x_2;z_1,z_2)+K_t(x_1,x_2;z_1,z_2)`$ |
| --- |
| $`+K_s(x_1,x_2;z_1,z_2)+K_1(x_1,x_2;z_1,z_2)+K_2(x_1,x_2;z_1,z_2)`$ |
(158)
This just is the kernel appearing in Eq. (2.38). The five terms on the RHS of Eq. (4.51), as respectively shown in Eqs. (4.12), (4.14), (4.17), (4.28) and (4.43), are only represented in terms of the quark, antiquark and gluon propagators and some kinds of three, four and five-line proper vertices and therefore exhibits a more specific structure of the kernel. The equivalence between the both expressions given in the preceding section and this section for the kernel will be illustrated in section 6 for the one-gluon exchange kernels. The exact proof of the equivalence has been done by the technique of irreducible decomposition of the Green’s functions. From the proof, we find that the expression described in this section can surely be obtained from the expression given in the preceding section.
## V Pauli-Schrödinger equation
As mentioned in Introduction, the D-S equations formulated in the Dirac spinor space may be reduced to an equivalent equations represented in the Pauli spinor space with the help of Dirac spinors. Let us start from the equation given in Eq. (2.44). The Dirac spinors are defined as
$$U(\stackrel{}{p})=\sqrt{\frac{\omega +m}{2\omega }}\left(\begin{array}{c}1\\ \frac{\stackrel{}{\sigma }\stackrel{}{p}}{\omega +m}\end{array}\right)$$
(159)
$$V(\stackrel{}{p})=\sqrt{\frac{\omega +m}{2\omega }}\left(\begin{array}{c}\frac{\stackrel{}{\sigma }\stackrel{}{p}}{\omega +m}\\ 1\end{array}\right)$$
(160)
where $`\stackrel{}{\sigma \text{ }}`$ are the Pauli matrices, $`\omega =(\stackrel{}{p}^2+m^2)^{1/2}`$, $`U(\stackrel{}{p})`$ and $`V(\stackrel{}{p})`$ are the positive energy and negative energy spinors respectively. They satisfy the orthonornality relations
$$\begin{array}{c}U^+(\stackrel{}{p})U(\stackrel{}{p})=V^+(\stackrel{}{p})V(\stackrel{}{p})=1\\ U^+(\stackrel{}{p})V(\stackrel{}{p})=V^+(\stackrel{}{p})U(\stackrel{}{p})=0\end{array}$$
(161)
and the completeness relation
$$\mathrm{\Lambda }^+(\stackrel{}{p})+\mathrm{\Lambda }^{}(\stackrel{}{p})=1$$
(162)
where $`\mathrm{\Lambda }^+(\stackrel{}{p})`$ and $`\mathrm{\Lambda }^{}(\stackrel{}{p})`$ are respectively the positive and negative energy state projection operators defined by
$$\mathrm{\Lambda }^+(\stackrel{}{p})=U(\stackrel{}{p})U^+(\stackrel{}{p}),\text{ }\mathrm{\Lambda }^{}(\stackrel{}{p})=V(\stackrel{}{p})V^+(\stackrel{}{p})$$
(163)
Define
$$W_a(\stackrel{}{p})=\left\{\begin{array}{c}U(\stackrel{}{p}),\text{ }if\text{ }a=+,\\ V(\stackrel{}{p}),\text{ }if\text{ }a=,\end{array}\right\}$$
(164)
then, the two fermion spinors can be written as
$$W_{ab}(\stackrel{}{P},\stackrel{}{q})=W_a(\stackrel{}{p_1})W_b(\stackrel{}{p_2})$$
(165)
With this definition, the completeness relation for two fermion spinors can be represented as
$$\underset{ab}{}W_{ab}(\stackrel{}{P},\stackrel{}{q})W_{ab}^+(\stackrel{}{P},\stackrel{}{q})=1$$
(166)
Premiltiplying Eq. (2.44) with $`W_{ab}^+(\stackrel{}{P},\stackrel{}{q})`$, sandwiching Eq. (5.8) between the kernel $`K(P,q,k)`$ and the amplitude $`\chi _{P\varsigma }(k)`$ on the RHS of Eq. (2.44) and applying the Dirac equation
$$h(\stackrel{}{p})W_a(\stackrel{}{p})=a\omega (\stackrel{}{p})W_a(\stackrel{}{p})$$
(167)
we obtain
$$\mathrm{\Delta }_{ab}(P,\stackrel{}{q})\varphi _{ab}(P,q)=\underset{cd}{}\frac{d^4k}{(2\pi )^4}K_{abcd}(P,q,k)\varphi _{cd}(P,k)$$
(168)
where $`a,b,c,d=\pm 1`$, $`P=(E,\stackrel{}{P})`$,
$$\varphi _{ab}(P,q)=W_{ab}^+(\stackrel{}{P},\stackrel{}{q})\chi _{P\varsigma }(q)$$
(169)
$$\mathrm{\Delta }_{ab}(P,\stackrel{}{q})=Ea\omega _1(\stackrel{}{p_1})b\omega _2(\stackrel{}{p_2})$$
(170)
$$K_{abcd}(P,q,k)=W_{ab}^+(\stackrel{}{P},\stackrel{}{q})K(P,q,k)W_{cd}(\stackrel{}{P},\stackrel{}{k})$$
(171)
Eq. (5.10) is a set of coupled equations satisfied by the amplitudes $`\varphi _{ab}(P,q)`$ each of which is represented in the Pauli spinor space and of dimension four. In the infinite-dimensional space of the momentum $`q`$ or $`k`$, according to $`ab=++`$ and $`ab++`$, Eq. (5.10) may be, in the matrix form, separately written as
$$\mathrm{\Delta }_{++}(p)\varphi _{++}(P)=K_{++++}(P)\varphi _{++}(P)+\underset{cd++}{}K_{++cd}(P)\varphi _{cd}(P)$$
(172)
and
$$\mathrm{\Delta }_{ab}(P)\varphi _{ab}(P)=K_{ab++}(P)\varphi _{++}(P)+\underset{cd++}{}K_{abcd}(P)\varphi _{cd}(P)$$
(173)
where $`ab++`$ and the terms related to $`\varphi _{++}(P)`$ have been separated out from the others. Furthermore, In the three-dimensional spinor space spanned by $`\varphi _{ab}(P)`$ with $`ab++`$, Eqs. (5.14) and (5.15) may be written in the full matrix form
$$\mathrm{\Delta }_+(P)\psi (P)=K_+(P)\psi (P)+\overline{K}^t(P)\varphi (P)$$
(174)
and
$$\varphi (P)=\overline{G}(P)\psi (P)+G(P)\varphi (P)$$
(175)
where $`\psi (P)=\varphi _{++}(P)`$, $`\mathrm{\Delta }_+(P)=\mathrm{\Delta }_{++}(P)`$, $`K_+(P)=K_{++++}(P)`$, while, $`\varphi (P)=\{\varphi _{ab}(P)\}`$, $`\overline{K}^t(P)=\{K_{++cd}(P)\}`$, $`\overline{G}(P)=\{K_{ab++}(P)/\mathrm{\Delta }_{ab}(P)\}`$ and $`G(P)=\{K_{abcd}(P)/\mathrm{\Delta }_{ab}(P)\}`$ represent the matrices in the three-dimensional spinor space. Solving the equation (5.17), we obtain
$$\varphi (P)=\frac{1}{1G(P)}\overline{G}(P)\psi (P)$$
(176)
Substituting the above expression into Eq. (5.16), we finally arrive at
$$\mathrm{\Delta }_+(P)\psi (P)=V(P)\psi (P)$$
(177)
where
$$V(P)=K_+(P)+\overline{K}^t(P)\frac{1}{1G(P)}\overline{G}(P)$$
(178)
which identifies itself with the interaction Hamiltonian. With the definition
$$\frac{1}{1G(P)}=\underset{n=0}{\overset{\mathrm{}}{}}G^{(n)}(P),$$
(179)
Eq. (5.20) can be written as
$$V(P)=\underset{n=0}{\overset{\mathrm{}}{}}V^{(n)}(P)$$
(180)
where
| $`V^{(0)}(P)=K_+(P),`$ |
| --- |
| $`V^{(1)}(P)=\overline{K}^t(P)\overline{G}(P),`$ |
| $`V^{(2)}(P)=\overline{K}^t(P)G(P)\overline{G}(P),`$ |
| $`\mathrm{}\mathrm{}`$ |
(181)
Written out explicitly, Eq. (5.19) reads
$$[E\omega _1(\stackrel{}{p_1})\omega _2(\stackrel{}{p_2})]\psi (P,q)=\frac{d^4k}{(2\pi )^4}V(P,q,k)\psi (P,k)$$
(182)
The terms in the interaction Hamiltonian in Eq. (5.23) are specified as
$$V^{(0)}(P,q,k)=K_{++++}(P,q,k)$$
(183)
$$V^{(1)}(P,q,k)=\underset{ab++}{}\frac{d^4l}{(2\pi )^4}\frac{K_{++ab}(P,q,l)K_{ab++}(P,l,k)}{Ea\overline{\omega }_1(\stackrel{}{l})b\overline{\omega }_2(\stackrel{}{l})}$$
(184)
$$V^{(2)}(P,q,k)=\underset{ab++}{}\underset{cd++}{}\frac{d^4l_1}{(2\pi )^4}\frac{d^4l_2}{(2\pi )^4}\frac{K_{++ab}(P,q,l_1)K_{abcd}(P,l_1,l_2)K_{cd++}(P,l_2,k)}{[Ea\overline{\omega }_1(\stackrel{}{l_1})b\overline{\omega }_2(\stackrel{}{l_1})][Ea\overline{\omega }_1(\stackrel{}{l_2})b\overline{\omega }_2(\stackrel{}{l_2})]}$$
(185)
and so on, where for simplicity of representation, we have defined $`\overline{\omega }_1(\stackrel{}{l})=\omega _1(\eta _1\stackrel{}{P}+\stackrel{}{l})`$ and $`\overline{\omega }_2(\stackrel{}{l})=\omega _2(\eta _2\stackrel{}{P}\stackrel{}{l})`$. In the center of mass frame, $`\overline{\omega }_i(\stackrel{}{l})=\omega _i(\stackrel{}{l})`$ ($`i=1,2`$). Eq. (5.24) is the equation satisfied by the positive energy state amplitude $`\psi (P,q)`$ which is of dimension four in the two-fermion Pauli spinor space. This is the reason why the above equation is called Pauli-Schrödinger (P-S) equation.
By the same procedure, the D-S equation in Eq. (2.45) may also be reduced to a corresponding P-S equation as represented in the following
$$[q_0\eta _2\omega _1(\stackrel{}{p_1})+\eta _1\omega _2(\stackrel{}{p_2})]\psi (P,q)=\frac{d^4k}{(2\pi )^4}\overline{V}(P,q,k)\psi (P,k)$$
(186)
where $`q_0`$ is the relative energy and $`\overline{V}(P,q,k)`$ is a kind of interaction Hamiltonian which can be written out from the expression of $`V(P,q,k)`$ by the replacement: $`K(P,q,k)\overline{K}(P,q,k)`$ and $`\mathrm{\Delta }_{ab}(P,\stackrel{}{q})`$ $`\overline{\mathrm{\Delta }}_{ab}(\stackrel{}{P},q)=q_0a\eta _2\omega _1(\stackrel{}{p_1})+b\eta _1\omega _2(\stackrel{}{p_2})`$. For the three-dimensional D-S equation, the P-S equation in Eq. (5.28) disappears. We are left only with a three-dimensional P-S equation derived from Eq. (2.56) such that
$$[E\omega _1(\stackrel{}{p_1})\omega _2(\stackrel{}{p_2})]\psi (P,\stackrel{}{q})=\frac{d^3k}{(2\pi )^3}V(P,\stackrel{}{q},\stackrel{}{k})\psi (P,\stackrel{}{k})$$
(187)
where the Hamiltonian $`V(P,\stackrel{}{q},\stackrel{}{k})`$ formally has the same expressions as written in Eqs. (5.22) and (5.25)-(5.27) except that the four-dimensional kernel $`K(P,q,k)`$ in those expressions is now replaced by the three-dimensional one $`K(P,\stackrel{}{q},\stackrel{}{k})`$ which is the Fourier transform of the kernel in Eq. (3.40)
It is worthy to point out that for a given kernel in the D-S equation, there are a series of terms (the ladder diagrams) to appear in the interaction Hamiltonian in the P-S equation. If the D-S equation with a given kernel could be solved, the contribution arising from a series of ladder diagrams characterized by the series of terms in the Hamiltonian are precisely taken into account. Another point we would like to stress is that as seen from Eqs. (5.26) and (5.27), the negative energy state only acts as intermediate states to appear in the interaction Hamiltonian. Particularly, for the bound state, the positive energy state does not appear in the intermediate states. While, for the scattering state P-S equation as discussed in Ref. , the intermediate states in the interaction Hamiltonian must include the positive energy state. In this case, the series expansion of the interaction Hamiltonian in Eq. (5.22) has an one-to-one correspondence with the perturbative expansion of the S-matrix. The above statement reveals an essential difference between the interactions taking place in the bound state and the scattering state.
## VI One-gluon exchange kernels
In this section, we limit ourself to give a brief derivation and description of the one gluon exchange kernel (OGEK). First we discuss the t-channel OGEK and the s-channel OGEK appearing in the four-dimensional D-S equation to illustrate the equivalence between the expressions of the interaction kernel derived in the sections 3 and 4. Then, we show the OGEK in the three-dimensional D-S equation and the corresponding Hamiltonian in the P-S equation.
### VI.1 The t-channel one-gluon exchange kernel
The exact form of the t-channel OGEK was represented in Eq. (4.14) with $`\mathrm{\Sigma }_\mu ^a(x_2x_1;z_1)`$ and $`\mathrm{\Sigma }_\mu ^{𝐜a}(x_1x_2;z_2)`$ given in Eqs. (4.6) and (4.9). In the lowest-order approximation, the propagators and the vertices in Eqs. (4.6) and (4.9) are taken respectively to be the free ones and the bare ones. The bare vertices are of the form
$$\begin{array}{c}\mathrm{\Gamma }^{b\nu }(u_1u_2,z_1)=ig\gamma ^\nu T^b\delta ^4(u_1u_2)\delta ^4(u_2z_1)\\ \mathrm{\Gamma }_c^{b\nu }(u_1u_2,z_2)=ig\gamma ^\nu \overline{T}^b\delta ^4(u_1u_2)\delta ^4(u_2z_2)\end{array}$$
(188)
With the vertices given above, the kernel in Eq. (4.14) becomes
| $`K_t^0(XY,x,y)=ig^2T^a\overline{T}^b\{\mathrm{\Delta }_{\mu \nu }^{ab}(x_2y_1)S_F(x_1y_1)\gamma _1^\mu \gamma _2^0\gamma _2^\nu \delta ^4(x_2y_2)`$ |
| --- |
| $`+\mathrm{\Delta }_{\mu \nu }^{ab}(x_1y_2)S_F^c(x_2y_2)\gamma _2^\nu \gamma _1^0\gamma _1^\mu \delta ^4(x_1y_1)\}`$ |
(189)
From now on, the $`S_F(xy)`$ and $`\mathrm{\Delta }_{\mu \nu }^{ab}(xy)`$ in the above are understood to be free propagators. By the Fourier transformation, we get in the momentum space that
$$K_t^0(P,q,k)=S(P,q)ig^2T^a\overline{T}^b\mathrm{\Delta }_{\mu \nu }^{ab}(qk)\gamma _1^\mu \gamma _2^\nu $$
(190)
where
$$S(P,q)=S_F(p_1)\gamma _2^0+S_F^c(p_2)\gamma _1^0=[\widehat{S}_F(p_1)+\widehat{S}_F^c(p_2)]\gamma _1^0\gamma _2^0$$
(191)
in which
$$S_F(p)=\widehat{S}_F(p)\gamma ^0$$
(192)
$$\widehat{S}_F(p)=\frac{1}{p_0h(\stackrel{}{p})+i\epsilon }=\frac{\mathrm{\Lambda }^+(\stackrel{}{p})}{p_0\omega (\stackrel{}{p})+i\epsilon }+\frac{\mathrm{\Lambda }^{}(\stackrel{}{p})}{p_0+\omega (\stackrel{}{p})i\epsilon }$$
(193)
here $`h(\stackrel{}{p})`$ is the free fermion Hamiltonian, $`\mathrm{\Lambda }^+(\stackrel{}{p})`$ and $`\mathrm{\Lambda }^{}(\stackrel{}{p})`$ were defined in Eq. (5.5).
Let us turn to derive the above kernel from the closed expression presented in Eq. (3.30). In the perturbative approximation of order $`g^2`$, only the first and second terms in Eq. (3.30) can contribute to the OGEK. In the first term which was defined in Eqs. (3.31), (3.23) and (3.21), there are four terms: two represent the self-energies of quark and antiquark and the other two are related to the t-channel one-gluon exchange interaction which is concerned here only. The $`\mathrm{\Lambda }_\mu ^a(x_2x_1,z_1)`$ and $`\mathrm{\Lambda }_\mu ^{𝐜a}(x_1x_2,z_2)`$ in Eq. (3.21) was respectively represented in Eqs. (4.5), (4.6), (4.8) and (4.9). When the vertices are taken to be the bare ones shown in Eq. (6.1), the terms contained in the $`(x_1,x_2;z_1,z_2)`$ which contributes to the OGEK can be written as
| $`H_1^t(x_1,x_2;z_1,z_2)=ig^2T^a\overline{T}^bd^4u\{\mathrm{\Delta }_{\mu \nu }^{ab}(x_2u)\gamma _2^0\gamma _2^\mu \gamma _2^0S_F(x_1u)\gamma _1^\nu S_F(uz_1)\delta ^4(x_2z_2)`$ |
| --- |
| $`+\mathrm{\Delta }_{\mu \nu }^{ab}(x_1u)\gamma _1^0\gamma _1^\mu \gamma _1^0S_F^c(x_2u)\gamma _2^\nu S_F^c(uz_2)\delta ^4(x_1z_1)\}`$ |
(194)
Next, we turn to the second term in Eq. ( 3.30) which was defined in Eqs. (3.32) and (3.14). Through a perturbative calculation of the Green’s function $`𝒢_{\mu \nu }^{ab}(x_i,z_jx_1,x_2;z_1,z_2)`$ defined in Eq. (3.14) or performing a decomposition of the Green’s function into the connected ones, one may find that there are a function $`i\mathrm{\Delta }_{\mu \nu }^{ab}(x_iz_j)S_F(x_1z_1)S_F^c(x_2z_2)`$ included in the function $`𝒢_{\mu \nu }^{ab}(x_i,z_jx_1,x_2;z_1,z_2)`$ which just is related to the t-channel OGEK when $`ij`$. Thus, according to Eq. (3.32), the term included in the $`𝒬(x_1,x_2;z_1,z_2)`$ which contributes to the OGEK may be written as:
| $`H_2^t(x_1,x_2;z_1,z_2)=ig^2T^a\overline{T}^b\{\mathrm{\Delta }_{\mu \nu }^{ab}(x_1z_2)\gamma _1^0\gamma _1^\mu S_F(x_1z_1)S_F^c(x_2z_2)\gamma _2^\nu \gamma _2^0`$ |
| --- |
| $`+\mathrm{\Delta }_{\mu \nu }^{ab}(x_2z_1)\gamma _2^0\gamma _2^\mu S_F^c(x_2z_2)S_F(x_1z_1)\gamma _1^\nu \gamma _1^0\}`$ |
(195)
Substituting Eqs. (6.7) and (6.8) into Eq. (3.30), we will obtain the expression of the t-channel OGEK $`K_t^0(XY,x,y)`$. By Fourier transformation, its expression given in the momentum space may be found to be
$$K_t^0(P,q,k)=\underset{i=1}{\overset{2}{}}H_i^t(P,q,k)S^1(P,k)$$
(196)
where
| $`H_1^t(P,q,k)=ig^2T^a\overline{T}^b\mathrm{\Delta }_{\mu \nu }^{ab}(qk)\{\gamma _2^0\gamma _2^\mu \gamma _2^0S_F(p_1)\gamma _1^\nu S_F(q_1)`$ |
| --- |
| $`+\gamma _1^0\gamma _1^\mu \gamma _1^0S_F^c(p_2)\gamma _2^\nu S_F^c(q_2)\}`$ |
(197)
and
| $`H_2^t(P,q,k)=ig^2T^a\overline{T}^b\mathrm{\Delta }_{\mu \nu }^{ab}(qk)\{\gamma _1^0\gamma _1^\mu S_F(q_1)S_F^c(p_2)\gamma _2^\nu \gamma _2^0`$ |
| --- |
| $`+\gamma _2^0\gamma _2^\mu S_F^c(q_2)S_F(p_1)\gamma _1^\nu \gamma _1^0\}`$ |
(198)
Employing the representation of fermion propagator denoted in Eqs. (6.5) and (6.6) and noticing
$$S^1(P,k)=\gamma _1^0\gamma _2^0[\widehat{S}_F(q_1)+\widehat{S}_F^c(q_2)]^1$$
(199)
one may exactly obtain from Eqs. (6.9)-(6.11) the expression denoted in Eq. (6.3). Thus, the equivalence between the both expressions of the D-S kernel derived in sections 3 and 4 is proved in the lowest order approximation.
Now let us focus on the three-dimensional t-channel OGEK which was derived for the first time in Ref. . For comparison with the four-dimensional kernel, it is necessary to give this kernel a further description based on the closed expression formulated in Eqs. (3.40)-(3.47). Analogous to the four-dimensional case, in the lowest order approximation, only the first two terms in Eq. (3.40) can contribute to the three-dimensional t-channel OGEK. Therefore, we can write
$$K_t^0(\stackrel{}{x_1},\stackrel{}{x_2};\stackrel{}{y_1},\stackrel{}{y_2};t_1t_2)=d^3z_1d^3z_2\underset{i=1}{\overset{2}{}}H_i^t(\stackrel{}{x_1},\stackrel{}{x_2};\stackrel{}{z_1},\stackrel{}{z_2};t_1t_2)S^1(\stackrel{}{z_1},\stackrel{}{z_2};\stackrel{}{y_1},\stackrel{}{y_2})$$
(200)
where $`H_1^t(\stackrel{}{x_1},\stackrel{}{x_2};\stackrel{}{z_1},\stackrel{}{z_2};t_1t_2)`$ arises from Eqs. (3.41) and (3.42) with the three-point Green’s functions in Eq. (3.42 ) being given by
| $`\mathrm{\Lambda }_\mu ^a(\stackrel{}{x_2}\stackrel{}{x_1},\stackrel{}{z_1};t_1t_2)=igd^3u𝑑u_0\mathrm{\Delta }_{\mu \nu }^{ab}(\stackrel{}{x_2}\stackrel{}{u};t_1u_0)`$ |
| --- |
| $`\times S_F(\stackrel{}{x_1}\stackrel{}{u};t_1u_0)\gamma ^\nu T^bS_F(\stackrel{}{u}\stackrel{}{z_1};u_0t_2)\}`$ |
(201)
| $`\mathrm{\Lambda }_\mu ^{𝐜a}(\stackrel{}{x_1}\stackrel{}{x_2},\stackrel{}{z_2};t_1t_2)=igd^3u𝑑u_0\mathrm{\Delta }_{\mu \nu }^{ab}(\stackrel{}{x_1}\stackrel{}{u};t_1u_0)`$ |
| --- |
| $`\times S_F^c(\stackrel{}{x_2}\stackrel{}{u};t_1u_0)\gamma ^\nu \overline{T}^bS_F^c(\stackrel{}{u}\stackrel{}{z_2};u_0t_2)\}`$ |
(202)
and $`H_2^t(\stackrel{}{x_1},\stackrel{}{x_2};\stackrel{}{z_1},\stackrel{}{z_2};t_1t_2)`$ is derived from Eqs. (3.43) and (3.44) when the terms $`i\mathrm{\Delta }_{\mu \nu }^{ab}(\stackrel{}{x_i}\stackrel{}{z_j};t_1t_2)S_F(\stackrel{}{x_1}\stackrel{}{z_1};t_1t_2)S_F^c(\stackrel{}{x_2}\stackrel{}{z_2};t_1t_2)`$ included in $`𝒢_{\mu \nu }^{ab}(\stackrel{}{x_i},\stackrel{}{z_j}\stackrel{}{x_1}_,\stackrel{}{x_2};\stackrel{}{z_1},\stackrel{}{z_2};E)`$ with $`ij`$ are taken into account only. By the Fourier transformation, it is not difficult to derive the following expression
$$K_t^0(\stackrel{}{P},\stackrel{}{q},\stackrel{}{k},E)=\underset{i=1}{\overset{2}{}}H_i^t(\stackrel{}{P},\stackrel{}{q},\stackrel{}{k},E)S^1(\stackrel{}{P},\stackrel{}{k})$$
(203)
where
| $`H_1^t(\stackrel{}{P},\stackrel{}{q},\stackrel{}{k},E)`$ |
| --- |
| $`=ig\frac{dq_0}{2\pi }\frac{dk_0}{2\pi }\{\mathrm{\Omega }_1^{a\mu }\gamma _1^0\mathrm{\Delta }_{\mu \nu }^{ab}(\stackrel{}{q}\stackrel{}{k};q_0k_0)S_F^c(\stackrel{}{p_2},q_0)\gamma ^\nu \overline{T}^bS_F^c(\stackrel{}{q_2},k_0)`$ |
| $`+\mathrm{\Omega }_2^{a\mu }\gamma _2^0\mathrm{\Delta }_{\mu \nu }^{ab}(\stackrel{}{k}\stackrel{}{q};k_0q_0)S_F(\stackrel{}{p_1},q_0)\gamma ^\nu T^bS_F(\stackrel{}{q_1},k_0)\}`$ |
(204)
and
| $`H_2^t(\stackrel{}{P},\stackrel{}{q},\stackrel{}{k},E)`$ |
| --- |
| $`=i\frac{dq_0}{2\pi }\frac{dk_0}{2\pi }\{\mathrm{\Omega }_1^{a\mu }\mathrm{\Delta }_{\mu \nu }^{ab}(\stackrel{}{q}\stackrel{}{k};Eq_0k_0)S_F(\stackrel{}{q_1},k_0)S_F^c(\stackrel{}{p_2},q_0)\overline{\mathrm{\Omega }}_2^{b\nu }`$ |
| $`+\mathrm{\Omega }_2^{a\mu }\mathrm{\Delta }_{\mu \nu }^{ab}(\stackrel{}{k}\stackrel{}{q};Ek_0q_0)S_F(\stackrel{}{p_1},q_0)S_F^c(\stackrel{}{q_2},k_0)\overline{\mathrm{\Omega }}_1^{b\nu }\}`$ |
(205)
The integrals over $`q_0`$ and $`k_0`$ can easily be calculated by applying the Cauchy theorem in the complex planes of $`q_0`$ and $`k_0`$. Since QCD is an unitary theory, the matrix element of the kernel between the spinor wave functions is independent of the gauge parameter. Therefore, we only need to show the result given in the Feynman gauge. In this gauge, noticing the representation of the gluon propagator
| $`\mathrm{\Delta }_{\mu \nu }^{ab}(Q)=\frac{\delta _{ab}g_{\mu \nu }}{Q_0^2\stackrel{}{Q}^2+i\epsilon }`$ |
| --- |
| $`=\frac{\delta _{ab}g_{\mu \nu }}{2\left|\stackrel{}{Q}\right|}[\frac{1}{Q_0\left|\stackrel{}{Q}\right|+i\epsilon }\frac{1}{Q_0+\left|\stackrel{}{Q}\right|i\epsilon }]`$ |
(206)
and the expression shown in Eqs. (6.5) and (6.6), it can be found that
| $`H_1^t(\stackrel{}{P},\stackrel{}{q},\stackrel{}{k},E)=\frac{ig^2T_1^a\overline{T}_2^a}{2\left|\stackrel{}{q}\stackrel{}{k}\right|}\{\frac{1}{\omega (\stackrel{}{p_2})+\omega (\stackrel{}{q_2})+\left|\stackrel{}{q}\stackrel{}{k}\right|}\gamma _1^0\gamma _1^\mu [\mathrm{\Lambda }^+(\stackrel{}{p_2})\gamma _2^0\gamma _{2\mu }\mathrm{\Lambda }^{}(\stackrel{}{q_2})`$ |
| --- |
| $`+\mathrm{\Lambda }^{}(\stackrel{}{p_2})\gamma _2^0\gamma _{2\mu }\mathrm{\Lambda }^+(\stackrel{}{q_2})]+\frac{1}{\omega (\stackrel{}{p_1})+\omega (\stackrel{}{q_1})+\left|\stackrel{}{q}\stackrel{}{k}\right|}`$ |
| $`\times \gamma _2^0\gamma _2^\mu [\mathrm{\Lambda }^+(\stackrel{}{p_1})\gamma _1^0\gamma _{1\mu }\mathrm{\Lambda }^{}(\stackrel{}{q_1})+\mathrm{\Lambda }^{}(\stackrel{}{p_1})\gamma _1^0\gamma _{1\mu }\mathrm{\Lambda }^+(\stackrel{}{q_1})]\}\gamma _1^0\gamma _2^0`$ |
(207)
and
| $`H_2^t(\stackrel{}{P},\stackrel{}{q},\stackrel{}{k},E)=\frac{ig^2T_1^a\overline{T}_2^a}{2\left|\stackrel{}{q}\stackrel{}{k}\right|}\{\frac{\mathrm{\Lambda }^+(\stackrel{}{p_2})\gamma _2^0\gamma _{2\mu }\gamma _1^0\gamma _1^\mu \mathrm{\Lambda }^+(\stackrel{}{q_1})}{E\omega (\stackrel{}{p_2})\omega (\stackrel{}{q_1})\left|\stackrel{}{q}\stackrel{}{k}\right|}`$ |
| --- |
| $`+\frac{\mathrm{\Lambda }^+(\stackrel{}{p_1})\gamma _1^0\gamma _1^\mu \gamma _2^0\gamma _{2\mu }\mathrm{\Lambda }^+(\stackrel{}{q_2})}{E\omega (\stackrel{}{p_1})\omega (\stackrel{}{q_2})\left|\stackrel{}{q}\stackrel{}{k}\right|}\frac{\mathrm{\Lambda }^{}(\stackrel{}{p_2})\gamma _2^0\gamma _{2\mu }\gamma _1^0\gamma _1^\mu \mathrm{\Lambda }^{}(\stackrel{}{q_1})}{E+\omega (\stackrel{}{p_2})+\omega (\stackrel{}{q_1})+\left|\stackrel{}{q}\stackrel{}{k}\right|}`$ |
| $`\frac{\mathrm{\Lambda }^{}(\stackrel{}{p_1})\gamma _1^0\gamma _1^\mu \gamma _2^0\gamma _{2\mu }\mathrm{\Lambda }^{}(\stackrel{}{q_2})}{E+\omega (\stackrel{}{p_1})+\omega (\stackrel{}{q_2})+\left|\stackrel{}{q}\stackrel{}{k}\right|}\}\gamma _1^0\gamma _2^0`$ |
(208)
It is seen that $`H_1^t(\stackrel{}{P},\stackrel{}{q},\stackrel{}{k},E)`$ is actually independent of the energy E. We would like to emphasize that the expressions in Eqs. (6.20) and (6.21) can more directly be obtained from Eqs. (6.10) and (6.11) by the following integration
$$H_i^t(\stackrel{}{P},\stackrel{}{q},\stackrel{}{k},E)=\frac{dq_0}{2\pi }\frac{dk_0}{2\pi }H_i^t(P,q,k),i=1,2.$$
(209)
Now, we discuss the inverse of the function $`S(\stackrel{}{P},\stackrel{}{k})`$ which is the Fourier transform of the function in Eq. (2.50). The equal-time propagators can be defined in such a manner
| $`S_F(\stackrel{}{x}\stackrel{}{y})=\frac{1}{2i}0^+\left|\psi (\stackrel{}{x},t)\overline{\psi }(\stackrel{}{y},t)\overline{\psi }^T(\stackrel{}{y},t)\psi ^T(\stackrel{}{x},t)\right|0^{}`$ |
| --- |
| $`=\frac{d^3p}{(2\pi )^3}S_F(\stackrel{}{p})e^{i\stackrel{}{p}(\stackrel{}{x}\stackrel{}{y})}`$ |
(210)
where
$$S_F(\stackrel{}{p})=\frac{1}{2i}\frac{h(\stackrel{}{p})}{\omega (\stackrel{}{p})}\gamma ^0$$
(211)
With this representation, the function $`S(\stackrel{}{P},\stackrel{}{k})`$ and its inverse, i.e. the three-dimensional counterparts of those in Eqs. (6.4) and (6.12) will be written as
$$S(\stackrel{}{P},\stackrel{}{k})=\frac{1}{2i}[\frac{h(\stackrel{}{q_1})}{\omega (\stackrel{}{q_1})}+\frac{h(\stackrel{}{q_2})}{\omega (\stackrel{}{q_2})}]\gamma _1^0\gamma _2^0$$
(212)
and
$$S^1(\stackrel{}{P},\stackrel{}{k})=2i\gamma _1^0\gamma _2^0[\frac{h(\stackrel{}{q_1})}{\omega (\stackrel{}{q_1})}+\frac{h(\stackrel{}{q_2})}{\omega (\stackrel{}{q_2})}]^1$$
(213)
When Eqs. (6.20), (6.21) and (6.26) are substituted into Eq. (6.16), one may write out explicitly the expression of the three-dimensional t-channel OGEK. On inserting this kernel into the first term of the effective interaction Hamiltonian denoted in Eq. (5.25) and employing the orthogonality relations of Dirac spinors and the Dirac equation, we are led to
$$V_t^{(0)}(\stackrel{}{P},\stackrel{}{q},\stackrel{}{k;E})=g^2T_1^a\overline{T}_2^a\mathrm{\Delta }(\stackrel{}{q}\stackrel{}{k};E)\overline{U}(\stackrel{}{p_1})\gamma _1^\mu U(\stackrel{}{q_1})\overline{U}(\stackrel{}{p_2})\gamma _{2\mu }U(\stackrel{}{q_2})$$
(214)
where $`U(\stackrel{}{q})`$ was represented in Eq. (5.1) and
| $`\mathrm{\Delta }(\stackrel{}{q}\stackrel{}{k};E)=\frac{1}{2\left|\stackrel{}{q}\stackrel{}{k}\right|}[\frac{1}{E\omega (\stackrel{}{p_2})\omega (\stackrel{}{q_1})\left|\stackrel{}{q}\stackrel{}{k}\right|}`$ |
| --- |
| $`+\frac{1}{E\omega (\stackrel{}{p_1})\omega (\stackrel{}{q_2})\left|\stackrel{}{q}\stackrel{}{k}\right|}]`$ |
(215)
is just the exact three-dimensional gluon propagator given in the Feynman gauge which is off-shell because the energy E is off-shell. It is noted here that the lowest order interaction Hamiltonian $`V^{(0)}(\stackrel{}{P},\stackrel{}{q},\stackrel{}{k;E})`$ is only given by the function $`H_2^t(\stackrel{}{P},\stackrel{}{q},\stackrel{}{k},E)`$ in Eq. (6.21) because the function $`H_1^t(\stackrel{}{P},\stackrel{}{q},\stackrel{}{k},E)`$ in Eq. (6.20) gives a vanishing contribution to the lowest order Hamiltonian.
### VI.2 The s-channel one-gluon exchange kernel
The four-dimensional s-channel OGEK was represented in Eq. (4.17). By means of the charge conjugation of the quark field, the vertex in Eq. (4.17) can be expressed as
$$\overline{\mathrm{\Gamma }}^{b\nu }(zz_1,z_2)_{\rho \sigma }=C_{\sigma \lambda }\mathrm{\Gamma }^{b\nu }(zz_2,z_1)_{\lambda \rho }$$
(216)
With this relation and the expression shown in Eq. (6.1), the kernel in the lowest-order approximation can be written as
| $`K_s^0(x_1,x_2;z_1,z_2)_{\alpha \beta \rho \sigma }=ig\{\mathrm{\Delta }_{\mu \nu }^{ab}(x_1z_1)(\mathrm{\Omega }_1^{a\mu })_{\alpha \gamma }S_F^{}(x_1x_2)_{\gamma \beta }`$ |
| --- |
| $`+\mathrm{\Delta }_{\mu \nu }^{ab}(x_2z_2)(\mathrm{\Omega }_2^{a\mu })_{\beta \lambda }S_F^{}(x_1x_2)_{\alpha \lambda }\}(C\gamma ^\nu T^b)_{\rho \sigma }\delta ^4(z_1z_2)`$ |
(217)
In the momentum space, it reads
$$K_s^0(P,q,k)_{\alpha \beta \rho \sigma }=ig\mathrm{\Delta }_{\mu \nu }^{ab}(P)L^{a\mu }(P,q)_{\alpha \beta }(C\gamma ^\nu T^b)_{\rho \sigma }$$
(218)
where
$$L^{a\mu }(P,q)_{\alpha \beta }=(\mathrm{\Omega }_1^{a\mu })_{\alpha \gamma }S_F^{}(p_2)_{\gamma \beta }+(\mathrm{\Omega }_2^{a\mu })_{\beta \lambda }S_F^{}(p_1)_{\alpha \lambda }$$
(219)
in which
$$S_F^{}(p)=d^4xS_F^{}(x)e^{iqx}=S_F(p)C^1=C^1S_F^c(p)^T$$
(220)
Now, let us derive the above kernel from the closed expression in Eq. (3.30). From the perturbative calculation, It can be found that in the lowest order approximation, only the second term in Eq. (3.30) can contribute to the s-channel OGEK because in the perturbative expansion of the Green’s function $`𝒢_{\mu \nu }^{ab}(x_i,z_jx_1,x_2;z_1,z_2)`$, there is a term $`i\mathrm{\Delta }_{\mu \nu }^{ab}(x_iz_j)S_F^{}(x_1x_2)\overline{S}_F^{}(z_1z_2)`$ which is merely related to the s-channel OGEK. Thus, the terms in the $`𝒬(x_1,x_2;z_1,z_2)`$ which contribute to the s-channel OGEK, according to Eq. (3.32) can be written as
$$H^s(x_1,x_2;z_1,z_2)=\underset{i,j=1}{\overset{2}{}}i\mathrm{\Delta }_{\mu \nu }^{ab}(x_iz_j)\mathrm{\Omega }_i^{a\mu }S_F^{}(x_1x_2)\overline{S}_F^{}(z_1z_2)\overline{\mathrm{\Omega }}_j^{b\nu }$$
(221)
Substituting the above expression into Eq. (3.30), in the momentum space, we have
$$K_s^0(P,q,k)=H^s(P,q,k)S^1(P,k)$$
(222)
where
$$H^s(P,q,k)_{\alpha \beta \lambda \delta }=i\mathrm{\Delta }_{\mu \nu }^{ab}(P)L^{a\mu }(P,q)_{\alpha \beta }\overline{L}^{b\nu }(P,k)_{\lambda \delta }$$
(223)
in which $`L^{a\mu }(P,q)_{\alpha \beta }`$ was given in (6.32) and
$$\overline{L}^{b\nu }(P,k)_{\lambda \delta }=[\overline{S}_F^{}(q_1)_{\lambda \tau }(\overline{\mathrm{\Omega }}_2^{a\mu })_{\tau \delta }+\overline{S}_F^{}(q_2)_{\tau \delta }(\overline{\mathrm{\Omega }}_1^{a\mu })_{\tau \lambda }]$$
(224)
here
$$\overline{S}_F^{}(q)=d^4x\overline{S}_F^{}(x)e^{iqx}=C^1S_F^c(q)=S_F^T(q)C^1$$
(225)
In light of the charge conjugation for the $`\gamma `$-matrix and for the propagators shown in Eqs. (6.33) and (6.38), it is easy to find
$$\overline{L}^{b\nu }(P,k)_{\lambda \delta }S^1(P,k)_{\lambda \delta \rho \sigma }=g(C\gamma ^\nu T^b)_{\rho \sigma }$$
(226)
With this relation, we see, the kernel in Eq. (6.35) is exactly equal to the one written in Eq. (6.31). This gives a further proof of the equivalence between the both expressions of the D-S kernel derived in sections 3 and 4. By the charge conjugation, it is not difficult to find
$$L^{a\mu }(P,q)_{\alpha \beta }=g\widehat{S}(P,q)_{\alpha \beta \lambda \tau }(C\gamma ^\mu T^a)_{\lambda \tau }$$
(227)
where
$$\widehat{S}(P,q)=\widehat{S}_F(p_1)+\widehat{S}_F^c(p_2)$$
(228)
Therefore, the kernel in Eq. (6.31) can be expressed as
$$K_s^0(P,q,k)=ig^2\mathrm{\Delta }_{\mu \nu }^{ab}(P)\widehat{S}(P,q)_{\alpha \beta \lambda \tau }(C\gamma ^\mu T^a)_{\lambda \tau }(C\gamma ^\nu T^b)_{\rho \sigma }$$
(229)
At the last of this section, we would like to discuss the three-dimensional form of the s-channel OGEK. In accordance with Eq. (3.40), this kernel is represented as
$$K_s^0(\stackrel{}{x_1},\stackrel{}{x_2};\stackrel{}{y_1},\stackrel{}{y_2};t_1t_2)=d^3z_1d^3z_2H^s(\stackrel{}{x_1},\stackrel{}{x_2};\stackrel{}{z_1},\stackrel{}{z_2};t_1t_2)S^1(\stackrel{}{z_1},\stackrel{}{z_2};\stackrel{}{y_1},\stackrel{}{y_2})$$
(230)
where $`H^s(\stackrel{}{x_1},\stackrel{}{x_2};\stackrel{}{z_1},\stackrel{}{z_2};t_1t_2)`$ can be written from Eq. (6.34) by setting $`x_1^0=x_2^0=t_1`$ and $`z_1^0=z_2^0=t_2`$ in the equal-time frame, that is
$$H^s(\stackrel{}{x_1},\stackrel{}{x_2};\stackrel{}{z_1},\stackrel{}{z_2};t_1t_2)=i\underset{i,j=1}{\overset{2}{}}\mathrm{\Delta }_{\mu \nu }^{ab}(\stackrel{}{x}_i\stackrel{}{z}_j;t_1t_2)\mathrm{\Omega }_i^{a\mu }S_F^{}(\stackrel{}{x}_1\stackrel{}{x}_2)\overline{S}_F^{}(\stackrel{}{z}_1\stackrel{}{z}_2)\overline{\mathrm{\Omega }}_j^{b\nu }$$
(231)
In the momentum space, it is of the form
$$H^s(\stackrel{}{P},\stackrel{}{q},\stackrel{}{k};E)_{\alpha \beta \lambda \delta }=i\mathrm{\Delta }_{\mu \nu }^{ab}(\stackrel{}{P},E)L^{a\mu }(\stackrel{}{P},\stackrel{}{q})_{\alpha \beta }\overline{L}^{b\nu }(\stackrel{}{P},\stackrel{}{k})_{\lambda \delta }$$
(232)
where
$$L^{a\mu }(\stackrel{}{P},\stackrel{}{q})_{\alpha \beta }=(\mathrm{\Omega }_1^{a\mu })_{\alpha \gamma }S_F^{}(\stackrel{}{p}_2)_{\gamma \beta }+(\mathrm{\Omega }_2^{a\mu })_{\beta \gamma }S_F^{}(\stackrel{}{p}_1)_{\alpha \gamma }$$
(233)
and
$$\overline{L}^{b\nu }(\stackrel{}{P},\stackrel{}{k})_{\lambda \delta }=[\overline{S}_F^{}(\stackrel{}{q}_1)_{\lambda \tau }(\overline{\mathrm{\Omega }}_2^{a\mu })_{\tau \delta }+\overline{S}_F^{}(\stackrel{}{q}_2)_{\tau \delta }(\overline{\mathrm{\Omega }}_1^{a\mu })_{\tau \lambda }]$$
(234)
which are the three-dimensional form of the functions in Eqs. (6.32) and (6.37). It is emphasized that in Eq. (6.45), only the gluon propagator is dependent on energy E, while, the fermion propagators are energy-independent. By the same charge conjugation transformations as shown in Eqs. (6.33) and (6.38), one may obtain a kernel similar to Eq. (6.42)
$$K_s^0(\stackrel{}{P},\stackrel{}{q},\stackrel{}{k};E)_{\alpha \beta \rho \sigma }=ig^2\mathrm{\Delta }_{\mu \nu }^{ab}(\stackrel{}{P},E)\widehat{S}(\stackrel{}{P},\stackrel{}{q})_{\alpha \beta \lambda \tau }(C\gamma ^\mu T^a)_{\lambda \tau }(C\gamma ^\nu T^b)_{\rho \sigma }$$
(235)
which may also be represented as
$$K_s^0(\stackrel{}{P},\stackrel{}{q},\stackrel{}{k};E)_{\alpha \beta \rho \sigma }=ig^2\mathrm{\Delta }_{\mu \nu }^{ab}(\stackrel{}{P},E)S(\stackrel{}{P},\stackrel{}{q})_{\alpha \beta \lambda \tau }(\gamma ^\mu CT^a)_{\lambda \tau }(C\gamma ^\nu T^b)_{\rho \sigma }$$
(236)
where
$$S(\stackrel{}{P},\stackrel{}{q})=\widehat{S}(\stackrel{}{P},\stackrel{}{q})\gamma _1^0\gamma _2^0$$
(237)
which is the three-dimensional form of Eq. (6.4).
In the P-S equation, similar to Eq. (6.27), the lowest order interaction Hamiltonian given by the kernel in Eq. (6.48), according to Eq. (5.25), will be written as
$$V_s^{(0)}(\stackrel{}{P},\stackrel{}{q},\stackrel{}{k};E)=g^2\mathrm{\Delta }_{\mu \nu }^{ab}(\stackrel{}{P},E)U_\alpha (\stackrel{}{p_1})^+U_\beta (\stackrel{}{p_2})^+(C\gamma ^\mu T^a)_{\alpha \beta }(C\gamma ^\nu T^b)_{\rho \sigma }U_\rho (\stackrel{}{q_1})U_\sigma (\stackrel{}{q_2})$$
(238)
It should be noted that the positive energy state Dirac spinors used here were defined in Eq. (5.1). The negative energy state spinor may be given by the charge conjugation relation: $`V(\stackrel{}{p})=CU(\stackrel{}{p})`$ here $`C=\gamma ^5\gamma ^0`$ . However, the matrix $`C`$ in Eq. (6.51) is defined by $`C=i\gamma ^2\gamma ^0`$ . Correspondingly, the charge conjugation relation between the spinor wave functions is given by $`v^s(\stackrel{}{p})=C\overline{u}^s(\stackrel{}{p})^T`$ where $`\overline{u}^s(\stackrel{}{p})=u^s(\stackrel{}{p})^+\gamma ^0`$ with $`u^s(\stackrel{}{p})`$ and $`v^s(\stackrel{}{p})`$ being the positive and negative energy spinor wave functions respectively and represented as
$$u^s(\stackrel{}{p})=\stackrel{~}{U}(\stackrel{}{p})\phi ^s(\stackrel{}{p}),v^s(\stackrel{}{p})=\stackrel{~}{V}(\stackrel{}{p})\chi ^s(\stackrel{}{p})$$
(239)
here $`\phi ^s(\stackrel{}{p})`$ and $`\chi ^s(\stackrel{}{p})`$ are the spin wave functions and
$$\stackrel{~}{U}(\stackrel{}{p})=\sqrt{\frac{\omega }{m}}U(\stackrel{}{p}),\stackrel{~}{V}(\stackrel{}{p})=\sqrt{\frac{\omega }{m}}\gamma ^0V(\stackrel{}{p})$$
(240)
Usually, the S-matrix element given by the kernel in Eq. (6.48) is represented by
$$T_s(\stackrel{}{P},\stackrel{}{q},\stackrel{}{k;E})=u_\alpha ^{s_1}(\stackrel{}{p_1})^+u_\beta ^{s_2}(\stackrel{}{p_2})^+K_s^0(\stackrel{}{P},\stackrel{}{q},\stackrel{}{k};E)_{\alpha \beta \rho \sigma }u_\rho ^{r_1}(\stackrel{}{q_1})u_\sigma ^{r_2}(\stackrel{}{q_2})$$
(241)
For later derivation, it is more convenient to use the expression of the kernel written in Eq. (6.49). On inserting Eq. (6.49) into Eq. (6.54) and noticing
$$u_\alpha ^{s_1}(\stackrel{}{p_1})^+u_\beta ^{s_2}(\stackrel{}{p_2})^+S(\stackrel{}{P},\stackrel{}{q})_{\alpha \beta \lambda \tau }=i\overline{u}_\lambda ^{s_1}(\stackrel{}{p_1})\overline{u}_\tau ^{s_2}(\stackrel{}{p_2})$$
(242)
which is obtained by applying the Dirac equation and
$$\begin{array}{c}\overline{u}_\beta ^{s_2}(\stackrel{}{p_2})(C\gamma ^\mu T^a)_{\alpha \beta }=(\gamma ^\mu T^a)_{\alpha \beta }v_\beta ^{s_2}(\stackrel{}{p_2}),\\ (C\gamma ^\nu T^b)_{\rho \sigma }u_\rho ^{r_1}(\stackrel{}{q_1})=\overline{v}_\rho ^{r_1}(\stackrel{}{q_1})(\gamma ^\nu T^b)_{\rho \sigma },\end{array}$$
(243)
one can get
$$T_s(\stackrel{}{P},\stackrel{}{q},\stackrel{}{k;E})=g^2\mathrm{\Delta }_{\mu \nu }^{ab}(\stackrel{}{P},E)\overline{u}^{s_1}(\stackrel{}{p_1})\gamma ^\mu T^av^{s_2}(\stackrel{}{p_2})\overline{v}^{r_1}(\stackrel{}{q_1})\gamma ^\nu T^bu^{r_2}(\stackrel{}{q_2})$$
(244)
This just is the S-matrix element for the one-gluon exchange interaction taking place in the s-channel. It is easy to verify that the above matrix element is independent of the gauge parameter. Therefore, we only need to work in the Feynman gauge. In this gauge,
$$\mathrm{\Delta }_{\mu \nu }^{ab}(\stackrel{}{P},E)=\frac{\delta _{ab}g_{\mu \nu }}{E^2\stackrel{}{P}^2+i\epsilon }$$
(245)
With this propagator, as shown in Ref. , by the charge conjugation and Fierz transformation, Eq. (6.57) can be represented in the form
$$T_s(\stackrel{}{P},\stackrel{}{q},\stackrel{}{k;E})=\left(\frac{\omega _1\omega _2}{m_1m_2}\right)^{^{1/2}}\phi _{s_1}^+\phi _{s_2}^+V_s^{(0)}(\stackrel{}{P},\stackrel{}{q},\stackrel{}{k;E})\phi _{r_1}\phi _{r_2}$$
(246)
where
$$V_s^{(0)}(\stackrel{}{P},\stackrel{}{q},\stackrel{}{k;E})=\frac{g^2\widehat{C}_s\widehat{F}_s}{E^2\stackrel{}{P}^2+i\epsilon }\overline{U}(\stackrel{}{p_1})\overline{U}(\stackrel{}{p_2})\mathrm{\Gamma }_{12}U(\stackrel{}{q_1})U(\stackrel{}{q_2})$$
(247)
is the interaction Hamiltonian occurring in the P-S equation in which the spinor is still defined in Eq. (5.1), $`\widehat{C}_s`$ is the color matrix
$$\widehat{C}_s=\frac{1}{24}(\lambda _1^a\lambda _2^a)^2$$
(248)
with $`\lambda _i^a`$ being the Gell-Mann matrices, $`\widehat{F}_s`$ is the flavor matrix which has an expression for flavor SU(2) such that
$$\widehat{F}_s=\frac{1}{2}(1\stackrel{}{\tau _1}\stackrel{}{\tau _2})$$
(249)
here $`\stackrel{}{\tau _i}`$ are isospin Pauli matrices and
$$\mathrm{\Gamma }_{12}=I_1I_2+\gamma _1^5\gamma _2^5\frac{1}{2}\gamma _1^\mu \gamma _{2\mu }+\frac{1}{2}(\gamma _1^5\gamma _1^\mu )(\gamma _2^5\gamma _{2\mu })$$
(250)
In the end, it is pointed out that since the matrix element of the color operator $`\widehat{C}_s`$ in the $`q\overline{q}`$ color singlet vanishes, the s-channel OGEK contributes nothing to the $`q\overline{q}`$ bound states. However, for many-quark-antiquark systems such as $`\pi \pi `$, $`K\overline{K}`$, $`\pi N`$, $`KN`$ systems and etc., the contribution of the s-channel OGEK is not negligible and plays an important role to the interations taking place in those systems.\[36- 38\].
## VII Discussions and remarks
In this paper, the D-S equation satisfied by the $`q\overline{q}`$ bound states has been derived from QCD and, especially, the interaction kernel in the equation has been given two equivalent closed expressions which were respectively derived by making use of the equations of motion obeyed by the Green’s functions and the irreducible decomposition of the Green’s functions. Since the B-S equation is commonly viewed as the correct equation for the bound state problem, it is natural to ask what is the relation between the D-S equation and the B-S equation? As shown in Ref. , the B-S equation may be derived from the D-S equation. In fact, when applying the Lehmann representation in Eq. (2.36) to Eqs. (2.16) and (2.17), by the same procedure stated in section.2, one may obtain two D-S equations as follows
$$(i_{x_1}m_1)\chi _{P\varsigma }(x_1,x_2)=d^4y_1d^4y_2\widehat{K}_1(x_1,x_2;y_1,y_2)\chi _{P\varsigma }(y_1,y_2)$$
(251)
$$(i_{x_2}m_2)\chi _{P\varsigma }(x_1,x_2)=d^4y_1d^4y_2\widehat{K}_2(x_1,x_2;y_1,y_2)\chi _{P\varsigma }(y_1,y_2)$$
(252)
where
$$\widehat{K}_i(x_1,x_2;y_1,y_2)=\gamma _i^0K_i(x_1,x_2;y_1,y_2),\text{ }i=1,2$$
(253)
in which $`K_i(x_1,x_2;y_1,y_2)`$ was given in Eq. (3.28). The D-S equations shown in Eqs. (2.38) and (2.39) may directly be written out from Eqs. (7.1) and (7.2).
Operating on Eq. (7.1) with $`(i_{x_2}m_2)`$ or on Eq. (7.2) with $`(i_{x_1}m_1)`$, we have
$$(i_{x_1}m_1)(i_{x_2}m_2)\chi _{P\varsigma }(x_1,x_2)=d^4y_1d^4y_2K_B(x_1,x_2;y_1,y_2)\chi _{P\varsigma }(y_1,y_2)$$
(254)
where
$$K_B(x_1,x_2;y_1,y_2)=(i_{x_1}m_1)\widehat{K}_2(x_1,x_2;y_1,y_2)=(i_{x_2}m_2)\widehat{K}_1(x_1,x_2;y_1,y_2)$$
(255)
is the B-S interaction kernel whose explicit expression was derived in Ref. . Acting on Eq. (7.4) with the inverse of the operator $`(i_{x_1}m_1)(i_{x_2}m_2)`$, the B-S equation will be recast in an integral equation
$$\chi _{P\varsigma }(x_1,x_2)=d^4y_1d^4y_2d^4z_1d^4z_2S_F^{(0)}(x_1z_1)S_F^{c(0)}(x_2z_2)K_B(z_1,z_2;y_1,y_2)\chi _{P\varsigma }(y_1,y_2)$$
(256)
where $`S_F^{(0)}(x_1z_1)`$ and $`S_F^{c(0)}(x_2z_2)`$ are the free propagators of quark and antiquark respectively. In the momentum space, it becomes
$$\chi _{P\varsigma }(q)=S_F^{(0)}(p_1)S_F^{c(0)}(p_2)\frac{d^4k}{(2\pi )^4}K_B(P,q,k)\chi _{P\varsigma }(k)$$
(257)
Conversely, if we act on Eq. (7.6) with the operators $`(i_{x_2}m_2)`$ and $`(i_{x_1}m_1)`$ respectively, the D-S equations in Eqs. (7.1) and (7.2) will be recovered. This shows that to get the B-S equation, we need merely to consider the D-S equation. It should be noted that the four-dimensional B-S equation can not directly be transformed into the three-dimensional D-S equation. In order to obtain a three-dimensional equation from the four-dimensional B-S equation, it is necessary to introduce a certain constraint condition on the relative time (or relative energy) as was done in an approximate manner such as the instantaneous approximation or the quasipotential approaches \[7-12\].
As stated above, the four-dimensional D-S equation and the corresponding B-S equation can be derived from one another. But, this does not mean that the D-S equation and the B-S equation are fully equivalent to each other, similar to the Dirac equation and the K-G equation which can also be derived from each other. As seen from Eq. (7.4), the B-S equation is a kind of second-order differential equation in the position space. Therefore, a solution to the equation depends on not only the amplitude at time origin, but also the time-differential of the amplitude at the time origin as in the case for K-G equation. This probably is the origin that causes the B-S equation to have the unphysical solutions of negative norm. In order to exclude the unphysical solutions, as mentioned before, the common procedure is to recast the four-dimensional B-S equation in a three-dimensional form by eliminating the relative time (or the relative energy) from the equation. Certainly, the three-dimensional equation, particularly, the exact version of the equation presented in the sections 2 and 3 is much convenient to use in solving the eigenvalue problem. However, since the three-dimensional equation loses the Lorentz-covariance of a relativistic dynamics, it is sometimes not suitable for carrying out extensive theoretical analyses, for instance, to perform the irreducible decomposition of the Green’s functions contained in the kernel given in Eq. (3.30). The decomposition can readily be done in the four-dimensional form as shown in section 4. At this point, we may ask whether the relativistic bound state problem can be solved Lorentz-covariantly in the Minkowski space without occurrence of the unphysical solutions? The answer should be positive because the Lorentz-covariance of the equation implies that one may work in any Lorentz frame and gets the same result. Let us turn to the D-S equations shown in Eqs. (2.41) and (2.42) which are represented in the position space. Clearly, the equation in Eq.(2.42) is a first-order differential equation of Schrödinger-type. One may first solve this equation to get an amplitude which describes the evolution of the amplitude with the relative time $`\tau `$. and then substitute this amplitude into Eq.(2.41) to solve the eigenvalue $`E`$ and the amplitude $`\chi _{P\varsigma }(x)`$. In solving these equations, we only need the initial conditions of the amplitude at the time origin without concerning the initial conditions of the time-differentials of the amplitude. Therefore, the unphysical solutions would not appear in this case. In this sense, we can say, the D-S equation derived in this paper, as it provides a new formulation of the relativistic equation for the two fermion bound system, gives a suitable prescription to solve the four-dimensional equation. Moreover, based on the relation denoted in Eq. (7.5), the B-S kernel may conveniently be evaluated from the D-S kernel. In comparison with the closed expression of the B-S kernel derived in Ref. , the D-S kernel shown in Eq. (3.30) is rather simpler. The main contribution to the D-S kernel is given by the Green’s function $`𝒢_{\mu \nu }^{ab}(x_i,y_jx_1,x_2;y_1,y_2)`$ written in Eq. (3.14). While, the B-S kernel concerns more complicated Green’s functions such as
| $`𝒢_{\mu \nu \lambda \tau }^{abcd}(x_1,x_2,y_1,y_2x_1,x_2;y_1,y_2)`$ |
| --- |
| =$`0^+|T\{N[𝐀_\mu ^a(x_1)𝐀_\nu ^b(x_2)\psi (x_1)\psi ^c(x_2)]N[𝐀_\lambda ^c(y_1)𝐀_\tau ^d(y_2)\overline{\psi }(y_1)\overline{\psi }^c(y_2)]|0^{}`$ |
(258)
which gives the major contribution to the B-S kernel. In particular, in comparison of the four-dimensional kernel represented in Eqs. (3.30)-(3.33) with the three-dimensional counterpart written in Eqs. (3.40-)-(3.43), we see, there is an one-to-one correspondence between the both kernels. Therefore, to calculate the three-dimensional kernel, one may first calculate the four-dimensional one and then convert it to the three-dimensional form according the correspondence relation between the both of them. Since the four-dimensional D-S equation is Lorentz-covariant, its kernel can conveniently be analyzed and calculated by means of the familiar technique developed in the covariant quantum field theory.
In the end, we would like to address that unlike the Dyson-Schwinger equation which contains an infinite set of equations, the D-S equation derived in this paper is of a closed form with a closed expression of the kernel as given in section 3 or section 4. The kernel can easily be calculated by the perturbation method. For example, in the perturbative calculation of the kernel given in section 3 which contains only a few types of Green’s functions, we only need the familiar perturbative expansions of the Green’s functions without concerning the calculation of other more-point Green’s functions as it is necessary to be done for the Dyson-Schwinger equation. Especially, each of the Green’s functions can be represented in the form of functional integral and is possible to be evaluated by a nonperturbative method as suggested, for example, by the lattice gauge theory. Therefore, the expression of the kernel given in this paper provides a new formalism for exploring the QCD nonperturbative effect and the quark confinement which are important for the formation of a $`q\overline{q}`$ bound state. In the ordinary quark potential model , the quark confinement is usually simulated by a linear potential which was suggested by the lattice computation of a Wilson loop and by the area law \[42-44\]. Obviously, this simulation is oversimplified. For the purpose of investigating the quark confinement, it is appropriate to start from the kernel given in this paper for the case that the quark and the antiquark have different flavors. In this case, all the Green’s functions become the ordinary ones as shown in Eqs. (2.5), (2.13) and (3.11). Since the kernel derived in this paper contains all the interactions taking place in the bound state and includes the color-spin matrices $`\mathrm{\Omega }_i^{a\mu }`$ defined in Eq. (2.23) in it, it is anticipated that a nonperturbative calculation of this kernel would give a sophisticated confining potential which includes not only its spatial form, but also its spin and color structures. This just is the advantage of the formalism of D-S equation presented in this paper.
## VIII Acknowledgment
This project was supposed by National Natural Science Foundation of China.
## IX Appendix A: Derivation of equations of motion satisfied by the Green’s functions
This appendix is used to derive the equations of motion satisfied by the quark-antiquark two and four-point Green’s functions. These equations may be derived from the following QCD generating functional .
$$Z[J,\overline{\eta },\eta ,\overline{\xi },\xi ]=\frac{1}{N}𝒟(A,\overline{\psi },\psi ,\overline{C},C)e^{iI}$$
(259)
where
$$I=d^4x[+J^{a\mu }A_\mu ^a+\overline{\eta }\psi +\overline{\psi }\eta +\overline{\xi }C+\overline{C}\xi ]$$
(260)
in which $``$ is the effective Lagrangian of QCD
$$=\overline{\psi }(im+g𝐀)\psi \frac{1}{4}F^{a\mu \nu }F_{\mu \nu }^a\frac{1}{2\alpha }(^\mu A_\mu ^a)^2+\overline{C}^a^\mu (D_\mu ^{ab}C^b)$$
(261)
here $`𝐀=\gamma ^\mu T^aA_\mu ^a`$ with $`A_\mu ^a`$ being the vector potentials of gluon fields,
$$F_{\mu \nu }^a=_\mu A_\nu ^a_\nu A_\mu ^a+gf^{abc}A_\mu ^bA_\nu ^c$$
(262)
are the strength tensors of the gluon field,
$$D_\mu ^{ab}=\delta ^{ab}_\mu gf^{abc}A_\mu ^c$$
(263)
are the covariant derivatives, $`\overline{C}^a,C^b`$ represent the ghost fields, and $`J^{a\mu },\overline{\eta },\eta ,\overline{\xi }`$ and $`\xi `$ denote the external sources coupled to the gluon, quark and ghost fields respectively. By the charge conjugation transformations shown in Eq. (2.2) for the quark fields and in the following for the external sources
$$\eta ^c=C\overline{\eta }^T,\overline{\eta }^c=\eta ^TC^1$$
(264)
it is easy to prove the relation
$$\overline{\psi }(im+g𝐀)\psi +\overline{\eta }\psi +\overline{\psi }\eta =\overline{\psi }^c(im+g\overline{𝐀})\psi ^c+\overline{\eta }^c\psi ^c+\overline{\psi }^c\eta ^c$$
(265)
where $`\overline{𝐀}=\gamma ^\mu \overline{T}^aA_\mu ^a`$ .
### IX.1 Equations of motion with respect to the coordinate $`x_1`$
Upon taking the functional derivative of the generating functional in Eq. (A.1) with respect to the field function $`\overline{\psi }_\alpha (x_1)`$ and considering
$$\frac{\delta Z}{\delta \overline{\psi }_\alpha (x_1)}=0$$
(266)
and
$$\frac{\delta I}{\delta \overline{\psi }_\alpha (x_1)}=\eta _\alpha (x_1)+[(i_{x_1}m_1)_{\alpha \gamma }+g𝐀(x_1)_{\alpha \gamma }]\psi _\gamma (x_1)$$
(267)
it can be found that
$$\{\eta _\alpha (x_1)+[(i_{x_1}m_1)_{\alpha \gamma }+(\mathrm{\Gamma }^{a\mu })_{\alpha \gamma }\frac{\delta }{i\delta J^{a\mu }(x_1)}]\frac{\delta }{i\delta \overline{\eta }_\gamma (x_1)}\}Z=0$$
(268)
where the fields $`A_\mu ^a(x_1)`$ and $`\psi _\gamma (x_1)`$ have been replaced by the derivatives of the generating functional with respect to the sources $`J^{a\mu }(x_1)`$ and $`\overline{\eta }_\gamma (x_1)`$ and each of the subscripts $`\alpha ,\beta `$ and $`\gamma `$ marks the components of color, flavor and spinor. Differentiating Eq. (A.10) with respect to the source $`\eta _\rho (y_1)`$ and then setting all the sources to vanish, we obtain the equation satisfied by the quark propagator
$$[(i_{x_1}m_1+\mathrm{\Sigma })S_F]_{\alpha \rho }(x_1,y_1)=\delta _{\alpha \rho }\delta ^4(x_1y_1)$$
(269)
where
| $`(\mathrm{\Sigma }S_F)_{\alpha \rho }(x_1,y_1)d^4z_1\mathrm{\Sigma }(x_1,z_1)_{\alpha \gamma }S_F(z_1y_1)_{\gamma \rho }`$ |
| --- |
| $`=(\mathrm{\Gamma }^{a\mu })_{\alpha \gamma }\mathrm{\Lambda }_\mu ^a(x_1x_1,y_1)_{\gamma \beta }`$ |
(270)
here $`\mathrm{\Sigma }(x_1,z_1)`$ stands for the quark proper self-energy and $`\mathrm{\Lambda }_\mu ^a(x_1x_1,y_1)_{\gamma \rho }`$ was defined in the first equality of Eq. (3.11).
Let us turn to derive the equation of motion satisfied by the Green’s function defined in Eq. (2.1). In doing this, we need first to derive the equations of motion obeyed by the Green’s function defined in Eq. (2.5). By successively differentiating Eq. (A.10) with respect to the sources $`\overline{\eta }_\beta ^c(x_2),\eta _\rho (y_1)`$ and $`\eta _\sigma ^c(y_2)`$, noticing the equality in Eq. (A.7) and the following nonvanishing derivatives
| $`\frac{\delta \eta _\alpha ^c(x)}{\delta \overline{\eta }_\beta (y)}=C_{\alpha \beta }\delta ^4(xy),\frac{\delta \overline{\eta }_\alpha ^c(x)}{\delta \eta _\beta (y)}=(C^1)_{\alpha \beta }\delta ^4(xy),`$ |
| --- |
| $`\frac{\delta \overline{\eta }_\alpha (x)}{\delta \eta _\beta ^c(y)}=(C^1)_{\alpha \beta }\delta ^4(xy),\frac{\delta \eta _\alpha (x)}{\delta \overline{\eta }_\beta ^c(y)}=C_{\alpha \beta }\delta ^4(xy)`$ |
(271)
we have
| $`\{C_{\alpha \beta }\delta ^4(x_1x_2)\frac{\delta }{i^3\delta \eta _\rho (y_1)\delta \eta _\sigma ^c(y_2)}+\delta _{\alpha \rho }\delta ^4(x_1y_1)\frac{\delta ^2}{i\delta \overline{\eta }_\beta ^c(x_2)\delta \eta _\sigma ^c(y_2)}`$ |
| --- |
| $`+\eta _\alpha (x_1)\frac{\delta ^3}{i\delta \overline{\eta }_\beta ^c(x_2)\delta \eta _\rho (y_1)\delta \eta _\sigma ^c(y_2)}[(i_{x_1}m_1)_{\alpha \gamma }+(\mathrm{\Gamma }^{a\mu })_{\alpha \gamma }\frac{\delta }{i\delta J^{a\mu }(x_1)}]`$ |
| $`\times \frac{\delta ^4}{\delta \overline{\eta }_\gamma (x_1)\delta \overline{\eta }_\beta ^c(x_2)\delta \eta _\rho (y_1)\delta \eta _\sigma ^c(y_2)}\}Z=0`$ |
(272)
When all the sources are set to be zero, one immediately derives Eq. (2.10) from Eq. (A.14).
It is noted that the first equation in Eq. (2.14) may directly derived from Eq. (A.11) by the charge conjugation transformation represented in Eq. (2.2) or by differentiating Eq. (A.10) with respect to the source $`\overline{\eta }_\beta ^c(x_2)`$.
### IX.2 Equations of motion with respect to the coordinate $`x_2`$
When taking the derivative of the generating functional in Eq. (A.1) with respect to the field variable $`\overline{\psi }_\beta ^c(x_2)`$ and noticing the relation in Eq. (A.7), by the same procedure as described in Eqs. (A.8)-(A.10), one may obtain
$$\{\eta _\beta ^c(x_2)+[(i_{x_2}m_2)_{\beta \lambda }+(\overline{\mathrm{\Gamma }}^{b\nu })_{\beta \lambda }\frac{\delta }{i\delta J^{b\nu }(x_2)}]\frac{\delta }{i\delta \overline{\eta }_\lambda ^c(x_2)}\}Z=0$$
(273)
Differentiating Eq. (A.15) with respect to $`\eta _\sigma ^c(y_2)`$ and then setting all the sources to be zero, one can get the equation for the antiquark propagator
$$[(i_{x_2}m_2+\mathrm{\Sigma }^c)S_F^c]_{\beta \sigma }(x_2,y_2)=\delta _{\beta \sigma }\delta ^4(x_2y_2)$$
(274)
where
| $`(\mathrm{\Sigma }^cS_F^c)_{\beta \sigma }(x_2,y_2)d^4z_2\mathrm{\Sigma }^c(x_2,z_2)_{\beta \lambda }S_F^c(z_2y_2)_{\lambda \sigma }`$ |
| --- |
| $`=(\overline{\mathrm{\Gamma }}^{b\nu })_{\beta \lambda }\mathrm{\Lambda }_\nu ^{cb}(x_2x_2,y_2)_{\lambda \sigma }`$ |
(275)
here $`\mathrm{\Sigma }^c(x_2,z_2)`$ is the antiquark proper self-energy and $`\mathrm{\Lambda }_\nu ^{𝐜b}(x_2x_2,y_2)_{\lambda \sigma }`$ was represented in the second equality of Eq. (3.12).
Similarly, Upon differentiating Eq. (A.15) with respect to the sources $`\overline{\eta }_\alpha (x_1),\eta _\rho (y_1)`$ and $`\eta _\sigma ^c(y_2)`$, one gets
| $`\{C_{\alpha \beta }\delta ^4(x_1x_2)\frac{\delta }{i^3\delta \eta _\rho (y_1)\delta \eta _\sigma ^c(y_2)}+\delta _{\beta \sigma }\delta ^4(x_2y_2)\frac{\delta ^2}{i\delta \overline{\eta }_\alpha (x_1)\delta \eta _\rho (y_2)}`$ |
| --- |
| $`+\eta _\beta ^c(x_2)\frac{\delta ^3}{i\delta \overline{\eta }_\alpha (x_1)\delta \eta _\rho (y_1)\delta \eta _\sigma ^c(y_2)}[(i_{x_2}m_2)_{\beta \lambda }+(\overline{\mathrm{\Gamma }}^{b\nu })_{\beta \lambda }\frac{\delta }{i\delta J^{b\nu }(x_2)}]`$ |
| $`\times \frac{\delta ^4}{\delta \overline{\eta }_\alpha (x_1)\delta \overline{\eta }_\lambda ^c(x_2)\delta \eta _\rho (y_1)\delta \eta _\sigma ^c(y_2)}\}Z=0`$ |
(276)
Setting all the sources to vanish, we directly obtain Eq. (2.11) from the above equation.
It is mentioned that the second equation in Eq. (2.14) may directly be derived from Eq. (A.16) by the charge conjugation transformation or by differentiating Eq. (A.15) with respect to the source $`\overline{\eta }_\alpha (x_1)`$.
### IX.3 Equations of motion with respect to the coordinate $`y_1`$
By taking the derivative of the generating functional in Eq. (A.1) with respect to the field variable $`\psi _\rho (y_1),`$ following the same procedure as deriving Eq. (A.10), one may get
$$\{\overline{\eta }_\rho (y_1)+\frac{\delta }{i\delta \eta _\tau (y_1)}[(i\stackrel{}{}_{y_1}+m_1)_{\tau \rho }(\mathrm{\Gamma }^{a\mu })_{\tau \rho }\frac{\delta }{i\delta J^{a\mu }(y_1)}]\}Z=0$$
(277)
On differentiating the above equation with respect to $`\overline{\eta }_\alpha (x_1)`$ and then turning off all the sources, we arrive at
$$[S_F(i\stackrel{}{}_{y_1}+m_1\mathrm{\Sigma })]_{\alpha \rho }(x_1,y_1)=\delta _{\alpha \rho }\delta ^4(x_1y_1)$$
(278)
where
| $`(S_F\mathrm{\Sigma })_{\alpha \rho }(x_1,y_1)d^4z_1S_F(x_1z_1)_{\alpha \tau }\mathrm{\Sigma }(z_1,y_1)_{\tau \rho }`$ |
| --- |
| $`=\mathrm{\Lambda }_\mu ^a(y_1x_1,y_1)_{\alpha \tau }(\mathrm{\Gamma }^{a\mu })_{\tau \rho }`$ |
(279)
If we differentiate Eq. (A.19) with respect to $`\eta _\sigma ^c(y_2)`$, after letting the sources to be vanishing, we get an equation satisfied by the propagator $`\overline{S}_F^{}(y_1y_2)`$ as written in the first equation in Eq. (3.3).
Now let us differentiate Eq. (A.19) with respect to $`\overline{\eta }_\alpha (x_1),\overline{\eta }_\beta ^c(x_2)`$ and $`\eta _\sigma ^c(y_2)`$ and then set all the sources but the source $`J`$ to be zero. By these operations, we get
| $`G(x_{1,}x_2;y_1,y_2)_{\alpha \beta \tau \sigma }^J(i\stackrel{}{}_{y_1}+m_1)_{\tau \rho }=\delta _{\alpha \rho }\delta ^4(x_1y_1)S_F^c(x_2y_2)_{\beta \sigma }^J`$ |
| --- |
| $`C_{\rho \sigma }\delta ^4(y_1y_2)S_F^{}(x_1x_2)_{\alpha \beta }^J+G_\mu ^a(y_1x_1,x_2;y_1,y_2)_{\alpha \beta \tau \sigma }^J(\mathrm{\Gamma }^{a\mu })_{\tau \rho }`$ |
(280)
where
$$G(x_{1,}x_2;y_1,y_2)^J=\frac{\delta ^4Z[J,\overline{\eta },\eta ,\overline{\xi },\xi ]}{\delta \overline{\eta }(x_1)\delta \overline{\eta }^c(x_2)\delta \eta (y_1)\delta \eta ^c(y_2)}_{\overline{\eta }=\overline{\eta }^c=\eta =\eta ^c=0}$$
(281)
$$S_F^c(x_2y_2)_{\beta \sigma }^J=\frac{\delta ^2Z[J,\overline{\eta },\eta ,\overline{\xi },\xi ]}{i\delta \overline{\eta }_\beta ^c(x_2)\delta \eta _\sigma ^c(y_2)}_{\overline{\eta }=\eta =\overline{\xi }=\xi =0}$$
(282)
$$S_F^{}(x_1x_2)_{\alpha \beta }^J==\frac{\delta ^2Z[J,\overline{\eta },\eta ,\overline{\xi },\xi ]}{i^3\delta \overline{\eta }_\alpha (x_1)\delta \overline{\eta }_\beta ^c(x_2)}_{\overline{\eta }=\eta =\overline{\xi }=\xi =0}$$
(283)
and
$$G_\mu ^a(y_ix_{1,}x_2;y_1,y_2)^J=\frac{\delta }{i\delta J^{a\mu }(y_i)}G(x_{1,}x_2;y_1,y_2)^J,i=1,2$$
(284)
Once we set $`J=0`$, Eq. (A.22) will give rise to Eq. (3.1). Furthermore, if we differentiate Eq. (A.22) with respect to $`J^{a\mu }(x_i)`$ and subsequently set $`J=0`$, noticing
$$G_{\mu \nu }^{ab}(x_i,y_jx_{1,}x_2;y_1,y_2)=\frac{\delta }{i\delta J^{a\mu }(x_i)i\delta J^{b\nu }(y_j)}G(x_{1,}x_2;y_1,y_2)^J_{J=0}$$
(285)
the equations in Eq. (3.8) will immediately be derived.
### IX.4 Equations of motion with respect to the coordinate $`y_2`$
To derive the equations of motion with respect to $`y_2,`$ we need to differentiate the generating functional with respect to the field $`\psi _\sigma ^c(y_2).`$ By the same procedure as formulated in Eqs. (A.8)-(A.10), we get
$$\{\overline{\eta }_\sigma ^c(y_2)+\frac{\delta }{i\delta \eta _\delta ^c(y_2)}[(i\stackrel{}{}_{y_2}+m_2)_{\delta \sigma }(\overline{\mathrm{\Gamma }}^{b\nu })_{\delta \sigma }\frac{\delta }{i\delta J^{b\nu }(y_2)}]\}Z=0$$
(286)
Then, the differentiation of the above equation with respect to $`\overline{\eta }_\beta ^c(x_2)`$ with setting all the sources to vanish subsequently will lead us to
$$[S_F^c(i\stackrel{}{}_{y_2}+m_2\mathrm{\Sigma }^c)]_{\beta \sigma }(x_2,y_2)=\delta _{\beta \sigma }\delta ^4(x_2y_2)$$
(287)
where
| $`(S_F^c\mathrm{\Sigma }^c)_{\beta \sigma }(x_2,y_2)d^4z_2S_F^c(x_2z_2)_{\beta \delta }\mathrm{\Sigma }^c(z_2,y_2)_{\delta \sigma }`$ |
| --- |
| $`=\mathrm{\Lambda }_\nu ^{𝐜b}(y_2x_2,y_2)_{\beta \delta }(\overline{\mathrm{\Gamma }}^{b\nu })_{\delta \sigma }`$ |
(288)
Upon differentiating Eq. (A.28) with respect to $`\eta _\rho (y_1)`$ and then turning off all the sources, one can obtain the second equation in Eq. (3.3) for the propagator $`\overline{S}_F^{}(y_1y_2)`$
Let us differentiate Eq. (A.28) with respect to $`\overline{\eta }_\alpha (x_1),\overline{\eta }_\beta ^c(x_2)`$ and $`\eta _\rho (y_1)`$ and set all the sources except for the $`J`$ to vanish. As a result, we get
| $`G(x_{1,}x_2;y_1,y_2)_{\alpha \beta \rho \delta }^J(i\stackrel{}{}_{y_2}+m_2)_{\delta \sigma }=\delta _{\beta \sigma }\delta ^4(x_2y_2)S_F(x_1y_1)_{\alpha \rho }^J`$ |
| --- |
| $`C_{\rho \sigma }\delta ^4(y_1y_2)S_F^{}(x_1x_2)_{\alpha \beta }^J+G_\nu ^b(y_2x_1,x_2;y_1,y_2)_{\alpha \beta \rho \delta }^J(\overline{\mathrm{\Gamma }}^{b\nu })_{\delta \sigma }`$ |
(289)
where
$$S_F(x_1y_1)_{\gamma \rho }^J=\frac{\delta ^2Z[J,\overline{\eta },\eta ,\overline{\xi },\xi ]}{i\delta \overline{\eta }_\gamma (x_1)\delta \eta _\rho (y_1)}_{\overline{\eta }=\eta =\overline{\xi }=\xi =0}$$
(290)
and the other Green’s functions in the presence of source $`J`$ were defined before. When we set $`J=0`$, Eq. (A.31) straightforwardly yields Eq. (3.2). Finally, on differentiating Eq. (A.31) with respect to $`J^{a\mu }(x_i)`$ and subsequently setting $`J=0`$, Eq. (3.9) will be derived.
## X Appendix B: One-particle irreducible decompositions of the connected Green’s functions
Let us begin with the relation between the generating functional for full Green’s functions $`Z[J,\overline{\eta },\eta ,\overline{\xi },\xi ]`$ and the one for connected Green’s functions $`W[J,\overline{\eta },\eta ,\overline{\xi },\xi ]`$
$$Z[J,\overline{\eta },\eta ,\overline{\xi },\xi ]=\mathrm{exp}\{iW[J,\overline{\eta },\eta ,\overline{\xi },\xi ]\}$$
(291)
Taking the derivatives of Eq. (B.1) with respect to the sources $`\overline{\eta }(x_1),\overline{\eta }^c(x_2),\eta (y_1)`$ and $`\eta ^c(y_2)`$ and then setting all the sources except for the source $`J`$ to be zero, one may obtain the following decomposition
| $`G(x_{1,}x_2;y_1,y_2)^J=G_c(x_{1,}x_2;y_1,y_2)^J+S_F(x_1y_1)^JS_F^c(x_2y_2)^J`$ |
| --- |
| $`S_F^{}(x_1x_2)^J\overline{S}_F^{}(y_1y_2)^J`$ |
(292)
where $`G(x_{1,}x_2;y_1,y_2)^J,`$ $`S_F^c(x_2y_2)^J`$, $`S_F^{}(x_1x_2)^J`$ and $`S_F(x_1y_1)^J`$ were defined in Eqs. (A.23)-(A.25) and (A.32) respectively, while, $`G_c(x_{1,}x_2;y_1,y_2)^J`$ and $`\overline{S}_F^{}(y_1y_2)^J`$ are defined by
$$G_c(x_{1,}x_2;y_1,y_2)^J=i\frac{\delta ^4W[J,\overline{\eta },\eta ,\overline{\xi },\xi ]}{\delta \overline{\eta }(x_1)\delta \overline{\eta }^c(x_2)\delta \eta (y_1)\delta \eta ^c(y_2)}_{\overline{\eta }=\overline{\eta }^c=\eta =\eta ^c=0},$$
(293)
and
$$\overline{S}_F^{}(y_1y_2)^J=\frac{\delta ^2W[J,\overline{\eta },\eta ,\overline{\xi },\xi ]}{i^2\delta \eta (y_1)\delta \eta ^c(y_2)}_{\overline{\eta }=\eta =\overline{\xi }=\xi =0}$$
(294)
When we set $`J=0`$, Eq. (B.2) will go over to the decomposition shown in Eq. (4.1). Differentiating Eq. (B.2) with respect to the source $`J^{a\mu }(x_i)`$, we have
| $`G_\mu ^a(x_ix_{1,}x_2;y_1,y_2)^J=G_{c\mu }^a(x_ix_{1,}x_2;y_1,y_2)^J`$ |
| --- |
| $`+\mathrm{\Lambda }_\mu ^a(x_ix_1;y_1)^JS_F^c(x_1y_1)^J+S_F(x_1y_1)^J\mathrm{\Lambda }_\mu ^{𝐜a}(x_ix_2;y_2)^J`$ |
| $`\mathrm{\Lambda }_\mu ^a(x_ix_1,x_2)^J\overline{S}_F^{}(y_1y_2)^JS_F^{}(x_1x_2)^J\overline{\mathrm{\Lambda }}_\mu ^a(x_iy_1,y_2)^J`$ |
(295)
where $`G_\mu ^a(x_ix_{1,}x_2;y_1,y_2)^J`$ was defined in Eq. (A.26) with $`y_i`$ being replaced by $`x_i`$ and the other functions are defined by
$$G_{c\mu }^a(x_ix_{1,}x_2;y_1,y_2)^J=\frac{\delta }{i\delta J^{a\mu }(x_i)}G_c(x_{1,}x_2;y_1,y_2)^J$$
(296)
$$\mathrm{\Lambda }_\mu ^a(x_ix_1;y_1)^J=\frac{\delta }{i\delta J^{a\mu }(x_i)}S_F(x_1y_1)^J$$
(297)
$$\mathrm{\Lambda }_\mu ^{𝐜a}(x_ix_1;y_1)^J=\frac{\delta }{i\delta J^{a\mu }(x_i)}S_F^c(x_1y_1)^J$$
(298)
$$\mathrm{\Lambda }_\mu ^a(x_ix_1,x_2)^J=\frac{\delta }{i\delta J^{a\mu }(x_i)}S_F^{}(x_1x_2)^J$$
(299)
and
$$\overline{\mathrm{\Lambda }}_\mu ^a(x_iy_1,y_2)^J=\frac{\delta }{i\delta J^{a\mu }(x_i)}\overline{S}_F^{}(y_1y_2)^J$$
(300)
Upon setting $`J=0,`$ Eq. (B.5) immediately gives rise to the decomposition in Eq. (4.3).
Now, let us proceed to carry out one-particle-irreducible decompositions of the connected Green’s functions on the RHS of Eq. (4.4). The decompositions are easily performed with the help of the Legendre transformation which is described by the relation between the generating functional of proper vertices $`\mathrm{\Gamma }`$ and the one for connected Green’s functions $`W`$
$$\mathrm{\Gamma }[A_\mu ^a,\overline{\psi },\psi ,\overline{C}^a,C^a]=W[J,\overline{\eta },\eta ,\overline{\xi },\xi ]d^4x[J^{a\mu }A_\mu ^a+\overline{\eta }\psi +\overline{\psi }\eta +\overline{\xi }C+\overline{C}\xi ]$$
(301)
and the relations between the field functions and the external sources
$$\psi (x)=\frac{\delta W}{\delta \overline{\eta }(x)},\overline{\psi }(x)=\frac{\delta W}{\delta \eta (x)},A_\mu ^a(x)=\frac{\delta W}{\delta J^{a\mu }(x)},C^a(x)=\frac{\delta W}{\delta \overline{\xi }^a(x)},\overline{C}^a(x)=\frac{\delta W}{\delta \xi ^a(x)}$$
(302)
$$\eta (x)=\frac{\delta \mathrm{\Gamma }}{\delta \overline{\psi }(x)},\overline{\eta }(x)=\frac{\delta \mathrm{\Gamma }}{\delta \psi (x)},J_\mu ^a(x)=\frac{\delta \mathrm{\Gamma }}{\delta A^{a\mu }(x)},\xi ^a(x)=\frac{\delta \mathrm{\Gamma }}{\delta \overline{C}^a(x)},\overline{\xi }^a(x)=\frac{\delta \mathrm{\Gamma }}{\delta C^a(x)}$$
(303)
where the field functions in Eq. (B.12) are all functionals of the external sources in Eq. (B.13) and, simultaneously, the sources in Eq. (B.13) are all functionals of the field functions in Eq. (B.12).
Taking the derivative of the both sides of the first equality in Eq. (B.12) with respect to $`\psi (y)`$ and employing the first relation in Eq. (B.13), one may get
$$d^4z\frac{\delta ^2\mathrm{\Gamma }}{\delta \psi (y)\delta \overline{\psi }(z)}\frac{\delta ^2W}{\delta \eta (z)\delta \overline{\eta }(x)}=d^4z\frac{\delta ^2W}{\delta \overline{\eta }(x)\delta \eta (z)}\frac{\delta ^2\mathrm{\Gamma }}{\delta \overline{\psi }(z)\delta \psi (y)}=\delta ^4(xy)$$
(304)
where we only keep the term on the RHS of Eq. (B.14) which is nonvanishing when the sources are set to vanish. In order to find the one-particle-irreducible decomposition for the quark-gluon three-point Green’s functions, one may differentiate Eq. (B.14) with respect to the source $`J^{a\mu }(x_i)`$ and then using Eq. (B.14) once again. By this procedure, it can be derived that
| $`\frac{\delta ^3W}{\delta J^{a\mu }(x_i)\delta \overline{\eta }(x_j)\delta \eta (y_k)}=d^4zd^4u_1d^4u_2\frac{\delta ^2W}{\delta J^{a\mu }(x_i)\delta J^{b\nu }(u_1)}\frac{\delta ^2W}{\delta \overline{\eta }(x_j)\delta \eta (u_2)}`$ |
| --- |
| $`\frac{\delta ^3\mathrm{\Gamma }}{\delta A_\nu ^b(u_1)\delta \overline{\psi }(u_2)\psi (z)}\frac{\delta ^2W}{\delta \overline{\eta }(z)\delta \eta (y_k)}`$ |
(305)
where the coordinates in Eq. (B.14) have been appropriately changed. When all the sources are set to be zero, noticing the definitions given in Eq. (A.32) where the $`Z`$ is replaced by $`iW`$ and in Eq. (B.7) as well as
$$\mathrm{\Delta }_{\mu \nu }^{ab}(x_iy_j)=\frac{\delta ^2W}{i^2\delta J^{a\mu }(x_i)\delta J^{b\nu }(y_j)}_{J=0}$$
(306)
$$\mathrm{\Gamma }^{b\nu }(u_1u_2,z)=i\frac{\delta ^3\mathrm{\Gamma }}{\delta A_\nu ^b(u_1)\delta \overline{\psi }(u_2)\delta \psi (z)}_{A=\overline{\psi }=\psi =0}$$
(307)
the decomposition shown in Eqs. (4.5) and (4.6) straightforwardly follows from Eq. (B.15). Analogously, if we replace $`\overline{\eta }(x_j)`$ and $`\eta (y_k)`$ by $`\overline{\eta }^c(x_j)`$ and $`\eta ^c(y_k)`$ in Eq. (B.15) and noticing
$$\mathrm{\Gamma }_c^{b\nu }(u_1u_2,z)=i\frac{\delta ^3\mathrm{\Gamma }}{\delta A_\nu ^b(u_1)\delta \overline{\psi }^c(u_2)\delta \psi ^c(z)}_{A=\overline{\psi }=\psi =0}$$
(308)
the decomposition shown in Eq. (4.8) and (4.9) will be derived. This decomposition may also be derived from Eq. (B.15) by the charge conjugation transformation for the quark fields. By this transformation, one may readily derive from Eq. (B.15) the decomposition denoted in Eq. (4.15) in which the gluon-quark-antiquark vertex is defined by
$$\overline{\mathrm{\Gamma }}^{b\nu }(zz_1,z_2)=i\frac{\delta ^3\mathrm{\Gamma }}{\delta A_\nu ^b(z)\delta \psi (z_1)\delta \psi ^c(z_2)}_{A=\overline{\psi }=\psi =0}$$
(309)
The one-particle-irreducible decomposition of the connected Green’s function $`G_c(x_{1,}x_2;y_1,y_2)`$ can be derived by the same procedure as obtaining Eq. (B.15). On differentiating Eq. (B.14) with respect to $`\overline{\eta }^c(x_2)`$ and $`\eta ^c(y_2)`$ and setting all the sources but the source $`J`$ to vanish, one may obtain
| $`G_c(x_{1,}x_2;y_1,y_2)^J=\underset{i=1}{\overset{2}{}}d^4u_id^4v_iS_F(x_1u_1)^JS_F^c(x_2u_2)^J`$ |
| --- |
| $`\times \mathrm{\Gamma }(u_1,u_2;v_1,v_2)^JS_F(v_1y_1)^JS_F^c(v_2y_2)^J`$ |
(310)
where the four-point connected Green’s function and the propagators given in the presence of the sources were defined before and the function $`\mathrm{\Gamma }(u_1,u_2;v_1,v_2)^J`$ is formally the same as that defined in Eqs. (4.19)-(4.21). When the source $`J`$ is turned off, Eq. (B.20) directly goes over to the decomposition in Eq. (4.18) with the vertices in Eqs. (4.19)-(4.21) being defined in Eqs. (B.17)-(B.19) and in the following
$$\mathrm{\Gamma }^{b\nu }(x_1u_1,v_1)=i\frac{\delta ^3\mathrm{\Gamma }}{\delta A_\nu ^b(x_1)\delta \overline{\psi }(u_1)\delta \overline{\psi }^c(v_2)}_{A=\overline{\psi }=\psi =0}$$
(311)
which is the charge conjugate to the vertex $`\overline{\mathrm{\Gamma }}^{b\nu }(zz_1,z_2)`$ as well as
$$\mathrm{\Gamma }_3(u_1,u_2;v_1,v_2)=i\frac{\delta ^4\mathrm{\Gamma }}{\delta \overline{\psi }(u_1)\delta \overline{\psi }^c(u_2)\delta \psi (v_1)\delta \psi ^c(v_2)}_{\overline{\psi }=\psi =\overline{\psi }^c=\psi ^c=0}$$
(312)
which is the quark-antiquark four-line proper vertex. It is emphasized here that the decomposition of the function $`G_c(x_{1,}x_2;y_1,y_2)`$ in the absence of the source $`J`$ has the same form as that given in the presence of $`J`$. This is because the Green’s function is defined only by the differentials with respect to the fermion fields as indicated in Eq. (B.3).
The one-particle irreducible decomposition of the Green’s function $`G_{c\mu }^a(x_ix_{1,}x_2;y_1,y_2)`$ may be derived by starting from the expression given in Eq. (B.15) with $`j,`$ $`k=1`$. By differentiating the both sides of Eq. (B.15) with respect to the sources $`\overline{\eta }^c(x_2)`$ and $`\eta ^c(y_2)`$ and then turning off all the external sources, one may obtain the decomposition of the function $`G_{c\mu }^a(x_ix_{1,}x_2;y_1,y_2)`$ as shown in Eqs. (4.22)-(4.25). Alternatively, the decomposition may also be obtained by starting with the expression written in Eq. (B.20). Substituting Eq. (B.20) into Eq. (B.6), then completing the differentiation with respect to the source $`J^{a\mu }(x_i)`$ and finally setting the source to vanish, one may also derive the irreducible decomposition of the function $`G_{c\mu }^a(x_ix_{1,}x_2;y_1,y_2)`$. In doing this, it is necessary to perform the differentiations of the fermion propagators with respect to the source $`J^{a\mu }(x_i)`$ as shown in Eqs. ( B.7)-(B.10) and use their decompositions presented in Eqs. (4.5)-(4.9). In addition, we need to carry out the differentiations of the gluon propagator and some vertices with respect to the source $`J^{a\mu }(x_i)`$ as shown below. For the gluon propagator defined in Eq. (B.16), from its representation in presence of the external source $`J`$ ,in the same way as deriving the decomposition represented in Eqs. (B.15), ( 4.5) and (4.6), one may obtain the one-particle irreducible decomposition of the gluon three-point Green’s function as follows:
$$\mathrm{\Lambda }_{\mu \rho \sigma }^{acd}(x_i,z_1,z_2)=\frac{\delta }{\delta J^{a\mu }(x_i))}\mathrm{\Delta }_{\rho \sigma }^{cd}(z_1z_2)^J_{J=0}=d^4zD_{\mu \nu }^{ab}(x_iz)\mathrm{\Pi }_{\rho \sigma }^{bcd,\nu }(z,z_1,z_2)$$
(313)
where $`D_{\mu \nu }^{ab}(x_iz)=i\mathrm{\Delta }_{\mu \nu }^{ab}(x_iz)`$ and $`\mathrm{\Pi }_{\rho \sigma }^{bcd,\nu }(z,z_1,z_2)`$ was represented in Eq. (4.32) with the gluon three-line proper vertex defined by
$$\mathrm{\Gamma }_{bcd}^{\nu \rho \sigma }(z,u_1,u_2)=i\frac{\delta ^3\mathrm{\Gamma }}{\delta A_\nu ^b(z)\delta A_\rho ^c(u_1)\delta A_\sigma ^d(u_2)}_{A=0}$$
(314)
For a proper vertex $`\mathrm{\Gamma }_\alpha (z_1,z_2,\mathrm{})`$ with $`\alpha `$ marking the other indices, its derivative with respect to the source $`J^{a\mu }(x_i)`$ can be represented as
$$\frac{\delta }{i\delta J^{a\mu }(x_i))}\mathrm{\Gamma }_\alpha (z_1,z_2,\mathrm{})^J_{J=0}=d^4zD_{\mu \nu }^{ab}(x_iz)\mathrm{\Gamma }_\alpha ^{b\nu }(z,z_1,z_2,\mathrm{})$$
(315)
where
$$\mathrm{\Gamma }_\alpha ^{b\nu }(z,z_1,z_2,\mathrm{})=\frac{\delta }{\delta A_\nu ^b(z)}\mathrm{\Gamma }_\alpha (z_1,z_2,\mathrm{})^J_{J=0}$$
(316)
According to the procedure stated above, it is not difficult to derive the expressions described in Eqs. (4.30)-(4.33). In the expressions, the vertices are defined as follows:
$$\mathrm{\Gamma }^{a\mu }(x_iu_1,u_2;v_1,v_2)=\frac{\delta \mathrm{\Gamma }(u_1,u_2;v_1,v_2)^J}{i\delta J^{a\mu }(x_i)}_{J=0}$$
(317)
$$\mathrm{\Gamma }_{\nu \lambda }^{bc}(z,z_1u_1,v_1)=i\frac{\delta ^4\mathrm{\Gamma }}{\delta A^{b\nu }(z)\delta A^{c\lambda }(z_1)\delta \overline{\psi }(u_1)\delta \psi (v_1)}_{A=\overline{\psi }=\psi =0}$$
(318)
$$\mathrm{\Gamma }_{𝐜\nu \lambda }^{bc}(z,z_1u_1,v_1)=i\frac{\delta ^4\mathrm{\Gamma }}{\delta A^{b\nu }(z)\delta A^{c\lambda }(z_1)\delta \overline{\psi }^c(u_1)\delta \psi ^c(v_1)}_{A=\overline{\psi }^c=\psi ^c=0}$$
(319)
$$\mathrm{\Gamma }_{\nu \lambda }^{bc}(z,z_1u_1,v_1)=i\frac{\delta ^4\mathrm{\Gamma }}{\delta A^{b\nu }(z)\delta A^{c\lambda }(z_1)\delta \overline{\psi }(u_1)\delta \overline{\psi }^c(v_1)}_{A=\overline{\psi }=\overline{\psi }^c=0}$$
(320)
$$\overline{\mathrm{\Gamma }}_{𝐜\nu \lambda }^{bc^{}}(z,z_1u_1,v_1)=i\frac{\delta ^4\mathrm{\Gamma }}{\delta A^{b\nu }(z)\delta A^{c\lambda }(z_1)\delta \psi (u_1)\delta \psi ^c(v_1)}_{A=\psi =\psi ^c=0}$$
(321)
Particularly, by the following differentiation
$$\mathrm{\Gamma }_3^{a\mu }(x_iu_1,u_2;v_1,v_2)=\frac{\delta \mathrm{\Gamma }_3(u_1,u_2;v_1,v_2)^J}{i\delta J^{a\mu }(x_i)}_{J=0}$$
(322)
it is easy to give the expression in Eq. (4.34) in which
$$\widehat{\mathrm{\Gamma }}^{b\nu }(zu_1,u_2;v_1,v_2)=i\frac{\delta ^4\mathrm{\Gamma }}{\delta A^{b\nu }(z)\delta \overline{\psi }(u_1)\delta \overline{\psi }^c(u_2)\delta \psi (v_1)\delta \psi ^c(v_2)}_{A=\overline{\psi }=\psi =\overline{\psi }^c=\psi ^c=0}$$
(323)
is the gluon-quark-antiquark five-line proper vertex.
## XI References |
warning/0506/hep-th0506130.html | ar5iv | text | # Untitled Document
hep-th/0506130 SU-ITP-05/22 SLAC-PUB-11283
The Tachyon at the End of the Universe
John McGreevy and Eva Silverstein
SLAC and Department of Physics, Stanford University, Stanford, CA 94305-4060
We show that a tachyon condensate phase replaces the spacelike singularity in certain cosmological and black hole spacetimes in string theory. We analyze explicitly a set of examples with flat spatial slices in various dimensions which have a winding tachyon condensate, using worldsheet path integral methods from Liouville theory. In a vacuum with no excitations above the tachyon background in the would-be singular region, we analyze the production of closed strings in the resulting state in the bulk of spacetime. We find a thermal result reminiscent of the Hartle-Hawking state, with tunably small energy density. The amplitudes exhibit a self-consistent truncation of support to the weakly-coupled small-tachyon region of spacetime. We argue that the background is accordingly robust against back reaction, and that the resulting string theory amplitudes are perturbatively finite, indicating a resolution of the singularity and a mechanism to start or end time in string theory. Finally, we discuss the generalization of these methods to examples with positively curved spatial slices.
July 2005
1. Introduction
Closed string tachyon condensation affects the dynamics of spacetime in interesting and tractable ways in many systems \[1--5\]. In this paper, we study circumstances in which closed string tachyon condensation plays a crucial role in the dynamics of a system which has a spacelike singularity at the level of general relativity. The singular region of the spacetime is replaced by a phase of tachyon condensate which lifts the closed string degrees of freedom, effectively ending ordinary spacetime.
We will focus primarily on a simple set of examples with shrinking circles, in which we can make explicit calculations exhibiting this effect. Before specializing to this, let us start by explaining the relevant structure of the stringy corrections to spacelike singularities appearing in a more general context. Much of this general discussion appeared earlier in .<sup>1</sup> The possibility of applying the worldsheet mass gap in higher dimensional generalizations of was also independently suggested by A. Adams and M. Headrick. Consider a general relativistic solution approaching a curvature singularity in the past or future. The metric is of the form
$$ds^2=G_{\mu \nu }dx^\mu dx^\nu =(dx^0)^2+R_i(x^0)^2d\mathrm{\Omega }_i^2+ds_{}^2$$
$`(1.1)`$
with $`R_i(x^0)0`$ for some $`i`$ at some finite time. Here $`\mathrm{\Omega }_i`$ describe spatial coordinates whose scale factor is varying in time and $`ds_{}^2`$ describes some transverse directions not directly participating in the time dependent physics.
In the large radius regime where general relativity applies, the background (1.1) is described by a worldsheet sigma model with action in conformal gauge
$$S_0\frac{1}{4\pi \alpha ^{}}d^2\sigma G_{\mu \nu }(X)_aX^\mu ^aX^\nu +\mathrm{fermions}+\mathrm{ghosts}$$
$`(1.2)`$
Here we are considering a type II or heterotic string with worldsheet supersymmetry in order to avoid bulk tachyons.
As the space shrinks in the past or future, at leading order in $`\alpha ^{}`$ (i.e. in GR) the corresponding sigma model kinetic terms for $`\mathrm{\Omega }`$ develop small coefficients, leading to strong coupling on the worldsheet. This raises the possibility of divergent amplitudes in the first quantized worldsheet path integral description from lack of suppression from the action. This would correspond also to the development of an effectively strong coupling in the spacetime theory as the size of the $`\mathrm{\Omega }`$ directions shrink.
However, there is more to the story in string theory. Let us first consider the sigma model on the angular geometry at fixed time $`X^0`$. When any of the radii $`R_i`$ in (1.1) is of order string scale, this strongly coupled sigma model is very different from the free flat-space theory. In particular, it can dynamically generate a mass gap in the IR \[7,,8\]. In such cases the quantum effective action in this matter sigma model has terms of the form $`d^2\sigma 𝒪_\mathrm{\Delta }\mathrm{\Lambda }^{2\mathrm{\Delta }}`$ where $`\mathrm{\Delta }<2`$ is the dimension of some relevant operator $`𝒪`$ and $`\mathrm{\Lambda }`$ is a mass scale. Hence the full string path integral (1.2) generates additional contributions to the worldsheet effective action of the form
$$S_T=d^2\sigma \mu f(X^0)𝒪_\mathrm{\Delta }(X_{},\mathrm{\Omega })+d^2\sigma \mathrm{\Phi }(X^0)R^{(2)}+\mathrm{fermions}$$
$`(1.3)`$
where $`f`$ has dimension $`2\mathrm{\Delta }`$ in the unperturbed sigma model (1.2). We will henceforth refer to such deformations as “tachyons”; in the simplest case of an $`S^1`$ spatial component the corresponding mode is a standard winding tachyon. Let us discuss the big bang case for definiteness: the system becomes weakly curved in the future (large positive $`X^0`$) and goes singular at some finite value of $`X^0`$ in the past. The contribution (1.3) goes to zero as $`X^0+\mathrm{}`$ since the sigma model is weakly coupled there. So at its onset the coefficient $`f`$ increases as $`X^0`$ decreases; i.e. the effects of the term (1.3) increase as we go back in time in the direction of the would-be big bang singularity of the GR solution (1.1). In the simple case we will study in detail in §2,3 below, a term of the form (1.3) will arise from winding tachyon condensation, and the operator $`f`$ will be of the form $`e^{\kappa X^0}`$ for real positive $`\kappa `$ in the big bang case.
This growth of (1.3) as we approach the singularity contrasts to the suppression of the original sigma model kinetic terms from the metric (1.2). In the Minkowski path integral, the growth of the term (1.3) serves to suppress fluctuations of the fields. This provides a possibility of curing – via perturbative string effects – the singular amplitudes predicted by a naive extrapolation of GR. We will see this occur explicitly in the examples we will study in detail in §3.
Fig. 1: In string theory (with string length scale $`l_s`$), a tachyon condensate phase replaces a spacelike singularity that would have been present at the level of general relativity.
Because of the mass gap in the matter sector and the effect of the deformation $`S_T`$ on the spacetime mass spectrum, the condensation of tachyons has long been heuristically argued to lift the string states and lead to a phase of “Nothing”\[9--14\]. In the examples where conical singularities resolve into flat space, this is borne out in detail, as the tip of the cone disappears in the region of tachyon condensation; a similar phenomenon was found for localized winding tachyons in . In the present work, we will use methods from Liouville theory (for a review see \[15,,16,,17\]). We employ and extend the methods of \[18--25\] to perform systematic string theoretic calculations of amplitudes exhibiting this effect in our temporally but not spatially localized case. We study a particular vacuum state, analogous to the Euclidean vacuum. In this vacuum the support of string theoretic amplitudes is restricted to the bulk region of spacetime in a way that we can derive from the zero mode integral of $`X^0`$ in the worldsheet path integral.
As discussed above, the metric coefficient $`G_{\mathrm{\Omega }\mathrm{\Omega }}=R(X^0)^2`$ in the worldsheet action $`S_0+S_T`$ (1.2)(1.3) goes to zero at finite $`X^0`$. In the models we consider below we will set up the system such that the approach to this singularity is parameterically slower than the timescale for the relevant term $`S_T`$ to become important and lift the closed string degrees of freedom. This will avoid strong dilaton effects as well as effects of the shrinking space, and provides a mechanism which is inherently perturbative in $`g_s`$. The amplitudes will exhibit limited support in the spacetime, contributing only in the bulk region away from the epoch when the couplings become important. This is similar to the situation in spacelike Liouville theory, where similar strong coupling effects are avoided by the presence of the Liouville wall, and similar computational methods apply, though the physical mechanism for suppressing amplitudes is different in the two cases. Hence a self-consistent perturbative analysis is available. Relatedly, black hole formation of the sort found in \[26,,27\] is evaded here: the tachyon lifts the degrees of freedom of the system before the Planckian regime is reached.
The observables of this theory are correlation functions of integrated vertex operators computed by the worldsheet path integral with semiclassical action $`S_0+S_T`$ (1.2)(1.3); let us now discuss their spacetime interpretation. As in Liouville theory, the form of these operators is known in the weakly curved bulk region where there is no tachyon condensate ($`X^0\mathrm{}`$ in the big bang case); there they asymptote in locally flat coordinates to operators of the form
$$V_{\stackrel{}{k},n}e^{i\stackrel{}{k}\stackrel{}{X}}e^{i\omega (\stackrel{}{k},n)X^0}\widehat{V}_n\mathrm{as}X^0\mathrm{}$$
$`(1.4)`$
where $`n`$ labels the string state with mass $`m_n`$ coming from oscillator excitations created by $`\widehat{V}_n`$, $`\stackrel{}{k}`$ its spatial momentum, and $`\omega ^2=\stackrel{}{k}^2+m_n^2`$.<sup>2</sup> Although we do not know the form of these operators in the regime where the corrections (1.3) become important, we do know their conformal dimensions by virtue of their form in the bulk region of the spacetime. This is as in Liouville theory, where one knows the operators and the stress tensor away from the Liouville wall, and hence the spectrum of dimensions. And as in Liouville theory, an important question which we will address is where the amplitudes built from these operators have their support. Integrated correlation functions of these operators have the interpretation as components of the state of the strings in the bulk region of spacetime in a basis of multiple free string modes. In our example below, we will focus on a vacuum with no excitations above the tachyon condensate in the would-be singular region, and compute the resulting state of perturbative strings in the bulk region. This is a string-theoretic analogue of the Hartle-Hawking State (equivalently, the Euclidean Vacuum) on our time-dependent background.
We will treat the condensing tachyon in string perturbation theory. We obtain a self-consistent analysis at weak string coupling, in systems with bulk supersymmetry, and with supersymmetry breaking near the would-be singular region similar to that expected in the early universe and inside black holes. Other interesting recent work on perturbative closed string calculations in time dependent backgrounds includes \[27,,28\]. In our case the tachyon condensation, related to the supersymmetry breaking of the time dependent background, plays a crucial role, in a way anticipated in .
It would be interesting to relate our analysis to other approaches based on non-perturbative formulations of the theory \[29--31\]. These approaches may provide a complete nonperturbative dual formulation of observables in spacetimes with singularities at the level of GR. On the other hand, the dictionary between the two sides is sometimes rather indirect as applied to approximately local processes on the gravity side. A useful feature of the current approach is that the tachyon condensation provides a direct gravity-side mechanism for quelling the singularity. It would be interesting to see how this information is encoded in the various dual descriptions.
Finally, analogously to the case of open string tachyons (for a review, see ), closed string tachyons may be a subject well-studied via closed string field theory; a candidate “nothing” state obtained from bosonic closed string bulk tachyon condensation was recently presented in . It is clear (as we will review as we go) that the physics of the tachyon condensate is stringy – low energy effective field theory is not sufficient. In the setup we consider here, perturbative methods using techniques from Liouville theory will suffice, but in more general situations the off shell methods of string field theory may be required.
In the next two sections we set up and analyze a class of realizations of the mechanism. In §4 we describe the generalization to positive spatial curvature, which is velocity-dominated. Philosophy-dominated comments are restricted to the concluding section.
2. Examples with winding tachyons
In this section, we will introduce the simplest backgrounds we will study; those with flat spatial slices which expand at a tunable rate. We will start with an example pertaining to 2+1 dimensional black holes (reducing to the 1+1 dimensional Milne spacetime inside), and then generalize to higher dimensional flat FRW cosmology with topologically nontrivial spatial slices and radiation.
2.1. The Milne Spacetime
Consider the Milne spacetime described by the metric
$$ds^2=(dx^0)^2+v^2(x^0)^2d\mathrm{\Omega }^2+d\stackrel{}{x}^2$$
$`(2.1)`$
For $`x^0>0`$, this solution describes a growing $`S^1`$ along the $`\mathrm{\Omega }`$ direction. At $`x^0=0`$ there is a spacelike big bang singularity, and general relativity breaks down. The evolution from $`x^0=\mathrm{}`$ to $`x^0=0`$ similarly describes an evolution toward a big crunch singularity. This geometry appears inside $`2+1`$ dimensional black holes, BTZ black holes in $`AdS_3`$.<sup>3</sup> “Whisker” regions with closed timelike curves also appear in the maximally extended spacetime; our methods here will also have the effect of excising these regions, as obtained in other examples in . We will show that for a wide class of string theories, the spacelike big bang or big crunch singularity (2.1) is evaded – the regime $`|vx^0|<l_s`$ is replaced by a phase of tachyon condensate.
In particular, to avoid bulk tachyons, consider type II, type I or heterotic string theory on the spacetime (2.1). Take antiperiodic boundary conditions around the $`\mathrm{\Omega }`$ circle for spacetime fermions. Further consider the regime of parameters where $`v1`$. In addition to providing the control we will require, the last two conditions correspond to those appropriate for small BTZ black holes which can form naturally from excitations in pure $`AdS_3`$ (which has antiperiodic boundary conditions for fermions around the contractible spatial circle surrounding the origin).
With these specifications, we can determine with control the spectrum of string theory on the spacetime (2.1) for $`x^00`$. In the regime
$$v^2(x^0)^2l_s^2$$
$`(2.2)`$
a closed string winding mode becomes tachyonic and hence important to the dynamics. The regime $`v|x^0|l_s`$ of the singularity in (2.1) is replaced by a phase of tachyon condensate. This offers a concrete avenue toward resolving a spacelike singularity in string theory, and a corresponding notion of how time can begin or end.
This in itself is worth emphasizing. The problem of bulk tachyon condensation is often motivated by the question of the vacuum structure of string theory. The present considerations provide an independent motivation for pursuing the physics of closed string tachyon condensation: it appears crucially in a string-corrected spacelike singularity. In our system here there is no tachyonic mode in the bulk of spacetime: for a semiinfinite range of time the system is perturbatively stable. That is, the tachyon phase is localized in time. As we will see, this provides significant control over the problem even though the condensation is not also localized in space.
2.2. Flat FRW with topology
Next let us set up a somewhat more realistic case which shares the essential features of the above example. Consider flat-sliced FRW cosmology with bulk metric
$$ds^2=(dx^0)^2+a^2(x^0)d\stackrel{}{x}^2+ds_{}^2$$
$`(2.3)`$
with $`\stackrel{}{x}`$ a 3-dimensional spatial vector and $`ds_{}^2`$ describing the extra dimensions. Let us consider some periodicity in the spatial directions $`\stackrel{}{x}`$: $`\stackrel{}{x}\stackrel{}{x}+\stackrel{}{L}_I`$; e.g. letting $`I`$ run from 1 to 3 produces a spatial torus (for simplicity let us take a rectangular torus). In real cosmology, such topology could well exist at sufficiently large scales (most generically well outside our horizon today due to inflation), but if present would play a role in the far past in the epoch of the would-be big bang singularity. (See e.g. for one recent discussion of spatial topology.).
Let us study the above system in the presence of a stress-energy source. For definiteness, consider a homogeneous bath of radiation. The Friedmann equation implies
$$a(x_0)=a_0\sqrt{x^0t_0}$$
$`(2.4)`$
where the coefficient $`a_0`$ can be tuned by dialing the amount of radiation.
In particular, as in the above example (2.2), we can choose the radiation density and hence $`a_0`$ so as to obtain a slow expansion of the toroidal radii $`RLa(x^0)`$ as the smallest radius passes through the string scale. Again considering antiperiodic boundary conditions for fermions along one or more of the 1-cycles of the torus, we then obtain in a controlled way a winding tachyon in the system as the radius $`RLa(x^0)`$ of a circle passes below the string scale. The would-be big bang singularity is again replaced by a tachyon condensate phase, whose consequences we will analyze in detail in the next section.
3. Examples with winding tachyons: some basic computations of observables
In this section, we develop a systematic computational scheme to compute physical observables in this system, assess back reaction, and test and make more precise the proposition that tachyon condensation lifts closed string excitations (leading to a phase we will refer to as a Nothing state).
3.1. Wick rotation
Let us start by defining the path integral via appropriate Wick rotation. In its original Lorentzian signature form, the tachyon term appears to increase the oscillations of the integrand, hence suppressing contributions in the region of the tachyon condensate. As is standard in quantum field theory, we will perform a Wick rotation to render the path integral manifestly convergent (up to, as we will see, divergences at exceptional momenta expected from the bulk S-matrix point of view). The path integral in conformal gauge includes an integral over the target space time variable $`X^0`$, which has a negative kinetic term in the worldsheet theory. Because this field also appears necessarily in the tachyon interaction term (which is proportional to $`e^{\kappa X^0}`$, specializing to the big bang case), we will find it convenient to Wick rotate the worldsheet theory to directly obtain real positive kinetic terms for $`X^0`$ without rotating the contour for $`X^0`$ integration; this will entail rotating the contours for the spatial target space coordinates as well as continuing $`\mu `$ in a way we will specify. (Alternatively one could rotate $`X^0`$ as is standardly done in the free theory, and continue in $`\kappa `$ at the same time, as in .)
Prelude: worldline quantum field theory
Before turning to the full string path integral, let us briefly describe a much simpler analogue of our system which arises in the worldline description of quantum field theory, as emphasized in . Consider a relativistic particle action
$$S=𝑑\tau \left((_\tau X^0)^2+(_\tau \stackrel{}{X})^2(m_0^2+\mu ^2e^{2\kappa X^0})\right)$$
$`(3.1)`$
where we have included a time-dependent mass squared term $`m^2(X^0)=m_0^2+\mu ^2e^{2\kappa X^0}`$.
For $`\mu ^2>0`$, this theory describes a particle with a time-dependent positive mass-squared that increases exponentially in the past $`X^0\mathrm{}`$. The potential term in the relativistic worldline action leads to a lifting of particles in the region where it becomes important. If one starts with none of these massive modes excited in the past, then the future state gets populated due to the time dependent mass. The Bogoliubov coefficient $`\beta _\stackrel{}{k}`$ describing mixing of positive and negative frequency modes has magnitude $`e^{\pi \omega /\kappa }`$ with $`\omega =\sqrt{\stackrel{}{k}^2+m_0^2}`$ the frequency of the particle modes in the region $`X^0+\mathrm{}`$. We will find similar features in our string theoretic examples, where the phase in which states are lifted replaces a spacelike singularity.
For $`\mu ^2<0`$, this theory describes a system with time-dependent negative mass squared. A particle with positive mass squared in the far future becomes tachyonic in the far past. One could formally start again with no excitations in the past, but this would be unnatural as the tachyonic modes there would condense.
In Lorentzian signature, the worldline path integral is
$$[dX]e^{iS}$$
$`(3.2)`$
If we continue $`\tau e^{i\gamma }\tau _\gamma `$ and $`\stackrel{}{X}e^{i\gamma }\stackrel{}{X}_\gamma `$, taking $`\gamma `$ continuously from $`0`$ to $`\pi /2`$, we obtain a path integral
$$[dX_E]\mathrm{exp}\left[𝑑\tau _E\left((_{\tau _E}X^0)^2+(_{\tau _E}\stackrel{}{X})^2+\mu ^2e^{2\kappa X^0}\right)\right].$$
$`(3.3)`$
Ambiguities in defining the $`X^0`$ integral correspond to choices of vacuum state. In order to obtain a convergent path integral, we can continue $`\mu ^2\mu ^2`$ (i.e. $`\mu e^{i\gamma }\mu `$) as we do the above Wick and contour rotation, compute the amplitudes, and then continue back. That is, our computation is related to one in a purely spacelike target space via a reflection of the potential term in the worldline theory.
This reflection appears also in the direct spacetime analysis of particle production in field theory with time dependent mass. The Heisenberg equation of motion satisfied by the Heisenberg picture fields in spacetime takes the form of a Schrödinger problem for each $`\stackrel{}{k}`$ mode. The effective potential in the Schrödinger problem is $`U_{eff}=(m^2(X^0)+\stackrel{}{k}^2)`$. This leads to highly oscillating mode solutions as $`X^0\mathrm{}`$, reflecting the exponentially increasing mass in the far past.
With this warmup, let us now turn to the case of string theory with a tachyon condensate, which has a worldsheet potential term analogous to the mass squared term in the field theory case. We will study both the heterotic and type II theories on our background.
Heterotic Theory
In the RNS description of the heterotic string, the worldsheet theory has local (0,1) supersymmetry. This case is in some ways the simplest for studying closed string tachyons using the string worldsheet description – unlike the bosonic theory, there is no tachyon in the bulk region where the $`S^1`$ is large; unlike the type II theory the worldsheet bosonic potential is automatically nonnegative classically (as in the open superstring theory).
There is a choice of discrete torsion in the heterotic theory on a space with shrinking Scherk-Schwarz circle. The background can be regarded as a $`\mathrm{ZZ}_2`$ orbifold of a circle by a shift halfway around combined with an action of $`(1)^F`$ where $`F`$ is the spacetime fermion number. Combining this $`\mathrm{ZZ}_2`$ with that of the left-moving fermions (in say the SO(32) Heterotic theory) yields two independent choices of action of the left moving GSO on the states of the Scherk-Schwarz twisted sector. A standard choice, giving rise to a Hagedorn tachyon, is to act trivially on the Scherk-Schwarz twisted sector; this yields a twisted sector tachyon made from momentum and winding modes .
This would also be the most natural choice for us in some sense, since the usual spacelike singularities in cosmology and inside black holes are a priori independent of Yang-Mills degrees of freedom. However because the winding+momentum twisted tachyon in the above case is a nonlocal operator on the worldsheet in both T-duality frames, we will make here a technically simpler choice. Namely, we can choose the discrete torsion such that the left-moving GSO projection acts nontrivially on the states of the twisted sector, yielding a twisted tachyon created by a left moving fermion and a winding operator.
In the heterotic theory we have target space coordinates given by (0,1) scalar superfields $`𝒳^\mu =X^\mu +\theta ^+\psi _+^\mu `$ and left moving fermion superfields $`\mathrm{\Psi }_{}^a=\psi _{}^a+\theta ^+F^a`$ containing auxiliary fields $`F^a`$. In terms of these fields we have a Lorentzian signature path integral
$$G(\{V_n\})[d𝒳][d\mathrm{\Psi }_{}][d(\mathrm{ghosts})]d(\mathrm{moduli})e^{iS}\underset{n}{}\left(i𝑑\sigma 𝑑\tau V_n[𝒳]\right)$$
$`(\mathcal{3.4})`$
where the semiclassical action is
$$\begin{array}{cc}\hfill iS=& id\sigma d\tau d\theta ^+(D_{\theta ^+}𝒳^\mu _{}𝒳^\nu G_{\mu \nu }(𝒳)\mu \mathrm{\Psi }_{}:e^{\kappa 𝒳^0}\mathrm{cos}(w\stackrel{~}{\mathrm{\Omega }}):\hfill \\ & +\mathrm{\Psi }_{}^aD_{\theta ^+}\mathrm{\Psi }_{}^a+(\mathrm{dilaton}))+iS_{\mathrm{ghost}}\hfill \end{array}$$
$`(3.5)`$
and $`V_n[𝒳]`$ are vertex operator insertions. Here $`\stackrel{~}{\mathrm{\Omega }}`$ is the T-dual of the coordinate $`\mathrm{\Omega }`$ on the smallest circle in the space (let us consider for genericity a rectangular torus, whose smallest cycle will play a leading role in the dynamics); $`\mathrm{cos}w\stackrel{~}{\mathrm{\Omega }}`$ is the winding operator for strings wrapped around the $`\mathrm{\Omega }`$ direction.
The case of no insertions corresponds to the vacuum amplitude $`Z`$. The fluctuations of the worldsheet fields in (3.4) generate corrections to the action (3.5); for example, the term proportional to $`\mu `$ coming from the tachyon condensate is marginal but not exactly marginal. This is similar to the form of the Liouville wall in Liouville field theory, which is a priori only semiclassically given by a pure exponential.
Similarly, the form of the vertex operators is known semiclassically. Because the bulk region of the geometry (2.1) is approximately flat space, we may identify the $`V_n`$ with operators of the form
$$V_{\stackrel{}{k},n}e^{i\stackrel{}{k}\stackrel{}{𝒳}}e^{i\omega (\stackrel{}{k},n)𝒳^0}\widehat{V}_n\mathrm{as}X^0\mathrm{}$$
$`(3.6)`$
where as in (1.4) we have pulled out the oscillator and ghost contributions into $`\widehat{V}`$.
Finally, at the semiclassical level the dilaton is also known: it goes to a constant
$$\mathrm{\Phi }\mathrm{\Phi }_0\mathrm{as}X^0+\mathrm{}.$$
$`(3.7)`$
In particular, the tachyon vertex operator in (3.5) is semiclassically marginal without an additional dilaton contribution (though not exactly marginal) and the metric terms solve Einstein’s equations. The path integral over fluctuations of the fields will generate corrections to these semiclassical statements (3.5)(3.6)(3.7).
The basic physical effect of the tachyon condensate is the oscillations it induces in the path integral, suggesting a suppression of contributions rather than a divergence in the would-be crunch region. For some computations, it is useful to Wick rotate in evaluating the path integral (though as we will see below some computations such as the 1-loop partition function could be evaluated more directly).
Let us Wick rotate the worldsheet time coordinate $`\tau `$, the spatial target space coordinates $`\stackrel{}{𝒳}(\sigma ,\tau )`$ (including $`\stackrel{~}{\mathrm{\Omega }}`$), and the parameters $`\mu `$ and $`\stackrel{}{k}`$ by
$$\tau e^{i\gamma }\tau _\gamma \stackrel{}{𝒳}e^{i\gamma }\stackrel{}{𝒳}_\gamma \mu =e^{i\gamma }\mu _\gamma \stackrel{}{k}=e^{i\gamma }\stackrel{}{k}_\gamma $$
$`(3.8)`$
where $`\gamma `$ is a phase which we will rotate from $`0`$ to $`\pi /2`$. This produces a Euclidean path integral for the worldsheet theory (where we label the quantities rotated to $`\gamma =\pi /2`$ by a subscript $`E`$)
$$G(\{V_n\})[d\stackrel{}{𝒳}_E][d𝒳^0][d\mathrm{\Psi }_{}][d(\mathrm{ghosts})]d(\mathrm{moduli})e^{S_E}\underset{n}{}(1)𝑑\sigma 𝑑\tau _EV_{n,i\stackrel{}{k}_E}[𝒳^0,i\stackrel{}{𝒳}_E]$$
$`(3.9)`$
with Euclidean action
$$\begin{array}{cc}\hfill S_E=& d\sigma d\tau _Ed\theta ^+(D_{\theta ^+}𝒳^0_{}𝒳^0+v^2(𝒳^0)^2D_{\theta ^+}\stackrel{~}{\mathrm{\Omega }}_E_{}\stackrel{~}{\mathrm{\Omega }}_E+G_{ij}D_{\theta ^+}𝒳_{,E}^i_{}𝒳_{,E}^j\hfill \\ & i\mu _Ee^{\kappa 𝒳^0}\mathrm{cosh}(w\stackrel{~}{\mathrm{\Omega }}_E)\hfill \\ & +\mathrm{\Psi }_{}^aD_{\theta ^+}\mathrm{\Psi }_{}^a+(\mathrm{dilaton}))+S_E(\mathrm{ghost})\hfill \end{array}$$
$`(3.10)`$
Here $`\stackrel{}{𝒳}_E(\mathrm{\Omega }_E,\stackrel{}{𝒳}_E)`$ refers to the worldsheet superfields corresponding to the spatial target space coordinates, and we have plugged in the spacetime metric (2.1).
The procedure (3.8) alone yields a formal path integral whose tachyon term is not marginal in the Euclidean theory. In order to retain semiclassical conformal invariance, we could include in (3.8) a linear dilaton coefficient: $`Q_\gamma \gamma (2/\pi )Q`$ where $`Q`$ corresponds to the coefficient of the linear dilaton which renders the tachyon vertex operator marginal at leading order. This is similar to part of the procedure followed in e.g. . Alternatively we can analyze the theory for $`Q=0`$, in which case the intermediate steps in the analysis involve path integral computations in a nonconformal Euclidean worldsheet theory, obtained by analytic continuation in $`\kappa `$ from a conformal Euclidean theory. As we will see, for our computations of the vacuum energy and of particle production, both procedures yield the same result.
In the specification of operators in the bulk region we have neglected terms proportional to the slow velocity $`v`$ by which the circles shrink in the bulk metric (2.3)(2.1). Relatedly, $`\mu `$ could depend weakly on the other spatial directions $`\stackrel{}{X}_{}`$; we will ignore this for the purposes of the current discussion though it is simple to incorporate.
In fact the small velocity approximation will play an important role more generally in our analysis of the approach to the singularity. As it stands, the path integral (3.9)(3.10) does not extend over all values of $`X^0`$: the metric term $`G_{\mathrm{\Omega }\mathrm{\Omega }}`$ of classical GR goes to zero in finite $`X^0`$ in the past. As we will see, for a constant radius circle of size $`L`$, a winding tachyon condensate will produce a truncation of the support of amplitudes to a range of $`X^0`$ of order $`(\mathrm{ln}\mu )/\kappa `$ in the region of the condensate. We can arrange the parameters in our worldsheet CFT such that the velocity is sufficiently small that this range of time is far smaller than the time it takes to reach the singularity starting from a circle of size $`L`$. The basic idea is that the effective Newton constant does get large as the space shrinks, but the effects of the tachyon kick in first. Namely, consider a winding tachyon which turns on when the circle size is $`L`$. The corresponding value of $`\kappa `$ is $`\kappa =\sqrt{1(L/l_s)^2}`$. Set $`v`$ such that
$$\frac{L}{v}\frac{\mathrm{ln}\mu }{\kappa }$$
$`(3.11)`$
This specification, combined with the $`(\mathrm{ln}\mu )/\kappa `$ truncation of the amplitudes’ support in the $`X^0`$ direction to be derived below, yields a self-consistent perturbative string analysis. Note that the worldsheet potential term in the heterotic theory is classically always non-negative.
From this well-defined path integral (3.9)(3.10) we will obtain the $`\mu `$ dependence of amplitudes using methods developed for Liouville theory which also apply to our theory. This will enable us to read off the effect of the tachyon on the support of amplitudes, and will determine the spectrum of particles produced due to the time dependence of the tachyon background.
Type II Theory
Let us next briefly include the type II version of the above formulas. In the type II theory, we have (1,1) scalar superfields $`𝒳^\mu =X^\mu +\theta ^+\psi _+^\mu +\theta ^{}\psi _{}^\mu +\theta ^+\theta ^{}F^\mu `$. In terms of these, we have a Lorentzian signature path integral
$$G(\{V_n\})[d𝒳][d(\mathrm{ghosts})]d(\mathrm{moduli})e^{iS}\underset{n}{}\left(i𝑑\sigma 𝑑\tau V_n[𝒳]\right)$$
$`(\mathcal{3.12})`$
where the semiclassical action is
$$\begin{array}{cc}\hfill iS=& id\sigma d\tau d\theta ^+d\theta ^{}(D_{\theta ^+}𝒳^\mu D_\theta ^{}𝒳^\nu G_{\mu \nu }(𝒳)\mu :e^{\kappa 𝒳^0}\mathrm{cos}(w\stackrel{~}{\mathrm{\Omega }}):\hfill \\ & +(\mathrm{dilaton}))+iS_{\mathrm{ghost}}\hfill \end{array}$$
$`(3.13)`$
and $`V_n[𝒳]`$ are vertex operator insertions. As in the heterotic case, the form of the vertex operators is known in the flat space region to be of the form (3.6). The dilaton is (3.7).
Let us Wick rotate the worldsheet time coordinate $`\tau `$, the spatial target space coordinates $`\stackrel{}{X}(\sigma ,\tau )`$ (including $`\stackrel{~}{\mathrm{\Omega }}`$), and the parameters $`\mu `$ and $`\stackrel{}{k}`$ by
$$\tau e^{i\gamma }\tau _\gamma \stackrel{}{𝒳}e^{i\gamma }\stackrel{}{𝒳}_\gamma \mu =e^{i\gamma }\mu _\gamma \stackrel{}{k}=e^{i\gamma }\stackrel{}{k}_\gamma $$
$`(3.14)`$
where $`\gamma `$ is a phase which we will rotate from $`0`$ to $`\pi /2`$. This produces a Euclidean path integral for the worldsheet theory (where we label the quantities rotated to $`\gamma =\pi /2`$ by a subscript $`E`$)
$$G(\{V_n\})[d\stackrel{}{𝒳}_E][d𝒳^0][d(\mathrm{ghosts})]d(\mathrm{moduli})e^{S_E}\underset{n}{}()𝑑\sigma 𝑑\tau _EV_{n,i\stackrel{}{k}_E}[𝒳^0,i\stackrel{}{𝒳}_E]$$
$`(3.15)`$
with Euclidean action
$$\begin{array}{cc}\hfill S_E=& d\sigma d\tau _Ed\theta ^+d\theta ^{}(D_{\theta ^+}𝒳^0D_\theta ^{}𝒳^0+v^2(𝒳^0)^2D_{\theta ^+}\stackrel{~}{\mathrm{\Omega }}D_\theta ^{}\stackrel{~}{\mathrm{\Omega }}+G_{ij}D_{\theta ^+}𝒳_{,E}^iD_\theta ^{}𝒳_{,E}^j\hfill \\ & i\mu _Ee^{\kappa 𝒳^0}\mathrm{cosh}(w\stackrel{~}{\mathrm{\Omega }}_E)+(\mathrm{dilaton}))+S_E(\mathrm{ghost}).\hfill \end{array}$$
$`(3.16)`$
Again, as discussed below (3.10), we could include a $`\gamma `$-dependent shift of the linear dilaton coefficient (which vanishes in our physical critical-dimension Lorentzian bulk theory) if we wish to maintain semiclassical conformal invariance in the intermediate (rotated) steps of the calculations.
3.2. Vacuum Amplitude and Back Reaction
In this subsection we will present computations exhibiting the effect we advertised above that the amplitudes will be limited in their support to the weakly-coupled bulk. Let us start with the vacuum amplitude. At one loop, this is defined by the amplitude (3.15) with no vertex operator insertions, evaluated on a genus one worldsheet; let us call this amplitude $`Z_1`$. In a time dependent background, one must specify the vacuum in which the amplitudes are defined (for example, one definition of the S matrix would involve in-vacuum to out-vacuum amplitudes). We will return to the question of the vacuum after computing the first quantized path integral defined above at this 1-loop order.
In the bulk, this quantity describes a trace over spacetime single-particle states. In our case, the integral over the zero mode of $`X^0`$ will work differently than in flat space, and we will determine from this the support of the amplitudes as well as the quantum corrections to the stress energy in spacetime. In particular, as in Liouville field theory, we will find this amplitude to be supported only in the bulk region where the tachyon condensate is small. This supports the interpretation of the tachyon condensate as lifting the closed string degrees of freedom. Further, with our asymptotic supersymmetry in the bulk region this also provides a useful bound on the back reaction in the model. Finally, the imaginary part of the amplitude will provide information about the vacuum with respect to which the amplitude is being computed from the spacetime point of view.
Following \[36,,15\], let us compute first the quantity $`Z_1/\mu `$ and perform the path integral by doing the integral over the $`X^0`$ zero modes first. That is, decompose
$$X^0X_0^0+\widehat{X}^0(\sigma ,\tau _E)$$
$`(3.17)`$
where $`\widehat{X}^0`$ contains the nonzero mode dependence on the worldsheet coordinates $`\sigma ,\tau _E`$.<sup>4</sup> The reader should be grateful that we are suppressing the atomic number and baryon number indices on $`{}_{0}{}^{0}X_{0}^{0}`$. The path integral measure then decomposes as $`[dX^0]=dX_0^0[d\widehat{X}^0]`$. We obtain for heterotic and type II respectively
$$\begin{array}{cc}\hfill \frac{Z_1^{(Het)}}{\mu _E}=[d\stackrel{}{𝒳}_E][d\mathrm{\Psi }_{}]& [d(\mathrm{ghosts})]d(\mathrm{moduli})[d\widehat{𝒳}^0]dX_0^0\hfill \\ & \left(𝑑\sigma 𝑑\tau _E𝑑\theta ^+\mathrm{\Psi }_{}e^{\kappa 𝒳^0}i\mathrm{cosh}(w\stackrel{~}{\mathrm{\Omega }}_E)\right)e^{S_E}\hfill \end{array}$$
$`(3.18)`$
$$\begin{array}{cc}\hfill \frac{Z_1^{(II)}}{\mu _E}=[d\stackrel{}{𝒳}_E]& [d(\mathrm{ghosts})]d(\mathrm{moduli})[d\widehat{𝒳}^0]dX_0^0\hfill \\ & \left(𝑑\sigma 𝑑\tau _E𝑑\theta ^+𝑑\theta ^{}e^{\kappa 𝒳^0}i\mathrm{cosh}(w\stackrel{~}{\mathrm{\Omega }}_E)\right)e^{S_E}\hfill \end{array}$$
$`(3.19)`$
Decomposing $`e^{\kappa 𝒳^0}=e^{\kappa X_0^0}e^{\kappa \widehat{𝒳}^0}`$, we can change variables in the zero mode integral to $`ye^{\kappa X_0^0}`$ and integrate from $`y=0`$ to $`y=\mathrm{}`$ as $`X_0^0`$ ranges from $`\mathrm{}`$ to $`\mathrm{}`$<sup>5</sup> Note that the support of the integrand is negligible in the added region $`X_0^0[\mathrm{},0]`$. . For each point in worldsheet field space, the zero mode integral is of the form
$$_0^{\mathrm{}}𝑑ye^{Cy}=\frac{1}{C}$$
$`(3.20)`$
where the coefficient $`C`$ is the nonzeromode part of the tachyon vertex operator in $`S_E`$, integrated over worldsheet superspace.
Integrating over $`\theta ^\pm `$ produces a worldsheet potential term contributing to $`C`$. For regions of field space where $`C`$ is positive, the integral (3.20) converges. For regions of negative $`C`$ the equation (3.20) gives a formal definition of the function by analytic continuation. However, it is important to keep in mind the physical distinction between these two cases. As discussed above in the quantum field theory case, when $`C`$ is positive this corresponds to a time dependent massing up of modes, while negative $`C`$ corresponds to time dependent tachyonic masses. In the latter case, the formal definition (3.20) describes an analytic continuation of an interesting physical IR divergence.
In the heterotic theory, this coefficient $`C`$ is nonnegative everywhere in field space for $`\mu ^2>0`$, at least classically. Hence the computation (3.20) applies directly.
In the type II theory, this coefficient can become negative near particular points in $`\stackrel{~}{\mathrm{\Omega }}`$ and $`\stackrel{}{X}_{}`$. In regions where the potential is positive, (3.20) applies, and as we will see leads to a truncation of the support of the closed string states. However, in regions where $`C`$ is negative, there are physical instabilities remaining. These localized instabilities we interpret as subcritical type 0 tachyons. In particular, in §4.1 we will see using linear sigma model techniques that the GSO projection acts on the corresponding subcritical theory as in type 0.
This analysis yields
$$\frac{Z_1}{\mu _E}=\frac{1}{\kappa \mu _E}\widehat{Z}_1$$
$`(3.21)`$
where $`\widehat{Z}_1`$ is the partition function in the free theory (with no tachyon term in the action) and with no integral over the zero mode of $`X^0`$. Referring to the functional measure for the rest of the modes (including all fields) as $`[d(\mathrm{fields})]^{}`$ this is
$$\widehat{Z}_1=[d(\mathrm{fields})]^{}[d(\mathrm{ghosts})]d(\mathrm{moduli})e^{\widehat{S}_E}$$
$`(3.22)`$
where $`\widehat{S}_E`$ is the action ((3.10) and (3.16) respectively for heterotic and type II) with $`\mu =0`$. Finally, integrating with respect to $`\mu `$ yields the result for the 1-loop partition function
$$Z_1=\frac{\mathrm{ln}(\mu _E/\mu _{})}{\kappa }\widehat{Z}_1.$$
$`(3.23)`$
Here $`\mu _{}=e^{\kappa X_{}^0}`$ where $`X_{}^0`$ is an IR cutoff on the $`X^0`$ field space in the free-field region. As discussed above, this is valid for regions of the worldsheet field space where the potential is positive, which is always true in the Heterotic case and true for most contributions in type II.
It is worth remarking here that this result follows simply from the exponential form of the potential in the tachyon vertex operator. The above derivation formally goes through whether rotated back to target space Minkowski signature or not, (since the zeromode does not appear in the kinetic term), reducing the amplitude to a tachyon-free computation $`\widehat{Z}`$. Relatedly it is insensitive to the issue of rotating in $`Q`$ as discussed below (3.10).
To interpret this result, recall that in a background of $`d`$-dimensional flat space, the partition function scales like the volume of spacetime: integration over the zero modes of the $`X^\mu `$ fields yields the factor $`\delta ^{(d)}(0)=V_d=V_0V_{d1}`$ where $`V_{d1}`$ is the volume of space and $`V_0`$ is the volume of the time direction. In our present case (3.23), the spatial extensivity reflected in the factor $`V_{d1}`$ is still present. But the volume of time $`V_0`$ has been truncated to $`\frac{1}{\kappa }\mathrm{ln}\mu /\mu _{}`$. This corresponds to the range of $`X^0`$ where the tachyon is absent. Again, this is similar to the situation in spacelike Liouville field theory , where the Liouville wall cuts off the support of the partition function.
This result has several implications. First, it provides a concrete verification that the string states are lifted in the tachyon phase, for positive worldsheet potential, supporting the interpretation of this phase as a theory of Nothing. Combined with the specification (3.11) this result justifies the use of the worldsheet path integral with metric coefficients going to zero in finite time, as the amplitudes are not supported in this region. Note that in particular all states are lifted – the would-be tachyon and graviton fluctuations are absent and hence back reaction is suppressed. As emphasized above, we are computing in a particular state analogous to the Euclidean or Hartle-Hawking vacuum, with no excitations in the far past. Putative vacua with excitations in the far past will be analyzed separately . The analagous context of quantum field theory with exponentially growing mass again suggests that the string states will become lifted (effectively go to infinite mass) in an appropriate sense; moreover an interesting BRST anomaly appears which may limit the consistency of such vacua in the string theory context. (See for some discussion of the issue in field theory).
Second, it indicates that the 1-loop vacuum energy is only supported in the bulk region of the spacetime. Because the asymptotic bulk region $`X^0\mathrm{}`$ is weakly coupled and weakly curved (in fact approximately supersymmetric), this means that back reaction is restricted to the intermediate region where the tachyon $`T`$ is of order 1.
Third, the imaginary part of the 1-loop partition function is significant and will provide an important check on the consistency of our computations. Recall that the analytic continuation (3.14) included a rotation $`\mu =e^{i\pi /2}\mu _E`$. This means that for real values of our original parameter $`\mu `$, the partition function has an imaginary part:
$$Z_1=\left(\frac{1}{\kappa }\mathrm{ln}\frac{\mu }{\mu _{}}i\frac{\pi }{2\kappa }\right)\widehat{Z}_1.$$
$`(3.24)`$
Note that the IR cutoff $`X_{}^0`$ is always real in our prescription. We will interpret this as indicating that the system is in a pure state with thermal occupation numbers corresponding to a temperature $`\kappa /\pi `$, a result which will also follow from an analysis of the 2 point amplitudes at genus zero. Let us return to this interpretation after explaining those computations in the next subsection.
In general, it would be interesting to unpack the 1-loop amplitudes in more detail, to follow the fate of the various closed string states and D-branes in our background. An important aspect of this is mode mixing induced by the tachyon operator: the oscillator modes in the bulk generally mix under the action of the tachyon term in the region where it is substantial.
It might be possible to analyze this using the ideas in . In the type II case, a similar compuational technique to the one we have described above applies to the amplitudes of open strings in this background, for example the 1-loop open string partition function. The closed string channel of such amplitudes describes the response of the would-be graviton and other closed string modes to D-brane sources. It is necessary to specify consistently the boundary conditions defining the D-branes in this background, but it seems likely that the $`X_0^0`$ integral will again reveal that these amplitudes are shut off in the tachyon phase. We note that in the spacelike Liouville theory, the ZZ-brane is localized under the tachyon barrier, and has a paucity of degrees of freedom. It cannot move; basically it can only decay. It would be very interesting to understand the conformal boundary states in the timelike case.
3.3. 2-point function, particle production, and Euclidean State
Let us now include vertex operator insertions. The $`\mu `$-dependence of amplitudes can be determined by a similar technique to that above. We analyze the derivative of the correlation function (3.15)<sup>6</sup> Note that as in LFT, we use the semiclassical form of the vertex operators and dilaton as well as of the the action. with respect to $`\mu _E`$ by doing the integral over $`X^0`$’s zero mode $`X_0^0`$ first. From that we can determine its dependence on $`\mu _E`$, and finally use (3.14) to determine its dependence on $`\mu `$.
This is similar to the above computation of the partition function, except now the integral over $`y=e^{\kappa X_0^0}`$ (which gave (3.20) in the case of the partition function) is of the form
$$\frac{G(\{V_{n,i\stackrel{}{k}_E}\})}{\mu _E}=[d\widehat{𝒳}^0][d\stackrel{}{𝒳}][d(\mathrm{ghost})]𝑑yy^{_ni\frac{\omega _n(\stackrel{}{k}_n)}{\kappa }}e^{Cy}e^{\widehat{S}_E}\frac{C}{\mu _E}$$
$`(3.25)`$
where $`C`$ is defined after (3.20). This yields a result for $`G(\{V_{n,i\stackrel{}{k}_E}\})`$ proportional to
$$\mu _E^{i_n\omega _n/\kappa }$$
$`(3.26)`$
times a complicated path integral over nonzero modes, which would be difficult to evaluate directly, but which is independent of $`\mu _E`$. As discussed below (3.10), we could rotate in a way that introduces a linear dilaton in the Euclidean continuation. This would shift the $`\omega `$-dependent exponent by a term linear in $`Q`$, which would rotate back to zero when we rotate $`\mu _E`$ back to $`\mu `$.
Fortunately, in the case of the 2-point function, we can use a simple aspect of the analytic continuation we used to define the path integral to determine the magnitude of the result. As explained for example in , the two-point function of two negative frequency modes in the bulk is
$$G(\stackrel{}{k},n;\stackrel{}{k}^{},n^{})=\delta _{nn^{}}\delta (\stackrel{}{k}\stackrel{}{k}^{})\frac{\beta _{\stackrel{}{k},n}}{\alpha _{\stackrel{}{k},n}}$$
$`(3.27)`$
where $`\alpha _{\stackrel{}{k},n}`$ and $`\beta _{\stackrel{}{k},n}`$ are the Bogoliubov coefficients describing the mixing of positive and negative frequency modes. This is the timelike Liouville analogue of the reflection coefficients describing the mixing of positive and negative momentum modes bouncing off a spacelike Liouville wall.
In fact, this relation is precise here, and we can determine the magnitude $`|\beta _\omega /\alpha _\omega |`$ as follows. As we discussed above for the partition function, after performing the Euclidean path integral defined via the rotations (3.14), we must continue back $`\mu i\mu `$ in order to obtain the amplitude of interest. The regions where the worldsheet potential is positive translate in the Euclidean path integral to a positive Liouville wall. For these regions, the Euclidean 2-point function is a reflection coefficient of magnitude 1. The continuation above (3.8)(3.14) in $`\mu `$,
$$\mu =e^{i\frac{\pi }{2}}\mu _E$$
$`(3.28)`$
therefore yields a 2-point function for the Lorentzian theory of magnitude
$$\left|\frac{\beta _{\stackrel{}{k},n}}{\alpha _{\stackrel{}{k},n}}\right|=e^{\omega (\stackrel{}{k},n)\pi /\kappa }.$$
$`(3.29)`$
Using the relations $`|\alpha _\omega |^2|\beta _\omega |^2=1`$ for bosonic and fermionic spacetime fields, and the fact that the number of particles produced $`N_{\stackrel{}{k},n}`$ is given by $`|\beta _{\stackrel{}{k},n}|^2`$, this result translates into a distribution of pairs of particles of a thermal form
$$N_{\stackrel{}{k},n}=\frac{1}{e^{2\pi \omega (\stackrel{}{k},n)/\kappa }1}.$$
$`(3.30)`$
This corresponds to a Boltzmann suppression of the distribution of pairs of particles (each pair having energy $`2\omega `$) by a temperature $`T=\kappa /\pi `$. This temperature can also be deduced from the imaginary part of the 1-loop partition function (3.24), providing a check on the calculations, as follows.
At the level of the genus zero two-point functions, the system is in a pure state whose phase information we have not computed, but whose number density is thermal. We would like to understand how the result (3.24) fits into this thermal interpretation. To explain this, for simplicity let us work in a field theory limit in the bulk, so that we can express the state in terms of the bulk Fock space of states built on the bulk vacuum $`|\mathrm{bulk}`$ killed by the annihilation operators $`a_\stackrel{}{k}`$. Ignoring interactions, the bulk state is of the form of a squeezed state
$$|\mathrm{\Psi }_0=𝒩\mathrm{exp}\left(\underset{\stackrel{}{k}}{}e^{i\gamma _\stackrel{}{k}}e^{\omega (\stackrel{}{k})\kappa /\pi }a_\stackrel{}{k}^{}a_\stackrel{}{k}^{}\right)|\mathrm{bulk}$$
$`(3.31)`$
where $`𝒩`$ is a normalization factor and the $`\gamma _\stackrel{}{k}`$ are the phases arising in the two point function. Expanding the exponent, this is of the form
$$|\mathrm{\Psi }_0=𝒩\underset{n}{}e^{i\gamma _n}e^{\frac{\beta _T}{2}E_n}|E_n$$
$`(3.32)`$
for some real phases $`\gamma _n`$, where $`n`$ indexes the Fock space states and $`E_n`$ the corresponding energies, and where $`\beta _T=\pi /\kappa `$ is the inverse temperature.
Switching back to the Schrodinger picture for simplicity, the bulk state can be thought of then as evolving as if the time evolution operator, normally $`U(t,0)=T\left(\mathrm{exp}\left(i_0^t𝑑t^{}H(t^{})\right)\right)`$, is now $`T\left(\mathrm{exp}\left(i_C𝑑t^{}H(t^{})\right)\right)`$ where the contour $`C`$ runs from 0 to $`t`$ along the real axis, and then runs vertically from $`t`$ to $`ti\beta _T/2`$.<sup>7</sup> This is the first half of the contour obtained in real time thermal field theory in a mixed thermal state; see for example section 5 of the second reference in for a recent discussion.
Now let us incorporate the 1-loop vacuum energy correction (3.24). This calculates the zero point energy contribution to the spacetime Hamiltonian, integrated over time. As we have just seen, in a state with thermal occupation numbers, the effective Hamiltonian evolution up to time $`t`$ is on a time contour with an extra segment from $`t`$ to $`ti\beta _T/2`$. The result (3.24) arises from the bulk vacuum result via exactly such a shift, with $`\beta _T=\pi /\kappa `$ corresponding to a temperature $`T=\kappa /\pi `$. Namely, as discussed in §3.2, the partition function is the summed zero-point energy in spacetime, times the volume of time: $`Re(Z)=\mathrm{\Lambda }V_{d1}V_0`$. Said differently, $`Z=𝑑t\mathrm{\Lambda }V_{d1}`$ where $`t`$ is the target-space time direction. The imaginary part of our amplitude (3.24) is obtained from the bulk vacuum by shifting the spacetime Hamiltonian evolution by $`\mathrm{\Lambda }V_{d1}X^0\mathrm{\Lambda }V_{d1}(X^0i\frac{\pi }{2\kappa })`$. So the same temperature arises in both the genus zero and the 1-loop corrected amplitudes in our prescription, providing a check.
Altogether, this yields the following simple result. Let us consider the big bang case, with the tachyon condensate turned on in the past. Modulo expected subcritical tachyons in the type II case, the closed string states are lifted in the far past (and in the type II case, we expect the subcritical tachyons to also condense and lift degrees of freedom). Start with no excitations above this tachyon background (perhaps a natural choice given the enormous effective masses in this region). The state in the bulk $`X^0\mathrm{}`$ region has a thermal distribution of pairs of particles (3.30), with temperature $`\kappa /\pi `$. These pairs are created during the phase where the tachyon condensate is order one<sup>8</sup> Indeed, the time-dependence of the Hamiltonian is only non-adiabatic $`1\frac{\dot{\omega }}{\omega ^2}=\frac{\mu ^2\kappa e^{\kappa t}}{\omega ^3}`$ in a small window of time near $`t1/\kappa `$. Similar suppression obtains for other measures of nonadiabaticity $`\frac{_t^n\omega }{\omega ^{n+1}}\stackrel{t\mathrm{}}{}(\kappa /\mu )^ne^{nt/2}`$ ., and hence the calculation is self-consistent if we tune the bare dilaton to weak coupling.
This choice of state is analogous to the Hartle-Hawking, or Euclidean, State in the theory of quantum fields on curved space, but it arises here in a perturbative string system via crucially stringy effects. In quantum field theory on curved space, the Euclidean vacuum is obtained by calculating Greens functions in the Euclidean continuation of the spacetime background (when it exists) and continuing them back to Lorentzian signature. In our case, a similar continuation has been made, but here the Euclidean system is a spacelike Liouville field theory. The choice of vacuum (nothing excited above the tachyon condensate) is natural from the point of view of the spacelike continuation, as it involves excitations of modes with only one sign of frequency in the tachyon phase; these Wick rotate to exponentially dying modes under the Liouville wall.
3.4. Singularity Structure
So far we focused on two particularly instructive physical quantities: the 1-loop partition function and the genus zero two-point function (Bogoliubov coefficients). Let us now determine the singularity structure of more general amplitudes. This is important in order to complete our assessment of the ability of the tachyon condensate to resolve the spacelike singularity. Namely, if the perturbative amplitudes are finite up to expected divergences associated with physical states (which we will make precise below), then we may conclude that the perturbative string theory is capable of resolving the singularity in the circumstances we have specified (in particular, at weak coupling).
N-point functions at genus zero
As discussed above, the genus zero two-point function describes particle production in the linearized spacetime theory. The singularity structure of general $`N`$-point amplitudes can be ascertained from the path integral (3.9)(3.15). In a nontrivial bulk vacuum, such as that derived above in the Euclidean vacuum (3.30), we are interested in a linear combination of vertex operators (3.6) comprised of $`\alpha _{\stackrel{}{k},n}`$ times a negative frequency component plus $`\beta _{\stackrel{}{k},n}`$ times a positive frequency component.
The path integral diverges when a bosonic degree of freedom can go off to infinity unobstructed by the $`e^{S_E}`$ factor. As discussed in the introduction, this situation appears in the big bang region of the spacetime in the naive extrapolation of GR, as the space shrinks and the kinetic terms in $`S_E`$ go away. In our case, where it is positive the tachyon term obstructs this divergence (everywhere in the Heterotic case, and away from the subcritical type 0 regions of the type II system).
There are divergences in the bulk region $`X^0\mathrm{}`$ that are expected in a time dependent S matrix. Generically, the vertex operators provide oscillations suppressing the path integral contribution in this region. However a divergence in $`X_0^0+\mathrm{}`$ appears when the frequencies are such as to cancel this oscillation:
$$\omega _n(\pm 1)_n=0$$
$`(3.33)`$
The $`\pm `$ sign here comes from the presence of both positive and negative frequency modes in the vertex operators.<sup>9</sup> This is another aspect analogous to the situation in Liouville theory (see eqn. (87) of ). These divergences correspond to expected divergences from physical intermediate states in time dependent systems (see e.g. chapter 9 for a discussion of this).
In \[20,,41,,22\] subtleties were encountered in a prescription for analytically continuing higher point amplitudes from spacelike Liouville theory. It will be interesting to study these issues in more detail in our case, keeping careful track of the vacuum in which the computations are being done in the various prescriptions for continuing the amplitudes.
Higher loops
Higher loop amplitudes arise from the path integral (3.9)(3.15) defined on a Riemann surface of higher genus $`h`$. These contain dependence on the dilaton $`\mathrm{\Phi }\mathrm{\Phi }_0+\widehat{\mathrm{\Phi }}`$ where $`\mathrm{\Phi }_0`$ is the constant value in the bulk region. This introduces a factor of $`e^{\mathrm{\Phi }_0(2h2)}`$ from the bare bulk string coupling as well as a contribution
$$S_\mathrm{\Phi }=_\mathrm{\Sigma }\widehat{}^{(2)}\widehat{\mathrm{\Phi }}[X^0]$$
$`(3.34)`$
(plus its supersymmetric completion in the Heterotic and Type II cases). Semiclassically $`\mathrm{\Phi }=\mathrm{\Phi }_0`$ (i.e. $`\widehat{\mathrm{\Phi }}=0`$) as discussed in (3.7). The dilaton will ultimately get sourced by the tachyon. The corresponding corrections will be generated by the worldsheet path integral, but are suppressed by powers of $`e^{\mathrm{\Phi }_0}`$. Moreover, as in our analysis of the 1-loop vacuum amplitude, the $`X_0^0`$ integral reveals that higher genus amplitudes have support limited to the weakly coupled bulk of spacetime.
4. Positively-curved spatial slices
In this section, we generalize our techniques to strings in geometries of the form (1.1) where the $`\mathrm{\Omega }`$ are coordinates on higher dimensional spheres. The worldsheet theory will be described by an $`O(N)`$ model at an energy scale related to $`X^0`$ in a way we will specify. In this case there is no topologically-stabilized winding tachyon<sup>10</sup> It might be interesting to consider examples of positively-curved spaces with nonzero $`\pi _1`$ such as $`S^N/\mathrm{\Gamma }`$ with freely acting $`\mathrm{\Gamma }`$. . The sigma model on spatial slices nevertheless develops a mass gap. We will frame this fact in terms of the discussion of §1, and investigate the extent to which it can be used to remove the singularity present in the GR approximation. Some aspects of the analysis of §3 persist. Unlike in the case of flat spatial slices, however, the back-reaction from the velocity of the radion will be harder to control in these examples.
4.1. The mass gap of the $`O(N)`$ model
Consider $`N`$ two-dimensional scalar fields arranged into an $`O(N)`$ vector $`\stackrel{}{n}`$. The partition function of the $`O(N)`$ model is
$$Z=[dn]e^{{\scriptscriptstyle d^2zR^2(_\mu \stackrel{}{n})^2}}\underset{z}{}\delta (n^2(z)1)$$
$`(4.1)`$
A nice way to see the mass term appear is to use a Lagrange multiplier to enforce the delta function localizing the path integral onto a sphere, and large $`N`$ to simplify the resulting dynamics (see e.g. ):
$$Z=[dn][d\lambda ]e^{{\scriptscriptstyle d^2z[R^2\stackrel{}{n}(^2+i\lambda )\stackrel{}{n}+i\lambda ]}}$$
$`(4.2)`$
where $`\lambda `$ is the Lagrange multiplier field introduced to represent the delta function. Now integrate out $`n`$:
$$Z=[d\lambda ]e^{N/2\mathrm{tr}\mathrm{ln}(^2+\lambda )+R^2{\scriptscriptstyle d^2z\lambda }}.$$
$`(4.3)`$
At large $`N`$, the $`\lambda `$ integral has a well-peaked saddle at
$$\lambda (x)=im^2$$
$`(4.4)`$
where the mass $`m`$ satisfies
$$R^2=N^\mathrm{\Lambda }\frac{d^2k}{(2\pi )^2}\frac{1}{k^2+m^2}=\frac{N}{2\pi }\mathrm{ln}\frac{\mathrm{\Lambda }}{m}.$$
$`(4.5)`$
Renormalize by defining the running coupling at the scale $`M`$ by
$$R^2(M)=R_0^2+\frac{N}{2\pi }\mathrm{ln}\mathrm{\Lambda }/M.$$
$`(4.6)`$
Plugging back into the action for $`n`$, we have a mass for the $`n`$-field which runs like
$$m=Me^{\frac{2\pi R^2}{N}}.$$
$`(4.7)`$
An alternative UV completion of the model which is sometimes more convenient (and easier to supersymmetrize) gives $`\lambda `$ a bare mass: add to the action
$$\delta S=a\lambda ^2.$$
for a large parameter $`a`$. Integrating out $`\lambda `$, this smoothens the delta function, and imposes the $`n^2=R^2`$ relation weakly in the UV by a quartic potential.
Supersymmetric $`O(N)`$ model
Since we wish to study string theories without bulk tachyons, we will need to understand the supersymmetric version of the model. A (1,1) supersymmetric version of the $`O(N)`$ model has an action
$$S=d^2\theta \left(ϵ_{\alpha \beta }D_\alpha nD_\beta n+\mathrm{\Lambda }(n^2R^2)+a\mathrm{\Lambda }^2\right);$$
$`\alpha ,\beta =\pm `$, $`\mathrm{\Lambda }=\lambda +\theta ^\alpha \psi _\alpha +\theta ^2F_\lambda `$ is now a Lagrange multiplier superfield, and $`D_\alpha =\frac{}{\theta ^\alpha }+i\theta ^\beta \sigma _{\beta \alpha }^\mu _\mu .`$ Note that the type II GSO symmetry acts on $`\mathrm{\Lambda }`$ as
$$(1)^{F_L}:\mathrm{\Lambda }\mathrm{\Lambda },(1)^{F_R}:\mathrm{\Lambda }\mathrm{\Lambda }.$$
$`(4.8)`$
The large $`N`$ physics is the same as in the bosonic case (see e.g. ), exhibiting a mass gap, except now there are two vacua for $`\mathrm{\Lambda }`$. When
$$\lambda =\pm m$$
$`(4.9)`$
the GSO symmetry is spontaneously broken; the two vacua are identified by the GSO projection. This is just as in the appendix of , and it results in a single type zero vacuum. This statement about the GSO projection applies to all $`N`$, including the case $`N=2`$ of a shrinking circle described in §2, §3. In particular, this justifies the comment made in the discussion above (3.21) that the regions of negative potential in the type II worldsheet have a type 0 subcritical GSO projection.
4.2. $`\text{ }\mathrm{C}\mathrm{IP}^1`$ Model
The $`(1,1)`$ sigma model on $`S^2`$ actually has $`(2,2)`$ supersymmetry. Consider a $`(2,2)`$ linear sigma model for it. There are two chiral superfields $`Z^i`$ each with charge one with respect to a single $`U(1)`$ vectormultiplet. The D-term equation is
$$0=\underset{i=1}{\overset{2}{}}|Z^i|^2\rho .$$
$`(4.10)`$
Below the scale $`e`$ of the gauge coupling, this model describes strings propagating on a 2-sphere of radius $`R=\sqrt{\rho \alpha ^{}}`$. The FI coupling $`\rho `$ flows logarithmically towards smaller values in the IR:
$$\rho (M)=\rho _02\mathrm{ln}\frac{M}{M_0}.$$
$`(4.11)`$
This breaking of scale invariance is in the same $`(2,2)`$ supermultiplet as an anomaly in the chiral $`U(1)`$ R-symmetry; only a $`\mathrm{ZZ}_2`$ subgroup of this latter group is a symmetry of the quantum theory (this is part of the GSO symmetry in type II theories).
Integrating out the chiral multiplets $`Z^i`$ leads to an effective twisted superpotential for the vectormultiplet scalar
$$\stackrel{~}{W}=2\mathrm{\Sigma }\mathrm{ln}\mathrm{\Sigma }t\mathrm{\Sigma }.$$
$`(4.12)`$
Mirror symmetry relates this to a model with one twisted chiral superfield $`Y`$, governed by a twisted superpotential
$$\stackrel{~}{W}=\mathrm{\Lambda }\left(e^Y+e^Y\right)$$
$`(4.13)`$
where $`\mathrm{\Lambda }=me^{t/2},t=\rho +i\vartheta `$. This effective twisted superpotential has isolated massive vacua.
Next let us discuss the GSO projection, to ensure that the relevant operator we are generating is present in the type II theory (as opposed to being a type 0 bulk tachyon). In the type II case, the twisted chiral superpotential must be odd under the chiral GSO actions. This is accomplished by
$$(1)^{F_L}:Y_2Y_2+\pi ,(1)^{F_R}:Y_2Y_2+\pi ,$$
$`(4.14)`$
where $`YY_1+iY_2`$, or simply $`\mathrm{\Sigma }\mathrm{\Sigma }`$.
The twisted superpotential (4.12) has two massive vacua
$$\sigma =\pm e^{t/2},$$
$`(4.15)`$
which are permuted by the GSO action. The condensate (4.13) is therefore not a bulk tachyon mode.
We can use this mirror description to further elucidate the physical interpretation of the condensate. It is invariant under the $`SU(2)SO(3)`$ rotations of the $`S^2`$. Since $`Y_2`$ is the variable T-dual to the phase of the $`Z`$s, from the point of view of the original linear model, (4.13) represents a condensation of winding modes. It is tempting to interpret this as a condensate resembling a ball of rubber-bands wrapping great circles of the small sphere.
A special RG trajectory
In the case $`N=3`$ of the two-sphere, where there is a two-cycle in the geometry, there are topologically-charged worldsheet instantons. The contribution to the sum over maps of the sector with winding number $`n`$ is weighted by $`e^{in\theta }`$ where $`\theta =_{S^2}B`$ is the period of the NSNS B-field through the two-sphere.
When $`\theta =\pi `$ this introduces wildly fluctuating signs in the path integral which can result in a critical theory in the IR. In fact, the model flows to the $`SU(2)`$ WZW model at level one, also known as a free boson at the self-dual radius. This model has a topologically-stabilized winding mode, which is exactly massless. At this point, the evolution may be glued onto the analysis of §2,3.
4.3. Coupling to string theory
Eternal nothingness is fine if you happen to be dressed for it.
-- Woody Allen
We need to check whether the mass gap whose origin we have reviewed takes effect before large curvature develops. In the example of §2,3, the rate of shrinking $`_tR`$ of the circle was a tunable parameter which we used to control the collapse. In this case, where the spatial curvature exerts a force on $`R`$, we will need to reevaluate the behavior of $`R(t)`$.
In order to do this, we begin at large radius, and use the fact that in this regime, the beta function equations for the worldsheet theory are the same as the gravity equations of motion.
In the case of positive spatial curvature, the Friedmann equation (the Hamiltonian constraint) requires a stress-energy source which dominates over the curvature contribution:
$$\left(\frac{\dot{R}}{R}\right)^2=\frac{1}{R^2}+G_N\rho $$
$`(4.16)`$
where $`\rho `$ is the energy density in non-geometrical sources. The curvature term $`1/R^2`$ alone, in the absence of the term from extra sources $`\rho `$, would not yield consistent initial data; instead the source term must dominate over the curvature term in the large radius general relativistic regime. This means that unlike the previous two cases of §2.1,2.2, we do not classically have a tunable parameter allowing us to slow down the approach to the would-be singularity.
Inclusion of matter
As we mentioned in §2.3, in the case of positive spatial curvature, the Friedmann equation (4.16) has no real solutions in the absence of matter. We will overcome this problem by including some nonzero radiation energy density on the RHS of (4.16). With $`\rho =x/R^N`$ ($`x`$ is a constant and $`N`$ is the number of spacetime dimensions participating in the FRW space), the maximum radius reached is
$$R_{\mathrm{max}}=l_Px^{\frac{1}{N2}},$$
$`(4.17)`$
where $`l_P=G_N^{\frac{1}{N2}}`$ is the $`N`$-dimensional Planck length. In the curvature-dominated regime,
$$R(t)R_{\mathrm{max}}^{\frac{N2}{N}}t^{\frac{2}{N}}.$$
$`(4.18)`$
Now we can estimate the time at which the mass gap takes effect. For convenience (as opposed to phenomenology), consider the case $`N=4`$, where (4.18) implies $`R(t)\sqrt{tR_{max}}`$. Semiclassically, the tachyon term in the worldsheet effective action (4.7) depends on time via the “tachyon” profile
$$𝒯(t)\mu e^{R^2/N\alpha ^{}}.$$
$`(4.19)`$
If we assume that the leading effect of the radiation that we added is to fix the evolution of the scale factor (i.e. that it does not couple significantly to the Liouville mode in any other way), we can make a similar estimate to those of the previous sections. In fact, for $`N=4`$, the zeromode integral over $`X^0`$ is of the same form as (3.23)
$$Z\mathrm{ln}\mu .$$
$`(4.20)`$
$`X^0`$ goes to zero at the would-be bang singularity. The range of $`X^0`$ for which the amplitudes have support is $`X^0>\frac{1}{T}\mathrm{ln}\mu /\mu ^{}`$. Increasing $`\mu `$ makes the range of $`X^0`$ support of amplitudes smaller. So if we take $`\mu `$ to be large, we can ensure that the lifting of modes occurs in a regime where the kinetic terms have not yet died as we approach the singularity.
Physically, this parameter $`\mu `$ determines the amplitude of the oscillating mode in the bulk and hence its initial behavior in its exponential regime. We are introducing a classical solution with a large amplitude condensate of tachyon even in the initial “bulk” region where the space is larger than string scale, but we expect that these modes will decay once the other states come down.
Particle Production
From here the analysis proceeds as in §3, but the absence of a tunably small rate of growth of the $`S^3`$ space leads to a much larger density of produced closed strings. In particular, we again obtain an effective temperature via the periodicity in imaginary time of the condensate. Using the definition
$$𝒯(t)e^{Tt}$$
for the effective temperature $`T`$, we find (again, for $`N=4`$)
$$T\frac{R_{max}}{\alpha ^{}}.$$
Thus, when the cosmology has a phase during which it is bigger than string scale, the effective temperature is larger than the Hagedorn temperature. This is in contrast to the tunably small value of $`\kappa `$ we obtained in (3.30) in the case of flat spatial slices.
The upshot is that in this case of positively curved spatial slices, although the mass gap lifts the would-be GR divergences in the worldsheet path integral, a new source of back reaction is generated through copious particle production. It is worth emphasizing that the GR solution alone will lead to particle production of momentum modes, whose back reaction may also correct the background in an important way. We leave this analysis and its potential application to Schwarzchild black hole physics to further work .
5. Discussion
Application to Black Hole Physics
Spacelike singularities appear inside generic black hole solutions of general relativity. The case of a shrinking $`S^2`$ described in §4.1 appears inside the horizon of the Schwarzchild black hole solution in four dimensions (with an additional spatial direction $`t`$ which is stretching at the same time)
$$ds^2=(1L_S/r)dt^2+\frac{dr^2}{1L_S/r}+r^2d\mathrm{\Omega }^2$$
$`(5.1)`$
where $`L_S`$ is the Schwarzchild radius. Inside the horizon ($`r<L_S`$), $`r`$ is a timelike coordinate.
When the $`S^2`$ parameterized by $`\mathrm{\Omega }`$ shrinks, the worldsheet path integral develops contributions arising from the mass gap of the corresponding sigma model as discussed in §4.1. It would be interesting to understand if this might clarify black hole dynamics . These results may also apply to the proposal of , where the possibility of postselection on a “nothing” state was explored. The unitarity required in may arise from the unitary evolution along the $`t`$ direction inside the horizon, generated by the momentum generator inside the horizon.
Other vacua and the shape of the S-matrix
We have focused on a vacuum with no extra excitations above the tachyon background in the initial state. This is motivated by the lifting of closed string degrees of freedom in the presence of the tachyon. However, it would be very interesting to understand if other states are allowed. The analogy with quantum field theory with exponentially growing mass suggests a similar “massing up” of closed string modes in tachyon backgrounds in more general states. Moreover some interesting potential consistency conditions which may restrict the range of allowed states suggest themselves. In particular, a propagator entering the tachyon phase (rather than annihilating against another mode) has the property that semiclassically its worldsheet time ends at a finite value as $`X^0`$ goes to infinity. A similar issue arises in formulating the S matrix in the case of pure quantum field theory with an exponentially growing mass (see for a discussion of the issues there); in a generic vacuum there is a sense in which time ends prematurely in this system due to the absence of well defined on shell asymptotic states. Some aspects of the specification of the spacetime quantum state from the worldsheet point of view have been discussed in \[35,,49\].
In particular, one of the main questions raised by spacelike singularities is that of predictivity. In field theory or GR on a background with a putative big bang singularity, the initial conditions on the fields in the bulk region are ambiguous. If there are other consistent states of the system involving some extra excitations introduced initially and becoming light as the tachyon turns off, then the singularity, while resolved, will not be arbitrarily predictive. It is important to understand the status of all possible states.
Big crunch
Our main computations were done in the vacuum discussed in §3 with no excitations above the tachyon condensate. In the case of the big bang, this is perhaps a natural choice of initial state. In the case of the big crunch, the methods employed in this paper are not yet sufficient to answer the question of what happens starting from an arbitrary initial state. For example, it is interesting to ask what happens if we start with no particles in the bulk. At the level of the genus zero diagrams, we can accomplish this by considering correlation functions of vertex operators which are nontrivial linear combinations of positive and negative frequency modes in the bulk. For the 1-loop and higher genus diagrams, it is an open question here (as in the case of open string tachyons) how the different vacua translate into different prescriptions for the worldsheet path integral.
One aspect of the system is pair production of winding modes themselves as they become massless ; this can drain energy from the rolling radius to some extent .
Negative spatial curvature
We have discussed the cases of $`k=1`$ and $`k=0`$. It is natural to ask about the case $`k=1`$ where the spatial sections have negative curvature. In this case, at large radius $`R`$, the system expands according to the simple relation $`RX^0`$. Localized tachyon dynamics in some such examples were discussed in .
The big bang singularity in the far past in this case is not related by RG flow toward the IR in the matter sector of the corresponding worldsheet sigma model. The direction of flow is opposite; the small radius big bang regime corresponds to the UV. Hence in this case, the big bang resolution may depend on the appropriate UV completion of the sigma model on negatively curved spatial slices.<sup>11</sup> One can alternatively add ingredients to metastabilize the system away from this difficult regime .
Cosmology
It will be interesting to see if these methods and results translate into concrete results for more realistic string cosmology. Inflation tends to dilute information about the big bang singularity, but depending on the level of predictivity of the singularity, it may nonetheless play a role. Stretched strings play an important role in our mechanism for resolving the singularity: perhaps there is some relation between them and late-time cosmic strings (whether inside or outside the horizon at late times).
Toward a theory of Nothing
It is the silence between the notes that makes the music;
it is the space between the bars that holds the tiger.
-- Anonymous
Our calculations using methods borrowed from Liouville theory exhibit truncation of support of amplitudes to the bulk of spacetime, and hence concretely support the notion of a “Nothing” phase in the regime of the tachyon condensate. Conversely, spacetime emerges as the tachyon turns off.
It would be very interesting to characterize this phase and its onset in more detail, for example by unpacking the partition function to analyze its individual contributions. Although perturbative methods exhibit the basic effect, perhaps there is some dual formulation for which the emergence of time as the tachyon turns off is also built in.
Acknowledgements
We would like to thank A. Adams, M. Berkooz, S. Kachru, X. Liu, A. Maloney, H. Ooguri, J. Polchinski, N. Seiberg, S. Shenker, A. Strominger for helpful discussions, and Gary Horowitz for early collaboration and very helpful discussions. We are supported in part by the DOE under contract DE-AC03-76SF00515 and by the NSF under contract 9870115.
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warning/0506/nucl-th0506058.html | ar5iv | text | # Hadron-string cascade versus hydrodynamics in Cu + Cu collisions at √𝑠_{𝑁𝑁}=200 GeV
## Abstract
Single particle spectra as well as elliptic flow in Cu+Cu collisions at $`\sqrt{s_{NN}}=200`$ GeV are investigated within a hadronic cascade model and an ideal hydrodynamic model. Pseudorapidity distribution and transverse momentum spectra for charged hadrons are surprisingly comparable between these two models. However, a large deviation is predicted for the elliptic flow. The forthcoming experimental data will clarify the transport and thermalization aspects of matter produced in Cu+Cu collisions.
One of the primary current interests in the Relativistic Heavy Ion Collider (RHIC) experiments is to explore the properties of QCD matter far from stable nuclei, especially the confirmation of the deconfined and thermalized matter, i.e. the quark gluon plasma (QGP), which has been predicted from the lattice QCD calculations lattice . While high and medium $`p_T`$ observables such as the parton energy loss e-loss and coalescence behavior of hadron elliptic flows Fries:2004ej are generally believed to give strong evidences of high dense matter formation, hadrons at these momenta are not necessarily formed from thermalized matter. Therefore, low $`p_T`$ observables are also important to confirm whether equilibrium is achieved or not.
Elliptic flow Ollitrault:1992bk is one of the promising observables to study the degree of thermalization for QCD matter produced in heavy ion collisions since it is believed to be sensitive to the properties of the matter at initial stages and the collision geometry Heiselberg:1998es ; Zhang:1999rs . Indeed, incident energy as well as impact parameter dependences of elliptic flow have been investigated extensively. Elliptic flow, i.e. the momentum anisotropy with respect to the reaction plane $`v_2=\mathrm{cos}(2\varphi )`$, has been measured in a wide energy range from GSI-SIS ($`E_{\mathrm{inc}}1A`$ GeV) sis , BNL-AGS ($`E_{\mathrm{inc}}=211A`$ GeV) ags , to CERN-SPS ($`E_{\mathrm{inc}}=40158A`$ GeV) NA49 , in addition to BNL-RHIC RHICflow . Measured collective flows are well reproduced by nuclear transport models assuming the momentum dependent nuclear mean-field at SIS to AGS ($`E_{\mathrm{inc}}0.211A`$GeV) da02 ; Giessen and SPS ($`E_{\mathrm{inc}}=40,158A`$GeV) isse05 energies, whereas elliptic flow at RHIC at mid-rapidity is underestimated in nonequilibrium transport models which do not include explicit partonic interactions Bleicher2000 ; Cassing03 ; Cassing04 . It is also reported that hadronic models explain elliptic flow only at low transverse momentum $`p_T1`$ GeV/$`c`$ at RHIC Burau:2004ev . Partonic interactions followed by quark coalescence hadronization mechanism are proposed in Ref. lin02 to account for the experimental data on elliptic flow. Note, however, that hadronic cascade models reproduce elliptic flow in forward/backward rapidity regions at RHIC Stoecker:2004qu .
On the other hand, in Au+Au collisions at RHIC energies the magnitude of $`v_2`$ and its transverse momentum $`p_T`$ and mass $`m`$ dependences are close to predictions based on ideal and non-dissipative hydrodynamics simulations around midrapidity ($`\eta \stackrel{<}{}\mathrm{\hspace{0.25em}1}`$), in the low transverse momentum region ($`p_T\stackrel{<}{}\mathrm{\hspace{0.25em}1}`$ GeV/$`c`$), and up to semicentral collisions ($`b\stackrel{<}{}\mathrm{\hspace{0.25em}5}`$ fm) Kolb:2000fh ; Hirano:2001eu . This is one of the main results which leads to a recent announcement of the discovery of perfect fluidity at RHIC BNL . (See Ref. HiranoGyulassy for recent reinterpretation of the RHIC data based on current hydrodynamic results.) Despite the apparent success near midrapidity at RHIC, ideal hydrodynamics overestimates the data at lower incident energies (SIS, AGS and SPS) as well as in forward/backward rapidity regions at RHIC probably due to the lack of dissipative effects.
We study Cu+Cu collisions at RHIC in the present work, which is a complementary study of elliptic flow in Au+Au collisions. The particle density and the size of the system are smaller in Cu+Cu collisions than in Au+Au collisions. So the reasonable agreement of hydrodynamic results with Au+Au data may be spoiled in Cu+Cu collisions and a non-equilibrium hadronic description can be relatively important even at RHIC energies. Therefore we employ both a hadronic transport model JAM and a hydrodynamic model to make predictions for elliptic flow in Cu+Cu collisions. Below we briefly summarize hadron-string cascade JAM jam and a hydrodynamic model Hirano:2004rs adopted in this paper.
A hadronic transport model JAM simulates nuclear collisions by the individual hadron-hadron collisions. Soft hadron productions in hadron-hadron scattering are modeled by the resonance and color string excitations. Hard partonic scattering is also included in line with HIJING hijing . Color strings decay into hadrons after their formation time $`(\tau 1`$ fm/$`c`$) according to the Lund string model PYTHIA Sjostrand . Hadrons within their formation time can scatter with other hadrons assuming the additive quark cross section. This simulates constituent quark collisions effectively which is known to be important at SPS energies urqmd . Therefore, matter initially created in collisions is represented by the many strings at RHIC, which means that there is no QGP in the model.
Default parameters in JAM are adopted in this work except for a little wider $`p_T`$ width in the string decay and a larger partonic minimum $`p_T`$ ($`p_0=2.7\text{GeV}/c`$) to fit charged hadron $`p_T`$ spectrum in $`pp`$ collisions at $`\sqrt{s_{NN}}=200`$ GeV. In addition to hadron-hadron collisions, nuclear mean field is incorporated in JAM and its effects are known to be important at AGS and SPS energies isse05 , but mean field is not expected to play major roles at RHIC. We have thus neglected nuclear mean field in this work. The detailed description of JAM can be found in Ref. jam .
Two of the authors (T.H. and Y.N.) have already developed another dynamical framework to describe three important aspects of relativistic heavy ion collisions Hirano:2004rs , namely color glass condensate (CGC) for collisions of two nuclei MV ; Iancu:2003xm ; KLN , hydrodynamics for space-time evolution of thermalized matter hydroreview , and jet quenching for high $`p_T`$ non-thermalized partons GLV . Along the line of these works, we use the same model in this study. However, our aim is to study the bulk properties of matter produced in Cu+Cu collisions. In this paper, we neither include jet components in this model nor discuss jet quenching, unlike a series of the previous work Hirano:2002sc . So hydrodynamic results to be presented below include purely boosted thermal components without any semi-hard components.
In Ref. Hirano:2004rs , a systematic hydrodynamic analysis in Au+Au collisions at $`\sqrt{s_{NN}}=200`$ GeV was performed by using initial conditions taken from the CGC picture for the colliding nuclei. In the conventional hydrodynamic calculations, one chooses initial condition for hydrodynamic equations and thermal freezeout temperature $`T^{\mathrm{th}}`$ so as to reproduce the observed particle spectra, such as (pseudo)rapidity distribution and transverse momentum distribution. So it is believed that hydrodynamics has a less predictive power compared with cascade models. However, if the initial particle production at high collisional energies is supposed to be universal as described by the CGC, hydrodynamics with CGC initial conditions can predict particle spectra. Here, we employ the IC-$`n`$, i.e. a prescription that the number density produced in a CGC collision is matched to the hydrodynamic initial condition Hirano:2004rs , to obtain the initial distribution of thermodynamic variables at the initial time $`\tau _0`$. Once the initial condition is obtained, one solves hydrodynamic equation $`_\mu T^{\mu \nu }=0`$ in the three-dimensional Bjorken coordinate $`(\tau ,\eta _s,x,y)`$ Hirano:2001eu . Here we neglect a dissipative effect and a finite (but probably tiny) baryon density. Assuming $`N_c=N_f=3`$ massless partonic gas, an ideal gas equation of state (EOS) with a bag constant $`B^{1/4}=247`$ MeV is employed in the QGP phase ($`T>T_c=170`$ MeV). We use a hadronic resonance gas model with all hadrons up to $`\mathrm{\Delta }(1232)`$ mass for later stages ($`T<T_c`$) of collisions. We take into account chemical freezeout separated from thermal freezeout Hirano:2002ds as required to obtain sufficient yields for heavier particles. Specifically, we assume that chemical freezeout temperature $`T^{\mathrm{ch}}=170`$ MeV and kinetic freezeout temperature $`T^{\mathrm{th}}=100`$ MeV. Note that the slope of $`p_T`$ spectra becomes insensitive to $`T^{\mathrm{th}}`$ (while $`v_2`$ becomes sensitive) when chemical freezeout is taken into account Hirano:2002ds . In the calculation of $`v_2(\eta )`$, we also use $`T^{\mathrm{th}}=160`$ MeV for comparison. If the strongly coupled QGP (sQGP) core expands as a perfect fluid and the hadronic corona does as a highly dissipative gas as suggested in Ref. HiranoGyulassy , the resultant $`v_2(\eta )`$ and $`v_2(p_T)`$ are expected to be frozen after hadronization Teaney:2000cw ; Teaney:2002aj due to the strong viscous effect. Moreover, $`T^{\mathrm{th}}`$ should be higher for a smaller size of the system Hung:1997du as observed in the centrality dependence of the $`p_T`$ spectra Adams:2003xp . So the freezeout picture in Cu+Cu collisions can be different from that in Au+Au collisions. In the following predictions for hydrodynamic elliptic flow, we show the results for $`T^{\mathrm{th}}=`$ 100 and 160 MeV. For further details of the hydrodynamic model used in this work, see Refs. Hirano:2004rs ; Hirano:2002sc ; Hirano:2002ds .
We first compare the bulk single particle spectra between JAM and hydrodynamics. We emphasize again that our hydrodynamic results are insensitive to a choice of $`T^{\mathrm{th}}`$ for transverse and rapidity distributions of charged hadrons. We show results of the pseudorapidity distribution $`dN/d\eta `$ for charged hadrons in Fig. 1 at impact parameters $`b=1,2`$ and $`5`$ fm. It is seen from this figure that the shape and the magnitude of the distributions from JAM are almost similar to those from hydrodynamics.
In Fig. 2, we compare JAM and hydrodynamic results of the $`p_T`$ spectra for charged hadrons at impact parameters of $`b=1,2`$ and 5 fm for $`|\eta |<0.33`$. Accidentally, these results agree well with each other in transverse momentum range of $`p_T<2`$ GeV/$`c`$. Deviation at higher transverse momentum is due to the lack of jet components in the hydrodynamic simulations.
At least within our models, two distinct pictures, i.e. pictures of coherent particle production via CGC combined with sequential sQGP expansion and of transports of secondary hadrons after hadron-hadron collisions summed up by an overlap region of colliding nuclei, are indistinguishable in the bulk single hadron distributions in Cu+Cu collisions. Note that free parameters in the “CGC+hydro” model has been fixed by fitting the charged multiplicity in Au+Au collisions at midrapidity. We also note that parameters in JAM are already fixed to fit the data in $`pp`$ collisions at $`\sqrt{s_{NN}}=200`$ GeV.
In Fig. 3, we show the impact parameter $`b`$ dependence of the elliptic flow $`v_2`$ at mid-rapidity for charged hadrons. In the hydrodynamic calculations, kinetic freezeout temperatures $`T^{\mathrm{th}}=100`$ MeV and 160 MeV are chosen. While single particle spectra from JAM and hydrodynamics look very similar, a clear difference of $`v_2(b)`$ is seen: $`v_2`$ grows almost linearly with $`b`$ in hydrodynamics, which is the same as the case in Au+Au collisions, while we find a peak at around $`b=6`$ fm in JAM and that the magnitude is only around 20% of the hydrodynamic prediction with $`T^{\mathrm{th}}=100`$ MeV. The two distinct pictures within our approach appear differently in the centrality dependence of elliptic flow. Due to the smaller initial energy density in Cu+Cu collisions compared to Au+Au collisions, the spatial anisotropy is still out-of-plane just after the hadronization and $`v_2`$ continues to be generated even in the late non-viscous hadronic stage in the ideal hydrodynamic simulation. The data is expected to be comparable with the result for $`T^{\mathrm{th}}`$ = 160 MeV if the initial energy density is large ($`e_01`$ GeV/fm<sup>3</sup>) and the equilibration time is small ($`\tau _01`$ fm/$`c`$) enough to create the sQGP phase in Cu+Cu collisions. On the contrary, one expects that it takes more time to reach equilibrium ($`\tau _0>1`$ fm/$`c`$) and that the system may not reach the equilibrated sQGP state since the system size and the produced particle number are small compared with those in Au+Au collisions. In that case, the data will be comparable with the result from JAM.
Pseudorapidity dependences of the elliptic flow from JAM and hydrodynamics are compared with each other to understand the longitudinal dynamics in Cu+Cu collisions in Fig. 4. In JAM, we find almost flat behavior of $`v_2(\eta )`$ around midrapidity ($`|\eta |<2`$), where the charged hadron $`\eta `$ distribution also shows flat behavior. In JAM, elliptic flow is slowly generated ($`t10\mathrm{f}\mathrm{m}`$) as the hadrons are formed from strings after some formation times. In the hydrodynamic calculations, we show the results for $`T^{\mathrm{th}}=100`$ and 160 MeV, which could be an upper and a lower limit of the ideal hydrodynamic prediction respectively. $`v_2(\eta )`$ for $`T^{\mathrm{th}}=100`$ MeV becomes a trapezoidal shape, which looks similar to the result in the previous hydrodynamic study in Au+Au collisions Hirano:2001eu ; Hirano:2002ds . $`v_2(\eta )`$ from ideal hydrodynamics for $`T^{\mathrm{th}}=160`$ MeV is also shown as a possible result for the situation HiranoGyulassy in which $`v_2`$ is generated by the perfect fluid of the sQGP core and is not generated significantly in the dissipative hadronic corona like the result from JAM. Indeed, $`v_2(\eta )`$ for $`T^{\mathrm{th}}=160`$ MeV appears to be a triangle shape which looks similar to the shape in Au+Au data observed by PHOBOS phobos02 .
In Fig. 5, we compare transverse momentum $`p_T`$ dependence of elliptic flow for charged hadrons. Hydrodynamic predictions are of course larger than the ones of JAM. In JAM, $`v_2`$ starts to be saturated at around 0.8 GeV, and the behavior is qualitatively similar to that in Au+Au collisions and another theoretical prediction in Cu+Cu collisions chen05 . It should be noted that we will also find a mass dependent saturating behavior of $`v_2(p_T)`$ when semihard components are combined with the hydrodynamic components Hirano:2004rs .
In summary, we have investigated low-$`p_T`$ observables in a hadron-string cascade model JAM jam and a hydrodynamical model Hirano:2004rs in Cu+Cu collisions at $`\sqrt{s_{NN}}=200`$ GeV. For $`dN/d\eta `$ and $`p_T`$-spectra for charged hadrons, we have obtained good agreement between JAM and hydrodynamics. However, clear deviations between model predictions are found in $`v_2`$ as a function of centrality, pseudorapidity $`\eta `$, and transverse momentum $`p_T`$ for charged hadrons. The lack of elliptic flow in hadronic transport models compared to the ideal hydrodynamic predictions is due to the initial particle production being performed by string decays which only generate a limited amount of transverse momentum for the produced particles in conjunction with the formation time for these hadrons. In that sense, the “EOS” of hadron-string cascade models in the very early stage is to be considered as a super-soft one and cannot generate sufficient pressure needed for elliptic flow to develop. In addition, even if full thermalization is achieved in the hadron cascade model, a higher viscosity in the hadron cascade model would yield a lower elliptic flow than the ideal non-viscous hydrodynamics in the hadron phase. Therefore, in order to interpret correctly the result, $`v_{2,\mathrm{cascade}}<v_{2,\mathrm{hydro}}`$, which has been seen also in SPS energies, we should study the dissipative effects carefully. Measurements of pseudorapidity and transverse momentum dependence of elliptic flow ($`v_2(b)`$, $`v_2(\eta )`$, and $`v_2(p_T)`$) in Cu+Cu collisions at RHIC will provide very important information for transport aspects of QCD matter in heavy ion collisions.
One of the authors (Y.N.) acknowledges discussions with M. Bleicher. The work was supported in part by the US-DOE under Grant No. DE-FG02-93ER40764 (T.H.), and by the Grant-in-Aid for Scientific Research No. 1554024 (A.O.) from the Ministry of Education, Science and Culture, Japan. |
warning/0506/math-ph0506016.html | ar5iv | text | # Rotation Numbers, Boundary Forces and Gap Labelling
## 1 Introduction
It is an interesting and well known observation that the boundary of a domain plays a prominent role both in mathematics and in physics. A case that comes immediately into mind is the theory of differential equations where the boundary conditions determine quite a lot of the whole solution. In a purely topological context the boundary may even determine the behaviour of the system in the bulk completely. A case like this was studied in \[KS04a, KS04b\] where a correspondance between bulk and boundary topological invariants for certain physical systems arising in solid state physics was found. This was mathematically based on $`K`$-theoretic and cyclic cohomological properties of the Wiener-Hopf extension of the $`C^{}`$-algebra of observables. In most applications we have in mind, this $`C^{}`$-algebra is obtained by considering the Schrödinger operator and its translates describing the $`1`$-particle approximation of the solid.
In this article we consider a simple example, a Schrödinger operator on the real line, where such a correspondance can be established more directly with the help of the Sturm-Liouville theorem. The $`K_0`$-theory gap labels (below referred to also as even $`K`$-gap labels) introduced by Bellissard et al. \[BLT85, Be92\] are bulk invariants. It is known that these are equal to the Johnson-Moser rotation numbers \[JM82\] the existing proof being essentially a corollary of the Sturm-Liouville theorem by which they are identified with the integrated density of states on the gaps. In the first part of the paper (Sections 2,3) we provide a direct identification of the Johnson-Moser rotation number (for energies in gaps) with a boundary invariant, here called the Dirichlet rotation number. This boundary invariant has a physical interpretation, namely as boundary force per unit energy. Moreover, it can be interpreted as a $`K_1`$-theory gap label (or odd $`K`$-gap label).
In the second part (Sections 4,5) we indicate how the equality between the $`K_0`$ and the $`K_1`$-theory gap labels also follows from the above-mentioned noncommutative topology of the Wiener Hopf extension. The advantage of this approach is that, unlike the definition of the geometrical rotation numbers and the Sturm-Liouville theorem, it is not restricted in dimension. We tend to think of the $`K_1`$-theory gap label, which is naturally defined in any dimension, as the operator algebraic formulation of the Johnson-Moser rotation number.
Whereas the first part is based on a single operator, although its translates play a fundamental role, we consider in the second part covariant families of operators indexed by the hull of the potential. This is the right framework for the use of ergodic theorems and noncommutative topology. The last section is mainly based on \[Kel\] and therefore held briefly.
## 2 Preliminaries
In this article we consider as in \[Jo86\] a one-dimensional Schrödinger operator $`H=^2+V`$ with (real) bounded potential which we assume (stricter as in \[Jo86\]) to be bounded differentiable. We also consider its translates $`H_\xi :=^2+V_\xi `$, $`V_\xi (x)=V(x+\xi )`$, and lateron its hull. The differential equation $`H\mathrm{\Psi }=E\mathrm{\Psi }`$ for complex valued functions $`\mathrm{\Psi }`$ over $``$ has for all $`E`$ two linear independent solutions but not all $`E`$ belong to the spectrum $`\sigma (H)`$ of $`H`$ as an operator acting on $`L^2()`$. In this situation the following property of solutions holds \[CL55\].
###### Theorem 1
If $`E\sigma (H)`$ there exist two real solutions $`\mathrm{\Psi }_+`$ and $`\mathrm{\Psi }_{}`$ of $`(HE)\mathrm{\Psi }=0`$, $`\mathrm{\Psi }_+`$ vanishing at $`\mathrm{}`$ and $`\mathrm{\Psi }_{}`$ vanishing at $`\mathrm{}`$. These solutions are linear independent and unique up to multiplication by a factor.
We mention as an aside that Johnson proves even exponential dichotomy for such energies \[Jo86\]. Clearly $`\sigma (H_\xi )=\sigma (H)`$ for all $`\xi `$.
We consider also the action of $`H_\xi `$ on $`L^2(^0)`$ with Dirichlet boundary conditions at the boundary. If we need to emphazise this we will also write $`\widehat{H}_\xi `$ for the half-sided operator. The spectrum is then no longer the same. Whereas the essential part of the spectrum of $`\widehat{H}_\xi `$ is contained in that of $`H_\xi `$ \[Jo86\] the half sided operator may have isolated eigenvalues in the gaps in $`\sigma (H_\xi )`$. Here a gap is a connected component of the complement of the spectrum, hence in particular an open set. $`E`$ is an eigenvalue of $`\widehat{H}_\xi `$ if $`(\widehat{H}_\xi E)\mathrm{\Psi }=\mathrm{\Psi }`$ for $`\mathrm{\Psi }L^2(^0)`$ which for $`E`$ in a gap of $`\sigma (H_\xi )`$ amounts to saying that the solution $`\mathrm{\Psi }_{}`$ of $`(H_\xi E)\mathrm{\Psi }_{}=0`$ from Theorem 1 satisfies in addition $`\mathrm{\Psi }_{}(0)=0`$.
###### Definition 1
We call $`E`$ a right Dirichlet value of $`H_\xi `$ if it is an eigenvalue of $`\widehat{H}_\xi `$.
We recall the important Sturm-Liouville theorem:
###### Theorem 2
Consider $`H:=^2+V`$ with (real) bounded continuous potential acting on $`L^2([a,b])`$ with Dirichlet boundary conditions. The spectrum is discrete and bounded from below. A real eigenfunction to the $`n`$th eigenvalue (counted from below) has exactly $`n1`$ zeroes in the interior $`(a,b)`$ of $`[a,b]`$.
## 3 Rotation numbers
The winding number of a continuous function $`f://`$ is intuitively speaking the number of times its graph wraps around the circle $`/`$. This is counted relative to the orientations induced by the order on $``$. Let $`\mathrm{\Lambda }=\{\mathrm{\Lambda }_n\}_n`$ be an increasing chain of compact intervals $`\mathrm{\Lambda }_n=[a_n,b_n]\mathrm{\Lambda }_{n+1}`$ whose union covers $``$. The quantity
$$\mathrm{\Lambda }(f):=\underset{n\mathrm{}}{lim}\frac{1}{b_na_n}_{a_n}^{b_n}f(x)𝑑x$$
is called the $`\mathrm{\Lambda }`$-mean of the function $`f:`$, existence of the limit assumed. Now let $`f:/`$ be continuous and choose a continuous extension $`\stackrel{~}{f}:`$. To define the rotation number of $`f`$ we consider the expression
$$\mathrm{rot}_\mathrm{\Lambda }(f)=\underset{n\mathrm{}}{lim}\frac{\stackrel{~}{f}(b_n)\stackrel{~}{f}(a_n)}{b_na_n}$$
which becomes the winding number of $`f`$ if $`f`$ is periodic of period $`1`$. The limit does not exist in general but if it does it is independent of the extension $`\stackrel{~}{f}`$. If $`f`$ is piecewise differentiable then $`\mathrm{rot}_\mathrm{\Lambda }(f)=\mathrm{\Lambda }(f^{})`$. Moreover, if $`U:`$ is a nowhere vanishing continuous piecewise differentiable function then we can consider the rotation number of its argument function which becomes
$$\mathrm{rot}_\mathrm{\Lambda }(\frac{\mathrm{arg}(U)}{2\pi })=\underset{n\mathrm{}}{lim}\frac{1}{2\pi i(b_na_n)}_{a_n}^{b_n}\frac{\overline{U}}{|U|}\left(\frac{U}{|U|}\right)^{}𝑑x$$
(1)
### 3.1 The Johnson-Moser rotation number
Johnson and Moser in \[JM82\] have defined rotation numbers for the Schrödinger operator $`H=^2+V`$ on the real line where $`V`$ is a real almost periodic potential. They are defined as follows: Let $`\mathrm{\Psi }(x)`$ be the nonzero real solution of $`(HE)\mathrm{\Psi }=0`$ which vanishes at $`\mathrm{}`$, then $`\mathrm{\Psi }^{}+i\mathrm{\Psi }:`$ is nowhere vanishing and
$$\alpha _\mathrm{\Lambda }(H,E):=2\mathrm{rot}_\mathrm{\Lambda }(\frac{\mathrm{arg}(\mathrm{\Psi }^{}+i\mathrm{\Psi })}{2\pi }).$$
(2)
(Our normalisation differs from that in \[JM82\] for later convenience.) For the class of potentials considered here the limit is indeed defined and even independent on the choice of $`\mathrm{\Lambda }`$, we will come back to that in Section 4.
Note that $`\alpha _\mathrm{\Lambda }(H,E)`$ has the following interpretations. If $`N(a,b;E)`$ denotes the number of zeroes of the above solution $`\mathrm{\Psi }`$ in $`[a,b]`$ then $`\alpha _\mathrm{\Lambda }(H,E)`$ is the $`\mathrm{\Lambda }`$-mean of the density of zeroes of $`\mathrm{\Psi }`$, namely one has
$$\alpha _\mathrm{\Lambda }(H,E)=\underset{n\mathrm{}}{lim}\frac{N(a_n,b_n;E)}{b_na_n}.$$
The integrated density of states of $`H`$ at $`E`$ is
$$\mathrm{IDS}_\mathrm{\Lambda }(H,E)=\underset{n\mathrm{}}{lim}\frac{1}{|\mathrm{\Lambda }_n|}\text{Tr}(P_E(H_{\mathrm{\Lambda }_n}))$$
(3)
provided the limit exists. Here $`|\mathrm{\Lambda }_n|=b_na_n`$ is the volume of $`\mathrm{\Lambda }_n`$, $`H_{\mathrm{\Lambda }_n}`$ the restriction of $`H`$ to $`\mathrm{\Lambda }_n`$ with Dirichlet boundary conditions and, for self-adjoint $`A`$, $`P_E(A)`$ is the spectral projection onto the spectral subspace of spectral values smaller or equal to $`E`$. It will be important that $`P(A)`$ is a continuous function of $`A`$ if $`E`$ is not in the spectrum of $`A`$. Since $`\text{Tr}(P_E(H_{\mathrm{\Lambda }_n}))`$ is the number of eigenfunctions of $`H_{\mathrm{\Lambda }_n}`$ to eigenvalue smaller or equal $`E`$ Theorem 2 implies
###### Corollary 1
$`\alpha _\mathrm{\Lambda }(H,E)=\mathrm{IDS}_\mathrm{\Lambda }(H,E)`$.
In particular, like the integrated density of states $`\alpha _\mathrm{\Lambda }(H,E)`$ is monotonically increasing in $`E`$ and constant on the gaps of the spectrum of $`H`$. It is moreover the same for all $`H_\xi `$.
### 3.2 The Dirichlet rotation number
We now consider the continuous 1-parameter family of operators $`\{H_\xi \}_\xi `$ with $`\xi `$ and $`H_\xi =^2+V_\xi `$, where $`V_\xi (x)=V(x+\xi )`$. We shall prove that the Johnson-Moser rotation number is a rotation number which is defined by right Dirichlet values as a function of $`\xi `$.
We choose a gap $`\mathrm{\Delta }`$ in $`\sigma (H_\xi )=\sigma (H)`$ for this section and define the set of right Dirichlet values in $`\mathrm{\Delta }`$
$$D_\xi (\mathrm{\Delta }):=\{\mu \mathrm{\Delta }|\mathrm{\Psi }:(H_\xi \mu )\mathrm{\Psi }=0\text{ and }\mathrm{\Psi }(0)=\mathrm{\Psi }(\mathrm{})=0\}.$$
Thus with respect to this choice of gap we can define
$$S(\mu ):=\{\eta |\mu D_\eta (\mathrm{\Delta })\}.$$
Suppose $`\mu D_\xi (\mathrm{\Delta })`$ for some $`\xi `$ (in particular, $`D_\xi (\mathrm{\Delta })\mathrm{}`$). Then there exists a non-zero solution $`(H_\xi \mu )\mathrm{\Psi }=0`$ satisfying $`\mathrm{\Psi }(0)=\mathrm{\Psi }(\mathrm{})=0`$. Let
$$Z(\mu ,\xi ):=\{x|\mathrm{\Psi }(x\xi )=0\}.$$
This set depends actually only on $`\mu `$, since $`\mathrm{\Psi }`$ is unique up to a multiplicative factor and we have:
###### Lemma 1
Let $`\xi `$ such that $`D_\xi (\mathrm{\Delta })\mathrm{}`$ and $`\mu D_\xi (\mathrm{\Delta })`$. Then $`S(\mu )=Z(\mu ,\xi )`$.
Proof: Let $`\mathrm{\Psi }`$ be a non-zero solution $`(H_\xi \mu )\mathrm{\Psi }=0`$ satisfying $`\mathrm{\Psi }(0)=\mathrm{\Psi }(\mathrm{})=0`$ and define $`\mathrm{\Psi }_\eta (x)=\mathrm{\Psi }(x+(\eta \xi ))`$. Then $`(H_\eta \mu )\mathrm{\Psi }_\eta =0`$ and $`\mathrm{\Psi }_\eta (\mathrm{})=0`$ for all $`\eta `$. Hence $`Z(\mu ,\xi )=\{\eta |\mathrm{\Psi }(\eta \xi )=0\}=\{\eta |\mathrm{\Psi }_\eta (0)=0\}S(\mu )`$.
For the opposite inclusion if $`\mu D_\eta (\mathrm{\Delta })`$, then there exists $`\mathrm{\Phi }`$ such that $`(H_\eta \mu )\mathrm{\Phi }=0`$ with $`\mathrm{\Phi }(0)=\mathrm{\Phi }(\mathrm{})=0`$. Define $`\mathrm{\Phi }_\xi (x+(\eta \xi ))=\mathrm{\Phi }(x)`$. Then $`(H_\xi \mu )\mathrm{\Phi }_\xi =0`$ with $`\mathrm{\Phi }_\xi (\mathrm{})=0`$. By Theorem 1, $`\mathrm{\Psi }=\lambda \mathrm{\Phi }_\xi `$ for some $`\lambda ^{}`$, which implies $`\mathrm{\Psi }(\eta \xi )=\lambda \mathrm{\Phi }(0)=0`$ and hence $`\eta Z(\mu ,\xi )`$, thus $`S(\mu )Z(\mu ,\xi )`$. $`\mathrm{}`$
Let $`\xi S(\mu )`$, $`\mu \mathrm{\Delta }`$. Since the spectrum of $`\widehat{H}_\xi `$ in the gap $`\mathrm{\Delta }`$ consists of isolated eigenvalues which are non-degenerate by Theorem 1 we can use perturbation theory to find a neighbourhood $`(\xi ϵ,\xi +ϵ)`$ and a differentiable function $`\xi \mu (\xi )`$ on this neighbourhood which is uniquely defined by the property that $`\mu (\xi )D_\xi (\mathrm{\Delta })`$. In fact, level-crossing of right Dirichlet values cannot occur in gaps, since it would lead to degeneracies. As in \[Ke04\] we see that its first derivative is strictly negative:
$$\frac{d\mu (\xi )}{d\xi }=_{\mathrm{}}^0𝑑x|\mathrm{\Psi }_\xi (x)|^2V_\xi ^{}=|\mathrm{\Psi }_\xi ^{}(0)|^2<0.$$
Here $`\mathrm{\Psi }_\xi `$ is a normalised eigenfunction of $`\widehat{H}_\xi `$. Thus around each value $`\xi `$ for which we find a right Dirichlet value in $`\mathrm{\Delta }`$ we have locally defined curves $`\mu (\xi )`$ which are strictly monotonically decreasing and non-intersecting. Since $`\widehat{H}_\xi `$ is norm-continuous in $`\xi `$ in the generalised sense, its spectrum $`\sigma (\widehat{H}_\xi )`$ is lower semi-continuous \[K\] in $`\xi `$ so that the curves $`\mu (\xi )`$ can be continued until they reach the boundary of $`\mathrm{\Delta }`$ or their limit at $`+\mathrm{}`$ or $`\mathrm{}`$, if it exists.
Let $`K`$ be the circle of complex numbers of modulus $`1`$. We define the function $`\stackrel{~}{\mu }:K`$ by
$$\stackrel{~}{\mu }(\xi )=\mathrm{exp}\left(2\pi i\underset{\mu D_\xi }{}\frac{\mu E_0}{|\mathrm{\Delta }|}\right)$$
where $`E_0=inf\mathrm{\Delta }`$ and $`|\mathrm{\Delta }|`$ is the width of $`\mathrm{\Delta }`$. Then $`\stackrel{~}{\mu }`$ is a continuous function which is differentiable at all points where none of the curves $`\mu (\xi )`$ touches the boundary.
###### Definition 2
The Dirichlet rotation number is
$$\beta _\mathrm{\Lambda }(H,\mathrm{\Delta }):=\mathrm{rot}_\mathrm{\Lambda }(\frac{\mathrm{arg}\stackrel{~}{\mu }}{2\pi }).$$
###### Lemma 2
If, for some $`\mu \mathrm{\Delta }`$, $`|S(\mu )|>1`$ then $`\mathrm{\Delta }`$ contains at most one right Dirichlet value of $`H_\xi `$.
Proof: We first remark that the same discussion can be performed for the left Dirichlet values of $`H_\xi `$, namely values $`E`$ for which exist $`\mathrm{\Psi }`$ solving $`(H_\xi E)\mathrm{\Psi }=0`$ with $`\mathrm{\Psi }(0)=\mathrm{\Psi }(+\mathrm{})=0`$. These similarily define locally curves $`\mu ^{}(\xi )`$ whose first derivative are now strictly positive. They can’t intersect with any of the curves $`\mu (\xi )`$, because a right Dirichlet value which is at the same time a left Dirichlet value must be a true eigenvalue of $`H`$. Let $`S^{}(\mu )`$ and $`Z^{}(\mu )`$ be defined as $`S(\mu )`$ and $`Z(\mu )`$ but for left Dirichlet values. We claim that between two points of $`S(\mu )`$ lies one point of $`S^{}(\mu )`$. This then implies the lemma, because if $`D_\xi `$ contained two points an elementary geometric argument shows that the curves defined by right Dirichlet values through these points necessarily have to intersect a curve defined by left Dirichlet values. To prove our claim we consider the analogous statement for $`Z(\mu )`$ and $`Z^{}(\mu )`$ and let $`\mathrm{\Psi }_\pm `$ be a real solution of $`(H_0\mu )\mathrm{\Psi }=0`$ with $`\mathrm{\Psi }_\pm (\pm \mathrm{})=0`$. Since $`\mu `$ is not an eigenvalue the Wronskian $`[\mathrm{\Psi }_+,\mathrm{\Psi }_{}]`$ which is always constant does not vanish. Furthermore, if $`\mathrm{\Psi }_+(x)=0`$ then $`\mathrm{\Psi }_{}(x)=[\mathrm{\Psi }_+,\mathrm{\Psi }_{}]/\mathrm{\Psi }_+^{}(x)`$. This expression changes sign between two consecutive zeroes of $`\mathrm{\Psi }_+`$ and hence $`\mathrm{\Psi }_{}`$ must have a zero in between. $`\mathrm{}`$
###### Remark 1
Under the hypothesis of the lemma the sum in the definition of $`\stackrel{~}{\mu }`$ contains at most one element. We believe that the result of the lemma is true under all circumstances.
###### Theorem 3
$`\alpha _\mathrm{\Lambda }(H,E)=\beta _\mathrm{\Lambda }(H,\mathrm{\Delta })`$.
Proof: By Lemma 1 $`\alpha _\mathrm{\Lambda }(H,\mu )`$ is the $`\mathrm{\Lambda }`$-mean of the density of $`S(\mu )`$. Suppose the hypothesis of the Lemma 2 holds. Then $`S(\mu )`$ can be identified with the set of intersection points between the constant curve $`\xi \mathrm{exp}2\pi i\frac{EE_0}{|\mathrm{\Delta }|}`$ and $`\stackrel{~}{\mu }(\xi )`$. Since $`\mu ^{}(\xi )<0`$ the $`\mathrm{\Lambda }`$-mean of the density of these intersection points is minus the rotation number of $`\frac{\mathrm{arg}\stackrel{~}{\mu }}{2\pi }`$.
Now suppose that $`S(\mu )`$ contains at most one element. Then $`\alpha _\mathrm{\Lambda }(H,\mu )=0`$. On the other hand, there can only be finitely many curves defined by right Dirichlet values. Since they intersect the constant curve $`\xi \mathrm{exp}2\pi i\frac{\mu E_0}{|\mathrm{\Delta }|}`$ only once, $`\beta _\mathrm{\Lambda }(H,\mathrm{\Delta })`$ must be $`0`$. $`\mathrm{}`$
Remark 1 An even nicer geometric picture arrises if we take into account also the left Dirichlet values of $`H_\xi `$ for the definition of $`\stackrel{~}{\mu }`$. For this purpose redefine $`\stackrel{~}{\mu }:K`$ by
$$\stackrel{~}{\mu }(\xi )=\mathrm{exp}\pi i\left(\underset{\mu D_\xi }{}\frac{\mu E_0}{|\mathrm{\Delta }|}\underset{\mu D_\xi ^{}}{}\frac{\mu E_0}{|\mathrm{\Delta }|}\right)$$
where $`D_\xi (\mathrm{\Delta })^{}`$ is the set of left Dirichlet values of $`H_\xi `$ in $`\mathrm{\Delta }`$. Then $`\stackrel{~}{\mu }`$ is as well a continuous piecewise differentiable function and $`\mathrm{rot}_\mathrm{\Lambda }(\frac{\mathrm{arg}\stackrel{~}{\mu }}{2\pi })`$ is the same number as before except that it yields the $`\mathrm{\Lambda }`$-mean of the winding per length of the Dirichlet values around a circle which is obtained from two copies of $`\mathrm{\Delta }`$ by identification of their boundary points. For periodic systems, this circle can be identified with the homology cycle corresponding to a gap in the complex spectral curve of $`H`$ \[BBEIM\] and so $`\beta _\mathrm{\Lambda }(H,\mathrm{\Delta })`$ is the winding number of the Dirichlet values around it. This is similar to Hatsugai’s interpretation of the edge Hall conductivity as a winding number (see \[Ha93\]). There the role of the parameter $`\xi `$ is played by the magnetic flux.
### 3.3 Odd $`K`$-gap labels and Dirichlet rotation numbers
We define another type of gap label which is formulated using operator traces and derivations instead of curves on topological spaces. It has its origin in an odd pairing between $`K`$-theory and cyclic cohomology.
We fix a gap $`\mathrm{\Delta }`$ in the spectrum of $`H`$ of length $`|\mathrm{\Delta }|`$ and set $`E_0=inf(\mathrm{\Delta })`$. Let $`P_\mathrm{\Delta }=P_\mathrm{\Delta }(\widehat{H}_\xi )`$ be the spectral projection of $`\widehat{H}_\xi `$ onto the energy interval $`\mathrm{\Delta }`$. Then
$$𝒰_\xi :=P_\mathrm{\Delta }e^{2i\pi \frac{\widehat{H}_\xi E_0}{|\mathrm{\Delta }|}}+1P_\mathrm{\Delta }$$
(4)
acts essentially as the unitary of time evolution by time $`\frac{1}{|\mathrm{\Delta }|}`$ on the eigenfunctions of $`\widehat{H}_\xi `$ in $`\mathrm{\Delta }`$. These eigenfunctions are all localised near the edge and therefore is the following expression a boundary quantity.
###### Definition 3
The odd $`K`$-gap label is
$$\mathrm{\Pi }_\mathrm{\Lambda }(H,\mathrm{\Delta })=\underset{n\mathrm{}}{lim}\frac{1}{2i\pi |b_na_n|}_{a_n}^{b_n}\text{Tr}[(𝒰_\xi ^{}1)_\xi 𝒰_\xi ]𝑑\xi $$
Where Tr is the standard operator trace on $`L^2()`$.
###### Theorem 4
$`\mathrm{\Pi }_\mathrm{\Lambda }(H,\mathrm{\Delta })=\beta _\mathrm{\Lambda }(H,\mathrm{\Delta })`$.
Proof: Note that the rank of $`P_\mathrm{\Delta }`$ is equal to $`|D_\xi (\mathrm{\Delta })|`$, the number of elements in $`D_\xi (\mathrm{\Delta })`$. Let us first suppose that this is either $`1`$ or $`0`$ which would be implied under the conditions of Lemma 2. Since $`𝒰_\xi ^{}1=P_\mathrm{\Delta }(e^{2i\pi \frac{\widehat{H}_\xi E_0}{|\mathrm{\Delta }|}}1)`$ we can express the trace using the normalised eigenfunctions $`\mathrm{\Psi }_\xi `$ of $`\widehat{H}_\xi `$ to $`\mu (\xi )`$, provided $`|D_\xi (\mathrm{\Delta })|=1`$,
$$\text{Tr}[(𝒰_\xi ^{}1)_\xi 𝒰_\xi ]_\xi 𝒰_\xi )|\mathrm{\Psi }_\xi =\mathrm{\Psi }_\xi |𝒰_\xi ^{}1|\mathrm{\Psi }_\xi \mathrm{\Psi }_\xi |_\xi 𝒰_\xi |\mathrm{\Psi }_\xi .$$
(5)
Substituting
$$\mathrm{\Psi }_\xi |_\xi 𝒰_\xi |\mathrm{\Psi }_\xi =_\xi \mathrm{\Psi }_\xi |𝒰_\xi |\mathrm{\Psi }_\xi =_\xi e^{2i\pi \frac{\mu (\xi )E_0}{|\mathrm{\Delta }|}}$$
in the previous expression we arrive at
$$\text{Tr}[(𝒰_\xi ^{}1)_\xi 𝒰_\xi ]=(e^{2i\pi \frac{\mu (\xi )E_0}{|\mathrm{\Delta }|}}1)_\xi e^{2i\pi \frac{\mu (\xi )E_0}{|\mathrm{\Delta }|}}.$$
Since $`𝒰_\xi ^{}1=0`$ if $`D_\xi (\mathrm{\Delta })=\mathrm{}`$ we have
$$\mathrm{\Pi }_\mathrm{\Lambda }(H,\mathrm{\Delta })=\underset{n\mathrm{}}{lim}\frac{1}{2i\pi |b_na_n|}_{a_n}^{b_n}(\overline{\stackrel{~}{\mu }(\xi )}1)\stackrel{~}{\mu }^{}(\xi )𝑑\xi =\frac{1}{2i\pi }\mathrm{\Lambda }(\overline{\stackrel{~}{\mu }}\stackrel{~}{\mu }^{})$$
(6)
which is the expression for $`\beta _\mathrm{\Lambda }(H,\mathrm{\Delta })`$.
If $`|D_\xi |>1`$ one has to replace the r.h.s. of (5) by a sum over eigenfunctions of $`\widehat{H}_\xi `$ and the calculation will be similar. $`\mathrm{}`$
### 3.4 Interpretation as boundary force per unit energy
We assume for simplicity $`|D_\xi |1`$. Then we obtain from (6)
$$\mathrm{\Pi }_\mathrm{\Lambda }(H,\mathrm{\Delta })=\underset{n\mathrm{}}{lim}\frac{1}{|b_na_n|}_{a_n}^{b_n}\mu ^{}(\xi )\frac{|D_\xi (\mathrm{\Delta })|}{|\mathrm{\Delta }|}𝑑\xi .$$
The r.h.s. is $`\frac{1}{|\mathrm{\Delta }|}`$ times the $`\mathrm{\Lambda }`$-mean of the expectation value of the gradient force w.r.t. the density matrix associated with the egde states in the gap. Since translating $`\widehat{H}_\xi `$ in $`\xi `$ is unitarily equivalent to translating the position of the boundary, $`\mathrm{\Pi }`$ can be seen as the force per unit energy the edge states in the gap of the system exhibit on the boundary \[Kel\].
## 4 Hulls and ergodic theorems
So far we have worked with a single potential and its translates. When completed w.r.t. a natural metric topology this set of translates yields a topological space, called the *hull* of the potential. As it has become apparent in recent years, many topological invariants of the physical system depend mainly on the topology of this hull with its $``$ action by translation of the potential. Besides, the use of invariant ergodic probability measures on the hull allows to tackle the problem of existence of the $`\mathrm{\Lambda }`$-means in a probabilistic sense. It is therefore most natural to interprete the results of the last section in the framework of $``$-actions on hulls. This allows for a generalisation to higher dimensional systems, to which the theorems of Section 2 do not extend.
Given a potential $`V`$ consider its hull
$$\mathrm{\Omega }=\overline{\{V_\xi |\xi \}},$$
which is the compactification of the set of translates of $`V`$ in the sense of \[Jo86, Be92\]. The action of $``$ by translation of the potential extends to an action on $`\mathrm{\Omega }`$ by homeomorphisms which we denote by $`\omega x\omega `$. The elements of $`\mathrm{\Omega }`$ may be identified with those real functions (potentials) which may be obtained as limits of sequences of translates of $`V`$. We shall write $`V_\omega `$ for the potential corresponding to $`\omega \mathrm{\Omega }`$. If $`\omega _0`$ is the point of $`\mathrm{\Omega }`$ corresponding to $`V`$ then $`V_\xi =V_{\xi \omega _0}`$. Also $`V_{y\omega }(x)=V_\omega (xy)`$ and so the family of Hamiltonians $`H_\omega =^2+V_\omega `$ is covariant in the sense that $`H_{x\omega }=U(x)H_\omega U^{}(x)`$ were $`U(x)`$ is the operator of translation by $`x`$. The bulk spectrum is by definition the union of their spectra.
The valididty of the following theorem, namely that $`\mathrm{\Omega }`$ carries an $``$-invariant ergodic probability measure, can be verified for many situations, see \[BHZ00\] for considerations relating it to the Gibbs measure.
###### Theorem 5
Suppose that $`(\mathrm{\Omega },)`$ carries an invariant ergodic probability measure $`𝐏`$. Let $`\mathrm{\Delta }`$ be a gap in the bulk spectrum and $`E\mathrm{\Delta }`$. Then almost surely (w.r.t. this measure) the limits to define $`\alpha _\mathrm{\Lambda }(H_\omega ,E)`$ and $`\mathrm{\Pi }_\mathrm{\Lambda }(H_\omega ,\mathrm{\Delta })`$ exist and are independent of $`\mathrm{\Lambda }`$ and $`\omega \mathrm{\Omega }`$. The almost sure value of $`\mathrm{\Pi }_\mathrm{\Lambda }`$ is the $`𝐏`$-average
$$\mathrm{\Pi }(\mathrm{\Delta })=\frac{1}{2i\pi }_\mathrm{\Omega }𝑑𝐏(\omega )\text{Tr}((𝒰_\omega ^{}1)\delta ^{}𝒰_\omega )$$
where $`(\delta ^{}f)(\omega )=\frac{df(t\omega )}{dt}|_{t=0}`$ and $`𝒰_\omega `$ is defined as in (4) with $`\widehat{H}_\omega `$ in place of $`\widehat{H}_\xi `$.
Proof: The crucial input is Birkhoff’s ergodic theorem which allows to replace
$$\underset{n\mathrm{}}{lim}\frac{1}{|\mathrm{\Lambda }_n|}_{\mathrm{\Lambda }_n}F(x\omega )𝑑x=_\mathrm{\Omega }𝑑𝐏F(\omega )$$
for almost all $`\omega `$ and any $`FL^1(\mathrm{\Omega },𝐏)`$. The corresponding construction for the rotation number $`\alpha `$ has been carried out in \[JM82\] for almost periodic potentials and for the more general set up in \[Jo86, Be92\]. For $`\mathrm{\Pi }_\mathrm{\Lambda }`$ the relevant function is $`F(\omega )=\text{Tr}((𝒰_\omega ^{}1)\delta ^{}𝒰_\omega )`$ which leads to the expression of the almost sure value of $`\mathrm{\Pi }_\mathrm{\Lambda }`$. $`\mathrm{}`$
## 5 $`K`$-theoretic interpretation
The dynamical system $`(\mathrm{\Omega },)`$ does not depend on the details of $`V`$, but only on its spatial structure (or what may be called its long range order). In fact, for systems whose atomic positions are described by Delone sets there are methods to construct the hull directly from this set, c.f. \[BHZ00, FHK02\]. The detailled form of the potential is rather encoded in a continuous function $`v:\mathrm{\Omega }`$ so that $`V_\omega (x)=v(x\omega )`$ is the potential corresponding to $`\omega `$. $`C(\mathrm{\Omega })`$ is thus the algebra of continuous potentials for a given spatial structure.
If one combines this algebra with the Weyl-algebra of rapidly decreasing functions of momentum operators one obtains the algebra of continuous observables which is the $`C^{}`$-crossed product $`C(\mathrm{\Omega })_\phi `$. It is the $`C^{}`$-closure of the convolution algebra of functions $`f:C(\mathrm{\Omega })`$ with product $`f_1f_2(x)=_{}𝑑yf_1(y)\phi _yf_2(xy)`$ and involution $`f^{}(x)=\phi _x\overline{f(x)}`$, where $`\phi _y(f)(\omega )=f(y\omega )`$. It has a faithful family of representations $`\{\pi _\omega \}_{\omega \mathrm{\Omega }}`$ on $`L^2()`$ by integral operators,
$$x|\pi _\omega (f)|y=f(yx)(x\omega ).$$
It has the following important property. For each continuous function $`F:`$ vanishing at $`0`$ and $`\mathrm{}`$ there exists an element $`\stackrel{~}{F}C(\mathrm{\Omega })_\phi `$ such that $`F(H_\omega )=\pi _\omega (\stackrel{~}{F})`$. Some of the topological properties of the family of Schrödinger operators $`\{H_\omega \}_{\omega \mathrm{\Omega }}`$ are therefore captured by the topology of the $`C^{}`$-algebra. The invariant measure $`𝐏`$ over $`\mathrm{\Omega }`$ gives rise to a trace $`𝒯:C(\mathrm{\Omega })_\phi `$, $`𝒯(f)=_\mathrm{\Omega }𝑑𝐏f(0)`$.
###### Theorem 6 (\[Be92\])
Let $`E`$ be in a gap of the bulk spectrum of $`\{H_\omega \}_{\omega \mathrm{\Omega }}`$ so that in particular there exists a projection $`\stackrel{~}{P}_EC(\mathrm{\Omega })_\phi `$ such that $`\pi _\omega (\stackrel{~}{P}_E)=P_E(H_\omega )`$ is the projection onto the spectral subspace of $`H_\omega `$ to energies below the gap. Suppose that the potential which gave rise to the hull $`\mathrm{\Omega }`$ is smooth. Then the almost sure value of $`\mathrm{IDS}_\mathrm{\Lambda }(H,E)`$ is $`\mathrm{IDS}(E):=𝒯(\stackrel{~}{P}_E)`$.
We mention that this result is more subtle than just an application of Birkhoff’s theorem and interpretating the result in $`C^{}`$-algebraic terms as it needs a Shubin type argument which holds for smooth potentials, namely
$$\underset{n\mathrm{}}{lim}\frac{1}{|\mathrm{\Lambda }_n|}(\text{Tr}(P_E(H_{\mathrm{\Lambda }_n})\text{Tr}(\chi _{\mathrm{\Lambda }_n}P_E(H)))=0.$$
The element $`\stackrel{~}{P}_E`$ is a projection. As any trace on a $`C^{}`$-algebra, $`𝒯`$ depends only on the homotopy class of $`\stackrel{~}{P}_E`$ in the set of projections of $`C(\mathrm{\Omega })_\phi `$. The even $`K`$-group $`K_0(C(\mathrm{\Omega }_\phi )`$ is constructed from homotopy classes of projections and the map on projections $`P𝒯(P)`$ induces a functional on this group, or stated differently, the elements of the $`K_0`$-group pair with $`𝒯`$. It is therefore reasonable to refer to $`𝒯(\stackrel{~}{P}_E)`$ as an *even* $`K`$-gap label (or $`K_0`$-theory gap label) of the gap. This is the $`K_0`$-theoretical gap labelling of \[BLT85, Be92\].
There is a similar identification of the odd $`K`$-gap label as the result of a functional applied to the odd $`K`$-group of a $`C^{}`$-algebra. This $`C^{}`$-algebra is the $`C^{}`$-algebra of observables on the half space near $`0`$, the position of the boundary. It turns out to be convenient to consider also the cases in which the boundary is at $`s0`$. We therefore consider the space $`\mathrm{\Omega }\times `$ with the product topology. This topological space, whose second component denotes the position of the boundary, carries an action of $``$ by translation of the potential and the boundary (so that their relative position remains the same). The relevant $`C^{}`$-algebra is then the crossed product (constructed as above) $`C_0(\mathrm{\Omega }\times )_{\stackrel{~}{\phi }}`$ with $`\stackrel{~}{\phi }_y(f)(\omega ,s)=f(y\omega ,s+y)`$. It has a family of representations $`\{\pi _{\omega ,s}\}_{\omega \mathrm{\Omega },s}`$ on $`L^2()`$ by integral operators,
$$x|\pi _{\omega ,s}(f)|y=f(yx)(x\omega ,sx).$$
It has the following important property: for each continuous function $`F:`$ vanishing at $`0`$ and $`\mathrm{}`$ and such that $`F(H_\omega )=0`$ for all $`\omega `$, there exists an element $`\widehat{F}C_0(\mathrm{\Omega }\times )_{\stackrel{~}{\phi }}`$ such that $`F(H_{\omega ,s})=\pi _{\omega ,s}(\widehat{F})`$, where $`H_{\omega ,s}`$ is the restriction of $`H_\omega `$ to $`^s`$ with Dirichlet boundary conditions at $`s`$. Let $`𝒰=\{𝒰_{\omega ,s}\}`$,
$$𝒰_{\omega ,s}:=P_\mathrm{\Delta }e^{2i\pi \frac{H_{\omega ,s}E_0}{|\mathrm{\Delta }|}}+1P_\mathrm{\Delta },$$
(7)
similar to (4). The product measure of $`𝐏`$ with the Lebesgue measure is an $``$-invariant measure on $`\mathrm{\Omega }\times `$ and defines a trace $`\widehat{𝒯}(f)=_\mathrm{\Omega }_{}𝑑𝐏𝑑sf(0).`$
###### Theorem 7 (\[Kel\])
Let $`\mathrm{\Delta }`$ be a gap in the bulk spectrum of $`\{H_\omega \}_{\omega \mathrm{\Omega }}`$. The almost sure value of $`\mathrm{\Pi }(\mathrm{\Delta })`$ is
$$\mathrm{\Pi }_\mathrm{\Lambda }(H,\mathrm{\Delta })=\mathrm{\Pi }(\mathrm{\Delta }):=\frac{1}{2i\pi }\widehat{𝒯}(\widehat{𝒰^{}1}\delta ^{}\widehat{𝒰1}).$$
The expression of the theorem depends only on the homotopy class of $`\widehat{𝒰1}+1`$ in the set of unitaries of (the unitization of) $`C_0(\mathrm{\Omega }\times )_{\stackrel{~}{\phi }}`$. The odd $`K`$-group $`K_1(C_0(\mathrm{\Omega }\times )_{\stackrel{~}{\phi }})`$ is constructed from homotopy classes of unitaries and the map on unitaries $`U\widehat{𝒯}((U^{}1)\delta ^{}U)`$ induces a functional on this group. It is therefore that we refer to $`\frac{1}{2i\pi }\widehat{𝒯}(\widehat{𝒰^{}1}\delta ^{}\widehat{𝒰1})`$ as an odd $`K`$-gap label of the gap.
The proof of the following theorem is based on the topology of the above $`C^{}`$-algebras.
###### Theorem 8 (\[Kel\])
$`𝒯(\stackrel{~}{P}_E)=\frac{1}{2i\pi }\widehat{𝒯}(\widehat{𝒰^{}1}\delta ^{}\widehat{𝒰1})`$. In other words, $`\mathrm{IDS}(E)=\mathrm{\Pi }(\mathrm{\Delta })`$, $`E\mathrm{\Delta }`$.
## 6 Conclusion and final remarks
We have discussed four quantities which serve as gap-labels for one-dimensional Schrödinger operators. They are all equal but their definition relies on different concepts. The Johnson-Moser rotation number $`\alpha `$ measures the mean oscillation of a single solution. The Dirichlet rotation number $`\beta `$ counts the mean winding of the eigenvalues of the halfsided operators around a circle compactification of the gap. $`\mathrm{\Pi }`$ and $`\mathrm{IDS}`$ are operator algebraic expressions with concrete physical interpretations, the boundary force per energy and the integrated density of states. Whereas the identities $`\alpha =\beta =\mathrm{\Pi }`$ are rather elementary, their identity with $`\mathrm{IDS}`$ is based on a fundamental theorem, the Sturm-Liouville theorem. We tend to think therefore of $`\mathrm{\Pi }`$ as the natural operator algebraic formulation of the Johnson-Moser rotation number and of Theorem 8 as an operator analog of the Sturm-Liouville theorem. The advantage is that $`\mathrm{\Pi }`$, $`\mathrm{IDS}`$ and Theorem 8 generalise naturally to higher dimensions \[Kel\]. In fact, the expression for $`\mathrm{IDS}`$ is the same as in (3) if one uses Føllner sequences $`\{\mathrm{\Lambda }_n\}_n`$ for $`^d`$. The expression of $`\mathrm{\Pi }_\mathrm{\Lambda }`$ in $`^d`$ requires a choice of a $`d1`$-dimensional subspace, the boundary, and so $`\widehat{H}_\xi `$ is the restriction of the Schrödinger operator $`H_\xi =\mathrm{\Sigma }_j_j^2+V_\xi `$, $`V_\xi (x)=V(x+\xi e_d)`$, to the half space $`^{d1}\times ^0`$ with Dirichlet boundary conditions. Then
$$\mathrm{\Pi }_\mathrm{\Lambda }=\underset{n\mathrm{}}{lim}\frac{1}{|\mathrm{\Sigma }_n|(b_na_n)}_{a_n}^{b_n}\text{Tr}((𝒰_{\xi ,\mathrm{\Sigma }_n}^{}1)_\xi 𝒰_{\xi ,\mathrm{\Sigma }_n})𝑑\xi ,$$
$$𝒰_{\xi ,\mathrm{\Sigma }_n}=P_\mathrm{\Delta }(\widehat{H}_{\xi ,\mathrm{\Sigma }_n})e^{2\pi i\frac{\widehat{H}_{\xi ,\mathrm{\Sigma }_n}E_0}{|\mathrm{\Delta }|}}+1P_\mathrm{\Delta }(\widehat{H}_{\xi ,\mathrm{\Sigma }_n}).$$
Here $`\mathrm{\Sigma }_n`$ is a Føllner sequence for the boundary and $`\widehat{H}_{\xi ,\mathrm{\Sigma }_n}`$ is the restriction of $`H_\xi `$ to $`\mathrm{\Sigma }_n\times ^0`$ with Dirichlet boundary conditions. We do not know of a direct link between this expression and the generalisation proposed by Johnson \[Jo91\] for odd-dimensional systems.
*Acknowledgements:*
The second author would like to thank EPSRC for financial support (contract number GR/R64995/01) and the University of Lyon I, Institute Girard Desargues, for its hospitality. |
warning/0506/math0506145.html | ar5iv | text | # Continuous and Tractable models for the Variation of Evolutionary Rates
## 1 Introduction
Understanding evolutionary rates and how they vary is one of the central concerns of molecular evolution. It has been clearly shown that inadequate models of rate variation, between lineages and between loci, can dramatically affect the accuracy of phylogenetic inference . The dependency of molecular dating on evolutionary rate models is even more critical: we will only obtain precise divergence time estimates from molecular data once we can model the rate at which sequences evolve .
Modelling the evolutionary rate is made difficult by the number and variety of factors influencing it. The base rate of mutation can vary because of changes in the accuracy of transcription machinery , DNA repair mechanisms , and metabolic rate . At the cellular level, selective pressures can lead to variation of rate between loci and over time, as evidenced by differential rates of the three codon position , the slower evolutionary rate of highly expressed genes , and the effect of tertiary structure on patterns of sequence conservation .
Selection also affects the evolutionary rate at the level of populations. For the most part, the only mutations that affect phylogenetics are those that are fixed in the population. Hence evolutionary rate is a combination of mutation rate and fixation rate. Fluctuations in population size, generation times, and environmental pressures affect fixation rates and thereby influence evolutionary rate .
Because of this complexity, the strategies employed for modelling evolutionary rate have tended to be statistical in nature. As with all statistical inference, there is an iterative sequence of model formulation, model assessment, and model improvement. The aim is to construct a model that accurately explains the observed variation but is as simple, and tractable, as possible.
Our goal in this paper is to derive a continuous model for rate evolution that avoids many of the problems of existing approaches. We base our model on the CIR process, a continuous Markov process that is widely used in finance to model interest rates . As we shall see, the model fits well into existing protocols for phylogenetic inference. The process has a stationary distribution given by a gamma distribution and yet, unlike the rates-across-sites (RAS) model of Uzzell and Corbin , the rate is allowed to vary along lineages. The CIR model adds only one parameter to the RAS model, and this parameter can be estimated directly from the index of dispersion or the autocorrelation (see below). Furthermore, we can derive exact transition probabilities when we incorporate CIR based rate variation into the standard models for sequence evolution.
The outline of the paper is as follows:
* In the following section we summarise the key characteristics of models for rate evolution, and show how existing models are classified with respect to these characteristics.
* In Section 3 we present the CIR model for rate evolution and discuss its basic properties.
* In Section 4 we derive transition probabilities for standard mutation models where the rate is described as a Markov process.
* In Section 5 we focus on the case where the rate is modelled by a CIR process.
* In Section 6 we extend this one step further to derive an expression for likelihood of a three-taxa tree using a mutation model with rate determined by the CIR process. We note that three-taxa trees are often used to study differences in evolutionary rate.
We conclude with an outline of future work and work in progress.
In a companion paper (in preparation) we describe the incorporation of this model into software for Bayesian phylogenetic inference, and use this to show how our model captures important information lost in standard RAS approaches.
## 2 Properties of models for rate variation
In this section we examine several important characteristics that can be used to distinguish, and choose between, different models for rate variation. We discuss how the different existing models fit into this scheme and summarise the differences between them in Table 1.
The rate of evolution for a given locus at time $`t0`$ is denoted by $`R_t`$. For each $`t>0`$, $`R_t`$ is a non-negative random variable, and different models of rate evolution give different distributions for the rates $`R_t`$, $`t0`$.
Here and throughout the paper we will restrict out attention to *Markov processes*. That is, for any $`t_1t_2t_3`$, we assume that $`R(t_3)`$ conditioned on $`R(t_2)`$ is independent of $`R(t_1)`$. In other words, the future depends on the past only through the present.
Property I: Continuous or Discontinuous Sample Paths
The first characteristic is whether sample paths of the process are continuous or discontinuous with respect to time. Typically, models with discontinous paths have rates $`R_t`$ that are constant except for discrete points in time at which there is a jump in the value (Figure 1-1(a)). If the number of possible values for the rate is finite, then the rate can easily be described as a continuous-time Markov chain with a infinitesimal rate matrix. For example, in the covarion process the basic rates are ‘off’ ($`R_t=0`$) or ‘on’ ($`R_t=1`$) and transitions occur between them at exponentially distributed random time intervals. Galtier generalizes this process to one with more than two possible states. In other models, the range of possible values for the rate is continuous, as in the model of Huelsenbeck , where a rate change event consists of multiplying the previous rate by a gamma random variable. The rate change events are still discrete and exponentially distributed.
There also are models that describes the rate as a continuous function with time, and the most important class of Markov processes with continuous paths are *diffusions* (Figure 1-2(a)). Examples include the CIR process presented here, the Ornstein-Uhlenbeck model of Aris-Brosou and Yang , and the log-normal model of Kishino et al. .
Finally, it is also possible for $`R_t`$ to make jumps in value at a discrete set of times while also changing continuously in between these points.
Property II: Long Term Behaviour and Ergodicity
The second property we consider is the distribution of $`R_t`$ as $`t`$ goes to infinity, that is, the distribution of the rate of evolution in the long term. Surprisingly, many models of rate evolution are very badly behaved in the limit.
One problematic class of processes that have already been applied to rates in phylogenetics is the martingales. We say that a Markov process is a *martingale* if, for all $`s,t0`$ we have $`\mathrm{E}[M_{t+s}|M_t]=M_t`$ . An example of of a Markov martingale is Brownian motion. As a result of this fairly innocuous looking condition, a martingale $`M_t`$ has the property that either $`\mathrm{E}[|M_t|]`$ is unbounded in time or $`M_t`$ converges to a random constant . Either possibility is undesirable from a modelling point of view. This may not be a problem if we only look at the process over a finite time, but neither is it particularly desirable. The processes of Kishino et al. and Huelsenbeck et al. all have the property that either $`R_t`$ or $`\mathrm{log}(R_t)`$ is a martingale.
At a purely theoretical level, we observe that an ever-increasing variance will result for almost any signal that is *only* driven by its initial value and a stochastic force, with no directional bias. The position of a particle subjected to a random force produced by collisions with other particles is a classical example of such a case. In our context, the effects on the evolutionary rate are not independent of the actual rate : whatever the theoretical framework we consider, a high evolutionary rate is not as likely to increase (or to stay in high values) as to go back to smaller values. The episodic evolution fits particularly well to this idea, where periods of drastic adaptation with high evolutionary rates are naturally followed by periods where a population is adapted and its genome evolves much more slowly. Even according to the neutral theory, as argued by Takahata , the overall dynamic of the rate should behave as a random function that takes high values whenever bottlenecks occur, and goes back to small values afterwards.
The concept of ergodicity naturally arises from this observation. We say that a Markov process is *ergodic* if for any initial rate $`R_0`$ the distribution of $`R_t`$ converges to a unique distribution as $`t`$ goes to infinity. The limiting distribution is known as the invariant or stationary distribution. Examples of ergodic processes include the OU process, the CIR model and (usually) the discrete space covarion and covarion-type models .
One possible way for a process to not be ergodic is if the distribution of $`R_t`$ does not converge for large $`t`$ for some initial rate $`R_0`$. This must be the case if $`R_t`$ is a martingale and does not converge to a constant, as is the case with Brownian motion. Another possibility is that $`R_t`$ converges to different stationary distributions for different values of $`R_0`$.
Property III: Tractability
A highly desirable feature of any model is its tractability, both mathematical (does there exist a closed formula?) and computational (can we compute probabilities efficiently?). Nowadays, Monte Carlo methods make it possible to use arbitrarily complex models: however, explicit analytical formulae allow for more efficient sampling .
There are several probability distribution functions that are important to have when working with rate processes. The most basic is the distribution of the rate $`R_t`$ given the rate at time $`t=0`$. This we have for the models and for the CIR model, but not for the models of .
In phylogenetics we incorporate the model for evolutionary rate into the mutation model for sequence evolution at a site. These interact to give a joint process $`(R_t,X_t)`$ for both the rate $`R_t`$ at time $`t`$ and the nucleotide or protein $`X_t`$ at time $`t`$. To evaluate the likelihood we require an expression for the joint conditional probability
$`P[X_t=j,R_t=s|X_0=i,R_0=r]`$ (1)
of going from one nucleotide (or amino acid) state and rate state to another pair of states. Even though it is sometimes possible to perform Monte Carlo computations to estimate this probability without a formula (as in ), having a formula will speed up the computations significantly without having to resort to approximations, as in .
Property IV: Autocovariance and dispersion
There is general agreement , , on the relevance of autocorrelation in the modelling of evolutionary rate. Broadly speaking, if the various causes that explain rate variation (generation time, population size, environmental fitness) vary with time, it should be reflected in rate variations. The extent to which the rate varies can be studied using the index of dispersion (Kimura, , Langley and Fitch ). Let $`N(t)`$ be the number of substitutions or mutations of a sequence over time $`t`$. The index of dispersion $`I(t)`$ is defined as
$`I(t)={\displaystyle \frac{\text{Var}[N(t)]}{\mathrm{E}[N(t)]}}.`$ (2)
This statistic can be estimated by comparing the number of substitutions that have accumulated in different lineages . The population genetics community has proposed different models to account for a high index of dispersion (, ), and any reasonable model should yield an index of dispersion of at least one.
The index of dispersion resulting from a particular model of rate variation is a function of the autocovariance of that model. The autocovariance for a process $`R_t`$ is defined by
$`\rho (t)=\text{Cov}(R_0,R_t).`$ (3)
For many processes we can derive an explicit formula for the autocovariance. If we assume that the substitutions occur according to a Poisson process with rate governed by our rate process (that is, the substitutions follow a doubly stochastic or Cox process, Section 4) and the rate process has autocovariance function $`\rho (t)`$ then
$`I(t)=1+{\displaystyle \frac{2_0^t\left(1\frac{s}{t}\right)\rho (s)𝑑s}{\mathrm{E}(R_t)}},`$ (4)
as stated by a theorem in , and the stationary index of dispersion is then
$`I(\mathrm{})=\underset{t\mathrm{}}{lim}I(t)=1+{\displaystyle \frac{2_0^{\mathrm{}}\rho (s)𝑑s}{\mu }},`$ (5)
provided that $`\mu `$, the stationary mean of the process $`R(t)`$, and the limit, exist. Note that if there is any variation in rate then the index of dispersion will be greater than one .
Some rate models in phylogenetics don’t model explicitly the rate, but instead assign a (fixed) rate to each branch, so that the expected number of substitutions on a particular branch is equal to its length times its assigned (constant) rate.
A close look at the log-normal model from Thorne et al. , which differs from their previous version in that the rate is explicitly modelled, we suggest that the rate has constant autocovariance, since this rate process is close to a transform of the Brownian motion, and Brownian motion has a constant autocovariance function. Put into equation (5), we see that the the index of dispersion diverges. This problematic result illustrates the necessity of a balance between the presence of autocorrelation on one side, and the decrease of autocorrelation on a large time scale.
Property V: Heterotachy or Homotachy
There are two general ways that models for evolutionary rate can be incorporated into phylogenetics. On one hand, we have rate variation among lineages that applies to all sites (or loci) together. This can be modelled by trees for which the paths from the root to the leaves have different lengths. The rate variation explains the extent to which the evolution of the sequences has violated the molecular clock. Alternatively, we can introduce a distinct rate process for each site or locus. This models heterotachy, where the lineage rate changes are site-specific . The transition probabilities that we derive in Section 4 can be applied to homotachous as well as heterotachous models.
## 3 A Continuous diffusion model for the evolutionary rate
A Markov process with continuous paths and satisfying some additional smoothness conditions on its transition probabilities is called a *diffusion*. There are many ways of specifying a diffusion process: perhaps the most intuitive one is by giving the probability distribution function (pdf) of $`R_t`$ given $`R_0=r_0`$, for arbitrary $`r_0`$.We denote this pdf by $`f_R[R_t|r_0]`$. For example, Brownian motion with parameter $`\sigma ^2`$ is defined by the condition that $`f_R[R_t|r_0]`$ is a normal density with mean $`r_0`$ and variance $`\sigma ^2t`$.
A mathematically convenient representation of a diffusion is by means of a *stochastic differential equation* (SDE). In the same way that a dynamical system can be defined as the solution of a differential equation, a diffusion process $`R_t`$ can be defined as the solution of an equation taking the general form (see p.61)
$`dR_t=\alpha (t,R_t)dt+\beta (t,R_t)dB_t.`$ (6)
Here, $`\alpha (t,R_t)`$ represents the deterministic effect on $`R_t`$, $`\beta (t,R_t)`$ the stochastic part, and $`dB_t`$ is an infinitesimal “random” increment. Brownian motion corresponds to the case when $`\alpha (t,R_t)=0`$ for all $`t`$, $`\beta (t,R_t)`$ is constant and the SDE becomes
$`dR_t=\sigma dB(t).`$
Note that if $`\beta (t,R_t)=0`$ for all $`t`$ and $`R_t`$ then (6) becomes a deterministic ordinary differential equation.
Going from an SDE such as (6) to a pdf for the diffusion involves solving a variable-coefficient second-order partial differential equation (PDE). For general functions $`\alpha `$ and $`\beta `$ this PDE has no analytic solution. There are very few diffusions known that have closed form equations for their pdfs, and even fewer of these are ergodic. The simplest ergodic diffusions with closed-form expressions for the pdf are the Ornstein-Uhlenbeck and the CIR (Cox-Ingersoll-Ross) processes.
The Ornstein Uhlenbeck (OU) process is described by the SDE
$$dR_t=bR_tdt+\sigma dB_t.$$
The pdf for $`R_t`$ given $`R_0=r_0`$ is the normal density with mean $`r_0e^{bt}`$ and variance $`\sigma ^2(1e^{2\theta t})`$. Its stationary distribution is normal with mean 0 and variance $`\sigma ^2`$. The OU process was used by Aris-Brosou and Yang to model evolutionary rates. However, the OU process can take on negative values, and it is not clear how it can be used directly without any transformation, such as a reflected OU or a squared OU. Aris-Brosou and Yang also proposed another model, the EXP (for exponential) model, defined as the following : the rate assigned to a branch is drawn from an exponential distribution with mean equal to the rate of its ancestral branch. It is then obvious that their EXP model was a martingale. They outlined that the OU model seemed to provide a better fit to their data than the EXP model. Even if the reason of this better fit is still to be investigated, it seems reasonable to think that the ergodic property of the OU model could be a important factor. They also mentioned that the $`\sigma ^2`$ parameter of the OU model was hard to infer, perhaps because the OU model has an insufficient number of free parameters.
The use of the CIR model solves the problem, since it is a generalization of the squared OU process, where the mean and the variance can be independently inferred by the addition of a third parameter. If the mean is fixed to one, we avoid any identifiability problem with branch lengths without fixing the variance, which can therefore be inferred as well as the autocorrelation.
The CIR process satisfies the SDE
$$dR_t=b(aR_t)dt+\sigma \sqrt{R_t}dB_t,$$
(7)
and the pdf $`f_R(R_t|r_0)`$ for $`R_t`$ given $`R_0=r_0`$ is a non-central $`\chi ^2`$ distribution with degree of freedom $`4ab/\sigma ^2`$ and parameter of non-centrality $`\frac{4br_0e^{bt}}{\sigma ^2(1e^{bt})}`$. Its mean and variance are equal to
$`\mathrm{E}[r_t]`$ $`=`$ $`r_0e^{bt}+a(1e^{bt})`$ (8)
$`\mathrm{Var}[r_t]`$ $`=`$ $`r_0{\displaystyle \frac{\sigma ^2}{b}}(e^{bt}e^{2bt})+{\displaystyle \frac{a\sigma ^2}{2b}}(1e^{bt})^2.`$ (9)
The stationary distribution of $`R_t`$ is a gamma distribution with shape parameter $`2ab/\sigma ^2`$ and scale parameter $`\sigma ^2/2b`$. Hence the mean of the stationary distribution is $`a`$ and the variance is $`\frac{a\sigma ^2}{2b}`$ .
Unlike an OU process, if $`r_0`$, $`a`$, and $`b`$ are all positive a CIR process is always non-negative. The square of an OU process is a special case of the CIR process. Furthermore, by multiplying $`R_t`$ by a constant in equation (7), we see that multiplying a CIR process by a positive constant gives another CIR process.
The covariance of the stationary CIR process can be exactly computed as
$`\rho (t)=\text{Cov}(R_0,R_t)={\displaystyle \frac{a\sigma ^2}{2b}}e^{bt}.`$ (10)
From this, (4) leads to a closed formula for the index of dispersion:
$`I_{CIR}(t)=1+{\displaystyle \frac{\sigma ^2}{b^3t}}(bt1+e^{bt}).`$
Thus
$$I_{CIR}(\mathrm{})=\underset{t\mathrm{}}{lim}I_{CIR}(t)=1+\frac{\sigma ^2}{b^2}.$$
(11)
In Section 2 we emphasized that the concept of autocovariance is close to the index of dispersion. As Zheng showed , the effect of complex infinitesimal rate matrices on the index of dispersion (with constant rate) is not likely to explain alone the observed large empirical values. If the rate varies, Cutler showed that an elevated index of dispersion can only be achieved if the time-scale of the rate process is approximately of the same magnitude as the substitution process itself. The CIR process provides the possibility to satisfy this property, while incorporating autocovariance. It is the consensus of these two ideas that should guide our choice for the rate of evolution.
From (7) we see that the CIR process possesses three parameters, namely the stationary mean $`a`$, the stationary variance $`\sigma ^2`$, and the intensity of the force that drives the process to its stationary distribution, $`b`$. The parameter $`b`$ determines how fast the process autocovariance goes to 0 as $`t`$ increases.
The three parameters of the CIR process can be quickly estimated from standard statistics in molecular evolution. The parameter $`a`$ is a scale parameter. It determines the expected rate at any time given no other information. Throughout the paper, we will assume that $`a=1`$, so that the model has an expected rate equal to one. This parallels the constraint that the gamma distribution has an expected rate equal to one in the Rate-Across-Site (RAS) model .
The CIR process has a stationary distribution given by a gamma distribution. To make the stationary distribution coincide with the gamma distribution of a RAS model with parameter $`\mathrm{\Gamma }`$ we choose $`\sigma `$ and $`b`$ such that
$$\mathrm{\Gamma }=\frac{\sigma ^2}{b}.$$
(12)
The stationary index of dispersion, $`I_{CIR}(\mathrm{})`$, can be estimated empirically . We can then use (11) and (12) to obtain the estimates
$`\widehat{b}`$ $`=`$ $`{\displaystyle \frac{\widehat{\mathrm{\Gamma }}}{\widehat{I}_{CIR}(\mathrm{})1}},`$
$`\widehat{\sigma }^2`$ $`=`$ $`{\displaystyle \frac{\widehat{\mathrm{\Gamma }}^2}{\widehat{I}_{CIR}(\mathrm{})1}}.`$
## 4 Mutation models with a rate process
The standard model for the substitution process at a particular locus is a continuous-time Markov chain. This kind of process is defined by a square matrix $`Q`$ called the *infinitesimal rate matrix*. Suppose, to begin, that there is a constant evolutionary rate $`r_0`$. As above, we let $`X_t`$ denote the state (e.g. amino acid) at time $`t`$. The transition probabilities are then given by
$$\mathrm{Pr}[X_t=j|X_0=i]=[e^{Qr_0t}]_{ij}.$$
(13)
We suppose that the process has a unique stationary distribution $`\pi `$, where $`\pi _j`$ is the stationary probability of state $`j`$ and
$$\pi _j=\underset{t\mathrm{}}{lim}\mathrm{Pr}[X_t=j|X_0=i]$$
for all $`i`$ and $`j`$. We assume that $`Q`$ has been normalised so that in the stationary distribution the expected number of mutations over time $`t`$ equals $`r_0t`$. Note that the transition probabilities (13) depend only on the product $`r_0t`$, so will be the same if we double the rate and halve the time, for example.
Suppose now that the rate is not constant, but instead varies according to some fixed function $`r_s`$, $`s0`$. Equation (13) then becomes
$$\mathrm{Pr}[X_t=j|X_0=i,r]=[e^{Q\tau _r}]_{ij}.$$
(14)
where
$$\tau _r=_{s=0}^{s=t}r_s𝑑s$$
is the area under the curve $`r_s`$.
In the models we will consider, the fixed function $`r=(r_t)_{t0}`$ is replaced by a random process $`R=(R_t)_{t0}`$ that is dependent only on the starting rate $`r_0`$. The integral
$$\tau _R=_{s=0}^{s=t}R_s𝑑s$$
(15)
is also random in this case; let $`g_R`$ denote its pdf. The transition probabilities can be determined from the expected value of (14) with $`\tau _r`$ replaced by the random variable $`\tau _R`$. By the law of total expectation, this simplifies to
$$\mathrm{Pr}[X_t=j|X_0=i]=_\tau [e^{Q\tau }]_{ij}g_R(\tau )𝑑\tau .$$
(16)
Let $`M(\eta )=\mathrm{E}_\tau [e^{\eta \tau _R}]`$ denote the moment generating function (mgf) for the random variable $`\tau _{sR}`$. Then (16) can be rewritten
$$\mathrm{Pr}[X_t=j|X_0=i]=[M(Q)]_{ij}$$
where the function $`M`$ is interpreted as a matrix function . We assume that $`Q`$ can be diagonalised as $`Q=V\mathrm{\Lambda }V^1`$, where $`\mathrm{\Lambda }=\text{diag}(\lambda _1,\mathrm{},\lambda _n)`$ is a diagonal matrix formed from the eigenvalues of $`Q`$. The matrix function $`M(Q)`$ can then be evaluated as $`M(Q)=VM(\mathrm{\Lambda })V^1`$, where
$$M(\mathrm{\Lambda })=\text{diag}(M(\lambda _1),\mathrm{},M(\lambda _n)).$$
See for a more details on matrix functions. The problem of determining pattern probabilities therefore boils down to the problem of determining the moment generating function of the integrated rate, $`\tau _R`$ (eqn. 15). Tuffley and Steel use this approach to derive distance estimates for the covarion process .
For applications in phylogenetics, we need the mgf of $`\tau _R`$ conditioned on just the starting rate, or both the starting and finishing rate. The mgf of $`\tau _R`$, conditioned on a starting rate of $`r_0`$, is then
$`M_{r_0}(\eta )`$ $`=`$ $`\mathrm{E}\left[\mathrm{exp}(\eta \tau _R)|R_0=r_0\right]`$ (17)
$`=`$ $`\mathrm{E}\left[\mathrm{exp}\left(\eta {\displaystyle _{s=0}^t}R_s𝑑s\right)|R_0=r_0\right].`$
As before, we let $`f_R(R_t|r_0)`$ denote the pdf of $`R_t`$ conditioned on $`R_0=r_0`$. Let $`\delta (x)`$ denote the Dirac delta function with $`\delta (0)=1`$ and $`\delta (x)=0`$ for all $`x0`$. The mgf of $`\tau _R`$ conditioned on both the starting and finishing rates is
$`M_{r_0,r_t}(\eta )`$ $`=`$ $`\mathrm{E}\left[\mathrm{exp}(\eta \tau _R)|R_0=r_0,R_t=r_t\right]`$ (18)
$`=`$ $`{\displaystyle \frac{1}{f_R(r_t|r_0)}}\mathrm{E}\left[\mathrm{exp}(\eta \tau _R)\delta (R_tr_t)|R_0=r_0\right]`$
$`=`$ $`{\displaystyle \frac{1}{f_R(r_t|r_0)}}\mathrm{E}\left[\mathrm{exp}\left(\eta {\displaystyle _{s=0}^t}R_s𝑑s\right)\delta (R_tr_t)|R_0=r_0\right].`$
Equations (17) and (18) hold irrespective of whether $`R`$ is discrete or continuous, a diffusion, jump process, or a continuous time Markov chain.
We note in passing that analytic formulae for $`M_{r_0}(\eta )`$ and $`M_{r_0,r_t}(\eta )`$ exist in the case that $`R`$ is a continuous time Markov chain, for example in the covarion-type model of Galtier . Suppose that the evolutionary rate switches between rate values $`g_1,g_2,\mathrm{}g_k`$ following a continuous time Markov chain with infinitesimal rate matrix $`G`$. Let $`D`$ be the $`k\times k`$ diagonal matrix with entries $`g_1,g_2,\mathrm{},g_k`$. A careful reworking of the proof of Theorem 1 in gives the mgf of $`\tau _R`$ conditioned on both the starting and finishing rate. The mgf for $`\tau _R`$ conditioned on $`r_0=g_i`$ is then
$$M_{g_i}=\underset{j=1}{\overset{k}{}}(e^{(G+\eta D)t})_{ij}$$
while the mgf of $`\tau _R`$ conditioned on $`r_0=g_i`$ and $`r_t=g_j`$ is
$$M_{g_i,g_j}=\frac{(e^{(G+\eta D)t})_{ij}}{(e^{Gt})_{ij}}.$$
This provides an independent derivation of the formula in for transition probabilities under a covarion-type model.
## 5 Moment generating functions and transition probabilities for the CIR model
In this section we derive expressions for the (joint) transition probabilities
$`\mathrm{Pr}[X_t=j|X_0=i,R_0=r_0].`$ (19)
and
$`\mathrm{Pr}[X_t=j,R_t=s|X_0=i,R_0=r_0]`$ (20)
As we have seen, to evaluate these probabilities we need to determine the moment generating functions (mgfs) defined in equations (17) and (18).
We use the Feynman-Kac formula to derive analytic formulae for $`M_{r_0}(\eta )`$ and $`M_{r_t,r_0}(\eta )`$ under the CIR model. Let $`g()`$ be a real-valued function. Define the function $`v=v(t,x)`$ by
$`v(t,x)=\mathrm{E}\left[\mathrm{exp}\left(\eta {\displaystyle _0^t}R(s)𝑑s\right)g(R_t)|R_0=x\right]`$ (21)
The Feynman-Kac formula asserts that $`v(t,x)`$ solves the following partial differential equation (PDE)
$`{\displaystyle \frac{}{t}}v(t,x)=b(1x){\displaystyle \frac{}{x}}v(t,x)+{\displaystyle \frac{1}{2}}\sigma ^2x(t,x){\displaystyle \frac{^2}{x^2}}v(t,x)+\eta xv`$ (22)
for $`t>0`$, $`x`$, and with boundary condition
$`v(0,x)=g(x)\text{for all }x\text{.}`$ (23)
We apply the methods in and to solve these pdes with the different boundary conditions.
First consider the case when we condition only on the initial rate, eqn. (17). To make (21) equal to (17) we set $`g(x)=1`$ for all $`x`$. The boundary condition (23) in this case becomes
$$v(0,x)=1\text{for all }x\text{.}.$$
With this boundary condition, the pde (22) has solution
$$v(t,x)=\mathrm{\Psi }(\eta ,t)e^{x\mathrm{\Xi }(\eta ,t)}$$
where
$`\mathrm{\Psi }(\eta ,t)`$ $`=`$ $`\left({\displaystyle \frac{\overline{b}e^{bt/2}}{\overline{b}\mathrm{cosh}(\overline{b}t/2)+b\mathrm{sinh}(\overline{b}t/2)}}\right)^{\frac{2b}{\sigma ^2}}`$ (24)
$`\mathrm{\Xi }(\eta ,t)`$ $`=`$ $`\left({\displaystyle \frac{2\eta \mathrm{sinh}(\overline{b}t/2)}{\overline{b}\mathrm{cosh}(\overline{b}t/2)+b\mathrm{sinh}(\overline{b}t/2)}}\right),`$ (25)
$`\overline{b}`$ $`=`$ $`\sqrt{b^22\eta \sigma ^2}`$ (26)
We therefore have
$`M_{r_0}(\eta )=\mathrm{\Psi }(\eta ,t)e^{r_0\mathrm{\Xi }(\eta ,t)}.`$ (27)
The case when both the starting and finished rates are specified is more complicated. From (18) the mgf $`M_{r_0,r_t}`$ can be written
$$M_{r_0,r_t}=\frac{1}{f_R(r_t|r_0)}v(t,x)$$
where, in this case, $`v(t,x)`$ is given by (21) with $`g(x)=\delta (xr_t)`$. The boundary condition (23) therefore becomes
$$v(0,x)=\delta (xr_t).$$
With this new boundary condition, the pde (22) has solution
$`v(t,x)`$ $`=`$ $`c\mathrm{exp}\left[{\displaystyle \frac{bt}{\sigma ^2}}(\overline{b}b)+{\displaystyle \frac{b\overline{b}}{\sigma ^2}}x{\displaystyle \frac{b+\overline{b}}{\sigma ^2}}r_tc(r_t+x)e^{\overline{b}t}\right]`$ (28)
$`\times \left({\displaystyle \frac{r_t}{xe^{\overline{b}t}}}\right)^{\frac{b}{\sigma ^2}1/2}I_{\frac{2b}{\sigma ^2}1}\left(2c\sqrt{xr_te^{\overline{b}t}}\right),`$
where
$`c`$ $`=`$ $`{\displaystyle \frac{2\overline{b}}{\sigma ^2(1e^{\overline{b}t})}}`$
$`\overline{b}`$ $`=`$ $`\sqrt{b^22\eta \sigma ^2}`$
and $`I_\nu (x)`$ is the modified Bessel function of the first kind with parameter $`\nu `$ . Hence the mgf conditioned on initial and final rates is given by
$`M_{r_0,r_t}(\eta )`$ $`=`$ $`c\mathrm{exp}\left[{\displaystyle \frac{bt}{\sigma ^2}}(\overline{b}b)+{\displaystyle \frac{b\overline{b}}{\sigma ^2}}r_0{\displaystyle \frac{b+\overline{b}}{\sigma ^2}}r_tc(r_t+r_0)e^{\overline{b}t}\right]`$
$`\times \left({\displaystyle \frac{r_t}{r_0e^{\overline{b}t}}}\right)^{\frac{b}{\sigma ^2}1/2}I_{\frac{2b}{\sigma ^2}1}\left(2c\sqrt{r_0r_te^{\overline{b}t}}\right){\displaystyle \frac{1}{f_R(r_t|R_0=r_0)}},`$
where $`c`$ and $`\overline{b}`$ are defined above and, from Section 3, $`f_R(R_t|R_0=r_0)`$ is the pdf for a non-central $`\chi ^2`$ distribution with degree of freedom $`4ab/\sigma ^2`$ and parameter of non-centrality $`\frac{4br_0e^{bt}}{\sigma ^2(1e^{bt})}`$,
Bringing everything together, we have our main result.
###### Theorem 1
Define $`P`$ by $`P_{ij}=\mathrm{Pr}[X_t=j,R_t=r_t|X_0=i,R_0=r_0]`$. Suppose that $`Q=V\mathrm{\Lambda }V^1`$ where $`\mathrm{\Lambda }`$ is a diagonal matrix containing the eigenvalues $`\lambda _1,\mathrm{},\lambda _n`$ of $`Q`$. Then
$`P=M(Q)=VM(\mathrm{\Lambda })V^1`$ (29)
where $`M(\mathrm{\Lambda })`$ is the diagonal matrix where, for all $`i`$,
$$M(\mathrm{\Lambda })_{ii}=M_{r_t,r_0}(\lambda _i).$$
## 6 Three taxa phylogenies
The simplest phylogeny for which we can distinguish between constant and variable evolutionary rates is a tree with three taxa. For this reason, there are many methods for testing, and estimating, rate variation that are based on three taxon analyses . Here we show that the likelihood for a three taxa tree, under the CIR model of rate variation, can be computed exactly. The problem for general phylogenies is more complex since we have to integrate out rates for the internal nodes. Here, we consider a heterotachous model, so that each site has its own rate history. Because the sites (and the rate at each site) evolve independently from each other, the likelihood of a sequence will be the product of all site-specific likelihoods. Therefore, we only require the likelihood computation for one site.
We recall that the stationary distribution of the CIR is a gamma distribution $`\mathrm{\Gamma }(\omega ,\nu )`$, where $`\omega =2b/\sigma ^2`$ and $`\nu =2b/\sigma ^2`$, i.e.
$`f_{R_0}(r)={\displaystyle \frac{\omega ^\nu }{\mathrm{\Gamma }(\nu )}}r^{\nu 1}e^{\omega r}.`$ (30)
Therefore the stationary mean and variance are $`1`$ and $`\sigma ^2/2b`$.
In order to get the transition probabilities, we will use the mgf of $`\tau _R`$ unconditioned on the final rate, given by equation (27). The transition probability matrix of the subsitution process, given initial rates, can be obtained by equations (27) and (29). Let $`\lambda _1,\mathrm{},\lambda _n`$ be the eigenvalues of $`Q`$. Using eigenvalue decomposition, we can find vectors $`𝐮^{(1)},\mathrm{},𝐮^{(n)}`$ and $`𝐯^{(1)},\mathrm{},𝐯^{(n)}`$ such that
$$\mathrm{Pr}[X_t=i|X_0=j,R_0=r_0]=\underset{k=1}{\overset{n}{}}𝐮_{j}^{(k)}{}_{}{}^{T}𝐯_i^{(k)}M_{r_0}(\lambda _k,t),$$
(31)
where we changed slightly our notation and explicitly wrote the dependency of $`M_{r_0}`$ on $`t`$.
Now consider the 3-taxa tree with branches of lengths $`t_1,t_2,t_3`$ leading to leaves labelled with states $`x_1,x_2,x_3`$ (Figure 2). If we condition on a rate $`r_0`$ and state $`x_0`$ at the root then the probability of observing $`x_1,x_2,x_3`$ at the leaves is given by
$`L(x_1,x_2,x_3|x_0,r_0)`$ $`=`$ $`P[X_{t_1}=x_1|x_0,r_0]P[X_{t_2}=x_2|x_0,r_0]P[X_{t_3}=x_3|x_0,r_0]`$ (32)
$`=`$ $`{\displaystyle \underset{i=1}{\overset{n}{}}}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \underset{k=1}{\overset{n}{}}}B_{ijk}M_{r_0}(\lambda _i,t_1)M_{r_0}(\lambda _j,t_2)M_{r_0}(\lambda _k,t_3)`$
where
$$B_{ijk}=𝐮_{x_0}^{(i)}𝐯_{x_1}^{(i)}𝐮_{x_0}^{(j)}𝐯_{x_2}^{(j)}𝐮_{x_0}^{(k)}𝐯_{x_3}^{(k)}.$$
The rate at the root is assumed to have the stationary distribution $`f_{R_0}`$ given by (30). The likelihood integrated with respect to $`r_0`$ is then
$`L(x_1,x_2,x_3|x_0)`$ $`=`$ $`{\displaystyle _{r_0}}L(x_1,x_2,x_3|x_0,r_0)f_{R_0}(r_0)𝑑r_0`$
which by (32) equals
$$\underset{i=1}{\overset{n}{}}\underset{j=1}{\overset{n}{}}\underset{k=1}{\overset{n}{}}B_{ijk}_{r_0}M_{r_0}(\lambda _i,t_1)M_{r_0}(\lambda _j,t_2)M_{r_0}(\lambda _k,t_3)f_{R_0}(r_0)𝑑r_0.$$
(33)
We now use the formula (27) for the mgf’s derived above.
$`M_{r_0}(\lambda _i,t_1)M_{r_0}(\lambda _j,t_2)M_{r_0}(\lambda _k,t_3)f_{R_0}(r_0)dr_0`$ (34)
$`=`$ $`\mathrm{\Psi }(\lambda _i,t_1)e^{r_0\mathrm{\Xi }(\lambda _i,t_1)}\mathrm{\Psi }(\lambda _j,t_2)e^{r_0\mathrm{\Xi }(\lambda _j,t_2)}\mathrm{\Psi }(\lambda _k,t_3)e^{r_0\mathrm{\Xi }(\lambda _k,t_3)}{\displaystyle \frac{\omega ^\nu }{\mathrm{\Gamma }(\nu )}}r_0^{\nu 1}e^{\omega r_0}`$
$`=`$ $`\mathrm{\Psi }(\lambda _i,t_1)\mathrm{\Psi }(\lambda _j,t_2)\mathrm{\Psi }(\lambda _k,t_3){\displaystyle \frac{\omega ^\nu }{\mathrm{\Gamma }(\nu )}}r_0^{\nu 1}e^{r_0\left(\omega +\mathrm{\Xi }(\lambda _i,t_1)+\mathrm{\Xi }(\lambda _j,t_2)+\mathrm{\Xi }(\lambda _k,t_3)\right)}`$
Using integration by parts, or simply using the fact that the gamma distribution integrates to 1, we get
$`{\displaystyle _{r_0}}M_{r_0}(\lambda _i,t_1)M_{r_0}(\lambda _j,t_2)M_{r_0}(\lambda _k,t_3)f_{R_0}(r_0)dr_0`$
$`=\mathrm{\Psi }(\lambda _i,t_1)\mathrm{\Psi }(\lambda _j,t_2)\mathrm{\Psi }(\lambda _k,t_3)\left({\displaystyle \frac{\omega }{\omega +\mathrm{\Xi }(\lambda _i,t_1)+\mathrm{\Xi }(\lambda _j,t_3)+\mathrm{\Xi }(\lambda _k,t_3)}}\right)^\nu .`$
Finally, we can substitute this back into (33) to obtain
$`L(x_1,x_2,x_3|x_0)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{n}{}}}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \underset{k=1}{\overset{n}{}}}B_{ijk}\mathrm{\Psi }(\lambda _i,t_1)\mathrm{\Psi }(\lambda _j,t_2)\mathrm{\Psi }(\lambda _k,t_3)`$
$`\times \left({\displaystyle \frac{\omega }{\omega +\mathrm{\Xi }(\lambda _i,t_1)+\mathrm{\Xi }(\lambda _j,t_3)+\mathrm{\Xi }(\lambda _k,t_3)}}\right)^\nu .`$
The formula extends immediately to phylogenies with $`n`$ leaves attached to the root, though the number of terms in the summation increases exponentially. Our approach has been to evaluate likelihoods on complete phylogenies using Monte-Carlo techniques, together with the exact transition probabilities derived here.
## 7 Discussion
### 7.1 Summary
We have shown how, given a few natural criteria for our model selection, the CIR appeared as the simplest continuous model that is at the same time ergodic, has a non-zero autocovariance function and that can account for an arbitrarily large index of dispersion. Moreover, we provided simple ways to estimate its parameters with the help of two observable statistics, namely the RAS gamma parameter and empirical index of dispersion. Another very interesting practical aspect of the CIR process is that it can be easily, and without approximations, implemented in the MCMC framework.
### 7.2 Future extensions
A possible future extension of our model could involve jump models, in which the rate path is discontinuous as in the continuous-time Markov chain, but also varies as diffusion between these discontinuities. However, the use of such a model implies the use of more parameters, and it may well be the case that the relative weakness of the rate of evolution signal cannot allow the use of more than two parameters, because of identifiability problems. |
warning/0506/hep-th0506097.html | ar5iv | text | # Membrane Instantons and de Sitter Vacua
## 1 Introduction
A central question in string theory is the existence and viability of “semi-realistic” four-dimensional ground states. In this context studying the vacuum structure arising from flux compactifications has recently attracted considerable attention. In particular, (KKLT) provided a qualitative picture for obtaining meta-stable de Sitter (dS) vacua from compactifications of the type IIB string, in which fluxes and non-perturbative instanton effects play a crucial role. In this paper we consider membrane instanton corrections arising in the compactification of the type IIA string on rigid Calabi-Yau threefolds (CY<sub>3</sub>) and show that including background fluxes and these non-perturbative corrections can provide another scenario to stabilize all hypermultiplet moduli at a meta-stable de Sitter vacuum.
The four-dimensional low-energy effective actions for string compactifications preserving some supersymmetry are supergravity actions coupled to matter multiplets. When fluxes are turned on, one typically obtains gauged supergravities with a potential for the scalar fields of the matter multiplets. The properties and extrema of such potentials are of great importance for string cosmology, since they determine the vacuum structure of the theory. In recent years, string theorists have searched intensively for models in which the potential admits vacua with a (small) positive cosmological constant. It turned out that it is fairly difficult to realize such vacua in string theory, as they can only be meta-stable (see e.g. for a review).
A qualitative picture on how such vacua can be obtained was given in in the context of type IIB flux compactifications on orientifolds. In this case the four-dimensional effective action has $`N=1`$ supersymmetry, and the potential is determined by a holomorphic superpotential. The KKLT scenario relies on three contributions to the superpotential: first there is a classical contribution coming from fluxes which stabilizes all moduli except the volume modulus which does not enter into a scalar potential of no-scale type. This modulus is then stabilized by a non-perturbative contribution to the potential due to D-instantons or gaugino condensation. These two ingredients stabilize all moduli in a supersymmetric AdS vacuum. In the third step a (small) positive energy contribution, as e.g. an anti-D3-brane, is added which lifts the AdS vacuum to a positive cosmological constant. Since its first proposal, possible realizations of this scenario either within type IIB orientifold compactifications or their F-theory descriptions have been studied intensively .
One of the goals of this paper is to provide an alternative scenario in the context of type IIA string theory compactified on a CY<sub>3</sub>. Without including background fluxes the LEEA arising from these compactifications is a four-dimensional $`N=2`$ supergravity action coupled to $`h_{1,1}`$ vector and $`h_{1,2}+1`$ hypermultiplets. There is no scalar potential and the scalars (moduli) of the theory parameterize flat directions. The coupling to $`N=2`$ supergravity requires the scalars of the hypermultiplets to parameterize a quaternion-Kähler manifold . The dilaton that controls the quantum corrections sits in a hypermultiplet (the universal hypermultiplet), and hence it is the quaternionic geometry that receives quantum corrections. Besides perturbative corrections, there are also non-perturbative instanton effects obtained by wrapping Euclidean D-branes around supersymmetric cycles of the internal manifold . From the counting of fermionic zero modes one can derive that they contribute to the low-energy effective action. In the KKLT models they contribute to the superpotential for the $`N=1`$ chiral multiplets, whereas in our case they correct the hypermultiplet scalar metric.
In this paper, we focus on the special case of the universal hypermultiplet, which can be obtained by compactifying on a rigid CY<sub>3</sub>, having $`h_{1,2}=0`$. We restrict ourselves to rigid CY manifolds, because we will be able to explicitly determine the instanton corrections in this special case only. The general situation when more complex structure moduli are present is technically more difficult because of the complicated nature of the quaternion-Kähler geometry. We believe, however, that our main conclusion will still persist in this case.
The classical quaternionic geometry of the universal hypermultiplet is well-known , and recently the perturbative corrections were found in , see also . Non-perturbatively, there are both membrane and NS fivebrane instanton corrections , but in this paper we shall consider membrane instantons only.<sup>1</sup><sup>1</sup>1For work on fivebrane instantons, we refer the reader to . Additional references on hypermultiplet moduli spaces and instantons are . Furthermore a program towards formulating an instanton calculus based on $`N=2`$ supersymmetric actions with Euclidean signature was started in . In this case the constraints from quaternionic geometry are captured by solutions of the three-dimensional Toda equation. This fact was, to our knowledge, first observed in (see also ). One of the main results of this paper is that we construct new solutions of the Toda equation that correspond to membrane instanton expansions. We have not uniquely fixed the solution, and at each order in the instanton expansion, there is still an undetermined integration coefficient that can in principle be computed in string theory. The solution of the Toda equation then determines the quaternion-Kähler (QK) metric in the ungauged supergravity effective action.<sup>2</sup><sup>2</sup>2Membrane instantons were also considered in , but our analysis below differs since we do *not* assume the existence of a rotational symmetry between the RR scalars in the UHM scalar metric. In fact our analysis will show that this isometry is broken. As we will show, our results are in complete agreement with the predictions made in .
Including background fluxes in the compactification leads to four-dimensional $`N=2`$ gauged supergravity where some isometries of the hypermultiplet scalar manifold are gauged .<sup>3</sup><sup>3</sup>3For an analysis on de Sitter vacua, purely in the context of N=2 supergravity, we refer to . This gauging induces a scalar potential in the LEEA which depends on the geometrical quantities of the QK space, such as the moment maps and the metric. It is therefore clear that the potential will receive quantum corrections, determined e.g. by the QK metric. We must be careful with this procedure, since isometries of the classical hypermultiplet moduli space can be broken by quantum corrections. This is already the case perturbatively . Non-perturbatively, isometries can get broken to discrete subgroups. To gauge an isometry in supergravity, the standard methods require an unbroken and continuous isometry. However, in the absence of fivebrane instantons, we explain how to find such an isometry, and moreover we show how the corresponding potential can be obtained from a flux compactification of the type discussed in .
In both the KKLT models with $`N=1`$, and as we will see, in our models with $`N=2`$, it is crucial to take into account the quantum corrections to the low-energy effective action. In particular, including the instanton corrections to the potential is an essential step for stabilizing the dilaton and finding meta-stable de Sitter vacua.<sup>4</sup><sup>4</sup>4Based on the instanton corrected UHM of , a similar analysis, also indicating the existence of meta-stable dS vacua, was performed in . This was the case in KKLT, and also applies to our models.<sup>5</sup><sup>5</sup>5Similar observations have also been made in heterotic M-theory (see e.g. ). In our set-up, we only study the hypermultiplet moduli in detail, and comment on the Kähler moduli at the end of the paper. In that case, the potential only depends on the hypermultiplet scalars and is determined by the solution of the Toda equation. As our solution still contains undetermined integration constants (which, in principle, should be determined by string theory), it is therefore perhaps not too surprising that one can choose coefficients that give de Sitter vacua. In a way, choosing these coefficients mimics stabilizing the volume modulus in the KKLT set-up.
The remainder of the paper is organized as follows. In section 2 we begin by describing the supergravity set-up for our investigations. The moduli space metric of the universal hypermultiplet is introduced and its possible quantum corrections are discussed qualitatively. We then show in section 3 how this metric fits into a general framework for four-dimensional QK geometries with one isometry, which are governed by the three-dimensional Toda equation. In section 4 we derive the leading terms of a solution to this equation describing non-perturbative quantum effects due to membrane instantons. Section 5 is devoted to a comparison of our results with string theory predictions on how these instanton corrections contribute to the four-fermion coupling; we shall find perfect agreement. Finally, in section 6 we investigate the effects of these corrections on the scalar potential that arises by gauging the one remaining isometry of the moduli space metric. It turns out that the undetermined parameters can be such that the potential develops a local meta-stable de Sitter minimum. After the conclusions we provide technical details in several appendices.
## 2 Supergravity description
For type IIA string theories compactified on a CY<sub>3</sub> manifold, the low-energy effective action is that of four-dimensional $`N=2`$ supergravity coupled to $`h_{1,1}`$ vector multiplets, $`h_{1,2}`$ hypermultiplets, and one tensor multiplet that contains the dilaton . In the case of a rigid CY<sub>3</sub>, there are no complex structure moduli: $`h_{1,2}=0`$. Suppressing the vector multiplets, the resulting four-dimensional low-energy effective action is that of a tensor multiplet coupled to $`N=2`$ supergravity, and the bosonic part of the Lagrangian at string tree-level is given by<sup>6</sup><sup>6</sup>6Throughout this paper, we work in units in which Newton’s constant $`\kappa ^2=2`$.
$`e^1_\mathrm{T}`$ $`=R{\displaystyle \frac{1}{2}}^\mu \varphi _\mu \varphi +{\displaystyle \frac{1}{2}}\mathrm{e}^{2\varphi }H^\mu H_\mu `$
$`{\displaystyle \frac{1}{4}}F^{\mu \nu }F_{\mu \nu }{\displaystyle \frac{1}{2}}\mathrm{e}^\varphi \left(^\mu \chi _\mu \chi +^\mu \phi _\mu \phi \right){\displaystyle \frac{1}{2}}H^\mu \left(\chi _\mu \phi \phi _\mu \chi \right),`$ (1)
where $`H^\mu =\frac{1}{6}\epsilon ^{\mu \nu \rho \sigma }H_{\nu \rho \sigma }`$ is the dual NS 2-form field strength. The first line comes from the NS sector in ten dimensions, and $`\varphi `$ together with $`H^\mu `$ forms an $`N=1`$ tensor multiplet. The second line descends from the RR sector. In particular, the graviphoton with field strength $`F_{\mu \nu }`$ descends from the ten-dimensional RR 1-form, and $`\phi `$ and $`\chi `$ can be combined into a complex scalar $`C`$ that descends from the holomorphic components of the RR 3-form with (complex) indices along the holomorphic 3-form of the CY<sub>3</sub>. Notice the presence of constant shift symmetries on both $`\chi `$ and $`\phi `$. Together with a rotation on $`\chi `$ and $`\phi `$ they form a three-dimensional subgroup of symmetries.
The tensor multiplet Lagrangian (2) is dual to the universal hypermultiplet. This can be seen by dualizing the 2-form into an axionic pseudoscalar field $`\sigma `$, after which one obtains (modulo a surface term)
$`e^1_{\mathrm{UH}}`$ $`=R{\displaystyle \frac{1}{4}}F^{\mu \nu }F_{\mu \nu }{\displaystyle \frac{1}{2}}^\mu \varphi _\mu \varphi {\displaystyle \frac{1}{2}}\mathrm{e}^\varphi \left(^\mu \chi _\mu \chi +^\mu \phi _\mu \phi \right)`$
$`{\displaystyle \frac{1}{2}}\mathrm{e}^{2\varphi }\left(_\mu \sigma +\chi _\mu \phi \right)^2.`$ (2)
The four scalars define the classical universal hypermultiplet at string tree-level, a non-linear sigma model with a quaternion-Kähler target space $`\mathrm{SU}(1,2)/\mathrm{U}(2)`$ . The metric can be written as
$$\mathrm{d}s^2=G_{AB}\mathrm{d}\varphi ^A\mathrm{d}\varphi ^B=\mathrm{d}\varphi ^2+\mathrm{e}^\varphi (\mathrm{d}\chi ^2+\mathrm{d}\phi ^2)+\mathrm{e}^{2\varphi }(\mathrm{d}\sigma +\chi \mathrm{d}\phi )^2.$$
(3)
This manifold has an $`\mathrm{SU}(1,2)`$ group of isometries, with a three-dimensional Heisenberg subalgebra that generates the following shifts on the fields,
$$\varphi \varphi ,\chi \chi +\gamma ,\phi \phi +\beta ,\sigma \sigma \alpha \gamma \phi ,$$
(4)
where $`\alpha `$, $`\beta `$, $`\gamma `$ are real (finite) parameters.
Quantum corrections, both perturbative and non-perturbative, will break some of the isometries and alter the classical moduli space of the universal hypermultiplet, while keeping the quaternion-Kähler property intact, as required by supersymmetry . At the perturbative level, a non-trivial one-loop correction modifies the low-energy tensor multiplet Lagrangian (2), as was shown in . After dualization, this corrects the universal hypermultiplet metric (3), while still preserving the isometries (4). More recently, this one-loop correction was written and analyzed in the language of projective superspace in , using the tools developed in .
At the non-perturbative level, there can be membrane and fivebrane instantons. The latter were analyzed in . Membrane instantons, which we are focussing on in this paper, arise from wrapping Euclidean D2-branes around three-cycles in the CY<sub>3</sub> . For rigid Calabi-Yau’s, there are two kind of membrane instantons, as there are two (supersymmetric) three-cycles that the membrane can wrap around. Correspondingly, there will be two membrane instanton charges. These instantons also have an effective supergravity description, as was shown in . The two instanton charges correspond to the shift symmetries in $`\chi `$ and $`\phi `$, as written down in (4). We denote these charges by $`Q_\chi `$ and $`Q_\phi `$ respectively. They can also be understood as being the charges of the corresponding dual 3-form field strengths that appear after dualizing one of the scalars $`\chi `$ or $`\phi `$ to a 2-form. Upon doing so, the tensor multiplet becomes a double-tensor multiplet, in which the instanton solution can be derived from a Bogomol’nyi equation . Following this procedure, it becomes clear that only one charge can be switched on simultaneously, either $`Q_\chi `$ or $`Q_\phi `$, depending on which scalar was dualized to a tensor. One cannot dualize both scalars to tensors, as the two shift symmetries on $`\chi `$ and $`\phi `$ do not commute. In section 5, we will rederive this property from a string theory perspective.
The instanton action is inversely proportional to the string coupling, which, in our conventions, is defined as
$$g_s=\mathrm{e}^{\varphi _{\mathrm{}}/2}.$$
(5)
The membrane instanton action, say for the $`\phi `$-instanton, then is
$$S_{\mathrm{inst}}=2\frac{|Q_\phi |}{g_s}+\mathrm{i}\phi Q_\phi .$$
(6)
The imaginary term comes from a surface term that arises upon dualizing the tensor to a scalar. It involves the zero mode of the dual scalar $`\phi `$, which can be identified with the value of the field at infinity. Its presence breaks the shift symmetry in $`\phi `$ to a discrete subgroup. A similar formula also holds for the $`\chi `$-instanton, by simply replacing $`\phi `$ by $`\chi `$. Notice also the factor $`2`$ in front of the real part of the instanton action (6). This will become important later.
To compare, the NS-fivebrane instanton action is inversely proportional to the square of the string coupling and, in the same normalization as above, has no factor of 2 in front . It has a theta-angle-like term proportional to the zero mode of $`\sigma `$. As long as we don’t switch on fivebrane instantons, the continuous shift symmetry in $`\sigma `$ will remain an exact symmetry. In other words, in the absence of fivebrane instantons, the quantum corrected universal hypermultiplet moduli space will be a quaternionic manifold with a (non-compact) U(1) isometry. Such manifolds have been classified by mathematicians in terms of a single function, as we describe in the next section.
## 3 Toda equation and universal hypermultiplet
As explained in the previous section, the effect of membrane instantons is to modify the hypermultiplet moduli space non-perturbatively, in a way consistent with the constraints from quaternion-Kähler (QK) geometry. In the absence of fivebranes the quaternionic manifold has an isometry that acts as a shift in the NS scalar $`\sigma `$. In this section, we discuss the geometry of QK manifolds with a U(1) isometry, and explain how the universal hypermultiplet fits into this framework.
### 3.1 The Przanowski-Tod metric
In Przanowski derived the general form of four-dimensional quaternion-Kähler metrics with (at least) one Killing vector. It was later rederived by Tod . The Przanowski-Tod (PT) metric in local coordinates $`(r,u,v,t)`$ reads
$$\mathrm{d}s^2=\frac{1}{r^2}\left[f\mathrm{d}r^2+f\mathrm{e}^h(\mathrm{d}u^2+\mathrm{d}v^2)+f^1(\mathrm{d}t+\mathrm{\Theta })^2\right].$$
(7)
The isometry acts as a shift in the coordinate $`t`$. The metric is determined in terms of one scalar function $`h(r,u,v)`$, which is subject to the three-dimensional Toda equation
$$(_u^2+_v^2)h+_r^2\mathrm{e}^h=0.$$
(8)
The function $`f(r,u,v)`$ is not independent, but related to $`h`$ through
$$f=\frac{3}{2\mathrm{\Lambda }}\left(2r_rh\right),$$
(9)
while the 1-form $`\mathrm{\Theta }(r,u,v)=\mathrm{\Theta }_r\mathrm{d}r+\mathrm{\Theta }_u\mathrm{d}u+\mathrm{\Theta }_v\mathrm{d}v`$ is a solution to the equation
$$\mathrm{d}\mathrm{\Theta }=(_uf\mathrm{d}v_vf\mathrm{d}u)\mathrm{d}r+_r(f\mathrm{e}^h)\mathrm{d}u\mathrm{d}v.$$
(10)
Manifolds with such a metric are Einstein with anti-selfdual Weyl tensor, and $`\mathrm{\Lambda }`$ in (9) is the target space cosmological constant, $`R_{AB}=\mathrm{\Lambda }G_{AB}`$.
As long as at least one isometry remains unbroken, the universal hypermultiplet moduli space metric (3) is of this form. Its Ricci tensor is found to be $`R_{AB}=(3/2)G_{AB}`$, thus $`\mathrm{\Lambda }=3/2`$ in our conventions.<sup>7</sup><sup>7</sup>7More on our conventions on quaternionic geometry can be found in appendix A.
It is quite remarkable that the non-perturbative structure of the universal hypermultiplet is fully encoded by the solutions of the Toda equation. This equation has been studied by mathematicians in the context of three-dimensional Einstein-Weyl spaces and hyperkähler manifolds (see also appendix B). More recently, a large class of solutions of the Toda equation was constructed by , see also . Unfortunately these do not seem to satisfy the boundary conditions required by our set-up, so in the next section we will construct new solutions that describe membrane instanton effects.
Integrable structures, including the Toda hierarchy, have also been discovered in topological string theory . Related to this, the Toda equation also appears in the non-perturbative description of the non-critical $`c=1`$ string theory . It would be interesting to better understand the connection, if any, to our work. Finally, we mention that the Toda equation also plays an important role in classifying BPS vacua in M-theory .
### 3.2 Symmetries, moment maps, and 4-fermion couplings
Clearly, the PT metric has a Killing vector $`_t`$ corresponding to a shift symmetry in $`t`$. In coordinates $`(r,u,v,t)`$, this Killing vector is given by
$$k^A=(\mathrm{\hspace{0.17em}0},\mathrm{\hspace{0.17em}0},\mathrm{\hspace{0.17em}0},e_0)^\mathrm{T},e_0.$$
(11)
The moment maps of the shift symmetry can be computed from (78), which we do in appendix A.2. The result is *independent* of the functions $`f`$, $`h`$ and $`\mathrm{\Theta }`$, and reads
$$P^1=0,P^2=0,P^3=\frac{e_0}{r}.$$
(12)
Furthermore, $`4n`$-dimensional quaternion-Kähler manifolds admit a completely symmetric rank four tensor $`𝒲_{\alpha \beta \gamma \delta }`$, where $`\alpha =1,\mathrm{},2n`$ labels the $`\mathrm{USp}(2n)`$ index that is part of the holonomy group of QK manifolds. This tensor can be constructed out of the Riemann curvature tensor; its definition and properties are discussed in , which we summarize in appendix A.3. In $`N=2`$ supergravity effective actions, the $`𝒲`$-tensor is contracted with four hyperinos; it will play an important role in section 5. In our case, the QK manifold is four-dimensional and hence $`n=1`$.
For the PT metric we carry out its construction in appendix A.3 and state here only the final result:
$`𝒲_{1111}`$ $`=4r^2f^3\mathrm{e}^h\left[f(_{\overline{z}}^2f_{\overline{z}}h_{\overline{z}}f)3(_{\overline{z}}f)^2\right]`$
$`𝒲_{2111}`$ $`=r^2f^3\mathrm{e}^{h/2}\left[2f_r_{\overline{z}}f3(f_rh+2_rf)_{\overline{z}}f+f^2_r_{\overline{z}}h\right]`$
$`𝒲_{2211}`$ $`=r^2f^3\left[f\left(r_r^3h(_rh)^2\right)4\mathrm{e}^h_zf_{\overline{z}}f+2(_rf)^2\right]`$
$`𝒲_{2221}`$ $`=r^2f^3\mathrm{e}^{h/2}\left[2f_r_zf3(f_rh+2_rf)_zf+f^2_r_zh\right]`$
$`𝒲_{2222}`$ $`=4r^2f^3\mathrm{e}^h\left[f(_z^2f_zh_zf)3(_zf)^2\right].`$ (13)
Here we have introduced the complex variable $`z=u+\mathrm{i}v`$ in order to write the components of $`𝒲_{\alpha \beta \gamma \delta }`$ in a compact way. We will use this tensor in a comparison of the properties of our instanton corrected universal hypermultiplet metric with the results for four-fermi correlation functions computed in string theory .
### 3.3 The universal hypermultiplet in the PT framework
To rewrite the metric (3) in the PT form, we have to identify the moduli of the universal hypermultiplet with the PT coordinates. This must be done consistently with the isometries, in particular with the shift symmetry in the coordinate $`t`$. From the Heisenberg algebra of isometries (4) it is apparent that one can choose to identify $`t`$ with either $`\sigma `$ or $`\phi `$. The shift symmetries are generated by the parameters $`\alpha `$ and $`\beta `$, respectively. This leads to two ‘dual’ representations of the PT metric that describe the same moduli space. We can call these bases the membrane and the fivebrane basis, respectively.
In the membrane basis, which is the relevant basis for our purposes, we identify the coordinate $`t`$ with $`\sigma `$, such that the $`\alpha `$-shift symmetry is manifest. This is because of the absence of fivebrane instantons, which would break the continuous $`\alpha `$-shift symmetry to a discrete subgroup . So, the coordinates can be chosen as
$$t=\sigma ,r=\mathrm{e}^\varphi ,u=\chi ,v=\phi .$$
(14)
In this basis, the classical moduli space metric of the universal hypermultiplet corresponds to the solution $`\mathrm{e}^h=r`$ of the Toda equation (8). It follows that $`f=1`$ and $`\mathrm{\Theta }=u\mathrm{d}v`$, the latter being defined only modulo an exact form.
As mentioned above, besides the instanton contributions that we want to determine in this paper, there are also perturbative quantum corrections to the moduli space metric . These can easily be incorporated in our approach: Observe that with $`h(r,u,v)`$ also $`h(r+c,u,v)`$ is a solution to the Toda equation for constant $`c`$. Applied to the classical solution $`\mathrm{e}^h=r`$, we obtain
$$\mathrm{e}^h=r+c,f=\frac{r+2c}{r+c},\mathrm{\Theta }=u\mathrm{d}v,$$
(15)
which turns out to describe the 1-loop (in the string frame) corrected metric of if we identify
$$c=\frac{4\zeta (2)\chi (X)}{(2\pi )^3}=\frac{1}{6\pi }(h_{1,1}h_{1,2}).$$
(16)
Here $`h_{1,1}`$ and $`h_{1,2}`$ are the Hodge numbers of the CY threefold $`X`$ on which the type IIA string has been compactified; for rigid CY’s, where $`h_{1,2}=0`$, we have the important bound $`c<0`$. The function $`h`$ in (15) is simply the general $`(u,v)`$-independent solution to the Toda equation (modulo a constant rescaling of $`r`$); in this sense the perturbative corrections appear naturally. The PT coordinate $`r`$ is related to $`\rho `$ in through $`r=\rho ^2c=\mathrm{e}^\varphi `$; the relation between the fields and PT coordinates receives no (perturbative) quantum corrections.
Note that if we consider $`c<0`$, the function $`f`$ in (15) becomes negative for $`r<2|c|`$, which results in a negative-definite metric (7). We thus have to restrict $`r`$ to the open interval $`2|c|<r<\mathrm{}`$.
## 4 Instanton corrections
In this section, we construct solutions to the Toda equation (8) that include an (infinite) series of exponential corrections describing the membrane instantons. As we have learned from the supergravity description, the real part of the instanton action is inversely proportional to the dilaton, which becomes the square root of the radial variable $`r`$. The precise form of the supergravity instanton action is given in (6). This motivates us to make a general ansatz of the form
$$\mathrm{e}^h=r+\underset{n1}{}\underset{m}{}f_{n,m}(u,v)r^{m/2+\alpha }\mathrm{e}^{2n\sqrt{r}}.$$
(17)
As explained in the previous section, one can shift the value of $`r`$ with a constant to construct a new solution. This will then include the perturbative one-loop correction of . The power series in $`r`$ in front of the exponent describes the perturbative corrections around the instantons. Using (5) we have that $`r^{m/2}=g_s^m`$, and the sum over $`m`$ is over the integers . At each instanton level $`n`$, there is a lowest value $`m_n`$ that defines the leading term in the expansion,
$$f_{n,m}(u,v)=0\text{for}m<m_n.$$
(18)
We have also introduced a parameter $`\alpha `$ which, without loss of generality, lies in the interval $`[0,1/2[`$. This leaves open the possibility that the leading term is not an integer power of $`g_s`$, as e.g. in . We will show later on that the Toda equation enforces $`\alpha =0`$. With the $`r`$-dependence made explicit, solving the Toda equation amounts to solving the differential equations for the functions $`f_{n,m}(u,v)`$. These are of the type of inhomogeneous Laplace equations, and we can solve them iteratively, order by order in $`n`$ and $`m`$, to any order needed.
To get some additional insight, we focus for a moment on the asymptotic (large $`r`$) behavior of the solution. We can then further specify the ansatz as
$$\mathrm{e}^h=r+A\mathrm{cos}(k_uu+k_vv)r^\beta \mathrm{e}^{2k\sqrt{r}},$$
(19)
with $`A`$ a normalization constant. One can now check that, to leading order, the Toda equation is satisfied for any value of $`\beta `$, provided that
$$k^2=k_u^2+k_v^2.$$
(20)
This asymptotic behavior indeed reproduces leading order charge $`k`$ instanton effects, including a one-loop correction in front of the exponent. The cosine in the ansatz (19) could also be replaced by a sine, or a linear combination. Rewriting them in terms of exponentials, one produces theta-angle like terms for both instantons and anti-instantons, depending on the signs of $`(k_u,k_v)`$. The relation (20) is completely consistent with the supergravity description of the instanton action (6), which describes the special case of either $`k_u=0`$ or $`k_v=0`$.
We now give a more complete analysis for solving the Toda equation, based on the general ansatz (17). This will enable us to determine the subleading corrections to the solution (19). More technical details are given in appendix C. For instance, in appendix C.1 we show that $`m_n2`$ for all $`n`$, and in appendix C.2 it is proven that $`\alpha =0`$.
We first bring the Toda equation into the equivalent form
$$\mathrm{e}^h\left(_u^2+_v^2+\mathrm{e}^h_r^2\right)\mathrm{e}^h(_u\mathrm{e}^h)^2(_v\mathrm{e}^h)^2=0,$$
(21)
such that it depends on $`h`$ only through $`\mathrm{e}^h`$. We then decompose this equation into $`N`$-instanton sectors, each containing a sum over all loop corrections,
$`0={\displaystyle \underset{n,m}{}}`$ $`r^{m/2}\mathrm{e}^{2n\sqrt{r}}\{(\mathrm{\Delta }+n^2)f_{n,m+2}+na_{m+2}f_{n,m+1}+b_{m+2}f_{n,m}`$
$`+{\displaystyle \underset{n^{},m^{}}{}}\mathrm{e}^{2n^{}\sqrt{r}}[2na_{m^{}+1}f_{n^{},mm^{}1}+2b_{m^{}+2}f_{n^{},mm^{}2}`$
$`+f_{n^{},mm^{}}(\mathrm{\Delta }+2n^2)f_{n^{},mm^{}}]f_{n,m^{}}`$
$`+{\displaystyle \underset{n^{},m^{}}{}}{\displaystyle \underset{n^{\prime \prime },m^{\prime \prime }}{}}\mathrm{e}^{2(n^{}+n^{\prime \prime })\sqrt{r}}f_{n,m^{}}f_{n^{},m^{\prime \prime }}[n^2f_{n^{\prime \prime },mm^{}m^{\prime \prime }2}`$
$`+na_{m^{}+1}f_{n^{\prime \prime },mm^{}m^{\prime \prime }3}+b_{m^{}+2}f_{n^{\prime \prime },mm^{}m^{\prime \prime }4}]\},`$ (22)
where $`=(_u,_v)`$, $`\mathrm{\Delta }=^2`$, and
$$a_m=\frac{1}{2}(2m1),b_m=\frac{1}{4}m(m2).$$
(23)
In the $`(N=1)`$-instanton sector only the single-sum terms contribute, while the double- and triple-sums have to be taken into account beginning with the $`(N=2)`$\- and $`(N=3)`$-instanton sectors, respectively.
### 4.1 The one-instanton sector
We start with $`N=1`$. In this sector the Toda equation requires at the $`m`$th loop order
$$(\mathrm{\Delta }+1)f_{1,m}+a_mf_{1,m1}+b_mf_{1,m2}=0,$$
(24)
It is convenient to first consider a one-dimensional truncation, where $`f_{1,m}=f_{1,m}(x)`$ with $`x\{u,v\}`$. In appendix C.3 we prove that the general one-dimensional solution of (24) is given by
$$f_{1,m}(x)=\mathrm{Re}\underset{s0}{}\frac{1}{s!(2)^s}k_{1,m}(s)G_s(x)$$
(25)
with recursively defined coefficients
$$k_{1,m}(s+1)=a_mk_{1,m1}(s)+b_mk_{1,m2}(s).$$
(26)
$`G_s(x)`$ are complex functions related to the spherical Bessel functions of the third kind; their precise definition can be found in appendix C.3. $`k_{1,m}(0)=A_{1,m}`$ are complex integration constants originating from the homogeneous part of (24). By definition of $`m_1`$ we have that $`A_{1,m}=0`$ for $`m<m_1`$, and from this it follows that $`k_{1,m}(s>mm_1)=0`$ by using (26). The highest $`x`$-monomial contained in $`f_{1,m}(x)`$ is then of order $`mm_1`$. Explicitly, the first two solutions read
$`f_{1,m_1}(x)`$ $`=\mathrm{Re}\left\{A_{1,m_1}\mathrm{e}^{\mathrm{i}x}\right\},`$
$`f_{1,m_1+1}(x)`$ $`=\mathrm{Re}\left\{A_{1,m_1+1}\mathrm{e}^{\mathrm{i}x}+\frac{1}{2}a_{m_1+1}A_{1,m_1}\mathrm{i}x\mathrm{e}^{\mathrm{i}x}\right\}.`$ (27)
We now extend the $`N=1`$ result to the general $`u,v`$ dependent solution. This can be done by Fourier transforming in the $`u,v`$ plane or, as we do below, by separation of variables. In both cases one finds a basis of solutions; the most general solution is then obtained by superposition. Using separation of variables, we find a basis and parameterize it by a continuous parameter $`\lambda `$.
Introducing $`\omega =\sqrt{1\lambda ^2}`$ with $`\lambda ^2`$ being real, the general solution can then be written as
$$f_{1,m}(u,v)=𝑑\lambda \mathrm{Re}\underset{s0}{}\frac{1}{s!(2\omega ^2)^s}k_{1,m}(s,u;\lambda )G_s(\omega v),$$
$$k_{1,m}(s+1,u;\lambda )=a_mk_{1,m1}(s,u;\lambda )+b_mk_{1,m2}(s,u;\lambda ),$$
(28)
where
$$k_{1,m}(0,u;\lambda )=B_{1,m}(\lambda )A_{1,m}(\lambda )\mathrm{e}^{\mathrm{i}\lambda u}.$$
(29)
Here $`A_{1,m}(\lambda )`$, $`B_{1,m}(\lambda )`$ are arbitrary complex integration functions which determine the “frequency spectrum” of the solution. The $`u`$-independent solution of the previous paragraph can then be obtained by setting
$$A_{1,m}(\lambda )=A_{1,m},B_{1,m}(\lambda )=\delta (\lambda ),$$
(30)
with $`A_{1,m}`$ being the corresponding integration constants.
One can now combine the general $`u`$-independent solution with the general $`v`$-independent solution. This can be done by taking the coefficient functions $`A_{1,m}(\lambda )`$ and $`B_{1,m}(\lambda )`$ to be peaked around $`\lambda =0`$ and $`\lambda =1`$. This is not the most general solution, but it is the preferred one that describes our physical problem. For general values $`0<\lambda <1`$, one generates products of exponents in $`u`$ and in $`v`$ that describe theta-angle-like terms in the supergravity instanton action where both $`\phi `$ and $`\chi `$ and their charges $`Q_\phi `$ and $`Q_\chi `$ are turned on. As we have argued at the end of section 2, this cannot be the case. Moreover, as we will see in section 5, string theory also predicts such terms to be absent. We therefore only take contributions from $`\lambda =0,1`$. This implies that $`\mathrm{e}^h`$, including the perturbative corrections and the instanton corrections arising in the one-instanton sector, can be completely expressed in terms of the one-dimensional solutions (25). Here making the substitution $`xu,v`$ describes the one-instanton contribution to $`\mathrm{e}^h`$ arising from a $`u,v`$-instanton, respectively. Also taking into account the perturbative corrections to the solution by shifting $`rr+c`$ we find
$$\mathrm{exp}[h(r,u,v)]=\mathrm{exp}[h_{\text{pert}}(r)]+\mathrm{exp}[h_{\text{1-inst}}(r,u)]+\mathrm{exp}[h_{\text{1-inst}}(r,v)]+\mathrm{},$$
(31)
where
$$\mathrm{exp}[h_{\text{pert}}(r)]=r+c,\mathrm{exp}[h_{\text{1-inst}}(r,u)]=\mathrm{e}^{2\sqrt{r+c}}\underset{mm_1}{}f_{1,m}(u)(r+c)^{m/2},$$
(32)
and similarly for $`h_{\text{1-inst}}(r,v)`$. The coefficients $`f_{1,m}(x)`$ are the one parameter solution (25) and the ellipses denote the contributions from higher order instanton corrections.
For later reference, we also give the leading order expression for $`\mathrm{e}^h`$ in the regime $`r1`$ (small string coupling). To leading order in the semi-classical approximation, the instanton solution (31) reads
$$\mathrm{e}^h=r+c+\frac{1}{2}r^{m_1/2}\left(A_{1,m_1}\mathrm{e}^{\mathrm{i}v}+A_{1,m_1}^{}\mathrm{e}^{\mathrm{i}v}+B_{1,m_1}\mathrm{e}^{\mathrm{i}u}+B_{1,m_1}^{}\mathrm{e}^{\mathrm{i}u}\right)\mathrm{e}^{2\sqrt{r}}+\mathrm{}.$$
(33)
Notice that we need to include both instantons and anti-instantons to obtain a real solution.
To find the leading-order instanton corrected hypermultiplet metric, we first compute the leading corrections to $`f`$ defined in $`(\text{9})`$:
$$f=\frac{r+2c}{r+c}+\frac{1}{2}r^{(m_1+1)/2}\left(A_{1,m_1}\mathrm{e}^{\mathrm{i}v}+A_{1,m_1}^{}\mathrm{e}^{\mathrm{i}v}+B_{1,m_1}\mathrm{e}^{\mathrm{i}u}+B_{1,m_1}^{}\mathrm{e}^{\mathrm{i}u}\right)\mathrm{e}^{2\sqrt{r}}+\mathrm{}.$$
(34)
Substituting this result into (10), one derives the leading corrections to the $`\mathrm{\Theta }`$ 1-form. Setting
$$\mathrm{\Theta }=u\mathrm{d}v+\mathrm{\Theta }_{\text{inst}},$$
(35)
these are given by
$$\mathrm{\Theta }_{\text{inst}}=r^{m_1/2}\mathrm{e}^{2\sqrt{r}}\mathrm{Im}\{A_1\mathrm{e}^{\mathrm{i}v}\mathrm{d}u+B_1\mathrm{e}^{\mathrm{i}u}\mathrm{d}v\}+\mathrm{}.$$
(36)
The leading order corrections to the hypermultiplet scalar metric are then obtained by plugging these expressions into the PT metric (7).
### 4.2 Higher instanton sectors
We now briefly discuss the $`N=2`$ sector. The Toda equation requires at this level
$`0`$ $`=(\mathrm{\Delta }+4)f_{2,m}+2a_mf_{2,m1}+b_mf_{2,m2}`$
$`+{\displaystyle \underset{m^{}}{}}[f_{1,mm^{}2}+a_{m^{}+1}f_{1,mm^{}3}+b_{m^{}+2}f_{1,mm^{}4}`$
$`f_{1,mm^{}2}]f_{1,m^{}},`$ (37)
where we have used (24) for $`\mathrm{\Delta }f_{1,m^{}}`$ in the double sum. We have not derived the general solution to these equations in closed form; the one-dimensional truncation, however, is straightforward to solve order by order in $`m`$. At lowest order<sup>8</sup><sup>8</sup>8In appendix C.2 we show that for $`nn^{}`$ it is $`2m_nm_n^{}`$. $`m_2`$ we have
$$(\mathrm{\Delta }+4)f_{2,m_2}+\delta _{m_2,2}\left[(f_{1,m_2})^2(f_{1,m_2})^2\right]=0.$$
(38)
Note that the inhomogeneous term is present only for the lowest possible value $`m_2=2`$. The one-dimensional truncation yields the equation
$$(_x^2+4)f_{2,m_2}(x)+\delta _{m_2,2}\mathrm{Re}\left\{A_{1,m_2}^2\mathrm{e}^{2\mathrm{i}x}\right\}=0,$$
(39)
where we have inserted the solution (27) for $`f_{1,m_1}(x)`$. The general solution then reads
$`f_{2,m_2}(x)`$ $`=\mathrm{Re}\left\{A_{2,m_2}\mathrm{e}^{2\mathrm{i}x}+\frac{1}{4}\delta _{m_2,2}A_{1,m_2}^2\mathrm{i}x\mathrm{e}^{2\mathrm{i}x}\right\}`$
$`=\mathrm{Re}\left\{A_{2,m_2}G_0(2x)\frac{1}{8}\delta _{m_2,2}A_{1,m_2}^2G_1(2x)\right\},`$ (40)
$`A_{2,m_2}`$ being a further complex integration constant.
The solution for $`m>m_2`$ can now be constructed by solving the appropriate equation arising from (4.2). Based on (28) we can also construct the general $`(u,v)`$-dependent solution for $`f_{2,m_2}(u,v)`$. The idea is to decompose the products of $`\mathrm{cos}(\lambda _1u)\mathrm{cos}(\lambda _2u)`$, etc., appearing in the inhomogeneous part of (38) into a sum of $`\mathrm{cos}`$ and $`\mathrm{sin}`$ terms using product formulae for two trigonometric functions. We can then construct the full inhomogeneous solution by superposing the inhomogeneous solutions for every term in the sum. We refrain from giving the result, however, since it is complicated and not particularly illuminating.
We conclude this subsection by giving an argument that the iterative solution devised above indeed gives rise to a consistent solution of the Toda equation. The general equations which determine a new $`f_{n,m}(u,v)`$ are two-dimensional Laplace equations to the eigenvalue $`n^2`$ coupled to an inhomogeneous term, which is completely determined by the $`f_{n,m}(u,v)`$’s obtained in the previous steps of the iteration procedure. These equations are readily solved, e.g., by applying a Fourier transformation. It then turns out that the iteration procedure is organized in such a way that any level in the perturbative expansion (4) determines one “new” $`f_{n,m}(u,v)`$, i.e., there are no further constraints on the $`f_{n,m}(u,v)`$ determined in the previous steps. This establishes that our perturbative approach indeed extends to a consistent solution of the Toda equation (8).
### 4.3 The fate of the Heisenberg algebra
Based on the Toda solution (31) we now discuss the breaking of the Heisenberg algebra (4) in the presence of membrane instantons. We start with the shift symmetry in the axion $`\sigma \sigma \alpha `$. By identifying $`t=\sigma `$, this shift corresponds to the isometry of the Tod metric, so that it cannot be broken by the instanton corrections.
Analyzing the $`\beta `$ and $`\gamma `$-shifts is more involved. Under the identification (14), the $`\beta `$-shift then acts as $`vv+\beta `$. Taking the leading order one-instanton solution (33)-(36), we find that $`\mathrm{e}^h`$ as well as the resulting functions $`f`$ and $`\mathrm{\Theta }`$ appearing in the metric depend on $`v`$ through $`\mathrm{e}^{\pm \mathrm{i}v}`$ or $`\mathrm{d}v`$ only. These theta-angle-like terms break the $`\beta `$-shift to the discrete symmetry group .<sup>9</sup><sup>9</sup>9This agrees with earlier observations made in . Going beyond the leading instanton corrections by taking into account higher loop corrections around the single instanton will, however, generically break the $`\beta `$-shift completely, due to the appearance of polynomials in $`v`$ multiplying the factors $`\mathrm{e}^{\pm \mathrm{i}v}`$. We point out, however, that by setting the integration constants multiplying the terms odd in $`v`$ to zero, there is still an unbroken <sub>2</sub> symmetry defined by $`vv`$, $`tt`$, interchanging $`v`$-instantons and anti-instantons.
To deduce the fate of the $`\gamma `$-shift, $`uu+\gamma `$, $`tt\gamma v`$, we first observe that $`tt\gamma v`$ implies that the combination $`\mathrm{d}t+u\mathrm{dv}`$ is invariant. Applying the same logic as for the $`\beta `$-shift above, we then find that the one-loop corrections of a single $`u`$-instanton break the $`\gamma `$-shift to the discrete symmetry , which will be generically broken by higher order terms appearing in the loop expansion. Similar to the $`\beta `$-shift, however, we can arrange the constants of integration appearing in the solution in such a way that there is also a <sub>2</sub> symmetry. We expect that these two <sub>2</sub> symmetries could play a prominent role when determining (some of) the coefficients appearing in the solution (31) from string theory.
## 5 Comparison to string theory
As mentioned above, instanton corrections to the moduli space metric also induce corrections to the 4-fermion couplings in the supergravity effective action, since they couple to the curvature of the moduli space metric. It is therefore desirable to have a microscopic string theory derivation that reproduces these instanton corrections. Using the work of Becker, Becker and Strominger (BBS) , this is possible, and we show in this section that there is a perfect agreement with string theory.
The reason why 4-fermi terms are the relevant objects to look at is that our membrane instantons break half of the supersymmetries. The resulting four fermionic zero modes then lead to non-vanishing 4-fermion correlation functions. This was already observed by BBS in a string theoretic setting. Here we compare our supergravity result for the instanton corrected 4-hyperino couplings with those derived in . Their analysis was actually set up by starting with CY<sub>3</sub> compactifications of M-theory, and then reducing to type IIA in ten dimensions. As we also explain below, the only modifications are in the appearance of the string coupling constant. This is also consistent with the supergravity analysis, since the hypermultiplet couplings to supergravity are (almost) identical in four and five space-time dimensions.
### 5.1 The string calculation
The relevant curvature tensor that is contracted with the 4-fermi terms is the totally symmetric $`𝒲_{\alpha \beta \gamma \delta }`$ tensor introduced in section 3.2 (see also appendix A.3). In the BBS paper, this tensor was denoted by $`_{IJKL}`$. For compactifications on rigid CY<sub>3</sub> yielding one (the universal) hypermultiplet, we expect that they agree up to normalization (which, to our knowledge, has not been computed in a string theory setting) and an $`\mathrm{USp}(2)\mathrm{SU}(2)`$ rotation of the fermion frame.
Instanton configurations are obtained by wrapping Euclidean membranes over a supersymmetric three-cycle $`𝒞_3`$. The effect of such an instanton is to yield a non-vanishing 4-fermi correlator that gives a contribution to the curvature tensor. In the M-theory set-up, this was found to be (see eq. (2.49) in ),
$$\mathrm{\Delta }_{𝒞_3}_{IJKL}=N^{}\mathrm{e}^{S_{\mathrm{inst}}}_{𝒞_3}d_I_{𝒞_3}d_J_{𝒞_3}d_K_{𝒞_3}d_L.$$
(41)
Here, $`N^{}`$ is an unspecified normalization factor which, in principle, could depend on the string coupling $`g_s`$. Furthermore,
$$S_{\mathrm{inst}}=\mathrm{e}^𝒦\left|_{𝒞_3}\mathrm{\Omega }\right|+\mathrm{i}_{𝒞_3}C_3$$
(42)
is the bosonic part of the instanton action, $`𝒞_3`$ denotes the supersymmetric cycle that is wrapped by the membrane, $`C_3`$ is the 3-form potential in $`D=11`$ supergravity, and the $`d_I`$ form a real symplectic basis of $`H^3(X,)`$. Finally, we have $`𝒦=1/2(𝒦_V𝒦_H)`$ with
$$𝒦_V=\mathrm{log}\left(\frac{4}{3}_X\widehat{J}\widehat{J}\widehat{J}\right),𝒦_H=\mathrm{log}\left(\mathrm{i}_X\mathrm{\Omega }\overline{\mathrm{\Omega }}\right),$$
(43)
where $`\widehat{J}`$ is the Kähler form and $`\mathrm{\Omega }`$ the holomorphic 3-form on the CY threefold $`X`$.
All these quantities can be expressed in terms of our variables for the universal hypermultiplet. For this we need the relation between the Tod variables $`(r,u,v,t)`$ and the fields appearing in the IIA superstring action of BBS . To establish these relations the references are useful.
In order to compactify the IIA string on a CY<sub>3</sub> manifold $`X`$, we introduce $`2(h_{1,2}+1)`$ harmonic 3-forms $`(\alpha _a,\beta ^a)`$, which form a real basis of $`H^3(X,)`$, with the usual normalization
$$_X\alpha _a\beta ^b=_X\beta ^b\alpha _a=\delta _a^b,_X\alpha _a\alpha _b=_X\beta ^a\beta ^b=0.$$
(44)
They correspond to the $`d^I`$ in (41). Furthermore, we introduce the canonical dual basis of real 3-cycles $`(𝒜^a,_a)`$ of $`H_3(X,)`$, satisfying <sup>10</sup><sup>10</sup>10In these relations, we have chosen a normalization in which the volume of the CY<sub>3</sub> is set to one.
$$_{𝒜^a}\alpha _b=__b\beta ^a=\delta _b^a,_{𝒜^a}\beta ^b=__a\alpha _b=0.$$
(45)
For rigid CYs, the index $`a`$ takes only the value $`0`$ and may be omitted. We can then use this basis to define the periods of the holomorphic 3-form $`\mathrm{\Omega }`$ of the CY<sub>3</sub> as
$$z^a=_{𝒜^a}\mathrm{\Omega },𝒢_a=__a\mathrm{\Omega },$$
(46)
in terms of which
$$\mathrm{\Omega }=z^a\alpha _a𝒢_a(z)\beta ^a.$$
(47)
Here $`z^a`$ are the complex structure moduli, and $`𝒢_a(z)`$ are derivatives of the prepotential of the special geometry which is parameterized by the $`z^a`$. The Kähler potential on the space of complex structure deformations is then given by
$$𝒦_H=\mathrm{log}\left(2\mathrm{Im}(\overline{z}^a𝒢_a)\right).$$
(48)
In the case of a rigid CY<sub>3</sub> we only have $`z^0`$ and $`𝒢_0`$. We can then choose the normalization of $`\mathrm{\Omega }`$ (which is defined up to a complex rescaling only) such that
$$z^0=1,𝒢_0=\mathrm{i}.$$
(49)
The phase of $`𝒢_0`$ is determined in such a way that $`𝒦_H`$ is real.
With these prerequisites it is now possible to determine the supersymmetric cycles of the CY<sub>3</sub>. In fact, we shall find that these are given by the cycles $`𝒜^a`$, $`_a`$ themselves. In it was shown that a supersymmetric cycle has to satisfy the following two (equivalent) conditions:
1. The pull-back of the embedding space’s Kähler form has to vanish.
2. The cycle has to be calibrated with respect to the holomorphic 3-form $`\mathrm{\Omega }`$, i.e., the volume form of the cycle has to be proportional to the pull-back of $`\mathrm{\Omega }`$ up to a complex phase factor. (This implies that a supersymmetric cycle has to be a special Lagrangian submanifold of $`X`$ .)
In order to show that the cycles $`𝒜^a`$, $`_a`$ satisfy the first condition, we can generalize the argument given in the example of , section 2.2. There, an isometry $`D`$ of the metric was employed which corresponds to complex conjugation. The Kähler metric on $`X`$ is invariant under $`D`$, as are the real cycles $`𝒜^a`$, $`_a`$, whereas the Kähler form associated with the Kähler metric reverses its sign,
$$D:\widehat{J}\widehat{J}.$$
(50)
On the other hand, the pullback of $`\widehat{J}`$ onto $`𝒜^a`$ or $`_a`$, respectively, must be invariant, which is only possible if $`\widehat{J}`$ vanishes on $`𝒜^a`$ or $`_a`$.
The second condition is satisfied in rigid CY<sub>3</sub> compactifications since $`\alpha `$ and $`\beta `$ correspond to the induced volume forms on $`𝒜`$ and $``$, respectively.
We can now use the dimensional reduction outlined in <sup>11</sup><sup>11</sup>11The various fields in , and our paper, respectively, differ in their normalizations. The UHM variables in are related to ours through $`\widehat{\varphi }=\varphi /2`$, $`\widehat{\xi }=\phi /\sqrt{2}`$, $`\widehat{\overline{\xi }}=\chi /\sqrt{2}`$, $`\widehat{a}=\sigma +\frac{1}{2}\chi \phi `$. The RR fields in are obtained from those in by multiplication with $`\sqrt{2}`$. to evaluate the integrals appearing in (41). We find that $`C_3`$ in is expanded as
$$C_3=c_3+v\alpha +u\beta +\mathrm{},$$
(51)
where $`c_3`$ is the space-time 3-form potential (which is non-dynamical in four dimensions) and the ellipses denote the omitted vector multiplet contributions.
To evaluate the Kähler potential appearing in (42), we note that the volume of the CY<sub>3</sub> manifold measured with the 11-dimensional supergravity metric is given by
$$\widehat{V}_6=\frac{1}{3!}_X\widehat{J}\widehat{J}\widehat{J}.$$
(52)
Upon reducing to the IIA supergravity action in the string frame using
$$\mathrm{d}\widehat{s}_{11}^2=\mathrm{e}^{2\varphi /3}\left(\mathrm{d}x_{11}+A_m\mathrm{d}x^m\right)^2+\mathrm{e}^{\varphi /3}\mathrm{d}s_{10}^2,$$
(53)
$`𝒦`$ acquires a non-trivial dependence on the dilaton. Since
$$\widehat{J}=\mathrm{i}\widehat{g}=\mathrm{e}^{\varphi /3}\mathrm{i}g=\mathrm{e}^{\varphi /3}J,$$
(54)
where $`J`$ is the Kähler form in the string frame, we find the relation $`\widehat{V}_6=\mathrm{e}^\varphi V_6`$. In our normalization (44),
$$\mathrm{i}_X\mathrm{\Omega }\overline{\mathrm{\Omega }}=2V_6,$$
(55)
we can evaluate the Kähler potential in (42):
$`\mathrm{exp}(𝒦)`$ $`=\mathrm{exp}\left({\displaystyle \frac{1}{2}}(𝒦_V𝒦_H)\right)`$
$`=\mathrm{exp}\left({\displaystyle \frac{1}{2}}\left(\mathrm{log}(8\mathrm{e}^\varphi V_6)+\mathrm{log}(2V_6)\right)\right)`$
$`=2\mathrm{e}^{\varphi /2}.`$ (56)
It is now also straightforward to evaluate the remaining integral
$$\left|_{𝒞_3}\mathrm{\Omega }\right|=\sqrt{m^2+n^2},$$
(57)
for $`𝒞_3=m𝒜+n`$. Notice, however, that under the condition that $`𝒜`$ and $``$ are calibrated one can show that the linear combination $`𝒞_3=m𝒜+n`$ is a calibrated cycle if and only if either $`m`$, $`n=0`$ or $`m=0`$, $`n`$. Therefore a membrane wrapping $`𝒜`$ and $``$ simultaneously is not a supersymmetric configuration and does not contribute to the instanton corrected metric. This is also reflected in the form of $`S_{\text{inst}}`$ given in , where instanton charges are linear: $`n+m`$, where $`m`$ and $`n`$ cannot be non-zero simultaneously.
Putting everything together, using $`r=\mathrm{e}^\varphi `$, we then obtain the instanton weight of a configuration where $`𝒞_3=m𝒜+n`$:
$$\mathrm{e}^{S_{\mathrm{inst}}}=\mathrm{e}^{2\sqrt{m^2+n^2}\sqrt{r}}\mathrm{e}^{\mathrm{i}mv+\mathrm{i}nu}.$$
(58)
In the framework of a rigid CY<sub>3</sub> compactification, the $`d^I`$, $`I=1,2`$, correspond to the harmonic three-forms $`\alpha `$, $`\beta `$, while the wrapped cycle $`𝒞_3`$ can be either the cycle $`𝒜`$ or $``$ introduced above. Using the relations (45) it is then straightforward to verify that the instanton corrections predicted by (41) enter into the components of $`\mathrm{\Delta }_{𝒞_3}_{IJKL}`$ with $`I=J=K=L=1,2`$ only.
As we will show in the next subsection, at leading order in $`m`$, $`n`$ this agrees with the results obtained by substituting our instanton corrected UHM into $`W_{\alpha \beta \gamma \delta }`$ and also fixes the coefficient functions in our solution in terms of free parameters.
### 5.2 Comparison with the instanton corrected PT metric
Evaluating $`𝒲_{\alpha \beta \gamma \delta }`$ for the universal hypermultiplet in the membrane base, we find that at the perturbative level (classical plus loop corrections) the only non-vanishing component is given by
$$𝒲_{2211}=\frac{r^3}{(r+2c)^3}.$$
(59)
We now substitute the instanton expansion (33) in the general $`𝒲`$-tensor given in (3.2). After subtracting the classical contribution, we obtain to lowest order in $`g_s`$
$`\mathrm{\Delta }𝒲_{1111}`$ $`=N\left[A_1\mathrm{e}^{\mathrm{i}v}+A_1^{}\mathrm{e}^{\mathrm{i}v}B_1\mathrm{e}^{\mathrm{i}u}B_1^{}\mathrm{e}^{\mathrm{i}u}\right]`$
$`\mathrm{\Delta }𝒲_{1112}`$ $`=N\left[A_1\mathrm{e}^{\mathrm{i}v}A_1^{}\mathrm{e}^{\mathrm{i}v}+\mathrm{i}(B_1\mathrm{e}^{\mathrm{i}u}B_1^{}\mathrm{e}^{\mathrm{i}u})\right]`$
$`\mathrm{\Delta }𝒲_{1122}`$ $`=N\left[A_1\mathrm{e}^{\mathrm{i}v}+A_1^{}\mathrm{e}^{\mathrm{i}v}+B_1\mathrm{e}^{\mathrm{i}u}+B_1^{}\mathrm{e}^{\mathrm{i}u}\right]`$
$`\mathrm{\Delta }𝒲_{1222}`$ $`=N\left[A_1\mathrm{e}^{\mathrm{i}v}A_1^{}\mathrm{e}^{\mathrm{i}v}\mathrm{i}(B_1\mathrm{e}^{\mathrm{i}u}B_1^{}\mathrm{e}^{\mathrm{i}u})\right]`$
$`\mathrm{\Delta }𝒲_{2222}`$ $`=N\left[A_1\mathrm{e}^{\mathrm{i}v}+A_1^{}\mathrm{e}^{\mathrm{i}v}B_1\mathrm{e}^{\mathrm{i}u}B_1^{}\mathrm{e}^{\mathrm{i}u}\right].`$ (60)
Here we have set $`N=(r+c)^{(1m_1)/2}\mathrm{e}^{2\sqrt{r+c}}`$. Comparing the $`r`$-dependence of $`N`$ with the one appearing in the normalization factor $`N^{}`$ then fixes the value of $`m_1`$. In particular, for $`N^{}`$ being an $`r`$-independent normalization constant, we obtain $`m_1=1`$.
At first sight the tensorial structure of (5.2) seems to disagree with (41), as the latter predicts instanton corrections to the components $`\mathrm{\Delta }𝒲_{1111}`$ and $`\mathrm{\Delta }𝒲_{2222}`$ only. In order to match these two results it is crucial to observe that (41) describes the corrections arising from a single membrane wrapping *one* supersymmetric cycle, while our expression already contains the “instanton sum” over the $`𝒜`$ and $``$ cycle. Furthermore, the fermion frame used by BBS does not necessarily agree with ours. However, these two frames can differ at most by a local SU(2) rotation<sup>12</sup><sup>12</sup>12The rotation group has to be compatible with the reality condition imposed on the pair of symplectic Majorana spinors coupling to $`𝒲`$ and preserve fermion bilinears. These conditions then lead to the fact that the most general transformation is given by SU(2). This is discussed in .. We parameterize this transformation by
$$U=\left(\begin{array}{cc}\mathrm{e}^{\mathrm{i}\xi }\mathrm{cos}\eta & \mathrm{e}^{\mathrm{i}\rho }\mathrm{sin}\eta \\ \mathrm{e}^{\mathrm{i}\rho }\mathrm{sin}\eta & \mathrm{e}^{\mathrm{i}\xi }\mathrm{cos}\eta \end{array}\right),$$
(61)
where the parameters $`\eta `$, $`\xi `$, $`\rho `$ can in principle depend on the scalars (this, however, will not be necessary). To follow BBS, we then consider the contribution arising from the $``$-instanton only. This requires setting $`A_1=0`$. Upon performing a global $`\mathrm{SU}(2)`$ rotation of the fermion frame with parameters $`\eta =\pi /4`$, $`\xi =\rho +\pi /2`$, we obtain
$$\mathrm{\Delta }\stackrel{~}{𝒲}_{1111}=4NB_1^{}\mathrm{e}^{\mathrm{i}u},\mathrm{\Delta }\stackrel{~}{𝒲}_{2222}=4NB_1\mathrm{e}^{\mathrm{i}u},$$
(62)
with the other components vanishing identically. (Note that the remaining free parameter in the transformation $`\rho `$ only induces a phase on the components of $`\mathrm{\Delta }\stackrel{~}{𝒲}_{\alpha \beta \gamma \delta }`$, which we set to zero for convenience.) Observe that this now matches the prediction of BBS. Likewise, we can consider the contribution of the $`𝒜`$-instanton by setting $`B_1=0`$. In this case the transformation $`\eta =\pi /4`$, $`\xi =\rho =0`$ leads to a correction
$$\mathrm{\Delta }\stackrel{~}{𝒲}_{1111}=4NA_1\mathrm{e}^{\mathrm{i}v},\mathrm{\Delta }\stackrel{~}{𝒲}_{2222}=4NA_1^{}\mathrm{e}^{\mathrm{i}v},$$
(63)
with the other components vanishing identically. This is again of the form predicted by BBS, even though in a different fermionic frame than (62). Summing these two corrections involves rotating some of the contributions into the proper fermionic frame. Hence, our result precisely agrees with the one obtained in .
In order to sum the two contributions, we go to the $``$-instanton frame, i.e. the frame in which we obtained (62), but now we also include $`A_1`$. The corrections to the $`𝒲`$-Tensor then are
$`\mathrm{\Delta }\stackrel{~}{𝒲}_{1111}`$ $`=N\left(A_1\mathrm{e}^{\mathrm{i}v}+A_1^{}\mathrm{e}^{\mathrm{i}v}+4B_1^{}\mathrm{e}^{\mathrm{i}u}\right)`$
$`\mathrm{\Delta }\stackrel{~}{𝒲}_{1112}`$ $`=N\left(A_1\mathrm{e}^{\mathrm{i}v}A_1^{}\mathrm{e}^{\mathrm{i}v}\right)`$
$`\mathrm{\Delta }\stackrel{~}{𝒲}_{1122}`$ $`=N\left(A_1\mathrm{e}^{\mathrm{i}v}+A_1^{}\mathrm{e}^{\mathrm{i}v}\right)`$
$`\mathrm{\Delta }\stackrel{~}{𝒲}_{1222}`$ $`=N\left(A_1\mathrm{e}^{\mathrm{i}v}A_1^{}\mathrm{e}^{\mathrm{i}v}\right)`$
$`\mathrm{\Delta }\stackrel{~}{𝒲}_{2222}`$ $`=N\left(A_1\mathrm{e}^{\mathrm{i}v}+A_1^{}\mathrm{e}^{\mathrm{i}v}+4B_1\mathrm{e}^{\mathrm{i}u}\right).`$ (64)
This result implies that the four fermionic zero modes $`\psi _1`$, $`\psi _2`$ arising from a membrane wrapping the $`𝒜`$ and the $``$ cycle, respectively, are not orthogonal. If in the $``$-frame we denote the two zero modes giving rise to the $`\mathrm{e}^{\mathrm{i}u}`$ corrections with $`\psi _1`$, then the corrections proportional to $`\mathrm{e}^{\mathrm{i}v}`$ arise from the zero modes $`\psi _1\psi _2`$. This zero mode configuration produces all the signs appearing in (5.2), since the $`𝒜`$-anti-instanton has its zero modes in $`\psi _1+\psi _2`$.
## 6 Constructing meta-stable de Sitter vacua
We now move on and study the properties of the scalar potential that arises from gauging the isometry of the Przanowski-Tod metric. We will find that inclusion of the instanton corrections obtained in section 4 will lead to the stabilization of all hypermultiplet moduli and also opens up the possibility for obtaining meta-stable dS vacua from string theory. In order to highlight the role of the membrane instanton corrections played in this construction, we will work with a toy model of a rigid CY<sub>3</sub> compactification where we have truncated all vector multiplets.<sup>13</sup><sup>13</sup>13Based on the results obtained in this section it is straightforward to adapt this model to any rigid CY<sub>3</sub> compactification by including the corresponding vector multiplet sector. This also allows to study more general gaugings by turning on background fluxes in the even cohomology classes of the CY<sub>3</sub>. The only remnant of the vector multiplet sector will be that we allow for a non-trivial (negative) one loop correction encoded by $`c`$.
### 6.1 Turning on background fluxes
In $`N=2`$ supergravity the scalar potential arises from gauging isometries of the scalar manifolds. In particular, the potential is completely fixed once these manifolds are chosen and the gauged isometries are specified.<sup>14</sup><sup>14</sup>14This is different from $`N=1`$ supergravity where the potential depends on an arbitrary holomorphic function, the superpotential. For the instanton corrected UHM metric the Heisenberg algebra of isometries present at the perturbative level is broken explicitly and only the shift symmetry in $`t`$ remains. The identification $`t=\sigma `$ (see discussion around (4)) shows that gauging this isometry corresponds to gauging the shift symmetry in the axion.
It is now natural to ask whether gauging this isometry has an interpretation in terms of the 10-dimensional CY<sub>3</sub> compactification. Indeed, comparing our gauging with the ones arising from flux compactifications of the IIA string , we find that it arises from a non-trivial space-filling 3-form part $`c_3`$ of the RR 3-form $`C_3`$. As pointed out in , this is equivalent to having non-trivial 6-form flux $`F_6`$ in the CY<sub>3</sub>. In four dimensions $`c_3`$ can be dualized to a constant $`e_0`$, which precisely leads to gauging the shift in $`t`$; the Killing vector encoding the gauging is then given by (11). In the absence of vector multiplets gauging this isometry induces the scalar potential (see appendix A for our normalization and conventions)
$$V=4G_{AB}k^Ak^B3\stackrel{}{P}\stackrel{}{P}=\frac{1}{r^2}\left(4f^13\right)e_0^2.$$
(65)
Here we have substituted the Killing vector (11) together with the corresponding moment map (12) and the PT metric (7) in the second step.
Before discussing the properties of this potential, let us make some remarks about the gauging. First, one might worry that the inclusion of $`F_6`$ background flux could induce tadpoles, which would render our model inconsistent. However, it was shown in that in a type IIA compactification only turning on $`F_0`$ and $`H_3`$ flux simultaneously gives rise to a tadpole condition, so that is not an issue here. Second, as discussed in for the type IIB case, including background fluxes can change the number of zero-modes that arise from a $`p`$-brane instanton. But since $`c_3`$ has no support on the wrapped cycle, we do not expect this to happen in our model, so that the membrane instantons will still lead to a correction of the four-fermi coupling. Finally, including background fluxes in a compactification in general leads to a backreaction on the geometry of the internal manifold. These backreactions are particularly relevant when looking for flux compactifications which preserve (some) supersymmetry and, at the same time, are consistent solutions of the 10-dimensional equations of motion. In the case of CY<sub>3</sub> compactifications with non-trivial fluxes this implies that the internal manifold should be a generalized CY<sub>3</sub> having an SU(3)-structure (see and references therein). We will here neglect this backreaction of the flux on the geometry in the following and tacitly assume that turning on fluxes will not drastically alter the instanton results derived for a rigid CY<sub>3</sub> manifold in the previous sections.
### 6.2 The perturbative potential
We now compute the scalar potential (65) for the one-loop metric (15). Setting $`e_0=1`$ (which does not affect the vacuum structure of the potential) we find
$$V_{\mathrm{pert}}=V_{\mathrm{class}}+V_{\text{loop}},$$
(66)
where
$$V_{\text{class}}=\frac{1}{r^2},V_{\text{loop}}=\frac{4c}{r^2(r+2c)}.$$
(67)
Fig. 1 displays $`V_{\mathrm{class}}`$ and $`V_{\mathrm{pert}}`$, respectively, for a “typical” value $`c=10`$.
The classical potential shows a typical runaway behavior in $`r`$. It is positive definite and diverges as $`r0`$. For increasing $`r`$, $`V_{\mathrm{class}}`$ decreases monotonically and there are no vacua except for the trivial one at $`r=\mathrm{}`$ (vanishing string coupling). This is shown in the left diagram of fig. 1.
Let us now add the $`V_{\text{loop}}`$-term to the scalar potential. The sign of this contribution crucially depends on the sign of $`c`$, or equivalently, on the Euler number of the Calabi-Yau.<sup>15</sup><sup>15</sup>15For $`c>0`$, $`V_{\mathrm{pert}}`$ goes to $`\mathrm{}`$ as $`r0`$. It then increases monotonically up to $`r=(1+\sqrt{5})c`$, where it has an unstable extremum, $`V_{\mathrm{pert}}|_{r=(1+\sqrt{5})c}>0`$. For $`r>(1+\sqrt{5})c`$ the potential decreases monotonically and approaches $`V_{\mathrm{pert}}0`$ for $`r\mathrm{}`$. The generic behavior of $`V_{\mathrm{pert}}`$ for $`c<0`$ is shown in the left diagram of fig. 1. Also in this case $`V_{\mathrm{pert}}=\mathrm{}`$ as $`r0`$. In the interval $`0<r<(1\sqrt{5})c`$ the potential increases monotonically. At $`r=(1\sqrt{5})c`$ we again find an unstable extremum $`V_{\mathrm{pert}}|_{r=(1\sqrt{5})c}=(\sqrt{5}+1)/(c^2(1+\sqrt{5})^2(3\sqrt{5}))<0`$. For $`(1\sqrt{5})c<r<2c`$ the potential decreases monotonically and we obtain a second singularity at $`r=2c`$. For $`r>2c`$, $`V_{\mathrm{pert}}`$ displays the runaway behavior already found in the classical case. Notice, however, that in the region $`0<r<2c`$ the perturbatively corrected metric (7) is no longer positive definite, so that this region does not belong to the moduli space of the universal hypermultiplet.
It is important to note that the perturbative potential is independent of $`(u,v,t)`$, so that these scalars correspond to flat directions. The status of the flat directions corresponding to $`(u,v)`$ and $`t`$, respectively, is quite different, however. This is due to the fact that we have gauged the shift symmetry in $`t`$ (the axion). Gauge invariance then requires that $`t`$ parameterizes a flat direction, which in turn implies that one can gauge away the scalar $`t`$, giving a mass to the gauge field, i.e., the vector field becomes massive by “eating” a scalar via the Stückelberg mechanism. Therefore, only $`u`$ and $`v`$ have to be stabilized by the potential in order to fix all moduli. As we will now show, this is readily achieved by including the leading membrane instanton corrections in the scalar potential.
### 6.3 The membrane-instanton contribution
We now demonstrate how the leading instanton correction can drastically alter the vacuum structure of our low-energy effective action. To illustrate this, let us consider the modifications arising from the $`v`$-instanton sector only, while the (equally important) terms coming from the $`u`$-instanton will be switched off for the sake of clarity. Eq. (31) indicates that the contributions stemming from the one $`(u,v)`$-instanton sector enter in exactly the same way. Hence, the corrections arising from the one $`u`$-instanton can be included by taking the $`v`$-dependent expressions given below, replacing $`v`$ by $`u`$ and adding these additional terms to the potential. Therefore, it is clear that our discussion for the $`v`$-modulus also applies to $`u`$. In particular, the stabilization of the $`v`$-modulus can trivially be extended to $`u`$ by including the $`u`$-instanton corrections as well. Furthermore, the existence of a meta-stable dS vacuum is not limited to the $`u`$-independent case and can also be obtained by including the $`u`$-dependent terms in the potential. This will shift the boundaries for the “dS window” discussed below to lower values of the integration constants.
Let us now compute the leading contribution of a single $`v`$-instanton to the scalar potential.
Substituting $`f`$ given in (34) into the potential (65), we find at leading order
$$V_{\text{1-inst}}=4r^{(m_1+5)/2}\left(\widehat{A}_{1,m_1}\mathrm{cos}(v)\stackrel{~}{A}_{1,m_1}\mathrm{sin}(v)\right)\mathrm{e}^{2\sqrt{r}}.$$
(68)
Here we have set $`A_{1,m_1}=\widehat{A}_{1,m_1}+\mathrm{i}\stackrel{~}{A}_{1,m_1}`$ and $`B_{1,m_1}=0`$. Adding this contribution to the perturbative potential (66), we then obtain in the semi-classical approximation
$$V_{\text{tot}}=V_{\text{class}}+V_{\text{loop}}+V_{\text{1-inst}}.$$
(69)
The most important change arising from including $`V_{\text{1-inst}}`$ in the potential is that the potential is no longer independent of $`v`$ (and, when including the $`u`$-instanton contribution, also of $`u`$). Therefore the *instanton correction lifts the $`u,v`$-degeneracy and provides a non-perturbative mechanism to stabilize these moduli.*
Based on (69) we can make the following additional observations: For $`r\mathrm{}`$ (vanishing $`g_s`$) all terms in $`V_{\text{tot}}`$ vanish, $`lim_r\mathrm{}V_{\text{tot}}=0`$. For $`r1`$, $`V_{\text{tot}}`$ is dominated by its classical piece $`V_{\mathrm{class}}0`$, so that $`V_{\text{tot}}`$ approaches zero from above. Furthermore, $`V_{\text{1-inst}}`$ has no poles except at $`r=0`$.<sup>16</sup><sup>16</sup>16In fact, this is an artefact of the expansion in (68). If we do not expand the denominator containing $`(r+2c)`$, $`V_{\text{1-inst}}`$ also develops a singularity at $`r=2c`$, which even dominates over the one in $`V_{\text{loop}}`$, and the potential is no longer bounded from below. Resolving this singularity presumably requires resumming the entire instanton expansion to obtain expressions which are valid at small values of $`r2c`$. This resummation is, however, beyond the scope of the present paper, and we will continue to work with the expanded expressions (68). Notice, however, that resolving singularities by non-perturbative effects has been shown to work in the context of the Coulomb branch of three-dimensional gauge theories with eight supercharges . In these cases the moduli space is hyperkähler instead of quaternion-Kähler. Hence, for $`r>0`$ the only divergence in $`V_{\text{tot}}`$ is contained in $`V_{\text{loop}}`$, which diverges at $`r=2c`$. As a result, $`V_{\text{tot}}`$ is bounded from below and diverges, $`V_{\text{tot}}=+\mathrm{}`$, as $`r2c`$.
Analyzing the vacuum structure arising from $`V_{\text{tot}}`$ analytically is rather difficult due to the transcendental nature of the potential. We therefore have analyzed the vacuum structure using numerical methods. Since we lack any better knowledge about the $`g_s`$-dependence of the instanton measure plus 1-loop determinant around the single $`v`$-instanton, we choose the lowest possible value $`m_1=2`$. As it turns out, the qualitative picture of the vacuum structure is not sensitive to this choice.
In order to further simplify the potential (69), we impose the discrete <sub>2</sub> symmetry $`vv`$, $`tt`$ discussed in subsection 4.3. This symmetry can be made manifest by setting $`A_{1,m_1}=A_{1,m_1}^{}`$ or, equivalently, $`\stackrel{~}{A}_{1,m_1}=0`$. Without loss of generality we can furthermore choose $`\widehat{A}_{1,m_1}`$ to be positive, as, in the leading order approximation, considering negative values of $`\widehat{A}_{1,m_1}`$ merely corresponds to shifting $`vv+\pi `$. This choice of parameters then implies that $`v=0`$ corresponds to a local minimum of the potential in the $`v`$-direction.
Depending on the value of the remaining free parameter $`\widehat{A}_{1,m_1}`$ we obtain three classes of vacuum structures<sup>17</sup><sup>17</sup>17Recall that we consider rigid CY’s, where $`c<0`$, and the region $`r>2c`$ only., which are separated by two ($`c`$-dependent) thresholds $`\widehat{A}_{1,m_1}=A_{\text{min}}`$ and $`\widehat{A}_{1,m_1}=A_{\text{max}}`$. For $`\widehat{A}_{1,m_1}<A_{\text{min}}`$ we find the runaway behavior present in the classical and perturbative potentials. In this case there is no vacuum, except for the trivial one at $`r=\mathrm{}`$. For $`A_{\text{max}}<\widehat{A}_{1,m_1}`$ on the other hand, we obtain a stable AdS vacuum which is separated from the runaway vacuum by a saddle point of the potential where $`V_{\text{tot}}|_{\text{saddle}}>0`$. In this case all hypermultiplet moduli can be stabilized in the AdS vacuum. The most interesting case, however, occurs for $`A_{\text{min}}<\widehat{A}_{1,m_1}<A_{\text{max}}`$. In this case the AdS vacuum is lifted to positive cosmological constant and one obtains a *meta-stable dS vacuum*. As in the AdS case, this dS vacuum is separated from the runaway vacuum by a saddle point of the potential where $`V_{\text{tot}}|_{\text{saddle}}>0`$. *This meta-stable dS vacuum stabilizes all the hypermultiplet moduli*.
Furthermore, one can verify that increasing $`\widehat{A}_{1,m_1}`$ results in the (A)dS vacuum moving closer to the singularity at $`r=2c`$. This implies that the vacuum value (i.e., the value for which the string coupling is weakest) of $`r`$ is obtained by setting $`\widehat{A}_{1,m_1}=A_{\text{min}}`$. In this case $`r`$ is stabilized in the meta-stable dS vacuum and we will denote its corresponding value by $`r_{\mathrm{dS}}`$. Table 1 then summarizes the values for $`A_{\text{min}}`$, $`A_{\text{max}}`$, and $`r_{\mathrm{dS}}`$ for two “typical” values $`c=6/\pi 1.9`$ and $`c=10`$, respectively. The former corresponds to the $`𝒵`$-manifold (see e.g. ), the prototype of a rigid CY<sub>3</sub> with $`h_{1,2}=0`$ and $`h_{1,1}=36`$, while $`c=10`$ reflects a fictional rigid CY<sub>3</sub> where $`h_{1,1}=𝒪(100)`$. Table 1 indicates that decreasing $`|c|`$ also decreases the values for $`A_{\text{min}}`$ and $`A_{\text{max}}`$, while the relative width of the “dS window”, $`(A_{\text{max}}A_{\text{min}})/(A_{\text{max}}+A_{\text{min}})`$, stays approximately constant. Furthermore, we observe that decreasing $`|c|`$ moves $`r_{\mathrm{dS}}`$ closer to the singularity at $`r=2c`$.
Figs. 2 and 3 show the typical shape of $`V_{\text{tot}}`$ in the dS phase. Here we chose $`c=10`$ and $`\widehat{A}_{1,m_1}=9867`$. Fig. 2 displays the $`(r,v)`$-dependence of the potential, illustrating that we have indeed a meta-stable dS vacuum. Fig. 3 depicts $`V_{\text{tot}}`$ in the $`r`$-direction for $`v=0`$.
Let us conclude this section with a remark on the periodicity of $`V_{\text{tot}}`$ in the $`v`$-direction, which arises from the oscillatory terms in $`V_{\text{1-inst}}`$. As discussed in subsection 4.3, this reflects the fact that at leading order in the instanton correction there is a residual symmetry arising from the broken $`\beta `$-shift. Strictly speaking, we then obtain an infinite number of copies of the (A)dS vacua found above. Including higher order subleading terms, however, will break this discrete symmetry completely, thereby lifting the degeneracy between these vacua. In order to decide on the fate of these vacua one would have to sum the whole instanton series, which is, however, beyond the scope of this paper. We have verified that, when performing the above analysis by taking into account the sub- and subsubleading contributions to the potential (69), one still has (for a suitable choice of integration constants) one meta-stable dS vacuum, while the other local minima generically become meta-stable AdS vacua. This analysis gives some evidence that the qualitative picture found above will remain valid after resumming the instanton series. In particular, we expect that the stabilization of the hypermultiplets will also be a feature of the complete instanton solution.
To make further progress on this issue, we have to improve our understanding on membrane instanton calculations beyond what has been done in . Ideally we would like to fix the numerical coefficients in our instanton expansion completely through the microscopic string theory description. For (some) more general CY<sub>3</sub> compactifications of the type IIA string this may be done using the duality to heterotic string theory on $`K3\times T^2`$ , where the hypermultiplet moduli space is classically exact in the string coupling constant. However, this duality generically involves more than one hypermultiplet, which requires a more generic setup than what is considered here. Furthermore, it would be interesting to investigate whether these coefficients have some deeper meaning in the context of topological string theory, analogous to the coefficients appearing in the D3-brane instanton corrections to orientifold compactifications of the type IIB string recently investigated in .
## 7 Discussion
Let us now discuss the relation between KKLT and our set-up. In order to stabilize all moduli and to obtain a meta-stable dS vacuum, KKLT proposed a three step procedure, where first all moduli apart from the dilaton were fixed by fluxes, second the dilaton is stabilized by non-perturbative instanton effects at an AdS vacuum, and finally a positive energy contribution (in form of anti-D3-branes) is added to lift this vacuum to a meta-stable dS vacuum. When including a space-filling RR 3-form flux in our case, the classical potential is positive definite and of runaway type. There is no vacuum, except the one at vanishing string coupling, and both RR scalars correspond to flat directions. This does not change when the perturbative corrections to the universal hypermultiplet found in are included. The picture changes completely when taking into account the leading membrane instanton corrections to the universal hypermultiplet. These corrections to the scalar metric lift the flat directions corresponding to the RR scalars, so that all the present moduli are fixed by the potential. Furthermore, by making a suitable choice of the numerical parameters corresponding to the one-loop determinant around a one-instanton background, the moduli can be stabilized in a meta-stable dS vacuum at small string coupling $`g_s1`$. Here, the appearance of the dS vacuum does not require to add a positive energy contribution (like anti-D3-branes) by hand, as this contribution is already provided by the background flux when taking the perturbative corrections to the hypermultiplet geometry into account. This picture is completely analogous to the one obtained for type the IIB orientifold compactifications studied in where it was found that the leading order $`\alpha ^{}`$ corrections to the Kähler potential together with the leading D3-brane instanton correction can also give rise to meta-stable dS vacua without the need of adding an additional positive energy contribution. On phenomenological grounds, it would be interesting to analyze the stability and lifetime of these dS vacua along the lines of .
In our model we have truncated the vector multiplets that arise in a realistic compactification of type IIA strings on a rigid Calabi-Yau manifold. In order to address the issue of stabilizing these moduli as well, it would be necessary to include them in the effective action. This is, however, beyond the scope of the present paper. In this context, let us remark that including these vector multiplets (whose scalar fields correspond to the complexified Kähler moduli of the compactification) also allows to consider more general fluxes, like e.g. 2- and 4-form fluxes related to the even homology cycles of the rigid Calabi-Yau manifold, which could then be used to stabilize these moduli as well. Previous investigations on this topic indicate that these moduli will likely be stabilized at special points of the Kähler moduli space where the Calabi-Yau geometry degenerates (as e.g. at conifold points), either through fluxes or non-perturbative effects arising from (Lorentzian) branes wrapping the degenerate cycles . In this context it is also interesting to note that instanton corrections of the type discussed in this paper play an important role in a proper understanding of string theory at these degeneration points . It would therefore be highly desirable to extend the present analysis to include vector multiplets.
Possible extensions of the present work would be to include also fivebrane instanton corrections to the universal hypermultiplet. Some results were already obtained in , and it would be challenging to combine them with the results obtained in this paper. Another direction to pursue is to consider more generic Calabi-Yau threefolds, in which one also has additional hypermultiplets. The quaternion-Kähler manifold would then be higher dimensional, and it would be interesting to study instanton effects in this context and see how they stabilize all the complex structure moduli. We leave this open for future research.
###### Acknowledgments.
We thank Lilia Anguelova for collaboration in the early stages of the project. We also thank Serguei Alexandrov, David Calderbank, Jan de Boer, Mathijs de Vroome and Albrecht Klemm for stimulating discussions. This project was initiated during the Simons Workshop in Mathematics and Physics, Stony Brook 2004. UT was supported by the DFG within the priority program SPP 1096 on string theory.
## Appendix A Notation and conventions
In order to study the scalar potential in gauged supergravity, we need to fix the conventions and normalizations of the various quaternionic quantities. We mainly follow the notation of , but with a few modifications on the conventions mentioned explicitly below.
### A.1 General properties of quaternion-Kähler geometry
The quaternionic structure is normalized such that<sup>18</sup><sup>18</sup>18The $`J^\mathrm{\Lambda }`$ defined here differ from by a minus sign.
$$J^\mathrm{\Lambda }J^\mathrm{\Sigma }=\delta ^{\mathrm{\Lambda }\mathrm{\Sigma }}\epsilon ^{\mathrm{\Lambda }\mathrm{\Sigma }\mathrm{\Pi }}J^\mathrm{\Pi },$$
(70)
with $`\mathrm{\Lambda },\mathrm{\Sigma },\mathrm{\Pi }=1,2,3`$. There exist quaternionic 1-form vielbeine $`V_i^\alpha `$, in terms of which the line element reads
$$\mathrm{d}s^2=G_{AB}\mathrm{d}\varphi ^A\mathrm{d}\varphi ^B=G_{\overline{\alpha }\beta }V_i^\beta \overline{V}^{i\overline{\alpha }}.$$
(71)
Here, $`G_{AB}`$ is the quaternionic metric, $`\overline{V}^{i\overline{\alpha }}`$ is the complex conjugate of $`V_i^\alpha `$, and $`G_{\overline{\alpha }\beta }`$ is the tangent space metric that appears in front of the kinetic terms of the fermions. The quaternionic 2-forms can then be written as
$$J^\mathrm{\Lambda }=\frac{\mathrm{i}}{2}G_{\overline{\alpha }\beta }V_i^\beta \overline{V}^{j\overline{\alpha }}(\tau ^\mathrm{\Lambda })^i{}_{j}{}^{},$$
(72)
where $`\tau ^\mathrm{\Lambda }`$ are the Pauli matrices. For $`4n`$-dimensional QK manifolds, the range of the indices is $`i=1,2`$, $`\alpha =1,\mathrm{},2n`$, and $`A=1,\mathrm{},4n`$.
Quaternion-Kähler manifolds are Einstein, and hence the Ricci tensor is proportional to the metric. Following , we have
$$R_{AB}=\frac{1}{4n}G_{AB}R,$$
(73)
where $`R`$ is the (constant) Ricci scalar. Furthermore, there exist SU(2) connection 1-forms $`\stackrel{}{𝒱}=\stackrel{}{𝒱}_A\mathrm{d}\varphi ^A`$ with SU(2) curvature<sup>19</sup><sup>19</sup>19The convention for the SU(2) connection and curvature is chosen to be the same as e.g. in . With respect to , our SU(2) connection is chosen (minus) twice the one in , and therefore also the SU(2) curvature is (minus) twice as large.
$$\stackrel{}{}\mathrm{d}\stackrel{}{𝒱}\frac{1}{2}\stackrel{}{𝒱}\times \stackrel{}{𝒱}.$$
(74)
The exterior derivative on $`\stackrel{}{}`$ yields the Bianchi identities
$$\mathrm{d}\stackrel{}{}=\stackrel{}{𝒱}\times \stackrel{}{}.$$
(75)
The relation between SU(2) curvature and quaternionic 2-forms reads
$$\stackrel{}{}=\nu \stackrel{}{J},\nu \frac{1}{4n(n+2)}R.$$
(76)
For the gauging, we need the conventions for the moment maps. They are defined from<sup>20</sup><sup>20</sup>20Our definition of the moment map is the same as in . This normalization is different from , and our moment maps are (minus) two times the ones defined in .
$$\stackrel{}{J}_{AB}k_I^B=D_A\stackrel{}{P}_I=(_A\stackrel{}{𝒱}_A\times )\stackrel{}{P}_I,$$
(77)
where $`I`$ labels the different isometries and $`D_A`$ is the SU(2) covariant derivative.
One can solve this relation for the moment maps to get
$$\stackrel{}{P}_I=\frac{1}{2n\nu }\stackrel{}{J}^A{}_{B}{}^{}D_{A}^{}k_I^B.$$
(78)
Notice that the right-hand side is independent of the metric, except for the factor $`\nu `$. Choosing this factor sets the scale of the metric, and for the universal hypermultiplet that we discuss below, we set the scale<sup>21</sup><sup>21</sup>21From the definition of $`\nu `$, it is clear that changing the QK metric $`G_{AB}\lambda G_{AB}`$ changes the value of $`\nu `$ according to $`\nu \lambda ^1\nu `$, while keeping the first relation in (76) invariant. such that $`\nu =1/2`$.
In supergravity, the value of $`\nu `$ is fixed in terms of the gravitational coupling constant. If we normalize the kinetic terms of the graviton and scalars in the supergravity action as
$$e^1_{\text{kin}}=\frac{1}{2\kappa ^2}R(e)\frac{1}{2}G_{AB}_\mu \varphi ^A^\mu \varphi ^B,$$
(79)
then local supersymmetry fixes $`\nu =\kappa ^2`$. This is in accordance with , and with after a rescaling of the metric $`G_{AB}`$ with a factor 1/2. For the universal hypermultiplet, we will work with conventions in which $`\nu =1/2`$, so we set $`\kappa ^2=1/2`$ below. To compare with , we first multiply the Lagrangian (79) by 2 and then set $`\kappa ^2=2`$.
We now include the scalar potential that arises after gauging a single isometry. The isometry can then be gauged by the graviphoton and in the absence of any further vector multiplets, the relevant terms in the Lagrangian are
$$e^1=\frac{1}{2\kappa ^2}R\frac{1}{2}G_{AB}D_\mu \varphi ^AD^\mu \varphi ^B\left(2\kappa ^2G_{AB}k^Ak^B3\stackrel{}{P}\stackrel{}{P}\right).$$
(80)
Here, $`D_\mu `$ is the covariant derivative with respect to the gauged isometry that corresponds to the Killing vector $`k^A`$. The factors of $`\kappa `$ appear on dimensional grounds, as one can easily verify. For $`\kappa ^2=2`$ this agrees precisely with the result in ; here, however, we set $`\kappa ^2=1/2`$.
Our conventions are chosen such that they naturally apply to the universal hypermultiplet metric and the conventions used in . At the classical level we have
$$\mathrm{d}s^2=G_{AB}\mathrm{d}\varphi ^A\mathrm{d}\varphi ^B=\mathrm{d}\varphi ^2+\mathrm{e}^\varphi (\mathrm{d}\chi ^2+\mathrm{d}\phi ^2)+\mathrm{e}^{2\varphi }(\mathrm{d}\sigma +\chi \mathrm{d}\phi )^2.$$
(81)
For the corresponding Ricci tensor we find
$$R_{AB}=\frac{3}{2}G_{AB}.$$
(82)
The Ricci scalar is then $`R=6`$ and therefore we have $`\nu =1/2`$. This implies that in these conventions we should set $`\kappa ^2=1/2`$, which is equivalent to a cosmological constant $`\mathrm{\Lambda }=3/2`$ on the quaternion-Kähler manifold.
### A.2 Quaternion-Kähler geometry of the PT metric
The quaternionic properties of the PT metric can be demonstrated by constructing the corresponding quaternionic 1-form vielbeine (71), which we parameterize as
$$V_i^\alpha =\left(\begin{array}{cc}\overline{a}& \overline{b}\\ b& a\end{array}\right).$$
(83)
Substituting this ansatz into (71), we obtain
$$\mathrm{d}s^2=a\overline{a}+b\overline{b}+\text{c.c.}.$$
(84)
Comparing this expression with the PT metric (7), we can choose
$$a=\frac{1}{\sqrt{2}r}\left(f^{1/2}\mathrm{d}r+\mathrm{i}f^{1/2}(\mathrm{d}t+\mathrm{\Theta })\right),b=\frac{1}{\sqrt{2}r}(f\mathrm{e}^h)^{1/2}\left(\mathrm{d}u+\mathrm{i}\mathrm{d}v\right).$$
(85)
The computation of the quaternionic 2-forms (72) then yields
$$J^1=\mathrm{i}(ab\overline{a}\overline{b}),J^2=ab+\overline{a}\overline{b},J^3=\mathrm{i}(a\overline{a}+b\overline{b}).$$
(86)
These satisfy the quaternionic algebra (70).
Using (76) and (75), we then determine the SU(2) connection for the PT metric,
$$𝒱^1=\frac{1}{r}\mathrm{e}^{h/2}\mathrm{d}v,𝒱^2=\frac{1}{r}\mathrm{e}^{h/2}\mathrm{d}u,$$
$$𝒱^3=\frac{1}{2r}(\mathrm{d}t+\mathrm{\Theta })\frac{1}{2}(_vh\mathrm{d}u_uh\mathrm{d}v).$$
(87)
The PT metric has a shift symmetry in $`t`$. In coordinates $`(r,u,v,t)`$ the corresponding Killing vector is given by
$$k^A=(\mathrm{\hspace{0.17em}0},\mathrm{\hspace{0.17em}0},\mathrm{\hspace{0.17em}0},e_0)^\mathrm{T}.$$
(88)
The moment maps of this shift symmetry can be computed from (78). The result is *independent* of the functions $`f`$, $`h`$, and $`\mathrm{\Theta }`$ and reads
$$P^1=0,P^2=0,P^3=\frac{e_0}{r}.$$
(89)
### A.3 The 4-fermion coupling of the PT metric
In order to make contact with the string calculation of , we need to construct the symmetric tensor $`𝒲_{\alpha \beta \gamma \delta }`$, which appears in the four-fermion term. This can be done along the lines outlined in . Note that the tensor $`\mathrm{\Omega }_{XYZW}`$ appearing in Bagger and Witten is totally symmetric for rigid supersymmetry, but *not* in the supergravity case.
The symmetric tensor $`𝒲_{\alpha \beta \gamma \delta }`$ can be obtained from the curvature decomposition
$$R_{ABCD}=\nu (R^{\mathrm{SU}(2)})_{ABCD}+\frac{1}{2}L_{DC}^{}{}_{}{}^{\alpha \beta }𝒲_{\alpha \beta \gamma \delta }L_{AB}^{}{}_{}{}^{\gamma \delta }.$$
(90)
Here,
$$(R^{\mathrm{SU}(2)})_{ABCD}=\frac{1}{2}g_{D[A}g_{B]C}+\frac{1}{2}J_{AB}^\mathrm{\Lambda }J_{DC}^\mathrm{\Lambda }\frac{1}{2}J_{D[A}^\mathrm{\Lambda }J_{B]C}^\mathrm{\Lambda },$$
(91)
and
$$L_{AB\alpha }^{}{}_{}{}^{\beta }=V_{Ai\alpha }\overline{V}_B^{i\beta }.$$
(92)
Eq. (90) can be solved for $`𝒲_{\alpha \beta \gamma \delta }`$ by using the inverse relation for $`L_{AB}^{}{}_{}{}^{\alpha \beta }`$:
$$\frac{1}{2}V_\gamma ^{iB}V_{i\delta }^AL_{AB}^{}{}_{}{}^{\alpha \beta }=\delta _\gamma ^\alpha \delta _\delta ^\beta .$$
(93)
The resulting expression for $`𝒲_{\alpha \beta \gamma \delta }`$ then reads:
$$𝒲_{\alpha \beta \gamma \delta }=\frac{1}{2}ϵ^{ij}ϵ^{kl}V_{i\delta }^AV_{j\gamma }^BV_{k\beta }^DV_{l\alpha }^C\left(R_{ABCD}\nu (R^{\mathrm{SU}(2)})_{ABCD}\right).$$
(94)
The components of $`𝒲_{\alpha \beta \gamma \delta }`$ can now be obtained by calculating $`R_{ABCD}`$ for the PT metric (7) and substituting the expressions for the vielbeins and complex structures obtained above in the corresponding definitions. In order to write the independent components of $`𝒲_{\alpha \beta \gamma \delta }`$ in a compact way, it is useful to introduce the complex variable $`z=u+\mathrm{i}v`$. The result is given in (3.2).
## Appendix B Tensor multiplet description
Consider the hypermultiplet Lagrangian based on the PT metric (7). It is interesting to write down the $`N=2`$ tensor multiplet Lagrangian obtained after dualizing the scalar $`t`$ into a 2-form gauge potential with field strength $`H_{\mu \nu \rho }`$. Using the results of , this Lagrangian can easily be read off,
$$_T=\frac{1}{2}r^2fH_\mu H^\mu \frac{1}{2}𝒢_{AB}_\mu \varphi ^A^\mu \varphi ^B\mathrm{\Theta }_AH^\mu _\mu \varphi ^A.$$
(95)
Here, $`\mathrm{\Theta }_A`$ are the three components of the one-form defined in (10), and $`𝒢_{AB}`$ is the metric on the manifold spanned by the three scalars $`(r,u,v)`$. The line element can be written as
$$\mathrm{d}s^2=\frac{f}{r^2}\left[\mathrm{d}r^2+\mathrm{e}^h\left(\mathrm{d}u^2+\mathrm{d}v^2\right)\right],$$
(96)
where $`\mathrm{e}^h`$ satisfies the Toda equation and $`f(r,u,v)`$ the constraint (9). This 3-dimensional geometry is related to Einstein-Weyl spaces, as explained in .
## Appendix C Details of the Toda solution
This appendix collects several technical details about the solution of the Toda equation constructed in section 4. We start by proving $`m_n2`$ in subsection C.1, while the proof for $`\alpha =0`$ is given in subsection C.2. The derivation of the one-instanton solution is given in subsection C.3.
### C.1 The lower bound on $`m_n`$
In this subsection we establish $`m_n2`$. Our starting point is the ansatz (17), which we substitute into the Toda equation (21). This results in the following power series expansion<sup>22</sup><sup>22</sup>22Here we have not performed the splitting into instanton sectors yet.
$`0=`$ $`{\displaystyle \underset{n,m}{}}r^{m/2+\alpha +1}\mathrm{e}^{2n\sqrt{r}}\left[(\mathrm{\Delta }+n^2)f_{n,m}+(na_{m+1}r^{1/2}+b_{m+2}r^1)f_{n,m}\right]`$
$`+{\displaystyle \underset{n,m}{}}{\displaystyle \underset{n^{},m^{}}{}}r^{(m+m^{})/2+2\alpha }\mathrm{e}^{2(n+n^{})\sqrt{r}}[f_{n^{},m^{}}(\mathrm{\Delta }+2n^2)f_{n,m}`$
$`f_{n,m}f_{n^{},m^{}}+2(a_{m+1}r^{1/2}+b_{m+2}r^1)f_{n,m}f_{n^{},m^{}}]`$
$`+{\displaystyle \underset{n,m}{}}{\displaystyle \underset{n^{},m^{}}{}}{\displaystyle \underset{n^{\prime \prime },m^{\prime \prime }}{}}r^{(m+m^{}+m^{\prime \prime })/2+3\alpha 1}\mathrm{e}^{2(n+n^{}+n^{\prime \prime })\sqrt{r}}f_{n,m}f_{n^{},m^{}}f_{n^{\prime \prime },m^{\prime \prime }}`$
$`\times \left[n^2+na_{m+1}r^{1/2}+b_{m+2}r^1\right],`$ (97)
where we have extended the definitions for $`a_m`$, $`b_m`$ given in (23) to non-zero $`\alpha `$:
$$a_m=\frac{1}{2}(2m4\alpha 1),b_m=\frac{1}{4}(m2\alpha )(m2\alpha 2).$$
(98)
In order to obtain a bound on $`m_n`$ (for which the $`f_{n,m_n}0`$), we extract the leading order contributions in the $`r`$-expansion arising from the single, double and triple sum in (C.1). Starting at $`n=1`$ and working iteratively towards higher values $`n=2,3,\mathrm{}`$, we find that at a fixed value of $`n`$ these contributions are proportional to
single sum $`r^{m_n/2+\alpha +1}`$
double sum $`r^{m_n+2\alpha }`$
triple sum $`r^{3m_n/2+3\alpha 1}.`$ (99)
Investigating the $`m_n`$-dependence of these relations, we find that for $`m_n3`$ the leading order term in $`r`$ arises from the triple sum, which decouples from all the other terms in (C.1).
We now assume that for a fixed value $`n`$ there exsists an $`f_{n,m_n}0`$ for $`m_n3`$. Extracting the equation leading in $`r`$ from (C.1), we find that
$$n^2f_{n,m_n}^3=0,m_n3,$$
(100)
which has $`f_{n,m_n}=0`$ as its only solution. Hence, we establish the lower bound
$$m_n2$$
(101)
for all values of $`n`$ or, equivalently, all instanton sectors.<sup>23</sup><sup>23</sup>23Notice that this argument is not quite sufficient to also fix $`\alpha =0`$, as for $`\alpha =1/4`$ the single and triple sums do not decouple, which has been crucial in establishing (100).
### C.2 Fixing the parameter $`\alpha `$
When making the ansatz (17) in order to describe membrane instanton corrections to the universal hypermultiplet, we included the parameter $`\alpha [0,1/2[`$ to allow for the possibility that the leading term in the instanton solution occurs with a fractional power of $`g_s`$. Based on the plausible assumption that the perturbation series around the instanton gives rise to a power series in $`g_s`$ (and not fractional powers thereof) we now give a proof that a consistent solution of the Toda equation requires $`\alpha =0`$.
Splitting (C.1) into instanton sectors gives us the following analogue of (4)
$`0={\displaystyle \underset{n,m}{}}`$ $`r^{m/2+\alpha }\mathrm{e}^{2n\sqrt{r}}\{(\mathrm{\Delta }+n^2)f_{n,m+2}+na_{m+2}f_{n,m+1}+b_{m+2}f_{n,m}`$
$`+{\displaystyle \underset{n^{},m^{}}{}}r^\alpha \mathrm{e}^{2n^{}\sqrt{r}}[\mathrm{\hspace{0.17em}2}na_{m^{}+1}f_{n^{},mm^{}1}+2b_{m^{}+2}f_{n^{},mm^{}2}`$
$`+f_{n^{},mm^{}}(\mathrm{\Delta }+2n^2)f_{n^{},mm^{}}]f_{n,m^{}}`$
$`+{\displaystyle \underset{n^{},m^{}}{}}{\displaystyle \underset{n^{\prime \prime },m^{\prime \prime }}{}}r^{2\alpha }\mathrm{e}^{2(n^{}+n^{\prime \prime })\sqrt{r}}f_{n,m^{}}f_{n^{},m^{\prime \prime }}[n^2f_{n^{\prime \prime },mm^{}m^{\prime \prime }2}`$
$`+na_{m^{}+1}f_{n^{\prime \prime },mm^{}m^{\prime \prime }3}+b_{m^{}+2}f_{n^{\prime \prime },mm^{}m^{\prime \prime }4}]\}.`$ (102)
Based on this equation we can now make several observations. First, we find that the $`N=1`$ sector of (C.2) still gives rise to (24), with the coefficients $`a_m`$, $`b_m`$ now replaced by (98). To lowest order, $`m=m_1`$, this is just the equation
$$(\mathrm{\Delta }+1)f_{1,m_1}(u,v)=0.$$
(103)
Second, we observe that the equation describing the $`N=2`$ sector is modified to
$`0`$ $`=(\mathrm{\Delta }+4)f_{2,m}+2a_mf_{2,m1}+b_mf_{2,m2}`$
$`+{\displaystyle \underset{m^{}}{}}r^\alpha \left[f_{1,mm^{}2}+a_{m^{}+1}f_{1,mm^{}3}+b_{m^{}+2}f_{1,mm^{}4}f_{1,mm^{}2}\right]f_{1,m^{}}.`$
Note that for $`\alpha =0`$ the sum appearing in the second line is just an inhomogeneous term to the equations determining $`f_{2,m}`$. For $`\alpha 0`$, however, the sum decouples due to the different powers in $`r`$. Therefore, in the case $`\alpha 0`$, the sum gives rise to an additional constraint equation, which is absent for $`\alpha =0`$. Since the sum contains the $`f_{1,m}`$ only, this additional relation imposes a restriction on the $`N=1`$ instanton solution. Upon using (103), this additional constraint reads, at the lowest level,
$$f_{1,m_1}^2(f_{1,m_1})^2=0.$$
(104)
For $`\alpha 0`$ a non-trivial 1-instanton solution has to satisfy both (103) and (104), so that for establishing $`\alpha =0`$ it suffices to show that these equations have no common non-trivial solution:
Suppose that $`f_{1,m_1}0`$, which by definition of $`f_{1,m_1}`$ has to hold. We then multiply (103) with $`f_{1,m_1}`$, giving
$$0=f_{1,m_1}\mathrm{\Delta }f_{1,m_1}+f_{1,m_1}^2=f_{1,m_1}\mathrm{\Delta }f_{1,m_1}+(f_{1,m_1})^2=\frac{1}{2}\mathrm{\Delta }f_{1,m_1}^2,$$
where we have used (104) in the first step. In terms of complex coordinates $`z=u+\mathrm{i}v`$ it is $`\mathrm{\Delta }=4_z_{\overline{z}}`$, and the general solution reads
$$f_{1,m_1}^2(z,\overline{z})=g(z)+\overline{g}(\overline{z}).$$
Substituting this back into (103), we find
$$0=(\mathrm{\Delta }+1)f_{1,m_1}=f_{1,m_1}^3\left[_zg(z)_{\overline{z}}\overline{g}(\overline{z})+(g(z)+\overline{g}(\overline{z}))^2\right],$$
which is equivalent to
$$_zg(z)_{\overline{z}}\overline{g}(\overline{z})=g(z)^2+2g(z)\overline{g}(\overline{z})+\overline{g}(\overline{z})^2.$$
Since the right-hand side of this expression contains terms which are (anti-) holomorphic, whereas the left-hand side does not, we find that the only solution is given by $`g(z)=\mathrm{i}c`$ with $`c`$ constant. Thus $`f_{1,m_1}=0`$, which contradicts our assumption and shows that the ansatz (17) does *not* give rise to a one-instanton sector if $`\alpha 0`$. Conversely, a non-trivial one-instanton sector exists for $`\alpha =0`$ only, which then fixes $`\alpha =0`$.
### C.3 The one-instanton solution
The general one-dimensional solution in the one-instanton sector was given in (25). The functions $`G_s(x)`$ introduced there are defined by
$$G_s(x)=x^{s+1}h_{s1}(x),$$
(105)
where $`h_s(x)=j_s(x)+\mathrm{i}y_s(x)`$ are the spherical Bessel functions of the third kind. For $`s0`$ the $`G_s(x)`$ have no poles. Explicitly, they read
$$G_0(x)=\mathrm{e}^{\mathrm{i}x},G_{s>0}(x)=2^s\mathrm{e}^{\mathrm{i}x}\underset{k=1}{\overset{s}{}}\frac{(2sk1)!}{(sk)!(k1)!}(2\mathrm{i}x)^k.$$
(106)
Using the properties
$$x^2h_s^{\prime \prime }+2xh_s^{}+\left[x^2s(s+1)\right]h_s=0,h_s^{}+\frac{s+1}{x}h_s=h_{s1},$$
(107)
we easily verify the relation
$$(_x^2+1)G_s(x)=2sG_{s1}(x).$$
(108)
The proof of (25) is now simple:
$`(_x^2+1)f_{1,m}(x)`$ $`=\mathrm{Re}{\displaystyle \underset{s0}{}}{\displaystyle \frac{1}{s!(2)^s}}k_{1,m}(s)(_x^2+1)G_s(x)`$
$`=\mathrm{Re}{\displaystyle \underset{s1}{}}{\displaystyle \frac{1}{(s1)!(2)^{s1}}}k_{1,m}(s)G_{s1}(x)`$
$`=\mathrm{Re}{\displaystyle \underset{s0}{}}{\displaystyle \frac{1}{s!(2)^s}}k_{1,m}(s+1)G_s(x)`$
$`=\mathrm{Re}{\displaystyle \underset{s0}{}}{\displaystyle \frac{1}{s!(2)^s}}\left[a_mk_{1,m1}(s)+b_mk_{1,m2}(s)\right]G_s(x)`$
$`=a_mf_{1,m1}(x)b_mf_{1,m2}(x).`$ (109)
For the general $`(u,v)`$-dependent solution given in (28), the proof is almost identical. |
warning/0506/math0506033.html | ar5iv | text | # LOSSES IN 𝑀/𝐺𝐼/𝑚/𝑛 QUEUES
## 1. Introduction
Analysis of loss queueing systems is very important from both the theoretical and practical points of view. While the multiserver loss queueing system $`M/GI/m/0`$ and its network extensions have been intensively studied (see the review paper of Kelly , the book of Ross and references in these sources), the information about $`M/GI/m/n`$ queueing systems ($`n1`$) is very scanty, because explicit results for characteristics of these queueing systems are unknown. (In the present paper, for multiserver queueing systems the notation $`M/GI/m/n`$ is used, where $`m`$ denotes the number of servers and $`n`$ denotes the number of waiting places. Another notation which is also acceptable in the literature is $`M/GI/m/m+n`$.)
From the practical point of view, $`M/GI/m/n`$ queueing systems serve as a model for telephone systems, where $`n`$ is the maximally possible number of calls that can wait in the line before their service start. The loss probability is one of the most significant performance characteristics. In the present paper, we study the expected number of losses during a busy period (the characteristic closely related to the stationary loss probability) under the assumption that the arrival rate ($`\lambda `$) is equal to the maximum service capacity ($`m\mu `$), which seems to be the most interesting from the theoretical point of view.
There are two main reasons for studying this case.
The first reason is that the case $`\lambda =m\mu `$ is a critical case for queueing systems with $`m`$ identical servers, i.e. the case associated with critically loaded systems. The theoretical and practical interest in studying heavily loaded loss systems is very high, and there are many results in the literature related to the analysis of the loss probability in heavily loaded systems. The asymptotic results for losses in heavily loaded single server systems ($`n\mathrm{}`$) such as $`M/GI/1/n`$ and $`GI/M/1/n`$ and for associated models of telecommunication systems and dams have been studied in , , , , and . Heavy-traffic analysis of losses in heavily loaded multiserver systems have been provided in , , and . The mathematical foundation of heavy traffic theory can be found in the textbook of Whitt . Although the case $`\lambda =m\mu `$ is idealistic, it enables us to understand the possible behaviour of the system in certain cases when the values $`\lambda `$ and $`m\mu `$ are close and approach one another as $`n`$ increases to infinity. (Obtaining nontrivial results in the cases $`\lambda <m\mu `$ and $`\lambda >m\mu `$ is a hard problem, so the analytic investigation of the aforementioned asymptotic behaviour as $`n`$ increases to infinity is difficult.)
The second reason is that $`\lambda =m\mu `$ is an interesting theoretical case associated with an extension of the following non-trivial property of the symmetric random walk. Let $`X_1`$, $`X_2`$, …, $`X_i`$, …, be a sequence of independent and identically distributed random variables taking the values $`\pm 1`$ with the equal probability $`\frac{1}{2}`$. Let $`S_0=0`$, and $`S_{i+1}=S_i+X_{i+1}`$, $`i0`$, be a symmetric random walk, and let $`t=\tau `$ be the first time instant after $`t=0`$ when this random walk returns to zero, i.e. $`S_\tau =0`$. It is known that the expected number of level-crossings through any level $`n1`$ (or $`n1`$) is equal to $`\frac{1}{2}`$ independently of that level. The mentioning of this fact (but in a slightly different formulation) can be found in Szekely , and its proof is given in Wolff , p.411. The reformulation of this fact in terms of queueing theory is as follows. Consider $`M/M/1/n`$ queueing system with equal arrival and service rates. For this system, the expected number of losses during a busy period is equal to 1 for all $`n0`$. It has been recently noticed that this property holds true for $`M/GI/1/n`$ queueing systems. Namely, it was shown in several recent papers (see Abramov , , , Righter , Wolff ), that under mutually equal expectations of interarrival and service time, the expected number of losses during a busy period is equal to 1 for all $`n0`$. Further extension of this property to queueing systems with batch arrivals have been given in Abramov , Wolff and Peköz, Righter and Xia . Applications of the aforementioned property of losses can be found in for analysis of lost messages in telecommunication systems and in for optimal control of large dams. Further relevant results associated with the properties of losses have been obtained in the paper by Peköz, Righter and Xia . They solved a characterization problem associated with the properties of losses in $`GI/M/1/n`$ queues and established similar properties for $`M/M/m/n`$ and $`M^X/M/m/n`$ queueing systems. Recently, a similar property related to consecutive losses in busy periods of $`M/GI/1/n`$ queueing systems has been discussed in . It follows from the results obtained in this paper that for $`M/GI/1/n`$ queueing systems with mutually equal expectations of interarrival and service times, the expected number of losses in series containing at least $`k>1`$ consecutive losses during a busy period generally depends on $`n`$. However, for $`M/M/1/n`$ queueing systems with equal arrival and service rates that expected number of consecutive losses during a busy period is the same constant (depending on the value $`k`$) for all $`n0`$.
The aim of the present paper is further theoretical contribution to this theory of losses, now to the theory of multiserver loss queueing systems. On the basis of the aforementioned results on losses in $`M/GI/1/n`$ and $`M/M/m/n`$ queueing systems we address the following open question. Does the result on losses in $`M/M/m/n`$ queueing systems remain true for those $`M/GI/m/n`$ too?
The answer on this question is not elementary. On one hand, under the assumption $`\lambda =m\mu `$ the expected numbers of losses in $`M/GI/m/0`$ and $`M/GI/m/n`$ queueing systems ($`m2`$ and $`n1`$) during their busy periods are different. A simple example for this confirmation can be built for $`M/GI/2/1`$ queueing systems having the service time distribution $`G(x)=1p\text{e}^{\mu _1x}q\text{e}^{\mu _2x}`$, $`p+q=1`$. The analysis of the stationary characteristics for these systems, resulting in an analysis of losses during a busy period, can be provided explicitly. Specifically, the structure of the $`9\times 9`$ Markov chain intensity matrix for the states of the Markov chain associated with an $`M/GI/2/1`$ queueing system shows a clear difference between the structure of the stationary probabilities in $`M/GI/2/1`$ queues and that in $`M/GI/2/0`$ queues given by the Erlang-Sevastyanov formulae. So, the parameters $`p`$, $`q`$, $`\mu _1`$ and $`\mu _2`$ can be chosen such that the expected number of losses during busy periods in these two queueing systems will be different.
On the other hand, the property of losses, which is similar to the aforementioned one, indeed holds. The correctness of this similar property for multiserver $`M/GI/m/n`$ queueing systems is proved in the present paper. Namely, we establish the following results.
Let $`L_{m,n}`$ denote the number of losses during a busy period of the $`M/GI/m/n`$ queueing system, let $`\lambda `$, $`\mu `$ be the arrival rate and, respectively, the reciprocal of the expected service time, and let $`m`$, $`n`$ denote the number of servers and, respectively, the number of waiting places. We will prove that, under the assumption $`\lambda =m\mu `$, the expected number of losses during a busy period of the $`M/GI/m/n`$ queueing system, $`\mathrm{E}L_{m,n}`$, is the same for all $`n1`$, which is not generally the same as that for the $`M/GI/m/0`$ loss queueing system (when $`n=0`$). In addition, if the probability distribution function of the service time belongs to the class NBU (New Better than Used), then $`\mathrm{E}L_{m,n}=\frac{cm^m}{m!}`$, where a constant $`c1`$ is independent of $`n1`$. In the opposite case of the NWU (New Worse than Used) service time distribution we correspondingly have $`\mathrm{E}L_{m,n}=\frac{cm^m}{m!}`$ with a constant $`c1`$ independent of $`n1`$ as well. (The constant $`c`$ becomes equal to 1 in the case of exponentially distributed service times.) Recall that a probability distribution function $`\mathrm{\Xi }(x)`$ of a nonnegative random variable is said to belong to the class NBU if for all $`x0`$ and $`y0`$ we have $`\overline{\mathrm{\Xi }}(x+y)\overline{\mathrm{\Xi }}(x)\overline{\mathrm{\Xi }}(y)`$, where $`\overline{\mathrm{\Xi }}(x)=1\mathrm{\Xi }(x)`$. If the opposite inequality holds, i.e. $`\overline{\mathrm{\Xi }}(x+y)\overline{\mathrm{\Xi }}(x)\overline{\mathrm{\Xi }}(y)`$, then $`\mathrm{\Xi }(x)`$ is said to belong to the class NWU.
The proof of the main results of this paper is based on an application of the level-crossing approach to the special type stationary processes. The construction of the level-crossings approach used in this paper is a substantially extended version of that used in the earlier papers by the author (e.g. , , , , and ) and by Pechinkin . It uses modern geometric methods of analysis and involves an algebraically close system of processes and a nontrivial construction of deleting intervals and merging the ends together with nontrivial applications of the PASTA property.
Throughout the paper, it is assumed that $`m2`$. (This is not the loss of generality since the case $`m=1`$ is known, see , and .)
The paper is organized as follows. In Section 2, which is the first part of the paper, $`M/M/m/n`$ queueing systems are studied. The results for $`M/M/m/n`$ queueing systems are then used in Section 3, which is the second part of the paper, in order to study $`M/GI/m/n`$ queueing systems. The study in both of Sections 2 and 3 is based on the level-crossing approach. The construction of level-crossings for $`M/M/m/n`$ queueing systems is then developed for $`M/GI/m/n`$ queueing systems as follows. The stationary processes associated with these queueing systems is considered, and the stochastic relations between the times spent in state $`m1`$ associated with $`m1`$ busy servers during a busy period of $`M/GI/m/n`$ ($`n1`$) and $`M/GI/m1/0`$ queueing systems are established. To prove these stochastic relations, some ideas from the paper of Pechinkin are involved to adapt and develop the level-crossing method for the problems of the present paper. The obtained stochastic relations are crucial, and they are then used to prove the main results of the paper in Section 4. In Section 5, possible development of the results for $`M^X/GI/m/n`$ queueing systems with batch arrivals is discussed.
## 2. The $`M/M/m/n`$ queueing system
In this section, the Markovian $`M/M/m/n`$ loss queueing system is studied with the aid of the level-crossings approach, in order to establish some relevant properties of this queueing system. Those properties are then developed for $`M/GI/m/n`$ queueing systems in the following sections.
Let $`f(j)`$, $`1jn+m+1`$, denote the number of customers arriving during a busy period who, upon their arrival, meet $`j1`$ customers in the system. It is clear that $`f(1)=1`$ with probability 1. Let $`t_{j,1}`$, $`t_{j,2}`$,…, $`t_{j,f(j)}`$ be the instants of arrival of these $`f(j)`$ customers, and let $`s_{j,1}`$, $`s_{j,2}`$,…, $`s_{j,f(j)}`$ be the instants of the service completions when there remain only $`j1`$ customers in the system. Notice, that $`t_{n+m+1,k}=s_{n+m+1,k}`$ for all $`k=1,2,\mathrm{},f(n+m+1)`$.
For $`1jn+m`$ let us consider the intervals
(2.1)
$$(t_{j,1},s_{j,1}],(t_{j,2},s_{j,2}],\mathrm{},(t_{j,f(j)},s_{j,f(j)}].$$
Then, by incrementing index $`j`$ we have the following intervals
(2.2)
$$(t_{j+1,1},s_{j+1,1}],(t_{j+1,2},s_{j+1,2}],\mathrm{},(t_{j+1,f(j+1)},s_{j+1,f(j+1)}].$$
Delete the intervals of (2.2) from those of (2.1) and merge the ends, that is each point $`t_{j+1,k}`$ with the corresponding point $`s_{j+1,k}`$, $`k=1,2,\mathrm{},f(j+1)`$ (see Figure 1).
Then $`f(j+1)`$ has the following properties. According to the property of the lack of memory of the exponential distribution, the residual service time for a service completion, after the procedure of deleting the interval and merging the ends as it is indicated above, remains exponentially distributed with parameter $`\mu \mathrm{min}(j,m)`$. Therefore, the number of points generated by merging the ends within the given interval $`(t_{j,1},s_{j,1}]`$ coincides in distribution with the number of arrivals of the Poisson process with rate $`\lambda `$ during an exponentially distributed service time with parameter $`\mu \mathrm{min}(j,m)`$. Namely, for $`1jm1`$ we obtain
$$\mathrm{E}\{f(j+1)|f(j)=1\}=\underset{u=1}{\overset{\mathrm{}}{}}u_0^{\mathrm{}}\text{e}^{\lambda x}\frac{(\lambda x)^u}{u!}j\mu \text{e}^{j\mu x}\text{d}x=\frac{\lambda }{j\mu }.$$
Considering now a random number $`f(j)`$ of intervals (2.1) we have
(2.3)
$$\mathrm{E}\{f(j+1)|f(j)\}=\frac{\lambda }{j\mu }f(j).$$
Analogously, denoting the load of the system by $`\rho =\frac{\lambda }{m\mu }`$, for $`mjm+n`$ we have
(2.4)
$$\mathrm{E}\{f(j+1)|f(j)\}=\frac{\lambda }{m\mu }f(j)=\rho f(j).$$
The properties (2.3) and (2.4) mean that the stochastic sequence
(2.5)
$$\{f(j+1)\left(\frac{\mu }{\lambda }\right)^j\underset{i=1}{\overset{j}{}}\mathrm{min}(i,m),_{j+1}\},_j=\sigma \{f(1),f(2),\mathrm{},f(j)\},$$
forms a martingale.
It follows from (2.5) that for $`0jm1`$
(2.6)
$$\mathrm{E}f(j+1)=\frac{\lambda ^j}{j!\mu ^j},$$
and for $`mjm+n`$
(2.7)
$$\mathrm{E}f(j+1)=\frac{\lambda ^m}{m!\mu ^m}\rho ^{jm}.$$
For example, when $`\rho =1`$ from (2.7) we obtain the particular case of the result of Peköz, Righter and Xia : $`\mathrm{E}L_{m,n}`$=$`\mathrm{E}f(n+m+1)=\frac{m^m}{m!}`$ for all $`n0`$, where $`L_{m,n}`$ denotes the number of losses during a busy period of the $`M/M/m/n`$ queueing system.
Next, let $`B(j)`$ be the period of time during a busy cycle of the $`M/M/m/n`$ queueing system when there are exactly $`j`$ customers in the system. For $`0jm+n`$ we have:
(2.8)
$$\lambda \mathrm{E}B(j)=\mathrm{E}f(j+1)=\{\begin{array}{cc}\frac{\lambda ^j}{j!\mu ^j},\hfill & \text{for}0jm1,\hfill \\ \frac{\lambda ^m}{m!\mu ^m}\rho ^{jm},\hfill & \text{for}mjm+n.\hfill \end{array}$$
Now, introduce the following notation. Let $`T_{m,n}`$ denote the length of a busy period of the $`M/M/m/n`$ queueing system, let $`T_{m,0}`$ denote the length of a busy period of the $`M/M/m/0`$ queueing system with the same arrival and service rates as in the initial $`M/M/m/n`$ queueing system, and let $`\zeta _n`$ denote the length of a busy period of $`M/M/1/n`$ queueing system with arrival rate $`\lambda `$ and service rate $`\mu m`$. From (2.6)-(2.8) for the expectation of a busy period of the $`M/M/m/n`$ queueing system we have
(2.9)
$$\mathrm{E}T_{m,n}=\underset{j=1}{\overset{n+m}{}}\mathrm{E}B(j)=\underset{j=1}{\overset{m1}{}}\frac{\lambda ^{j1}}{j!\mu ^j}+\frac{\lambda ^{m1}}{m!\mu ^m}\underset{j=0}{\overset{n}{}}\rho ^j.$$
In turn, for the expectation of a busy period of the $`M/M/m/0`$ queueing system we have
(2.10)
$$\mathrm{E}T_{m,0}=\underset{j=1}{\overset{m}{}}\mathrm{E}B(j)=\underset{j=1}{\overset{m}{}}\frac{\lambda ^{j1}}{j!\mu ^j},$$
where (2.10) is the particular case of (2.9) where $`n=0`$.
It is clear that $`T_{m,n}`$ contains one busy period $`T_{m1,0}`$, where the subscript $`m1`$ underlines that there are $`m1`$ servers, and a random number of independent busy periods, which will be called orbital busy periods. Denote an orbital busy period by $`\zeta _n`$. (It is assumed that an orbital busy period $`\zeta _n`$ starts at instant when an arriving customer finds $`m1`$ servers busy and occupies the $`m`$th server, and it finishes at the instant when after a service completion there at the first time remain only $`m1`$ busy servers.) Therefore, denoting the independent sequence of identically distributed orbital busy periods by $`\zeta _n^{(1)}`$, $`\zeta _n^{(2)}`$,…, we have
(2.11)
$$T_{m,n}\stackrel{d}{=}T_{m1,0}+\underset{i=1}{\overset{\kappa }{}}\zeta _n^{(i)},$$
where $`\kappa `$ is the random number of the aforementioned orbital busy periods and $`\stackrel{d}{=}`$ means an equality in distribution. It follows from (2.9), (2.10) and (2.11)
(2.12)
$$\mathrm{E}\underset{i=1}{\overset{\kappa }{}}\zeta _n^{(i)}=\frac{\lambda ^{m1}}{m!\mu ^m}\underset{j=0}{\overset{n}{}}\rho ^j.$$
On the other hand, the expectation of an orbital busy period $`\zeta _n`$ is
$$\mathrm{E}\zeta _n=\frac{1}{m\mu }\underset{j=0}{\overset{n}{}}\rho ^j$$
(this can be easily checked, for example, by the level-crossings method , and an application of Wald’s identity , p. 384), and we obtain
(2.13)
$$\mathrm{E}\kappa =\frac{\lambda ^{m1}}{(m1)!\mu ^{m1}}.$$
Thus, $`\mathrm{E}\kappa `$ coincides with the expectation of the number of losses during a busy period in the $`M/M/m/0`$ queueing system. In the case $`\rho =1`$ we have $`\mathrm{E}\kappa =\frac{m^m}{m!}`$.
## 3. $`M/GI/m/n`$ queueing systems
In this section, the inequalities between the times spent in the state $`m1`$ in the $`M/GI/m/n`$ ($`n1`$) and $`M/GI/m/0`$ queueing systems during their busy periods are derived.
Consider two queueing systems: $`M/GI/m/n`$ ($`n1`$) and $`M/GI/m/0`$ both having the same arrival rate $`\lambda `$ and probability distribution function of a service time $`G(x)`$, $`\frac{1}{\mu }=_0^{\mathrm{}}x\text{d}G(x)<\mathrm{}`$. Let $`T_{m,n}(m1)`$ denote the time spent in the state $`m1`$ during its busy period (i.e. the total time during a busy period when $`m1`$ servers are occupied) of the $`M/GI/m/n`$ queueing system, and let $`T_{m,0}(m1)`$ have the same meaning for the $`M/GI/m/0`$ queueing system.
We prove the following lemma.
###### Lemma 3.1.
Under the assumption that the service time distribution $`G(x)`$ belongs to the class NBU (NWU),
(3.1)
$$T_{m,n}(m1)_{st}(\text{resp.}_{st})T_{m,0}(m1),$$
###### Proof.
. The proof of the lemma is relatively long. In order to make it transparent and easily readable we strongly indicate the steps of this proof given by several propositions (properties). There are also six figures (Figures 2-7) illustrating the constructions in the proof. Each of these figures contain two graphs. The first (upper) of them indicates the initial (or intermediate) possible path of the process (sometimes two-dimensional), while the second (lower) one indicates the part of the path of one or two-dimensional process after a time scaling or specific transformation (e.g. in Figure 5). Arc braces in the graphs indicate the intervals that should be deleted and their ends merged.
Two-dimensional processes are shown as parallel graphs. For example, there are two parallel processes in Figure 3 which are shown in the upper graph, and there are two parallel processes which are shown in the lower graph. The same is in Figures 4, 6 and 7.
For the purpose of the present paper we use strictly stationary processes of order 1 or strictly 1-stationary processes. Recall the definition of a strictly stationary process of order $`n`$ (see , p.206).
###### Definition 3.2.
The process $`\xi (t)`$ is said to be strictly stationary of order $`n`$ or strictly $`n`$-stationary, if for a given positive $`n<\mathrm{}`$, any $`h`$ and $`t_1`$, $`t_2`$,…, $`t_n`$ the random vectors
$$(\xi (t_1),\xi (t_2),\mathrm{},\xi (t_n))\text{and}(\xi (t_1+h),\xi (t_2+h),\mathrm{},\xi (t_n+h))$$
have identical joint distributions.
If $`n=1`$ then we have strictly 1-stationary processes satisfying the property:
$$\mathrm{P}\{\xi (t)x\}=\mathrm{P}\{\xi (t+h)x\}.$$
The probability distribution function $`\mathrm{P}\{\xi (t)x\}`$ in this case will be called limiting stationary distribution.
The class of strictly 1-stationary processes is wider than the class of strictly stationary processes, where it is required that for all finite dimensional distributions
$$\mathrm{P}\{(\xi (t_1),\xi (t_2),\mathrm{},\xi (t_k))B_k\}=\mathrm{P}\{(\xi (t_1+h),\xi (t_2+h),\mathrm{},\xi (t_k+h))B_k\},$$
for any $`h`$ and any Borel set $`B_k^k`$. The reason of using strictly 1-stationary processes rather than strictly stationary processes themselves is that, the operation of deleting intervals and merging the ends is algebraically close with respect to strictly 1-stationary processes, and it is not closed with respect to strictly stationary processes. The last means that if $`\xi (t)`$ is a strictly 1-stationary process, then for any $`h>0`$ and arbitrary $`t_0`$ the new process
$$\xi _1(t)=\{\begin{array}{cc}\xi (t),\hfill & \text{if}tt_0,\hfill \\ \xi (t+h),\hfill & \text{if}t>t_0\hfill \end{array}$$
is also strictly 1-stationary and has the same one-dimensional distribution as the original process $`\xi (t)`$. The similar property is not longer valid for strictly stationary processes. If $`\xi (t)`$ is a strictly stationary process, then generally $`\xi _1(t)`$ is not strictly stationary.
In the following the prefix ‘strictly’ will be omitted, so strictly stationary and strictly 1-stationary processes will be correspondingly called stationary and 1-stationary processes.
Let us introduce $`m`$ independent and identically distributed stationary renewal processes (denoted below $`𝐱_m(t)`$) with a renewal period having the probability distribution function $`G(x)`$.
On the basis of these renewal processes we build the stationary $`m`$-dimensional Markov process $`𝐱_m(t)=\{\xi _1(t),\xi _2(t),\mathrm{},\xi _m(t)\}`$, the coordinates $`\xi _k(t)`$, $`k=`$1,2,…, $`m`$ of which are the residual times to the next renewal times in time moment $`t`$, following in an ascending order.
Let us now consider the two $`(m+1)`$-dimensional Markov processes corresponding to the $`M/GI/m/n`$ ($`n1`$) and $`M/GI/m/0`$ queueing systems, which are denoted by $`𝐲_{m,n}(t)`$ and $`𝐲_{m,0}(t)`$. Let $`Q_{m,n}(t)`$ denote the stationary queue-length process (the number of customers in the system) of the $`M/GI/m/n`$ queueing system, and let $`Q_{m,0}(t)`$ denote the stationary queue-length process corresponding to the $`M/GI/m/0`$ queueing system. We have:
$$𝐲_{m,n}(t)=\{\eta _1^{(m,n)}(t),\eta _2^{(m,n)}(t),\mathrm{},\eta _m^{(m,n)}(t),Q_{m,n}(t)\},$$
where
$$\{\eta _{m\nu _{m,n}(t)+1}^{(m,n)}(t),\eta _{m\nu _{m,n}(t)+2}^{(m,n)}(t),\mathrm{},\eta _m^{(m,n)}(t)\}$$
are the ordered residual service times corresponding to $`\nu _{m,n}(t)=\mathrm{min}\{m,Q_{m,n}(t)\}`$ customers in service in time $`t`$, and
$$\{\eta _1^{(m,n)}(t),\eta _2^{(m,n)}(t),\mathrm{},\eta _{m\nu _{m,n}(t)}^{(m,n)}(t)\}=\{0,0,\mathrm{},0\}$$
all are zeros. Analogously,
$$𝐲_{m,0}(t)=\{\eta _1^{(m,0)}(t),\eta _2^{(m,0)}(t),\mathrm{},\eta _m^{(m,0)}(t),Q_{m,0}(t)\},$$
only replacing the index $`n`$ with 0.
Let us delete all time intervals of the process $`𝐲_{m,n}(t)`$ related to the $`M/GI/m/n`$ queueing system ($`n1`$) where there are more than $`m1`$ or less than $`m1`$ customers and merge the ends. Remove the last component of the obtained process which is trivially equal to $`m1`$. Then we get the new ($`m1`$)-dimensional Markov process (in the following the prefix ‘Markov’ will be omitted and only used in the places where it is meaningful):
$$\widehat{𝐲}_{m1,n}(t)=\{\widehat{\eta }_1^{(m,n)}(t),\widehat{\eta }_2^{(m,n)}(t),\mathrm{},\widehat{\eta }_{m1}^{(m,n)}(t)\},$$
the components of which are now denoted by hat. All of the components of this vector are 1-stationary, which is a consequence of the existence of the limiting stationary probabilities of the processes $`\eta _j^{(m,n)}(t)`$, $`j=1,2,\mathrm{},m`$ (e.g. Takács ) and consequently those of the processes $`\widehat{\eta }_j^{(m,n)}(t)`$, $`j=1,2,\mathrm{},m`$. The joint limiting stationary distribution of the process $`\widehat{𝐲}_{m1,n}(t)`$ can be obtained by conditioning of that of the processes $`𝐲_{m,n}(t)`$ given $`Q_{m,n}(t)=m1`$.
The similar operation of deleting intervals and merging the ends, where there are less than $`m1`$ customers in the system, for the process $`𝐲_{m1,0}(t)`$ is used. We correspondingly have
$$\widehat{𝐲}_{m1,0}(t)=\{\widehat{\eta }_1^{(m,0)}(t),\widehat{\eta }_2^{(m,0)}(t),\mathrm{},\widehat{\eta }_{m1}^{(m,0)}(t)\}.$$
We establish the following elementary property related to the $`M/GI/1/n`$ queueing systems, $`n`$=0,1,…
###### Property 3.3.
(3.2)
$$\mathrm{P}\left\{\eta _1^{(1,n)}(t)B_1\right|Q_{1,n}(t)1\}=\mathrm{P}\{𝐱_1(t)B_1\},$$
for any Borel set $`B_1^1`$.
###### Proof.
Delete all of the intervals where the server is free and merge the corresponding ends (see Figure 2). Then in the new time scale, the processes all are structured as a stationary renewal process with the length of a period having the probability distribution function $`G(x)`$. Therefore (3.2) follows. ∎
In order to establish similar properties in the case $`m=2`$ let us first study the properties of 1-stationary processes and explain the construction of tagged server station which is substantially used in our construction throughout the paper.
Properties of 1-stationary processes. Recall (see Definition 3.2) that if $`\xi (t)`$ is a 1-stationary process, then for any $`h`$ and $`t_0`$ the probability distributions of $`\xi (t_0)`$ and $`\xi (t_0+h)`$ are the same. The result remains correct (due to the total probability formula) if $`h`$ is replaced by random variable $`\vartheta `$ with some given probability distribution, which is assumed to be independent of the process $`\xi `$. Namely, we have:
(3.3) $`\mathrm{P}\{\xi (t_0+\vartheta )x\}`$ $`={\displaystyle _{\mathrm{}}^{\mathrm{}}}\mathrm{P}\{\xi (t_0+h)x\}\text{d}\mathrm{P}\{\vartheta h\}`$
$`=\mathrm{P}\{\xi (t_0)x\}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\text{d}\mathrm{P}\{\vartheta h\}`$
$`=\mathrm{P}\{\xi (t_0)x\}.`$
That is, $`\xi (t_0)`$ and $`\xi (t_0+\vartheta )`$ have the same distribution.
The above property will be used for the following construction of the sequence of 1-stationary processes $`\xi ^{(1)}(t)`$, $`\xi ^{(2)}(t)`$, …, having identical one-dimensional distributions.
Let $`\xi ^{(0)}(t)=\xi (t)`$ be a 1-stationary process, let $`t_1`$ be an arbitrary point, and let $`\vartheta _1`$ be a random variable with some given probability distribution, which is independent of the process $`\xi ^{(0)}(t)`$. Let us build a new process $`\xi ^{(1)}(t)`$ as follows. Put
(3.4)
$$\xi ^{(1)}(t)=\{\begin{array}{cc}\xi ^{(0)}(t),\hfill & \text{for all}t<t_1,\hfill \\ \xi ^{(0)}(t+\vartheta _1),\hfill & \text{for all}tt_1.\hfill \end{array}$$
Since the probability distributions of $`\xi (t)`$ and $`\xi (t+\vartheta _1)`$ are the same for all $`tt_1`$, then the processes $`\xi (t)`$ and $`\xi ^{(1)}(t)`$ have the same one-dimensional distributions, and $`\xi ^{(1)}(t)`$ is a 1-stationary process as well.
With a new point $`t_2`$ and a random variable $`\vartheta _2`$, which is assumed to be independent of the process $`\xi ^{(0)}(t)`$ and random variable $`\vartheta _1`$ (therefore, it is also independent of the process $`\xi ^{(1)}(t)`$) by the same manner one can build the new 1-stationary process $`\xi ^{(2)}(t)`$. Specifically,
(3.5)
$$\xi ^{(2)}(t)=\{\begin{array}{cc}\xi ^{(1)}(t),\hfill & \text{for all}t<t_2,\hfill \\ \xi ^{(1)}(t+\vartheta _2),\hfill & \text{for all}tt_2.\hfill \end{array}$$
The new process $`\xi ^{(2)}(t)`$ has the same one-dimensional distribution as the processes $`\xi ^{(0)}(t)`$ and $`\xi ^{(1)}(t)`$. The procedure can be infinitely continued, and one can obtain the infinite family of 1-stationary processes, having the same one-dimensional distribution.
The points $`t_1`$, $`t_2`$,…in the above construction are assumed to be some fixed (non-random) points. However, the construction also remains correct in the case of random points $`t_0`$, $`t_1`$,…of Poisson process, since according to the PASTA property the limiting stationary distribution of a 1-stationary process in a point of a Poisson arrival coincides with the limiting stationary distribution of the same 1-stationary process in an arbitrary non-random point. Furthermore, the aforementioned property of process remains correct when the random points $`t_0`$, $`t_1`$,…are the points of the process which is not necessarily Poisson but belongs to the special class of processes that contains Poisson. In this case the property is called ASTA (e.g. ).
1-stationary Poisson process. Consider an important particular case when the process $`\xi (t)`$ is Poisson. Let $`\xi ^{(0)}(t)=\xi (t)`$. Then the process $`\xi ^{(1)}(t)`$ that obtained by (3.4) is no longer Poisson. Its limiting stationary distribution is the same as that of the original process $`\xi (t)`$, but the joint distributions of this process given in different points $`s`$ and $`t`$ distinguish from those of the original process $`\xi (t)`$.
The process $`\xi ^{(1)}(t)`$ will be called 1-stationary Poisson process or simply 1-Poisson. Clearly, that the further processes such as $`\xi ^{(2)}(t)`$, $`\xi ^{(3)}(t)`$, …that obtained similarly to the procedure in (3.4), (3.5) all are 1-Poisson with the same limiting stationary distribution. According to the above construction, a 1-Poisson process is obtained by deleting intervals and merging the ends of an original Poisson process. Therefore, a sequence of 1-Poisson arrival time instants is a scaled subsequence of those instants of the ordinary Poisson arrivals. Hence, for 1-Poisson process the ASTA property is satisfied, i.e. 1-Poisson arrivals see time averages exactly as those Poisson arrivals.
Tagged server station. Consider a stationary queueing system $`M/GI/m/n`$, which is referred to as main server station, and in addition to this queueing system introduce another one containing a server station in order to register specific arrivals, for example losses or, say, customers waiting their service in the main system. This server station is called tagged server station. The main idea of introducing tagged server stations is to decompose the main system as follows. Assume that along with a Poisson stream of arrivals of customers occupying servers in the main system, there is another stream of arrivals of customers in the tagged server system. For instance, the losses in the main system can be supposed to occupy the tagged server station. Although the stream of these losses is not Poisson (see e.g. , p. 83 or , p. 320), it is shown later that it is 1-Poisson. Therefore, the original system is decomposed into smaller systems with the same (1-Poisson) type of input stream. It is worth noting that only one dimensional distributions of 1-Poisson process are the same for all of them that generated similarly to the procedure in (3.4), (3.5). However, the two-dimensional distributions are distinct in general.
In fact, applications of a tagged server station is wider than that, and its aim is a proper decomposition of the original system into the main and tagged systems for further study of the properties of losses.
Another idea of using tagged server stations is a proper application of the ASTA property as follows. At the moment of arrival of a customer in the tagged server station, the stationary characteristics in the main server station remain the same. Specifically, the distributions of residual service times in servers of the main station at the moment of arrival of a customer in the tagged station coincides with the usual stationary distributions of these residual service times.
Let us now formulate and prove a property similar to Property 3.3 for $`m=2`$. We have the following.
###### Property 3.4.
For the $`M/GI/2/0`$ queueing system we have:
(3.6)
$$\mathrm{P}\{\widehat{𝐲}_{1,0}(t)B_1\}=\mathrm{P}\{𝐱_1(t)B_1\},$$
###### Proof.
In order to simplify the explanation in this case, let us consider two auxiliary stationary one-dimensional processes $`\zeta _{1,0}(t)`$ and $`\zeta _{2,0}(t)`$. The first process describes a residual service time in the first server, and the second one describes a residual service time in the second server. If the $`i`$th server ($`i=1,2`$) is free in time $`t`$, then we set $`\zeta _{i,0}(t):=0`$.
Our further convention is that the first server is a tagged server. We assume that if at the moment of arrival of a customer both of the servers are free, he/she occupies the first server. Clearly that this assumption is not a loss of generality. For instance, if we assume that both of the servers are equivalent and can be occupied with the equal probability $`\frac{1}{2}`$, then an occupied server (let it be the first) can be called tagged. In another busy period start an arriving customer can occupy the second server. It this case, nothing is changed if the servers will be renumbered, and the occupied server will be numbered as first and called tagged.
Our main idea is a decomposition of the stationary $`M/GI/2/0`$ queueing system into two systems and study the properties of stationary (1-stationary) processes $`\zeta _{1,0}(t)`$ and $`\zeta _{2,0}(t)`$. The arrival stream to the tagged system is Poisson, so the first system is $`M/GI/1/0`$, while the second one is denoted $`/GI/1/0`$, where $``$ in the first place of the notation stands for the input process in the second system, which is the output (loss) stream in the first one. Clearly, that an arriving customer is arranged to the second queueing system if and only if at the moment of his/her arrival the tagged system is occupied. Therefore, let us delete all the intervals when the tagged system is empty and merge the ends. In this case, the tagged system becomes an ordinary renewal process, and the stream of arrivals to the second queueing system becomes 1-Poisson rather then Poisson (because after deleting intervals and merging the ends in the new time scale the original Poisson process is transformed into 1-Poisson). Therefore the second system now can be re-denoted by $`\stackrel{~}{M}/GI/1/0`$, where $`\stackrel{~}{M}`$ in the first place of the notation stands for 1-Poisson input and replaces the initially written symbol $``$.
Thus, the $`M/GI/2/0`$ queueing system is decomposed into the $`M/GI/1/0`$ and $`\stackrel{~}{M}/GI/1/0`$ queueing systems. Clearly, that without loss of generality one can assume that the original arrival stream is 1-Poisson rather than Poisson, i.e. the original queueing system is $`\stackrel{~}{M}/GI/2/0`$, and it is decomposed into two $`\stackrel{~}{M}/GI/1/0`$ queueing systems. The last note is important for the further extension of the result for the systems $`M/GI/m/0`$ (or generally $`\stackrel{~}{M}/GI/m/0`$) having $`m>2`$ servers.
Let $`\tau `$ be the time moment when an arriving customer occupies the tagged server station. According to the ASTA property,
(3.7)
$$\mathrm{P}\{\zeta _{2,0}(\tau )x\}=\mathrm{P}\{\zeta _{2,0}(t)x\},$$
where $`t`$ is an arbitrary fixed point, and the probability distribution function of $`\zeta _{2,0}(t)`$ in this point coincides with the distribution of residual service time in specific $`\stackrel{~}{M}/GI/1/0`$ system with some specific value of parameter of 1-Poisson process, which is not important here. On the other hand, the process $`\zeta _{2,0}(t)`$ is stationary and Markov. Therefore from (3.7) for any $`h>0`$ we have
(3.8) $`\mathrm{P}\{\zeta _{2,0}(\tau +h)x\}`$ $`=\mathrm{P}\{\zeta _{2,0}(t+h)x\}=\mathrm{P}\{\zeta _{2,0}(t)x\}.`$
Let $`\chi `$ denotes the service time of the customer, who arrives at the time moment $`\tau `$ occupying the tagged server station. Our challenge is to prove that
(3.9)
$$\mathrm{P}\{\zeta _{2,0}(\tau +\chi )x\}=\mathrm{P}\{\zeta _{2,0}(t)x|\zeta _{1,0}(t)>0\}.$$
Instead of the original processes $`\zeta _{i,0}(t)`$, $`i=1,2`$, consider another processes $`\stackrel{~}{\zeta }_{i,0}(t)`$, which are obtained by deleting the intervals where the tagged server is free, and merging the ends. Then, $`\stackrel{~}{\zeta }_{1,0}(t)`$ is a renewal process, and the 1-stationary process $`\stackrel{~}{\zeta }_{2,0}(t)`$ and the random variable $`\chi `$ (the length of a service time in the tagged server that starts at moment $`\tau `$) are independent. Hence, for any event $`\{\chi =h\}`$ according to the properties of 1-stationary processes we have
(3.10) $`\mathrm{P}\{\stackrel{~}{\zeta }_{2,0}(\tau +\chi )x|\chi =h\}`$ $`=\mathrm{P}\{\stackrel{~}{\zeta }_{2,0}(\tau )x\},`$
and, due to the total probability formula from (3.10) we have
(3.11)
$$\mathrm{P}\{\stackrel{~}{\zeta }_{2,0}(\tau +\chi )x\}=\mathrm{P}\{\stackrel{~}{\zeta }_{2,0}(\tau )x\}.$$
The only difference between (3.11) and the basic property (3.3) is that the time moment $`\tau `$ is random, while $`t_0`$ is not. However keeping in mind (3.8), this modified equation (3.11) follows by the same derivation as in (3.3).
Hence, from (3.11),
(3.12)
$$\mathrm{P}\{\stackrel{~}{\zeta }_{2,0}(\tau +\chi )x\}=\mathrm{P}\{\stackrel{~}{\zeta }_{2,0}(\tau )x\}=\mathrm{P}\{\zeta _{2,0}(t)x|\zeta _{1,0}(t)>0\},$$
and since $`\mathrm{P}\{\stackrel{~}{\zeta }_{2,0}(\tau +\chi )x\}=\mathrm{P}\{\zeta _{2,0}(\tau +\chi )x\}`$ we finally arrive at (3.9).
As well, noticing that
$$\mathrm{P}\{\stackrel{~}{\zeta }_{2,0}(\tau )x\}=\mathrm{P}\{\zeta _{2,0}(\tau )x\},$$
from (3.11) and (3.7) we also have
(3.13)
$$\mathrm{P}\{\zeta _{2,0}(\tau +\chi )x\}=\mathrm{P}\{\zeta _{2,0}(\tau )x\}=\mathrm{P}\{\zeta _{2,0}(t)x\}.$$
Similarly to (3.13) one can prove
(3.14)
$$\mathrm{P}\{\zeta _{1,0}(\tau +\chi )x\}=\mathrm{P}\{\zeta _{1,0}(\tau )x\}=\mathrm{P}\{\zeta _{1,0}(t)x\},$$
where $`\tau `$ is the moment of arrival of a customer, who at this moment $`\tau `$ occupies the second server, and $`\chi `$ denotes his/her service time. Relations (3.14) can be proved with the aid of the same construction of deleting intervals and merging the ends, but now in the second server. So, combining (3.13) and (3.14) we arrive at the following fact. In any arrival or service completion time instant in one server, the residual service time in another server has the same stationary distribution.
This fact is used in the constructions below.
Now consider the stationary $`M/GI/2/0`$ queueing system, in which both servers are equivalent in the sense that if at the moment of arrival of a customer both servers are free, then a customer can occupy each of servers with the equal probability $`\frac{1}{2}`$. In this case the both of the processes $`\zeta _{1,0}(t)`$ and $`\zeta _{2,0}(t)`$ have the same distribution.
Let us delete the time intervals where the both servers are simultaneously free, and merge the corresponding ends (see Figure 3). The new processes are denoted by $`\stackrel{~}{\zeta }_{1,0}(t)`$ and $`\stackrel{~}{\zeta }_{2,0}(t)`$, and both of them have the same equivalent distribution. (We use the same notation as in the construction above believing that it is not confusing for readers.) This two-dimensional 1-stationary process characterizes the system where in any time $`t`$ at least one of two servers is busy. Consider the event of arrival of a customer in a stationary system at the moment when only one of two servers is busy. Let $`\tau ^{}`$ be the moment of this arrival, and let $`\tau ^{}`$ denote the moment of the first service completion in one of two servers following after the moment $`\tau ^{}`$. Then, at the endpoint $`\tau ^{}`$ of the interval $`[\tau ^{},\tau ^{})`$ the distribution of the residual service time will be the same as that at the moment $`\tau ^{}`$ (due to the established fact that at the end of a service completion in one server, the distribution of a residual service time in another server must coincide with the stationary distribution of a residual service time and due to the fact that both servers are equivalent.)
The additional details here are as follows. There can be different events associated with the points $`\tau ^{}`$ and $`\tau ^{}`$. For example, at the moment $`\tau ^{}`$ an arriving customer can be accepted by one of the servers, while the service completion at the moment $`\tau ^{}`$ can be either in the same server of in another server. If time moments $`\tau ^{}`$ and $`\tau ^{}`$ are associated with the same server (for example, the moment of service start and service completion in the first server) then we speak about residual service times in another server (in this example - the second server). If time moments $`\tau ^{}`$ and $`\tau ^{}`$ are associated with different servers (say, $`\tau ^{}`$ is the service start in the first server, but $`\tau ^{}`$ is the service completion in the second one), then we speak about residual service times in different servers (in this specific case we speak about residual service time in the second server at the time moment $`\tau ^{}`$ and residual service time in the first server at the time moment $`\tau ^{}`$). However, according to the earlier result, it does not matter which specific event of these mentioned occurs. The only fact, that the stationary distribution of a residual service time in a given server must be the same for all time moments of arrival and service completion occurring in another server and vice versa, is used.
Deleting the interval \[$`\tau ^{}`$, $`\tau ^{}`$) and merging the ends $`\tau ^{}`$ and $`\tau ^{}`$ (see Figure 4) we obtain the following structure of the 1-stationary process $`\widehat{𝐲}_{1,0}(t)`$.
In the points where idle intervals are deleted and the ends are merged we have renewal points: one of periods is finished and another is started. In the other points where the intervals of type \[$`\tau ^{}`$, $`\tau ^{}`$) are deleted and their ends are merged we have the points of ‘interrupted’ renewal processes. In this ‘interrupted’ renewal process the point $`\tau ^{}`$ is a point of 1-Poisson arrival, and, according to ASTA, the distribution in this point in the server that continue to serve a customer coincides with the stationary distribution of the residual service time. In the other point $`\tau ^{}`$, which is the point of a service completion, the distribution in this point in the server that continue to serve a customer coincides with the stationary distribution of a residual service time as well. Therefore, in the point of the interruption (which is a point of discontinuity) the residual service time distribution coincides with the stationary distribution of a residual service time, i.e. with the distribution of $`𝐱_1(t)`$. (Notice, that the intervals of type \[$`\tau ^{},\tau ^{}`$) are an analogue of the intervals \[$`s_{1,k},t_{1,k}`$) considered in the Markovian case in Section 2.)
By amalgamating the residual service times of the first and second servers given in the lower graph in Figure 4, one can built a typical one-dimensional 1-stationary process $`\widehat{𝐲}_{1,0}(t)`$, the limiting stationary distribution of which coincides with that of $`𝐱_1(t)`$. (see Figure 5).
Therefore the processes $`\widehat{𝐲}_{1,0}(t)`$ and $`𝐱_1(t)`$ have the identical one-dimensional distribution, and relation (3.6) follows. ∎
Let us develop Property 3.4 to the case $`m=3`$ and then to the case of an arbitrary $`m>1`$ for the $`M/GI/m/0`$ queueing systems. Namely, we have the following.
###### Property 3.5.
For the $`M/GI/m/0`$ queueing system we have:
(3.15)
$$\mathrm{P}\{\widehat{𝐲}_{m1,0}(t)B_{m1}\}=\mathrm{P}\{𝐱_{m1}(t)B_{m1}\},$$
where $`B_{m1}`$ is an arbitrary Borel set of $`^{m1}`$.
###### Proof.
The proof will be concentrated in the case $`m=3`$ for the 1-stationary process $`\widehat{𝐲}_{2,0}(t)`$, which is associated with the paths of the $`M/GI/3/0`$ queueing system where only two servers are busy. Then the result will be concluded for an arbitrary $`m2`$ by induction.
Prior studying this case, we first study the specific case of the $`M/GI/2/0`$ queueing system by considering the paths when the both servers are busy. Then using the arguments of the proof of Property 3.4 enables us to extend that specific result related to the $`M/GI/2/0`$ queueing systems to the 1-stationary process $`\widehat{𝐲}_{2,0}(t)`$ of the $`M/GI/3/0`$ queueing system.
As in the proof of Property 3.4 in the specific case of the $`M/GI/2/0`$ queueing system considered here, we will study the stationary one-dimensional processes $`\zeta _{1,0}(t)`$ and $`\zeta _{2,0}(t)`$. However the idea of the present proof generally differs from that of the proof of Property 3.4. Here we do not call the first (or second) server a tagged server station to use decomposition. We simply use the fact established in the proof of Property 3.4 that at the moment of arrival or service completion of a customer in one server, the distribution of a residual service time in another server will coincide with the stationary distribution of a residual service time in this server. (The same idea has been used in the proof of Property 3.4.)
The present proof explicitly uses the fact that the class of 1-stationary processes is algebraically closed with respect to the operations of deleting intervals and merging the ends, which was mentioned before.
Let us delete the idle intervals of the process $`\zeta _{1,0}(t)`$ and merge the ends. Then we get a stationary renewal process as in the above case $`m=1`$ (Property 3.3).
After deleting the same time intervals in the second stationary process $`\zeta _{2,0}(t)`$ and merging the ends, the process will be transformed as follows. Let $`t^{}`$ be a moment of 1-Poisson arrival when a customer occupies the first server. (Recall that owing to the known properties of 1-Poisson process, the stream of arrival to each of $`i`$ servers ($`i=1,2`$) is 1-Poisson.) Then, according to the ASTA property, $`\zeta _{2,0}(t^{})=\zeta _{2,0}(t)`$ in distribution. Therefore after deleting all of the idle intervals of the second server and merging the ends, after the first time scaling (i.e. removing corresponding time intervals, see Figure 6) instead of the initial 1-stationary process $`\zeta _{2,0}(t)`$ we obtain the new 1-stationary process with the equivalent one-dimensional distribution. This process is denoted by $`\stackrel{~}{\zeta }_{2,0}(t)`$.
Notice, that the process $`\stackrel{~}{\zeta }_{2,0}(t)`$ is obtained from the process $`\zeta _{2,0}(t)`$ by constructing a sequence of 1-stationary processes described above.
Then we have the two-dimensional process the first component of which is $`𝐱_1(t)`$ and the second one is $`\stackrel{~}{\zeta }_{2,0}(t)`$. For our convenience this first component is provided with upper index, and the two-dimensional vector looks now as $`\{𝐱_1^{(1)}(t),\stackrel{~}{\zeta }_{2,0}(t)\}`$.
Let us repeat the above procedure, deleting the remaining idle intervals of the second server and merging the ends. We get the 1-stationary process being equivalent in the distribution to the stationary renewal process $`𝐱_1(t)`$, which is denoted now $`𝐱_1^{(2)}(t)`$.
Upon this (final) time scaling the first process $`𝐱_1^{(1)}(t)`$ is transformed as follows. Let $`t^{}`$ be a random point of 1-Poisson arrival when the second server is occupied. Applying the ASTA property once again, for the first component of the process we obtain that $`𝐱_1^{(1)}(t^{})`$ coincides in one-dimensional distribution with $`𝐱_1^{(1)}(t)`$. Thus, after deleting the entire idle intervals and merging the ends, we finally obtain the two-dimensional process $`\{𝐱_1^{(1)}(t),𝐱_1^{(2)}(t)\}`$ each component of which has the same one-dimensional distribution as this of the process $`𝐱_1(t)`$. The dynamic of this time scaling is shown in Figure 7.
For our further purpose, the independence of the processes $`𝐱_1^{(1)}(t)`$ and $`𝐱_1^{(2)}(t)`$ is needed. The constructions in this paper enables us to prove this independence. However, the independence of $`𝐱_1^{(1)}(t)`$ and $`𝐱_1^{(2)}(t)`$ follows automatically from the known results by Takács and the easiest way is to follow a result of that paper. Namely, it follows from formulae (6) and (7) on page 72, that the joint conditional stationary distribution of residual service times given that $`k`$ servers are busy coincides with the stationary distribution of $`𝐱_k(t)`$, which in turn is the product of the stationary distributions of $`𝐱_1(t)`$. In particular,
(3.16)
$$\mathrm{P}\left\{𝐱_1^{(1)}(t)x_1,𝐱_1^{(2)}(t)x_2\right\}=\mathrm{P}\left\{𝐱_1(t)x_1\right\}\mathrm{P}\left\{𝐱_1(t)x_2\right\}.$$
Now, using the arguments of the proof of Property 3.4 one can easily extend the result obtained now for $`M/GI/2/0`$ queueing system to the $`M/GI/3/0`$ queueing system, and thus prove (3.15) for the $`M/GI/3/0`$ queueing system.
Similarly to the proof of Property 3.4, let us introduce the processes $`\zeta _{1,0}(t)`$, $`\zeta _{2,0}(t)`$ and $`\zeta _{3,0}(t)`$ of the residual service times in the first, second and third servers correspondingly. These processes all are assumed to have the same stationary distribution of residual times, which respects to the scheme where an arriving customer occupies one of available free servers with equal probability. Let us call a server of the $`M/GI/3/0`$ queueing system that occupied at the moment of busy period start a tagged server station. So, we decompose the original system into the $`\stackrel{~}{M}/GI/2/0`$ and tagged queueing system $`\stackrel{~}{M}/GI/1/0`$. However, it is shown above that $`\stackrel{~}{M}/GI/2/0`$ can be decomposed into two $`\stackrel{~}{M}/GI/1/0`$ queueing systems, where after the procedure of deleting idle intervals and merging the ends we obtain the process having the same stationary distribution as that of the process $`𝐱_2(t)`$. This stationary distribution remains the same in all random points of arrivals and service completions in the tagged service station. So, after deleting intervals and merging the ends in the tagged service station, in a new time scaling we arrive at the same stationary distribution as that of the process $`𝐱_3(t)`$. So, the result for $`m=3`$ follows.
This induction becomes clear for an arbitrary $`m2`$ as well, where the original $`\stackrel{~}{M}/GI/m/0`$ system can be decomposed into the $`\stackrel{~}{M}/GI/m1/0`$ queueing system and a tagged server station $`\stackrel{~}{M}/GI/1/0`$. ∎
Now we will establish a connection between the processes $`\widehat{𝐲}_{m1,n}(t)`$, $`\widehat{𝐲}_{m1,0}(t)`$ and $`𝐱_{m1}(t)`$. We start from the case $`m=2`$.
###### Property 3.6.
Under the assumption that the probability distribution function $`G(x)`$ belongs to the class NBU (NWU) we have
(3.17)
$$\mathrm{P}\{\widehat{𝐲}_{1,n}(t)x\}(\text{resp.})\mathrm{P}\{𝐱_1(t)x\}.$$
###### Proof.
Along with the 1-stationary processes $`\widehat{𝐲}_{1,n}(t)`$, let us introduce another 1-stationary processes $`\widehat{𝐘}_{2,n}(t)`$. This last process is related to the same $`\stackrel{~}{M}/GI/2/n`$ queueing system as the process $`\widehat{𝐲}_{1,n}(t)`$, and is obtained by deleting time intervals when there are more than two or less than two customers in the system, and merging the ends.
Using the same arguments as in the proof of Property 3.5 one can prove that the components of this process are generated by independent 1-stationary processes and having the same distribution. Indeed, involving as earlier in the proof of Property 3.4 the processes $`\zeta _1(t)`$ and $`\zeta _2(t)`$ having the same distribution, one can delete intervals where the system is empty and merge the ends. Apparently, the new processes $`\stackrel{~}{\zeta }_1(t)`$ and $`\stackrel{~}{\zeta }_2(t)`$ obtained after this procedure have the same stationary distribution. (However, it is shown later that the limiting stationary distribution of these one-dimensional processes differs from such the distribution obtained after the similar procedure for the $`\stackrel{~}{M}/GI/2/0`$ system, and, therefore, its one-dimensional distribution distinguishes from that of the process $`𝐱_1(t)`$.)
Let us go back to the initial process $`𝐲_{2,n}(t)`$, $`n1`$, to delete the time intervals where the both servers are free and merge the corresponding ends. We also remove the last component corresponding to the queue-length $`Q_{2,n}(t)`$. (The exact value of the queue-length $`Q_{2,n}(t)`$ is irrelevant here and is not used in our analysis.) In the new time scale we obtain the two-dimensional process $`\stackrel{~}{𝐲}_{2,n}(t)`$.
Similarly to the proof of relation (3.6) we have time moments $`\tau ^{}`$ and $`\tau ^{}`$. The first of them is a moment of arrival of a customer at the system with one busy and one free server, and the second one is the following after $`\tau ^{}`$ moment of service completion of a customer when there remain one busy server only. The time interval \[$`\tau ^{}`$, $`\tau ^{}`$) is an orbital busy period. (The concept of orbital busy period is defined in Section 2 for Markovian systems. For $`M/GI/m/n`$ queueing systems this concept is the same.) It can contain queueing customers waiting for their service. Let $`t^{\mathrm{begin}}`$ be a moment of arrival of a customer during the orbital busy period \[$`\tau ^{}`$, $`\tau ^{}`$) who occupies a waiting place, and let $`t^{\mathrm{end}}`$ be the following after $`t^{\mathrm{begin}}`$ moment of time when after the service completion the queue space becomes empty again. A period of time \[$`t^{\mathrm{begin}}`$, $`t^{\mathrm{end}}`$) is called queueing period. (Note that for the M/GI/$`m`$/$`n`$ queueing system, the intervals of type \[$`\tau ^{}`$, $`\tau ^{}`$) are an analogue of the intervals of type \[$`s_{m1,k},t_{m1,k}`$) in the Markovian queueing system $`M/M/m/n`$, and the intervals of type \[$`t^{\mathrm{begin}},t^{\mathrm{end}}`$) are an analogue of the intervals of type \[$`s_{m,k},t_{m,k}`$) in that Markovian queueing system $`M/M/m/n`$.)
All customers of queueing periods, i.e. those arrived during orbital busy period can be considered as customers arriving in a tagged server station. At the moment of $`t^{\mathrm{begin}}`$, which is an instant of a Poisson arrival, the two-dimensional distribution of the random vector $`\stackrel{~}{𝐲}_{2,n}(t^{\mathrm{begin}})`$ coincides with the stationary distribution of the random vector $`\widehat{𝐘}_{2,n}(t)`$. However, in the point $`t^{\mathrm{end}}`$, the probability distribution of $`\stackrel{~}{𝐲}_{2,n}(t^{\mathrm{end}})`$ is different from the stationary distribution of $`\widehat{𝐘}_{2,n}(t)`$, because this specific time instant $`t^{\mathrm{end}}`$ coincides with a service beginning in one of servers of the main system. Therefore, deleting the interval \[$`t^{\mathrm{begin}},t^{\mathrm{end}}`$) and merging the end leads to the change of the distribution.
More specifically, at the time instant $`t^{\mathrm{end}}`$ one of the components of the vector $`\widehat{𝐘}_{2,n}(t)`$, say the first one, is a random variable having the probability distribution $`G(x)`$. Then, another component, i.e. the second one, because of the aforementioned properties of 1-stationary processes, coincides in distribution with $`\stackrel{~}{\zeta }_{1,n}`$ (or $`\stackrel{~}{\zeta }_{2,n}`$), which is a component of the stationary process $`\widehat{𝐘}_{2,n}(t)`$. Indeed, let customers arriving in a busy system and waiting in the queue be assigned to the tagged server station. At the moment of 1-Poisson arrival $`t^{\mathrm{begin}}`$ of a customer in the tagged server station, the two-dimensional Markov process associated with the main queueing system has the same distribution as the vector $`\widehat{𝐘}_{2,n}(t)`$, i.e. two-dimensional distribution coinciding with the joint distribution of $`\stackrel{~}{\zeta }_{1,n}`$ and $`\stackrel{~}{\zeta }_{2,n}`$. Then, at the moment of the service completion $`t^{\mathrm{end}}`$, which coincides with the moment of service completion in one of two servers, the probability distribution function of the residual service time in another server, where the service is being continued, coincides with the distribution of a component of the vector $`\widehat{𝐘}_{2,n}(t)`$, i.e. with the distribution of $`\stackrel{~}{\zeta }_{1,n}`$.
If the probability distribution function $`G(x)`$ belongs to the class NBU, then the 1-stationary process $`\stackrel{~}{𝐲}_{2,n}(t)`$ satisfies the property $`\stackrel{~}{𝐲}_{2,n}(t^{\mathrm{begin}})_{st}\stackrel{~}{𝐲}_{2,n}(t^{\mathrm{end}})`$. If $`G(x)`$ belongs to the class NWU, then the opposite inequality holds: $`\stackrel{~}{𝐲}_{2,n}(t^{\mathrm{end}})_{st}\stackrel{~}{𝐲}_{2,n}(t^{\mathrm{begin}})`$. (The stochastic inequality between vectors means the stochastic inequality between their corresponding components.)
The above stochastic inequalities are between random values of the process $`\stackrel{~}{𝐲}_{2,n}(t)`$ in the different time instants $`t^{\mathrm{begin}}`$ and $`t^{\mathrm{end}}`$. Our further task is to compare two different processes $`\stackrel{~}{𝐲}_{2,n}(t)`$ and $`\stackrel{~}{𝐲}_{2,0}(t)`$. The first of these processes is associated with the $`M/GI/m/n`$ queueing system, while the second one is associated with the $`M/GI/m/0`$ queueing system. The idea of comparison is very simple. Suppose that both queueing system are started at zero, i.e. consider the paths of these system when the both of them are not in steady state, and compare the Markov processes associated with these system. For the non-stationary processes we will use the same notation $`\stackrel{~}{𝐲}_{2,n}(t)`$ and $`\stackrel{~}{𝐲}_{2,0}(t)`$ understanding that it is spoken about usual (not stationary) Markov processes. The notation for time moments such as $`t^{\mathrm{begin}}`$ and $`t^{\mathrm{end}}`$ is now associated with these usual (i.e. non-stationary) processes as well. We will consider the Markov processes associated with $`M/GI/2/n`$ and $`M/GI/2/0`$ queueing systems on the same probability space. In the time interval $`[0,t^{\mathrm{begin}})`$ the paths of the Markov processes $`\stackrel{~}{𝐲}_{2,n}(t)`$ and $`\stackrel{~}{𝐲}_{2,0}(t)`$ coincide ($`n0`$). However, after deleting the interval $`[t^{\mathrm{begin}},t^{\mathrm{end}})`$ and merging the ends, then in the end point $`t^{\mathrm{begin}}`$ the values of the processes $`\stackrel{~}{𝐲}_{2,n}(t)`$ and $`\stackrel{~}{𝐲}_{2,0}(t)`$ will be different. Indeed, in the case of the process $`\stackrel{~}{𝐲}_{2,0}(t)`$, which is associated with the $`M/GI/m/0`$ queueing system, $`t^{\mathrm{begin}}`$ and $`t^{\mathrm{end}}`$ is the same point, and the value of Markov processes will be the same after replacing the points $`t^{\mathrm{begin}}`$ with $`t^{\mathrm{end}}`$. However, in the case of the process $`\stackrel{~}{𝐲}_{2,n}(t)`$ associated with $`M/GI/m/n`$ queueing system, the values in these points will be different with probability 1, and consequently, because of the inequality $`\stackrel{~}{𝐲}_{2,n}(t^{\mathrm{begin}})_{st}\stackrel{~}{𝐲}_{2,n}(t^{\mathrm{end}})`$ we have $`\stackrel{~}{𝐲}_{2,0}(t^{\mathrm{begin}})_{st}\stackrel{~}{𝐲}_{2,n}(t^{\mathrm{begin}})`$ in the case when $`G(x)`$ belongs to the class NBU. If $`G(x)`$ belongs to the class NWU, we have the opposite inequality: $`\stackrel{~}{𝐲}_{2,n}(t^{\mathrm{begin}})_{st}\stackrel{~}{𝐲}_{2,0}(t^{\mathrm{begin}})`$.
Therefore, after deleting all the intervals of the type $`[t^{\mathrm{begin}},t^{\mathrm{end}})`$ from the original Markov process we obtain new Markov process, and in the case when $`G(x)`$ belongs either to the class NBU or to the class NWU one can apply the theorem of Kalmykov (see also ) to compare these two Markov processes. In the case where $`G(x)`$ belongs to the class NBU, all the path of the Markov process, associated with $`M/GI/2/n`$ is not smaller (in stochastic sense) than that path of the Markov process, associated with $`M/GI/2/0`$. If $`G(x)`$ belongs to the class NWU, then the opposite stochastic inequality holds between two different Markov processes. Apparently, the same stochastic inequalities remain correct if we speak about stationary Markov processes. Nothing is changed if we let $`t`$ to increase to infinity and arrive at stationary distributions. So, under the assumption that $`G(x)`$ belongs to the class NBU, for the stationary processes we obtain $`\widehat{𝐲}_{1,0}(t)_{st}\widehat{𝐲}_{1,n}(t)`$. In other words, due to the fact that $`\widehat{𝐲}_{1,0}(t)=_{st}𝐱_1(t)`$, we obtain that $`𝐱_1(t)_{st}\widehat{𝐲}_{1,n}(t)`$. In the case where $`G(x)`$ belongs to the class NWU, the opposite inequality holds.
The arguments of the proof given for $`m=2`$ remain correct for an arbitrary $`m2`$. The proof given by induction uses decomposition of the original system into the main system and a tagged server station as above. The further arguments for stochastic comparison of Markov processes are also easily extended for the case of an arbitrary $`m2`$. ∎
From the above results for the Markov processes the statement of the lemma follows. The stochastic inequalities between $`T_{m,n}(m1)`$ and $`T_{m,0}(m1)`$ follow by the coupling arguments. The lemma is completely proved. ∎
## 4. Theorems on losses in $`M/GI/m/n`$ queueing systems
The results obtained in the previous section enable us to establish theorems for the number of losses in $`M/GI/m/n`$ queueing systems during their busy periods.
###### Theorem 4.1.
Under the assumption $`\lambda =m\mu `$, the expected number of losses during a busy period of the $`M/GI/m/n`$ queueing system is the same for all $`n1`$.
###### Proof.
Consider the system $`M/GI/m/n`$ under the assumption $`\lambda =m\mu `$, and similarly to the construction in the proof of Lemma 3.1 let us delete all the intervals where the number of customers in the system is less than $`m`$, and merge the corresponding ends. The process obtained is denoted $`\widehat{𝐲}_{m,n}(t)`$. This is the 1-stationary process of orbital busy periods.
The stationary departure process, together with the arrival 1-Poisson process of rate $`\lambda `$ and the number of waiting places $`n`$ describes the stationary $`M/G/1/n`$ queue-length process (with generally dependent service times). As soon as a busy period is finished (in our case it is an orbital busy period, see Section 2 for the definition), the system immediately starts a new busy period by attaching a new customer into the system. This unusual situation arises because of the construction of the process. There are no idle periods, and servers all are continuously busy. Thus, the busy period, which is considered here, is one of the busy periods attached one after another.
Let $`T`$ be a large period of time, and during that time there are $`K(T)`$ busy periods of the $`M/G/1/n`$ queueing system (which does not contain idle times as mentioned). Let $`L(T)`$ and $`\nu (T)`$ denote the number of lost and served customers during time $`T`$. We have the formula:
(4.1)
$$\underset{T\mathrm{}}{lim}\frac{1}{\mathrm{E}K(T)}\left(\mathrm{E}L(T)+\mathrm{E}\nu (T)\right)=\underset{T\mathrm{}}{lim}\frac{1}{\mathrm{E}K(T)}\left(\lambda T+\mathrm{E}K(T)\right),$$
the proof of which is given below.
Relationship (4.1) has the following explanation. The left-hand side term $`\mathrm{E}L(T)+\mathrm{E}\nu (T)`$ is the expectation of the number of lost customers plus the expectation of the number of served customers during time $`T`$, and the right-hand side term $`\lambda T+\mathrm{E}K(T)`$ is the expectation of the number of arrivals during time $`T`$ plus the expected number of attached customers.
Relationship (4.1) can be proved by renewal arguments as follows.
There are $`m`$ independent copies $`𝐱_1^{(1)}(t)`$, $`𝐱_1^{(2)}(t)`$, …, $`𝐱_1^{(m)}(t)`$ of the stationary renewal process, which model the process $`\widehat{𝐲}_{m,n}(t)`$. (In fact, we have $`m`$ 1-stationary processes, which have the same distributions as $`m`$ independent renewal processes $`𝐱_1^{(1)}(t)`$, $`𝐱_1^{(2)}(t)`$, …, $`𝐱_1^{(m)}(t)`$.) Let $`1im`$, and let $`C_1`$, $`C_2`$,…$`C_{K_i(T)}`$ be such points of busy period starts associated with the renewal process $`𝐱_1^{(i)}(t)`$ (one of those $`m`$ independent and identically distributed renewal processes), where $`K_i(T)`$ denotes the total number of these regeneration point indexed by $`i`$. Denote also by $`z_1`$, $`z_2`$, …, $`z_{K_i(T)}`$ the corresponding lengths of busy periods, by $`\mathrm{}_1`$, $`\mathrm{}_2`$,…,$`\mathrm{}_{K_i(T)}`$ the corresponding number of losses during these $`K_i`$ busy periods, and by $`n_1`$, $`n_2`$,…,$`n_{K_i(T)}`$ the corresponding number of served customers during these busy periods. Let $`T_i=z_1+z_2+\mathrm{}+z_{K_i(T)}`$, let $`L_i(T)=\mathrm{}_1+\mathrm{}_2+\mathrm{}+\mathrm{}_{K_i(T)}`$ and let $`\nu _i(T)=n_1+n_2+\mathrm{}+n_{K_i(T)}`$.
Since at the moments $`C_1`$, $`C_2`$,…$`C_{K_i(T)}`$ of the busy period starts the distribution of the above stationary Markov process of residual times is the same, then the numbers of losses $`\mathrm{}_1`$, $`\mathrm{}_2`$,…,$`\mathrm{}_{K_i(T)}`$ and, respectively, the numbers of served customers $`n_1`$, $`n_2`$,…,$`n_{K_i(T)}`$ during each of these busy periods have the same distributions, and one can apply the renewal reward theorem.
By the renewal reward theorem we have:
(4.2)
$$\underset{T\mathrm{}}{lim}\frac{1}{m\mathrm{E}K_i(T)}\left(\mathrm{E}L_i(T)+\mathrm{E}\nu _i(T)\right)=\underset{T\mathrm{}}{lim}\frac{1}{m\mathrm{E}K_i(T)}\left(\lambda \mathrm{E}T_i+m\mathrm{E}K_i(T)\right).$$
Taking into account that
$$\underset{T\mathrm{}}{lim}\frac{\mathrm{E}K(T)}{\mathrm{E}K_i(T)}=m,$$
$$\underset{T\mathrm{}}{lim}\frac{\mathrm{E}L(T)}{\mathrm{E}L_i(T)}=1,$$
$$\underset{T\mathrm{}}{lim}\frac{\mathrm{E}\nu (T)}{\mathrm{E}\nu _i(T)}=1,$$
and
$$\underset{T\mathrm{}}{lim}\frac{\mathrm{E}T_i}{T}=1,$$
because of the correspondence between the left- and right-hand sides, from (4.2) we arrive at (4.1). Thus, we bypass the fact that the times between departures are dependent, and thus (4.1) is actually obtained by application of the renewal reward theorem by a usual manner, like in the case of independent times between departures (see e.g. Ross , Karlin and Taylor ).
Together with (4.1) we have
(4.3)
$$\underset{T\mathrm{}}{lim}\frac{1}{\mathrm{E}K(T)}\mathrm{E}\nu (T)=\underset{T\mathrm{}}{lim}\frac{1}{\mathrm{E}K(T)}m\mu T.$$
Let us now introduce the following notation. Let $`\zeta _n`$ denote the length of an orbital busy period, and let $`L_n`$ and $`\nu _n`$ correspondingly denote the number of lost and served customers during that orbital busy period. Using the arguments of , we prove that $`\mathrm{E}L_n=1`$ for all $`n1`$.
Indeed, from (4.1) and (4.3) we have the equations:
(4.4)
$$\mathrm{E}L_n+\mathrm{E}\nu _n=\lambda \mathrm{E}\zeta _n+1,$$
(4.5)
$$\mathrm{E}\nu _n=m\mu \mathrm{E}\zeta _n.$$
The substitution $`\lambda =m\mu `$ into the system of equations (4.4) and (4.5) yields $`\mathrm{E}L_n=1`$.
Hence, during an orbital busy period there is exactly one lost customer in average for any $`n0`$. To finish the proof we need in a deeper analysis. First, we should find the expected number of queueing periods during one orbital busy period. For this purpose, one can use the similar construction by deleting all the intervals when the number of customers in the system is not greater than $`m`$, and merge the corresponding ends. The obtained process is denoted $`\widehat{𝐲}_{m+1,n}(t)`$, and this is one stationary process of queueing periods following one after another.
The structure of the process $`\widehat{𝐲}_{m+1,n}(t)`$ is similar to that of the process $`\widehat{𝐲}_{m,n}(t)`$. The process $`\widehat{𝐲}_{m+1,n}(t)`$ describes the stationary $`M/G/1/n1`$ queueing system, the service times of which are generally dependent. As soon as one busy period in this system is finished, a new customer starting a new busy period is immediately attached into the system. Thus, the only difference between the processes $`\widehat{𝐲}_{m,n}(t)`$ and $`\widehat{𝐲}_{m+1,n}(t)`$ is that the numbers of waiting places differ by value of parameter $`n`$. Therefore, using the similar notation and arguments, one arrive at the conclusion that the expected number of losses per queueing period is equal to 1 as well. Therefore, in long-run period of time, the number of queueing periods and orbital busy periods is the same in average. So, there is exactly one queueing period per orbital busy period in average.
Therefore, during a long-run period the number of events that the different Markov processes $`\widehat{𝐲}_{m1,n}(t)`$ change their values after deleting queueing periods and merging the ends (as exactly explained in the proof of Lemma 3.1) is the same in average for all $`n1`$, and the stationary characteristics of all of these Markov processes $`\widehat{𝐲}_{m1,n}(t)`$, given for different values $`n`$=1,2,…, are the same. Hence, the expectation $`\mathrm{E}T_{m,n}(m1)`$ is the same for all $`n`$=1,2,…as well. (Recall that $`T_{m,n}(m1)`$ denote the total time during a busy period when $`m1`$ servers are occupied.)
Hence, using Wald’s identity connecting $`\mathrm{E}T_{m,n}(m1)`$ with $`\mathrm{E}L_{m,n}`$ (the expected number of losses during a busy period) we arrive at the desired result, since $`\mathrm{E}T_{m,n}(m1)`$ and the expectation of the number of orbital busy periods during a busy period of $`M/GI/m/n`$ both are the same for all $`n1`$. ∎
Application of Lemma 3.1 and the arguments of Theorem 4.1 enables us to prove the following result.
###### Theorem 4.2.
Let $`\lambda =m\mu `$. Then, under the assumption that $`G(x)`$ belongs to the class NBU, for the number of losses in $`M/GI/m/n`$ queueing systems, $`n1`$, we have
(4.6)
$$\mathrm{E}L_{m,n}=\frac{cm^m}{m!},$$
where the constant $`c1`$ depends on $`m`$ and the probability distribution $`G(x)`$ but is independent of $`n`$.
Under the assumption that $`G(x)`$ belongs to the class NWU we have (4.6) but with constant $`c1`$.
###### Proof.
Notice first, that for the expected number of losses in $`M/GI/m/0`$ queueing systems we have
$$\mathrm{E}L_{m,0}=\frac{m^m}{m!}$$
This result follows immediately from the Erlang-Sevastyanov formula , so that the expected number of losses during a busy period of the $`M/GI/m/0`$ queueing system is the same that this of the $`M/M/m/0`$ queueing system. The expected number of losses during a busy period of the $`M/M/m/0`$ queueing system, $`\mathrm{E}L_{m,0}=\frac{m^m}{m!}`$, is also derived in Section 2.
In the case where $`G(x)`$ belongs to the class NBU according to Lemma 3.1 we have $`\mathrm{E}T_{m,n}(m1)\mathrm{E}T_{m,0}(m1)`$, and therefore, the expected number of orbital busy periods in the $`M/GI/m/n`$ queueing system ($`n1`$) is not smaller than this in the $`M/GI/m/0`$ queueing system. Therefore, repeating the proof of Theorem 4.1 leads to the inequality $`\mathrm{E}L_{m,n}\mathrm{E}L_{m,0}`$ and consequently to the desired result. If $`G(x)`$ belongs to the class NWU, then we have the opposite inequalities, and finally the corresponding result stated in the formulation of the theorem. ∎
## 5. Batch arrivals
The case of batch arrivals is completely analogous to the case of ordinary (non-batch) arrivals. In the case of a Markovian $`M^X/M/m/n`$ queueing system one can also apply the level-crossing method to obtain equations analogous to (2.6)-(2.13). The same arguments as in Sections 3 and 4 in an extended form can be used for $`M^X/GI/m/0`$ queueing systems. |
warning/0506/hep-ph0506250.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The experimental and theoretical status of $`\mathrm{\Theta }^+`$-pentaquark remains controversial , . The QCD sum rules (SRs) have shown to be a very powerful tool for the investigation of the properties of conventional and exotic multiquark hadronic states . Several attempts to describe the properties of $`\mathrm{\Theta }^+`$ pentaquark using SRs have appeared . However, these calculations were either restricted to low dimension operators or they used interpolating currents which did not have the most suitable quantum numbers to project onto the $`\mathrm{\Theta }^+`$ .
The dynamics associated with the instanton, the ’t Hooft interaction, has been successful in understanding the spectroscopy of four-quark and H-dibaryon states and it was important for the spectroscopy of the pentaquark . Moreover, instantons are crucial for understanding chiral symmetry breaking in the strong interactions and lie at the basis of chiral soliton model for baryons, which has predicted the pentaquarks and their peculiar properties, e.g. small widths and masses .
In the SR calculations thus far, the contribution from so-called direct instantons, has not been investigated <sup>1</sup><sup>1</sup>1The direct instanton contribution of ref. should vanish due to the Pauli principle for the quarks in the instanton field.. It is well known that direct instantons play an important role in the SRs calculations to determine the properties of the pseudoscalar mesons and the nucleon octet baryons . We showed, in a mixed model-SR calculation, that they might be also important for the pentaquarks .
We perform a calculation for the pentaquark SRs which takes into account operators up to dimension $`d=13`$ and direct instanton contributions, and leads to evidence for a positive parity state whose mass is close to the observed mass.
## 2 The standard OPE contribution to the sum rules
The QCD sum rule approach starts from the correlator of some relevant current,
$$\mathrm{\Pi }(q^2)=id^4xe^{iqx}0|T\eta _\mathrm{\Theta }(x)\overline{\eta }_\mathrm{\Theta }(0)|0=\widehat{q}\mathrm{\Pi }_1(q^2)+\mathrm{\Pi }_2(q^2).$$
(1)
Here $`\eta _\mathrm{\Theta }`$ represents a current with non vanishing projection onto the pentaquark state. We use the conventional notation $`\widehat{q}=\gamma q`$.
The use of the narrow resonance approximation,
$$\mathrm{Im}\mathrm{\Pi }(q^2)=\pi \lambda _\mathrm{\Theta }^2(\widehat{q}+M_\mathrm{\Theta })\delta (q^2M_\mathrm{\Theta }^2)+\theta (q^2s_0^2)[\widehat{q}\mathrm{Im}\mathrm{\Pi }_1(q^2)+\mathrm{Im}\mathrm{\Pi }_2(q^2)],$$
(2)
where $`M_\mathrm{\Theta }`$ is the mass of the pentaquark, $`\lambda _\mathrm{\Theta }`$ its residue, $`s_0`$ the threshold, and the appropriate dispersion relations lead to the so-called chirality even
$$\frac{1}{\pi }_0^{s_0^2}𝑑s^2e^{s^2/M^2}\mathrm{Im}\mathrm{\Pi }_1^{OPE}(s^2)=\lambda _\mathrm{\Theta }^2e^{M_\mathrm{\Theta }^2/M^2},$$
(3)
and chirality odd
$$\frac{1}{\pi }_0^{s_0^2}𝑑s^2e^{s^2/M^2}\mathrm{Im}\mathrm{\Pi }_2^{OPE}(s^2)=\lambda _\mathrm{\Theta }^2M_\mathrm{\Theta }e^{M_\mathrm{\Theta }^2/M^2}$$
(4)
sum rules.
Our choice of current in the pentaquark correlator is
$$\eta _\mathrm{\Theta }^A=\frac{1}{4\sqrt{2}}ϵ_{afg}ϵ_{abc}ϵ_{bde}[(u_d^TCd_e)\gamma _5C\overline{s}_c^T][u_f^TC\gamma _5d_g],$$
(5)
whose structure corresponds to the A–state of refs. , which consists of the product of a scalar $`ud\overline{s}`$ triquark and a pseudoscalar $`ud`$–diquark. It can be easily seen, that this current has the same structure as that of ref. except for the $`\gamma _5`$ in front of the strange quark field. The consequence of this similarity is that for the chirality odd sum rule our results become identical to their results, if we restrict the calculation to low dimension operators, take into account our different normalization and an additional negative sign due to negative intrinsic parity of our current, Eq. (5).
We have also considered the current,
$$\eta _\mathrm{\Theta }^B=\frac{1}{4\sqrt{6}}ϵ_{acd}[(u_a^TC\gamma _\mu d_b+u_b^TC\gamma _\mu d_a)\gamma _5\gamma ^\mu C\overline{s}_b^T][u_c^TC\gamma _5d_d],$$
(6)
which corresponds to the B–state of refs. and contains a vector $`ud\overline{s}`$ triquark and a scalar $`ud`$–diquark. This current has also negative intrinsic parity. However, our analysis has shown that the B current coupling with $`\mathrm{\Theta }^+`$ is weak and therefore no definitive conclusion about the values of the mass and residue can be drawn from the consideration of its correlator. We will therefore not discuss it further.
Let us proceed with the analysis of the chirality odd SR for the pentaquark, Eq. (4). This SR is directly related to the mass of the state and usually is more stable than the chirality even SR for the ground state baryons and triquark $`ud\overline{s}`$ states .
After Borel transforming the pentaquark correlator has dimension $`d=13`$. The calculation of the chirality odd sum rule will be performed taking into account operators up to dimension $`d=13`$. To this order we obtain good stability in the OPE result. The operators of higher dimension ($`d>13`$) appear in the sum rule multiplied by inverse powers of the Borel mass, thus, in the interesting Borel mass region, their contribution is small and can be safely neglected.
With our interpolating current, the relevant trace to the chirality odd sum rule is expressed as
$`\mathrm{Tr}0|T\eta _\mathrm{\Theta }(x)\overline{\eta }_\mathrm{\Theta }(0)|0_{odd}=(i)^5{\displaystyle \frac{1}{32}}ϵ_{abc}ϵ_{bde}ϵ_{afg}ϵ_{a^{}b^{}c^{}}ϵ_{b^{}d^{}e^{}}ϵ_{a^{}f^{}g^{}}\mathrm{Tr}[\gamma _5CS_{c^{}c}^{s,T}(x)C\gamma _5]`$
$`\times (\mathrm{Tr}[CS_{dd^{}}^{u,T}CS_{ee^{}}^d]\mathrm{Tr}[CS_{gg^{}}^{d,T}C\gamma _5S_{ff^{}}^u\gamma _5]+\mathrm{Tr}[CS_{df^{}}^{u,T}CS_{eg^{}}^d\gamma _5]\mathrm{Tr}[CS_{ge^{}}^{d,T}C\gamma _5S_{fd^{}}^u]`$
$`\mathrm{Tr}[CS_{df^{}}^{u,T}CS_{ee^{}}^dCS_{fd^{}}^{u,T}C\gamma _5S_{gg^{}}^d\gamma _5]\mathrm{Tr}[CS_{dd^{}}^{u,T}CS_{eg^{}}^d\gamma _5CS_{ff^{}}^{u,T}C\gamma _5S_{ge^{}}^d])`$ (7)
where the superscripts on the quark propagator mean the quark flavor and $`a,b,c,\mathrm{},`$ are the color indices. In Fig. 1 the diagrams which contribute to the chirality odd SR up to $`d=13`$ are shown. In order to calculate the correlator to a certain order we need to consider the quark propagator to the appropriate dimension. In Fig. 2 we show the corresponding OPE diagrams for the quark propagator which lead to
$`S_{ab}^q(x)`$ $`=`$ $`i0|Tq_a(x)\overline{q}_b(0)|0`$ (8)
$`=`$ $`\delta _{ab}(\widehat{x}F_1^q+F_2^q)i\stackrel{~}{g}G_{ab}^{\mu \nu }{\displaystyle \frac{1}{x^2}}(\widehat{x}\sigma _{\mu \nu }+\sigma _{\mu \nu }\widehat{x})`$
$`m_q\stackrel{~}{g}G_{ab}^{\mu \nu }\sigma _{\mu \nu }\left(\mathrm{ln}(x^2\mathrm{\Lambda }^2/4)+2\gamma _{EM}\right),`$
$`a,b`$ are the color indices and $`\stackrel{~}{g}=g_c/32\pi ^2`$. The two functions are given by
$`F_1^q`$ $`=`$ $`{\displaystyle \frac{1}{2\pi ^2x^4}}+{\displaystyle \frac{m_q\overline{q}q}{48}}+i{\displaystyle \frac{m_qx^2}{2^73^2}}g_c\overline{q}\sigma Gq,`$
$`F_2^q`$ $`=`$ $`i{\displaystyle \frac{m_q}{4\pi ^2x^2}}+i{\displaystyle \frac{\overline{q}q}{12}}{\displaystyle \frac{x^2}{192}}g_c\overline{q}\sigma Gq+i{\displaystyle \frac{g_c^2x^4}{2^93^3}}\overline{q}qG^2`$ (9)
$`+i{\displaystyle \frac{m_qg_c^2}{2^93\pi ^2}}G^2x^2\left(\mathrm{ln}(x^2\mathrm{\Lambda }^2/4)+2\gamma _{EM}{\displaystyle \frac{2}{3}}\right),`$
where $`\gamma _{EM}`$ is the Euler–Mascheroni constant and we take $`\mathrm{\Lambda }=500`$ MeV . Note, that for massless $`u`$, $`d`$ quarks, $`F_i^u=F_i^d`$.
Our result for the SR including operators up to dimension $`d=13`$ has the form
$`{\displaystyle \frac{1}{4}}[{\displaystyle \frac{1}{15}}m_sM^{12}E_5|_{(a)}{\displaystyle \frac{2}{15}}f_saM^{10}E_4|_{(b)}+{\displaystyle \frac{1}{6}}f_sm_0^2aM^8E_3|_{(c)}`$
$`{\displaystyle \frac{1}{12}}bm_sM^8E_3|_{(d)}{\displaystyle \frac{1}{12}}bm_sM^8W_3|_{(e)}{\displaystyle \frac{4}{27}}bf_saM^6E_2|_{(f)}`$
$`+{\displaystyle \frac{1}{12}}f_sm_0^2abM^4E_1|_{(g)}{\displaystyle \frac{1}{12}}m_sba^2M^2E_0|_{(h)}+{\displaystyle \frac{8}{27}}m_sa^4|_{(i)}`$
$`{\displaystyle \frac{1}{72}}bf_sa^3|_{(j)}+{\displaystyle \frac{1}{48}}m_sm_0^2ba^2|_{(k)}]=\stackrel{~}{\lambda }^2_\mathrm{\Theta }M_\mathrm{\Theta }e^{M_\mathrm{\Theta }^2/M^2}.`$ (10)
Each term corresponds to a diagram in Fig. 1. The residue is defined by $`\stackrel{~}{\lambda }_\mathrm{\Theta }=(4\pi )^4\lambda _\mathrm{\Theta }`$. The contributions from the continuum are given by the following functions :
$`E_n(M)`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Gamma }(n+1)M^{2n+2}}}{\displaystyle _0^{s_0^2}}𝑑xe^{x/M^2}x^n,`$
$`W_n(M)`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Gamma }(n+1)M^{2n+2}}}{\displaystyle _0^{s_0^2}}dxe^{x/M^2}x^n(2\mathrm{ln}(x/\mathrm{\Lambda }^2)+\mathrm{ln}\pi `$ (11)
$`+\psi (n+1)+\psi (n+2)+2\gamma _{EM}{\displaystyle \frac{2}{3}})`$
with $`\psi (n)=1+1/2+\mathrm{}+1/(n1)\gamma _{EM}.`$ The numerical values of various quantities in the sum rule will be given in below.
We have checked that our result Eq. (10) taking into account only contributions of operators with dimensions $`d9`$ is in agreement with the previous calculations , . We should mention that due to specific spin and color structure of our current Eq. (5) the potentially important contributions from operators $`\overline{\psi }\psi ^3`$, $`\overline{\psi }\psi ^2\overline{\psi }g_c\sigma G\psi `$, $`m_sg_c\overline{q}\sigma Gq^2`$, and $`\overline{s}sg_c\overline{q}\sigma Gq^2`$ can not appear in chirality odd sum rules. The absence of the contribution of these operators can be seen by direct analysis of the terms in Eq. (7). We would like to mention that their appearance depends strongly on the structures of the interpolating current. For example, in the recent paper , where another pentaquark current has been used, it was shown that these operators give non-vanishing contribution. We do not include the contributions proportional to $`g_c^3GGG`$ and $`g_c^2GG^2`$, which in the OPE are related to terms of higher orders in the expansion in the strong coupling constant, and therefore their contributions are expected to be very small for the light quark systems (see for example the discussion in ref. on the three-gluon condensate contribution). This statement is in the agreement with a general observation that pure gluonic operators are not very important in QCD SRs for hadrons consisting of light $`u,d,`$ and $`s`$ quarks , .
## 3 The direct instanton contribution to the sum rule
In addition to contributions of power type, arising from the OPE expansion, there are exponential contributions coming from direct instantons contributions to the correlators . They can be calculated by using the following formula for the quark propagator in the instanton background in the regular gauge
$$S_{ab}^{q,inst}(x,y)=A_q(x,y)\gamma _\mu \gamma _\nu (1+\gamma _5)(U\tau _\mu ^{}\tau _\nu ^+U^{})_{ab},$$
(12)
where
$$A_q(x,y)=i\frac{\rho ^2}{16\pi ^2m_q^{}}\varphi (xz_0)\varphi (yz_0)$$
and
$$\varphi (xz_0)=\frac{1}{[(xz_0)^2+\rho ^2]^{3/2}}.$$
Note that $`\rho `$ stands for the instanton size and $`z_0`$ the center of the instanton ; $`U`$ represents the color orientation matrix of the instanton in $`SU(3)_c`$ and $`\tau _{\mu ,\nu }`$ are $`SU(2)_c`$ matrices ; $`m_q^{}=m_{cur}^q2\pi ^2\rho _c^2\overline{q}q/3`$ is the effective quark mass in the instanton vacuum and $`m_{cur}^q`$ the current quark mass. The final result should be multiplied by a factor of two to take into account the anti-instanton contribution, and has to be integrated over the instanton density.
To leading order in the instanton density, the direct instanton contributions arise from two body $`ud`$, $`u\overline{s}`$, $`d\overline{s}`$ and three body $`ud\overline{s}`$ quark zero mode propagators in the correlator Eq.(1), as shown in Fig. 3.
The final result for two body instanton contribution is
$`\mathrm{\Pi }_2(M)`$ $`=`$ $`{\displaystyle \frac{n_{eff}\rho _c^4\overline{q}q}{2^63\pi ^8m_q^{}m_s^{}}}\widehat{B}[f_6(Q)],`$ (13)
where Shuryak’s instanton liquid model for QCD vacuum with density $`n(\rho )=n_{eff}\delta (\rho \rho _c)`$ has been used and $`\widehat{B}[f_6(Q)]`$ is the Borel transform of $`f_6(Q)`$ which is defined by
$$f_6(Q)=d^4z_0d^4x\frac{e^{iqx}}{x^6[z_0^2+\rho _c^2]^3[(xz_0)^2+\rho _c^2]^3},$$
(14)
where $`\rho _c`$ is the average instanton size. There are two types of singularities in Eq.(14). One of them is related to the pole at the origin $`x^2=0`$, the other is due to the pole at finite distance from origin $`x^2\rho _c^2`$. The pole at $`x^2=0`$ produces, after Fourier transforming, power terms in $`1/Q^n`$ in addition to the exponential type direct instanton contributions $`\mathrm{exp}(Q\rho _c)`$, arising from finite distances. One should carefully subtract that contribution to avoid double counting with the standard OPE terms. We follow the procedure suggested in ref. for the analysis of direct instanton contributions to heavy quark decay. More specifically, for a general integral
$$\mathrm{\Pi }_{ins}=d^4xd^4z_0e^{iqx}\frac{S(x)}{x^{2n}((xz_0)^2+\rho _c^2)^\alpha (z_0^2+\rho _c^2)^\beta },$$
(15)
where $`S(x)`$ contains no singularities for complex $`x_\mu `$, we use Feynman’s parameterization
$$\mathrm{\Pi }_{ins}=\frac{\mathrm{\Gamma }(\alpha +\beta )}{\mathrm{\Gamma }(\alpha )\mathrm{\Gamma }(\beta )}d^4xd^4z_1e^{iqx}\frac{S(x)}{x^{2n}}_0^1𝑑t\frac{t^{\alpha 1}(1t)^{\beta 1}}{(t(1t)x^2+z_1^2+\rho _c^2)^{\alpha +\beta }},$$
(16)
where $`z_1=z_0tx`$, but consider only the contribution from the pole at
$$x^2=(z_1^2+\rho _c^2)[t(1t)]^1.$$
(17)
The Borel transform of the function $`f_6`$ is given by
$`\widehat{B}[f_6(Q)]`$ $`=`$ $`{\displaystyle \frac{\pi ^4M^{12}}{2^{13}}}{\displaystyle _0^1}dt{\displaystyle _0^1}dy{\displaystyle \frac{e^{M^2\rho _c^2/(4ty(1y))}}{y^2(1y)^2}}(X^2+5X^3+10X^4`$ (18)
$`+10X^5+5X^6+X^7),`$
where $`X=(1t)/t`$. Note that only the contribution from the pole at finite quark separation has been considered.
We have also performed the calculation of the three body contributions induced by instantons, Fig. 3b, and found that they vanish as the result of the cancellation between diagrams with different $`ud\overline{s}`$ combinations in the pentaquark current. The two body $`ud`$ instanton contribution also cancels. The only non vanishing contributions arise from the two body $`u\overline{s}`$ and $`d\overline{s}`$ terms of Eq. (13). The behaviour of the different contributions is associated to the Dirac structure of our current Eq. (5) which includes both scalar and pseudoscalar $`ud`$ diquarks with equal weights. The instanton induced contribution is very sensitive to the parity of the state and it flips sign when the parity of the state changes. Of course, the three-body instanton induced forces might give non-zero contribution for other choices of pentaquark currents, for example, for currents with derivatives.
We should mention that three body instanton terms induce forces which give non zero contributions to the mass of some specific triquark $`ud\overline{s}`$ pentaquark clusters and furthermore, they give non vanishing contribution to the $`\mathrm{\Theta }`$ mass within the bag model .
## 4 Numerical analysis
We use the following values for parameters at the normalization point 2 GeV (see also recent discussion about uncertainties in values of various condensates in )
$`\overline{u}u=(243\mathrm{MeV})^3{\displaystyle \frac{a}{(2\pi )^2}},`$
$`b=g_c^2G^2=0.88\mathrm{GeV}^4,`$
$`ig_c\overline{u}\sigma Gu=m_0^2\overline{u}u=0.8\mathrm{GeV}^2\overline{u}u,`$
$`{\displaystyle \frac{\overline{s}s}{\overline{u}u}}={\displaystyle \frac{\overline{s}\sigma Gs}{\overline{u}\sigma Gu}}=f_s=0.8,`$
$`m_s=111\mathrm{MeV},`$ (19)
with $`\rho _c=1.6`$ GeV<sup>-1</sup> for the average instanton size in the QCD vacuum. In the numerical estimate of the direct instanton contribution the relations of the instanton liquid model
$$\frac{n_{eff}}{m_q^2}=\frac{3}{4\pi ^2\rho _c^2},\frac{m_q^{}}{m_s^{}}=\frac{1}{f_s\frac{3m_s}{2\pi ^2\rho _c^2\overline{q}q}},$$
(20)
are used.
In Figs. 4 and 5 the result of the calculation of the pentaquark mass and residue within the standard OPE expansion for the different orders in operator dimensions is shown. In Figs. 6 and 7 the mass and residue of the pentaquark as a function of the value of the Borel parameter with direct instanton contributions are shown. In Fig. 8 we present the results of the calculation of the OPE and the direct instanton contributions to the left-hand side of the SR, Eq.(4). All curves are given for a value of the threshold $`s_0=2`$ GeV. We chose this value of threshold because the stability was best. From the fit of the sum rules we arrive at the following values for pentaquark mass: $`M_{\mathrm{\Theta }^+}=1.66`$ GeV for $`d=7`$, $`M_{\mathrm{\Theta }^+}=1.75`$GeV for $`d=9`$, $`M_{\mathrm{\Theta }^+}=1.73`$GeV for $`d=11`$, and $`M_{\mathrm{\Theta }^+}=1.75`$ GeV for $`d=13`$ <sup>2</sup><sup>2</sup>2Once the instanton contribution is included the stability in the Borel parameter for the SR up to $`d=5`$ operators disappears (see Fig. 6). Therefore it is not possible to extract the value of pentaquark mass from the SR with only up to $`d=5`$ operators..
One important result of our calculation is in the change of the sign of the squared of the residue when increasing the dimension of the operators which contribute to the OPE. Thus, for $`d=5`$ the sign is positive, while it becomes negative for higher dimensions. In particular, the contribution from the dimension $`d=7`$ operators is crucial for inverting the sign. Due to the negative intrinsic parity of our current Eq. (5), the negative (positive) sign of the squared of the residue implies positive (negative) parity for the state. Therefore, our final result for the residue presented in Fig. 7 shows that one can arrive to the wrong conclusion about the parity of the pentaquark state , if one takes the decision based on only the contributions of low dimension ($`d=5`$) operators. We also stress the necessity to include high dimension operators to get good convergence for the sum rule. It is evident that this effect is directly related to the high dimension of the pentaquark current.
Once we include the contributions from the high dimension operators and the instantons, our result for the $`\mathrm{\Theta }^+`$ pentaquark mass, $`M_{\mathrm{\Theta }^+}1.75`$ GeV, is higher than was given by previous SR calculations but still in rough agreement with the available experimental data, if one admits about 10% accuracy in the predictions of the SR approach due to uncertainties in the values of the various condensates, the mass of the strange quark, the contribution from higher dimension operators $`d>13`$, higher order pQCD corrections, etc. Furthermore, some additional effects such as the mixing between various pentaquark states , which are beyond the scope of the present paper, might give some additional contribution to the mass of the $`\mathrm{\Theta }^+`$.
We also note that in our calculation the pentaquark has positive parity in agreement with the soliton model prediction . Our estimate for direct instanton contribution is done within Shuryak’s instanton liquid model. We have found that the instanton contribution for the full SR is rather small, but can give a large contribution to it when one considers operators only up to dimension $`d=5`$ (see Figs. 4,6). The smallness of the instanton contribution to the full SR is mainly related to the large mass of the Borel parameter $`M1.7`$ GeV, where we obtain the plateau of stability (Fig. 6). In this region the instanton contribution is small in comparison with the contribution from the high dimension operators in the OPE (Fig. 8).
There is a significant dependence of our results on the value of threshold. This is a common feature in all the studies about the properties of the pentaquark within the QCD sum rule approach. In our case, we have chosen $`s_0=2`$ GeV to satisfy the physical requirement of having a large stability plateau.
In summary, we have shown the analysis of the QCD sum rules for the $`\mathrm{\Theta }^+`$ pentaquark current including high dimension operators in the OPE and direct instanton contributions. Our results conclude that the role of the high dimension operators is important for obtaining a positive parity for pentaquark state. Our calculation though produces a bound state whose mass is higher than the experimental observation. More sophisticated models and probably states mixing might reduce the obtained value to the observed one.
## 5 Acknowledgments
We are grateful to A.E. Dorokhov, S.V. Esaibegian, S.V. Mikhailov, A.G. Oganesian, and A.A. Pivovarov for very useful discussions. This work was supported by grants by MCyT-FIS2004-05616-C02-01 and GV-GRUPOS03/094 (VV), by Ministerio de Educacion y Ciencia of Spain (SEEU-SB2002-0009) (HJL), and by Brain Pool program of Korea Research Foundation through KOFST, grant 042T-1-1 and in part by RFBR-03-02-17291, RFBR-04-02-16445 (NIK). NIK is very grateful to the School of Physics, SNU for their warm hospitality during the final stage of this work. |
warning/0506/astro-ph0506227.html | ar5iv | text | # X-ray Spectral Evolution of Her X-1 in a Low State and the Following Short High State
## 1 Introduction
Her X-1, being an eclipsing accretion powered pulsar consisting of a $`1.24`$ s period pulsar in a $`1.7`$ day circular orbit, was discovered in 1972 with Uhuru observations (Tananbaum et al. 1972). Soon after its discovery, its optical companion was identified to be a blue variable thirteenth magnitude star HZ Her (Davidsen al. 1972). The Her X-1/Hz Her binary system was shown to display a 35 day long cycle of high and low X-ray flux states. There are two high states, main high and short high, within a single 35 day cycle lasting roughly about 10 and 5 days respectively and separated by $`10`$ day long low states (Giacconi et al. 1973; Scott & Leahy 1999). Superposed on this cycle are eclipses of the neutron star by the companion once per orbital period. While 35 day cycle is mostly coherent since the discovery of the source, there have been four occasions so far when the high states have either failed to turn-on or main high state flux has been reduced (i.e. anomalous low states, see Parmar et al. 1985; Vrtilek et al. 1994; Oosterbroek et al. 2000; Boyd et al. 2004).
The nature of the 35 day cycle is generally attributed to the periodic variable obscuration of the emission region by a precessing accretion disc viewed nearly edge-on (e.g. Petterson 1975; Scott& Leahy 1999; Scott, Leahy, Wilson 2000; Leahy 2004). It has been claimed that there is no strict periodicity of this precession (Ogelman 1987): ”35 day cycles” generally have durations of 20, 20.5 and 21 orbital cycles with equal statistical probabilities which may lead to the idea that the precession and orbital cycles are related physically (Scott & Leahy 1999). Baykal et al. (1993) revealed that the statistical interpretation of turn-on behaviour is consistent with a white-noise process in the first derivative of the 35 days phase fluctuations. Other manifestations of 35 day cycle include spectral variations in 35 days (Mihara et al. 1991; Choi et al. 1994B; Leahy 1995; Leahy 2001; Zane et al. 2004), 35 day X-ray pulse profile evolution (Scott et al. 2000), optical pulsations occurring at certain 35 day and orbital phases (Middleditch 1983), and systematic 35 day variations in the optical light curve (Deeter et al. 1976; Gerend & Boynton 1976).
Distinct features of high states of Her X-1 in X-ray band are the preeclipse X-ray absorption dips occurring at $`0.8`$ orbital phase (Giacconi et al. 1973; Crosa & Boynton 1980), and anomalous dips (Giacconi et al. 1973; Scott & Leahy 1999) occurring at $`0.4`$ orbital phase. Both types of dips are thought to be due to absorption of cold matter on the line of sight. The occurrence period of the X-ray dips is thought to be related to the beat period of 35 days and 1.7 days ($`1.65`$ days), thus they progress to earlier orbital phases within main high and short high states. More detailed features of preeclipse and anomalous dips of Her X-1 were discussed by Leahy et al. (1994), Reynolds & Parmar (1995), Leahy (1997), Shakura et al. (1998) and Stelzer et al. (1999).
Although considerable changes in the X-ray spectrum in 35 day cycle have been found to be evident, the X-ray spectrum in different 35 day phases can be fitted using a constant-emission model with varying absorption (e.g. Mihara et al. 1991; Leahy 2001). This is another supporting fact that 35 day cycle is due to the precession of an accretion disc that periodically obscures the neutron star beam. In other words, spectral variations in X-ray band are not primarily due to the variations of the primary outgoing beam of the neutron star but rather due to changes in absorption.
Being the first accretion powered pulsar found to have cyclotron absorption feature (Truemper et al. 1978), X-ray spectrum of Her X-1 shows this absorption feature around $`40`$ keV which indicates a surface magnetic field of $`(3.54.5)\times 10^{12}`$ Gauss. There is no evidence of a second harmonic around $`80`$ keV (Coburn et al. 2002).
Iron line complex at $`6.5`$ keV is evident in the X-ray spectrum of Her X-1 (Leahy 2001 and references therein). Furthermore, Endo, Nagase & Mihara (2000) were able to resolve the feature into two discrete emission lines at $`6.4`$ keV and $`6.7`$ keV. A distinct Fe XXVI line at $`7`$ keV was detected with XMM-Newton throughout 35 day cycle except the main high (Zane et al. 2004). Variations in iron line complex feature over the 35 day period and and the spin period have been evident (e.g. Choi et al. 1994), while variations of this complex feature with the orbital phase have also been studied (e.g. Zane et al. 2004).
In this paper, results of spectral analysis of December 2001 low state and short high state RXTE data of Her X-1 are presented. This set of observations enables, for the first time, frequent continuous monitoring (111 pointings in $`8.5`$ days) of the source with RXTE including $`1.7`$ days ($`1`$ orbital period) long low state part and the following $`6.8`$ days ($`4`$ orbital periods) long short high state part. Section 2 is an overview about the instruments and the observations. In Section 3, we present the analysis of the observations. In Section 4, results are discussed.
## 2 Instruments and Observations
The RXTE observations used in this paper cover the period from 18 to 26 December 2001 (MJD 52261-52269) with a total exposure of 281.2 ksec (Figure 1). These observations have common proposal ID 60017. The results presented here are based on data collected with the Proportional Counter Array (PCA, Jahoda et al., 1996) and the High Energy X-ray Timing Experiment (HEXTE, Rothschild et al. 1998). The PCA instrument consists of an array of 5 proportional counters (PCU) operating in the 2-60 keV energy range, with a total effective area of approximately 6250 cm<sup>2</sup> and a field of view of $`1^{}`$ FWHM. Although the number of active PCU’s varied between 2 and 5 during the observations, observations after 2000 May 13 belong to the observational epoch for which background level for one of the PCUs (PCU0) increased due to the fact that this PCU started to operate without a propane layer. The latest combined background models (CM) are used together with FTOOLS 5.3 to estimate the appropriate PCA background. The HEXTE instrument consists of two independent clusters of detectors, each cluster containing four NaI(TI)/CsI(Na) PHOSWICH scintillation counters (one of the detectors in cluster 2 is not used for spectral information) sharing a common $`1^{}`$ FWHM. The field of view of each cluster is switched on and off source to provide background measurements. The net open area of the seven detectors used for spectroscopy is 1400 cm<sup>2</sup>. Each detector covers the energy range 15-250 keV.
Spectral analysis and error estimation of the spectral parameters are performed using XSPEC. The overall 3-60 keV PCA-HEXTE short high spectrum is extracted. In addition to the overall spectrum, individual PCA spectra are extracted in 3-20 keV.
## 3 Spectral Analysis
2001 December RXTE observation of Her X-1 consists of one low state and four short high state orbits with a total exposure of 281.2 ksec. We extract 111 individual PCA spectra in 3-20 keV energy range to investigate the spectral evolution throughout the whole dataset. For the individual PCA spectra, energies lower than 3 keV are ignored due to uncertainties in background modeling while energies higher than 20 keV are ignored as a result of poor counting statistics. Then, overall 3-60 keV PCA-HEXTE spectrum from the short high state orbits is extracted (Figure 2). We exclude the 3-60 keV spectrum of the low state orbit since the count rate coming from the channels corresponding to $`20`$ keV is not statistically significant for both PCA and HEXTE. Additional 2% systematic error (see Wilms et al. 1999; Coburn et al. 2000) is added to the statistical errors for both the individual spectra and 3-60 keV PCA-HEXTE spectrum.
It is found that the model consisting of an absorbed power law (Morrison & McCammon, 1983) with an high energy cut-off (White, Swank, Holt 1983) and an iron line feature modeled as a Gaussian is not able to fit the data well. This is expected since two component spectral model with a basic emission spectrum and an absorbed component models or partial covering soft absorption models have been found to fit the main high, short high and low state spectrum of Her X-1 (Mihara et al. 1991; Oosterbroek et al. 1997; Coburn et al. 2000; Oosterbroek et al. 2000; Choi et al. 1994A,B; Leahy 2001). We have found that the model containing partial cold absorber (i.e the model contains both overall absorption represented with the wabs model and partial absorption represented with the pcfabs in XSPEC) fit the 3-20 keV PCA spectra well for both low state and short high state parts of the observation.
In Table 1, we present sample results of our spectral fits to low state, peak of short high state, anomalous dip and preeclipse dip data. In Figure 3, we present evolution of spectral parameters during the low state and short high state orbits.
To fit the overall 3-60 keV PCA-HEXTE short high state spectrum, we used the same spectral model with an additional cyclotron absorption feature at $`42`$ keV (Table 2; Figure 2). Her X-1 is the first X-ray pulsar to have shown to have this spectral feature (Truemper et al. 1978). Gruber et al. (2001) showed that the cyclotron resonance energy of Her X-1 changed from $`34`$ to $`42`$ keV sometime between 1991 and 1993.
## 4 Discussion and Conclusion
We analyzed spectra obtained from continuous RXTE monitoring of $`8.5`$ days of Her X-1. This rather frequent continuous monitoring, for the first time, enabled the study of the spectral evolution of low state, short high peaks and short high dips within a relatively compact observation period.
For the short high state part, stability of the power law index, cut-off energy and e-folding energy throughout the whole monitoring may be an indication of the fact that accretion geometry and mass accretion rate does not change considerably. This idea is in accord with the fact that the X-ray flux changes in Her X-1 should primarily be due to the opaque obscuration of the pulsar beam, while there should not be no significant variation in the flux emerging initially from the magnetic poles of the pulsar.
We have found some evidences of differences between anomalous dips and preeclipse dips in short high state orbits. First marginal evidence is the higher ratio of unabsorbed flux to absorbed flux for more than $`50`$% percent of preeclipse dips compared to anomalous dips (left panel of Figure 4). For the rest of the preeclipse dip regions, this ratio is similar to that of anomalous dips. We consider this ratio as a measure of sum of absorptions coming from the total absorption and partial covering absorption (i.e. greater value of this ratio means more absorption). Then, this marginal evidence shows that absorption is more dominant in preeclipse dips compared to anomalous dips. This result is a confirmation of the previous EXOSAT (Reynolds and Parmar 1995) and GINGA (Leahy 1997) results. From EXOSAT observations, it was found that Hydrogen column density was about a few times greater in the preeclipse dips compared to that in the anomalous dips, and that partial covering fractions were higher in preeclipse dips compared to those in the anomalous dips. Similarly, GINGA observations of the absorption dips revealed that $`n_H`$ was almost an order higher in the preeclipse dips compared to that in anomalous dips. This difference of absorption in two types of absorption dips was explained as a sign of the fact that the intersection of the outer rim of the accretion disc with the line of sight is responsible for the both types of absorption dips, and different absorptions for both types of dips is only due to the difference in orbital phase for these dips (Crosa and Boynton, 1980; Leahy 1997). Another evidence supporting this interpretation was the lack of 35 day dependence on the variation of absorption in absorption dips, which is in accord with the fact that the outer edge of the accretion disc, being not responsible of 35 day cycle, should not be affected by 35 day cycle of precession (Leahy 1997). From the left panel of Figure 4, we have also found that unabsorbed flux is higher in anomalous dips compared to preeclipse dips and that the ratio of unabsorbed flux to absorbed flux decreases with increasing unabsorbed flux independently for both preeclipse dips and anomalous dips. Higher unabsorbed flux in anomalous dips is consistent with the general orbital phase dependence of X-ray flux.
On the other hand, decreasing ratio of unabsorbed flux to absorbed flux (i.e decreasing absorption) with increasing unabsorbed flux can be understood by the fact that the X-ray flux variation of Her X-1 is the result of opaque obscuration rather than soft absorption (e.g. Still et al. 2001). Variations of unabsorbed flux should be primarily due to the change in opaque obscuration. Opaque obscuration is result of either obscuration of the accretion disc or the eclipsing stellar companion. This relation, therefore, shows that the regions causing opaque obscuration and soft absorption are not geometrically far away from each other.
Low states in 35 day cycle of Her X-1 is generally believed to be composed of the phases in which X-rays from the neutron star is obscured (Petterson 1975; Scott& Leahy 1999; Scott, Leahy, Wilson 2000; Leahy 2004). The resultant X-ray flux for both normal low states and anomalous low states is a combination of scattered and reflected components and has been found to fit well except for normalizations and excess absorption to high state spectrum models (e.g. Coburn et al. 2000). Coburn et al. (2000) also found that absorption can be modeled with a cold partial absorber with $`n_H`$ $`5\times 10^{23}`$ cm<sup>-2</sup> and partial covering fraction of $`0.7`$ for both normal and anomalous low states. Results of our low state fits, giving values $`(510)\times 10^{23}`$ cm<sup>-2</sup> and $`0.30.8`$ for $`n_H`$ and partial covering fraction respectively, are in agreement with those results.
From the right panel of Figure 4, iron line flux is correlated, in general, with the X-ray flux for both short high state and low state observations. However, there is a marginal evidence of higher iron line flux for the anomalous dips. Increase in the strength of the iron line feature in the anomalous dips is a natural consequence of the fact that we have some of the X-ray flux coming from the neutron star passing through the outer rim of the accretion disc which is possibly enhancing this iron line feature, so the resultant iron line flux increases in the anomalous dips. Although the origin of the anomalous dips and preeclipse dips are similar, the iron line flux is not significantly higher for the preeclipse dips. This may be related to orbital phase difference between anomalous dips and preeclipse dips. Future observations of preeclipse dips of Her X-1 will be useful to study the variation of iron line feature.
Peak energy of the iron line feature is, in general, at $`6.6`$ keV for the short high state. From Figure 3, the peak energy tends to increase at the initial short high turn-on (around $`1.9`$ days), anomalous dip of the first short high orbit and preeclipse dips. The iron line feature is also found to peak at higher energies ($`6.8`$ keV) in the low state. Iron line energy peaking at higher energies is a sign of the fact that this feature does not consist of only the K$`\alpha `$ Fe emission line. In fact, it has recently been noted that K$`\alpha `$ Fe line emission tends to be less powerful in short high and low states compared to the main high state (Zane et al. 2004). Zane et al. (2004), using XMM-Newton observations, also has been able to resolve a secondary iron line feature (Fe XXVI) peaking at $`6.97.0`$ keV which is present only in low state and short high state. Spectral resolution of RXTE-PCA is not good enough to resolve these lines but the higher peak energies of the resultant iron line feature may be a clue of the presence of iron line components other than K$`\alpha `$ emission line.
Acknowledgments
We acknowledge D.M. Scott for his guidance and we thank him for the useful discussions. S. C.İ. acknowledges the Integrated Doctorate Programme scholarship from the Scientific and Technical Research Council of Turkey (TÜBİTAK). |
warning/0506/hep-ex0506023.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The study of the radiative kaon decays can give valuable information on the kaon structure and allows for good test of theories describing hadron interactions and decays, like Chiral Perturbation Theory (ChPT). Until now the studies of the radiative K<sub>l3</sub>-decays are restricted by the decay modes with electrons in the final state or by the studies of K<sub>L</sub> decays . Only one paper , dated by 1973, published the upper limit on the branching of K$`{}_{\mu 3\gamma }{}^{}{}_{}{}^{+}`$-decay.
The interest to the study of K<sub>l3γ</sub> decays is further enhanced by the theoretical proposals to search for effects of new physics using T-odd kinematical variable $`\xi =\stackrel{}{p}_\gamma \left(\stackrel{}{p}_l\times \stackrel{}{p}_\pi \right)/m_K^3`$ . In the standard model the expected asymmetry for K$`{}_{\mu 3\gamma }{}^{}{}_{}{}^{}`$ decay
A$`\left(\xi \right)=\frac{N\left(\xi >0\right)N\left(\xi <0\right)}{N\left(\xi >0\right)+N\left(\xi <0\right)}`$
is at the level $`1.14\times 10^4`$ , whereas in the extensions of the standard model it can achieve $`2.6\times 10^4`$ .
In this paper we present first observation of the radiative K$`{}_{\mu 3}{}^{}{}_{}{}^{}`$ decay. The experimental setup and event selection are described in section 2. The results of the analysis are presented in section 3, where we, first, show the presence of the signal for the photon energy in the kaon rest frame E$`{}_{\gamma }{}^{}{}_{}{}^{}`$ below 60 MeV, measure the branching ratio for the region 5 $`<`$ E$`{}_{\gamma }{}^{}{}_{}{}^{}`$ $`<`$ 30 MeV, measure the asymmetries in this region and finally measure the branching ratio for the region 30 $`<`$ E$`{}_{\gamma }{}^{}{}_{}{}^{}`$ $`<`$ 60 MeV. Our conclusions are given in the last section.
## 2 Experimental setup and event selection
The experiment is performed at the IHEP 70 GeV proton synchrotron U-70. The ISTRA+ spectrometer has been described in detail in recent papers on $`K_{e3}`$ , $`K_{\mu 3}`$ and $`\pi ^{}\pi ^{}\pi ^{}`$ decays . Here we recall briefly the characteristics relevant to our analysis. The ISTRA+ setup is located in a negative unseparated secondary beam line 4A of U-70. The beam momentum is $`25`$ GeV/c with $`\mathrm{\Delta }p/p1.5\%`$. The admixture of $`K^{}`$ in the beam is $`3\%`$, the beam intensity is $`310^6`$ per 1.9 sec U-70 spill. A schematic view of the ISTRA+ setup is shown in Fig. 1. The beam particles are deflected by the magnet M<sub>1</sub> and are measured by four proportional chambers BPC<sub>1</sub>—BPC<sub>4</sub> with 1 mm wire spacing, the kaon identification is done by three threshold Cerenkov counters Č<sub>0</sub>—Č<sub>2</sub>. The 9 meter long vacuum decay volume is surrounded by eight lead glass rings used to veto low energy photons. The 72-cell lead-glass calorimeter SP<sub>2</sub> plays the same role. The decay products are deflected in the magnet M<sub>2</sub> with 1 Tm field integral and are measured with 2 mm step proportional chambers PC<sub>1</sub>—PC<sub>3</sub>, with 1 cm cell drift chambers DC<sub>1</sub>—DC<sub>3</sub> and, finally, with 2 cm diameter drift tubes DT<sub>1</sub>—DT<sub>4</sub>. The wide aperture threshold Cerenkov counters Č<sub>3</sub> ,Č<sub>4</sub> , filled with He, serve to trigger electrons and are not used in the present measurement. SP<sub>1</sub> is a 576-cell lead-glass calorimeter, followed by HC, a scintillator-iron sampling hadron calorimeter. MH is a 11x11 cell scintillating hodoscope, used to improve the time resolution of the tracking system, MuH is a 7x7 cell muon hodoscope.
The trigger is provided by scintillation counters S<sub>1</sub>—S<sub>5</sub>, beam Cerenkov counters and by the analog sum of amplitudes from last dynodes of the SP<sub>1</sub> : T=S$`{}_{1}{}^{}\text{S}_2\text{S}_3\overline{\text{S}}_4\text{Č}_1\overline{\text{Č}}_2\overline{\text{Č}}_3\overline{\text{S}}_5\mathrm{\Sigma }\left(\text{SP}_1\right)`$, here S<sub>4</sub> is a scintillation counter with a hole to suppress the beam halo, S<sub>5</sub> is a counter downstream of the setup at the beam focus, $`\mathrm{\Sigma }`$(SP$`{}_{1}{}^{})`$ requires that the analog sum to be larger than the MIP signal.
During the physics run in November-December 2001 350 million trigger events were collected with high beam intensity. This information is complemented by 150 M Monte Carlo (MC) events generated using Geant3 for the dominant $`K^{}`$ decay modes, 100 M of them are the mixture of the dominant decay modes with the branchings exceeding 1 %, 30 M are the decays K$`{}_{}{}^{}\mu ^{}\pi ^{}\nu (`$K<sub>μ3</sub>) and 20 M are the decays K$`{}_{}{}^{}\pi ^{}\pi ^{}\pi ^{}(`$K<sub>π3</sub>). Signal efficiency has been estimated using 5.7 M MC events of the radiative $`K_{\mu 3}`$ decay weighted with the matrix element, calculated in the leading approximation (up to terms of $`O\left(p^4\right)`$) of chiral perturbation theory
Some information on the data processing and reconstruction procedures is given in , here we briefly mention the details relevant for present analysis.
The muon identification (see ) is based on the information from the SP<sub>1</sub> and the HC. The energy deposition in the SP<sub>1</sub> is required to be compatible with the MIP signal in order to suppress charged pions and electrons. The sum of the signals in the HC cells associated with charged track is required to be compatible with the MIP signal. The muon selection is further enhanced by the requirement that the ratio $`r_3`$ of the HC energy in last three layers to the total HC energy exceeds 5 %. The used cut values are the same as in .
Events with one reconstructed charged track and three reconstructed showers in the calorimeter SP<sub>1</sub> are selected. We require the effective mass m$`\left(\gamma \gamma \right)`$ to be within $`\pm 40`$ MeV/c<sup>2</sup> from m$`_\pi ^{}`$. In the following analysis the central $`\pm 20`$ MeV/c<sup>2</sup> band is used for signal search and the side bands $`95115`$ MeV/c<sup>2</sup> and $`155175`$ MeV/c<sup>2</sup> are used for background studies. We require also the reconstructed $`z`$-coordinate of the vertex to be below 1650 cm. 183672 events have been selected and written to miniDST’s using the above cuts with relaxed cut on $`r_3`$ to be above 1 %.
## 3 Evidence for signal and measurements of the branching ratios
A set of cuts is developed to suppress backgrounds to the K<sub>μ3γ</sub> decay and/or to do data cleaning:
0) We select events with good charged track having two reconstucted ($`xz`$ and $`yz`$) projections and the number of hits in the MH below 5. We require also that the missing mass squared to the ($`\mu ^{}\pi ^{}\gamma `$) system abs(m$`{}_{}{}^{2}\left(\mu ^{}\pi ^{}\gamma \right))<0.05`$ (GeV/c<sup>2</sup>)<sup>2</sup>. 50804 events have survived these cuts.
1) Events with the reconstructed vertex inside the interval $`400<z<1600`$ cm are selected.
2) The measured missing energy $`E_{mis}=E_{beam}E_\mu E_\pi ^{}E_\gamma `$ is required to be above zero.
3) We require the effective mass M$`\left(\gamma \gamma \right)`$ to be within $`\pm 20`$ MeV/c<sup>2</sup> from m$`_\pi ^{}`$.
4) We require also that the missing mass squared to the $`\pi ^{}\pi ^{}`$ system is below 0.025 (GeV/c<sup>2</sup>)<sup>2</sup> ( m$`{}_{}{}^{2}\left(\pi ^{}\pi ^{}\right)<0.025`$).
5) The events with missing momentum pointing to the SP<sub>1</sub> working aperture are selected in order to suppress possible $`\pi ^{}\pi ^{}\gamma `$ background ( $`6<r<60`$ cm, here $`r`$ is the distance between the impact point of the missing momentum and the SP<sub>1</sub> center in the SP<sub>1</sub> transverse plane).
6) We require the photon energy E$`{}_{\gamma }{}^{}{}_{}{}^{}`$ in the kaon rest frame to be below 60 MeV.
The remaining $`K_{\pi 2}`$ decays are suppressed by requirements:
7) $`\mathrm{cos}\left(\theta \right)>0.96`$ , where $`\theta `$ is the angle between $`\pi ^{}`$ and $`\pi ^{}`$ in the kaon rest frame;
8) $`\phi <3.0`$, where $`\phi `$ is the angle between $`\pi ^{}`$ and $`\pi ^{}`$ in the laboratory frame in the plane perpendicular to the beam momentum.
9) We require also the absence of the signal above the threshold in the calorimeter SP<sub>2</sub>.
We look for a signal in the distributions over the effective mass M$`\left(\mu ^{}\pi ^{}\gamma \nu \right)`$, where $`\nu `$ four-momentum is calculated using the measured missing momentum and assuming m$`{}_{\nu }{}^{}=0`$, and in the distributions of the missing mass squared to the $`\left(\mu ^{}\pi ^{}\gamma \right)`$-system, m$`{}_{}{}^{2}\left(\mu ^{}\pi ^{}\gamma \right)`$. Effective mass spectra for cut levels 1, 4, 6 and 9 are shown in Fig. 2. These spectra show the evidence for peak at m<sub>K</sub> after the cut on the photon energy in the kaon rest frame.
### 3.1 The region below 30 MeV.
We have found that the signal is clearly seen for E$`{}_{\gamma }{}^{}<30`$ MeV and the background in this region is dominated by K<sub>μ3</sub> decays (with an accidental extra photon) and K<sub>π3</sub> decays. The main MC sample of 100 M events (with the natural mixture of the dominant decay modes) has been found to be insufficient for estimates of the background shapes, therefore specialized MC samples of 20 M K<sub>π3</sub> and 30 M K<sub>μ3</sub> events have been used. The background has been divided into three contributions:
1) Non-$`\pi ^{}`$ contribution has been estimated using tails of the M$`\left(\gamma \gamma \right)`$ distribution for real data, see Fig. 3 .
2) K<sub>π3</sub> contribution has been approximated by the form given by specialized MC sample, its normalization has been fixed using the observed K<sub>π3</sub> signal in the m$`{}_{}{}^{2}\left(\pi ^{}\pi ^{}\right)`$ distribution for selected events.
3) K<sub>μ3</sub> contribution has been approximated by the form from specialized MC sample, its normalization has been kept free.
The shapes for all three background contributions have been found using the histogram smoothing by the HQUAD routine from the HBOOK package . The signal has been parametrized by the sum of two Gaussians with widths and relative fractions fixed at the values given by the signal MC sample.
Results of the fits are illustrated respectively in Fig. 4 and Fig. 5 for the distributions over M$`\left(\mu \pi \gamma \nu \right)`$ and m$`{}_{}{}^{2}\left(\mu \pi \gamma \right)`$. First parameter here (and in the following) is the number of observed events, second parameter is the position of the peak , and three last parameters are the respective normalization factors of the K<sub>π3</sub>, K<sub>μ3</sub> and non-$`\pi ^{}`$ contributions.
The number of observed events is equal to $`383.7\pm 40.9`$ in Fig. 4 and to $`412.9\pm 36.2`$ in Fig. 5. The difference between these two values (29.2) is our estimate of the systematics caused by the imprecise knowledge of the backgrounds. The results of fits with the polynomial parametrization of the background lie also within this uncertainty.
The K<sub>μ3</sub> decay has been used for the normalization. The number of K<sub>μ3</sub> events in the region 400$`<z<1600`$ cm, corrected for the geometrical acceptance and the measurement efficiency, has been found using the cuts described in . It is equal to $`N(`$K$`{}_{\mu 3}{}^{})=5536000`$. Independent normalization of using K<sub>π2</sub> decays gives the branching ratio lower by 6.2%. The origin of this difference is explained mainly by the trigger bias. This difference has been taken into account in our final estimates of the systematic uncertainties.
The signal efficiency has been found from the signal MC weighted with the matrix element, calculated within O$`\left(p^4\right)`$ ChPT approximation. It is equal to 2.6 %. The signal efficiency has been calculated using the cut E$`{}_{\gamma }{}^{}>5`$ MeV, the cut value is our detection threshold explained by the beam momentum and the threshold in the energy of the SP<sub>1</sub> showers equal to $`0.5`$ GeV.
The measured branching ratio is equal to BR = ($`8.82\pm 0.94(`$stat$`)\pm 0.86(`$syst$`\left)\right)\times 10^5`$. This should be compared with theoretical prediction of 6.86$`\times 10^5`$. The ratio R=BR(K$`{}_{\mu 3\gamma }{}^{})`$/BR(K$`{}_{\mu 3}{}^{})`$ is equal to R = ($`2.70\pm 0.29(`$stat$`)\pm 0.26(`$syst$`\left)\right)\times 10^3`$. This should be compared with theoretical prediction of 2.1$`\times 10^3`$. In the transformations from the ratio R to the branching ratio we use the branching ratio BR(K$`{}_{\mu 3}{}^{})=3.27\%`$ .
### 3.2 Asymmetries for the region $`5<`$E$`{}_{\gamma }{}^{}<30`$ MeV.
For this region we have measured the asymmetry of photon emission towards muon direction in the kaon rest frame $`A\left(\mathrm{cos}\theta _{\mu \gamma }^{}\right)`$
A$`\left(\mathrm{cos}\theta _{\mu \gamma }^{}\right)=\frac{N\left(\mathrm{cos}\theta _{\mu \gamma }^{}>0\right)N\left(\mathrm{cos}\theta _{\mu \gamma }^{}<0\right)}{N\left(\mathrm{cos}\theta _{\mu \gamma }^{}>0\right)+N\left(\mathrm{cos}\theta _{\mu \gamma }^{}<0\right)}`$
and the asymmetry in $`\xi `$, A$`\left(\xi \right)`$. The effective mass spectra are shown in Fig. 6 for positive and negative $`\mathrm{cos}\theta _{\mu \gamma }^{}`$ separately and in Fig. 7 for positive and negative $`\xi `$.
The asymmetry in $`\mathrm{cos}\theta _{\mu \gamma }^{}`$ is equal to $`0.09\pm 0.14`$, this value is below the theoretical expectation, equal to 0.35, by 2 standard deviations.
The asymmetry in $`\xi `$ is equal to $`0.03\pm 0.13`$. Of course our statistics is insufficient to test the theoretical predictions .
### 3.3 The region $`30<`$E$`{}_{\gamma }{}^{}<60`$ MeV.
For this region we see strong K<sub>π3</sub> background and residual K<sub>π2</sub> background. These backgrounds have been suppressed by additional cut 0.1$`<p^{}\left(\pi ^{}\right)<0.185`$ MeV/c. The parametrization is illustrated in Fig. 8 for the effective mass M$`\left(\mu ^{}\pi ^{}\gamma \nu \right)`$ spectrum. The signal efficiency was found to be 6.30 %. The number of observed events is equal to $`152.7\pm 23.06(`$stat$`)\pm 32.2(`$syst$`)`$. The systematics in the number of observed events has been calculated in the same way as in the Section 3.1. The branching ratio BR = ($`1.46\pm 0.22(`$stat$`)\pm 0.32(`$syst$`\left)\right)\times 10^5`$, is compatible with the theoretical expectation $`1.53\times 10^5`$. The ratio R = ($`4.48\pm 0.68(`$stat$`)\pm 0.99(`$syst$`\left)\right)\times 10^4`$, is compatible with the expectation $`4.67\times 10^4`$.
## 4 Conclusions
Our conclusions are the following.
$``$ First observation of the radiative kaon decay K$`{}_{\mu 3\gamma }{}^{}{}_{}{}^{}`$ is presented.
$``$ The measured ratio R = BR(K$`{}_{\mu 3\gamma }{}^{})/`$BR(K$`{}_{\mu 3}{}^{})`$ for the region $`5<`$E$`{}_{\gamma }{}^{}<30`$ MeV is equal to 0.270$`\pm 0.029(`$stat$`)\pm 0.026(`$syst$`)\%`$. This is consistent with theoretical prediction equal to 0.21 %.
$``$ The measured ratio R for the region $`30<`$E$`{}_{\gamma }{}^{}<60`$ MeV is equal to ($`4.48\pm 0.68(`$stat$`)\pm 0.99(`$syst$`\left)\right)\times 10^4`$, this value is compatible with theoretical prediction equal to $`4.67\times 10^4`$.
$``$ The measured asymmetry in the T-odd variable $`\xi `$ for the region $`5<`$E$`{}_{\gamma }{}^{}<30`$ MeV is equal to $`0.03\pm 0.13`$.
$``$ The measured asymmetry in the $`\mathrm{cos}\theta _{\mu \gamma }^{}`$ is equal to $`0.093\pm 0.141`$, this value is two standard deviations away from the theoretical prediction equal to 0.354.
The work is supported in part by the RFBR grant N03-02-16330(IHEP group) and RFBR grant N03-0216135(INR group). We are indebted to V. Braguta for giving us routine for O$`\left(p^4\right)`$ ChPT matrix element calculations. |
warning/0506/cond-mat0506269.html | ar5iv | text | # Local Features of the Fermi Surface Curvature and the Anomalous Skin Effect in Metals
## I I. Introduction
It is well known that electromagnetic waves incident at the surface of a metal cannot penetrate deeply inside. Actually, the field inside the metal vanishes at the distances of the order of $`\delta `$ from the surface. This effect is called the skin effect, and the characteristic depth $`\delta `$ is called the skin depth. The suppression of the electromagnetic field inside the metal originates from the response of conduction electrons, and it occurs when the frequency $`\omega `$ of the incident wave is smaller than the electrons plasma frequency $`\omega _p.`$ The latter is the characteristic frequency for the response of the conduction electrons system to an external disturbance. When $`\omega >\omega _p`$ the electrons are too slow to respond, and the electromagnetic field penetrates into the metal without decay. Due to the skin effect the incident electromagnetic field could affect condition electrons only when they move inside the layer of the thickness $`\delta `$ near the metal surface. The skin depth depends on the electric conductivity of the metal $`\sigma `$ and on the frequency $`\omega `$ of the incident wave as well. Increase in $`\sigma `$ and/or $`\omega `$ leads to the decrease in the skin depth. At high frequencies $`\tau ^1\omega \omega _p(\tau `$ is the scattering time for conduction electrons) and low temperatures, $`\delta `$ may become smaller than the electrons mean free path $`l.`$ When the condition $`\delta <l`$ is satisfied the effect is referred to as the anomalous skin effect. At the anomalous skin effect the response of a metal to an incident electromagnetic wave is determined with the electrons moving in the skin layer nearly in parallel with the surface of the metal sample. These “efficient” electrons are associated with a few small “effective segments” on the Fermi surface (FS). The remaining electrons stay in the skin layer only for a very short while which prevents them from responding to the electromagnetic field.
A theory of the anomalous skin effect in metals was first proposed more than five decades ago by A. B. Pippard 1 and G. E. Reuter and E. H. Zondheimer 2 , and R. B. Dingle 3 using an isotropic model for a metal. The main results of these studies were presented in some books where high frequency phenomena in metals were discussed 4 . Then the theory was further developed to make it applicable to realistic metals with anisotropic Fermi surfaces 5 ; 6 ; 7 ; 8 ; 9 . It became clear that the response of conduction electrons to an external electromagnetic field under the anomalous skin effect depends on the Fermi surface (FS) geometry, especially its Gaussian curvature $`K(𝐩)=1/R_1(𝐩)R_2(𝐩),`$ where $`R_{1,2}(𝐩)`$ are the principal radii of curvature. For the most of real metals FSs are complex in shape, and their curvature turns zero at some points. These points could be partitioned in two classes. First, these exist zero curvature points where only one of the principal radii has a singularity, whereas another one remains finite. Usually, such points are combined in lines of zero curvature. The latter are either inflection lines or they label positions of nearly cylindrical strips on the FSs. Also, some points could be found where both principal radii tend to infinity. These points are set out separately, and the FSs are flattened in their vicinities.
When a FS includes points of zero curvature it leads to an enhancement of the contribution from the neighborhoods of these points to the electron density of states (DOS) on the FS. Normally, this enhanced contribution is small compared to the main term of the DOS which originates from the major part of the FS. Therefore it cannot produce noticeable changes in the response of the metal when all segments of the FS contribute essentially equally. However, when the curvature turns zero at some points on an “effective” part of the FS, it can give a sensible enhancement in the number of efficient electrons and, in consequence, a pronounced change in the response of the metal to the disturbance.
It has been shown that when the FS includes nearly cylindrical and/or flattened segments, noticeable changes may be observed in the frequency and temperature dependencies of sound dispersion and absorption 10 ; 11 ; 12 ; 13 . Also, the shape and amplitude of quantum oscillations in various characteristics of a metal could be affected by of the FS local geometry in the vicinities of the extremal cross sections. Qualitative anomalies in the de Haas-van Alphen oscillations associated with cylindrical pieces of the FSs were considered in Refs. 14 ; 15 . Similar anomalies in quantum oscillations in the static elastic constants and the velocity of sound were analyzed in 16 ; 17 .
Here, we concentrate on the analysis of possible manifestations of the FS local geometry in the surface impedance of a metal at the anomalous skin effect. In this case the main contribution to the surface impedance of a metal originates from electrons moving nearly in parallel with the surface of the metal. These electrons are efficient quasiparticles, and they belong to the “effective” part of the FS. The effects of the FS geometry on the metal response at the anomalous skin effect were analyzed before 9 ; 18 adopting some simplified models for the FS. The purpose of the present work is to carry out a general analysis whose results are independent on particularities in energy-momentum relations and could be applied to a broad class of metals.
## II II. results and discussion
We consider a metal filling the half-space $`z<0`$. A plane electromagnetic wave is incident on the metal surface making a right angle with the latter. To analyze the response of the metal to the wave we calculate the surface impedance:
$$Z_{\alpha \beta }=E_\alpha (0)/_0^{\mathrm{}}J_\beta (z)𝑑z$$
(1)
Here, $`\alpha ,\beta =x,y;E_\alpha (z)`$ and $`J_\beta (z)`$ are the components of the electric field $`𝐄`$ and electric current density $`𝐉,`$ respectively. Considering the anomalous skin effect we can limit our analysis to the case of specular reflection of electrons from the surface. Then the surface impedance tensor has the form:
$$Z_{\alpha \beta }=\frac{8i\omega }{c^2}_0^{\mathrm{}}\left(\frac{4\pi i\omega }{c^2}\sigma q^2I\right)_{\alpha \beta }^1𝑑q.$$
(2)
Here, $`\omega `$ and $`𝐪`$ are the frequency and the wave vector of the incident wave, respectively $`(𝐪=(0,0,q));\sigma `$ is the electron conductivity tensor, and $`I_{\alpha \beta }=\delta _{\alpha \beta }.`$
To proceed we assume that the FS has a mirror symmetry in a momentum space relative to a plane $`p_z=0.`$ To simplify calculations of the electron conductivity we divide each sheet of the FS in segments in such a way that the momentum $`𝐩`$ is a one-to-one function of the electron velocity $`𝐯`$ over a segment. The segments may coincide with the FS sheets. Also, it could happen that some sheets include a few segments. This depends on the FS shape. In calculation of the conductivity we carry out integration over each segment using spherical coordinates in the velocity space, namely, the velocity magnitude at the $`j`$-th segment $`v_j,`$ and the spherical angles $`\theta ,\phi .`$ So, the element of the surface area is given by the expression: $`dA_j=\mathrm{sin}\theta d\theta d\phi /|K_j(\theta ,\phi )|`$ where $`K_j(\theta ,\phi )`$ is the Gaussian curvature of the $`j`$-th FS segment. Summing up contributions from all these segments we obtain:
$`\sigma _{\alpha \beta }(\omega ,q)={\displaystyle \frac{ie^2}{4\pi ^3\mathrm{}^3q}}{\displaystyle \underset{j}{}}{\displaystyle 𝑑\phi }`$
$`\times {\displaystyle }{\displaystyle \frac{n_\alpha n_\beta \mathrm{sin}\theta d\theta }{\left|K_j(\theta ,\phi )\right|\left[(\omega +i/\tau )/qv_j\mathrm{cos}\theta \right]}}.`$ (3)
Here, $`n_{\alpha ,\beta }=v_{j\alpha ,\beta }/v_j,`$ and $`\tau `$ is the electron scattering time. The limits in the integrals over $`\theta ,\phi `$ are determined with the shape of the segments. We remark, however, that the effective strips on the FS are determined by the condition $`v_z0`$ for efficient electrons move in parallel with the metal surface at $`z=0.`$ Therefore, the upper limit in the integral over $`\theta `$ in the terms corresponding to the segments including the effective strips must equal $`\pi /2.`$ In the following calculations we omit the term $`i/\tau `$ in the Eq. (3) assuming $`\omega \tau 1`$ which is typical for the anomalous skin effect in good metals. Using Eq. (3) we can easily write out the expressions for the conductivity components. We have:
$`\sigma _{xx}(\omega ,q)={\displaystyle \frac{ie^2}{4\pi ^3\mathrm{}^3q}}{\displaystyle \underset{j}{}}{\displaystyle 𝑑\phi }`$
$`\times {\displaystyle }{\displaystyle \frac{d\theta \mathrm{cos}^2\phi \mathrm{sin}^3\theta }{\left|K_j(\theta ,\phi )\right|\left[(\omega +i/\tau )/qv_j\mathrm{cos}\theta \right]}}.`$ (4)
Another conductivity component $`\sigma _{yy}`$ is described with the similar expression where $`\mathrm{cos}^2\phi `$ in the integrand numerator is replaced by $`\mathrm{sin}^2\phi .`$ In further calculations we assume for simplicity that the chosen $`z`$ axis coincides with a high symmetry axis for the FS, so that both conductivity and surface impedance tensors are diagonalized.
The main contribution to the surface impedance under the anomalous skin effect comes from the region of large $`q`$ where $`\omega /qv1.`$ To calculate the corresponding asymptotic expressions for the conductivity components we expand the integrand in the Eq. (4) in powers of $`\omega /qv`$. Then we can write the well known result for the principal term in the expansion of the conductivity component $`\sigma _{xx}(\omega ,q)`$:
$$\sigma _0(q)=\frac{e^2}{4\pi ^3\mathrm{}^3q}\underset{l}{}𝑑\phi \frac{\mathrm{cos}^2\phi }{\left|K_l(\pi /2,\phi )\right|}\frac{e^2}{4\pi \mathrm{}^3q}p_0^2.$$
(5)
The same asymptotics could be obtained for $`\sigma _{yy},`$ so the indices are omitted for simplicity here and in following expressions. Summation over $`l`$ is carried out over all segments of the FS containing effective strips which correspond to $`\theta =\pi /2(v_z=0)`$ and the curvature $`K_l(\pi /2,\phi )`$ is supposed to take finite and nonzero value at any point of any effective strip. For a spherical FS $`p_0`$ equals the Fermi momentum $`p_F.`$ In realistic metals the two are not equal but have the same order of magnitude. In general, $`p_0`$ is determined by the Eq. (5). Then we can calculate the next term in the expansion of conductivity in powers of $`\omega /qv.`$ For a FS whose curvature everywhere is finite and nonzero we arrive at the result
$$\sigma _1(\omega ,q)=\sigma _0(q)\frac{i\omega }{qv_0}.$$
(6)
Here, the velocity $`v_0`$ has the order of the Fermi velocity $`v_F:`$
$`{\displaystyle \frac{1}{v_0}}`$ $`=`$ $`{\displaystyle \frac{2}{\pi ^2p_0^2}}{\displaystyle \underset{l}{}}{\displaystyle 𝑑\phi _{\alpha _l}^{\pi /2}\frac{d\theta \mathrm{cos}^2\phi \mathrm{sin}\phi }{\mathrm{cos}^2\theta }}`$
$`\times `$ $`\left[{\displaystyle \frac{1+\mathrm{cos}^2\theta /\mathrm{cos}^2\alpha _l}{|K_l(\pi /2,\phi )|v_l(\pi /2,\phi )}}{\displaystyle \frac{\mathrm{sin}^2\theta }{|K_l(\theta ,\phi )|v_l(\theta ,\phi )}}\right]`$
$``$ $`{\displaystyle \frac{2}{\pi ^2p_0^2}}{\displaystyle \underset{jl}{}}{\displaystyle 𝑑\phi _{\theta <\pi /2}\frac{d\theta \mathrm{cos}^2\phi \mathrm{sin}^3\theta }{|K_j(\theta ,\phi )|v_j(\theta ,\phi )\mathrm{cos}^2\theta }}.`$
Here, the lower limit $`\alpha _l`$ in the integral over $`\theta `$ in the first term takes on values determined by the FS shape $`(\alpha _l<\pi /2).`$ The second term corresponds to the contribution from the FS segments which do not include effective lines. For the spherical FS we have $`v_0=\pi v_F/4.`$
When the curvature at any effective line turns zero, it changes the asymptotics for the conductivity, as we show below. First, we assume that the curvature becomes zero at a whole effective line passing through one of the segments of the Fermi surface. Keeping in mind that $`z`$ axis in the chosen reference system runs along a high order symmetry axis we can present the relevant energy-momentum relation in the form:
$$E(𝐩)=E_1(p_x,p_y)+E_2(p_z).$$
(8)
Near the effective line where $`v_zE_2/p_z=0`$ we can approximate $`E_2(p_z)`$ as follows:
$$E_2(p_z)E_0\left(\frac{p_zp^{}}{p^{}}\right)^{2l};l1.$$
(9)
Here, $`E_0,p^{}`$ have the dimensions of the energy and momentum, respectively; $`p_z=p^{}`$ corresponds to the effective line.
In the vicinity of the effective line we can write the following expression for the curvature of the FS corresponding to the Eq. (8):
$$K(𝐩)=\frac{1}{v^4}\frac{v_z}{p_z}\left(v_y^2\frac{v_x}{p_x}+v_x^2\frac{v_y}{p_y}2v_xv_y\frac{v_x}{p_y}\right).$$
(10)
The value of the curvature at $`v_z=0`$ is determined by the factor $`v_z/p_z[(p_zp^{})/p^{}]^{2l2}.`$ The curvature becomes zero at the effective line when $`l>1.`$ Expressing this factor as a function of velocity (which is necessary to carry out integration over a region in the velocity space) we get $`v_z/p_z(v_z)^\beta `$ where $`\beta =1+1/(2l1).`$
So, we can use the following approximation for the curvature $`K(\theta ,\phi )`$ at $`\theta \pi /2:`$
$$K(\theta ,\phi )=W(\theta ,\phi )(\mathrm{cos}\theta )^\beta ,$$
(11)
In this expression, the function $`W(\theta ,\phi )`$ everywhere assumes finite and nonzero values, and the exponent $`\beta `$ takes on negative values which correspond to the line of zero curvature at $`\theta =\pi /2.`$ In the close vicinity of this line the FS is nearly cylindrical in shape. The closer $`\beta `$ to $`1,`$ the closer to a cylinder is the effective strip on the FS. The contribution to the conductivity from the nearly cylindrical segment on the FS is given by:
$`\sigma _a(\omega ,q)={\displaystyle \frac{ie^2\omega }{2\pi ^3\mathrm{}^3q}}[{\displaystyle }d\phi {\displaystyle _\alpha ^{\pi /2}}d\theta ({\displaystyle \frac{\mathrm{sin}^2\theta }{\left|W(\theta ,\phi )\right|v(\theta ,\phi )}}`$
$``$ $`{\displaystyle \frac{1}{\left|W(\pi /2,\phi )\right|v(\theta ,\phi )}}){\displaystyle \frac{\mathrm{cos}^2\phi \mathrm{sin}\theta (\mathrm{cos}\theta )^\beta }{\left(\omega /qv(\theta ,\phi )\right)^2\mathrm{cos}^2\theta }}`$
$`+`$ $`{\displaystyle }{\displaystyle \frac{d\phi \mathrm{cos}^2\phi }{\left|W(\pi /2,\phi )\right|v(\pi /2,\phi )}}{\displaystyle _\alpha ^{\pi /2}}{\displaystyle \frac{d\theta \mathrm{sin}\theta (\mathrm{cos}\theta )^\beta }{\left(\omega /qv(\theta ,\phi )\right)^2\mathrm{cos}^2\theta }}].`$
Using this asymptotic expression we can calculate the “anomalous” contribution to the conductivity $`\sigma _a(\omega ,q)`$ for small $`\omega /qv.`$ Introducing the largest magnitude of the velocity on the effective line $`v_a`$ we have:
$`\sigma _a(\omega ,q)`$ $`=`$ $`\rho \sigma _0(q)\left({\displaystyle \frac{\omega }{qv_a}}\right)^\beta \left[1i\mathrm{tan}\left({\displaystyle \frac{\pi \beta }{2}}\right)\right],`$ (13)
$`\rho `$ $``$ $`{\displaystyle \frac{1}{\pi p_0^2}}{\displaystyle \frac{d\phi \mathrm{cos}^2\phi }{|W(\pi /2,\phi )|}},`$ (14)
Comparing Eq. (14) with the definition for $`p_0^2`$ introduced earlier by Eq. (5) we see that $`\rho `$ is a dimensionless factor whose value is determined with the relative number of the “effective” electrons concentrated at the nearly cylindrical effective segment.
The value of the contribution to the conductivity from the ”anomalous” effective strip depends on the character of the curvature anomaly at given strip, and on the relative number of effective electrons concentrated here as shown in the Fig. 1. In this figure we display plots of $`|\sigma (\omega ,q)/\sigma _0(q)||1+\sigma _a(\omega ,q)/\sigma _0(q)|`$ versus $`\omega /qv.`$ When the parameter $`\rho `$ takes on values of the order or greater than 0.1 (the number of effective electrons associated with the anomalous sections on the FS is comparable to the total number of the effective electrons), the term $`\sigma _a(\omega ,q)`$ can predominate over $`\sigma _0(q)`$ and determine the conductivity value at large $`q.`$ This occurs when the shape parameter $`\beta `$ accepts values not too close to zero, and the curvature anomaly at the effective line is well pronounced. When either $`\rho `$ or $`\beta `$ or both are very small in magnitude, the main approximation to the conductivity is described with Eq. 5 as well as for a metal whose FS curvature is everywhere nonzero. Nevertheless, in such cases the term $`\sigma _a(\omega ,q)`$ also is important for it gives the first correction to the principal term in the expression for the conductivity.
Also, the anomalous contribution to the conductivity could appear when the FS is flattened at some points belonging to an effective segment. To avoid lengthy calculations we illustrate the effect of such points on the conductivity using a simple expression representing the energy momentum relation near the point of flattening $`_0(p_1,0,0):`$
$$E(𝐩)=\frac{p_1^2}{2m_1}\left(\frac{p_x^2}{p_1^2}\right)+\frac{p_2^2}{m_2}\left(\frac{p_y^2+p_z^2}{p_2^2}\right)^l$$
(15)
where $`p_1,p_2`$ have dimensions of momentum. When $`l=1`$ this expression corresponds to the ellipsoidal FS, and $`m_1,m_2`$ are the principal values of the effective mass tensor. The FS curvature equals:
$`K(𝐩)={\displaystyle \frac{l}{m_2v^4}}\left({\displaystyle \frac{p_y^2+p_z^2}{p_2^2}}\right)^{l1}`$
$`\times \left[{\displaystyle \frac{1}{m_1}}(v_y^2+v_z^2)+v_x^2{\displaystyle \frac{l(2l1)}{m_2}}\left({\displaystyle \frac{p_y^2+p_z^2}{p_2^2}}\right)^{l1}\right].`$ (16)
For $`l>1`$ the curvatures of both principal cross sections of the FS become zero at the point $`(p_1,0,0)`$ indicating the FS local flattening.
Turning to the spherical coordinates in the velocity space we can rewrite the expression (16) in the form:
$$K(\theta ,\phi )=W(\theta ,\phi )(\mathrm{cos}^2\theta +\mathrm{sin}^2\phi )^{(1\beta )/2}$$
(17)
where the shape parameter $`\beta =1+2/(2l1).`$ When $`l>1,`$ the FS curvature becomes zero at $`\theta =\pi /2,\phi =0`$ which correspond to the point $`_0.`$ The parameter $`\beta `$ takes on values from the interval $`(1,1)`$, and the more pronounced is the FS flattening near the point $``$ ( the greater is the value of $`l)`$ the closer is $`\beta `$ to $`1.`$ The “anomalous” contribution to the conductivity originating from the flattened segment of the FS has the form similar to Eq. (13), namely:
$$\sigma _a(\omega ,q)=\mu \sigma _0(q)\left(\frac{\omega }{qv(\pi /2,0)}\right)^\beta \left[1i\mathrm{tan}\left(\frac{\pi \beta }{2}\right)\right].$$
(18)
Here, $`\mu `$ is a small dimensionless factor proportional to the relative number of conduction electrons associated with the flattened part of the FS. Due to the smallness of $`\mu `$ the term (17) may be significant only when $`\beta 0(l1.5).`$ Otherwise, it could be neglected.
Now, we proceed in calculations of the surface impedance given by the expression (2). Under anomalous skin effect conditions the impedance can be represented as an expansion in inverse powers of the anomaly parameter $`(\xi 1).`$ Representing the conductivity as the sum of terms (5) and (6), we can calculate two first terms in the expansion of the surface impedance in inverse powers of the anomaly parameter:
$`Z`$ $``$ $`RiH={\displaystyle \frac{8i\omega }{c^2}}\delta {\displaystyle _0^{\mathrm{}}}𝑑t{\displaystyle \frac{1}{1it^3(1+it/\xi )}}`$
$``$ $`Z_0\left({\displaystyle \frac{\omega }{\omega _0}}\right)^{2/3}\left[1i\sqrt{3}{\displaystyle \frac{2}{3}}\left({\displaystyle \frac{\omega }{\omega _0}}\right)^{2/3}\left(1+\sqrt{3}\right)\right]`$
where $`\delta =(c^2\mathrm{}^3/e^2p_0^2\omega )^{1/3}`$ is the skin depth,
$$Z_0=\frac{8\pi }{3\sqrt{3}}\frac{v_0}{c^2};\xi =\frac{v_0}{\omega \delta }\left(\frac{\omega }{\omega _0}\right)^{2/3}$$
is the anomaly parameter, and frequency $`\omega _0`$ equals:
$$\omega _0=\left(\frac{v_0}{\mathrm{}}\right)^{3/2}\frac{ep_0}{c}.$$
(20)
Keeping in mind that $`v_0v_F`$ and $`p_0p_F`$ we can roughly estimate the characteristic frequency $`\omega _0.`$ In good metals the electron density has the order of $`10^{21}10^{22}`$ cm<sup>-3</sup>, so $`\omega _010^{12}10^{13}s^1.`$ This is significantly smaller that the plasma frequency $`\omega _p`$ which in good metals is of the order of $`10^{15}10^{16}s^1.`$ As one would expect, the inequality $`\omega \omega _0(\xi 1)`$ agrees with the general requirement on frequencies $`\omega \omega _p,`$ and could be satisfied at $`\omega 10^{10}10^{11}s^1.`$
The expression of the form (19) was first obtained by R. B. Dingle (see 3 ) within the isotropic model of metal. Later it was generalized to be applied to realistic metals, assuming that their FSs do not include segments of zero curvature 8 . For such FSs the frequency dependence of the surface impedance has the same character, as for a Fermi sphere. The main approximation of the surface impedance is proportional to $`\omega ^{2/3}`$ and the first correction to it is proportional to $`\omega ^{4/3}.`$
When the FS includes locally flattened or nearly cylindrical segments the asymptotic expression for the surface impedance changes. The effect of this anomalous local geometry of the FS on the impedance is especially strong for $`\beta <0.5,\rho \stackrel{>}{}0.1.`$ Under these conditions the “anomalous” contribution dominates over the other terms in the expression for conductivity and determines the principal term of the surface impedance. As a result we have:
$$ZZ_0\zeta (\beta )\left(\frac{\omega }{\omega _0}\right)^{2/(3+\beta )}$$
(21)
where
$`\zeta (\beta )`$ $`=`$ $`{\displaystyle \frac{3\sqrt{3}\rho ^{1/3+\beta }}{3+\beta }}\left(\mathrm{cos}{\displaystyle \frac{\pi \beta }{2}}\right)^{1/(3+\beta )}`$ (22)
$`\times `$ $`\left[\mathrm{cot}\left({\displaystyle \frac{\pi }{3+\beta }}\right)i\right].`$
The surface impedance described with (21) differs in magnitude from that of a conventional metal whose FS does not include zero curvature segments. Frequency dependence of the surface impedance also changes as shown in the Fig. 2. Now it is proportional to $`\omega ^{2/(\beta +3)}.`$ For a nearly cylindrical effective strip the exponent $`2/(\beta +3)`$ varies in the interval $`(0.8,1)`$ where the value $`1`$ corresponds to the precisely cylindrical strip. So, when the significant part of effective electrons is associated with a nearly cylindrical effective strip, the impedance should slower increase with increase of frequency than in a “conventional” case.
Now we consider more realistic case when either $`\rho `$ or $`\beta `$ or both take on values close to zero (a zero curvature segment on the effective part of the FS is narrow and/or the curvature anomaly is only moderately pronounced). In this case the anomalous contribution (13) is the first correction to the main approximation for the conductivity, and it determines the first correction to the approximation for the surface impedance:
$$Z=Z_0\left(\frac{\omega }{\omega _0}\right)^{2/3}(1i\sqrt{3})+\mathrm{\Delta }Z.$$
(23)
Here,
$`\mathrm{\Delta }Z`$ $``$ $`\mathrm{\Delta }Ri\mathrm{\Delta }HZ_0\eta (\beta )\left({\displaystyle \frac{\omega }{\omega _0}}\right)^{2(\beta +1)/3};`$ (24)
$`\eta (\beta )`$ $`=`$ $`{\displaystyle \frac{\rho (\beta +1)}{\sqrt{3}}}{\displaystyle \frac{1}{\mathrm{cos}(\pi \beta /2)}}\left\{\mathrm{cot}\left({\displaystyle \frac{\pi (\beta +1)}{3}}\right)+i\right\}.`$
We can apply this result (24) to describe the contribution to the surface impedance from narrow or weakly developed nearly cylindrical strip and also from a point of flattening located on the effective segment of the FS. The correction to the main approximation of the surface impedance now is proportional to $`\omega ^{2(\beta +1)/3}`$ as we show in the Fig. 3. The presented analysis shows that the surface impedance of a semiinfinite metal whose FS has points or lines of zero curvature can be described by the formulas (21)–(24). The obtained asymptotic expressions indicate that an anomaly of curvature on an effective line changes frequency dependence of the surface impedance, and under certain conditions it can essentially change its magnitude. This follows from the discussed above relation between the curvature of the FS and the number of effective electrons.
## III III. conclusion
The concept of Fermi surface is recognized as one of the most meaningful concepts in condensed matter physics, providing an excellent insight in the main electronic properties of conventional metals and other materials with metallic-like conductivity. Extensive studies based on experimental data concerning effects responsive to the structure of electronic spectra in metals and using advanced computational methods were carried out to restore the FS geometries. These efforts were resulted in the impressive mapping of the FSs of conventional metals. However, in the course of these studies a comparatively little attention was paid to fine local features in the FS geometries including zero curvature lines and/or points of flattening. These local curvature anomalies do not significantly affect main geometrical characteristics of FSs (such as connectivity, locations of open orbits, sizes and arrangements of sheets) which are usually determined from the standard experiments. Therefore these local features could be easily missed when a FS is restored, if one does not expect them to be present, and does not pay special attention to keep them in the resulting FS image. So, it is important to explore possible experimental manifestations of the FS curvature anomalies. Adoption of the phenomenological models is justified in these studies as far as these models are based on reasonable assumptions concerning the FS geometry. Actually, phenomenological models were commonly used to develop the theory of “standard” effects such as de Haas-van Alphen effect, which were (and are) employed as tools to obtain informations concerning FSs shapes 19 . This approach to fermiology does not contradict that one based on electron band structure calculations. It supplements the latter. Such supplementing analysis could bring new insight in the physical nature and origin of some physical effects including these considered in the present work, and show their usefulness in studies of the FSs geometries.
As for the particular models adopted in the present work they may be reasonably justified within a nearly free electron approximation. Adopting the nearly free-electrons approach we arrive at the energy-momentum relation for conduction electrons:
$$E=\frac{𝐤^\mathrm{𝟐}}{2m}+\frac{𝐠^\mathrm{𝟐}}{2m}\frac{1}{m}\sqrt{(𝐤𝐠)^\mathrm{𝟐}+m^2V^2},$$
(26)
where $`m`$ is the effective mass, $`𝐤=𝐠𝐩;𝐠=\mathrm{}𝐆/2;𝐆`$ is a reciprocal lattice wave vector; $`V`$ is Fourier component of the potential energy of electron in the lattice field which corresponds to the vector $`𝐆.`$ Within the nearly free-electron model the energy $`V`$ is assumed to be small compared to the Fermi energy $`E_F,`$ so we introduce a small parameter $`ϵ=\sqrt{V/E_F}`$.
The corresponding FS looks like a sphere with “knobs” located at those segments which are close to the boundaries of the Brillouin zone. Inflection lines of zero curvature pass along the boundaries between the knobs and the main body of the FS. A FS segment including a knob and its vicinity is axially symmetric, and the symmetry axis is directed along the corresponding reciprocal lattice vector. In further analysis we single out such segment to consider it separately. For certainty we choose the coordinate system whose $`\mathrm{`}\mathrm{`}x\mathrm{"}`$ axis is directed along the reciprocal lattice vector. Within the chosen segment the FS curvature is described with the expression:
$$K=\frac{m^2v_x^2+p_{}^2dv_x/dp_x}{(p_{}^2+m^2v_x^2)^2}$$
(27)
where $`p_{}^2=p_y^2+p_z^2.`$
Equating the FS curvature to zero, and using the energy-momentum relation (26) we find the values $`p_{x0}`$ and $`p_0`$ corresponding to the inflection line. We get:
$$p_{x0}=p_F(1ϵ^2/\sqrt{2}),p_0p_Fϵ,$$
(28)
where $`p_F`$ is the radius of the original Fermi sphere. Now, we can expand the variable $`p_x`$ in powers of $`(p_{}p_0)`$ near the zero curvature line. Taking into account that $`d^2p_x/dp_{}^2`$ turns zero at points belonging to the inflection line, and keeping the lowest-order terms in the expansion, we obtain:
$$p_xp_{x0}\frac{ϵ}{\sqrt{2}}(p_{}p_0)\frac{p_F}{\sqrt{2}ϵ}\left(\frac{p_{}p_0}{p_F}\right)^3.$$
(29)
Substituting this approximation into Eq. (26) we arrive at the following energy-momentum relation:
$$E(𝐩)=\frac{p_x^2}{2m}+\frac{2}{ϵ}\frac{p_F^2}{2m}\left(\frac{p_{}p_0}{p_F}\right)^3.$$
(30)
The latter could be employed near the zero curvature line where $`p_0p_F.`$ Omitting $`p_0`$ we arrive at the energy-momentum relation of the form (15) where $`l=3/2.`$ Also, we can compare the equations describing cross-sections $`p_y=0`$ of the FS corresponding to the Eq. (30) and our phenomenological model (8), (9). Again, we see that these equations are in agreement with each other provided that $`l=3/2.`$
There is an experimental evidence that ”necks” connecting quasispherical pieces of the FS of copper include nearly cylindrical belts 19 . It is also likely that the FS of gold possesses the same geometrical features for it closely resembles that of copper. As for possible flattening of the FS, experiments of 20 ; 21 ; 22 on the cyclotron resonance in a magnetic field normal to the metal surface give grounds to conjecture that such anomalies could be found on the FSs of cadmium, zinc and even potassium. Another group of materials where we can expect the FS curvature anomalies to be manifested includes layered structures with metallic-type conductivity (e.g. $`\alpha (BEDTTTF)_2Mhg(SCH)_4`$ group of organic metals). Fermi surfaces of these materials are sets of rippled cylinders, isolated or connected by links. There exists experimental evidence that the quasi-two-dimensional FSs of some organic metals include segments with zero curvature 23 . Also, recent investigations give grounds to expect the FSs of some new conducting materials include flattened segments 24 ; 25 ; 26 .
The most important result of the present work is that it shows how such fine geometric features as points of flattening and/or zero curvature lines could be manifested in experiments on the anomalous skin effect. It is shown that when the FS includes nearly cylindrical segments or it is flattened at some points, qualitative changes may occur in frequency dependencies of the surface impedance under the anomalous skin effect. Being observed in experiments, such unusual frequency dependencies would indicate the presence of zero-curvature lines and points on the FS, and display their location. Also, analyzing these frequency dependencies, the shape parameter $`\beta `$ could be found giving additional information on the FSs local structure. This information may be used in further studies of the FSs geometries.
## IV Acknowledgments
The author thanks G. M. Zimbovsky for help with the manuscript. This work was supported by NSF Advance program SBE-0123654, NSF-PREM 0353730, and PR Space Grant NGTS/40091. |
warning/0506/cond-mat0506801.html | ar5iv | text | # Condensation and vortex formation in Bose-gas upon cooling
## I Introduction
The ideas of the kinetics of phase transitions have been thoroughly developed for first-order phase transitions and envisage the existence of the metastable phase itself and an equilibrium critical nucleus. The corresponding theory was worked out in Bec , Zel and described in detail in Lan . However, theoretical concepts concerning the kinetics of second-order phase transitions, where these two facts do not exist, have been developed insufficiently. In the Lif was proposed a certain special model for the formation of an ordered phase after the fast phase-transition stage in the ”short-range” order in the presence of only two types of ordering.
The interest in the problem of a phase transition upon a fast change in external parameters (e.g., temperature) has been aroused in connection with the cosmological ideas of the Big Bang, where the rapidly expanding Universe must be cooled and pass through a series of phase transformations accompanied by a change in the symmetry of physical fields Zell . It was proposed that the kinetics of these transformations can be modeled in condensed matter Kib .
In the Zur was proposed a theory of the second-order phase transition upon a rapid change in temperature in liquid $`{}_{}{}^{4}\mathrm{He}`$. The main assumption in the proposed mechanism is about the ”critical retardation” of all processes in the vicinity of the transition temperature and ”fast” formation of the nuclei of a new phase upon the subsequent cooling. This gives rise to a large number of defects on the order of the number of fluctuations far above the transition point.
However, no retardation in the formation of a new phase has been detected experimentally; the critical retardation is associated with the duration of the equilibration process at macroscopic distances, which is insignificant for the nonuniform process of formation of a new phase.
In this work, we consider the transition to a new phase via the evolution of fluctuations on scales much smaller than the correlation length, which can occur quite rapidly even in the vicinity of the critical temperature. The transition kinetics in this case is found to be directly related to the cooling process itself. In our preliminary paper BIS this approach was proposed for the specific problem of Bose condensation of a weakly interacting Bose gas. The appropriate set of equations governing critical fluctuations was derived. In the present work a detailed analysis of the main equations is performed including numerical calculations and a generalization to the case of the two-dimensional exciton gas. We will consider the formation of a condensate in the model of a weakly nonideal Bose gas with external cooling and demonstrate an analogy with first-order phase transitions.
A similar approach to the problem of the wave nucleation rate due to thermal noise was considered in Henry . But the problem of the present work requires the noise to be connected to the random thermal fluxes. In this case it is appropriate to use a more general approach for the instanton formation namely local Hamilton equations rather than the Lagrangian equations used in Henry (see also Marder ). In Cher the problem of large negative gradients in Burgers turbulence was considered, which has some resemblance to the differential equations discussed in our work, but the boundary problem is quite different.
## II Dilute Bose-gas upon cooling
The standard theory of a weakly nonideal Bose gas involves a Hamiltonian of the form :
$$\widehat{H}=\underset{p}{}\frac{\widehat{p}^2}{2m}\widehat{a_p}^+\widehat{a_p}+\frac{2\pi \mathrm{}^2a_0}{m}\underset{p}{}\widehat{a_{p_4}}^+\widehat{a_{p_3}}^+\widehat{a_{p_2}}\widehat{a_{p_1}},$$
(1)
where $`a_0`$ – is the scattering amplitude having the atomic scale and $`m`$ – is the atomic mass. The properties of such gas for a small density $`n`$ (determined by the gas parameter $`\eta =na_{0}^{}{}_{}{}^{3}1`$ ) are close to the properties of an ideal Bose gas with the transition temperature Ld :
$$T_c=\frac{3.31}{\sqrt{2}}\frac{\mathrm{}^2n^{2/3}}{m}.$$
(2)
At temperatures below the transition point, the ideal Bose gas has a pressure depending only on the temperature:
$$P_{id}=0.0851\frac{m^{3/2}T^{5/2}}{\mathrm{}^3},$$
(3)
which corresponds to zero isothermal sound velocity.
Considering the finite scattering amplitude we can write the qualitative equation of state below the transition point as
$$P=P_{id}(T)+\frac{\mathrm{}^2a_0n^2}{m}.$$
(4)
We omitted the insignificant constant factor in the second term.
The entire kinetics is essentially determined by the Bose-gas cooling mechanism. We will consider a simple model where the Bose gas is in a certain solid matrix with which it only slightly interacts. Such a situation may take place, for example, for the exciton gas in a crystal. The crystal can be rapidly cooled to a low temperature; in this case, the Bose-gas cooling proceeds via phonon emission. Assuming that the heat capacity of the crystal is large compared to the Bose gas, we can disregard the presence of thermal phonons in the crystal and their effect on the Bose gas. As a result, we obtain a uniform energy-loss mechanism, which is described by a phenomenological quantity $`T/\tau _{ph}`$. The other models of cooling necessitate the analysis of heat transfer at the sample boundaries, which is a much more complicated problem. The loss rate $`1/\tau _{ph}`$ is determined by the collisions of particles with each other and by the interaction with phonons, which will be regarded as weak:
$$1/\tau _{ph}1/\tau _{tr}.$$
Since $`1/\tau _{ph}nv_T\sigma _{ph}`$ , $`1/\tau _{tr}nv_Ta_{0}^{}{}_{}{}^{2}`$ ($`v_T`$ s the thermal velocity), this means that
$$\sigma _{ph}a_{0}^{}{}_{}{}^{2},$$
where $`\sigma _{ph}`$ is the cross section for scattering with phonon emission, which corresponds to the weak interaction of Bose gas with the crystal.
In view of the smallness of quantity $`1/\tau _{ph}`$ the evolution of the Bose system is slow; in particular, we assume that the acoustic wavelength $`c\tau _{ph}v_T\tau _{ph}`$ is large compared to the characteristic length $`\sqrt{\chi \tau _{ph}}`$ , where $`\chi `$ \- is the thermal diffusivity:
$$\frac{\sqrt{\chi \tau _{ph}}}{v_T\tau _{ph}}\sqrt{\frac{l^2}{v_{T}^{}{}_{}{}^{2}\tau _{tr}\tau _{ph}}}\sqrt{\frac{\tau _{tr}}{\tau _{ph}}}1$$
(5)
($`l`$ stands for the mean free path). This makes it possible to assume that the fluctuation evolution occurs at a constant pressure that coincides with the thermodynamically equilibrium pressure.
It follows from Eq.(4) that the density variation $`\delta n`$ in the fluctuation region is related to a change in temperature by
$$\frac{\delta n}{n}=\frac{\delta T}{T}\frac{1}{\eta ^{\frac{1}{3}}}.$$
(6)
The relative density fluctuation is large compared to the relative temperature fluctuation in the temperature range $`T<T_c`$. This leads to a rapid increase in the reciprocal phonon time
$$\delta \frac{1}{\tau _{ph}}\frac{\delta T}{T}\frac{1}{\eta ^{\frac{1}{3}}}\frac{1}{\tau _{ph}^{}{}_{}{}^{0}}$$
(7)
(where $`1/\tau _{ph}^{}{}_{}{}^{0}=nv_T\sigma _{ph}`$) with decreasing temperature and enhancement of cooling in the fluctuation region. For this reason, we will disregard the phonon emission in the region far from the developed fluctuations, assuming that
$$\frac{1}{\tau _{ph}}\frac{\delta T}{T_c}\frac{1}{\eta ^{\frac{1}{3}}}\frac{1}{\tau _{ph}^{}{}_{}{}^{0}}U(T_cT),$$
(8)
where $`U(T_cT)=1`$ for $`\delta T=TT_c<0`$ and $`U(T_cT)=0`$ for $`TT_c>0`$.
This allows us to consider the problem of fluctuation kinetics within the framework of the theory of hydrodynamic fluctuations by supplementing the hydrodynamic equations with the energy flux carried away as a result of phonon emission:
$$\frac{T}{\tau _{ph}^{}{}_{}{}^{0}}\frac{\delta n}{n}=\frac{T_cT}{\tau _{ph}}U(T_cT),$$
(9)
where $`1/\tau _{ph}=(1/\tau _{ph}^{}{}_{}{}^{0})(1/\eta ^{1/3})`$. In view of the constancy of pressure, we can describe the evolution of temperature fluctuations by the heat conduction equation
$$nc_p\left(\frac{T}{t}+\stackrel{}{v}\frac{T}{\stackrel{}{r}}\right)=(\varkappa T)+\frac{TT_c}{\tau _{ph}}nc_pU(T_cT),$$
(10)
where the energy flux carried away by phonons is added, $`\varkappa `$ – is the heat conductivity, $`c_p`$ – is the specific heat per particle under a constant pressure. This equation contains the drift term with mass velocity $`v(𝐫)`$, which appears due to the high density in the fluctuation core. In the following analysis, this term will be omitted as a higher-order term in fluctuation. We are interested in the temperature-field fluctuations and their time evolution. To analyze these fluctuations, we must introduce random heat fluxes $`𝐪`$ Lan , Liff , i.e., the Langevin term $`𝐪`$. These fluxes are delta-correlated (i.e., correlated at distances and time intervals smaller than the hydrodynamic scales). In the case considered, this is ensured by the fact that the time $`\tau _{ph}`$ and distance $`\sqrt{\chi \tau _{ph}}`$ ($`\chi `$ – is thermal diffusivity) are larger than the microscopic characteristics.
The probability $`W_t(T(r))`$ of realizing the given configuration $`T(r)`$ fluctuation field at time $`t`$ obeys the Fokker-Planck equation in variational derivatives Kl
$`{\displaystyle \frac{}{t}}W={\displaystyle }{\displaystyle \frac{\delta }{\delta T(r)}}[{\displaystyle \frac{\chi T_{\mathrm{}}^{}{}_{}{}^{2}}{nc_p}}^2{\displaystyle \frac{\delta }{\delta T(r)}}W`$
$`+[\chi ^2T+U(T_cT){\displaystyle \frac{TT_c}{\tau _{ph}}}]W]d^3r.`$ (11)
In the absence of the interaction with phonons, the stationary solution to this equation coincides with the result obtained in the thermodynamic theory of fluctuations. The quantity
$$\chi ^2T+U(T_cT)\frac{TT_c}{\tau _{ph}}=\frac{T}{t}$$
(12)
is the temperature-variation rate upon the deviation from the mean value $`T=T_{\mathrm{}}`$.
We assume that fluctuations occur at a fixed temperature $`T_{\mathrm{}}>T_c`$. Fluctuations with $`\mathrm{}T=TT_{\mathrm{}}T_{\mathrm{}}`$ occur quite frequently and are characterized by a certain (in fact, stationary) spatial distribution that determines the value of $`W_t(T)`$. In view of the normalization, the latter quantity gives the number of small fluctuations in a unit volume. However, rare large-amplitude fluctuations with $`TT_cT_{\mathrm{}}`$, $`T<T_c`$, also sometimes occur, initiating the effective cooling by phonons, so that the fluctuation becomes irreversible and the nucleus of a new phase appears. Our goal is to calculate the probability of such fluctuations in a unit volume per unit time. Since they are infrequent and the distribution at small $`T_{\mathrm{}}T`$ is stationary, one can use the method of characteristics to determine the exponentially low probability of formation of such a nucleus (instanton for the Fokker-Planck equation). An important difference from the theory of nucleation in the first-order phase transition is that the probability of instanton formation in this case is determined by the cooling process.
## III A toy model
To clarify the situation, let us consider the instanton solution in the case of one degree of freedom, for which the Fokker-Planck equation has the form
$$\frac{W}{t}=\frac{}{x}\left(D\frac{W}{x}vW\right),$$
(13)
where $`D`$ is the constant diffusion coefficient and $`v(x)`$ is the macroscopic variation rate of the quantity $`x`$ with allowance for its relaxation upon the deviation from equilibrium and for an external effect (analogue of phonon emission). Setting $`W=e^S`$, and assuming that the moduli of S and its first derivative are large, we obtain, to leading terms, the equation
$$\frac{S}{t}=D\left(\frac{S}{x}\right)^2v\frac{S}{x}\frac{v}{x}.$$
(14)
This is the Hamilton-Jacobi equation with the Hamiltonian ($`S/x=p`$)
$$H(\frac{S}{x},x)=Dp^2+pv+\frac{v}{x}.$$
(15)
The Hamilton equations are the characteristics of this equation in partial derivatives,
$$\frac{dx}{dt}=2Dp+v,$$
(16)
$$\frac{dp}{dt}=\frac{dv}{dx}p\frac{d^2v}{d^2x}.$$
(17)
The contribution of the velocity divergence to the Hamiltonian is significant only in the vicinity of the point $`v=0`$. We are interested in the special solution that passes through the equilibrium point $`p=0`$, $`v=0`$. In the 1D Fokker-Planck equation, one can eliminate the term with a first derivative by substitution; in this case, we have an analogy with quantum mechanics and can use the well-known results. Nevertheless, we will use direct estimates in the vicinity of $`v=0`$.
In the Hamilton equation, the energy is conserved. In view of the smallness of the divergence term, this gives $`H=Dp^2+pv=0`$, whence $`p=v/D`$ and
$$S=\frac{v^2}{D}𝑑t=\underset{0}{\overset{x^{}}{}}\frac{v}{D}𝑑x.$$
(18)
We assume that the velocity $`v(x)`$ is a convex-down function with two zeros (stable at zero and unstable at $`x^{}`$ ($`x^{}>0`$). Such a shape of the function $`v(x)`$ is ensured by the entire cooling process, including phonon emission. For $`x>x^{}`$, the solution tends to larger values of x, while the action is gathered from zero to $`x^{}`$, where $`v<0`$. In the vicinity of $`x^{}`$, we must take into account the quantity $`dv/dx`$ . For large values of $`x`$, $`p^2`$ can be ignored, yielding $`p(dv/dx)/v`$,
$$SS_0\mathrm{ln}(v/v_0),$$
(19)
where $`v_0`$ is the effective velocity in the region where the solutions for $`x<x^{}`$ and $`x>x^{}`$ match. The solution $`S_0`$ itself has the form $`e^Sv_0e^{S_0}/v`$, and current $`jv_0e^{S_0}`$. One can estimate the value of $`v_0`$, assuming that all terms in the Hamiltonian $`H`$ are of the same order of magnitude:
$$Dp^2vp\frac{dv}{dx}\frac{v_{max}}{x^{}}$$
(20)
which gives
$$v_0\sqrt{\frac{D|v_{max}|}{x^{}}}\frac{|v_{max}|}{\sqrt{|S_0|}}.$$
(21)
In the many-dimensional case, the situation is the same,
$$\frac{dx^i}{dt}=2D^{ij}p_j+v^i,$$
(22)
$$\frac{dp_i}{dt}=\frac{v^k}{x^i}p_k\frac{(\mathrm{div}\stackrel{}{v})}{x^i},$$
(23)
where $`p=0`$ at the beginning and $`p0`$ at the end of the trajectory. Consequently, $`|p|`$ reaches its maximal value somewhere on the trajectory. At this point, the matrix $`v_i/x_k`$ has one zero eigenvalue and $`𝐩`$ is tangent to the corresponding eigenvector; subsequently, the trajectory passes to the neighborhood of the point corresponding to zero velocity $`v`$. This leads to the definition of the critical fluctuation (instanton) as a solution passing through the point $`𝐱=𝐩=0`$, whereupon $`𝐩0`$ for $`|x|\mathrm{}`$ as $`1/v`$, retaining the probability flux at a constant level.
## IV Optimal fluctuation
An analogous procedure can be carried out for the field as well. In this case, the Hamiltonian has the form, in accordance with Eq.(II),
$$H=p(𝐫)\left[\frac{\chi T_{\mathrm{}}^{}{}_{}{}^{2}}{nc_p}^2p(𝐫)+\chi ^2T+U(T_cT)\frac{TT_c}{\tau _{ph}}\right]d^3r$$
(24)
with the Hamilton equations
$$\frac{T}{t}=\frac{2\chi T_{\mathrm{}}^{}{}_{}{}^{2}}{nc_p}^2p(𝐫)+\chi ^2T+U(T_cT)\frac{TT_c}{\tau _{ph}},$$
(25)
$$\frac{p}{t}=\chi ^2p\frac{p}{\tau _{ph}}U(T_cT).$$
(26)
Here, $`p=\delta S/\delta T(𝐫)`$. Eqs. (25), (26) define the critical fluctuation and can be reduced to dimensionless variables by the substitutions $`\xi =r/\sqrt{\chi \tau _{ph}}`$, $`\tau =t/\tau _{ph}`$,
$$\mathrm{\Theta }=\frac{TT_c}{T_{\mathrm{}}T_c},p=\frac{nc_p(T_{\mathrm{}}T_c)\mathrm{\Pi }}{T_{}^{2}{}_{\mathrm{}}{}^{}},$$
where $`\mathrm{\Theta }`$ and $`\mathrm{\Pi }`$ are the new dimensionless fields. In this case, the dimensionless equations have the form
$$\frac{\mathrm{\Theta }}{\tau }=^2\mathrm{\Theta }+\mathrm{\Theta }U(\mathrm{\Theta })+2^2\mathrm{\Pi },$$
(27)
$$\frac{\mathrm{\Pi }}{\tau }=^2\mathrm{\Pi }\mathrm{\Pi }U(\mathrm{\Theta }).$$
(28)
The solution should fulfill the conditions $`\mathrm{\Pi }_\xi \mathrm{}0`$, $`\mathrm{\Pi }_\tau \mathrm{}0`$, $`\mathrm{\Theta }_\tau \mathrm{}1`$, $`\mathrm{\Theta }_\xi \mathrm{}1`$ and pass through the neighborhood of $`\mathrm{\Theta }/\tau 0`$, $`\mathrm{\Pi }0`$ at $`\tau \tau ^{}`$. Later the fluctuation is developed by cooling, while the random fluxes can be neglected and $`\mathrm{\Theta }/\tau ^2\mathrm{\Theta }+\mathrm{\Theta }U(\mathrm{\Theta })`$. With exponential precision we can assume that $`\mathrm{\Theta }/\tau 0`$ at $`\tau +\mathrm{}`$, and $`\mathrm{\Theta }`$ tends to the stationary solution $`\mathrm{\Theta }_{st}`$ of the thermal diffusion equation, Eq.(27):
$`\mathrm{\Theta }_{st}=\mathrm{sin}(\xi )/\xi ,\xi <\pi ,`$
$`\mathrm{\Theta }_{st}=1\pi /\xi ,\xi >\pi .`$ (29)
The difficulties with the numerical solution of this boundary value problem are due to the instability of Eq.(28) for ascending time whereas Eq.(27) is unstable for descending time. Therefore, it is impossible to find numerically the solution of the Cauchy problem in either direction of time. We briefly describe our adopted procedure. At early stages of the evolution, when $`\mathrm{\Theta }>0`$ everywhere in space, it is easy to check the validity of the relation
$$\mathrm{\Theta }=1\mathrm{\Pi },$$
(30)
which coincides with the thermodynamical theory of temperature fluctuations. The function $`\mathrm{\Pi }`$ grows with time according to Eq. (28) (thermal diffusion equation with negative time derivative). We can assume that at $`\tau =0`$ the maximum of $`\mathrm{\Pi }`$ will reach 1. After this Eq. (30) is no longer valid. We can consider Eq. (28) as a Schroedinger equation with imaginary time and $`\mathrm{\Pi }0`$ at infinite time. At large $`\tau `$, the function $`\mathrm{\Theta }`$ will be close to the stationary solution, Eq.(IV), which becomes zero at $`\xi =\pi `$. This means that $`\mathrm{\Pi }`$, at large $`\tau `$, has an asymptotic proportional to
$$\mathrm{\Pi }_{inf}=\mathrm{exp}(|\lambda |\tau )\mathrm{\Psi }_\lambda ,$$
(31)
where $`\mathrm{\Psi }_\lambda `$ corresponds to the eigenfunction with the negative eigenvalue, $`\lambda =0.4576`$, according to the Schroedinger equation with the potential $`U(\pi \xi )`$ (all other states will grow with $`\tau `$).
Let us denote by $`r^{}(\tau )`$ the space-point for which $`\mathrm{\Theta }=0`$ at time $`\tau `$. The curve $`r^{}(\tau )`$ is a function of $`\tau `$ starting at small $`\tau `$ as square root of $`\tau `$ (because $`\mathrm{\Theta }`$ has a minimum at $`\xi =0`$) and tending exponentially to $`\pi `$ at large $`\tau `$ (according to the Schroedinger-equation analogy, Eq.(31)). We don’t know the exact form of $`r^{}(\tau )`$ but we can choose some probe function with the same asymptotic behavior. Having such a probe function, we can numerically solve Eq.(28) integrating it backward in time (it is unstable while integrating it forward in time) with the condition $`\mathrm{\Pi }_{inf}=\alpha \mathrm{\Psi }_\lambda `$ (of form Eq. (31)) (the prefactor $`\alpha `$ should be chosen to satisfy Eq. (30) at $`\tau =0`$).
Using this function $`\mathrm{\Pi }`$, we can numerically solve Eq. (27) in two space-time regions independently. The first one is the internal region $`\xi <r^{}(\tau )`$ the other one is $`\xi >r^{}(\tau )`$ (the external region). $`\mathrm{\Theta }1`$ for $`\xi \mathrm{}`$ (for the external region) and $`\mathrm{\Theta }(r^{})=0`$ (for the both regions). For the exact function $`r^{}`$ the space derivative $`\mathrm{\Theta }^{^{}}=\mathrm{\Theta }/\xi `$ will be continuous. For our probe function $`r^{}`$ there will be some jump of the space derivative $`\mathrm{\Theta }^{^{}}`$ at $`r^{}(\tau )`$. Thus, after performing the calculations, we will have in general a nonzero jump-function
$$\mathrm{\Delta }\mathrm{\Theta }^{^{}}=\mathrm{\Theta }^{^{}}(r^{}0)\mathrm{\Theta }^{^{}}(r^{}+0)$$
depending on the choice of the curve $`r^{}`$. Afterwords we should search for a more exact $`r^{}`$ in order to minimize $`max|\mathrm{\Delta }\mathrm{\Theta }^{^{}}|`$. Some steps within the framework of this procedure have been performed.
We use a space-time grid 2000x2000. The grid spacing was taken as $`\delta \xi =0.015`$ for the space coordinate and $`\delta \tau =0.01`$ for the time coordinate. First, we perform calculations backward in time for $`\mathrm{\Pi }`$ using a standard implicit numerical scheme where the space derivatives are calculated for the final time of each time step. There are some modifications of the space grid in the vicinity of $`r=r^{}(t)`$, for a better finite difference representation of the Laplacian. Then we use the analogous scheme for Eq.(27) and find $`\mathrm{\Theta }`$ in the internal and external regions integrating forward in time and obtain $`\mathrm{\Delta }\mathrm{\Theta }^{^{}}`$. We have defined our probe function $`r^{}(t)`$ by a number of parameters. To obtain an exact solution we need, of course, an infinite number of parameters. In practice, for a reasonable precision we only need a few. The simplest form which obeys the asymptotic behavior is
$`r^{}=\alpha \sqrt{\tau }+\beta \tau ,\tau <\tau _0,`$
$`r^{}=\pi \delta \mathrm{exp}(|\lambda |\tau ),\tau >\tau _0.`$
The parameters should be chosen such that $`r^{}`$ is continuous and smooth at $`\tau =\tau _0`$. This means that we have two free parameters in this case (e.g. $`\alpha `$ and $`\tau _0`$). After minimization of $`max|\mathrm{\Delta }\mathrm{\Theta }^{^{}}|`$ with respect to these two parameters we find right and left derivatives $`\mathrm{\Theta }^{^{}}(r^{}\pm 0)`$ (see Fiq.1).
We can improve our results adding a new term $`\gamma \tau ^2`$:
$`r^{}=\alpha \sqrt{\tau }+\beta \tau +\gamma \tau ^2,\tau <\tau _0,`$
$`r^{}=\pi \delta \mathrm{exp}(|\lambda |\tau ),\tau >\tau _0.`$ (32)
We have found that for $`\alpha =2.28`$, $`\gamma =0.016`$, $`\tau _0=2.55`$ there is an acceptable minimum of $`|\mathrm{\Delta }\mathrm{\Theta }^{^{}}|`$ (see Fig.2). The corresponding jump-functions $`\mathrm{\Delta }\mathrm{\Theta }^{^{}}`$ for the two and three parameter cases (for comparison) are plotted in Fig.3. One can see that the inclusion of this additional term decreases the maximum deviation, $`max[|\mathrm{\Delta }\mathrm{\Theta }^{^{}}|]`$, by a factor of 3.
## V The optimal fluctuation probability
The solution of Eqs. (27),(28) allows us to calculate the action
$`S`$ $`=`$ $`{\displaystyle p\frac{T}{t}d^3r𝑑t}{\displaystyle H𝑑t}`$
$`=`$ $`s_0{\displaystyle \frac{nc_p(T_{\mathrm{}}T_c)^2}{2T_{}^{2}{}_{\mathrm{}}{}^{}}}(\sqrt{\chi \tau _{ph}})^3.`$ (33)
The negative constant $`s_0`$ is the dimensionless action
$$s_0=\underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑\tau _\tau \mathrm{\Theta }\mathrm{\Pi }d^3\xi \underset{\mathrm{}}{\overset{\mathrm{}}{}}H_1𝑑\tau ,$$
(34)
where the space integral is taken over the whole space. Here
$$H_1=\mathrm{\Pi }(𝐫)\left[^2\mathrm{\Pi }(𝐫)+^2\mathrm{\Theta }+U(\mathrm{\Theta })\mathrm{\Theta }\right]d^3\xi $$
(35)
is the dimensionless Hamiltonian, Eq.(24). The negative constant $`s_0`$ is a universal number corresponding to the largest action $`S_0`$ and is independent of the values of physical constants and the difference $`T_{\mathrm{}}T_c`$.
From Eq.(27) and from the expression for the Hamiltonian, Eq. (35), it is easy to see that
$$s_0=\underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑\tau \mathrm{\Pi }^2\mathrm{\Pi }d^3r$$
for a smooth solution of Eqs.(27), (28). In our case, the jump of $`\mathrm{\Theta }`$ derivatives should be taken into account which modifies the action :
$$s_0=\underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑\tau \mathrm{\Pi }^2\mathrm{\Pi }d^3\xi \underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑\tau \mathrm{\Delta }\mathrm{\Theta }^{^{}}\mathrm{\Pi }𝑑S_r^{}.$$
The surface integral should be taken over a sphere of radius $`r^{}`$. Using Eq. (28),
$$^2\mathrm{\Pi }=_\tau \mathrm{\Pi }U(r^{}\xi )\mathrm{\Pi },$$
we finally find for the action
$$s_0=\underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑\tau \underset{\xi <r^{}}{}\mathrm{\Pi }^2d^3\xi \underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑t\mathrm{\Delta }\mathrm{\Theta }^{^{}}\mathrm{\Pi }𝑑S_r^{}$$
(36)
Substituting the results of the numerical calculations we get
$$s_0=100.73.$$
Test results for a smaller grid spacing, $`\delta \xi =0.003`$, give
$$s_0=100.23$$
with a negligible jump-function, $`\mathrm{\Delta }\mathrm{\Theta }^{^{}}`$. These results are in agreement with the naive estimate from Eq.(36): The characteristic scale of $`\mathrm{\Pi }^2`$ is 1 and the action is proportional to the volume of a sphere of radius $`\pi `$, which is about 100.
To estimate the temperature-variation rate, one can take
$$|v_{max}|=\frac{T_{\mathrm{}}T_c}{\tau _{ph}}(\sqrt{\chi \tau _{ph}})^3n.$$
(37)
In this case, in accordance with Eq. (21), one can write for the probability flux in the transition region
$$j\frac{T_{\mathrm{}}(n\left(\sqrt{\chi \tau _{ph}}\right)^3)^{\frac{1}{2}}}{\tau _{ph}c_p}e^S\nu .$$
(38)
The constant $`\nu `$ cannot be estimated from the theory of hydrodynamic fluctuations PiT . This quantity gives the number of small equilibrium fluctuations with $`\delta TT`$ in a unit volume on the atomic scale. As an estimate, we can use the relation $`\nu =n/T_{\mathrm{}}`$ . Thus, the number of critical fluctuations per unit time in a unit volume is
$$\frac{dN}{dt}\frac{n(n\left(\sqrt{\chi \tau _{ph}}\right)^3)^{1/2}}{\tau _{ph}c_p}\mathrm{exp}\left[s_0\frac{nc_p(T_{\mathrm{}}T_c)^2}{2T_{}^{2}{}_{\mathrm{}}{}^{}}(\sqrt{\chi \tau _{ph}})^3\right].$$
(39)
The nuclei of a new phase are intensively formed as $`T_{\mathrm{}}`$ approaches $`T_c`$ and then grow rapidly. We have considered the initial phase of critical-fluctuation growth and restricted our analysis to the heat transfer via heat conduction, disregarding superfluidity effects at this stage. This approximation can be justified by the fact that the largest contribution to the action comes from the region lying far from the region of low velocities $`v`$, where $`p(𝐫)=\delta S/\delta T(𝐫)`$ becomes small and the fluctuation contribution can be neglected because $`Wv^1e^{S_0}`$.
## VI Bose condensation in the two-dimensional case
Most experiments on Bose-condensation in exciton systems were done in the two-dimensional (2D) case because these excitons are more stable But ; Tim . Condensation into the superfluid state of interwell exctions in AlAs/GaAs structures was supposedly observed in But . In these experiments excitons were obtained in a 2D quantum well and were cooled by phonon emission into the 3D volume of the surrounding semiconductor SVI . We can expect our theory to be suitable to explain the formation of a condensate in this case. In such systems the Bose-gas of excitons is dilute ($`na_{0}^{}{}_{}{}^{2}1`$) and has a long lifetime. There is no a true condensate at any nonzero temperature but it was shown in Pop that there is a superfluid transition at
$$T_c=\frac{2\pi n\mathrm{}^2}{m\mathrm{ln}\mathrm{ln}(1/na_{0}^{}{}_{}{}^{2})}.$$
(40)
According to Pop the pressure has the form:
$$P=\frac{2\pi n^2}{m\mathrm{ln}(1/na_{0}^{}{}_{}{}^{2})}+\frac{mT^2\zeta (2)}{2\pi \mathrm{}^2}.$$
(41)
As it was pointed out in fish these results hold only under the condition,
$$\mathrm{ln}\mathrm{ln}(1/na_{0}^{}{}_{}{}^{2})1.$$
(42)
Using the analogy to the 3D case, we consider temperature fluctuations at constant pressure and generalize the results of previous sections to the 2D case. In these fluctuations below $`T_c`$ there is also an increase of the reciprocal phonon time, and in the 2D case Eq.(6) should be replaced by
$$\frac{\delta n}{n}=\frac{\delta T}{T}\frac{\mathrm{ln}(1/(na_{0}^{}{}_{}{}^{2}))}{(\mathrm{ln}\mathrm{ln}(1/(na_{0}^{}{}_{}{}^{2})))^2}.$$
(43)
The large parameter $`(\mathrm{ln}(1/(na_{0}^{}{}_{}{}^{2}))/(\mathrm{ln}\mathrm{ln}(1/(na_{0}^{}{}_{}{}^{2})))^2`$ plays the role of $`1/\eta ^{\frac{1}{3}}`$ in Eq. (6). Eventually, we arrive to the same set of equations, Eqs. (27-28), as in the 3D case. Thus, the number of critical fluctuations per unit time in a unit volume in the 2D case is
$$\frac{dN}{dt}\frac{n(n\chi \tau _{ph})^{1/2}}{\tau _{ph}c_p}\mathrm{exp}\left[s_0\frac{nc_p(T_{\mathrm{}}T_c)^2}{2T_{}^{2}{}_{\mathrm{}}{}^{}}(\chi \tau _{ph})\right].$$
(44)
Numerical calculations in the 2D case are somewhat more complicated in comparison to the 3D case. There is no a stationary solution like Eq.(IV). However, if we consider the system of a finite size $`R`$, there will be an analog to Eq.(IV):
$`\mathrm{\Theta }_{stat}={\displaystyle \frac{J_0(\xi )}{\lambda _{}^{1}{}_{0}{}^{}J_1(\lambda _{}^{1}{}_{0}{}^{})\mathrm{ln}(\frac{R}{\lambda _{}^{1}{}_{0}{}^{}})}},\xi <\lambda _{}^{1}{}_{0}{}^{},`$
$`\mathrm{\Theta }_{stat}={\displaystyle \frac{\mathrm{ln}(\frac{\xi }{\lambda ^1}_0)}{\mathrm{ln}(\frac{R}{\lambda ^1}_0)}},\xi >\lambda _{}^{1}{}_{0}{}^{}.`$ (45)
Here $`J_n`$ is a Bessel function and $`\lambda _{}^{1}{}_{0}{}^{}`$ is the first zero of $`J_0(x)`$. We see that the negative values of $`\mathrm{\Theta }`$ are logarithmically suppressed. This is not very important because the system after passing in the vicinity of the stationary point evolves further due to the divergence term (which was neglected in Eq.(24) and Eqs.(27),(28)). This term in Eq. (24) has the form
$$\mathrm{\Delta }H=\frac{U(T_cT)}{\tau _{ph}}d^2r.$$
The Hamilton equations, Eqs.(27),(28), should also be modified:
$$\frac{\mathrm{\Theta }}{\tau }=^2\mathrm{\Theta }+\mathrm{\Theta }U(\mathrm{\Theta })+2^2\mathrm{\Pi },$$
(46)
$`{\displaystyle \frac{\mathrm{\Pi }}{\tau }}=`$ $``$ $`^2\mathrm{\Pi }\mathrm{\Pi }U(\mathrm{\Theta })+`$ (47)
$`\delta (\mathrm{\Theta }){\displaystyle \frac{T_{\mathrm{}}^{}{}_{}{}^{2}}{(T_{\mathrm{}}T_c)^2nc_p\chi \tau _{ph}}}.`$
Here $`\delta (\mathrm{\Theta })`$ is a delta-function. As in the toy-model, this term starts to play an important role at large $`\tau `$ while at small and intermediate $`\tau `$ it is unimportant. There is no need to take this term into account in our calculations of the action in the 3D case because it is a higher-order quasiclassical correction for the 3D instanton. But in the 2D case, we can consider the finite-time evolution due to this term. The large-time cutoff can be estimated by assuming that at this moment the divergence term becomes of the same order of magnitude as the main terms in Eq.(47). We know the large-time asymptotics of $`\mathrm{\Pi }`$ in a finite-size system of radius $`R`$,
$$\mathrm{\Pi }\alpha \mathrm{\Psi }_\lambda (r)\mathrm{exp}(|\lambda |\tau ),$$
where $`\mathrm{\Psi }_\lambda (r)`$ is the eigenfunction with the negative eigenvalue $`\lambda `$ of the 2D Schroedinger equation with the potential $`U(\lambda _{}^{1}{}_{0}{}^{}\xi )`$. Using this asym
ptotics and comparing terms in the equation for the norm, $`\mathrm{\Pi }=\mathrm{\Pi }^2d^2r`$,
$$_\tau \frac{\mathrm{\Pi }}{2}=|\lambda |\mathrm{\Pi }+\mathrm{\Pi }(\lambda _{}^{1}{}_{0}{}^{})\frac{T_{\mathrm{}}^{}{}_{}{}^{2}}{(T_{\mathrm{}}T_c)^2nc_p\chi \tau _{ph}},$$
we can estimate the large-time cutoff as:
$$\mathrm{\Delta }\tau \frac{1}{|\lambda |}\mathrm{ln}\left(\alpha |\lambda |\frac{(T_{\mathrm{}}T_c)^2nc_p\chi \tau _{ph}}{T_{\mathrm{}}^{}{}_{}{}^{2}}\right).$$
The space position of the thermal front, that corresponds to this time, is
$$\mathrm{\Delta }r=\sqrt{\mathrm{\Delta }\tau }.$$
If we choose the point $`R`$, where $`\mathrm{\Theta }(R)=1`$, at the distance $`\mathrm{\Delta }r`$ from the point $`r^{}`$ (where $`\mathrm{\Theta }(r^{})=0`$ at the largest time), we can assume that the value of the total action will be almost independent of the exact position of $`R`$. We have performed a number of runs with different values of $`R`$ and have found no essential difference in the action. The resulting constant is
$$s_0=13.6$$
The same naive estimate as in 3d can be performed. The action now is proportional $`\pi \left(\lambda _{}^{1}{}_{0}{}^{}\right)^2`$ which is about 13.
## VII The later stage of instanton growth
Analysis of the subsequent growth of the instanton requires the solution of hydrodynamic equations for a superfluid liquid, because a superfluid core appears in the developing fluctuation. We will qualitatively consider the phenomena that arise in this case. Proceeding from the assumption that the value of $`\tau _{ph}`$ is high, we assume that the motion in this region is quasi-stationary and tuned due to the slow cooling by phonons. We will use the hydrodynamic equations for a superfluid liquid in the vicinity of the transition point in the form proposed by Khalatnikov Hal :
$$\frac{𝐯_s}{t}=\left(\frac{v_s^2}{2}+\mu +\mu _s\right),$$
$$\frac{\rho }{t}+\mathrm{div}(\rho _s𝐯_s+\rho _n𝐯_n)=0,$$
$$\frac{}{t}\left(\rho _sv_s^i+\rho _nv_n^i\right)+\frac{}{x^k}\left(\rho _nv_n^iv_n^k+\rho _sv_s^iv_s^k+P\delta ^{ik}\right)=0,$$
$$T\frac{(n\sigma )}{t}+T\mathrm{div}(n\sigma 𝐯_n)$$
$$=\frac{2\mathrm{\Lambda }m}{\mathrm{}}\left[\mu _s+\frac{(𝐯_n𝐯_s)^2}{2}\right]^2\rho _s\frac{\rho _sc_pT}{\tau _{ph}},$$
$$\frac{\rho _s}{t}+\mathrm{div}\rho _s𝐯_s=\frac{2\mathrm{\Lambda }m}{\mathrm{}}\left[\mu _s+\frac{(𝐯_n𝐯_s)^2}{2}\right]\rho _s.$$
Here, $`\sigma `$ is the entropy per particle, n is the number of particles per unit volume, and $`\rho `$ is the density. The subscripts $`n,s`$ correspond to the normal and superfluid components, respectively; the constant $`\mathrm{\Lambda }`$ is the relaxation parameter; and we introduced the term that accounts for the phonon-induced energy removal in the equation for entropy. Here, the specific chemical potential $`\mu _s`$ for the superfluid density should ensure the condensate equilibrium density that is obtained by equating to zero the relaxation right-hand side of the equation for $`\rho _s`$. In our model of a weakly nonideal Bose gas, we can define phenomenologically
$$\mu _s=\frac{\mathrm{}^2a_0}{m^2}\left[(nn(T))\right]+\frac{\mathrm{}^2a_0}{m^3}\rho _s,$$
(48)
so that $`\rho _s=m(nn(T))=m\delta n`$ in the equilibrium. Here, $`n(T)`$ is the number of particles outside the condensate. We assume that the quantity $`\mathrm{\Lambda }m/\mathrm{}`$ is large and $`T_cT`$ is large enough for the approximate equality $`\mu _s+v_s^2/20`$ to be satisfied (we disregard quantity $`v_n`$, which is small compared to $`v_s`$); this gives
$$\frac{\rho _s}{m}=\delta n\frac{v_s^2m^2}{2\mathrm{}^2a_0}=\delta n\left(1\frac{v_s^2m^2}{2\mathrm{}^2a_0\delta n}\right).$$
(49)
In this case, it follows from the hydrodynamic equations that $`\mu \mu (P,T)=const`$,
$$T\frac{(n\sigma )}{t}+T\mathrm{div}(n\sigma 𝐯_n)=\frac{c_pT\rho _s}{m\tau _{ph}},$$
(50)
and the momentum conservation law gives
$$\frac{}{t}\left(\rho _sv_s^i+\rho _nv_n^i\right)+\frac{}{x^k}\left(\rho _nv_n^iv_n^k+\rho _sv_s^iv_s^k+P\delta ^{ik}\right)=0$$
In view of the smallness of $`v_s`$ compared to the sound velocity and the smallness of $`\rho _s`$, we will neglect these corrections to pressure $`PP_0`$. In this case, only the equation for entropy is significant. Assuming that the derivative $`(n\sigma )/t`$ is small, according to the assumption that the process is quasi-stationary (low temperature-variation rate) , we find that the stationary regime $`\mathrm{div}(\sigma n𝐯_𝐧)=c_p\rho _s/\tau _{ph}`$ should approximately take place and that the mass flux should be zero ($`\rho _sv_s+\rho _nv_n=0`$). Considering that $`\rho _n\rho `$ we obtain the equation
$$\sigma \mathrm{div}(\rho _sv_s)=\frac{c_p\rho _s}{\tau _{ph}}$$
(51)
where $`\sigma `$ is the entropy per particle. This equation determines the heat transfer in the fluctuation superfluid core. Using Eq.(49) we obtain
$$\sigma \frac{1}{r^2}\frac{}{r}r^2(1\frac{v_s^2}{u^2})v_s(1\frac{v_s^2}{u^2})\frac{c_p}{\tau _{ph}}=0,$$
$$u^2=\frac{2\mathrm{}^2a_0\delta n}{m^2}.$$
(52)
By introducing the dimensional distance $`\xi =c_pr/\sigma u\tau _{ph}`$ and $`v=v_s/u`$ we arrive at the equation
$$\frac{v}{\xi }=\frac{(1v^2)(1\frac{2}{\xi }v)}{13v^2}$$
(53)
The singular points of this differential equation are
$$\xi =0,v=0\text{and}\xi =\frac{2}{\sqrt{3}},v=\frac{1}{\sqrt{3}},$$
the latter point being a focus with the eigenvalues $`\lambda =1\pm i\sqrt{5}`$ Since the velocity $`v`$ must vanish at $`\xi =0`$, $`v\xi /3`$ for small $`\xi `$ and increases faster than by the linear law, with the derivative $`\frac{d\xi }{dv}`$ vanishing at $`v=\frac{1}{\sqrt{3}}`$ and at a certain $`\xi =\xi _c`$, whereupon the derivatives assume negative values upon the further increase in $`v`$. Thus, a regular superfluid flow cannot be continued after the point $`\xi _c`$ (the constant is on the order of unity and can be determined numerically). The critical radius
$$r_c\frac{\sigma u\tau _{ph}}{c_p}=\frac{\sigma }{c_p}\sqrt{2\frac{\delta n}{n}\eta ^{1/3}\frac{\tau _{ph}}{\tau _{tr}}}\sqrt{\chi \tau _{ph}}$$
(54)
can be smaller than $`\sqrt{\chi \tau _{ph}}`$ ; it should also be noted that $`1v^2>0`$(i.e., a singularity appears in the superfluid core). This singularity indicates that the quasi-stationarity conditions are violated at $`\xi \xi _c`$, and a complex nonstationary superfluid flow with the intense vortex formation in an instanton should appear upon the transition to the normal liquid at $`T>T_c`$. Similar effects are observed in a superfluid liquid in the gravitational field, where $`T_c`$ is a function of one (vertical) coordinate and a fixed heat flux from the superfluid to the normal liquid takes place Ak . We are dealing with a similar situation arising due to the nonuniform cooling as the critical temperature in the superfluid nucleus is approached. The results of numerical calculation Wen and experimental data Fen ,Bad indicate the formation of a ”vortex” superfluid phase with a higher but finite thermal conductivity without a superfluid transport. The mechanism of vortex formation and the vortex phase of this kind have been poorly studied both theoretically and experimentally.
## VIII conclusion.
Thus, we have shown that, in contrast to Zur , a transition to the superfluid phase can occur through an independent growth of critical fluctuations (instantons) at temperatures above the critical point $`(T>T_c)`$ immediately in the course of external cooling. These fluctuations subsequently transform into macroscopic formations. The growth of the nucleus of the superfluid state is accompanied by vortex generation in its external part. Consequently, vortex defects appear both due to the independent nucleation with an arbitrary phase upon cooling (the Zeldovich-Kibble hypothesis) and directly during the growth of each superfluid nucleus. This vortex-generation mechanism during the growth of an instanton significantly differs from the mechanism determined in Ar , where the existence of a superfluid flow interacting with the heated normal regions was presumed. In Ar , an attempt was made to explain the results of experiments Ru , in which $`{}_{}{}^{3}\mathrm{He}`$ was irradiated by neutrons. As a result, some regions heated to temperatures above $`T_c`$ appeared. These regions were cooled by the surrounding superfluid $`{}_{}{}^{3}\mathrm{He}`$, and the formation of vortices was detected. Thus, nonuniform cooling took place that differs considerably from the model used in our study. In the critical fluctuation considered here, heating takes place due to its nonsuperfluid surroundings. Consequently, it is advantageous for the fluctuation to preserve its spherical symmetry to reduce this heating. In the case of cooling of a heated region with superfluid surroundings Ru , the interface must obviously be unstable against its shape distortions, because this leads to a faster cooling. However, the stability, as well as the phase-transition mechanism itself, under such conditions (which, in contrast to Ar , are not associated with the existence of an external superfluid flow) calls for detailed investigations.
Probably, analogous schemes can be developed for the kinetics of various other phase transitions in the presence of the external cooling with the scenario essentially given by the toy model described in the text.
## IX ACKNOWLEDGMENTS
We are grateful to V.V. Lebedev, I.V. Kolokolov and H. Müller-Krumbhaar for discussions. This study was supported by the President of the Russian Federation (grant no. NSh-1715.2003.3) in Support of Young Russian Scientists and Leading Scientific Schools, the program Quantum Macrophysics of the Presidium of the Russian Academy of Sciences, and RFFI grants no. 03-02-16012, 05-02-16553. |
warning/0506/hep-th0506025.html | ar5iv | text | # Scale Invariant Hairy Black Holes
## 1 Introduction and discussion
The system of Gravity coupled to scalar fields has recently been under considerable scrutiny. Asymptotically AdS Hairy black holes have been shown to exist in . The issue of defining meaningful conserved charges has been considered in . Earlier references include . Within the AdS/CFT correspondence, the coupling of scalar matter was considered in .
Our aim in this paper is to make some general remarks on the structure of hairy black holes in three dimensions. Our key ingredient is the existence of a scale symmetry in the reduced action governing the black hole ansatz. This symmetry exists for any potential $`V(\varphi )`$ and provides, via Noether’s theorem, a radially conserved charge. We use this charge to find a relationship between the black hole parameters at infinity with those at the horizon.
Our main result is the following. Let $`M`$, $`J`$ and $`S`$ be the total energy, angular momentum, and entropy of a black hole solution with some non-zero scalar field $`\varphi `$. Let $`T`$ and $`\mathrm{\Omega }`$ be the black hole’s temperature and angular velocity. Assuming that the matter field is finite at the horizon and vanishes at infinity, it follows that these parameters must satisfy the three-dimensional Smarr relation,
$$M=\frac{1}{2}TS\frac{1}{2}\mathrm{\Phi }Q\mathrm{\Omega }J.$$
(1)
The remarkable aspect of this result is its universality. In fact the scalar field and its potential play no role. The only condition on the matter field is that it must be finite everywhere, and zero asymptotically. Of course this imposes non trivial constraints on the class of potentials being considered, which must elude the no-hair theorems. But, if the black hole exists, then it must satisfy (1).
The first law of black hole thermodynamics,
$$\delta M=T\delta S\mathrm{\Phi }\delta Q\mathrm{\Omega }\delta J,$$
(2)
is also valid in this theory. Inverting the Smarr relation (1) one finds that $`S(M,J,Q)`$ must be a homogeneous<sup>1</sup><sup>1</sup>1I terms of the total mass, the homogeneity property reads $`M(\sigma S,\sigma ^2J,\sigma ^2Q)=\sigma ^2M(S,J,Q)`$ function of degree 1/2 of its arguments, $`S(\sigma ^2M,\sigma ^2J,\sigma Q)=\sigma S(M,J,Q)`$. This is certainly true for the vacuum BTZ back hole. Our result implies that hairy black holes, regardless of the potential chosen, satisfy the same scaling relation.
A remark is in order here: the homogeneity property of $`S(M,J,Q)`$ is not a consequence of simple dimensional analysis and scaling arguments as is the case e.g. for the Kerr-Newman metrics, c.f. . This is due to the presence of an additional dimensionful parameter, the curvature radius of the AdS space-time or, equivalently, the cosmological constant<sup>2</sup><sup>2</sup>2Generalized Smarr relations have been considered in . . The reason why (1) holds nevertheless, even in the presence of scalar hair, is the scaling symmetry and the associated radially conserved charge.
In the non-rotating neutral case, $`J=Q=0`$, we can use (1) and (2) to find the general expression for the temperature of non-rotating black holes,
$$T=\kappa M^{1/2}$$
(3)
where $`\kappa `$ is a constant with no variation. This means, in particular, that for any potential $`V(\varphi )`$ the specific heat of the black hole is positive.
It is interesting to compare (3) with the result reported in . In three dimensions, considered the potential,
$$V=\frac{1}{8}\left(\mathrm{cosh}^6\varphi +\nu \mathrm{sinh}^6\varphi \right),$$
(4)
where $`\nu `$ is a real parameter. An explicit black hole configuration was displayed, whose temperature as a function of the total energy follows the general law (3), and $`\kappa `$ becomes a complicated function of the parameter $`\nu `$.
The three dimensional structure can be generalized to higher-dimensional black holes with toroidal topology , as well as to black holes on flat branes . This is analyzed in Sec. 8.
After this work was completed we became aware of in which the Smarr relation for hairy black holes in three-dimensions was also found. The parametrization of the reduced action used in this reference is very different from ours, and the relevant symmetry is an $`SL(2,R)`$ group rather than the scaling symmetry which we have employed.
## 2 Reduced action in d=3 and scaling symmetry
Consider the action describing three-dimensional gravity coupled to a scalar field $`\varphi `$,
$$I=\frac{1}{16\pi G}\left(R8g^{\mu \nu }_\mu \varphi _\nu \varphi 16V(\varphi )\right)\sqrt{g}d^3x.$$
(5)
We assume that $`V(\varphi )`$ has a non-zero negative value at $`\varphi =0`$, such that the gravitational background is anti-de Sitter space.
We shall first consider non-rotating solutions. The generalization can be done straightforwardly and will be indicated in Sec. 6. Consider static, spherically symmetric solutions of the form
$$ds^2=\gamma (r)^2h(r)dt^2+\frac{dr^2}{h(r)}+r^2d\phi ^2,\varphi =\varphi (r).$$
(6)
Solutions of this form include, for example, black holes and soliton solutions. The solitons are relevant for AdS/CFT applications, as recently considered in . We shall concentrate in this paper on black holes. One can write a reduced action for this problem,
$$I[h,\gamma ,\varphi ]=\frac{(t_2t_1)}{8G}𝑑r\gamma \left(h^{}+8rh\varphi ^2+16rV(\varphi )\right)+B.$$
(7)
where $`B`$ is a boundary term that we shall consider below. The equations of motion are,
$`h^{}+8rh\varphi ^2+16rV`$ $`=`$ $`0,`$
$`\gamma ^{}+8r\gamma \varphi ^2`$ $`=`$ $`0,`$ (8)
$`(r\gamma h\varphi ^{})^{}+r\gamma V_{,\varphi }`$ $`=`$ $`0.`$
They can be shown to be consistent with the original Einstein equations.
The key observation is that the action (7) is invariant under the scale transformations<sup>3</sup><sup>3</sup>3In the matter-free case this scale symmetry was already observed in .
$`\stackrel{~}{r}`$ $`=`$ $`\sigma r,`$ (9)
$`\stackrel{~}{h}(\stackrel{~}{r})`$ $`=`$ $`\sigma ^2h(r),`$ (10)
$`\stackrel{~}{\gamma }(\stackrel{~}{r})`$ $`=`$ $`\sigma ^2\gamma (r),`$ (11)
$`\stackrel{~}{\varphi }(\stackrel{~}{r})`$ $`=`$ $`\varphi (r),`$ (12)
with $`\sigma `$ a (positive) constant.
By direct application of Noether’s theorem to the above symmetry one finds that the combination
$`C=\gamma \left(h+{\displaystyle \frac{1}{2}}rh^{}+8r^2h\varphi ^2\right)`$ (13)
is conserved, $`C^{}=0`$. One can in fact prove this directly from the equations (8). A crucial property of this conservation law is that it holds for any potential $`V(\varphi )`$. This will allow us to make general statements about the nature of 3d black holes coupled to scalar fields.
Our strategy is now the following. Since $`C`$ does not depend on $`r`$ we can use it to find a relationship between the asymptotic parameters $`M,\beta `$ and the horizon $`r_+`$. As we shall see, this relation is precisely the Smarr relation (1). But before we can state this result, we need to find an expression for the energy of this system.
## 3 Energy and entropy
The analysis in this section assumes a generic potential. For some specific cases, as masses saturating the BF bound, a separate analysis may be needed.
The boundary term $`B`$ that appears in (7) is fixed by the condition that, upon varying the action, all boundary terms cancel for a set of given boundary conditions. At this point we shall switch to the Euclidean formalism, and interpret the on-shell action as the free energy of the thermodynamical system . The Euclidean action $`I_\mathrm{E}`$ is the same as (7), except that $`(t_2t_1)=1`$ and an overall sign, such that the weight in the functional integral is $`e^{I_\mathrm{E}}`$.
In the Euclidean Hamiltonian formalism, the boundary consist of two disconnected pieces, one in the asymptotic region $`r\mathrm{}`$ and the other at the horizon. The boundary term $`B`$ is specified by the condition,
$`\delta B`$ $`=`$ $`{\displaystyle \frac{1}{8G}}\gamma (\delta h+16rh\varphi ^{}\delta \varphi )|_{r=\mathrm{}}+{\displaystyle \frac{1}{8G}}\gamma \delta h|_{r=r_+},`$ (14)
where the horizon is defined by the equation $`h(r_+)=0`$. We assume that all fields are regular there.
Assuming that the matter field vanishes at infinity, Eqns. (8) imply that, asymptotically, $`\gamma ^{}=0`$. We write $`\gamma (\mathrm{})=\beta `$, where $`\beta `$ is a constant equal to the Euclidean period at infinity<sup>4</sup><sup>4</sup>4Note that solutions of the form $`\gamma \mathrm{log}(r)`$ will not occur for a generic potential..
The boundary term now has the form $`\delta B=\beta \delta M\delta S`$ where the variation of mass and entropy are given by,
$`\delta M`$ $`=`$ $`{\displaystyle \frac{1}{8G}}(\delta h+16rh\varphi ^{}\delta \varphi )|_{r=\mathrm{}},`$ (15)
$`\delta S`$ $`=`$ $`{\displaystyle \frac{1}{8G}}\gamma \delta h|_{r=r_+},`$ (16)
As usual in the Hamiltonian formalism the entropy comes from the variation of the action at the horizon. Our task now is to identify the actual values of $`S`$ and $`M`$.
The boundary term at the horizon gives the usual Bekenstein-Hawking entropy without any modifications. In fact, from $`h(r_+)=0`$ and $`(h+\delta h)(r_++\delta r_+)=0`$ it follows that $`\delta h(r_+)=h^{}(r_+)\delta r_+`$, as long as $`h^{}(r_+)0`$, which is satisfied for non-extreme black holes. In addition, the value of $`\gamma `$ at the horizon cannot be arbitrary. To avoid conical singularities at $`r=r_+`$ one must impose
$$\gamma (r_+)h^{}(r_+)=4\pi .$$
(17)
These two conditions allow us to identify $`S`$ as
$$S=\frac{2\pi r_+}{4G},$$
(18)
just as in the matter-free system.
We now turn to the problem of integrating (15) to extract the value of $`M`$. This problem is more subtle because we have not specified the potential. We shall integrate (15) by using again the scale invariance discussed above, which maps solutions to solutions.
The idea is the following. The functions $`h`$ and $`\varphi `$ have scaling dimensions 2 and 0, respectively. From (15) we conclude that $`M`$ must have scaling dimension 2. This means that under the scale variations of $`h`$ and $`\varphi `$,
$`\delta h`$ $`=`$ $`\delta \sigma (rh^{}+2h),`$ (19)
$`\delta \varphi `$ $`=`$ $`r\delta \sigma \varphi ^{},`$ (20)
the corresponding variation of $`M`$ satisfies
$$\delta M=2\delta \sigma M.$$
(21)
We now replace (19) and (20) in (15) and, comparing with (21), we obtain the desired formula for $`M`$<sup>5</sup><sup>5</sup>5The relationship between asymptotic functional variations and scale transformations can be checked explicitly in some examples. For the BTZ black hole with $`h(r)=r^2+h_0`$ one has $`\delta h=\delta h_0`$. The constant $`h_0`$ has scaling dimension 2, $`\delta h_0=2\delta \sigma h_0`$. One can check that in fact $`\delta \sigma (rh^{}+2h)=\delta h_0`$. In the system studied in , the asymptotic solution is $`hr^2+4Br3(1+\nu )B^2`$. It is direct to check that $`\delta h=\delta \sigma (rh^{}+2h)`$ with $`\delta B=\delta \sigma B`$, as claimed. Note finally that this correspondence fails in higher dimensional gravity. See Sec. 8 for details on this case.,
$$M=\frac{1}{8G}\left(h+\frac{1}{2}rh^{}+8r^2h\varphi ^2\right).$$
(22)
Before explaining and discussing the validity of this formula let us check that it gives the right results in known cases. For a BTZ black hole, $`h=r^28Gm`$ and $`\varphi =0`$. One finds $`M=m`$, as expected. A less trivial example is the exact hairy black hole solution found in with $`h=r^2+4Br3(1+\nu )B^2+𝒪(1/r)`$ and $`\varphi =(B/r)^{1/2}2/3(B/r)^{3/2}+𝒪(1/r^{5/2})`$. Replacing this field in (22) one obtains
$$M=\frac{3(1+\nu )B^2}{8G},$$
(23)
in full agreement with .
Now, some comments on the derivation and validity of (22) are necessary. The variations (19,20) do not explore the full set of asymptotic solutions. In fact, (19) and (20) represent a 1-parameter ($`\sigma `$) set of variations. On the other hand, the equations are of first order for $`h(r)`$, and second order for $`\varphi (r)`$ and the full space of solutions has three parameters. The key step is that since $`M`$ is a “function of state” (exact differential), its value does not depend on the path chosen and in this sense the formula (22) is the correct one. However, we must now make sure that $`\delta M`$, as given in (15), is actually an exact differential. An equivalent way of stating this is that the existence of a well-defined variational principle requires $`B`$ in (7) to exist, not just $`\delta B`$.
We do not need to worry about the first term in (15), $`\delta h`$, which is exact. The second piece, $`rh\varphi ^{}\delta \varphi `$, needs a separate analysis. For a generic potential the asymptotic form of the scalar field on AdS is,
$$\varphi =\frac{a}{r^{\mathrm{}_{}}}+\frac{b}{r^\mathrm{}_+}+\mathrm{},$$
(24)
where, for static black holes, $`a`$ and $`b`$ are arbitrary constants and represent the two degrees of freedom associated to $`\varphi `$. The exponents $`\mathrm{}_\pm `$ are the solutions to a quadratic equation and satisfy $`\mathrm{}_++\mathrm{}_{}=2.`$ We assume that both are positive.
Plugging (24) into (15) one finds finite terms of the form $`f(a,b)\delta a+g(a,b)\delta b`$. In order to write these terms as total variations (to achieve path independence) one needs to assume a relationship between $`a`$ and $`b`$. This restriction on the space of solutions is generic and was also found in .
The particular choice considered in these references (generalized here to arbitrary $`\mathrm{}_\pm )`$ is
$$b=\eta a^\frac{\mathrm{}_+}{\mathrm{}_{}},\delta \eta =0,$$
(25)
where $`\eta `$ is held fixed. This choice is consistent with the full anti-de Sitter asymptotic group, although this will not be relevant for our discussion<sup>6</sup><sup>6</sup>6Note that this particular choice is by no means the most general. For solitonic solutions, as in , $`a`$ and $`b`$ become related in a different way. We shall consider solitons in this theory elsewhere..
For our purposes, the choice (25) is singled out by demanding scale invariance of the asymptotic solution. In fact, once a relationship between $`a`$ and $`b`$ is assumed, the only function $`b=b(a)`$ consistent with (20) is precisely (25). We conclude that on the space of solutions satisfying the boundary conditions (25), the formula (22) for $`M`$ is correct.
Finally, we point out that the remarkable cancelations of divergent pieces in the total mass $`M`$, discovered in , can be seen in this case from a different perspective. Note that, up to the factor $`\gamma (r)`$ which becomes a constant at infinity, $`M`$ is exactly equal to the scale charge $`C`$ displayed in (13). Since $`C`$ does not depend on $`r`$, it cannot diverge; the total mass is then finite.
## 4 Thermodynamics of the hairy black hole
### 4.1 The First Law
The first law for our class of black hole solutions can be checked by standard Hamiltonian arguments. The form of the action, derived in the previous section, after all boundary terms have been included is,
$$I[\beta ]=𝑑r\gamma +\beta MS(r_+),$$
(26)
where $`M`$ is given in (22) and $`S`$ is the usual entropy in three dimensions, given in (18). $`=0`$ is one of the equations of motion. By construction, this action has an extremum when evaluated on solutions with $`\beta `$ fixed. The on-shell value of $`I`$ only depends on $`\beta `$. The value of $`M`$ is such that $`I`$ has an extremum.
Since the bulk contribution is proportional to a constraint, the on-shell value of the action is
$$I[\beta ]=\beta MS(M)$$
(27)
where $`r_+`$ is written as a function of the total energy $`M`$ using the solution<sup>7</sup><sup>7</sup>7There is a non-trivial assumption here, namely, that $`r_+`$ depends only on $`M`$, and not on the scalar field parameter $`a`$. We prove in the next section that black holes exists only for special values of $`a=a(M)`$, and hence $`a`$ is not an independent parameter.. $`M`$ is not fixed but has to be chosen such that $`I`$ has an extremum, that is, the first law is satisfied,
$$\beta \delta M=\delta S.$$
(28)
### 4.2 The Smarr relation
We are now ready to prove our main result. We go back to the expression for the scaling charge $`C`$ given in (13). Comparing (13) with (22) we conclude that the scaling charge is proportional to the total mass. Evaluating $`C`$ at infinity we get the exact relation,
$$C=8\beta MG.$$
(29)
On the other hand, since $`C`$ is $`r`$independent we can also evaluate it at the horizon $`h(r_+)=0`$ to get
$$C=2\pi r_+.$$
(30)
Here we have used the condition of absence of conical singularities (17). Comparing the values of $`C`$ at infinity and at the horizon we find the equation,
$$\beta M=\frac{1}{2}\frac{2\pi r_+}{4G}$$
(31)
representing the non-rotating version of (1). This relation is satisfied for any black hole solution with or without scalar field. Of course, this is also true for the BTZ vacuum black hole, as can be readily checked. The rotating case, leading to (1), will be indicated in Sec. 6.
### 4.3 Positivity of energy
We prove now that a hairy black hole can exist only if the total mass $`M`$ is positive. To see this we first note that the field $`\gamma (r)`$ does not change sign in the whole range $`r_+r\mathrm{}`$. In fact, directly from the equations of motion (8) we can write the formal solution
$$\gamma (r)=\gamma _0e^{^r𝑑s\mathrm{\hspace{0.17em}8}s\varphi ^{}(s)^2}$$
(32)
where $`\gamma _0`$ is an arbitrary integration constant. This expression for $`\gamma `$ is manifestly positive, if $`\gamma _0`$ is positive. Now, the scaling charge evaluated at the horizon and at infinity gives the equation (we relax here the condition (17) and consider either Minkowskian or Euclidean signature)
$$16G\gamma _{\mathrm{}}M=\gamma _+h_+^{}r_+$$
(33)
where the subscript <sub>+</sub> indicates the corresponding function evaluated at $`r_+.`$ The function $`h(r)`$ must be positive outside the horizon, and vanishes at $`r_+`$. This means that $`h_+^{}>0`$. Since $`\gamma (r)`$ does not change sign and $`r_+`$ is positive, we conclude that this equation requires
$$M>0.$$
(34)
### 4.4 Temperature and specific heat
Combining (31) with (18) and the first law we derive the general relation
$$8MG=\kappa _0^2r_+^2,$$
(35)
where $`\kappa _0`$ is an arbitrary (dimensionless) integration constant, with no variation. This constant cannot be computed from this analysis and depends on the details of the potential, as well as all other fixed parameters. For example, for the BTZ black hole $`\kappa _0=1`$ while for the potential (4) considered in one finds that $`\kappa _0`$ becomes a complicated function of $`\nu `$.
All thermodynamical properties can now be extracted, for example the temperature as a function of the mass gives,
$$\frac{1}{\beta }=T=\kappa _0\frac{\sqrt{2MG}}{\pi },$$
(36)
as announced. We can also check that the specific heat,
$$c=\frac{M}{T}=\frac{2\pi r_+}{4G},$$
(37)
becomes equal to the entropy (this was also noted by in their particular example).
## 5 A closer look at the hairy black hole
A hairy black hole is, by definition, a solution to the Einstein + matter system displaying a regular horizon. In particular, the value of the matter field $`\varphi (r)`$ at the horizon must be finite. We have argued in the previous section that, if the black hole exist, then the Smarr (1) relation is satisfied. However, we have said very little about the conditions for the existence of a black hole.
The condition of regularity at the horizon imposes constraints on the solutions which can be analyzed using scale invariance. In this section we will prove that for a given value of $`\eta `$ (see (25)), the values of the total mass $`M`$ and the parameter $`a`$ have to be fine tuned in order to have a regular black hole. This means that, apart from $`\eta `$ which acts as an external parameter with no variation, the only free parameter in the black hole spectrum is the total energy $`M`$.
Consider the set of equations of motion (8). We would like to find a solution displaying a regular event horizon $`h(r_+)=0`$. At the point $`r=r_+`$, the matter field, and its derivatives up to some sufficiently high order, must be finite. In particular,
$$\varphi (r_+)=\varphi _0.$$
(38)
Define the new field
$$\chi (r)=\varphi (r)\varphi _0$$
(39)
which vanishes at the horizon, $`\chi (r_+)=0`$. Near the horizon the new field $`\chi `$ is small and hence, without specifying $`V(\varphi )`$, we write the near horizon series,
$$V(\chi )=v_0+v_1\chi +v_2\chi ^2+\mathrm{}$$
(40)
where the constants $`v_i`$ depend on the potential $`V`$ and $`\varphi _0`$.
Under these conditions, the fields $`h,\gamma ,\varphi `$ have the following series expansions near the horizon
$`h`$ $`=`$ $`h_1(rr_+)+h_2(rr_+)^2+\mathrm{}`$ (41)
$`\gamma `$ $`=`$ $`\gamma _0+\gamma _1(rr_+)+\gamma _2(rr_+)^2+\mathrm{}`$ (42)
$`\chi `$ $`=`$ $`\chi _1(rr_+)+\chi _2(rr_+)^2+\mathrm{}`$ (43)
Recall that in the Euclidean formalism the values of $`h^{}`$ and $`\gamma `$ at $`r=r_+`$ are linked by (17), that is $`h_1\gamma _0=4\pi `$. Our conclusions, however, do not depend on the signature.
We have assumed that no fractional powers or logs are present because they would induce divergences in the derivatives of the fields.
We now plug this series expansion into the equations of motion and solve for the coefficients order by order. This is a straightforward exercise that we do not display here. The important comment is that all coefficients are fixed in terms of $`\varphi _0`$ and $`r_+`$ (recall that $`\varphi _0`$ enters in the coefficient $`v_0`$ in the series (40) for the potential). There are thus 2 arbitrary constants at the horizon:
$$\text{Horizon data: }\{\varphi _0,r_+\},$$
(44)
as opposed to the series analysis at infinity with,
$$\text{Asymptotic data: }\{\eta ,a,M\}.$$
(45)
What happens here is that the series expansion (42) is not the most general one. There exists other solutions with logs or fractional powers (probably depending on the potential), which are not contained in the regular ansatz. <sup>8</sup><sup>8</sup>8This has also been remarked in .
We conclude that if one integrates from infinity to the horizon, the values of $`a,\eta `$ and $`M`$ must be fine tuned in order to reach a regular event horizon. Conversely, if one integrates from the horizon, prescribing the values of $`r_+`$ and $`\varphi _0`$, one gets at infinity a surface in the $`\eta ,a,M`$ space. Actually, we can say something else. We shall prove now that $`\eta `$ only depends on the value of $`\varphi _0`$, and not on $`r_+`$,
$`\eta =\eta (\varphi _0).`$ (46)
To see this, suppose we are given a solution to the equations of motion, $`h(r),\gamma (r),\varphi (r)`$ displaying a regular event horizon. Using scale invariance we can provide immediately another exact solution to the equations by the simple transformation
$`\stackrel{~}{h}(r)`$ $`=`$ $`\sigma ^2h(r/\sigma )`$ (47)
$`\stackrel{~}{\gamma }(r)`$ $`=`$ $`\sigma ^2\gamma (r/\sigma )`$ (48)
$`\stackrel{~}{\varphi }(r)`$ $`=`$ $`\varphi (r/\sigma )`$ (49)
The new solution is a different one! If the horizon in the first solution was at $`r=r_+`$, then the location of the horizon in the second solution is at
$$\stackrel{~}{r}_+=\sigma r_+.$$
(50)
In fact, $`\stackrel{~}{h}(\stackrel{~}{r}_+)=0.`$ This means that acting with scale transformations, we can cover all possible values of $`r_+`$. On the other hand, the value of $`\varphi _0`$ remains unchanged since
$$\stackrel{~}{\varphi }(\stackrel{~}{r}_+)=\varphi (r_+)=\varphi _0.$$
(51)
Acting with scale transformations we thus cover all solutions with a given value of $`\varphi _0`$. Now, scale transformations act on the asymptotic parameters leaving $`\eta `$ invariant. We thus conclude that the asymptotic parameter $`\eta `$ is in one-to-one correspondence with the value of $`\varphi `$ at the horizon
$$\varphi _0\eta .$$
(52)
For a given value of $`\eta `$, the value of $`\varphi _0`$ is determined. In the example of , $`\eta =2/3`$ and $`\varphi _0=\mathrm{tanh}^1(1/\sqrt{3}).`$
Recall that $`\eta `$ is fixed in the action principle, and acts as an external parameter. For fixed $`\eta `$ (and hence $`\varphi _0`$), the remaining degrees of freedom are $`M`$ and $`a`$, at infinity, and $`r_+`$ at the horizon. This means that if one integrates from the horizon, varying the values of $`r_+`$, one obtains at infinity a curve in the $`M,a`$ plane. As we have shown this curve will cover only the $`M>0`$ half plane. Of course, for different values of $`\eta `$, the curve changes.
## 6 Adding angular momentum
We will now extend the discussion of the thermodynamics to black holes with angular momentum. This requires a change of the ansatz for the metric (6) to
$$ds_\mathrm{E}^2=\gamma (r)^2h(r)dt^2+\frac{dr^2}{h(r)}+r^2(d\phi +n(r)dt)^2,\varphi =\varphi (r).$$
(53)
The reduced action is
$$I[h,\gamma ,n,\varphi ]=\frac{1}{8G}𝑑r\left\{\gamma \left(\frac{2p^2}{r^3}+h^{}+8rh\varphi ^2+16rV\right)+2np^{}\right\}+B.$$
(54)
Here $`p=\pi _\phi ^r=\frac{r^3}{2\gamma }n^{}`$. The bulk term of the action vanishes on-shell. The equation of motion for $`n`$ gives $`p=\mathrm{const}.`$. By shifting the angular coordinate we arrange for $`n(r_+)=0`$.
The action is invariant under (9-12), augmented by
$`\stackrel{~}{p}(\stackrel{~}{r})`$ $`=`$ $`\sigma ^2p(r),`$ (55)
$`\stackrel{~}{n}(\stackrel{~}{r})`$ $`=`$ $`\sigma ^2n(r).`$ (56)
This leads to the following radially conserved Noether charge
$$C=\gamma \left(h+\frac{1}{2}h^{}r+8r^2h\varphi ^2\right)2np.$$
(57)
One checks that indeed $`C^{}=0`$ by virtue of the equations of motion.
The boundary terms $`B`$ must again be chosen such that $`\delta B`$ cancels the boundary terms which appear when one extremizes the action. One finds
$`\delta B`$ $`=`$ $`{\displaystyle \frac{1}{8G}}\{(\beta (\delta h+16rh\varphi ^{}\delta \varphi )+2n\delta p\}|_{r=\mathrm{}}+{\displaystyle \frac{1}{8G}}(\gamma \delta h+2n\delta p)|_{r=r_+}`$ (58)
$``$ $`\beta (\delta M+\mathrm{\Omega }\delta J)\delta S`$
Here we have used that $`h(r_+)=0`$ and the definitions $`\gamma (\mathrm{})\beta ,n(\mathrm{})\beta \mathrm{\Omega }`$. It follows from the equations of motion that $`\beta `$ and $`\mathrm{\Omega }`$ are finite. The first two terms are the contribution from $`r=\mathrm{}`$, the last term is the contribution from the horizon. Replacing once more the functional variations by those which follow from the scaling properties of the fields, combined with the fact that $`M`$ and $`J`$ have weight two, one finds
$$C=8G\beta (M+\mathrm{\Omega }J)$$
(59)
From the contribution at $`r=r_+`$ we find again Eq.(18), i.e. $`S=\frac{2\pi r_+}{4G}`$.
In order to find a relation between $`M,J`$ and $`S`$, we use the fact that $`C`$ is a constant. While its expression at $`r=\mathrm{}`$ was used to relate it to $`M`$ and $`J`$, we now use its expression at the horizon to relate it to $`S`$. This leads to
$$\beta (M+\mathrm{\Omega }J)=\frac{1}{2}S.$$
(60)
as promised. One easily verifies that this relation is satisfied for the rotating BTZ black hole.
## 7 Adding electric charge
In this section we will generalize the previous discussion to rotating electrically charged black holes in the presence of a charged scalar field.
Starting point is the covariant action
$$I=\frac{1}{16\pi G}d^3x\sqrt{g}(R16g^{\mu \nu }D_\mu \varphi (D_\nu \varphi )^{}16V4F^{\mu \nu }F_{\mu \nu })$$
(61)
with $`D_\mu \varphi =_\mu +iqA_\mu \varphi `$ where $`q`$ is the electric charge of thes calar field. For the metric and the scalar field we make again the ansatz eq.(53), for the gauge field we also assume $`A_\mu =A_\mu (r)`$. This leads to the reduced euclidean action
$`I_\mathrm{E}^{\mathrm{red}}`$ $`=`$ $`{\displaystyle \frac{1}{8G}}{\displaystyle }dr\{{\displaystyle \frac{r^3n^2}{2\gamma }}+h^{}\gamma +16r\gamma V+{\displaystyle \frac{16}{r}}q^2\gamma |\varphi |^2A^2+16rh\gamma |\varphi ^{}+iqA_r\varphi |^2`$ (62)
$`+{\displaystyle \frac{8}{r}}h\gamma A^2+{\displaystyle \frac{16r}{\gamma h}}q^2|\varphi |^2(nA\mathrm{\Phi })^2+{\displaystyle \frac{8r}{\gamma }}(\mathrm{\Phi }^{}nA^{})^2\}+B`$
Here we have introduced the notation $`\mathrm{\Phi }=A_t,A=A_\phi `$. We have checked that the equations of motion derived from this reduced action imply those derive from the covariant action.
It is convenient to introduce the following momenta:
$`p`$ $`=`$ $`{\displaystyle \frac{1}{2\gamma }}r^3n^{}`$ (63)
$``$ $`=`$ $`{\displaystyle \frac{16r}{\gamma }}(\mathrm{\Phi }^{}nA^{})`$
where, as before, $`p=\pi _\phi ^r`$ whereas $``$ is proportional to the radial component of the electric field. In terms of these the reduced action reads
$`I_\mathrm{E}^{\mathrm{red}}`$ $`=`$ $`{\displaystyle }dr\{\gamma [{\displaystyle \frac{2p^2}{r^3}}+h^{}+16rV+16rh|\varphi ^{}+iqA_r\varphi |^2+{\displaystyle \frac{16}{r}}q^2A^2|\varphi |^2+{\displaystyle \frac{8h}{r}}A^2{\displaystyle \frac{1}{32r}}^2]`$ (65)
$`+{\displaystyle \frac{16r}{\gamma h}}q^2|\varphi |^2(nA\mathrm{\Phi })^2n(2p^{}A^{})+\mathrm{\Phi }^{}\}+B`$
We will not write down the equations of motion but want to mention that, even not completely obvious but nevertheless true, the reduced action vanishes on-shell.
The action is invariant under (9-12,55,56) and
$`\stackrel{~}{}(\stackrel{~}{r})`$ $`=`$ $`\sigma (r),`$ (66)
$`\stackrel{~}{A}(\stackrel{~}{r})`$ $`=`$ $`\sigma A,`$ (67)
$`\stackrel{~}{\mathrm{\Phi }}(\stackrel{~}{r})`$ $`=`$ $`\sigma ^1\mathrm{\Phi }(r),`$ (68)
$`\stackrel{~}{A}_r(\stackrel{~}{r})`$ $`=`$ $`\sigma ^1A_r(\sigma ).`$ (69)
The conserved Noether charge is
$`C`$ $`=`$ $`\gamma \left[h^{}+32rh|\varphi ^{}|^2+16iqrhA_r(\varphi \varphi _{}^{}{}_{}{}^{}\varphi ^{}\varphi ^{})+{\displaystyle \frac{16h}{r}}A^2\right]r`$
$`n(2p^{}A^{})r+\mathrm{\Phi }^{}r2\gamma h{\displaystyle \frac{16}{r}}h\gamma AA^{}+4np\mathrm{\Phi }nA`$
Note that we have not yet fixed the gauge. We could set $`A_r=0`$ and $`\mathrm{\Phi }(r_+)=0`$, however for the following discussion only the latter gauge condition, which ensures that the time component of the gauge field is regular at the horizon, will be relevant. We will, as before, also use $`n(r_+)=0=h(r_+)`$.
To derive the Smarr relation, we proceed as in the previous sections. Using the scaling relations one derives
$$\delta B=\frac{1}{8G}C\delta \sigma $$
(71)
On the other hand
$`\delta B|_{\mathrm{}}`$ $`=`$ $`\beta (\delta M+\mathrm{\Omega }\delta J+\mathrm{\Phi }\delta Q)`$
$`=`$ $`2\beta (M+\mathrm{\Omega }J+{\displaystyle \frac{1}{2}}\mathrm{\Phi }Q)\delta \sigma `$
from which we obtain
$$C=16G\beta (M+\mathrm{\Omega }J+\frac{1}{2}\mathrm{\Phi }Q).$$
(73)
The relative factor $`\frac{1}{2}`$ comes from the fact that $`Q=|_{r=\mathrm{}}`$ has scaling charge one while $`M`$ and $`J`$ have scaling charge two. We also have, as before
$$\delta B|_{r_+}=\gamma \delta h|_{r_+}$$
(74)
leading to
$$S=\frac{2\pi r_+}{4G}$$
(75)
Using that $`C(r_+)=r_+\gamma _+h_+^{}=4\pi r_+`$ we finally find the Smarr relation
$$\frac{1}{2}S=\beta (M+\mathrm{\Omega }J+\frac{1}{2}\mathrm{\Phi }Q).$$
(76)
## 8 D=4
In four dimensions the equations of motion have a similar structure although there are important differences. For reasons which will become clear very soon, we make the general ansatz for the metric
$$ds^2=\gamma ^2hdt^2+\frac{dr^2}{h}+r^2d\mathrm{\Omega }_k$$
(77)
where the “sphere” $`d\mathrm{\Omega }_k`$ is either a 2-sphere, a 2-torus or a higher genus surface,
$$d\mathrm{\Omega }_k=\{\begin{array}{cc}d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2,\hfill & k=1\hfill \\ dx^2+dy^2,\hfill & k=0\hfill \\ d\theta ^2+\mathrm{sinh}^2\theta d\varphi ^2,\hfill & k=1.\hfill \end{array}$$
(78)
Black holes with unusual topologies were first discussed in .
The ansatz (77) leads to the reduced action
$$I[h,\gamma ,\varphi ]=\frac{(t_2t_1)}{8\stackrel{~}{G}}𝑑r\gamma \left(rh^{}+hk+8r^2h\varphi ^2+16r^2V(\varphi )\right)+B.$$
(79)
We have introduced the notation $`\stackrel{~}{G}=\frac{4\pi G}{V_k}`$ with $`V_k=𝑑\mathrm{\Omega }_k`$. The horizon area is then $`V_kr_+^2`$. Varying $`\gamma ,h,\varphi `$ one obtains the equations of motion
$`rh^{}+hk+8r^2h\varphi ^2+16r^2V(\varphi )`$ $`=`$ $`0,`$ (80)
$`\gamma ^{}+8r\gamma \varphi ^2`$ $`=`$ $`0,`$ (81)
$`(r^2\gamma h\varphi ^{})^{}+r^2\gamma V_{,\varphi }`$ $`=`$ $`0.`$ (82)
These equations are similar to those in three dimensions, (8), except for the constant $`k`$ appearing in (80). This constant, which is a fixed number associated to the sphere’s curvature, spoils scale invariance<sup>9</sup><sup>9</sup>9Note that if one replaces $`\gamma =\lambda ^{}`$ and varies $`\lambda `$, the piece $`\lambda ^{}k`$ is a boundary term and the action becomes scale invariant. The space of solutions has an extra integration constant, and for particular values of that constant, the original equations are recovered. The main obstruction to follow up this idea is the relativistic version of the modified equations of motion..
However, for the torus topology, $`k=0`$, the equations are scale invariant and we can immediately generalize the discussion from $`d=3`$ to $`d=4`$<sup>10</sup><sup>10</sup>10The generalization to arbitrary $`d`$ is straightforward if we take for $`d\mathrm{\Omega }_{\mathrm{d}2}`$ the volume element of a flat torus., in particular due to scale invariance there is a radially conserved charge. Due to the invariance of the action under the replacements $`(r,h(r),\gamma (r),\varphi (r))(\stackrel{~}{r},\stackrel{~}{h}(\stackrel{~}{r}),\stackrel{~}{\gamma }(\stackrel{~}{r}),\stackrel{~}{\varphi }(\stackrel{~}{r}))`$ with (c.f. (9-12))
$`\stackrel{~}{r}`$ $`=`$ $`\sigma r`$ (83)
$`\stackrel{~}{h}(\stackrel{~}{r})`$ $`=`$ $`\sigma ^2h(r)`$ (84)
$`\stackrel{~}{\gamma }(\stackrel{~}{r})`$ $`=`$ $`\sigma ^3\gamma (r)`$ (85)
$`\stackrel{~}{\varphi }(\stackrel{~}{r})`$ $`=`$ $`\varphi (r)`$ (86)
one finds
$$C=\gamma \left(r^2h^{}2h+8r^3h\varphi ^2\right)$$
(87)
with $`C^{}=0`$ by virtue of the equations of motion.
Eq.(18) for the entropy is now
$$S=\frac{4\pi r_+^2}{4\stackrel{~}{G}}=\frac{V_0r_+^2}{4G}.$$
(88)
and
$$\delta M=\frac{V_0}{8\pi G}\left(r\delta h+8r^2h\varphi ^{}\delta \varphi \right)|_{r=\mathrm{}}$$
(89)
Using (86,84) and the fact that $`M`$ now has scaling weight three, one finds from (89) and from comparing with (87) the relation
$$C=\frac{24\pi G\beta }{V_0}M.$$
(90)
On the other hand, evaluation of $`C`$ at the horizon gives
$$C=4\pi r_+^2.$$
(91)
Comparison of (90) with (91) leads to the relation
$$\beta M=\frac{V_0r_+^2}{6G}.$$
(92)
In place of (35) one now finds
$$8\pi GM=V_0\kappa _0^3r_+^3$$
(93)
and the specific heat can be computed to be twice the entropy and for the temperature one finds
$$T=\frac{3}{2}\pi ^{2/3}\kappa _0^2\left(\frac{GM}{V_0}\right)^{1/3}.$$
(94)
We stress once more that these results are valid for arbitrary potentials as long as they lead to a solution for the scalar field which vanishes asymptotically. The specific form of the potential only enters through the integration constant $`\kappa _0`$.
The proof of positivity of $`M`$ proceeds in exactly the same way as in $`d=3`$. It depends crucially on the existence of the scaling charge, i.e. on considering the case $`k=0`$. In fact, negative mass hairy black holes for $`k=1`$ have recently been constructed in .
It is now straightforward to check that for a constant potential $`V=3`$, i.e. in the presence of a cosmological constant but no scalar field, one finds the above results with $`\kappa _0=1`$ as one easily verifies given the explicit solution
$$h=r^2\frac{8\pi Gm}{V_0r}$$
(95)
## 9 Acknowledgments
The authors gratefully acknowledge comments on the manuscript by S. Carlip and G. Horowitz, and G. Clement for pointing out an important misprint in the first version. We would also like to thank the Centro de Estudios Científicos in Valdivia, Chile and in particular J. Zanelli for hospitality during the initial stages of this work. M.B. was partially supported by FONDECYT (Chile) grants # 1020832 and # 7020832. M.B. also thanks the Albert Einstein Institute, Potsdam, and ULB, Brussels for their kind hospitality during the course of this work, and G. Barnich and P. Liendo for useful conversations. S.T. thanks the Physics Departments of P. Universidad Católica de Chile and of the Universidad Autónoma in Madrid for hospitality, and A. Schwimmer and M. Volkov for useful discussions. His work is partially supported by GIF – the German-Israeli Foundation for Scientific Reasearch and the EU-RTN network Constituents, Fundamental Forces and Symmetries of the Universe (MRTN-CT-2004-005104). |
warning/0506/astro-ph0506628.html | ar5iv | text | # X-ray and Near-infrared Studies of a Star-forming Cloud; L1448
## 1 INTRODUCTION
L1448 is a dark cloud (Lynds, 1962) located at the western edge of the Perseus molecular cloud complex at a distance of $``$300 pc. It has a size of $``$15′$`\times `$7′ in a C<sup>18</sup>O map and a mass of $``$100 $`M_{}`$ (Bachiller & Cernicharo, 1986). The region has been intensively studied in wavelengths longer than far-infrared focusing mainly on a giant molecular outflow, and only few studies in the near-infrared (NIR) broad-band (Bally et al., 1993) and none in the X-ray band have been presented, both of which are currently strong measures to identify stellar constituents of star-forming clouds. The infrared excess due to emission from circumstellar material and the X-ray brightness are two major observational characteristics of young stellar objects (YSOs), which can be respectively revealed through NIR broad-band photometry in the J, H, and K<sub>S</sub> bands and X-ray imaging observations. The purpose of this paper is to report the results of the first X-ray imaging-spectroscopy study of L1448 aided by follow-up NIR broad-band imaging observations. Previous X-ray studies in the Perseus cloud complex can be found in Preibisch et al. (1996); Preibisch & Zinnecker (2001, 2002, 2004) for IC 348, Getman et al. (2002); Preibisch (2003) for NGC 1333, and Yamauchi et al. (2001) for Barnard 1 and IRAS 03282$`+`$3035.
The most magnificent feature in L1448 is a parsec-scale bipolar outflow emanating from the core of the cloud, which is one of the highest-velocity, best-collimated, and youngest outflows known (Davis & Smith, 1995). The outflow was discovered by Bachiller et al. (1990) in a molecular emission line image. The exciting source of the outflow was promptly detected in the millimeter and centimeter continua (Bachiller et al., 1991; Curiel et al., 1990), which is referred as L1448 mm or L1448 C(enter). Barsony et al. (1998) obtained a spectral energy distribution (SED) in far-infrared to millimeter wavelengths and established the class 0 status of this source.
The Infrared Astronomical Satellite (IRAS) detected three sources (IRAS 03220$`+`$3035, 03222$`+`$3034, and 03225$`+`$3034) in this region, which are called L1448 IRS 1, 2, and 3. IRS 1 has a visual counterpart with $`m_V`$ $`=`$ 19.5 mag (Cohen & Kuhi, 1979). From its strong H$`\alpha `$ emission and excess emission in the IRAS band (Cohen & Kuhi, 1979; de Grijp et al., 1987), it is probably a classical T Tauri star or a class I protostar in the cloud. IRS 2 is a class 0 source confirmed by O’Linger et al. (1999) based on an SED using IRAS and SCUBA data. This source is associated with various signatures of extreme youth, including CO and H<sub>2</sub> molecular outflows (O’Linger et al., 1999; Eislöffel, 2000), Herbig-Haro objects (Bally et al., 1997), an H<sub>2</sub>O maser source (Anglada et al., 1989), and a centimeter continuum emission (Anglada & Rodríguez, 2002). IRS 3, also known as L1448 N(orth), is resolved into three components; IRS 3(A), 3(B), and 3(C), also known respectively as L1448 N(A), N(B), and NW. The former two constitute a binary separated by $``$ 2100 AU. Classification of the individual binary components in millimeter to far-infrared observations had been difficult due to crowding. However, recent mid-infrared (MIR) observations clearly resolved IRS 3(A) and 3(B) with an N-band detection only from 3(A). Based on this, Ciardi et al. (2003) argued that 3(B) is a class 0 protostar whereas 3(A) might be at the transition phase between class 0 and I stages. IRS 3(C), away from the binary by $``$20″, is another class 0 source (Barsony et al., 1998).
## 2 OBSERVATIONS
The X-ray observation was carried out on 2004 March 15 using ACIS (Advanced CCD Imaging Spectrometer; Garmire et al. 2003) on-board Chandra (Weisskopf et al., 2002). Four ACIS-I chips (I0, I1, I2, and I3) covered a $``$17′$`\times `$17′ region aimed at R. A. $`=`$ 3<sup>h</sup>25<sup>m</sup>36.4<sup>s</sup> and Decl. $`=`$ 30°44′58″ (J2000.0) for an exposure time of 67.9 ks. ACIS has sensitivity in the 0.5–9.0 keV band with a resolution of $``$150 eV at 6 keV and a superb point spread function (PSF) radius of $``$0.5″ at the on-axis position. The data were taken with the very faint telemetry mode and the timed exposure CCD operation with a frame time of 3.2 s.
The NIR observation was conducted on 2003 December 16 using FLAMINGOS (Elston et al., 2003) on the Cassegrain focus of the 4 m telescope at the Kitt Peak National Observatory. FLAMINGOS uses a HgCdTe HAWAII-2 array which has a format of 2048$`\times `$2048 with a pixel scale of $``$0.316″ pixel<sup>-1</sup>. With dithering observations of an amplitude of 15″, we obtained J-, H-, and K<sub>S</sub>-band images that cover a $``$11′$`\times `$11′ area centered close to the ACIS aim point at R. A. $`=`$ 3<sup>h</sup>25<sup>m</sup>35.1<sup>s</sup> and Decl. $`=`$ 30°45′19″ (J2000.0). The total exposure time was 13.5, 7, and 17 minutes respectively in the J, H, and K<sub>S</sub> bands. The seeing was $``$1″. Figures 1 (a) and (b) show the X-ray and NIR intensity maps. The outflow from L1448 mm appears clearly in the K<sub>S</sub> band (Fig. 1b) presumably due to H<sub>2</sub> $`v`$=1–0 S(1), vibrational-rotational transition line emission by excited hydrogen molecules.
## 3 DATA REDUCTION & ANALYSIS
### 3.1 X-ray Data
We reprocessed the Level 1 data distributed by the Chandra X-ray Center to obtain our X-ray event list. The latest calibration results (CALDB 2.28) were incorporated and the background rejection algorithm specific to data taken with the very faint mode was applied. Events were further cleaned by removing cosmic-ray afterglows and applying filters based on the event grade, status, and good-time intervals.
We detected sources using a wavelet technique in CIAO<sup>1</sup><sup>1</sup>1See http://asc.harvard.edu/ciao/ for detail. from three images of different energy bands; soft (0.5–2.0 keV), hard (2.0–9.0 keV), and total (0.5–9.0 keV). For each detected source, we used ACIS Extract<sup>2</sup><sup>2</sup>2See http://www.astro.psu.edu/xray/docs/TARA/ae\_users\_guide.html for detail. for a systematic event extraction, background subtraction, calculation of instrumental responses, and construction and binning of spectra and light curves. Source events were extracted from a 90% encircled energy polygon of the 1.5 keV PSF, whereas background events were locally extracted from an annular region around the source. The positions of sources were fine-tuned by correlating the event distribution with the PSF. As the X-ray source distribution is relatively sparse in the study field (Fig. 1a), we do not suffer any overlaps of extraction regions. The details of the procedure are described in Getman et al. (2005).
For each source, we derived X-ray photometry information consisting of the source position and its uncertainty, source count (0.5–8.0 keV), net count rate (NCR), probability of no source (PNS), photometry significance (PS), and mean energy (ME). The net count ($`C_{\mathrm{net}}`$) is calculated from the number of counts ($`C_{\mathrm{src}}`$ and $`C_{\mathrm{bkg}}`$) and the extraction area ($`A_{\mathrm{src}}`$ and $`A_{\mathrm{bkg}}`$) of the source and background regions as $`C_{\mathrm{net}}=C_{\mathrm{src}}C_{\mathrm{bkg}}\times (A_{\mathrm{src}}/A_{\mathrm{bkg}})`$, which is divided by the effective exposure time to derive NCR. PNS is an index for the detection significance, with which we test the null hypothesis that all detected counts are explained by background fluctuation;
$$\mathrm{PNS}=1_0^{C_{\mathrm{src}}1}𝑑NP(C_{\mathrm{bkg}}\left(\frac{A_{\mathrm{src}}}{A_{\mathrm{bkg}}}\right),N),$$
(1)
where $`P(\lambda ,N)`$ represents a Poisson probability distribution function of a mean $`\lambda `$ to have $`N`$ counts. PS is a metric for the reliability of photometry defined as
$$\mathrm{PS}=\frac{C_{\mathrm{net}}}{\sqrt{(\mathrm{\Delta }C_{\mathrm{src}})^2+(\mathrm{\Delta }C_{\mathrm{bkg}})^2\times (A_{\mathrm{src}}/A_{\mathrm{bkg}})^2}},$$
(2)
where the approximation by Gehrels (1986) is used to estimate $`\mathrm{\Delta }C_{\mathrm{src}}=1+\sqrt{C_{\mathrm{src}}+0.75}`$ and $`\mathrm{\Delta }C_{\mathrm{bkg}}=1+\sqrt{C_{\mathrm{bkg}}+0.75}`$. ME is a metric to represent the spectral hardness of sources. The value is the average incident energy ($`E`$) weighted by $`C_{\mathrm{net}}`$ of each energy bin and is defined as
$$\mathrm{ME}=\frac{𝑑EEC_{\mathrm{net}}(E)}{𝑑EC_{\mathrm{net}}(E)}.$$
(3)
Table 1 shows the results of the ACIS photometry, consisting of 72 sources with high significance both in the detection (PNS $``$ 1$`\times `$10<sup>-3</sup>) and the photometry (PS $``$ 2).
For each source, flux variability was tested using the K-S test and the Bayesian block segmentation technique (Scargle, 1998). We recognized flux to be variable if the number of Bayesian blocks exceeds two and the K-S probability is less than 1%. Four sources (Nos. 4, 7, 34, and 69) have variable flux, which are marked with daggers in Table 1.
For 18 sources with more than 40 counts, we performed spectral analysis. The spectra were fit with a one-temperature thin-thermal plasma model (the APEC model; Smith et al. 2001) with the metallicity value fixed to 0.3 solar to determine the best-fit absorption ($`N_\mathrm{H}`$), plasma temperature ($`k_\mathrm{B}T`$), and X-ray luminosity ($`L_\mathrm{X}`$) values. They were also fit with a power-law model to derive the best-fit $`N_\mathrm{H}`$, index ($`\mathrm{\Gamma }`$), and absorption-corrected flux ($`F_\mathrm{X}^{\mathrm{cor}}`$) values. Table 2 shows the results of 13 and 11 successful fits with the thermal and power-law models with a reasonable value of $`k_\mathrm{B}T`$$`<`$ 10 keV and $`\mathrm{\Gamma }<`$ 3.
### 3.2 NIR Data
All FLAMINGOS frames were reduced following the standard procedures using IRAF; i.e., subtraction of dark current, flat-fielding using a dome flat, subtraction of sky using median sky, removal of bad pixels, and trimming of edges. Several frames with readout failure were discarded. Artificial signals were masked out, which we found at pixels to the right and left of, or above and below bright sources by $``$128 pixels. Nebulosity extending larger than the dithering amplitude may remain in the median-sky-subtracted images; it does not affect our photometry. Referring to sources in the 2MASS all-sky survey catalog (Skrutskie et al., 1997), all FLAMINGOS frames were corrected for astrometry by a low-order polynomial function to compensate for the image distortion. They were also corrected for photometry to match 2MASS magnitudes. After these corrections, the frames were combined into the final J-, H-, K<sub>S</sub>-band images.
We extracted sources from the K<sub>S</sub>-band image (Fig. 1b) above a 5 $`\sigma `$ level using SExtractor (Bertin & Arnouts, 1996). The number of false positive detections is statistically negligible. After manually removing suspicious sources (introduced mostly by misidentifying outflow knots as point sources) and sources at the edges, we obtained 294 detections. We measured J-, H-, and K<sub>S</sub>-band magnitudes at their positions (Table 3). We matched the FLAMINGOS and the 2MASS sources and found that all the 67 2MASS sources in the FLAMINGOS field of view have FLAMINGOS counterparts.
Figure 2 shows the astrometric accuracy of the FLAMINGOS sources, where the difference of FLAMINGOS and 2MASS positions ($`\mathrm{\Delta }`$R. A. and $`\mathrm{\Delta }`$Decl.) are plotted for the 67 FLAMINGOS–2MASS counterpart pairs. The root mean square of $`\mathrm{\Delta }`$R. A. and $`\mathrm{\Delta }`$Decl. are 0.18″ and 0.33″, indicating that the FLAMINGOS position accuracy is $``$0.4″ ($``$1.2 FLAMINGOS pixel). The accuracy is better for most sources, as the values are exaggerated by three sources at the top left in Figure 2, all of which are located too far away from the image center to be corrected for distortion.
We estimated the detection completeness of the FLAMINGOS images by embedding artificial sources and detecting them with the same source detection algorithm. The detection rates at different magnitude bins were derived to estimate the 90% completeness limit as $``$18.0 mag, $``$16.5 mag, and $``$17.0 mag in the J, H, and K<sub>S</sub> bands.
### 3.3 NIR and Optical Identification of X-ray Sources
Seven X-ray sources were found to have FLAMINGOS counterparts within 1.5″ and five additional sources to have 2MASS counterparts outside of the FLAMINGOS view. X-ray positions in Table 1 are shifted by $``$0.3″ to match the 2MASS astrometry. The root mean square of the positional difference between 2MASS–X-ray pairs is $``$0.6″ ($``$1.3 ACIS pixel). Ten of the NIR-identified X-ray sources are optically identified with the USNO-B catalog (Monet et al., 2003). None of the NIR-unidentified X-ray sources were identified in the optical catalog. No X-ray emission was found from Herbig-Haro objects (Bally et al., 1997) or H<sub>2</sub> knots (Davis & Smith, 1995).
## 4 DISCUSSION
### 4.1 Nature of X-ray Sources
Two independent quantities from the X-ray photometry (ME and NCR) combined with the NIR identification give a clue to classify and reveal the nature of the 72 X-ray sources. Figure 3 shows a scatter plot of ME and NCR. Most of the X-ray sources with NIR counterparts (filled symbols) are distributed in the limited ME range of $``$2.5 keV, while those without NIR counterparts (open symbols) are distributed in the limited ME and NCR ranges of $``$2.0 keV and $``$1.0 ks<sup>-1</sup>. This implies that the X-ray sources consist of two major classes of different nature.
The separation of the two classes becomes more apparent when NIR flux is taken into account. Figure 4 shows a scatter plot of ME and the ratio of NIR flux in the K<sub>S</sub> band to X-ray flux in 0.5–8.0 keV band. Here, the X-ray flux ($`F_\mathrm{X}`$) is calculated by
$$F_\mathrm{X}=\frac{𝑑EEC_{\mathrm{net}}(E)}{\mathrm{EA}t_{\mathrm{exp}}}=\frac{\mathrm{ME}\times \mathrm{NCR}}{\mathrm{EA}},$$
(4)
where $`t_{\mathrm{exp}}`$ is the exposure time and $`\mathrm{EA}`$ is the effective area of the mirror averaged over the energy. The X-ray sources with NIR counterparts (large filled symbols) and those without FLAMINGOS counterparts (large open squares) are clearly separated into two groups. One group (group A) is comprised of X-ray source Nos. 4, 6, 7, 14, 15, 23, 30, 42, and 48 in the upper left quarter, while the other (group B) includes Nos. 11, 16, 39, and FLAMINGOS-unidentified X-ray sources (large open squares) at the bottom. NIR-unidentified X-ray sources outside of the FLAMINGOS field of view (small open symbols) are hard to classify into either group as the K<sub>S</sub>-band upper limit of 2MASS is not deep enough. This illustrates the importance of our NIR observation deeper than 2MASS, even if it results in no detection from X-ray sources.
We conclude that the sources in group A are mostly YSOs that belong to L1448, while those in group B are background AGNs. We note that NIR identification of an X-ray source in star-forming regions does not always guarantee its stellar origin because AGNs, which constitutes the largest contaminant, can be detected even at the 2MASS depth. Instead, we employ the NIR to X-ray flux ratio as an effective measure to separate these two classes. The optical to soft X-ray flux ratio was used for a similar classification purpose in X-ray blank-sky surveys by Einstein and ROSAT satellites (Maccacaro et al., 1988; Krautter et al., 1999). We use the NIR to total band X-ray flux to be suitable for star-forming regions, which are routinely highly obscured. In (low-mass) main sequence and pre–main-sequence (PMS) sources, X-ray emission originates from stellar coronae while NIR emission originates from photosphere and from circumstellar matter. In AGNs, on the other hand, X-ray emission is from accretion disks around the central black hole while NIR is from hot dust and jets. The different origins bring different typical values of NIR to X-ray flux ratio in stars and AGNs.
Several lines of evidence support this conclusion. First, all sources with $``$40 counts in group A (Nos. 4, 6, 15, 23, 42, and 48) have a successful fit with thermal spectral models (Table 2) but not with power-law models with reasonable parameters. On the other hand, all sources with $``$40 counts in group B (Nos. 11, 16, 50, and 53) have a successful fit with power-law models. Three (Nos. 11, 50, and 53) of them are not fit with thermal models (Table 2). The clear distinction in X-ray spectral features is consistent with the idea that (1) these groups of sources are different in nature, (2) group A sources are stellar X-ray emission, which typically show thermal spectra of $`k_\mathrm{B}T`$$``$ 1 keV, and (3) group B sources are AGNs, which typically show power-law spectra of $`\mathrm{\Gamma }`$1.7. This conversely enables us to use X-ray spectral features to classify sources for which it is unclear which group to belong to. Nos. 54 and 72 are group B sources for successful fits only by the power-law model (Table 2).
The second argument is the X-ray flux variability of a typical time scale of $``$10 ks. Stellar X-ray emission from coronae is often characterized by X-ray flares due to magnetic reconnection, whereas no such short time-scale variability is expected from AGNs. Sources with X-ray variability are seen in group A (Nos. 4 and 7) and not in group B. This conversely indicates that sources with similar variability (Nos. 34 and 69) belong to group A.
Third, NIR magnitudes of the X-ray sources in group A are distributed in $`m_{K_S}`$ $``$ 9–13 mag. Contamination by extragalactic sources in this magnitude range is negligible as the expected number of galaxies in the FLAMINGOS field of view is $``$1 in $`m_{K_S}`$ $``$ 14 mag (Tokunaga, 2000). The color-magnitude diagram (Fig. 5) shows that group A sources have consistent NIR magnitudes for low mass or very low mass stars at a distance of 300 pc. Sources in group B, on the contrary, have $`m_{K_S}`$ $``$ 15 mag (Tables 1 and 3), which is too faint to be stars at that distance.
Fourth, the NIR to X-ray flux ratio of group A sources ($`F_{K_\mathrm{S}}`$/$`F_\mathrm{X}`$ $``$ 10<sup>1</sup>–10<sup>4</sup>) is consistent with typical values of the X-ray to the bolometric luminosity ratio of $`L_\mathrm{X}`$/$`L_{\mathrm{bol}}`$ $``$ 10<sup>-5</sup>–10<sup>-2</sup> (Preibisch et al., 2005) of PMS sources. Assuming an age of 2 Myr (the same with IC 348; Muench et al. 2003), and using $`L_{K_\mathrm{S}}`$/$`L_{\mathrm{bol}}`$ $``$ 4–9$`\times `$10<sup>-2</sup> (0.08–1.4 $`M_{}`$ PMS source; Baraffe et al. 1998; D’Antona & Mazzitelli 1994) and $`L_{K_\mathrm{S}}`$/$`L_\mathrm{X}`$ $``$ $`F_{K_\mathrm{S}}`$/$`F_\mathrm{X}`$ $``$ 5$`\times `$10<sup>2</sup> (the median value of group A sources), we obtain the value of $`L_\mathrm{X}`$/$`L_{\mathrm{bol}}`$ $``$ 0.8–1.8$`\times `$10<sup>-4</sup> for group A sources. A similar $`L_\mathrm{X}`$/$`L_{\mathrm{bol}}`$ value is obtained if we assume that these sources are main sequence stars (Tokunaga, 2000; Drilling & Landolt, 2000), which is higher than their typical values (10<sup>-7</sup>–10<sup>-5</sup>; Preibisch et al. 2005). We thus consider that most of the group A sources are YSOs in the cloud and not main sequence stars in the line of sight. The distinction between YSO and main sequence natures should be confirmed with follow-up spectroscopic studies. Group B sources, on the other hand, have $`F_{K_\mathrm{S}}`$/$`F_\mathrm{X}`$ values of 10<sup>-1</sup>–10<sup>1</sup>, which is consistent with the values seen in an X-ray-selected sample of AGNs (Watanabe et al., 2004).
Figure 6 shows a color-color diagram of NIR-identified X-ray (squares) and FLAMINGOS (circles) sources. All NIR-identified X-ray sources in group A (Nos. 4, 6, 7, 14, 15, 23, 30, 42, and 48) have NIR colors consistent with being reddened weak-line T Tauri or main sequence stars. None of them shows apparent NIR excess emission, which does not exclude the possibility that they are classical T Tauri stars or protostars, as K<sub>S</sub>-band observations are not completely sensitive to excess emission from circumstellar disks (e.g., Haisch et al. 2001). The only source with significant J-, H-, and K<sub>S</sub>-band detections in group B (No. 16) has an inconsistently red color for reddened stars of $`m_J`$$`m_{K_S}`$ $`=`$ 2.0 mag, which is commonly seen among 2MASS red AGNs (Cutri et al., 2001). The density of such sources is $``$0.5 arcdeg<sup>-2</sup> at K<sub>S</sub> $`<`$ 14.5 mag (Cutri et al., 2001), making the probability $``$5% to have a 2MASS red AGN in an ACIS field.
In total, twelve X-ray sources (Nos. 4, 6, 7, 14, 15, 23, 30, 34, 42, 48, 68, and 69) have X-ray and NIR features consistent of being low-mass YSOs in this cloud. Two of them are known sources; No. 7 is IRS 1 and No. 42 is an A5 type star (Roeser & Bastian, 1988) that consists of a double (Worley & Douglass, 1997). Main sequence A type stars are not expected to have X-ray emission due to the lack of any plasma creation mechanisms (Schmitt et al., 1993). The X-ray emission in No. 42 may come from a low-mass companion, which mimics a $`F_{K_\mathrm{S}}`$/$`F_\mathrm{X}`$ value consistent with a PMS source (Fig. 5). The remaining ten sources are candidates of new YSOs in L1448. The brightness histogram of the rest of the X-ray sources is consistent with a log $`N`$–log $`S`$ relation of AGNs (Moretti et al., 2003).
### 4.2 X-ray and NIR Emission from Embedded Sources
#### 4.2.1 Infrared Excess Sources
From Figure 6, ten FLAMINGOS sources (large open circles) are found to have NIR excess emission and thus are candidates for new embedded sources. Daggers are given for these sources in Table 3. Among them, FLAMINGOS Nos. 23 and 149 are promising YSO candidates with estimated magnitude uncertainty of $``$ 0.05 mag.
#### 4.2.2 IRS 1
Both X-ray and NIR emission was found from IRS 1 (ACIS No. 7 in Table 1 and FLAMINGOS No. 2 in Table 3 or 2MASS 03250943$`+`$3046215). The X-ray and NIR properties are consistent with the nature of this source as a classical T Tauri star or a class I protostar; (1) it shows X-ray flux variation commonly seen in X-ray emission of stellar magnetic activity origin, (2) its ME and $`F_{K\mathrm{s}}`$/$`F_\mathrm{X}`$ values indicate this source to be a YSO (Fig. 4), and (3) a nebulous structure is seen in the FLAMINGOS K<sub>S</sub>-band image (Fig. 1b). The NIR magnitudes of this source ($`m_J`$ $``$ 12.5 mag, $`m_H`$ $``$ 11.0 mag, and $`m_{K_S}`$ $``$ 10.1 mag) correspond to an object with the mass of 0.5–1.0 $`M_{}`$ at 300 pc obscured by $`A_\mathrm{V}`$ 7 mag.
#### 4.2.3 IRS 3(A)
No emission is seen at the position of IRS 3(A) in our NIR images. The 5$`\sigma `$ level of the local sky noise corresponds to $``$18.8 mag, $``$17.7 mag, and $``$17.5 mag in our J-, H-, and K<sub>S</sub>-band images. The SED slope defined as $`\alpha =(d\mathrm{log}\lambda F_\lambda )/(d\mathrm{log}\lambda `$) between 2.2 µm (K<sub>S</sub> band) and 10 µm (N band) is commonly used to classify YSOs (Wilking et al., 1989), using the fact that the 2–10 µm SEDs at a younger evolutionally stage are more dominated by emission from circumstellar material. The lower limit of $`\alpha `$ of IRS 3(A) is $``$3.2 from the non-detection in the NIR bands and a MIR detection in the N band (Ciardi et al., 2003). Here, we calculated the $`\alpha `$ value as the index between H and N bands, as the K<sub>S</sub>-band image is noisy at the position due to the emission arising from the outflow of IRS 3(A) itself. The value thus is a lower limit also in the sense that H-band emission is subject to contamination from reflected photospheric emission.
Greene et al. (1994) gave a qualitative classification of YSOs using $`\alpha `$, where $`\alpha <1.6`$ for class III, $`1.6<\alpha <0.3`$ for class II, and $`\alpha >0.3`$ for class I sources. Almost all known class I protostars have $`\alpha `$ 2 with the only exception of WL 22 ($`\alpha =`$ 3) in $`\rho `$ Ophiuchi dark cloud (Wilking et al., 1989). Sources with a steeper index than $``$ 3 have been rarely known, as they are quite difficult to detect in K band images that traditionally pilot protostar searches. However, recent progress of sensitive MIR imaging both by ground- and space-based telescopes is changing the situation by reporting several protostars with very steep spectra with indices of $`\alpha >3`$ (Cep E by Noriega-Crespo et al. 2004 and source X<sub>E</sub> in R CrA by Hamaguchi et al. 2005 and private communications with K. Nedachi). The index of $``$3.2 of IRS 3(A) places this source in the class of very steep spectrum protostars. More sources in this emerging class are expected to be discovered with Spitzer observations of near-by star-forming regions. Some of them are class 0 protostars (e.g., Cep E mm; Noriega-Crespo et al. 2004) and others are K-band-unidentified class I protostars, or in between.
An unusual concentration of three X-ray photons within a $``$0.5″ radius ($``$90% encircled energy radius at $``$0.4′ off-axis angle) can be found at the position of IRS 3(A). No such concentration is seen within 1′ of the on-axis position. Although it was not recognized as a source in the initial X-ray source search using the wavelet technique, its PNS value of 5$`\times `$10<sup>-4</sup> indicates the existence of an X-ray source of $`F_\mathrm{X}`$ $``$ 5$`\times `$10<sup>-16</sup> ergs s<sup>-1</sup> cm<sup>-2</sup>. The facts that the position of the concentration is consistent with that of IRS 3(A), that the energy of all photons is concentrated in the hard (2–5 keV) band, and that the arrival times of all photons fall in a limited time range of $``$15 ks in the $``$68 ks observation, which is indicative of a flare, may further strengthen the claim. There is no indication of temporal non-uniformity of background counts. A similar weak but plausible X-ray emission was found in an embedded source (FIR4) in NGC 2024 by Chandra (Skinner et al., 2003).
Given the capability of Chandra and MIR telescopes (Spitzer and ground-based facilities), it is now becoming possible to investigate X-ray emission from protostars with a very steep spectrum ($`\alpha >3`$) even without K<sub>S</sub>-band identifications. If IRS 3(A) is indeed an X-ray source, it has the steepest spectral index among YSOs with X-ray emission. Although it is not yet certain whether a steeper $`\alpha `$ value indicates a younger age of protostars, X-ray observations for steep spectrum sources are a key to answer how early X-ray emission starts and how intense it is at the earliest stage of star formation. These are important questions to understand star formation because the coupling between gas and magnetic fields can be largely affected by X-ray ionization (Feigelson & Montmerle, 1999).
#### 4.2.4 Class 0 Sources
No significant X-ray or NIR emission was found from the four known class 0 sources; L1448 mm, IRS 2, IRS 3(B), and IRS 3(C). The upper limit (PNS $`=1\times 10^3`$) of the X-ray flux in the 0.5–8.0 keV is $``$0.5–1.0$`\times `$10<sup>-15</sup> ergs s cm<sup>-2</sup>, whereas the upper limits (5 $`\sigma `$) in the NIR regime are $``$18.8 mag (J), $``$17.7 mag (H), and $``$17.5 mag (K<sub>S</sub>), except for L1448 mm and IRS 3(B) in the K<sub>S</sub>-band, which are contaminated by local diffuse emission.
The lack of X-ray detection from class 0 sources is widely reported. Froebrich (2005) compiled 28 well documented bona-fide class 0 sources. Sixteen of them were observed by Chandra and nine in references are reported to have no X-ray detection (L1448 NW and C; this work, NGC 1333 I2, I4A, and I4B; Getman et al. 2002, OMC 3 MMS6; Tsujimoto et al. 2002, HH 24 MMS and HH 25 MMS; Simon et al. 2004, and VLA 1623; Imanishi 2003). An additional six sources are reported with no detection in XMM-Newton references (IRAS 03256$`+`$3055; Preibisch & Zinnecker 2002, HH 211 MM; Preibisch & Zinnecker 2004, Serpens S68N, SMM2, SMM3, and SMM4; Preibisch 2004). Considering that these sources are deeply embedded, the lack of X-ray detection does not necessarily indicate that they do not have any X-ray-emitting activities, although they are still difficult for X-ray as well as MIR investigations with current technologies.
The authors express gratitude for Kentaro Motohara for his help in the observation at Kitt Peak, Pat Broos for his advice in X-ray data reduction using ACIS Extract and proofreading, and Eric D. Feigelson for discussion. We acknowledge financial supported by the Japan Society for the Promotion of Science (M. T.), a Grant-in-Aid for Scientific Research of the Ministry of Education, Culture, Sports, Science and Technology (No. 15740120; Y. T.), a Chuo University Grant for Special Research (Y. T.), the Saneyoshi Scholarship Foundation (Y. T.), and the Chandra guest observer grant (SAO G04–5009X). FLAMINGOS was designed and constructed by the IR instrumentation group (PI: R. Elston) at the University of Florida, Department of Astronomy, with support from NSF grant AST97-31180 and Kitt Peak National Observatory. IRAF is distributed by the National Optical Astronomy Observatories, which are operated by the Association of Universities for Research in Astronomy, Inc., under cooperative agreement with the National Science Foundation. This publication makes use of data products from the Two Micron All Sky Survey, which is a joint project of the University of Massachusetts and the Infrared Processing and Analysis Center/California Institute of Technology, funded by the National Aeronautics and Space Administration and the National Science Foundation. Facilities: CXO(ACIS), KPNO:4m(FLAMINGOS) |
warning/0506/nucl-ex0506003.html | ar5iv | text | # New precision determination of 𝑔_𝑃 and 𝐺_𝐹: the MuXperiments at PSI
## 1 Precision determination of $`G_F`$ \- the MuLAN experiment
We have seen impressive advances in our precise knowledge of many parameters defining the electroweak interaction within the Standard Model. However, the value of one of the most fundamental weak parameters, the Fermi coupling constant $`G_F`$, has not been improved in over two decades (see Fig.1). Usually $`G_F`$ is determined via a measurement of the muon lifetime $`\tau _\mu `$
$$\frac{\mathit{1}}{\tau _\mu }=\frac{G_F^\mathit{2}m_\mu ^\mathit{5}}{\mathit{192}\pi ^\mathit{3}}F\left(\frac{m_e^\mathit{2}}{m_\mu ^\mathit{2}}\right)\left\{\mathit{1}+\frac{\mathit{3}}{\mathit{5}}\frac{m_\mu ^\mathit{2}}{M_W^\mathit{2}}\right\}\left(\mathit{1}+\Delta _{\mathrm{𝑄𝐸𝐷}}(\alpha _{m_\mu })\right),$$
(1)
with $`F(x)=18x12x^2`$ ln$`x+8x^3x^4`$ . The QED corrections within the Fermi Model, $`\Delta _{\mathrm{𝑄𝐸𝐷}}`$, are included in this definition.
Within the Standard Model one can derive the relation
$$G_F=\frac{\pi \alpha (\mathit{0})}{\sqrt{\mathit{2}}M_W^\mathit{2}\left(\mathit{1}\frac{M_W^\mathit{2}}{M_Z^\mathit{2}}\right)}(\mathit{1}+\Delta _r),$$
(2)
with weak radiative corrections being summarized in $`\Delta _r`$ . The calculated quantity $`\Delta _r`$ depends on the entire set of input parameters, e.g. $`M_Z`$, $`M_{Higgs}`$, $`m_{top}`$, $`\alpha `$. Recently, these calculations were improved to the sub-ppm level by including numerically important QCD and electroweak higher-order terms up to 2-loop level. Therefore a precise comparison between theoretical and experimental values, e.g. for $`M_W`$ (Eq. 2) is possible . Consequently, $`G_F`$ sets important constraints on the Standard Model and SUSY parameters. Furthermore, $`G_F`$ sets the weak scale and is intimately connected to the vacuum expectation value of the Higgs field. The best possible experimental measurement of $`G_F`$ at the present technological limit is therefore highly desirable, as the 18 ppm precision limit on the PDG average on $`\tau _\mu `$ is dominated by experimental counting statistics.
The MuLAN experiment (Muon Lifetime ANalysis), intents to measure a total of $`10^{12}`$ $`\mu ^+`$ decays, in order to achieve a 1 ppm statistical error in the lifetime. Since the status report on the MuXprogram in we have achieved substantial progress. A modification of the continuous high intensity muon beam line at the Paul Scherrer Institut was necessary to enable the collection of $`10^{12}`$ events within a reasonable time. We have built an electrostatic kicker which applies an artificial time structure to the DC beam in the $`\pi `$E3 area and found a kickable beam tune which provides up to 8 MHz of muons. Following a 5 $`\mu `$s muon collection period in the target, the kicker deflects the beam for 27 $`\mu `$s while muon decays are measured.
MuLAN is designed to minimize the systematic errors in several ways:
$``$ Muon polarization: The beam muons are highly polarized, and the preferential emission of decay positrons in muon spin direction could cause a position- and time-dependent positron detection efficiency as polarized muons rotate in an external magnetic field. We are currently investigating two specific targets: i) Arnokrome-3 (Fig.2b) is a proprietary chromium-cobalt-iron alloy sheet, which, due to an internal field of a few Tesla, precesses muons very fast with respect to muon decay. Therefore polarization effects are negligible. ii) A solid sulfur target which maximizes the depolarization of the beam muons. It is placed in a homogeneous 120 Gauss magnetic field which causes a fast visible muon rotation and allows us to fit the corresponding decay positron asymmetry. We are presently investigating samples of $`10^{10}`$ decay positrons from each target to select the optimal material choice. Additionally, a polarization-preserving silver target is being used for control purposes.
MuLAN’s highly modular detector (Fig.2a) of 174 coincident scintillator tile pairs in “soccer ball” geometry allows us to compare opposite counters, thus strongly reducing precession effects in the count rate sum.
$``$ “Sneaky muons:” A fast thin entrance muon counter (EMC) records beam muons and looks for muons sneaking in during the measurement interval. A magnet positioned behind the EMC precesses the tiny fraction of muons stopped in the detector materials, otherwise they too could cause small detection inefficiencies.
$``$ Off-target muon stopping: Muon stops before the target are minimized by decreasing the materials in the muon path. Consequently we installed a helium bag (Fig.2b) and used very thin mylar windows and EMC materials.
$``$ Pile-up: The high detector modularity and the fast scintillator response time reduce pile-up. Additional time resolution will be gained via new 500 MHz waveform digitizers (WFD) presently under construction, which will provide a double pulse resolution better than 4 ns. Final WFD implementation to all detector channels is planned for 2005.
The MuLAN detector (Fig.2a) was successfully commissioned in 2004 and yielded its first physics data. Fig.3 shows a 10 minute snapshot from our AK3 target data. We used multi-hit TDCs for detector readout. The muon accumulation time and the decay recording time are indicated. Our present analysis goal with this data is a 5 ppm precision determination of $`G_F`$. We intend to collect the full statistics in 2006.
## 2 Precision determination of $`g_P`$ \- the MuCAP experiment
The $`VA`$ description of weak interactions has been tested to a high precision. Processes involving structureless fermions, e.g. muon decay, show equal vector (V) and axial-vector (A) coupling. In $`\beta `$-decay as well as in nuclear muon capture on the proton
$`\mu ^{}+p\nu _\mu +n`$ , the axial coupling $`G_A`$ is modified due to hadronic effects caused by the involved nucleon. Muon capture occurs at higher four-momentum transfer $`q=0.88m_\mu ^2`$ than $`\beta `$-decay. Lorentz invariance constrains the corresponding weak current matrix elements to six independent terms,
$`V_\mu `$ $`=`$ $`G_V(q^2)\gamma _\mu {\displaystyle \frac{iG_M(q^2)}{2m_N}}\sigma _{\mu \nu }q^\nu +{\displaystyle \frac{G_S(q^2)}{m_\mu }}q_\mu `$ (3)
$`A_\mu `$ $`=`$ $`G_A(q^2)\gamma _\mu \gamma _5+{\displaystyle \frac{G_P(q^2)}{m_\mu }}\gamma _5q_\mu +{\displaystyle \frac{iG_T(q^2)}{2m_N}}\sigma _{\mu \nu }q^\nu \gamma _5,`$ (4)
with corresponding weak form factors $`G_i`$ (<sub>i</sub> = scalar, pseudoscalar, vector, axial-vector, tensor, weak magnetism); mass of the nucleon $`m_N`$ and muon $`m_\mu `$. Because of G-symmetry $`G_S`$ and $`G_T`$ vanish . Due to the momentum dependence, only $`G_A`$ and $`G_V`$ contribute in $`\beta `$-decay at very low $`q^2`$. Nuclear muon capture is the process most sensitive to $`G_P`$. Therefore, $`G_P(0.88m_\mu ^2)`$ is dubbed induced pseudoscalar coupling constant $`g_P`$. While the values of $`G_V`$, $`G_A`$ and $`G_M`$ are established on the $`10^3`$ to $`10^4`$ level , the situation is totally different for the induced pseudoscalar $`g_P`$.
The theoretical view, historically based on PCAC and pion pole dominance, and recently strictly derived within chiral perturbation theory ($`\chi `$PT) , is remarkably precise:
$`G_P(q^2)`$ $`=`$ $`{\displaystyle \frac{2m_\mu g_{\pi NN}F_\pi }{m_\pi ^2q^2}}{\displaystyle \frac{1}{3c^2\mathrm{}^2}}G_A(0)m_\mu m_Nr_A^2,`$ (5)
$`g_P`$ $`=`$ $`(8.74\pm 0.23)(0.48\pm 0.02)=8.23\pm 0.23,`$ (6)
depending on the exact values of the pion-nucleon coupling constant $`g_{\pi NN}`$ and the mean axial radius of the nucleon $`r_A`$. The Standard Model based calculation of the singlet muon capture rate by Govaerts and Lucio-Martinez has reached 0.55% precision. This precision in calculation will allow a high precision measurement to distinguish between the pion pole contribution to $`g_P`$ and the correction term. Moreover, such a measurement will also set tight limits on various theoretical scenarios beyond the Standard Model .
The present experimental knowledge of $`g_P`$ is unsatisfying, and discrepancies cause considerable debate. Determinations via ordinary muon capture in hydrogen (OMC) , <sup>3</sup>He and larger nuclei essentially confirm the theory result. However the precision of the latter is troubled by model dependencies. A radiative muon capture on the proton (RMC) experiment , which measured an additionally very rarely emitted high energy $`\gamma `$ ray in conjunction with the muon capture process, yielded a different result<sup>1</sup><sup>1</sup>1While the RMC process has a 10<sup>5</sup> times lower branching than OMC, the emitted $`\gamma `$ can have energies up to 100 MeV. Therefore these $`\gamma `$’s come closer to the pion pole and the measurement is in principle four times more sensitive to $`g_P`$ than OMC.. The present most likely explanation lies in the insufficient knowledge of the complex kinetics of negative muons in hydrogen. The $`\mu p`$ atom is formed in statistically populated singlet and triplet states, and the muon capture rate from these states differ by a factor of $``$60 due to the strong spin dependence of weak interactions. An exact knowledge of the spin populations is therefore mandatory for the interpretation of a measurement. The initial $`\mu p`$ populations are modified via spin flip and molecular formation processes, which eventually yield a fraction of muonic hydrogen molecules ($`pp\mu `$). The respective rates for both types of process scale with density, and hence are high in the liquid hydrogen targets which were used in Refs. and . However, a large $`\mu `$-molecular population causes a large correction to the lifetime and hence a corresponding uncertainty. Specifically the ortho to para spin-flip rate in $`pp\mu `$ molecules $`\lambda _{OP}`$ (Fig.4b) is a prime suspect to cause the experimental discrepancy in $`g_P`$. Figure 4, updated from , shows the $`\lambda _{OP}`$ dependence of the OMC and RMC results. The controversial $`\lambda _{OP}`$ values are also shown. Two experimental values from Saclay and TRIUMF , which were obtained together within the same experiments performing the OMC and RMC measurements, strongly disagree, and the only theoretical calculation does not clarify the situation. It is evident that using the $`\lambda _{OP}`$ rate from the same measurements brings the $`g_P`$ value in agreement with theory. However, the use of the other $`\lambda _{OP}`$ value enlarges the disagreement to the prediction or lowers it to an unphysical negative value. Clearly only a new determination of $`g_P`$ independent of $`\lambda _{OP}`$ can resolve this situation.
The experimental principle of the MuCAP (Muon CAPture) experiment is based on the measurement and comparison of the decay time of positive and negative muons in hydrogen. The MuCAP experiment is designed to overcome the multiple difficult problems of previous experiments. The important conceptual advantage of MuCAP is the selection of target hydrogen at gaseous density (10 bar at room temperature), which minimizes the kinetics dependence of the result on $`g_P`$ as shown in Fig.4a. At low densities, muon capture occurs almost exclusively from the singlet state on the proton, and even a large estimated error on $`\lambda _{OP}`$ results only in a systematic error on the 10 ppm level. The full setup is shown in Fig.5. The active gaseous hydrogen target, a time projection chamber (TPC), allows for a full 3-dimensional reconstruction of the muon path to its stopping point and therefore a selection of clean muon stops away from walls and wires. The observed stopping distribution of muons is shown in Fig.6a.
The TPC also detects muon capture events on impurity atoms ($`Z>1`$) via the very large signals generated from capture products. Thus the TPC serves also as a very sensitive impurity monitor. The high rates of muon transfer to and nuclear capture on high-$`Z`$ atoms (Fig.4b) can cause a deflection of the exponential lifetime even at very low impurity concentrations as these rates are typically orders of magnitude higher than muon decay. In order to minimize this effect, target purity requirements are very stringent, with the goal to be on the 10<sup>-9</sup> contamination level for the sum of high-$`Z`$ atoms. Additionally, the exact knowledge of these contaminations is necessary to calculate the correction factor to the lifetime. Consequently, the hydrogen gas after production in electrolysis is filled via a palladium filter and continuously run through a Zeolite based purification system (CHUPS). The CHUPS system is specifically designed to maintain the hydrogen flow with negligible variations in density or hence in TPC gain. In fall 2004 we maintained clean target conditions (as low as 70 ppb impurities) for over 5 weeks.
MuCAP fully separates the muon and electron detectors to avoid dangerous cross-correlations. Decay electron times are measured in the scintillator hodoscope (eSC) surrounding the hydrogen vessel. Two cylindrical wire chambers track the electron back in 3 dimensions to its $`\mu `$-stop origin, thus largely reducing the background.
The impact parameter, defined as the minimal distance between detected muon stop and electron track, serves as an important handle on a very subtle systematic effect on the lifetime: the diffusion of muons which transfer to deuterium, an isotope always present in hydrogen. Although we are using special deuterium-depleted hydrogen, (“protium”), with deuterium content as low as 1.5 ppm, this deuterium concentration is still high enough to cause a visible effect in our setup via electron tracks with dislocated origin. This dislocation, due to $`\mu d`$ diffusion over macroscopic distances is possible because of a Ramsauer-Townsend minimum in the $`\mu d+p`$ scattering cross-section at low energies. Such $`\mu d`$’s can hit a wall or wire, transfer and then undergo unobserved nuclear capture outside the sensitive volume. The electron wire chamber tracking identifies such events and the corresponding lifetime dependence on the impact parameter will allow us to determine the deuterium contamination in situ . Additionally, we have been developing a precision trace deuterium detection method via $`pd\mu `$ fusion .
As in MuLAN, MuCAP applies a magnetic field to control muon spin rotation effects in the $`\mu ^+`$ measurement. A water-cooled aluminum coil was developed to minimize decay electron scattering.
Eventually MuCAP needs 10<sup>10</sup> cleanly observed decay electrons and positrons to statistically reach the goal of 1 % in the capture rate, reflecting 10 ppm in the respective muon lifetimes. This is possible in several weeks of running in the fully pile-up protected mode. In the future we also intend to use the MuLAN kicked beamline with the MuCAP experiment and operate it in a “muon on request mode” .
Fig.7 shows a preliminary lifetime plot from our fall 2004 data. One can clearly see the huge improvement in background reduction due to the implementation of the two cylindrical wire chambers. The shown $``$2$`\times `$10<sup>9</sup> “clean” $`\mu `$ decay events are presently being analyzed, and a first result on $`g_P`$ is expected in late 2005.
## 3 Summary
Both experiments, MuCAP and MuLan, are on the way to reach ultimate precision in their respective measurements. First physics results are expected in late 2005.
### Acknowledgments
I would like to express my cordial thanks to all colleagues in the MuLAN and MuCAP collaboration for creating a great scientific working environment and for their personal friendship. This work was supported by the US Department of Energy and the National Science Foundation. |
warning/0506/hep-ex0506068.html | ar5iv | text | # Charm Hadrons from Fragmentation and 𝐵 decays in 𝑒⁺𝑒⁻ Annihilation at √𝑠=10.6 GeV
## I Introduction
Over recent years perturbative quantum chromodynamics (pQCD) has shown impressive agreement with various inclusive measurements at $`e^+e^{}`$ colliders at many center-of-mass energies (CME) ranging from 14 GeV up to 206 GeV. These measurements utilised variables called event shapes or jet rates, see OPAL\_PETRAqcd for such an analysis. These are inclusive variables, whose values are calculated from the four-momenta of all particles in an event.
Other properties, such as the momentum spectra of charged or neutral particles, have also been measured, but their prediction has proven to be more difficult. The necessary calculations have to cover the entire energy range from the production of the partons at the CME down to the scale of the hadron masses (typically $`1\mathrm{GeV}/c^2`$), at which hadronisation occurs. Typically, powers of the form $`\mathrm{log}(Q^2/m^2)`$ arise when quark masses are taken into account, making pQCD calculations difficult to interpret.
Attempts have been made to extend the applicable range of pQCD to lower scales. These attempts have to be validated, for example by comparing so-called fragmentation functions. Due to the scaling violation of QCD, a fundamental property of this theory, the fragmentation function for a given particle depends explicitly on the CME. This energy dependence must follow the Dokshitzer-Gribov-Lipatov-Altarelli-Parisi (DGLAP) DGLAP evolution equations.
Thus, the fragmentation functions have to be properly evolved. Monte Carlo (MC) generators which include this scaling can be used instead of analytical evolution. Common MC generators which include this scaling are JETSET JETSET , (its variant) PYTHIA PYTHIA and HERWIG HERWIG .
These MC generators are also needed to model hadronisation, the transition of partons into hadrons, which cannot be calculated from first principles within QCD. Various models are implemented in MC generators. These can be distinguished by comparing the (identified) heavy hadron momentum spectra predicted by each model to the spectra seen in data.
Fragmentation functions for heavy quarks are attractive both experimentally and theoretically. Concerning theory, mass effects in the matrix elements only have to be considered for the heavy quark; in the limit of $`m_{light}0`$, a pQCD calculation based on an effective Lagrangian reduces the complexity of the calculation compared to the case of light quark fragmentation.
Experimentally, it is important to measure heavy quark fragmentation functions as their shapes are different from the corresponding functions for light quarks; such a measurement is furthermore straightforward, as very often hadrons containing heavy quarks can easily be identified. Since the production of heavy quarks is strongly suppressed in both the perturbative splitting of one parton to many partons (the so-called “parton shower”) and in hadronisation, a heavy quark found in an event will most likely be produced in the primary interaction.
At LEP and SLD, $`b`$ quark fragmentation functions have been measured with high precision ALEPHb ; DELPHIb ; OPALb ; SLDb . These measurements found that these fragmentation functions are in fact close to the ones of light quarks, suggesting that one combined model for all five flavours might describe the measured momentum spectra better than functions which have been introduced for heavy quarks alone. These collaborations have also published measurements of $`c`$ quark fragmentation functions ALEPHc ; OPALc , but with large statistical uncertainties due to the small product of the branching fraction and reconstruction efficiency for the various final states. Some commonly used fragmentation functions are described by the models of Peterson et al. Peterson , of Kartvelishvili et al. Kartvelishvili and of Collins and Spiller ColSpi , as well as by the models of the Lund group Lundsymm and one of its variants by Bowler Bowler .
For charm quark fragmentation functions at lower energies, the most recent published results for $`D^0`$, $`D^+`$, $`D^0`$ and $`D^+`$ are those of CLEO CLEOprel . The analysis presented here has better statistical precision as the data sample is five times larger. Other measurements are more than 10 years old CLEOc ; ARGUSc ; their data sample is over three orders of magnitude smaller than that used in this analysis. The systematic uncertainties are reduced significantly and are comparable to those in CLEOprel . For a recent review of fragmentation function measurements and theory, see Biebel .
A measurement of $`D^0`$ and $`D^+`$ performed by the same experiment on the same data set allows for an easy comparison of charged meson production rates and momentum spectra, as well as a comparison of the momentum-dependent production of secondary- to primary-produced mesons. The measurement of the excited states $`D^0`$ and $`D^+`$ allows the determination of the feed-down contribution to the ground states $`D^0`$ and $`D^+`$ and also a momentum-dependent determination of $`V/(V+P)`$, the ratio of the production rates of vector and the sum of vector and pseudo-scalar mesons. A comparison between $`D_s^+`$ production, and the production of $`D^0`$ and $`D^+`$, can be used to determine the fraction of $`s`$ quark production in hadronisation. Comparing the results for the $`\mathrm{\Lambda }_c^+`$ to those of the $`D`$ mesons makes a study of the baryon production mechanism possible.
In addition to charm fragmentation in the $`e^+e^{}`$$``$$`c\overline{c}`$ continuum, charmed hadrons in $`e^+e^{}`$ annihilation events can be produced in decays of b-hadrons. The dataset for this analysis includes events above the production threshold for $`B\overline{B}`$ pairs, at the $`\mathrm{{\rm Y}}(4\mathrm{S})`$ resonance, so the lower momentum hadrons include contributions from $`B^0`$ and $`B^+`$ decays. This allows a measurement of the production rate of charmed hadrons in B-meson decay.
## II Data Sample and Event Selection
This analysis uses data recorded at the Belle detector at the KEKB accelerator. The KEKB $`e^+e^{}`$ collider is a pair of storage rings for electrons and positrons with asymmetric energies, $`8.0\mathrm{GeV}`$ ($`e^{}`$) and $`3.5\mathrm{GeV}`$ ($`e^+`$), and a single intersection point with a 22 mrad crossing angle. The beam energies are tuned to produce an available CME of $`\sqrt{s}=10.58\mathrm{GeV}`$, corresponding to the mass of the $`\mathrm{{\rm Y}}(4\mathrm{S})`$. A detailed description can be found in KEKB .
The Belle detector covers a solid angle of almost $`4\pi `$. Closest to the interaction point is a high resolution silicon micro-vertex detector (SVD). It is surrounded by the central drift chamber (CDC). Two dedicated particle identification systems, the aerogel Čerenkov counter (ACC) and the time-of-flight system (TOF), are mounted between the CDC and the CsI(Tl) crystal electromagnetic calorimeter (ECL). All these sub-detectors are located inside a super-conducting coil that provides a magnetic field of 1.5 T. The return yoke of the coil is instrumented as a $`K_L^0`$ and $`\mu `$ detector. A detailed description can be found in BelleDet .
This analysis uses 87.7 $`\mathrm{fb}^1`$ of $`e^+e^{}`$ annihilation data taken at the $`\mathrm{{\rm Y}}(4\mathrm{S})`$ resonance at $`\sqrt{s}=10.58\mathrm{GeV}`$ (“on-resonance data”), above the production threshold for $`B\overline{B}`$ pairs. Additional 15.0 $`\mathrm{fb}^1`$ are taken 60 $`\mathrm{MeV}`$ below the resonance at $`\sqrt{s}=10.52\mathrm{GeV}`$ (“continuum data”), which is also below the production threshold for $`B\overline{B}`$ pairs. Hadronic events are selected as described in HadronB . The selection efficiency of events originating from light quarks ($`d`$, $`u`$ and $`s`$) passing this hadronic preselection has been estimated to be 84.0%, using $`9.6\times 10^6`$ MC events. For $`c`$ quarks, the efficiency has been determined with $`6.6\times 10^6`$ MC events to be 93.0%. The light quark sample contains almost no true candidates, reflecting the small rate for gluon splitting into open charm states, i.e. two mesons containing $`c`$ quarks.
To estimate the efficiency of reconstructing charmed hadrons and to correct for distortions due to the finite acceptance of the detector, MC samples of $`e^+e^{}c\overline{c}`$ events corresponding to a data luminosity of 217 $`\mathrm{fb}^1`$ (approximately 2 1/2 times the on-resonance data), and $`e^+e^{}q\overline{q}`$ ($`q=u`$, $`d`$ and $`s`$) events corresponding to 18 $`\mathrm{fb}^1`$ (approximately 1.2 times the continuum data), have been studied. The MC samples were generated using the QQ98 generator QQ98 employing the Peterson fragmentation function for $`c`$ quarks and were processed through a detailed detector simulation based on GEANT 3.21 GEANT . This sample will be referred to as the generic sample. Special samples of several million $`e^+e^{}c\overline{c}`$ events were generated with the EvtGen EvtGen generator using the Peterson as well as the Bowler fragmentation functions and were also run through the detector simulation. These samples will be referred to as reweighted samples; see Section VI.3 for details about the reweighting procedure. For each charmed hadron used in this analysis, a sample was generated where that hadron was forced to decay in the same channel as later reconstructed. These samples were reconstructed using the same procedures as for data.
### II.1 Particle Identification
To minimise possible kinematic biases due to tight selection criteria for identified particles, only loose cuts on the particle identification of the stable particles have been applied. All particles with mean lifetime longer than $`100\mathrm{ps}`$ have been called “stable”. Apart from reducing a potential kinematic bias, this increased the reconstruction efficiency at the cost of introducing more background, especially in the low momentum region.
In general, the identification for each track was based on one or more likelihood ratios, which combined the information from the time-of-flight and Čerenkov counters and the energy loss dE/dx in the drift chamber. Pions and kaons were separated by a single likelihood ratio $`(K)/((K)+(\pi ))`$. Charged particles were identified as pions if this ratio was less than 0.95 and as kaons if this ratio was larger than 0.05. This overlap allowed a charged particle to be identified as both a pion and a kaon, potentially resulting in identifying a mother (candidate) particle as its own anti-particle (i.e., a $`D^0K^{}\pi ^+`$ decay could be identified as a $`\overline{D}{}_{}{}^{0}\pi ^{}K^+`$ decay), and therefore overestimating the number of candidates. As this misidentification was only possible for neutral particles, an additional systematic uncertainty has been assigned for the $`D^0`$ and $`D^0`$; see section IV for details.
For proton identification, similar likelihood ratios were required to fulfil $`(p)/((p)+(\pi ))>0.6`$ and $`(p)/((p)+(K))>0.6`$. For the $`\pi ^0`$, photon candidates with energies above 30 MeV were combined to form a $`\pi ^0`$ candidate. Under the assumption that the $`\pi ^0`$ candidate decayed at the interaction point, it was required to have an invariant mass consistent with the $`\pi ^0`$ mass.
The efficiencies $`ϵ`$ and misidentification probabilities $`f`$ for tracks from signal candidates under these cuts have been measured in data, and are listed in Table 1; in all cases except the proton, $`ϵ>95\%`$ and $`f26\%`$. For kaons and pions the efficiencies and misidentification probabilities have been estimated in bins of the particle’s momentum from $`D^+`$ and subsequent $`D^0`$$``$$`K^{}`$$`\pi ^+`$ decays; for protons, $`\mathrm{\Lambda }`$ decays have been used. The observed momentum spectra in data have been used to derive the listed numbers.
In addition to the requirements on the particle identification, all tracks had to be consistent with coming from the interaction point (IP). For the slow pion from the $`D^+`$$``$$`D^0`$$`\pi ^+`$ decay, all track quality and particle identification requirements were removed to increase the efficiency.
### II.2 Reconstruction of charmed Hadrons
The reconstructed hadron decay chains used in this analysis are the following:
$`D^0K^{}\pi ^+`$, $`D^+K^{}\pi ^+\pi ^+`$, $`D_s^+\varphi \pi ^+`$ $`(\varphi K^+K^{})`$, $`\mathrm{\Lambda }_c^+p^+K^{}\pi ^+`$,
$`D^+D^0\pi ^+`$ $`(D^0K^{}\pi ^+)`$, $`D^+D^+\pi ^0`$ $`(D^+K^{}\pi ^+\pi ^+)`$ and
$`D^0D^0\pi ^0`$ $`(D^0K^{}\pi ^+)`$.
The inclusion of charge-conjugate modes is implied throughout this paper and natural units are used throughout. For all charmed ground state hadrons, candidates whose masses were within 50 $`\mathrm{MeV}/c^2`$ of their respective nominal mass were considered. For the intermediate $`D^0`$ and $`D^+`$ coming from the excited states $`D^0`$ and $`D^+`$ a mass window of 15 $`\mathrm{MeV}/c^2`$ around the nominal masses of the $`D^0`$ and the $`D^+`$was chosen. Additionally, the selection window for the two excited states was tightened to 15 $`\mathrm{MeV}/c^2`$ around the nominal mass difference between the excited meson and the $`D^0`$ or $`D^+`$. For the intermediate $`\varphi `$ from the $`D_s^+`$ decay, the mass window was chosen to be 7 $`\mathrm{MeV}/c^2`$. Multiple candidates for each particle and anti-particle were removed by a best candidate selection. Most false $`D^0`$ and $`D^+`$ candidates were formed from a true $`D^0`$ and a random slow pion. Therefore, the slow pions were used to determine the best candidate. For the neutral slow pion, the smallest $`\chi ^2`$ of the vertex fit was used. For the charged slow pion, the smallest distance to the IP of all hits used in the reconstruction was used. For all other charmed mesons, the selection was based on the particle identification of the kaon. In the rare case that multiple candidates were formed with the same kaon, the first candidate was randomly chosen.
## III Analysis Procedure
There are two variables commonly used in the measurements of fragmentation functions. These are the scaled energy $`\mathrm{x}_\mathrm{E}=\mathrm{E}_{\mathrm{candidate}}/\mathrm{E}_{\mathrm{candidate}}^{\mathrm{MAX}}`$ and the scaled momentum $`\mathrm{x}_\mathrm{P}=|\stackrel{}{\mathrm{p}}_{\mathrm{candidate}}|/|\stackrel{}{\mathrm{p}}_{\mathrm{candidate}}^{\mathrm{MAX}}|`$, where $`\mathrm{E}_{\mathrm{candidate}}^{\mathrm{MAX}}=\sqrt{s}/2`$, $`|\stackrel{}{\mathrm{p}}_{\mathrm{candidate}}^{\mathrm{MAX}}|=\sqrt{\mathrm{s}/4\mathrm{m}_\mathrm{H}^2}`$, and $`\mathrm{m}_\mathrm{H}`$ denotes the mass of the charmed hadron. For $`b`$ quarks at higher CMEs, the scaled energy $`\mathrm{x}_\mathrm{E}`$ is often used. In this case, the mass of the $`B`$ hadron reduces only slightly the allowed range at small $`\mathrm{x}_\mathrm{E}`$. For charmed hadrons at 10.58 $`\mathrm{GeV}`$ the range of $`\mathrm{x}_\mathrm{E}`$ is significantly reduced; hence $`\mathrm{x}_\mathrm{P}`$ is prefered and will be used in this analysis. Unless otherwise stated, all variables are given in the $`e^+e^{}`$ rest frame, taking into account the different beam energies for the on-resonance and the continuum samples.
For various bins in the range from 0.0 to 1.1 in the scaled momentum $`\mathrm{x}_\mathrm{P}`$, the signal yield has been determined from a fit to the mass or mass difference distributions of all candidates within the aforementioned selection windows. The finite momentum resolution of the detector can result in events being recorded in the region above the naïve limit of $`\mathrm{x}_\mathrm{P}`$=1, however, in the case of $`D^{}`$, the principal contribution is due to the process $`e^+e^{}D^{}D`$. See Section V.3 for details.
A bin width of 0.02 in $`\mathrm{x}_\mathrm{P}`$ has been chosen for all particles as a compromise between the statistical precision in each bin and the momentum resolution, which is a factor of two smaller. Additionally, to investigate the high $`\mathrm{x}_\mathrm{P}`$ region around and above the naïve limit of $`\mathrm{x}_\mathrm{P}`$=1, the bin width has been decreased to 0.01; an expanded view of the region $`0.90<\mathrm{x}_\mathrm{P}<1.05`$ with this binning will be discussed in Section V.3. Since this decreased bin width is still larger than, but comparable to, the momentum resolution, an unfolding using the singular-value-decomposition (SVD) approach SVD was tried in addition to the normal bin-by-bin correction and is discussed in Section V.3.
The mass or mass difference distributions were parametrized by a single Gaussian, except for the $`D^+`$$``$$`D^0`$$`\pi ^+`$ decay channel where a double Gaussian was employed. For the mass distributions, the background was parametrized by a quadratic function; for $`\mathrm{x}_\mathrm{P}>0.9`$ a linear function was found to be sufficient to fit the considerably lower background. For the mass difference distributions of the excited $`D`$ mesons, a phase-space-like function $`f(\mathrm{\Delta }m)=a(\mathrm{\Delta }m\mathrm{\Delta }M_0)^b`$ was used with $`a`$ and $`b`$ being free parameters and $`\mathrm{\Delta }M_0`$ the nominal difference between the mass of the excited mother particle and that of the ground state charm meson.
For all charmed hadrons, the mean mass $`m_i`$ and the width $`\mu _i`$ of the signal Gaussian was fitted separately for MC, continuum and on-resonance data. For these fits, $`\mathrm{x}_\mathrm{P}`$ was divided into 4 bins from $`0.2<\mathrm{x}_\mathrm{P}<1.0`$ with a constant bin size of 0.2. In a second fit, two quadratic functions $`m_i(\mathrm{x}_\mathrm{P})`$ and $`\mu _i(\mathrm{x}_\mathrm{P})`$ were fitted to the results of the first fit in these four bins. For the distributions with a bin width of 0.02 and 0.01 in $`\mathrm{x}_\mathrm{P}`$, the mean and width parameters in the fit were fixed to the values of the quadratic functions $`m_i(\mathrm{x}_\mathrm{P})`$ and $`\mu _i(\mathrm{x}_\mathrm{P})`$ for the appropriate $`\mathrm{x}_\mathrm{P}`$ value.
For the $`D^+`$$``$$`D^0`$$`\pi ^+`$ decay mode full correlations between the two Gaussians of the signal function were taken into account when determining the fit yield.
When combining the on-resonance data with the continuum data, two corrections have been applied to the on-resonance data. After normalising using the integrated luminosities of the respective samples, the naïve $`1/s`$ dependence on the total hadronic cross section has been taken out by multiplying the distributions of the on-resonance sample by the square of the ratio of the CME’s, namely by $`(10.58\mathrm{GeV}/10.52\mathrm{GeV})^2`$. Second, from MC an additional correction of $`+0.27\%`$ due to different initial state radiation (ISR) at the two energy points has been applied to the on-resonance samples. This correction was based on a MC study of the total cross sections at these two energy points.
### III.1 $`\mathrm{x}_\mathrm{P}`$-dependent Mass Fits
Fig. 1 and Fig. 2 show the mass distributions of all charmed hadrons reconstructed in this analysis for two representive bins in $`\mathrm{x}_\mathrm{P}`$. The $`\mathrm{x}_\mathrm{P}`$ bins shown are $`0.28<\mathrm{x}_\mathrm{P}<0.30`$ in Fig. 1 and $`0.68<\mathrm{x}_\mathrm{P}<0.70`$ in Fig. 2. They represent a low $`\mathrm{x}_\mathrm{P}`$ bin with higher background and a bin close to the maximum of the $`\mathrm{x}_\mathrm{P}`$ distribution with less background, respectively. All mass (mass difference) distributions have been shifted by their nominal mass (mass difference) to center the peaks at zero. See table 2 for the masses used. Note that the scale on the $`y`$-axis does not start at zero in the upper four plots in Fig. 1 and Fig. 2.
For $`0.28<\mathrm{x}_\mathrm{P}<0.30`$ (shown in Fig. 1), the mass distributions for the $`D^0`$ and $`D^+`$ ground states and the mass difference distributions for the excited states show clear peaks at the expected value for signal. Compared to higher $`\mathrm{x}_\mathrm{P}`$ values, the background is higher due to a larger amount of combinatorial background, and the signal-to-background ratio is lower. At higher $`\mathrm{x}_\mathrm{P}`$ values, such as those shown in Fig. 2 $`(0.68<\mathrm{x}_\mathrm{P}<0.70)`$, the background is considerably reduced, whereas the signal yield is enhanced. This significantly increased the signal-to-background ratio.
### III.2 Raw Signal Yield
Fig. 3 shows the signal yields as a function of $`\mathrm{x}_\mathrm{P}`$ for all charmed hadrons, not corrected for the reconstruction efficiencies and for the branching fractions, denoted with “B” in the plots. For all particles, the contribution from $`B`$ decays is clearly visible in the low $`\mathrm{x}_\mathrm{P}`$ range, which is $`\mathrm{x}_\mathrm{P}<0.5`$ for all charmed mesons containing a light quark as the spectator. For $`D_s^+`$ from $`B`$ decays, the upper bound is approximately $`\mathrm{x}_\mathrm{P}0.4`$, reflecting the energy required to produce an additional strange quark. Contributions from the $`bu`$ transition, where the $`D_s^+`$ is formed at the upper vertex, can populate the region up to $`\mathrm{x}_\mathrm{P}=0.5`$, but are strongly suppressed. For the $`\mathrm{\Lambda }_c^+`$, the only baryon reconstructed in this analysis, the upper bound is further decreased to approximately $`\mathrm{x}_\mathrm{P}0.37`$, due to the production of an additional anti-baryon.
All distributions peak around $`\mathrm{x}_\mathrm{P}0.60.7`$ and show similar shapes.
### III.3 Efficiency Correction
The efficiencies were determined from MC and are defined as the appropriate raw signal yield (determined by the same procedure as for data) divided by the generated MC $`\mathrm{x}_\mathrm{P}`$ distribution. The seven histograms in Fig. 4 show the $`\mathrm{x}_\mathrm{P}`$-dependent efficiency of each charmed hadron used in this analysis for continuum data and on-resonance data. The $`D^0`$ efficiency is close to 50% and almost constant over the entire $`\mathrm{x}_\mathrm{P}`$ range. The efficiency for $`D^+`$$``$$`D^0`$$`\pi ^+`$ approaches the $`D^0`$ efficiency at high values of $`\mathrm{x}_\mathrm{P}`$ and diminishes at lower values of $`\mathrm{x}_\mathrm{P}`$, reflecting the reduced efficiency of reconstructing low-momentum pions. The two $`D^{}`$ decay modes that include a neutral slow pion show a different behaviour: the efficiencies stay constant over a wide range of about $`0.3<\mathrm{x}_\mathrm{P}<1`$ and below $`0.3<\mathrm{x}_\mathrm{P}`$ the efficiency increases for $`\mathrm{x}_\mathrm{P}`$$``$0 due to the increasing reconstruction efficiency for slow $`\pi ^0`$. The reconstruction efficiencies for the three-particle decay modes do not show a strong dependence on $`\mathrm{x}_\mathrm{P}`$, slightly varying between 15% and 20% for the $`D_s^+`$ and remaining constant at about 30% for the $`\mathrm{\Lambda }_c^+`$. The decreasing efficiency for particles at values close to the kinematic limit is an artefact of the decreasing statistics in all generic MC samples. The reweighted samples, which were generated with a different fragmentation function than the generic samples, contain significantly more events in the very high $`\mathrm{x}_\mathrm{P}`$ region and do not show such behaviour. This difference between the two efficiency estimates was added to the systematic uncertainty.
The efficiency is a function of the production angle, which differs for charmed hadrons from $`B`$ decays and from continuum events. For the on-resonance samples, the efficiency has been determined by a luminosity-weighted mixture of charmed MC and dedicated samples containing decays of charged and neutral $`B`$ mesons. For the continuum sample, only charmed MC was used.
In data, it was verified that $`D^+`$ produced in $`e^+e^{}`$ annihilation are unpolarised by verifying that the distribution of the cosine of the helicity angle is flat. The helicity angle is defined as the angle between the slow charged pion in the $`D^+`$ rest frame and the flight direction of the $`D^+`$ in the center of mass system of the event. Because the efficiency for $`D^+`$$``$$`D^0`$$`\pi ^+`$ strongly depends on the momentum distribution of the slow $`\pi ^+`$, which in turn depends on the helicity angle, polarised $`D^+`$ can introduce a bias into the efficiency correction.
## IV Systematic Uncertainties
Various sources of systematic uncertainties have been considered:
Uncertainties due to tracking were estimated to be 1% per track using a sample of partially reconstructed $`D^+`$ decays. As the uncertainty increased at very low momentum, the estimated momentum-dependent uncertainty of the slow charged pion was folded with the observed momentum spectrum. The systematic uncertainty due to the slow neutral pion detection efficiency was assessed by examining the differences in the shapes of the fragmentation function of the two $`D^+`$ decay modes, $`D^+D^0\pi ^+`$ and $`D^+D^+\pi ^0`$.
Uncertainties due to the modeling of ISR in the MC were determined by restricting the longitudinal momentum in the laboratory frame of all candidates to $`p_z^{lab}>0`$ only. This cut preferentially removed events with ISR photons in the negative $`z`$ direction, potentially introducing an artificial asymmetry. The $`z`$ direction is defined as being anti-parallel to the positron beam, which coincides up to corrections due to the crossing angle with the boost vector into the $`e^+e^{}`$ rest-frame.
The cut on the likelihood ratios for kaon and proton candidates was tightened to 0.2 and 0.8, respectively, and the difference was taken into account in the systematic uncertainty.
Potential differences between the actual signal shape and the fitting function were estimated by determining the signal yield with a counting method instead of using the fit. Here, the number of entries in the mass (mass difference) distribution was counted in a window about one third the size of the total 50 $`\mathrm{MeV}/c^2`$ (15 $`\mathrm{MeV}/c^2`$) window around the peak position, corresponding to roughly three times the resolution. The number of background events was subtracted after integrating the background function of the standard fit within the same mass window.
An additional flavour assignment systematic uncertainty was taken into account for the neutral states $`D^0`$ and $`D^0`$. The loose cuts on the charged pion and kaon particle identification allowed a $`D^0`$ to be identified as a $`\overline{D}^0`$: the flavour of the $`D^0`$ from $`D^+`$ decays was identified by the charge of the slow pion (except for a small contribution from doubly-Cabibbo-suppressed decays). In the MC sample, the likelihood ratio of the pion candidate was larger than that of the kaon candidate for 1.3% of all $`D^0`$ candidates; the corresponding fraction was determined to be 1.1% for $`D^+`$ decays. The statistical uncertainties on these numbers are less than 0.05%. Accordingly, a difference of 0.2% was assigned as the uncertainty of the flavour assignment due to the overlap of the pion and kaon likelihoods of the particle identification.
The luminosity of the data sample was determined to have an uncertainty of about 1.4%. A corresponding scale uncertainty of 1.4% was assigned to the normalisation of the shape. It has been checked that the normalisation of the fragmentation functions of the on-resonance and continuum sample agree with each other, and their difference of 0.94% is well within the scale uncertainty.
Finally, the reconstruction efficiencies of the generic and the reweighted samples differed slightly. This small difference was added to the systematic uncertainty.
All systematic uncertainties were added in quadrature to give the total systematic uncertainty.
## V Results
In this section, various results for the charmed hadrons are presented.
### V.1 $`\mathrm{x}_\mathrm{P}`$ Distributions
Fig. 5 shows the efficiency-corrected $`\mathrm{x}_\mathrm{P}`$ distributions for the different particles for $`e^+e^{}`$ annihilation events, i.e. spectra of hadrons formed in the fragmentation of charm quarks. Above $`\mathrm{x}_\mathrm{P}>0.5`$, the differential $`\mathrm{x}_\mathrm{P}`$ distributions of the on-resonance sample and the continuum sample have been combined by a weighted average, where the inverse of the squared statistical uncertainty was used as the weight. As the systematic uncertainties for both samples are highly correlated, the larger uncertainty of the on-resonance and the continuum samples was used for the combined sample.
As most efficiencies do not depend strongly upon $`\mathrm{x}_\mathrm{P}`$, the shapes of the efficiency-corrected distributions are similar to those of the uncorrected distributions. All distributions peak around $`\mathrm{x}_\mathrm{P}0.60.7`$. To determine the peak position, a direct fit of the data to the Peterson fragmentation function was tried. The shape of the data agreed very poorly as this model does not include gluon radiation or decays from higher resonances. Therefore, a Gaussian function was used to determine the peak position. The fit ranges were chosen from $`\mathrm{x}_\mathrm{P}=0.4`$$`0.8`$. The results of the fits are listed in Table 3, together with the statistical and systematic uncertainties. The statistical uncertainty was determined by the RMS of the distribution divided by $`\sqrt{N}`$.
### V.2 Average Number of Charmed Hadrons per $`B`$ Decay
The $`\mathrm{x}_\mathrm{P}`$ distributions of the on-resonance and continuum samples differ in the contribution from $`B`$ decays for $`\mathrm{x}_\mathrm{P}<0.5`$. Fig. 6 shows this difference: the differential $`\mathrm{x}_\mathrm{P}`$ distribution of the continuum sample was subtracted from that of the on-resonance sample. Thus, up to statistical fluctuations it contains only contributions from decays of $`B`$ mesons.
Table 4 lists the average number of charmed hadrons per $`B`$ meson decay together with the present world average PDG2004 . In order to determine the average number of charmed hadrons produced per $`B`$ decay, we take the difference between the production rate in the on-resonance and the continuum sample and normalise by the $`B`$ meson production cross section, which is estimated to be $`(1.073\pm 0.019)\mathrm{nb}`$ based on the measured luminosity and the measured number of $`B\overline{B}`$ pairs in this sample. Note that this visible production cross-section depends strongly upon the energy spread of the accelerator. The uncertainties in Table 4 are from the limited statistics (first), the systematics as discussed in Section IV (second), and the luminosity measurement and the uncertainties on the branching fractions (third). Note that the luminosity measurement and the determination of the number of $`B\overline{B}`$ are strongly correlated. Both values agree well within one standard deviation with each other, only the average number of produced $`D^0`$’s here is lower by about one standard deviation and is closer to that of $`D^+`$’s.
The small bump seen in the $`\mathrm{x}_\mathrm{P}`$ distributions of the charmed mesons except the $`D_s^+`$ at $`\mathrm{x}_\mathrm{P}=0.35`$ is due to two body decays of the $`B`$ mesons such as $`BD^0D^{()}`$ in case of the $`D^0`$.
### V.3 High $`\mathrm{x}_\mathrm{P}`$ Region
An expanded view of the high $`\mathrm{x}_\mathrm{P}`$ region is shown in Fig. 7. The downward triangles show the efficiency-corrected data; the upward triangles show the corrected and unfolded data.
Unfolding was done using the singular-value-decomposition (SVD) method SVD . From MC, we determined the response matrix of the detector for producing for a certain true input value of $`\mathrm{x}_{\mathrm{P}}^{}{}_{,true}{}^{}`$ a measured value of $`\mathrm{x}_{\mathrm{P}}^{}{}_{,measured}{}^{}`$. This matrix was decomposed using the SVD into two orthogonal and one diagonal matrices which can easily be inverted. Inverting the diagonal matrix was limited by a criteria defined in SVD to contain only elements, which are of statistical significance.
The hatched histogram show the only process $`e^+e^{}D^+D^{}`$, the open histogram shows the sum of the previous process and $`e^+e^{}D^+D^{()}`$.
The $`\mathrm{x}_\mathrm{P}`$ distributions for the ground states $`D^0`$, $`D^+`$, $`D_s^+`$ and $`\mathrm{\Lambda }_c^+`$ extend up to the naïve kinematic endpoint $`\mathrm{x}_\mathrm{P}=1`$ and no significant number of events are present for $`\mathrm{x}_\mathrm{P}>1`$.
All three $`\mathrm{x}_\mathrm{P}`$ distributions for the excited $`D`$ mesons, however, show an enhancement at $`\mathrm{x}_\mathrm{P}>1`$. These events above $`\mathrm{x}_\mathrm{P}=1`$ correspond to events of the processes $`e^+e^{}D^+D^{()}`$ or $`e^+e^{}D^0D^{()0}`$ and are in good agreement with the measured cross sections Uglov of $`0.55\pm 0.03\pm 0.05\mathrm{pb}`$ for $`e^+e^{}D^+D^{}`$ and $`0.62\pm 0.03\pm 0.06\mathrm{pb}`$ for $`e^+e^{}D^+D^{}`$. Note that these events populate $`\mathrm{x}_\mathrm{P}>1`$ only because of the use of the simplified upper limit $`|\stackrel{}{\mathrm{p}}_{\mathrm{candidate}}^{\mathrm{MAX}}|`$, for producing two $`\mathrm{D}^{}`$ mesons. A background fluctuation producing an artificial peak is unlikely for two reasons. First, at high $`\mathrm{x}_\mathrm{P}`$, the background is negligible, and second, the unfolding procedure tends to identify signals at the edge of a distribution as statistical fluctuations rather than real signals, thus decreasing the significance of the signals.
### V.4 Total Production Cross-Section
The total production cross-section is given by the integral of the $`\mathrm{x}_\mathrm{P}`$ distribution. This integral was determined for the continuum sample using the current value of the world average product branching fraction of each particle, see Table 2 and PDG2004 . The results are listed in Table 5, where the third error component reflects the uncertainty on the product branching fraction.
The results by CLEO CLEOprel given in the last column used their own branching fractions, which differ slightly from the world averages used here. The results, however, agree well with each other. Another measurement by BaBar BaBar\_Ds is given in the same column. The total production cross-section for the $`D^0`$ differs only slightly from that of the $`D^+`$. This can be understood as resulting from different feed-down contributions from higher resonances.
### V.5 Mean Values for $`\mathrm{x}_\mathrm{P}`$ and Moments
In addition to the peak position for the seven $`\mathrm{x}_\mathrm{P}`$ distributions, the mean and higher moments of these distributions were determined from distributions in $`(\mathrm{x}_\mathrm{P})^n`$ with a bin width of 0.02 in $`\mathrm{x}_\mathrm{P}`$ and a bin-by-bin efficiency correction was applied. The $`n^{th}`$ moment was determined by the mean of the efficiency corrected distributions in $`(\mathrm{x}_\mathrm{P})^n`$, and its statistical uncertainty was determined by $`\sigma /\sqrt{N_0}`$, where $`N_0`$ is the number of entries in the uncorrected $`(\mathrm{x}_\mathrm{P})^n`$ distribution. Tables 6 and 7 show the moments for the different decay modes.
### V.6 Production Angle
Taking the interference between the exchange of virtual photons and $`\mathrm{Z}`$ bosons into account, the differential cross section for $`e^+e^{}c\overline{c}`$ is modified from the $`1+\mathrm{cos}^2\theta `$ form of the Born amplitude for pure photon exchange:
$$\frac{d\sigma }{d\mathrm{cos}\theta }=\frac{3}{8}(1+\mathrm{cos}^2\theta )\sigma _T+\frac{3}{4}\mathrm{sin}^2\theta \sigma _L+\frac{3}{4}\mathrm{cos}\theta \sigma _A$$
(1)
Here, $`\theta `$ describes the angle between the incoming electron beam and the outgoing hadron containing the charmed quark, as measured in the CM frame. The term $`\sigma _T`$ describes the contribution of pair production of spin-1/2 particles from transverse polarised vector bosons, the term $`\sigma _L`$ the contribution from longitudinal polarised vector bosons and the term $`\sigma _A`$ denotes the parity violating asymmetry due to the interference between $`\mathrm{Z}`$ bosons and virtual photons.
The $`𝒦𝒦`$ MC generator KK was used to predict the production angle distributions for the different charmed hadrons. This MC generator includes interference between initial and final state radiation (ISR and FSR) as well as electro-weak corrections. $`10^8`$ $`c\overline{c}`$ events were generated with $`𝒦𝒦`$ and hadronised with the PYTHIA generator.
For the generated events, $`\mathrm{x}_\mathrm{P}`$ was divided into 20 bins of equal width and a three-parameter fit to the production angle was performed:
$$f(\theta ,\mathrm{x}_\mathrm{P})=a_0(\mathrm{x}_\mathrm{P})\left[\frac{3}{8}(1+\mathrm{cos}^2\theta )+\frac{3}{4}r_S(\mathrm{x}_\mathrm{P})\mathrm{sin}^2\theta +\frac{3}{4}r_C(\mathrm{x}_\mathrm{P})\mathrm{cos}\theta \right]$$
(2)
where $`a_0`$ is the normalisation and $`r_S`$ and $`r_C`$ are the relative contributions for the $`\mathrm{sin}^2\theta `$ and the $`\mathrm{cos}\theta `$ terms, respectively. The results of these fits are shown in Fig. 8 as the solid line, together with the measured data points.
For data, the signal yield was determined in bins of $`\mathrm{x}_\mathrm{P}`$ and $`\mathrm{cos}\theta `$, where $`\theta `$ is the production angle of the charmed hadron. It should be noted that here, the efficiency correction depends on $`\mathrm{x}_\mathrm{P}`$ and $`\mathrm{cos}\theta `$.
$`\mathrm{x}_\mathrm{P}`$ was divided into only 4 bins: $`\mathrm{x}_\mathrm{P}<0.3`$, $`0.3<\mathrm{x}_\mathrm{P}<0.5`$, $`0.5<\mathrm{x}_\mathrm{P}<0.7`$ and $`0.7<\mathrm{x}_\mathrm{P}`$. The boundaries were chosen in order to roughly equalise the number of candidates per bin. Two bins each were chosen below and above $`\mathrm{x}_\mathrm{P}=0.5`$, which is the upper kinematic limit for hadrons from $`B`$ decays.
$`\mathrm{cos}\theta `$ was divided into 20 bins. In each bin of $`\mathrm{x}_\mathrm{P}`$ and $`\mathrm{cos}\theta `$, the efficiency corrected signal yield in the mass or mass difference distributions was fitted separately for the on-resonance and continuum samples. The same signal and background functions were used as in the fits which depended only on $`\mathrm{x}_\mathrm{P}`$.
In a second step, a three-parameter fit (similar to Eq. 2) to the signal yields was performed in bins of $`\mathrm{x}_\mathrm{P}`$. The fit range was restricted to $`0.8<\mathrm{cos}\theta <0.8`$. As a systematic check, the range was tightened to $`0.7<\mathrm{cos}\theta <0.7`$.
No significant deviation of the $`r_S`$ and the $`r_C`$ parameters from their expectation was found for the continuum sample. The expectation from the $`𝒦𝒦`$ generator was of the same order as the statistical uncertainties except for the $`r_S`$ term in the lowest $`\mathrm{x}_\mathrm{P}`$ bin, where gluon radiation introduces a longitudinal momentum component and therefore smears out the initial distribution of the production angle. This smearing introduces a significant $`\mathrm{sin}^2\theta `$ term, to which the measured values in this regime agree well. As the number of entries in the low $`\mathrm{x}_\mathrm{P}`$ bins also diminish, the statistical uncertainties increase to roughly the same size as the expected effect. The fitted values for $`r_S`$ and $`r_C`$ in the lowest $`\mathrm{x}_\mathrm{P}`$ bin in the continuum samples suffer from very low statistics (see Fig. 3). For the $`D_s^+`$ and the $`\mathrm{\Lambda }_c^+`$, the lowest bins in the continuum samples have been neglected as the numbers of entries in these bins were too low to perform an angular analysis.
For the on-resonance sample with higher statistics, the $`r_S`$ terms significantly deviate from zero for $`\mathrm{x}_\mathrm{P}<0.5`$, as $`B`$ decays with different production angle distributions also contribute. The $`r_C`$ term is again consistent with both zero and the expectation from the MC generator. Tables 8 and 9 list the measured values for the $`r_S`$ and $`r_C`$ values for the continuum and the on-resonance data. The systematic uncertainties here include the uncertainties discussed in the standard analysis as well as the uncertainty due to the restricted fit range.
## VI Interpretation
### VI.1 Contributions from Higher Resonances
Contributions from excited states have been considered only in the $`\mathrm{x}_\mathrm{P}`$ distributions of the $`D^0`$ and the $`D^+`$ and for these, only contributions from $`D^0`$ and $`D^+`$ were considered. For higher resonances, both production cross-section and branching fractions of e.g. $`D^{}`$ have large uncertainties and have been neglected. In order to reduce the statistical uncertainty, a MC-based correction was applied: Three large samples of several million MC events were generated without a detailed detector simulation. These samples were required to contain the decay modes $`D^+`$$``$$`D^0`$$`\pi ^+`$, $`D^+`$$``$$`D^+`$$`\pi ^0`$ and $`D^0`$$``$$`D^0`$$`\pi ^0`$, respectively. For these events, the $`\mathrm{x}_\mathrm{P}`$ of the parent $`D^+`$/$`D^0`$ vs. the $`\mathrm{x}_\mathrm{P}`$ of the daughter $`D^0`$/$`D^+`$ were stored in a two dimensional matrix. The measured and efficiency corrected $`D^0`$ and $`D^+`$ $`\mathrm{x}_\mathrm{P}`$ distributions were multiplied with this matrix in order to estimate the $`\mathrm{x}_\mathrm{P}`$ distribution for all $`D^0`$’s and $`D^+`$’s coming from $`D^+`$/$`D^0`$ decays.
The two plots at the top of Fig. 9 show the contributions of $`D^+`$ and $`D^0`$ decays to the $`D^0`$ fragmentation function (left), and the contribution of $`D^+`$ decays to the $`D^+`$ fragmentation function (right). These plots are not corrected for the branching fraction of the $`D`$ decay. The bottom plot in Fig. 9 shows primary $`D^0`$ and $`D^+`$ fragmentation spectra: the total $`D^0`$ ($`D^+`$) spectrum minus the contribution from $`D^0`$ and $`D^+`$ ($`D^+`$only) decays. The difference of 13% between the sum of primary produced $`D^0`$ and $`D^+`$ should be compared to the 6.5% relative uncertainty in the $`D^+`$$``$$`K^{}`$$`\pi ^+`$$`\pi ^+`$ branching fraction. Also, as only the contribution from $`D^{}`$ decays has been considered, the remaining difference may be due to the contribution of other resonances. In the generic MC sample, where primary $`D^0`$’s and $`D^+`$’s were produced in equal amounts, there was an excess of 6% in the production of $`D^0`$ (compared to $`D^+`$) mesons in the decay of resonances other than $`D^+`$.
### VI.2 Ratios
Comparisons of production rates for various particles are useful for understanding the dynamics of fragmentation, as systematic errors cancel in the ratio. In this section we present ratios of both integrated cross-sections and cross-sections as a function of $`\mathrm{x}_\mathrm{P}`$, to characterise general properties of fragmentation and to test the agreement between MC simulation and data.
Table 10 presents three ratios of total production cross sections. Since the production of $`D^+`$ and $`D^0`$ is included in the total $`D^+`$ and $`D^0`$ production rate (all $`D^+`$ decay to either $`D^+`$ or $`D^0`$, and all $`D^0`$ to $`D^0`$), the ratio of $`D^{}`$ to $`D`$ production measures $`V/(V+P)`$, the probability of producing a vector charmed meson. (Here we write $`V`$ for vector and $`P`$ for pseudo-scalar meson production rates.) A correction is necessary to account for higher resonances decaying directly to $`D^{+,0}`$. For example, based on the measured production rates of the $`D_1(2420)`$ and $`D_2^{}(2460)`$ CLEO\_Dstst mesons, known branching fractions PDG2004 , and isospin relations, we find a correction of $`(3.7\pm 3.3)\%`$ to the first ratio. In principle, further corrections due to decays of broad $`D^{}`$ states and charmed-strange mesons are also required. However, no corrections have been applied to the values presented in Table 10.
Similarly, the second ratio measures the production rate of charmed-strange mesons as a fraction of all charmed mesons, up to corrections for $`D_{s1}(2536)`$ and $`D_{s2}^{}(2573)`$ decays. The third ratio measures the production rate of charmed baryons relative to that of charmed mesons, excluding the charmed-strange baryon states. For comparison, see Gladilin .
Ratios of production rates as a function of $`\mathrm{x}_\mathrm{P}`$ allow momentum-dependent effects in fragmentation to be studied, although contributions from decays of higher states also appear. Fig. 10 shows the following five ratios as a function of $`\mathrm{x}_\mathrm{P}`$, for both on-resonance and continuum data:
* $`\mathrm{x}_\mathrm{P}(D^+)/\mathrm{x}_\mathrm{P}(D_{prim}^+)`$, sensitive to the production rate of vector relative to pseudo-scalar mesons;
* $`\mathrm{x}_\mathrm{P}(D_{prim}^0)/\mathrm{x}_\mathrm{P}(D_{prim}^+)`$, sensitive to charged relative to neutral pseudo-scalar production;
* $`\mathrm{x}_\mathrm{P}(D_s^+)/\mathrm{x}_\mathrm{P}(D_{prim}^+)`$, sensitive to the production of strange quarks;
* $`\mathrm{x}_\mathrm{P}(\mathrm{\Lambda }_c^+)/\mathrm{x}_\mathrm{P}(D_{prim}^+)`$, sensitive to the production of baryons relative to mesons;
* $`\mathrm{x}_\mathrm{P}(D^0)/\mathrm{x}_\mathrm{P}(D^+)`$, the relative production rate of the vector mesons.
The suffix “prim” denotes $`\mathrm{x}_\mathrm{P}`$ distributions corrected for the contributions from $`D^{}`$ decays; $`D^+`$ production has been chosen as the denominator in (a)–(d), as this correction is smaller than that for $`D^0`$. No other corrections for excited states have been applied.
The production ratios in Fig. 10(a) and (b) are similar for on-resonance and continuum data. In Fig. 10(c), the contribution of $`B`$ meson decays to $`D_s^+`$ production can be clearly seen in the low-$`\mathrm{x}_\mathrm{P}`$ region. In Fig. 10(d), baryon production in $`B`$ decays is seen to be suppressed below $`\mathrm{x}_\mathrm{P}0.4`$. As $`\mathrm{x}_\mathrm{P}`$ approaches unity, the $`\mathrm{\Lambda }_c^+/D^+`$ production ratio goes to zero, consistent with the conservation of baryon number.
Four similar ratios are shown in Fig. 11(a)–(d) for both continuum data (full squares) and MC simulations, to test the performance of the MC for various fragmentation function parameters. In these plots, the total $`D^+`$ production rate, without $`D^{}`$ subtraction, is used in the denominator of the ratios. The open histograms show the generic MC sample, which agrees with the data only for the highest values of $`\mathrm{x}_\mathrm{P}`$ of the distributions in Fig. 11(a) and (b), but fails to describe the data distributions at lower values. The open squares show a second MC sample generated with the Bowler fragmentation function, which shows a similar behaviour.
Noting that the parameter PARJ(13) in PYTHIA gives the probability for a charmed hadron produced in fragmentation to have spin one, 50 MC samples of $`10^7`$ events each were generated, with PARJ(13) values ranging from $`0.25`$ to the default value $`0.75`$ given by spin counting. These samples were generated in addition to the “reweighted samples” used for more refined MC comparisons as described in the next section. A reduced chi-squared $`\stackrel{~}{\chi }^2`$ was calculated for these samples and the measured and corrected ratios, and a fourth-order polynomial in PARJ(13) was fitted to the results. The minimum $`\stackrel{~}{\chi }^2`$ was found to occur at $`\text{PARJ(13)}=0.592\pm 0.021`$ for the $`\mathrm{x}_\mathrm{P}(D^+)/\mathrm{x}_\mathrm{P}(D^+)`$ ratio, and $`0.592\pm 0.046`$ for the $`\mathrm{x}_\mathrm{P}(D^0)/\mathrm{x}_\mathrm{P}(D^+)`$ ratio, where the uncertainties denote the $`1\sigma `$ range in the fitted polynomial. We note that models of hadron production more sophisticated than spin counting predict values for PARJ(13) below $`0.75`$; see VoverPV and references therein.
In Fig. 11(a)–(d) a third MC sample generated with the Bowler fragmentation function, and $`\text{PARJ(13)}=0.59`$, is shown with closed triangles. This sample and the data agree well within the error bars over almost the entire range in Fig. 11(a) and (b). All three MC samples fail to describe the ratios in Fig. 11(c) and (d): both the endpoints and the shape disagree. The difference between the MC samples is small compared to the discrepancy with data for $`D_s^+/D^+`$ production in Fig. 11(c); while Bowler fragmentation (open squares and triangles) gives an improved description of $`\mathrm{\Lambda }_c^+/D^+`$ production in Fig. 11(d), the agreement is still poor. There are no obvious parameters in the MC which can affect these two ratios in such a way as to improve the agreement between data and MC.
### VI.3 Comparison of $`\mathrm{x}_\mathrm{P}`$ distributions with predictions from MC generators
The models used by MC generators are based on simplified assumptions and require input from experiment: this is reflected in the models’ input parameters.
The commonly used JETSET/PYTHIA generators are based on the Lund or string model, in which a coloured string is expanded between two emerging partons. The energy stored in the string increases with increasing distance, and eventually allows the production of a new quark anti-quark pair. The quark (anti-quark) then produces a meson together with the initial anti-quark (quark). The energy distribution of the new quark or anti-quark is described by a fragmentation function. Various fragmentation function models have been published; see Table 11 for a summary. These models depend upon up to two independent variables, these are the transverse mass $`m_{}=\sqrt{m^2+p_{}^2}`$ of the newly created hadron, and $`z`$, the fraction of the longitudinal energy $`E+p_{}`$ which the meson inherits from the initial quark.
Not all models listed in Table 11 are implemented in the JETSET/PYTHIA generator. In order to be able to compare all models to data, a reweighting technique has been applied. Here, several million events referred to as “reweighted samples” have been generated, allowing a more elaborate comparison than described in Section VI.2. For these events, $`z`$ and $`p_{}`$ were stored together with the event. This allowed each event to be reweighted in order to mimic any other fragmentation function. Scans through the parameter space of the five listed fragmentation functions have been performed on these special samples. This analysis was performed on five different hadrons; $`D^0`$$``$$`D^0`$$`\pi ^0`$ and $`D^+`$$``$$`D^+`$$`\pi ^0`$ have been omitted because of the large systematic uncertainty due to the detection efficiency of the slow neutral pion.
For data and MC, the $`\mathrm{x}_\mathrm{P}`$ distributions were compared using uncorrected data and the reweighted special MC samples after full detector simulation. A $`\chi ^2`$ was calculated based upon the distribution of the reweighted special sample and the measured data distribution. Only statistical uncertainties in each $`\mathrm{x}_\mathrm{P}`$ bin were taken into account and only bins which contained entries in data or MC were included. The number of bins minus the number of parameters of the fragmentation function was used as the number of degrees of freedom ($`d.o.f.`$). The weights in the reweighting procedure were constructed in such a way that the number of events before and after reweighting stayed constant. This way the total value of the $`\chi ^2`$ becomes dependent on the size of the data and MC samples; the relative $`\chi ^2`$ values, however, allow a direct comparison between the different fragmentation functions.
Table 12 shows the $`\chi _{min}^2/d.o.f.`$ for all five particles and five fragmentation functions. For all five particles a similar trend is visible. The Bowler model in general agrees best with the data. The Lund models shows a similar performance in describing the spectra, its $`\chi _{min}^2/d.o.f.`$ being by factors of 2–3 better than the next best model. For $`D^+`$ and $`D^+`$ the $`\chi _{min}^2/d.o.f.`$ is slightly worse than for the Bowler model. In the minimum of the $`\chi ^2/d.o.f.`$ distributions the $`a`$ parameter deviates strongly from the default for most of the particles. As the Lund model is employed for fragmentation of all flavour species, such a large change in the parameter would also change the particle spectrum of light mesons. Therefore, further tuning of the second parameter to the Lund fragmentation function has been omitted.
The models by Collins and Spiller and by Kartvelishvili show a similar $`\chi _{min}^2`$ for all particles, about factors of two to three worse than that of the best models. The last model, that of Peterson, shows the worst agreement with a reduced $`\chi _{min}^2`$ of 15 and well above, ruling out this model for describing data at this CME.
The input parameters for the fragmentation functions at the minimum of the $`\chi ^2/d.o.f.`$ distributions are listed in Table 13.
In summary, the Bowler model shows the best agreement between data and MC, however, large differences are still present. These differences might be resolved by adjusting other parameters of the generators as well, but such a task is out of scope for this analysis. The Lund model shows the second best agreement. The models by Kartvelishvili and by Collins and Spiller show larger deviations and the commonly used model by Peterson shows the worst agreement between data and MC.
## VII Summary
A new determination of the charm fragmentation function at a CME close to the $`\mathrm{{\rm Y}}(4\mathrm{S})`$ resonance has been presented. The measured $`\mathrm{x}_\mathrm{P}`$ spectra have been compared to those of five different parametrisations in MC via a reweighting procedure, and the best input parameters have been found. The best agreement between data and MC has been found for the Bowler model and the Lund model. Additionally, the peak positions and the first six moments of the $`\mathrm{x}_\mathrm{P}`$ distributions have been measured. These measurements will allow detailed comparisons between experiment and theory. The total production cross-section, as well as $`\mathrm{x}_\mathrm{P}`$ dependent ratios of the fragmentation functions, place stringent tests on existing MC generators, which so far completely fail to describe the $`\mathrm{x}_\mathrm{P}`$ dependent ratios of $`\mathrm{x}_\mathrm{P}(D_s^+)/\mathrm{x}_\mathrm{P}(D^+)`$ and $`\mathrm{x}_\mathrm{P}(\mathrm{\Lambda }_c^+)/\mathrm{x}_\mathrm{P}(D^+)`$. For the first time, the production rates of $`D^+`$ and $`D^0`$ excluding the decay of $`D^{}`$ mesons have been measured. They were found to agree reasonably well with each other.
The efficiency corrected data points will be made available via download in the Durham HEP REACTION DATA DataBase DurhamHEP . It is presented in a different way as shown this article. Separate sets of the continuum and the on-resonance samples are given as $`sBd\sigma /d\mathrm{x}_\mathrm{P}`$, i.e. scaled by the nominal center-of-mass energies of 10.52 $`\mathrm{GeV}`$ and 10.58 $`\mathrm{GeV}`$, respectively, and not corrected for the branching ratios. The on-resonance data includes the additional correction for ISR of $`+0.27\%`$, see Section III for details.
## Acknowledgements
We thank O. Biebel and A. Mitov for discussion. We thank the KEKB group for the excellent operation of the accelerator, the KEK cryogenics group for the efficient operation of the solenoid, and the KEK computer group and the National Institute of Informatics for valuable computing and Super-SINET network support. We acknowledge support from the Ministry of Education, Culture, Sports, Science, and Technology of Japan and the Japan Society for the Promotion of Science; the Australian Research Council and the Australian Department of Education, Science and Training; the National Science Foundation of China under contract No. 10175071; the Department of Science and Technology of India; the BK21 program of the Ministry of Education of Korea and the CHEP SRC program of the Korea Science and Engineering Foundation; the Polish State Committee for Scientific Research under contract No. 2P03B 01324; the Ministry of Science and Technology of the Russian Federation; the Ministry of Higher Education, Science and Technology of the Republic of Slovenia; the Swiss National Science Foundation; the National Science Council and the Ministry of Education of Taiwan; and the U.S. Department of Energy. |
warning/0506/astro-ph0506579.html | ar5iv | text | # Dynamical Friction and Cooling Flows in Galaxy Clusters
## 1 Introduction
Galaxy clusters are the largest gravitationally bound structures in the universe. They typically contain several hundreds to thousands of galaxies orbiting in a gravitational potential well formed primarily by dark matter (e.g., Bahcall 1977). They are also filled with hot gas with $`T210`$ keV that loses thermal energy prolifically by emitting X-rays. In the absence of any heat sources, the radiative cooling in the cores of rich clusters would result in a cooling flow in which gas settles in the gravitational potential and drops out as cold condensations (e.g., Fabian 1994). However, recent high resolution Chandra and XMM-Newton observations have not shown the expected signature of gas cooling below one-third of the virial temperature (Blanton et al., 2001; Peterson et al., 2001, 2003; Tamura et al., 2001; Johnstone et al., 2002; Allen et al., 2001). This strongly suggests that there must be some source of heat that manages to balance the radiative cooling and thus prevents the mass dropout at the central regions of galaxy clusters. Most theoretical models proposed so far for cluster heating fall into the following two groups: (1) energy injection via radiations, jets, outflows, bubbles, or sound waves from a central active galactic nucleus (Ciotti & Ostriker 2001; Churazov et al. 2002; Kaiser & Binney 2003; Ruszkowski et al. 2004; Roychowdhury et al. 2004; Vecchia et al. 2004 and references therein); (2) diffusive transport of heat from the outer regions of the cluster to the center via conduction (Tucker & Rosner, 1983; Bregman & David, 1988; Narayan & Medvedev, 2001; Voigt et al., 2002; Zakamska & Narayan, 2003; Kim & Narayan, 2003a) or turbulent mixing (Cho et al., 2003; Kim & Narayan, 2003b; Voigt & Fabian, 2004; Dennis & Chandran, 2005).
Although less well studied than those mentioned above, there are clearly other energy sources in the cluster environments. These include kinetic energy in the orbital motions of galaxies (Miller, 1986; Just et al., 1990), gravitational potential energy of the gas (Markevitch et al., 2001; Fabian, 2003), feedback from intracluster supernovae (Domainko et al., 2004), dissipation of turbulent energy driven by infall of subclusters (Fujita et al., 2004), etc. Amount of available energy in each source is comparable to or larger than thermal energy of the intracluster medium (ICM), so that radiative cooling would be easily offset if there exist physical mechanisms that can tap the energy from the sources and convert it into thermal energy in the ICM.
One such process is dynamical fiction (DF) which occurs when a galaxy moving around a cluster center interacts with its gravitationally induced wake in the ICM (e.g., Dokuchaev 1964; Hunt 1971; Ruderman & Spiegel 1971; Rephaeli & Salpeter 1980). Because of the gravitational drag, the galaxy loses some of its kinetic energy. Turbulence in the ICM generated by the superposition of the wakes of many galaxies or by other processes will absorb the lost kinetic energy and, in the presence of tangled magnetic fields and/or viscosity, turns it into heat at the dissipation scales (e.g., Deiss & Just 1996). Miller (1986) and Just et al. (1990) independently estimated the heating rates by DF of galaxies and found that the DF-induced heating can fairly well balance radiative cooling of a rich cluster, provided the mass-to-light ratio is about 20. Using updated data for the properties of galaxies and gas in the Perseus cluster and employing Monte-Carlo approaches, El-Zant, Kim, & Kamionkowski (2004, hereafter EKK04) very recently showed that the total power generated by DF within the cooling radius is comparable to radiative loss from that region if the average mass-to-light ratio of galaxies is about 10. They also noted that this mechanism is self-regulating in the sense that gas would not be overheated since the DF of galaxies becomes ineffective when it attains high enough temperature for which galaxy motions become subsonic (see also Miller 1986).
While the results of the aforementioned work suggest that the total power supplied by the motions of member galaxies is comparable to X-ray luminosity of a typical rich cluster, it is sill uncertain whether it can balance the radiative cooling locally as well. If the DF of galaxy motions is to serve as a primary heat supplier in clusters in equilibrium, observed density and temperature profiles of clusters should be a natural consequence of the local heat balance between radiative cooling and DF-induced heating. Since optically-thin, X-ray emitting gas is prone to thermal instability (e.g., Fabian 1994), it is also an interesting question whether the DF-induced heating is able to suppress thermal instability completely or at least lengthen its growth time to the level comparable to the ages of clusters; even an equilibrium cluster would otherwise be subject to a mass dropout at its center (Kim & Narayan, 2003a).
In this paper, we take one step further from Miller (1986) and EKK04 to investigate equilibrium cluster models in which DF-induced heating is a main heating mechanism. We will first construct density and temperature distributions of the gas in clusters that are in strict hydrostatic and thermal equilibrium and compare them with observed profiles. We note that recent X-ray data from BeppoSAX, Chandra, and XMM-Newton observations indicate that gas temperature in rich clusters rapidly increases with radius for $`r\stackrel{<}{}(0.050.08)r_{180}`$, remains roughly isothermal up to $`r0.2r_{180}`$, and then gradually declines farther out (Ettori et al. 2000; De Grandi & Molendi 2002; Piffaretti et al. 2005; Vikhlinin et al. 2005, see also Markevitch et al. 1998 for ASCA results). Here, $`r_{180}`$ is a virial radius the interior of which has the mean density of gas equal to 180 times the critical density. While declining temperature profiles at large radii may provide important clues to cosmological structure formation, they appear not to be directly relevant to the cooling flow problem. This is because gas density in clusters decreases as or faster than $`r^2`$ at the outer parts, so that the radiative cooling time in the $`r\stackrel{>}{}0.2r_{180}`$ regions are much longer than the Hubble time. Indeed, the results of numerical simulations for cluster formation show that the shapes of declining temperature profiles are essentially independent of the presence of radiative cooling and/or supernova feedback (e.g., Loken et al. 2002; Motl et al. 2004). For this reason, when we compare our model results with observations, we pay attention to density and temperature profiles only in cooling regions with $`r\stackrel{<}{}0.2r_{180}`$ (typically $`0.51`$ Mpc) beyond which the effect of radiative cooling is not serious and thus heating mechanisms are not required.
The remainder of this paper is organized as follows: In §2 we evaluate the total heating rate resulting from the DF of galaxies using the formula given by Ostriker (1999) for gravitational drag force in a gaseous medium. While Miller (1986) and EKK04 approximated galaxy motions as being highly supersonic, we allow for both subsonic and supersonic motions of galaxies. In §3 we calculate equilibrium density and temperature profiles of galaxy clusters by assuming that the DF-induced heating is deposited at the locations of galaxies. Effects of distributed heating of DF and the radial mass variation of galaxies are discussed in §4. In §5 we analyze local thermal stability of the gas while in §6 we present the results of numerical hydrodynamic simulations that investigate the evolution of initially out of equilibrium configurations. Finally, in §7 we conclude this work and discuss other potential consequences of DF of galaxies in clusters.
## 2 Total Heating Rate
Galaxies orbiting around the center of a galaxy cluster gravitationally induce wakes in the ICM. The gravitational interaction between the galaxies and wakes causes the former to lose their orbital kinetic energy and converts it into thermal energy of the ICM via either compressional heating (including shock) or turbulent dissipation of the gas kinetic energy in the wakes. Although the notion that DF of galaxies can generate heat in the hot ICM has been well recognized and widely studied, most of previous work on this subject assumed galaxies in highly supersonic motions, estimated only total heating rate, and compared it with total radiative loss rate in clusters (Ruderman & Spiegel 1971; Rephaeli & Salpeter 1980; Miller 1986; EKK04). However, the motions of galaxies are only slightly supersonic, with an average Mach number of about 1.5 (e.g., Sarazin 1988; Faltenbacher et al. 2004), and even become subsonic at the outer parts of clusters. In this paper we consider both subsonic and supersonic galaxy motions and adopt the general formula of Ostriker (1999) for the DF force in a gaseous medium.
Following Ostriker (1999), we consider a galaxy of mass $`M_g`$ moving at velocity $`𝐯`$ through a uniform medium with density $`\rho `$ and sound speed $`c_s`$. The dynamical-friction force that the galaxy experiences is given by
$$𝐅_{\mathrm{DF}}=\frac{4\pi \rho (GM_g)^2I}{v^3}𝐯,$$
(1)
where the efficiency factor $`I`$ is defined by
$`I\{\begin{array}{cc}\frac{1}{2}\mathrm{ln}\left(\frac{1+}{1}\right),\hfill & <1\hfill \\ \frac{1}{2}\mathrm{ln}\left(1^2\right)+\mathrm{ln}\left(vt/r_{\mathrm{min}}\right),\hfill & >1,\hfill \end{array}`$ (4)
with $`v/c_s`$ being the Mach number of the galaxy motion and $`r_{\mathrm{min}}`$ the effective size of a galaxy (Ostriker, 1999). Note that for $`1`$, equations (1) and (4) with $`vt=r_{\mathrm{max}}`$ corresponding to the maximum system size, become identical to the Chandrasekhar (1943) formula for the drag force due to collisionless particles. Although equation (1) is valid for a perturber moving on rectilinear trajectory through a uniform-density medium, numerical simulations show that it is a good approximation even in a radially-stratified, spherical system (Sánchez-Salcedo & Brandenburg, 2001).
We assume that the orbital velocities of $`N_g`$ galaxies that contribute to heating of the gas are described by the Maxwellian distribution,
$$f(v)=\frac{4\pi N_g}{(2\pi \sigma _r^2)^{3/2}}v^2e^{v^2/(2\sigma _r^2)},$$
(5)
with a one-dimensional velocity dispersion $`\sigma _r`$. The total heating rate due to the DF is then given by
$$\dot{E}=N_g𝐅_{\mathrm{DF}}𝐯=\frac{4\pi \rho (GM_g)^2N_g}{c_s}\frac{I}{},$$
(6)
where the angular brackets denote an average over equation (5).
Using equations (4) and (5), one can evaluate $`I/`$ numerically. Figure 1 plots $`I/`$ as functions of $`m\sigma _r/c_s`$ for some values of $`\mathrm{ln}(vt/r_{\mathrm{min}})`$. For supersonic cases, the density perturbations in the wake are highly asymmetric and the regions influenced by the perturber shrink as $`m`$ increases. This results in a smaller heating rate for larger $`m`$. In this case, one can show $`I/(2/\pi )^{1/2}m^1\mathrm{exp}(0.5/m^2)\mathrm{ln}(vt/r_{\mathrm{min}})`$ for $`m1`$. For a subsonically moving perturber, on the other hand, the perturbed density in the front and back sides of the perturber becomes more or less symmetric, producing a smaller heating rate with decreasing $`m`$. In the limit of vanishingly small $`m`$, $`I/m^2`$. In general, $`I/`$ peaks at $`m=1`$, as Figure 1 indicates.
Taking $`vt=1`$ Mpc from a typical size of clusters and $`r_{\mathrm{min}}=10`$ kpc as the galaxy size, $`\mathrm{ln}(vt/r_{\mathrm{min}})4.6`$. Since $`m`$ usually varies from 0.8 to 3, $`I/2`$. Therefore, we have the total heating rate
$$\dot{E}4\times 10^{44}\mathrm{ergs}\mathrm{s}^1\left(\frac{N_g}{500}\right)\left(\frac{M_g}{10^{11}M_{}}\right)^2\left(\frac{n_e}{0.01\mathrm{cm}^3}\right)\left(\frac{T}{5\mathrm{keV}}\right)^{1/2}\left(\frac{I/}{2}\right),$$
(7)
where $`n_e`$ and $`T`$ denote the electron number density and gas temperature, respectively, and we adopt the solar abundance. Notice that the rate of gas heating given in equation (7) is similar to typical X-ray luminosity of rich clusters (e.g., Sarazin, 1988; Rosati, Borgani, & Norman, 2002), implying that there is a sufficient amount of energy available in the orbital motions of galaxies. This is essentially what led EKK04 to conclude that the dynamical-friction coupling between cluster galaxies and gas can provide thermal energy enough to compensate for radiative loss. However, it is still questionable whether DF of galaxies heats the gas in a right manner. That is, are the observed density and temperature profiles of the ICM a direct consequence of energy balance between heating by DF and radiative cooling? Is the intracluster gas heated by DF thermally stable? In what follows, we shall study models of clusters with DF-induced heating in detail by assuming hydrostatic equilibrium and thermal energy balance.
## 3 Model With Localized Heating
DF of galaxies is mediated by gravity which is a long-range force, so that heating of gas due to a single galaxy is likely to be well distributed throughout its wake. The functional form of heat distribution is unknown and its derivation may require numerical simulations, which are beyond the scope of the present work. In this section, we instead make a simplifying assumption that heat is deposited locally at the galaxy position. The effects of heat distribution will be discussed in §4.
### 3.1 Local Heating Function
Consider a spherically symmetric distribution of galaxies, with the number density given by $`n_g(r)`$ at radius $`r`$. The local heating rate per unit volume due to DF is given by
$$\dot{e}(r)=\frac{4\pi \rho (GM_g)^2I/}{c_s}n_g(r),$$
(8)
where we assume equal galaxy mass. To find $`n_g(r)`$, we assume that the orbits of galaxies are isotropic and isothermal. The usual Jeans equation (e.g., eq. \[4-55\] of Binney & Tremaine 1987) then reads
$$\sigma _r^2\frac{d\mathrm{ln}n_g}{dr}=\frac{d\mathrm{\Phi }}{dr},$$
(9)
where $`\mathrm{\Phi }`$ is the gravitational potential.
Under the NFW distribution of dark matter
$$\rho _{\mathrm{NFW}}=\frac{M_0/2\pi }{r(r+r_s)^2},$$
(10)
with a scale radius $`r_s`$ and a characteristic mass $`M_0`$, the gravitational acceleration is given by
$$\frac{d\mathrm{\Phi }}{dr}=\frac{2GM_0}{r_s^2}\left[\frac{\mathrm{ln}(1+x)}{x^2}\frac{1}{x(1+x)}\right],$$
(11)
where $`xr/r_s`$ is the dimensionless radius (Navarro, Frenk, & White, 1997; Klypin et al., 2001). Combining equations (9) and (11) together, one can easily find
$$n_g(x)=n_g(0)g(x)=N_g\frac{g(x)}{r_s^3g(x^{})d^3x^{}},$$
(12)
where
$$g(x)\left[\frac{(1+x)^{1/x}}{e}\right]^\eta ,$$
(13)
with the dimensionless parameter $`\eta `$ defined by
$$\eta \frac{2GM_0}{r_s\sigma _r^2}.$$
(14)
Note that $`\eta `$ measures the ratio of the gravitational potential energy of a galaxy in a cluster to that galaxy’s kinetic energy at $`r=r_s`$.
While the functional form of $`g(x)`$ looks strange, it behaves quite well near $`x0`$ and gives density profiles similar to the isothermal $`\beta `$-model of gas distributions (Makino et al., 1998) or to the King model of the observed galaxy distributions (e.g., Girardi et al., 1998). For example, Girardi et al. (1998) found that the best-fit distribution of galaxies in A1795 is $`n_g(r)=n_g(0)[1+(r/R_c)^2]^{1.27}`$ with a core radius $`R_c=43`$ pc, which is plotted as a solid line in Figure 2. On the other hand, X-ray and optical observations indicate $`M_0=6.6\times 10^{14}M_{}`$, $`r_s=460`$ kpc, and $`\sigma _r800\mathrm{km}\mathrm{s}^1`$ for A1795 (Girardi et al., 1998; Ettori et al., 2002), corresponding to $`\eta 19`$. Figure 2 plots as a dotted line the $`g(x)`$ curve with $`\eta =19`$, which is in good agreement with the observed galaxy distribution.
We take $`2r_s`$ as the outer boundary of a cluster. Inserting equation (12) into equation (8), one obtains
$`\dot{e}(x)`$ $`=`$ $`{\displaystyle \frac{4\pi \rho (GM_g)^2N_g}{c_sr_s^3g(x^{})d^3x^{}}}{\displaystyle \frac{I}{}}g(x)`$ (15)
$`=`$ $`1.7\times 10^{25}r_{s,460}^3N_{g,500}M_{g,11}^2n_eT_{\mathrm{keV}}^{1/2}I/g(x)\mathrm{erg}\mathrm{cm}^3\mathrm{s}^1,`$
where $`r_{s,460}=r_s/(460`$ kpc), $`N_{g,500}=N_g/500`$, $`M_{g,11}=M_g/(10^{11}M_{})`$, and $`T_{\mathrm{keV}}`$ is the temperature of the gas in units of keV. Equation (15) is our desired equation for the volume heating rate due to the DF of galaxies. Note that $`I/2`$ at $`m1`$ and depends on temperature.
### 3.2 Equilibrium Model
We look for cluster density and temperature distributions in which the hot gas satisfies both hydrostatic equilibrium and energy balance between radiative cooling and DF-induced heating. For thermal bremsstrahlung, the rate of energy loss per unit volume is given by
$$j=7.2\times 10^{24}n_e^2T_{\mathrm{keV}}^{1/2}\mathrm{erg}\mathrm{cm}^3\mathrm{s}^1,$$
(16)
(Rybicki & Lightman, 1979; Zakamska & Narayan, 2003). From equations (15) and (16), the condition of thermal energy balance yields
$$n_eT_{\mathrm{keV}}=2.4\times 10^2r_{s,460}^3N_{g,500}M_{g,11}^2I/g(x).$$
(17)
Neglecting the weak temperature dependence of $`I/`$ and assuming that $`M_g`$ does not vary with radius, equation (17) states that the gas pressure in equilibrium should trace the distribution of galaxies.
Hydrostatic equilibrium of the gas requires
$$\frac{c_0^2}{n_e}\frac{d(n_eT_{\mathrm{keV}})}{dr}=\frac{d\mathrm{\Phi }}{dr},$$
(18)
where $`c_0=390\mathrm{km}\mathrm{s}^1`$ is an isothermal sound speed at $`T_{\mathrm{keV}}=1`$. Substituting equations (11) and (17) into equation (18) and using equation (13), we obtain
$$\frac{d}{dx}\mathrm{ln}I/=\eta \left(\frac{\sigma _r^2}{c_0^2T_{\mathrm{keV}}}1\right)\left[\frac{\mathrm{ln}(1+x)}{x^2}\frac{1}{x(1+x)}\right].$$
(19)
Since $`I/`$ depends on $`T_{\mathrm{keV}}`$, equation (19) can be integrated as an initial value problem to yield $`T_{\mathrm{keV}}(x)`$. The equilibrium profile for $`n_e(x)`$ can then be found from equation (17). Since $`n_e(x)`$ depends rather sensitively on less-well known quantities such and $`N_g`$ and $`M_g^2`$ and is easily influenced by the mass profile of galaxies, we concentrate on $`T_{\mathrm{keV}}(x)`$ that is independent of the properties of galaxies other than $`\sigma _r`$.
Equation (19) has a trivial isothermal solution
$$T_{\mathrm{iso}}=6.5\left(\frac{\sigma _r}{10^3\mathrm{km}\mathrm{s}^1}\right)^2\mathrm{keV},$$
(20)
independent of the radius. The corresponding density profile is
$$n_e=7.4\times 10^3r_{s,460}^3\sigma _{r,3}^2N_{g,500}M_{g,11}^2g(x),$$
(21)
where $`\sigma _{r,3}=\sigma _r/(10^3\mathrm{km}\mathrm{s}^1)`$, indicating that the gas density follows the galaxy distribution exactly. Considering uncertainties in the values of $`r_{s,460}`$, $`\sigma _{v,1000}`$, $`N_{g,500}`$, and $`M_{g,11}`$,<sup>1</sup><sup>1</sup>1For example, the average mass of galaxies would increase up to $`5\times 10^{11}M_{}`$ if the contribution from dark halos is taken into account (e.g., Zentner & Bullock, 2003). equation (21) gives $`n_e0.010.2\mathrm{cm}^3`$ at the cluster centers, which is not much different from observed values. Note that the Mach number corresponding to $`T_{\mathrm{iso}}`$ is $`m=\gamma ^{1/2}`$. For later purposes, we define the transonic temperature
$$T_{\mathrm{sonic}}\frac{1}{\gamma }T_{\mathrm{iso}},$$
(22)
corresponding to unity Mach number of galaxy motions.
The isothermal solution (20) is consistent with what has long been known as the scaling relation between the gas temperature and the velocity dispersion of galaxies for many clusters (e.g., Mushotzky 1984; Edge & Stewart 1991; Lubin & Bahcall 1993; Wu et al. 1999, see also Sarazin 1988; Rosati, Borgani, & Norman 2002 for review). It explains why the values of the $`\beta `$ $`(\mu m_p\sigma _r^2/kT)`$ parameter in the $`\beta `$ models for gas distributions are close to unity (Lubin & Bahcall, 1993; Bahcall & Lubin, 1994; Girardi et al., 1996b, 1998). Physically, the scaling relation implies that the plasma and the galaxies are well relaxed under the common gravitational potential. It is uncertain whether the observed scaling law results truly from the DF coupling between the ICM and galaxies, but equation (20) suggests that the latter certainly makes the former tighter.
Although gas in the regions surrounding cooling cores of clusters has nearly constant temperatures close to the virial values, recent X-ray observations exhibit positive radial gradients of gas temperatures at the central $`300`$ kpc regions (Peterson et al., 2001, 2003), while showing declining temperature distributions at far outer regions (Markevitch et al., 1998; De Grandi & Molendi, 2002; Piffaretti et al., 2005; Vikhlinin et al., 2005). Thus, the isothermal solution cannot describe the observed temperature distributions for the whole range of radii. We are particularly interested in the regions within and adjacent to the cooling radius in which gas density is high enough to experience significant radiative cooling. To check whether the general solutions of equation (19) produce temperature and density distributions similar to the observed in such regions, we solve it numerically. We adopt $`\eta =19`$, $`r_s=460`$ kpc, $`N_g=500`$, $`M_g=10^{11}M_{}`$, and $`\sigma _r=10^3\mathrm{km}\mathrm{s}^1`$ corresponding to $`T_{\mathrm{iso}}=6.5`$ keV. We choose a value for the central temperature $`T(0)`$ and then integrate equation (19) from $`r=0`$ to $`2r_s`$. The resulting temperature and density profiles for a few selected values of $`T(0)`$ are plotted in Figure 3.
The spatial behavior of the solutions depends critically on $`T(0)`$. When $`T(0)>T_{\mathrm{sonic}}`$, galaxy motions are subsonic and $`dI//dm>0`$ from Figure 1. If $`T(0)>T_{\mathrm{iso}}`$ (or $`m<\gamma ^{1/2}`$), equation (19) gives $`dI//dr=(m/2T)(dI//dm)(dT/dr)>0`$ for $`T(0)>T_{\mathrm{iso}}`$, implying $`dT/dr<0`$. As temperature monotonically decreases with radius, $`m`$ increases until it reaches the isothermal value $`\gamma ^{1/2}`$ where $`T=T_{\mathrm{iso}}`$. One can similarly show that $`dT/dr>0`$ for $`T_{\mathrm{iso}}>T(0)>T_{\mathrm{sonic}}`$, and thus temperature slowly increases toward $`T_{\mathrm{iso}}`$. As long as $`T(0)>T_{\mathrm{sonic}}`$, the solutions asymptote to the isothermal ones (eqs. and ) as the radius increases. On the other hand, $`dT/dr<0`$ when $`T(0)<T_{\mathrm{sonic}}`$ (or $`m>1`$). Since $`m`$ increases with decreasing $`T`$ and since $`dI//dr`$ does not change its sign for $`m>1`$, $`dT/dr<0`$ is satisfied for all radii. The decrease of temperature is much faster than that of $`g(x)`$, resulting in a unrealistic distribution of electron number density that increases with radius.
As Figure 3 shows, the local heat balance with DF-induced heating yields rising temperature profiles, if and only if $`T_{\mathrm{sonic}}<T<T_{\mathrm{iso}}`$. While this is a tantalizing result considering the central depression of temperatures seen in cooling-flow clusters, the required range of temperatures is very narrow. If we identify $`T_{\mathrm{iso}}`$ with the virial temperature, the central temperatures in our equilibrium models must be larger than $`T_{\mathrm{sonic}}=0.6T_{\mathrm{iso}}`$ (for $`\gamma =5/3`$), which is about twice larger than the observed values of $`T(0)(0.30.4)T_{\mathrm{iso}}`$ (e.g., Peterson et al. 2001, 2003; Allen et al. 2001). By having too stringent upper and lower limits of temperatures, therefore, heating by DF of galaxies alone is unlikely to explain observed temperature and density distributions of galaxy clusters. We note however that the tight range of temperatures for the rising temperature profiles may be caused by the basic assumptions of this section, namely that DF-induced heating of the ICM is all localized at the positions of galaxies, and that all galaxies have equal mass. We relax these assumptions in the next section.
## 4 Effects of Distributed Heating
In the presence of viscosity and/or turbulent magnetic fields in the ICM, the kinetic energy lost by a galaxy via DF will eventually be converted into heat, rather than bulk motions distributed onto the gravitational wakes of the galaxy. Deiss & Just (1996) used a quasi-linear fluctuation theory to derive the spatial structure (in Fourier space) of mechanical heating in a turbulent medium self-consistently driven by DF of many galaxies. To our knowledge, there is no published work that study spatial heat distribution (in real space) caused by DF. If ICM heating by DF of galaxies occurs through turbulent dissipation of gas kinetic energy, finding the functional form of the heat distribution would not be viable unless the characteristics of ICM turbulence and related processes are prescribed (Deiss, Just, & Kegel, 1990). Instead of attempting to derive a realistic heat distribution, in this work we simply adopt a Gaussian function to parametrize the spatial extent to which heat is distributed<sup>2</sup><sup>2</sup>2We also tried with logarithmic heat distribution functions used in EKK04 and found that results are similar to those with Gaussian functions presented in this section.. Our aim in this section is to examine the effects of heat distribution on equilibrium structures in comparison with the localized heating models. Since masses of galaxies inside a cluster are likely to vary with the radius within the cluster, we also allow for spatially varying galaxy masses.
Assuming that heat distribution follows a Gaussian profile centered at the location of the galaxy, we write the heating rate per unit volume as
$$\dot{e}(r)=\frac{4\pi \rho G^2I/}{c_s}h(r),$$
(23)
where
$$h(r)\frac{1}{\pi ^{1/2}l_h}n_g(r^{})M_g(r^{})^2e^{(rr^{})^2/l_h^2}𝑑r^{},$$
(24)
is a convolution of $`n_gM_g^2`$ with the Gaussian function with scale length $`l_h`$. In writing equations (23) and (24), we assume that density and temperature of the gas do not change significantly inside the wake of a galaxy, which is acceptable within the framework of equation (1). The condition of hydrostatic equilibrium and thermal energy balance is then reduced to
$$\frac{d}{dx}\mathrm{ln}I/=\eta \frac{\sigma _r^2}{c_0^2T_{\mathrm{keV}}}\left[\frac{\mathrm{ln}(1+x)}{x^2}\frac{1}{x(1+x)}\right]\frac{d\mathrm{ln}h(x)}{dx}.$$
(25)
It can be easily verified that equation (25) recovers equation (19) in the limit of $`l_h0`$.
Figure 4 plots the solutions of equation (25) for the case of constant $`M_g`$. More widely distributed heating is equivalent to localized heating with a flatter distribution of galaxy number density, causing the $`|d\mathrm{ln}h/dx|`$ term in equation (25) to be smaller. Consequently, an equilibrium cluster with a larger value of $`l_h`$ tends to make $`I/`$ decreases faster with increasing $`r`$. This implies that for $`T(0)<T_{\mathrm{sonic}}`$ ($`T(0)>T_{\mathrm{sonic}}`$), the temperature in models with non-zero $`l_h`$ decreases (increases) more rapidly than in the $`l_h=0`$ counterpart, although the difference between models with $`l_h/r_s=0.1`$ and 1.0 is negligible for small $`r`$. Note that even when $`T(0)>T_{\mathrm{iso}}`$, for which localized-heating models have declining temperature profiles, equilibrium temperature in distributed-heating models increases in the regions with $`r<l_h`$ and eventually converges to $`T_{\mathrm{iso}}`$ at $`r\stackrel{>}{}(23)l_h`$. We see that distributed heating does not change the condition $`T_{\mathrm{sonic}}<T<T_{\mathrm{iso}}`$ for the existence of rising temperature distributions. The distributed-heating models do not work better than the localized-heating model in terms of producing temperature distributions similar to observations; it in fact aggravates the situation by making all of the gas in clusters nearly isothermal except only in the central $`1`$ kpc regions.
Another factor that may affect the equilibrium structure is the radial mass variation of cluster galaxies. Although galaxies located near the central parts of clusters tend to have a large fraction of luminous matter, (e.g., Biviano et al. 2002; Mercurio et al. 2003), the strong tidal stripping of dark-matter halos is likely to cause the total (luminous + dark) masses of individual galaxies to be smaller toward the cluster center (Zentner & Bullock 2003; EKK04; Nagai & Kravtsov 2005). Motivated by this consideration, we adopt several algebraic forms of $`M_g(r)`$ that either monotonically decrease or increase by about a factor of 20 or less from the cluster center to $`r=2r_s`$, and then solve equation (25). The resulting temperature profiles are almost identical to those of the constant-mass cases presented in Figure 4. Because the distribution of galaxy masses appears logarithmically in equation (25), it does not considerably affect temperature structures, although it dramatically changes electron number-density profiles.
## 5 Thermal Stability
We now discuss thermal stability of a system in which radiative cooling is locally balanced by DF-mediated heating. For spatially localized heating, the net loss function $`\rho `$ is given by
$$\rho =j\dot{e}=\alpha \rho ^2T^{1/2}\beta \rho T^{1/2+p},$$
(26)
where the power index $`p`$ accounts for the local temperature dependence of $`I/`$ in a piecewise manner, and the positive constants $`\alpha `$ and $`\beta `$ contain all information on the cluster properties other than gas density and temperature. Figure 5 plots $`p=d\mathrm{ln}I//d\mathrm{ln}T`$ as a function of $`m`$, which gives $`p0`$ at around the transonic temperature ($`m1`$), while $`p0.30.5`$ for supersonic temperatures ($`m>1`$). The curves asymptote to $`1`$ for $`m1`$ and to $`0.5`$ for $`m1`$, as the asymptotic formulae given in §2 suggest. Note that $`p`$ is insensitive to the choice of $`\mathrm{ln}(vt/r_{\mathrm{min}})`$ as long as $`m>0.8`$.
Thermal stability of a system can be easily checked by using the generalized Field criterion which states that the system is thermally unstable to isobaric perturbations if
$$\frac{(/T)}{T}|_P=\frac{1}{\rho T}\left[\frac{(\rho )}{T}|_\rho \frac{\rho }{T}\frac{(\rho )}{\rho }|_T\right]=p\frac{\alpha \rho }{T^{3/2}}<0,$$
(27)
where the second equality assumes $`\rho =0`$ corresponding to thermal equilibrium (Field, 1965; Balbus, 1986). Since $`p>0`$ for $`m>1`$ (and $`p<0`$ for $`m<1`$), this implies that the ICM is thermally unstable (stable) if heating is caused preferentially by galaxies moving at a supersonic (subsonic) speed. Dense inner parts of galaxy clusters, where the cooling time is less than the Hubble time, are filled with low-temperature gas such that galaxy motions are readily supersonic, implying that DF-induced heating is unable to quench thermal instability in those cooling cores. Whether thermal instability has important dynamical consequences on cluster evolution can be judged by considering its growth time which amounts to
$$t_{\mathrm{grow}}=\frac{\gamma }{\gamma 1}\frac{P}{\rho T^2}\frac{(/T)}{T}|_P^1=0.96p^1\mathrm{Gyr}\left(\frac{n_e}{0.05\mathrm{cm}^3}\right)^1\left(\frac{k_BT}{2\mathrm{keV}}\right)^{1/2},$$
(28)
(cf. Kim & Narayan 2003a). Even if cooling cores are thermally unstable, therefore, the virulence of thermal instability may or may not manifest itself during the lifetime of clusters depending on the value of $`p`$. For example, for a cluster with $`\sigma _r=10^3\mathrm{km}\mathrm{s}^1`$ and $`T=2`$ keV at the central regions, $`m=1.4`$ and $`p=0.23`$ from Figure 5, giving $`t_{\mathrm{grow}}=4.2`$ Gyr. This is almost comparable to the ages of clusters since the last major merger event ($`7`$ Gyr for massive clusters; Kitayama & Suto 1996), suggesting that thermally instability may be dynamically unimportant for practical purposes.
The inability of DF-induced heating to suppress thermal instability seems not to be a fatal problem as long as the following two conditions are met: i) cooling cores are in thermal equilibrium, a basic premise of the stability analysis; and ii) galaxies near the central parts move at near-transonic speeds ($`m\stackrel{<}{}2`$), so as to have a sufficiently small value of $`p`$. While the second condition appears to be easily satisfied in clusters, the results of the preceding section suggests that the first condition may not be so if DF-induced heating is to balance radiative cooling for the whole range of radii. Since the cooling time outside the cooling radius is longer than the Hubble time, one may argue that material beyond the cooling radius does not need any heating source and that it is sufficient for DF of galaxies to supply heat to gas only in cooling cores. We address this possibility in the next section.
As a final remark of this section, we compare our results with previous work. Formally, equation (27) implies that gas subject to DF-induced heating has a critical temperature above (below) which it becomes thermally stable (unstable) and that this critical temperature corresponds exactly to unity Mach number of galaxy motions. This result is quite different from those of previous work which made some approximations on the heating function. For instance, Miller (1986) argued that gas with drag heating is always thermally unstable and rapidly turns into a multi-phase medium, a consequence of his neglecting the temperature dependence of the heating function (corresponding to $`p=1/2`$ in eq. ). Although Bregman & David (1989) found a similar critical gas temperature, they considered only the supersonic regime and used an approximate heating function from the shock jump conditions, which makes their critical temperature smaller than ours by a factor of about three. Our result, based on the general formula of Ostriker (1999) for the drag force applicable for both supersonic and subsonic cases, represents the situation better.
## 6 Time Dependent Approach
We have shown in the preceding sections that an equilibrium cluster subject to DF-mediated heating has a temperature profile that either decreases (for $`T>T_{\mathrm{iso}}`$ or $`T<T_{\mathrm{sonic}}`$) or increases (for $`T_{\mathrm{sonic}}<T<T_{\mathrm{iso}}`$) with radius. While the radially increasing temperature profile is attractive, the requirement that thermal energy is balanced locally in the entire regions out to $`2r_s`$ causes the central gas temperature to be no less than 0.6 times the virial temperature, about a factor two larger than observations. Since radiative cooling is important only in the dense central parts inside the cooling radius, however, it would be sufficient for DF of galaxies to provide heat only in such a cooling core instead of the entire regions. In addition, thermal instability would not be an issue provided that gas is in thermal equilibrium and that galaxy motions are near transonic (§5). Therefore, it may be possible for a cluster to start from an arbitrary non-equilibrium state and evolve slowly under both radiative cooling and DF-mediated heating, ending up with an equilibrium core in which thermal instability has a very long growth time. The resulting temperature profile, albeit consistent with observations, need not represent an equilibrium for the whole range of radii.
To test this idea, we have run a number of one-dimensional hydrodynamic simulations in spherical polar coordinates. We set up a logarithmically spaced radial grid, with 400 zones, from 1 kpc to 1 Mpc, and explicitly implement the net cooling-heating function described by equation (26). We also include the effect of gravitational potential due to passive dark matter using equation (11) with $`M_0=6.6\times 10^{14}M_{}`$ and $`r_s=460`$ kpc; the DF-coupling between gas and the dark matter and back reaction on the latter component are not taken into account. As initial conditions, we consider spherically symmetric clusters and take $`n_e=n_e(0)g(x)\mathrm{cm}^3`$ with the central electron density $`n_e(0)=0.002`$, 0.005, 0.01, 0.02 cm<sup>-3</sup>, and fix $`T=6.5`$ keV independent of the radius for all models. We adopt $`r_s=460`$ kpc, $`N_g=500`$, and $`M_g=10^{11}M_{}`$, and consider both localized heating and distributed heating with $`l_h/r_s=`$0.1 and 1. Note that the corresponding equilibrium density profile has $`n_e(0)=0.0074`$ cm<sup>-3</sup> when heating is localized. We adopt the outflow boundary conditions for scalar variables (i.e., density, energy, etc.) that assign the same values in the ghost zones as in the corresponding active zones. For the radial velocity, we allow it to vary as a linear function of radius at the inner boundary, while fixing it to be zero at the outer boundary (e.g., Kim & Narayan 2003a). Using the ZEUS hydrodynamic code (Stone & Norman, 1992), we solve the basic hydrodynamic equations and follow the nonlinear evolution of each model cluster. The resulting radial profiles of electron number density and temperature of model A with $`n_e(0)=0.02\mathrm{cm}^3`$ and model B with $`n_e(0)=0.005\mathrm{cm}^3`$, both of which assume spatially localized heating, are shown in Figure 6 for a few time epochs.
Model A, which initially has larger density than the equilibrium value everywhere, immediately develops radial mass inflows. As time evolves, the temperature drops and the density increases. Compared to pure cooling cases which fairly well maintain near isobaric conditions (e.g., Kim & Narayan 2003a), DF-induced heating is found to cause thermal pressure to increase with density as $`P\rho ^{0.3}`$. Since $`jT^{2.4}`$ and $`\dot{e}T^{1.9+p}`$ with $`p>0`$ for supersonic galaxy motions, however, the system is thermally unstable and cooling occurs at a much faster rate as temperature decreases. Model A experiences a cooling catastrophe in less than 4 Gyr. On the other hand, model B with initial overheating becomes hotter and less dense, which in turn tends to increase the ratio of heating rate to the cooling rate for $`p\stackrel{>}{}0.5`$. As the gas temperature increases, the Mach number of galaxy motions become smaller, reducing the value of $`p`$. The cooling rate will therefore eventually exceed the heating rate and the cluster evolution will be reversed. In the case of model B, this turnaround will occur at $`t=70`$ Gyr when $`T18`$ keV (or $`m0.5`$) is reached.
The evolution of models with different density and different length scales for heat distribution is qualitatively similar to those of models A and B, namely that clusters catastrophically evolve toward vanishing central temperature if cooling dominates initially, while clusters with initial excessive heating heat up steadily. This suggests that DF-induced heating does not naturally lead non-equilibrium clusters to thermally stable, equilibrium cores. As shown in §3 and §4, DF of galaxies does not explain the observed temperature distributions of clusters if the condition of thermal energy balance is imposed for all radii. Although it is enough for DF-induced heating to balance radiative cooling only in cooling cores, such cooling cores do not naturally form from non-equilibrium states. Unless the properties and galaxies and ICM are fine tuned, small departures from an equilibrium state rapidly evolve into an extreme configuration. Therefore, we conclude that DF-induced heating alone is not likely to account for the absence of cold gas in the centers of galaxy clusters.
Although heating by DF from galaxies does not appear to provide a complete solution to the cooling flow problem, we see that the DF-induced heating can still offset a considerable amount of radiative cooling. The isobaric cooling time is given by
$$t_{\mathrm{cool}}=\frac{\gamma }{\gamma 1}\left(\frac{P}{\rho }\right)=\frac{0.96}{1/(n_eT_{\mathrm{keV}})}\left(\frac{n_e}{0.05\mathrm{cm}^3}\right)^1\left(\frac{kT}{2\mathrm{keV}}\right)^{1/2}\mathrm{Gyr},$$
(29)
where $`0.024r_{s,460}^3N_{g,500}M_{g,11}^2I/`$ represents the contribution of the DF-induced heating. Figure 7 plots as solid lines the cooling time as a function of $`n_e`$ for $`N_{g,500}M_{g,11}^2=1`$ or as a function of $`N_{g,500}M_{g,11}^2`$ with varying density. The cases without any heating are compared as dotted lines. The time epochs when clusters experience cooling catastrophe in the numerical simulations are marked in Figure 7$`a`$ as solid circle, triangle, cross, and open circle for no heating, the distributed heating with $`l_h/r=1.0`$ and 0.1, and the localized heating cases, respectively. The numerical results are in good agreement with the theoretical prediction. Models with distributed heating tend to have longer cooling time. When $`N_{g,500}M_{g,11}^2=1`$ and $`n_e(0)=0.05\mathrm{cm}^3`$, Figure 7$`a`$ shows that DF-induced heating lengthens the cooling time by about 13% compared to the pure cooling case. However, as Figure 7$`b`$ indicates, the offset of cooling by DF-induced heating is increasingly larger as $`N_{g,500}M_{g,11}^2`$ increases. The increment of the cooling time could be as large as 140% when $`n_e(0)0.030.1\mathrm{cm}^3`$ and $`N_{g,500}M_{g,11}37`$, suggesting that thermal energy supplied by DF of galaxies is by no means non-negligible.
## 7 Conclusions and Discussion
Friction of galaxy motions via the gravitational interaction with their own gravitationally induced wake in the ICM has often been invoked as an efficient heating mechanism of the gas (e.g., Miller 1986; Just et al. 1990; Deiss & Just 1996; EKK04). This idea of DF-mediated heating of the ICM is quite attractive because there is sufficient energy available in galaxy motions and the mechanism is self-regulating; it operates effectively only when the temperature of gas is in a certain range, which happens to be the typical gas temperatures in the cooling cores of galaxy clusters. In this paper, we take one step further from Miller (1986) and EKK04 to calculate equilibrium density and temperature profiles of the hot gas heated by DF of galaxies. Instead of restricting to cases where galaxy motions are all supersonic, we use the general formula derived by Ostriker (1999) for the drag force that takes allowance for both supersonic and subsonic galaxy motions in a gaseous medium. We show that the total heating rate due to the DF of galaxies in a typical rich cluster is comparable to its total X-ray luminosity, confirming the results of Miller (1986) and EKK04 (see §2).
Next, we derive the local heating function (eq. ) assuming that the orbits of galaxies are isotropic and isothermal under the NFW distribution of dark matter and that the kinetic energy lost by a galaxy is deposited into heat at the location of the galaxy. The condition that the gas is in hydrostatic equilibrium and maintains energy balance requires the temperature profile to be one of the following three kinds: isothermal, with $`T=T_{\mathrm{iso}}`$; a decreasing profile with radius when $`T<T_{\mathrm{sonic}}`$ or $`T>T_{\mathrm{iso}}`$; an increasing profile when $`T_{\mathrm{sonic}}<T<T_{\mathrm{iso}}`$, where $`T_{\mathrm{iso}}`$ and $`T_{\mathrm{sonic}}`$ denote the temperatures corresponding to unity Mach number of galaxy motions with respect to the isothermal and adiabatic sound speeds, respectively (eqs. and ). The isothermal solution is interesting because it quite well describes the observed scaling relationship among clusters properties. Although the rising temperature solution is attractive since clusters usually show a temperatures drop at the central regions, it strictly requires the central temperature to be no lower than 0.6 times the virial temperature, which is roughly twice smaller than the observed values (§3). We also consider cases in which DF-induced heating is distributed in space according to a Gaussian form, and/or the masses of cluster galaxies vary over radius, and find that the stringent limit of $`T>T_{\mathrm{sonic}}=0.6T_{\mathrm{iso}}`$ for rising temperature distributions remains unchanged (§4).
Using the local heating function we have derived, we analyze thermal stability of the gas subject to radiative cooling and DF-mediated heating (§5). When galaxy motions are subsonic, the heating rate that deceases steeply with temperature suppresses thermal instability completely. On the other hand, supersonic galaxy motions for which the heating rate is relatively insensitive to temperature are unable to stop the growth of thermal instability. The growth time in the presence of DF-induced heating is at least twice that in the pure cooling case, and becomes progressively longer as the Mach number of galaxies decreases. When galaxies move at slightly supersonic velocities with Mach number less than 2, which is very likely at cluster centers, the growth time becomes comparable to the ages of clusters. This implies that thermal instability, even if operates, may have insignificant dynamical consequences on cluster evolution.
Noting that regions outside the cooling radius need not be in strict thermal equilibrium, we look for a possibility using numerical hydrodynamic simulations that clusters evolve from an arbitrary non-equilibrium state and form current cooling cores in which DF-induced heating balances radiative cooling (§6). We find that clusters that were initially dominated by cooling unavoidably develop radial mass inflows and decreases their central temperatures in a runaway fashion, whereas clusters with initial overheating slowly heat up and result in radially decreasing temperature profiles. Equilibrium solutions therefore do not appear to form an attracting set for galaxy and gas configurations, suggesting that when DF from galaxies is the sole heating source, it is extremely difficult to obtain equilibrium cores by smoothly evolving non-equilibrium clusters (even if in some cases the cooling catastrophe can be deferred so as to occur on a longer timescale).
Putting together all the results of this paper we conclude that DF of galaxies alone, albeit an interesting heating mechanism with a lot of available energy, cannot be the lone heating agency to balance radiative cooling in rich galaxy clusters. We nonetheless note that the heating due DF of galaxies could considerably lengthen the cooling time, depending on the value of $`N_gM_g^2`$, and thus should not be neglected in energetics of galaxy clusters.
One of the key assumptions made in this paper is that all the kinetic energy lost by galaxies via DF is transferred to the thermal energy of the surrounding gas. The rationale behind this assumption is that the superposition of the wakes produces turbulence in the ICM and the kinetic energy injected into the ICM turbulence at saturation cascades down along the Kolmogorov-like energy spectrum, and eventually transforms into heat through viscous dissipation at small scales. Another possibility is that a large fraction of the kinetic energy of galaxies is used to merely enhance the level of the turbulence instead of being converted into heat. This may occur when the turbulence is not fully developed yet.
The characteristics of the ICM turbulence is not well known, yet observations and numerical simulations suggest an average velocity dispersion $`v_{\mathrm{turb}}200400\mathrm{km}\mathrm{s}^1`$ on scales of $`\lambda 520`$ kpc (Ricker & Sarazin, 2001; Carilli & Taylor, 2002; Churazov et al., 2004; Schuecker et al., 2004; Faltenbacher et al., 2004). The associated turbulent kinetic energy is $`(1/2)M_{\mathrm{gas}}v_{\mathrm{turb}}^29\times 10^{61}`$ erg for the total mass $`M_{\mathrm{gas}}10^{14}M_{}`$ of the ICM in a rich cluster. Let us assume that this energy is supplied solely by DF of galaxies at a rate given by equation (7). If the level of turbulence keeps increasing without dissipation, it would take about 7 Gyr for the DF of galaxies to feed the ICM turbulence to the observed level, which is almost comparable to the lifetime of the cluster. Conversely, if the turbulence is fully developed and in a steady state such that the energy injection rate by DF is equal to the rate of dissipation, it would have $`v_{\mathrm{turb}}=(2\lambda \dot{E}/M_{\mathrm{gas}})^{1/3}60\mathrm{km}\mathrm{s}^1`$ for $`\lambda 20`$ kpc, which is too small to explain the observations. All of these imply that the contribution from the DF of galaxies to the ICM turbulence is not considerable (see also Sánchez-Salcedo et al. 1999). If the dissipation of the (fully developed) turbulence is to provide enough heat to balance radiative cooling, as suggested by Churazov et al. (2004) and Fujita et al. (2004), the energy injection should not be entirely due to the DF of galaxies; it requires other energy sources, including jets from active galactic nuclei and mergers with smaller groups or clusters.
The assertion that the conversion of galaxy kinetic energy into thermal energy of the gas may not be so effective is supported by the results of Faltenbacher et al. (2004) who studied using numerical simulations cluster formation in $`\mathrm{\Lambda }`$CDM cosmology, with the DF of galaxies included explicitly. They found that (1) the motions of galaxies at the present epoch are slightly supersonic, with an average Mach number of $`1.4`$; (2) gas within the virial radius has a velocity dispersion $`v_{\mathrm{turb}}(0.30.5)c_s`$ resulting probably from infall motions of galaxies and small groups, with small contribution from the DF of galaxies; and (3) the clusters still suffer from the cooling catastrophe. Although higher-resolution, more-controlled simulations are required to make a definitive statement, the last point of their results suggests that the system is thermally unstable or heating by DF may be quite inefficient — for, as we have seen in the previous section, systems evolved from arbitrary initial conditions do not generally tend to the equilibrium solutions, which do not therefore form an attracting set; and for some initial parameters the delay in the cooling catastrophe that the DF mechanism leads to is not significant.
We finally note that several issues of potential importance were not investigated in this paper. One of these relates to the galaxy velocity anisotropy. It is easy to show that, for isothermal distribution with constant anisotropy, it is possible to obtain equilibrium solutions with smaller central gas temperatures, more in line with observations. However these tend to have unrealistic galaxy number density distributions. The equilibrium and stability of gas configurations in more realistic anisotropic models have not been investigated. We have adopted the general formula of Ostriker (1999) for the drag force on a galaxy moving at an arbitrary speed. This is an improvement over previous studies (e.g., Miller 1986; EKK04) that considered only supersonic galaxy motions. While its explicit dependence on the Mach number enabled us to explore temperature profiles of equilibrium clusters, the formula still assumes a galaxy moving in a straight-line orbit through a uniform gaseous medium. Clearly, the ICM is radially stratified and galaxies follow curved rather than rectilinear trajectories. In a collisionless system, Just & Peñarrubia (2004) recently showed that the density inhomogeneity reduces the Coulomb logarithm of the Chandrasekhar formula by limiting the the maximum impact parameter to the local scale length of density variation (see also Hashimoto et al. 2003; Spinnato et al. 2003), and induces an additional drag force in the lateral direction of the galaxy motion (see also Binney 1977). They further showed that the reduction of the Coulomb logarithm makes the orbital energy loss 15% less effective, inhibiting the circularization of the orbit, while the additional tangential drag force has a negligible effect on the orbital evolution. Similar effects are expected to occur in a gaseous medium, yet their consequences are not known quantitatively. A related important question is how heat (or turbulent) energy is distributed in the wakes that form in the turbulent, inhomogeneous, magnetized ICM. A real assessment of the effects of DF to dynamical evolution of galaxy clusters requires answers to the above questions, which will direct our future research.
We are pleased to thank P. Goldreich, M. G. Lee, C. McKee, R. Narayan, and E. Ostriker for many helpful discussions and communications. We also thank an anonymous referee for useful comments and suggestions. W.-T. Kim was supported in part by Korea Science and Engineering Foundation (KOSEF) grant R01-2004-000-10490-0 at the SNU and in part by NSF grant AST 0307433 at the CfA. M. Kamionkowski was supported by NASA NNG05GF69G at Caltech. |
warning/0506/quant-ph0506079.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The quantum entropy and entanglement is a vital feature of quantum information. It has important applications for quantum communication and quantum computation, for example, quantum teleportation, massive parallelism of quantum computation and quantum cryptographic schemes \[1-3\]. Therefore, it is very essential and interesting how to measure the entanglement of quantum states. Recently much attention has been focused on the entanglement of the field and atom when the system starts from a pure state \[4-16\]. Also, in the context of the initial mixed state some studies have been reported \[17-20\]. In this context it was shown that calculating the partial entropies of the field or the atom can be used as an operational measure of the entanglement degree of the generated quantum state. One finds that the higher the entropy, the greater the entanglement. Starting from an initial atom-field product state one can find perfectly entangled states between field and atom at certain later times even for initial coherent states with large photon number \[4-8\]. However, the time evolution of the field (atomic) entropy reflects the time evolution of the degree of entanglement only if one deals with a pure state of the system with zero total entropy.
The quest for proper entanglement measures has received much attention in recent years . From the identification and study of properties of such measures a gain of insight into the nature of entanglement is expected. In turn, their computation for particular states provide us with an account of the resources present in those states. Because one needs to understand the best way to benchmark states for quantum information protocols, here we examine the quantum partial entropies in the atom-field interaction for more general entangled states. From a practical point of view, an implementation of the quantum entropy will be used to measure the entanglement degree when the atom is assumed to be in its excited state and the field initially is in a coherent state, superposition state and a statistical mixture of two coherent states. As far as we are aware in the previous investigations, that have dealt with the present problem, the initial system density matrix is taken to be a product of two states of the factored form. The atom is often taken to be in the excited pure state or mixed state and the radiation field is taken to being a pure state density matrix. So, the present task is a nontrivial issue, since we look at the mixed state entanglement from other direction taking into account the entropy of the field not equal to the entropy of the atom, in this case. To overcome such a difficulty, we employ the quantum von-Neumann entropy to measure the entropy of the atom while a numerical method will be used to calculate the quantum entropy of the field.
The material of this paper is arranged as follows. In section 2, we find the exact solution of the system and write the expressions for the final state vector at any time $`t>0`$. We investigate the quantum field (atomic) entropy and the atom-field entanglement in section 3. Finally, numerical results and conclusions are provided in section 4.
## 2 The model
The system we will consider here consists of a two-level atom interacting with a single-mode quantized field via k-quanta processes. The Hamiltonian in the rotating wave approximation , can be written as $`(\mathrm{}=1)`$:
$$\widehat{H}=\widehat{H}_A+\widehat{H}_F+\widehat{H}_{in},$$
(1)
where
$`\widehat{H}_F`$ $`=`$ $`\omega \widehat{a}^{}\widehat{a},`$
$`\widehat{H}_A`$ $`=`$ $`{\displaystyle \frac{\omega _{}}{2}}\widehat{\sigma }_z,`$
$`\widehat{H}_{in}`$ $`=`$ $`\widehat{a}^{}\widehat{a}(\beta _1|gg|+\beta _2|ee|)+\lambda (\widehat{a}^k\widehat{\sigma }_{}+\widehat{a}^k\widehat{\sigma }_+),`$ (2)
where $`\omega `$ is the field frequency and $`\omega _{}`$ is the transition frequency between the excited and ground states of the atom. We denote by $`\widehat{a}`$ and $`\widehat{a}^{}`$ the annihilation and the creation operators of the cavity field respectively. $`\beta _1`$ and $`\beta _2`$ are parameters describing the intensity-dependent Stark shifts of the two levels that are due to the virtual transition to the intermediate relay level, $`\lambda `$ is the effective coupling constant, $`\widehat{\sigma }_z`$ is the population inversion operator, and $`\widehat{\sigma }_\pm `$ are the ”spin flip” operators, with the detuning parameter $`\mathrm{\Delta }=\omega _{}k\omega `$.
Let us consider, the atom starting in its excited state $`|e`$, i. e.,
$$\rho _0^a=|ee|,$$
(3)
and we are going to assume that the initial single mode electromagnetic field inside the cavity is in a superposition state of the kind:
$$\rho _0^f=\frac{1}{A}\left(|\alpha \alpha |+r^2|\alpha \alpha |+r|\alpha \alpha |+r|\alpha \alpha |\right),$$
(4)
where $`A=\left(1+r^2+2r\mathrm{exp}(2\alpha ^2)\right),`$with $`\alpha `$ real. The parameter $`r`$ takes the values $`1`$, $`0`$ and $`1`$, which corresponds to an odd coherent state, a coherent state and an even coherent state, respectively. While for certain classes of states a superpositions of coherent states, methods solely based on linear optical elements like beam splitters and photodetections could be found , an implementation covering other classes of entangled states remains a challenge.
Also we want to see how different would the behavior of the system be if the input state is a statistical mixture of states $`|\alpha `$ and $`|\alpha `$, i.e.,
$$\rho _0^f=\frac{1}{2}\left(|\alpha \alpha |+|\alpha \alpha |\right).$$
(5)
It is to be noted that when we put $`r=0`$ in equation (4) we get the same result as in Ref. 13. It is expedient to expand the atom-field state in terms of the dressed states:
$`|\mathrm{\Psi }_+^{(n)}`$ $`=`$ $`\mathrm{sin}\theta _n|n,e+\mathrm{cos}\theta _n|n+k,g,`$
$`|\mathrm{\Psi }_{}^{(n)}`$ $`=`$ $`\mathrm{cos}\theta _n|n,e\mathrm{sin}\theta _n|n+k,g,`$ (6)
which are the eigenstates of the interaction Hamiltonian, where
$`\widehat{H}|s,g`$ $`=`$ $`E_0|s,g,0s<k`$
$`\widehat{H}|\mathrm{\Psi }_\pm ^{(n)}`$ $`=`$ $`E_\pm ^{(n)}|\mathrm{\Psi }_\pm ^{(n)},`$ (7)
with the eigenvalues $`E_0`$ and $`E_\pm ^{(n)}`$
$`E_\pm ^{(n)}`$ $`=`$ $`\omega (n+{\displaystyle \frac{k}{2}})+{\displaystyle \frac{\omega _0}{2}}++{\displaystyle \frac{1}{2}}[n\beta _2+\beta _1(n+k)]\pm \mu _n,`$
$`E_0`$ $`=`$ $`\left(s\beta _1{\displaystyle \frac{\mathrm{\Delta }}{2}}\right),`$ (8)
where
$`\mu _n`$ $`=`$ $`\sqrt{\nu _n^2+\tau _n^2},`$
$`\nu _n`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }}{2}}+{\displaystyle \frac{1}{2}}(\beta _2n\beta _1(n+k)),`$
$`\tau _n`$ $`=`$ $`\lambda \sqrt{{\displaystyle \frac{(n+k)!}{n!}}}.`$ (9)
$`\mu _n`$ is a modified Rabi frequency. The angle $`\theta _n`$ is given by
$$\theta _n=\mathrm{sin}^1\left(\frac{\tau _n}{\sqrt{(\nu _n\mu _n)^2+\tau _n^2}}\right).$$
(10)
The unitary operator $`\widehat{U}_t`$ can be written as
$`\widehat{U}_t`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left\{\mathrm{exp}(itE_+^{(n)})|\mathrm{\Psi }_+^{(n)}\mathrm{\Psi }_+^{(n)}|+\mathrm{exp}(itE_{}^{(n)})|\mathrm{\Psi }_{}^{(n)}\mathrm{\Psi }_{}^{(n)}|\right\}`$ (11)
$`+{\displaystyle \underset{s=0}{\overset{k1}{}}}\mathrm{exp}(itE_0)|s,gg,s|.`$
Despite being straightforwardly solvable in this way, the JC-model is well-known for the fact that the time-evolution of most expectation values is usually expressible only in series form. Having obtained the explicit forms of the unitary operator $`\widehat{U}_t`$, for the system under consideration then the eigenvalues and the eigenfunctions can be used to discuss many features concerning the field or the atom.
Bearing these facts in mind we find that the evolution operator $`\widehat{U}_t`$ takes the next from
$$\rho _t=\left(\begin{array}{c}\rho _1\\ \rho _3\end{array}\begin{array}{c}\rho _2\\ \rho _4\end{array}\right),$$
where $`\left(\rho _i\right)_{nm}=n|\rho _i|m,`$ $`i=1,2,3,4.`$ $`\left(\rho _1\right)_{nm}=A_n(t)A_m^{}(t),`$ $`\left(\rho _2\right)_{nm}=A_n(t)B_{mk}^{}(t),`$ $`\left(\rho _3\right)_{nm}=B_{nk}(t)A_m^{}(t),`$ and $`\left(\rho _4\right)_{nm}=B_{nk}(t)B_{mk}^{}(t).`$ The coefficients $`A_n(t)`$ and $`B_n(t)`$ are given by
$`A_n(t)`$ $`=`$ $`q_nC_n\mathrm{exp}[i\lambda t\delta _+(n)]\left(\mathrm{cos}\lambda t\mu _ni\eta _n{\displaystyle \frac{\mathrm{sin}\lambda t\mu _n}{\mu _n}}\right),`$
$`B_n(t)`$ $`=`$ $`iq_nC_n\nu _n\mathrm{exp}[i\lambda t\delta _+(n)]{\displaystyle \frac{\mathrm{sin}\lambda t\mu _n}{\mu _n}},`$
$`R^2`$ $`=`$ $`\sqrt{\beta _1/\beta _2},\eta _n={\displaystyle \frac{\delta }{2}}+\delta _{}(n),\delta ={\displaystyle \frac{\mathrm{\Delta }}{\lambda }},`$
$`\delta _\pm (n)`$ $`=`$ $`\{\begin{array}{c}\frac{1}{2R}[n\pm R^2(n+k)],\\ 0\end{array}\begin{array}{c}whenR0\\ when\beta _i=0,\end{array}`$ (16)
where
$$C_n=[\frac{1}{\sqrt{A}}(1+r(1)^n)],$$
for the initial condition (3), while for the initial condition (4) is
$$C_n=[\frac{1}{\sqrt{2}}(\delta _i+(1)^n\delta _j)],$$
(18)
with $`\delta _i`$, $`\delta _j`$ satisfying the two following condition, $`(a)`$ $`\delta _i=\delta _j=(\delta _i)^2=(\delta _j)^2=1`$, and $`(b)`$ $`\delta _i.\delta _j=0`$,
$$|\alpha =\underset{n=0}{\overset{\mathrm{}}{}}q_n|n=\underset{n=0}{\overset{\mathrm{}}{}}e^{\alpha ^2/2}\frac{\alpha ^n}{\sqrt{n!}}|n.$$
(19)
With the final state obtained, any property related to the atom or the field can be calculated. Employing the reduced density operator for the atom or the field, we investigate the properties of the entropies ($`S_a,S_f`$) and hence entanglement.
## 3 Entropy and subentropy
There is growing interest in the roles of nonadditive measures in quantum information theory. Inadequacy of the additive Shannon von-Neumann entropy as a measure of the information content of a quantum system has been pointed out . Also, there is a theoretical observation that the measure of quantum entanglement may not be additive. Despite the fact that the basic idea of quantum entanglement was acknowledged almost as soon as quantum theory was discovered, it is only in the last few years, that consideration has been given to finding mathematical methods to generally quantify entanglement. In the case of a pure quantum state of two subsystems, a number of widely accepted measures of entanglement are known. However, the question of quantifying the degree of entanglement for general mixed states is still under discussion. Let us now briefly repeat some of the key underlying definitions. The entropy $`S`$ of a quantum-mechanical system described by the density operator $`\widehat{\rho }`$ is defined as follows:
$$S=Tr\{\widehat{\rho }\mathrm{ln}\widehat{\rho }\},$$
(20)
where we have set the Boltzmann constant $`K`$ equal to unity. If $`\widehat{\rho }`$ describes a pure state, then $`S=0`$, and if $`\widehat{\rho }`$ describes a mixed state, then $`S0`$. Entropies of the atomic and field sub-systems are defined by the corresponding reduced density operators:
$$S_{a(f)}=Tr_{a(f)}\{\widehat{\rho }_{a(f)}ln\rho _{a(f)}\}.$$
(21)
Taking the partial trace over the field, the reduced atomic matrix can be written as
$`\rho _t^a`$ $`=`$ $`Tr_f(\rho )=\left(\begin{array}{c}\rho _{ee}\\ \rho _{ge}\end{array}\begin{array}{c}\rho _{eg}\\ \rho _{gg}\end{array}\right)`$ (26)
$`=`$ $`\left(\begin{array}{c}\underset{n=0}{\overset{\mathrm{}}{}}|A_n(t)|^2\\ \\ \underset{n=0}{\overset{\mathrm{}}{}}B_{n+k}(t)A_n^{}(t)\end{array}\begin{array}{c}\underset{n=0}{\overset{\mathrm{}}{}}A_n(t)B_{n+k}^{}(t)\\ \\ \underset{n=0}{\overset{\mathrm{}}{}}|B_{n+k}(t)|^2\end{array}\right).`$ (33)
Thus we rigorously obtain the quantum atomic entropy in the following form
$$S(\rho _t^a)=\lambda _+^a(t)\mathrm{log}\lambda _+^a(t)\lambda _{}^a(t)\mathrm{log}\lambda _{}^a(t),$$
(34)
where $`\lambda _i^a(t)`$ is given by
$$\lambda _\pm ^a(t)=\frac{1}{2}\left\{1\pm \sqrt{\left(2\rho _{ee}(t)1\right)^2+4|\rho _{eg}(t)|^2}\right\}.$$
(35)
In this case, the probability of finding the atom in its excited or ground states are expressed as the diagonal element of the reduced atomic density matrix, i.e.,
$$\rho _{ii}(t)=i|\rho _t^a|i,i=e,g$$
(36)
and the off-diagonal element $`\rho _{eg}(t)`$ is given by
$$\rho _{ij}(t)=i|\rho _t^a|j,i=e,g.$$
(37)
Taking the partial trace over the atomic system, we obtain the reduced density operator in the form
$$\rho _t^f=tr_A\rho (t),$$
(38)
with its ($`\rho _t^f)_{nm}`$ element given by
$$(\rho _t^f)_{nm}=A_n(t)A_m^{}(t)+B_{nk}(t)B_{mk}^{}(t).$$
(39)
From this equation, it is difficult to obtain the eigenvalues of the reduced density operator for the field, in this paper we will evaluate them numerically.
## 4 Results and conclusion
We study the temporal behavior of the atom-field system in the JC-model for the cavity-field prepared initially in different forms. As an example we may consider a simple initial condition for the atom to be in the excited state and the field in a coherent state or a superposition of the coherent state or in a statistical mixture of two coherent states (equations (3) and (4)). Among the family of mixed quantum mechanical states, special status should be accorded to those for a given value of the entropy and have the largest possible degree of entanglement. The reason for this is that such states can be regarded as mixed-state generalizations of Bell states, the latter being known to be the maximally entangled two-qubit pure states. Hence, this kind of mixed states could be expected to provide useful resources for quantum information processing. At this end, we have the plot of the quantum partial entropies $`(S_a,S_f)`$ relative to these different initial states of the atom-field, as a function of the scaled time $`\lambda t/\pi ,`$ taking into account the two-photon process $`(k=2),`$ in order to investigate the Stark shift effects.
We assume a fixed value of the initial mean number of quanta $`\overline{n}=16`$ and different values of Stark shift parameter $`R`$ (namely, $`\beta _1=\beta _2=0,`$ i.e. in the absence of Stark shift in figure 1, $`R=0.5`$ in figure 2 and $`R=0.3`$ for figure 3). Furthermore, the detuning parameter is taken to be zero, and in figure 1$`a`$ we set $`r=0`$ (coherent state), figure 1$`b,`$ we set $`r=1`$ (even coherent state) and figure 1$`c`$ (a statistical mixture state). In the absence of the Stark shift, it is observed that the quantum field entropy and the quantum atomic entropy have the same values due to the initial coherent state $`r=0,`$ (see figure 1a). This behavior is similar to that obtained in the standard two-photon two-level systems obtained previously (see for example \[4-5\]). It is observed that the entropy evolves with a period $`\pi /\lambda `$, when $`t=n\pi /\lambda ,n=0,1,2,\mathrm{})`$, the quantum field entropy evolves to zero and the field is completely disentangled from the atom, while for $`t=(n+\frac{1}{2})\pi /\lambda `$, it evolves to the maximum value, and the field is strongly entangled with the atom.
The situation is completely changed when we consider an even coherent state i.e $`r=1`$. Although the quantum field entropy is still equal the atomic entropy, the quantum entropy in this case has minimum values at $`\lambda t=\frac{\pi n}{2}`$, $`(n=0,1,2,\mathrm{})`$ i.e at half of the revival time, also instead of the two pre-minimum values observed in figure 1a we have here only one pre-minimum value between each two consecutive minima. This effect is due to the interference between the two coherent states in the superposition, and can be understood by looking at the photon number distribution of the initial field. Hence this signal gives a clear measure of the remaining degree of coherence between the two components of the Schrodinger cat state, while the signals present in both cases are due to intrinsic revivals of each component individually. Because of this enhancement it is possible to have the generation of well-defined Schrödinger cat-like states during the evolution of the field in the two-photon process case model . We would like to remark that the approach to a pure state at half of the revival time occurs in the ordinary JC-model , if we start with the field in a pure state.
Let us now come to specific numerical examples to investigate the influence of the statistical mixture on the evolution of the quantum field entropy and the quantum atomic entropy (see figure 1c). An understanding of interaction between an atom and an electromagnetic field has been possible in recent years through the introduction of the statistical mixture state picture . By looking to the statistical mixture state one has a clear physical understanding of what are the parameters involved in such expression and what is going to neglect in order to go from a pure state to a mixed state. On the other words, the state of the initial field is an equally-weighted statistical mixture of two coherent states, which is a special class of the Schrodinger cat state. In this case, as it has been already discussed , the quantum field entropy $`S_f(t)`$ is greater than the quantum atomic entropy $`S_a(t)`$ (see figure 1c), $`S_a(t)`$ reach its maximum values at half of the revival time, while $`S_f(t)`$ evolves to minimum values. We note a little deviation from ordinary Rabi oscillations, due to the statistical mixture case. It is seen that entanglement evolves depending on the initial preparation of atoms, however from figures 1a and 1c we see that the quantum atomic entropy has similar behavior in both superposition and statistical mixture of coherent state.
For comparison purposes, we have chosen to set some different values of the Stark shift parameter R and the other parameters are the same as in figure 1. The outcome is presented in figure 2 (where $`R=0.5`$) and figure 3 (where $`R=0.3`$). We note a stronger modulation in the oscillations and a clear departure from ordinary Rabi oscillations being verified as the Stark shift parameter takes values far from the unity (see figures 2 and 3). Further studies about such a situation have been carried out and are not displayed here. It is interesting to refer here to the fact that the Stark shift creates an effective intensity dependent detuning $`\mathrm{\Delta }_N=\beta _2\beta _1`$ . When $`\beta _2=\beta _1,\mathrm{\Delta }_N=0`$, in this case, the Stark shift does not affect the time evolution of the quantum entropy. As is visible from the figures, the effects of the dynamic Stark shift are more pronounced when $`R`$ deviates from unity. Interestingly, when $`R`$ is decreases, the values of the maximum entropy are decreased. Periodic models therefore may be more robust in this sense. Also, with decreasing the parameter $`R,`$ the evolution period of the entropy as well as the subentropies is decreases (see figure 3). The sensitivity becomes even more clearly visible when we take small velues of the Strak shift parameter. It is worth mentioning that the Strak shift effect has the same impact for both the field entropy and the atomic entropy.
When the atomic entropy is calculated, we note that the terms involved are of the forms $`U(t)\alpha |U(t)\alpha `$ , $`U(t)\alpha |U(t)\alpha `$ and $`|U(t)\alpha |U(t)\alpha |^2`$ for the case of the mixture. These terms in particular do not differs from those of the case of the coherent state which has been discussed earlier. Therefore the temporal evolution of the entropy for the atomic system alone in the case of the Schrödinger cat sate ($`S`$) mimics the evolution of the entropy in the pure coherent state as can be seen from comparing figures (1c) and (1a). However when we consider the entropy for the field we note that terms of the form $`\pm U(t)\alpha |\pm U(t)\alpha `$ with all combinations. These terms are the ones that appear in the case of the superposition of the two states ($`r=1`$ in equation 4). Hence the resemblance between the figures for the entropy of the field in the mixed state of figure (1c) and the case of the initial superposed states of figure (1b). This may demonstrate relevance of investigating the entropies of the subsystems and their relation to entanglement.
In summary, we have shown in this paper that the final analytical expression of the composite density matrix along with its overlap matrix elements can be used to obtain the quantum field entropy and the quantum atomic entropy. This is accomplished by choosing to study the system in the representation in which the marginal initial density matrices are assumed to be in a coherent state, superposition states and statistical mixture states of two coherent states. We present different numerical examples to elucidate the effects of these different settings. Explicit computations are presented for different values of the Stark shift parameter. Our results show that the superposition of coherent states and Stark shift play an important role in the evolution of the quantum entropies in the two-quanta JC-model. In the coherent state, the quantum field (atom) entropy reaches its maximum values at the half of the revival time and the field and the atom are strongly entangled, while in the even coherent state, the entropies at the same time evolve to the minimum values (zero) and the field and the atom are strongly disentangled. In the statistical mixture state, ($`S_aS_f`$), the entropy for the atom reach to the maximum values at half of the revival time, while $`S_f`$ evolves to its minimum values. The significant effect of the Stark shift parameter appears when $`R`$ deviates from unity. The more $`R`$ deviates the more the two systems are weakly entangled. |
warning/0506/gr-qc0506129.html | ar5iv | text | # Quantum evaporation of a naked singularity
## Abstract
We investigate here quantum effects in gravitational collapse of a scalar field model which classically leads to a naked singularity. We show that non-perturbative semi-classical modifications near the singularity, based on loop quantum gravity, give rise to a strong outward flux of energy. This leads to the dissolution of the collapsing cloud before the singularity can form. Quantum gravitational effects thus censor naked singularities by avoiding their formation. Further, quantum gravity induced mass flux has a distinct feature which may lead to a novel observable signature in astrophysical bursts.
preprint: IGPG-05/6-8
Naked singularities are one of the most exotic objects predicted by classical general relativity. Unlike their black hole siblings, they can be in principle directly observed by an external observer. There have been many investigations which show that given the initial density and pressure profiles for a matter cloud, there are classes of collapse evolutions that lead to naked singularity formation (see e.g. review for some recent reviews), subject to an energy condition and astrophysically reasonable equations of state such as dust, perfect fluids and such others. This has led to extensive debates on their existence, with a popular idea being cosmic censorship conjectures which forbid classical nakedness CCC . Since naked singularities originate in the regime where classical general relativity is expected to be replaced by quantum gravity, it has remained an outstanding problem whether a quantum theory of gravity resolves their formation. Also, with the lack of observable signatures from the Planck regime, naked singularities could in fact be a boon for a quantum theory of gravity. Because, the singularity being visible, any quantum gravitational signature originating in the ultra-high curvature regime near a classical singularity can in principle be observed, thus providing us a rare test for quantum gravity.
One of the non-perturbative quantizations of gravity is loop quantum gravity lqg\_review whose key predictions include Bekenstein-Hawking entropy formula bek\_hawking . Its application to symmetry reduced mini-superspace quantization of homogeneous spacetimes is called loop quantum cosmology martin whose success includes resolution of the big bang singularity bigbang , initial conditions for inflation superinflation ; inflationcmb , and possible observable signatures in cosmic microwave background radiation inflationcmb . These techniques have also been applied to resolve black hole singularity in a scalar field collapse scenario bhole .
Since the dynamics of a generic collapse is very complicated and tools to address such a problem in quantum gravity are still under development, it is useful to work with a simple collapse scenario as of a scalar field. It serves as a good toy model to gain insights on the role of quantum gravity effects at the late stages of gravitational collapse. Existence of naked singularities in these models is well-known scalar and one of the simplest setting is to consider an initial configuration of a homogeneous and isotropic scalar field $`\mathrm{\Phi }=\mathrm{\Phi }(t)`$ with a potential $`V(\mathrm{\Phi })`$ (given by eq.(6)) and the canonical momentum $`P_\mathrm{\Phi }`$. In this case it has been shown that fate of the singularity being naked or covered depends on the rate of gravitational collapse JG1 . For an appropriately chosen potential, formation of trapped surfaces can be avoided even as the collapse progresses, resulting in a naked singularity with an outward energy flux, in principle observable. Since the interior of homogeneous scalar field collapse is described by a Friedmann-Robertson-Walker (FRW) metric, techniques of loop quantum cosmology can be used to investigate the way quantum gravity modifies the collapse.
Let us consider the classical collapse of a homogeneous scalar field $`\mathrm{\Phi }(t)`$ with potential $`V(\mathrm{\Phi })`$ and the canonical momentum for the marginally bound $`(k=0)`$ case. The interior metric is given by
$$ds^2=dt^2+a^2(t)\left[dr^2+r^2d\mathrm{\Omega }^2\right]$$
(1)
with classical energy density and pressure of the scalar field,
$$\rho (t)=\dot{\mathrm{\Phi }}^2/2+V(\mathrm{\Phi }),p(t)=\dot{\mathrm{\Phi }}^2/2V(\mathrm{\Phi }).$$
(2)
The dynamical evolution of the system is obtained from the Einstein equations which yield JG1
$$\dot{R}^2R=F(t,r),\rho =F_{,r}/\kappa aR^2,p=\dot{F}/\kappa R^2\dot{R}$$
(3)
Here $`\kappa =8\pi G`$, and $`F(t,r)=(\kappa /3)\rho (t)r^3a^3`$ has interpretation of the mass function of the collapsing cloud, with $`F0`$ and $`R(t,r)=ra(t)`$ is the area radius of a shell labeled by comoving coordinate $`r`$. In a continual collapse the area radius of a shell at a constant value of comoving radius $`r`$ decreases monotonically. The spacetime region is trapped or otherwise, depending on the value of mass function. If $`F`$ is greater (less) than $`R`$, the the region is trapped (untrapped). The boundary of the trapped region is given by $`F=R`$.
The collapsing interior can be matched at some suitable boundary $`r=r_b`$ to a generalized Vaidya exterior geometry, given as wang ,
$$ds^2=(12M(r_v,v)/r_v)dv^22dvdr_v+r_v^2d\mathrm{\Omega }^2.$$
(4)
The Israel-Darmois conditions then lead to JG1 ; wang $`r_ba(t)=r_v(v)`$, $`F(t,r_b)=2M(r_v,v)`$ and
$$M(r_v,v)_{,r_v}=F/2r_ba+r_b^2a\ddot{a}.$$
(5)
The form of the potential that leads to a naked singularity is determined as follows. The energy density of scalar field can be written in a generic form as $`\rho =l^{n4}a^n`$, where $`n>0`$ and $`l`$ is a proportionality constant. Using energy conservation equation, this leads to the pressure $`p=[(n3)/3]l^{n4}a^n`$. On subsituting eq.(2) in these we obtain JG1
$$\mathrm{\Phi }=\sqrt{n/\kappa }\mathrm{ln}a,V(\mathrm{\Phi })=(1n/6)l^{n4}e^{\sqrt{\kappa n}\mathrm{\Phi }}.$$
(6)
Then it is easily seen that $`F/R=(\kappa /3)l^{n4}a^{2n}r^2`$. Thus in the collapsing phase as $`a0`$, whether or not the trapped surfaces form is determined by the value of $`n`$. It is straightforward to check that for $`0<n<2`$, if no trapped surfaces exist initially then no trapped surfaces would form till the epoch $`a(t)=0`$ JG1 , with $`a(t)=\left(1nt/2\sqrt{3}\right)^{2/n}`$.
The absence of trapped surfaces is accompanied by a negative pressure implying that for a constant value of the comoving coordinate $`r`$, $`\dot{F}`$ is negative and so the mass contained in the cloud of that radius keeps decreasing. This leads to a classical outward energy flux. As the collapse proceeds, the scale factor vanishes in finite time and physical densities blow up, leading to a naked singularity. Since no trapped surfaces form during collapse, the outward energy flux shall in principle be observable. However, near the singularity when energy density is close to Planckian values, this classical picture has to be modified and we need to investigate the scenario incorporating quantum gravity modifications to the classical dynamics.
Let us hence consider the non-perturbative semi-classical modifications based on loop quantum gravity for the interior. The underlying geometry for the FRW spacetime in loop quantum cosmology is discrete and both the scale factor and the inverse scale factor operators have discrete eigenvalues Bohr . In particular, there exists a critical scale $`a_{}=\sqrt{j\gamma /3}\mathrm{}_\mathrm{P}`$ below which the eigenvalues of the inverse scale factor become proportional to the positive powers of scale factor. Here $`\gamma 0.2375`$ is the Barbero-Immirzi parameter bek\_hawking , $`\mathrm{}_\mathrm{P}`$ is Planck length and $`j`$ is a half-integer free parameter which arises because inverse scale factor operator is computed by tracing over SU(2) holonomies in an irreducible spin $`j`$ representation. The value of this parameter is arbitrary and shall be constrained by phenomenological considerations.
The change in behavior of the classical geometrical density ($`1/a^3`$) for scales $`aa_{}`$, can be well approximated by superinflation
$$d_j(a)=D(q)a^3,q:=a^2/a_{}^2,a_{}:=\sqrt{j\gamma /3}\mathrm{}_\mathrm{P}$$
(7)
with
$`D(q)=(8/77)^6q^{3/2}\{7[(q+1)^{11/4}|q1|^{11/4}]`$
$`11q[(q+1)^{7/4}\mathrm{sgn}(q1)|q1|^{7/4}]\}^6.`$ (8)
For $`aa_{}`$, $`d_j(a/a_{})^{15}a^3`$ and for $`aa_{}`$ it behaves classically with $`d_ja^3`$. The scale at which transition in the behavior of the geometrical density takes place is determined by the parameter $`j`$.
At the fundamental level the dynamics in the loop quantum regime is discrete, however, recent invsetigations pertaining to the evolution of coherent states have shown that for scales $`a_0=\sqrt{\gamma }\mathrm{}_\mathrm{P}aa_{}=\sqrt{j\gamma /3}\mathrm{}_\mathrm{P}`$, dynamics can be described by modifications to Friedmann dynamics on a continuous spacetime time with the modified matter Hamiltonian
$$_\mathrm{\Phi }=d_j(a)P_\mathrm{\Phi }^2/2+a^3V(\mathrm{\Phi })$$
(9)
and the modified Friedmann equation
$$\dot{a}^2/a^2=(\kappa /3)(\dot{\mathrm{\Phi }}^2/2D+V(\mathrm{\Phi }))$$
(10)
which is obtained by the vanishing of the total Hamiltonian constraint and the Hamilton’s equations: $`\dot{\mathrm{\Phi }}=d_j(a)P_\mathrm{\Phi },\dot{P}_\mathrm{\Phi }=a^3V_{,\mathrm{\Phi }}(\mathrm{\Phi })`$ superinflation . These also lead to the modified Klein-Gordon equation
$$\ddot{\mathrm{\Phi }}+\left(3\dot{a}/a\dot{D}(q)/D(q)\right)\dot{\mathrm{\Phi }}+D(q)V_{,\mathrm{\Phi }}(\mathrm{\Phi })=0.$$
(11)
Since at classical scales ($`aa_{}`$) $`D1`$, the modified dynamical equations reduce to the standard Friedmann dynamical equations. For scales $`aa_{}`$, the $`\dot{\mathrm{\Phi }}`$ term acts like a frictional term for a collapsing phase. We note that since semi-classical modifications for inhomogeneous case are still not known, we cannot do a complete quantum analysis of interior and exterior. The exterior is assumed to remain classical. Further, as a continuous spacetime can be approximated till scale factor $`a_0`$, the matching of interior and exterior spacetimes remains valid during the semi-classical evolution.
The modified energy density and pressure of the scalar field in the semi-classical regime can be similarly obtained from the eigenvalues of density operator and using the stress-energy conservation equation density
$$\rho _{\mathrm{eff}}=d_j(a)_\mathrm{\Phi }=\dot{\mathrm{\Phi }}^2/2+D(q)V(\mathrm{\Phi })$$
(12)
and
$$p_{\mathrm{eff}}=\left[1\frac{2}{3}\frac{1}{(\dot{a}/a)}\frac{\dot{D}(q)}{D(q)}\right]\frac{\dot{\mathrm{\Phi }}^2}{2}D(q)V(\mathrm{\Phi })\frac{\dot{D}(q)}{3(\dot{a}/a)}V(\mathrm{\Phi }).$$
(13)
It is then straightforward to check that $`p_{\mathrm{eff}}`$ is generically negative for $`aa_{}`$ and for $`aa_{}`$ it becomes very strong. For example, at $`aa_0`$, $`p_{\mathrm{eff}}9\rho _{\mathrm{eff}}`$. This is much stronger than its classical counterpart $`p=[(n3)/3]\rho `$ with $`0<n<2`$. Thus we expect a strong burst of outward energy flux in the semi-classical regime. Further, for $`aa_{}`$, $`D(q)1`$ and the Klein-Gordon equation yields $`\dot{\mathrm{\Phi }}a^{12}`$. Hence from the eq. (12) we easily see that the effective density, instead of blowing up, becomes extremely small and remains finite.
The modified mass function of the collapsing cloud can be evaluated using eq.(3) and eq.(10),
$$F=(\kappa /3)(d_j^1\dot{\mathrm{\Phi }}^2/2+a^3V(\mathrm{\Phi }))r^3.$$
(14)
In the regime $`aa_0`$, $`d_j^1\dot{\mathrm{\Phi }}^2`$ becomes proportional to $`a^{12}`$, the potential term becomes negligible and thus the mass function becomes vanishingly small at small scale factors.
The picture emerging from loop quantum modifications to collapse is thus following.
$``$ Before the area radius of the collapsing shell reaches $`R_{}=ra_{}`$ at $`t=t_{}`$, collapse proceeds as per classical dynamics and as smaller scale factors are approached $`\dot{\mathrm{\Phi }}`$ and the energy density $`\rho a^n`$ increase. The mass function is proportional to $`a^{n3}`$ and (as $`0<n<2`$) it decreases with decreasing scale factor so there is a mass loss to the exterior, which is also understood from existence of negative classical pressure.
$``$ As the collapsing cloud reaches $`R_{}`$, the geometric density classically given by $`a^3`$, modifies to $`d_j`$ and the dynamics is governed by the modified Friedmann and Klein-Gordon equations. The scalar field which experienced anti-friction in classical regime, now experiences friction leading to decrease of $`\dot{\mathrm{\Phi }}`$.
$``$ The slowing down of $`\mathrm{\Phi }`$ decreases the rate of collapse and formation of singularity is delayed. Eventually when scale factor becomes smaller than $`a_0`$ this leads to breakdown of continuum spacetime approximation and semi-classical dynamics. Discrete quantum geometry emerges at this scale time and the dynamics can only be described by quantum difference equation. The naked singularity is thus avoided till the scale factor at which a continuous spacetime exists.
We show the evolution of area radius in time as collapse proceeds in Fig.1. The semi-classical evolution (solid curve) closely follows classical trajectory (dashed) till the time $`t_{}`$. Within a finite time after $`t_{}`$, the classical collapse leads to a vanishing $`R`$ and naked singularity. However, the area radius never vanishes in the loop modified semi-classical dynamics and the naked singularity does not form as long as the continuum spacetime approximation holds. The inset of Fig.1 shows the evolution of energy density in Planck units. Classical energy density (dashed curve) blows up whereas it remains finite and in fact decreases in the semi-classical regime.
The phenomena of delay and avoidance of the naked singularity in continuous spacetime is accompanied by a burst of matter to the exterior. If the mass function at scales $`aa_{}`$ is $`F_i`$ and its difference with mass of the cloud for $`a<a_{}`$ is $`\mathrm{\Delta }F=F_iF`$, then the mass loss can be computed as
$$\frac{\mathrm{\Delta }F}{F(a_i)}=\left[1\frac{\rho _{\mathrm{eff}}d_j^1}{l^{n4}a_i^{3n}}\right].$$
(15)
For $`a<a_{}`$, as the scale factor decreases, the energy density and mass in the interior decrease and the negative pressure strongly increases. This leads to a strong burst of matter. The absence of trapped surfaces enables the quantum gravity induced burst to propagate via the generalized Vaidya exterior to an observer at infinity. The evolution of mass function is shown in Fig.2. In the semi-classical regime, $`\mathrm{\Delta }F/F_i`$ approaches unity very rapidly. This feature is independent of the choice of parameter $`j`$. The choice of potential causes mass loss to exterior in classical collapse also, but it is much smaller and in any case the classical description cannot be trusted at energy density greater than Planck, when we must consider quantum effects as above.
Interestingly, for a given collapsing configuration, the scale at which the strong outward flux initiates depends on the loop parameter $`j`$ which controls $`a_{}`$. If $`j`$ is large then burst occurs at an earlier area radius and vice versa. The inset of Fig.2 shows the mass loss ratio for different values of $`j`$. For all choices, $`\mathrm{\Delta }F/F_i1`$, but the outgoing flux profile changes. The loop quantum burst has a distinct signature, at $`aa_{}`$ the flux decreases for a short period and then rapidly increases. Since the causal structure of classical spacetime is such that trapped surface formation is avoided, this quantum gravitational signature can be in principle observed by an external observer as a slight dimming and subsequent brightening of the collapsing star. This peculiar phenomena is directly related to the peak in the function $`d_j(a)`$, and depends solely on the value of parameter $`j`$. If we compare this to other phenomenological applications superinflation ; inflationcmb ; bhole , this effect could not be masked by the role of other loop quantum parameters in a more general setting. This phenomena is thus a direct probe to measure $`j`$ and an observer can estimate the loop quantum parameter $`j`$ by observing the flux profile of the burst based on this mechanism and measuring the variation in luminosity of the collapsing cloud.
During such a burst most of the mass is ejected and this may dissolve the singularity. Thus non-perturbative semi-classical modifications may not allow formation of naked singularity as the collapsing cloud evaporates away due to super-negative pressures in the late regime. It has been demonstrated that these super-negative pressures would exist for arbitrary matter configurations density which implies that results obtained here would hold even in a more general setting JG . Loop quantum effects then imply a quantum gravitational cosmic censorship, alleviating the naked singularity problem. We note that the semi-classical effects do not show that the singularity is absent, it is only avoided till scale factor $`a_0`$, below which the semi-classical dynamics and matching may break down. If for a given choice of initial data, semi-classical dynamics is unable to completely dissolve the singularity, the final fate of naked singularity must be decided by using full quantum evolution. Even in such cases we have valuable insights from semi-classical loop quantum effects with the possibility of phenomenologically constraining the $`j`$ parameter.
In the toy model considered, we showed that the classical outcome and evolution of collapse is radically altered by the non-perturbative modifications to the dynamics. Our considerations are of course within the mini-superspace setting, and the general case of inhomogeneities and anisotropies remains open. However, the possibility of such observable signatures in astrophysical bursts, as originating from quantum gravity regime near singularity is intriguing, indicating that gravitational collapse scenario can be used as probes to test quantum gravity models.
Acknowledgments: We thank A. Ashtekar, M. Bojowald and R. Maartens for useful comments. PSJ thanks Bharat Kapadia for various discussions. PS is supported by Eberly research funds of Penn State and NSF grant PHY-00-90091. |
warning/0506/hep-ph0506057.html | ar5iv | text | # Five-dimensional Yang-Mills-Einstein supergravity on orbifolds: Parity assignments
## I Introduction
Phenomenological field-theoretic model building has recently refocused on scenarios in which the universe appears higher-dimensional above some energy scale to examine new electroweak and supersymmetry breaking schemes ewsusy and new strong-electroweak unification scenarios guts . In Freund:1982pf , it was pointed out that the estimated scale of strong-electroweak unification was around the energy scale where a Kaluza-Klein type universe may not be able to be approximated by a 4D theory, in which case grand unification would occur in higher dimensions. In Ant , it was suggested that the size of an extra dimension could be much larger (TeV scale) within the framework of perturbative string theory (one of the motivations was to tie this scale to the $`𝒩=1`$ supersymmetry breaking scale). Subsequently, the Horava-Witten (HW) scenario HW and Randall-Sundrum (RS) scenarios RS served as the most recent revival of interest in intermediate scale five-dimensional spacetimes in which the ground state spacetime is a singular space that is isomorphic to a spacetime with boundaries. In particular, the “orbifold” spacetime $`_4\times S^1/\mathrm{\Gamma }`$ (upstairs picture) corresponds to a manifold with two boundaries separated by a 5D bulk (downstairs picture). This spacetime can resolve a number of issues in the supersymmetric Standard Model and supersymmetric Grand Unified Theories (GUTs).
Since an orbifold spacetime is singular, field theories are not well-defined on it, which requires some further interpretation (in the downstairs picture, the boundary is sharp). Supergravity admits solitonic solutions that could ultimately be interpreted as the boundaries of these theories; these solutions are domain walls with some thickness, smoothing out the singular nature of sharp boundaries. Supergravity in orbifold scenarios is interesting for other reasons: as the familiar argument goes, the gauge couplings in supersymmetric versions of the Standard Model unify sin2:91 at a scale close to the scale at which quantum gravitational effects are expected to be non-negligible. Since there is a significant interpolation involved in these suggestive results, one might look for unification of gauge and gravitational couplings. Therefore, the next step in bottom-up model building is to consider supergravity versions of the interesting scenarios.
Orbifold constructions have been performed many times in the literature for both rigidly and locally supersymmetric field theories; for examples of the former, see ewsusy ; guts ; for the latter, see JB ; offshell ; YL:03dec ; ZGAZ:04jul ; AK04 . The generic result is a theory with 4D boundaries with 1/2 supersymmetry and broken gauge group; the low energy description is a 4D effective theory for each boundary. However, a systematic classification of the types of boundary conditions available via parity assignments has not been performed for Yang-Mills-Einstein supergravity theories (YMESGTs) coupled to vector, hyper-, and tensor multiplets (usually, orbifold supergravity theories are considered in the context of HW or (original) RS type scenarios in which the Standard Model gauge, and sometimes matter, fields are supported only on the boundaries). In this paper, we aim to provide a more complete list of options, including a careful look at the tensor sector, for the low energy spectrum via parity assignments in the simple case of the $`S^1/_2`$ orbifold. We do not presume boundary-localized field content, which would follow from the appropriate coupling of 3-branes on the orbifold fixed-planes (or other massless fields arising in an M-theoretic framework). We will then extend our results to the case $`S^1/(_2\times _2)`$. Some of the results are generic to theories with super-Yang-Mills coupled to hypermultiplets, while others are unique to supergravity. While orbifold-GUTs are a main motivation, the results are not restricted to these scenarios. As a novel example of 5D GUT in the framework of supergravity, we illustrate some of the parity assignments using an $`SU(5,1)`$ gauging.
In the next section, we will make generic remarks about field theory on orbifolds. In section III, we list the $`𝒩=1`$ supermultiplets that have propagating modes on the orbifold fixed planes of the spacetime $`M_4\times S^1/_2`$, as dictated by consistent $`_2`$-parity assignments. We first discuss the pure Yang-Mills-Einstein and hypermultiplet sectors in turn. The tensor sector, discussed in section III.3, is subtler and has not been discussed in the literature much. In section IV, we list the $`_2`$-parity assignments for objects other than fields appearing in the Lagrangian; in particular, spacetime- and field-independent quantities may be assigned odd parities, without the need to assume that there is a formulation of the theory in which an odd-parity field is responsible for this. In section VI, we extend the results to the case in which the fifth dimension is $`S^1/(_2\times _2)`$, which can be more phenomenologically interesting. We end with conclusions (section VII) and some future directions (section VIII). Appendix A contains some conventions used in this paper, while appendix B contains a discussion of the fermionic parity assignments.
## II Supergravity on $`S^1/\mathrm{\Gamma }`$
In modeling five-dimensional spacetimes with four-dimensional boundaries, we can choose a particular construction using a spacetime of the form $`M_4\times S^1/\mathrm{\Gamma }`$, where $`\mathrm{\Gamma }`$ is a discrete group that acts non-freely on the circle <sup>1</sup><sup>1</sup>1We’ll refer to this as an orbifold spacetime, even though orbifolds are usually defined as spaces with singularities of codimension greater than one.. The orbifold $`_4\times S^1/\mathrm{\Gamma }`$, which has two 4D fixed planes, is isomorphic to a groundstate spacetime $`_4\times `$, where $``$ is a closed interval so that this is a spacetime with two disjoint 4D boundaries. Instead of considering the 5D theories with these spacetime boundaries (the downstairs picture), it is often convenient to compactify the 5D theory on $`S^1`$, followed by assignment of $`\mathrm{\Gamma }`$-parities to quantities in the theory (the upstairs picture). Since the Lagrangian (density) must be even under $`_2`$, there are constraints on the relative parities of the fields and other objects. Further constraints are imposed by the consistency of local coordinate transformations, supersymmetry transformations and gauge transformations. The choice of $`\mathrm{\Gamma }`$ reflects different classes of boundary conditions from the downstairs point of view. We will first consider the simplest case $`\mathrm{\Gamma }=_2`$, which results in a theory with equivalent spectra and interactions at the two $`\mathrm{\Gamma }`$ fixed points.
The choice of the way $`_2`$ acts on quantities in the theory reflects a particular set of consistent boundary conditions. The $`_2`$ acts as reflections on the $`S^1`$ covering space $`[\pi R,\pi R]`$ (where $`\{\pi R\}\{\pi R\}`$), with fixed points at $`\{0\},\{\pi R\}`$. To leave the space $`_4\times S^1/_2`$ invariant under this action, the coordinate functions, basis vectors, basis 1-forms, and metric components have
$$P(x^\mu ;_\mu ;dx^\mu )=+1P(x^5;_5;dx^5)=1$$
$$P(\widehat{g}_{\mu \nu };\widehat{g}_{\mathrm{\hspace{0.17em}55}})=+1P(\widehat{g}_{\mu 5})=1,$$
where $`P(\mathrm{\Phi })`$ denotes the $`_2`$ parity of the object $`\mathrm{\Phi }`$. The fixed planes are non-oriented.
Fields carrying internal quantum numbers are sections of a fiber bundle, with spacetime being the base space. In such a situation, it makes sense for the action of $`_2`$ to be lifted from the base space to the total space orbifolds . There are a number of ways to perform this lift, corresponding to various classes of boundary conditions. Just as the $`_2`$ action on the covering space $`S^1`$ results in a singular space $`S^1/_2`$, the $`_2`$ action on the total space will, in general, change the structure of the fibers over the base space.
Objects other than fields appearing in the Lagrangian generally carry representation indices of the gauge group. Such quantities are structures appearing in the gauge bundle, and are therefore generally affected by modifications of the gauge bundle resulting after $`_2`$ action (even if they are field independent). This gives meaning to assigning these objects $`_2`$ parities.
Although physical states on $`^4\times S^1/_2`$ must be even under $`_2`$-action, the field operators can carry even or odd parity. We will take the expansion of an odd parity field on the orbifold to have an $`n`$th term of the form
$$\begin{array}{cc}\hfill \mathrm{\Phi }_{}^{(n)}(x^\mu ,x^5)=& A_n\mathrm{\Phi }_{}^{(n)}(x^\mu )\mathrm{sin}(nx^5/R)\hfill \\ & +B_n\mathrm{\Phi }_{}^{(n)}(x^\mu )ϵ(x^5)\mathrm{cos}(nx^5/R),\hfill \end{array}$$
(1)
where $`ϵ(x^5)`$ is $`+1`$ for $`(\pi R,0)`$ and $`1`$ for $`(0,\pi R)`$; the $`\mathrm{\Phi }_{}^{(n)}(x^\mu )`$ are even; and $`A_n,B_n`$ are normalization factors.
To avoid $`\delta ^{\mathrm{\hspace{0.17em}2}}`$ factors in the Lagrangian ($`\delta `$ being the Dirac distribution), bosonic fields cannot have $`ϵ(x^5)`$ factors in the above expansion. Therefore, we impose the condition $`B_n=0`$ for odd bosonic fields, which therefore vanish on the orbifold fixed planes.
On the other hand, the equations of motion and kinetic terms for fermionic fields involve a single derivative, so $`ϵ(x_5)`$ factors are allowed (they will give rise to $`\delta (x_5)`$ factors in the equations of motion and upstairs picture Lagrangian). Therefore, fermionic fields on $`S^1/_2`$ are not well-defined on the upstairs picture fixed planes.
Odd-parity field-independent objects $`C_{J_1\mathrm{}J_n}^{I_1\mathrm{}I_n}`$ carrying gauge indices can generally be redefined as either $`ϵ(x^5)C_{J_1\mathrm{}J_n}^{I_1\mathrm{}I_n}`$ or $`\kappa (x^5)C_{J_1\mathrm{}J_n}^{I_1\mathrm{}I_n}`$, where $`C_{J_1\mathrm{}J_n}^{I_1\mathrm{}I_n}`$ now has even parity, and $`\kappa (x^5)`$ is
$$\kappa (x^5)=\{\begin{array}{cc}0& \text{for }x^5=\pi R\hfill \\ 1& \text{for }\pi R<x^5<0\hfill \\ 0& \text{for }x^5=0\hfill \\ +1& \text{for }\mathrm{\hspace{0.33em}0}<x^5<\pi R\hfill \end{array}$$
(2)
However, consistency may require one or the other.
## III Supermultiplets appearing in $`\mathrm{\Gamma }=_2`$ case
### III.1 Pure Yang-Mills-Einstein supergravity
An $`𝒩=2`$ 5D Maxwell-Einstein supergravity theory (MESGT) MESGT consists of a minimal supergravity multiplet and $`n_V`$ vector supermultiplets. The total field content is
$$\{e_{\widehat{\mu }}^{\widehat{m}},\mathrm{\Psi }_{\widehat{\mu }}^i,A_{\widehat{\mu }}^I,\lambda ^{i\stackrel{~}{p}},\varphi ^{\stackrel{~}{x}}\},$$
where $`I=(0,1,\mathrm{},n_V)`$ labels the graviphoton and vector fields from the $`n_V`$ vector multiplets; $`i=(1,2)`$ is an $`SU(2)_R`$ index; and $`\stackrel{~}{p}=(1,\mathrm{},n_V)`$ and $`\stackrel{~}{x}=(1,\mathrm{},n_V)`$ label the fermions and scalars from the $`n_V`$ vector multiplets. The scalar fields parametrize an $`n_V`$-dimensional real Riemannian manifold $`_R`$, so the indices $`\stackrel{~}{p},\stackrel{~}{q},\mathrm{}`$ and $`\stackrel{~}{x},\stackrel{~}{y},\mathrm{}`$ may also be viewed as flat and curved indices of $`_R`$, respectively. The supersymmetry parameters $`ϵ^i`$, the gravitini $`\mathrm{\Psi }_\mu ^i`$, and the spin-1/2 fields $`\lambda ^{\stackrel{~}{p}i}`$ are 5D symplectic-Majorana spinors (see appendix A), which can be written in 2-component spinor notation as
$$ϵ^1=\left(\begin{array}{c}\eta \\ e\zeta ^{}\end{array}\right)ϵ^2=\left(\begin{array}{c}\zeta \\ e\eta ^{}\end{array}\right)$$
$$\mathrm{\Psi }_\mu ^1=\left(\begin{array}{c}\alpha _\mu \\ e\beta _\mu ^{}\end{array}\right)\mathrm{\Psi }_\mu ^2=\left(\begin{array}{c}\beta _\mu \\ e\alpha _\mu ^{}\end{array}\right)$$
(3)
$$\lambda ^{\stackrel{~}{p}\mathrm{\hspace{0.33em}1}}=\left(\begin{array}{c}\delta ^{\stackrel{~}{p}}\\ e\gamma ^{\stackrel{~}{p}}\end{array}\right)\lambda ^{\stackrel{~}{p}\mathrm{\hspace{0.33em}2}}=\left(\begin{array}{c}\gamma ^{\stackrel{~}{p}}\\ e\delta ^{\stackrel{~}{p}}\end{array}\right).$$
Introducing $`(n_V+1)`$ parameters $`\xi ^I(\varphi )`$ depending on the scalar fields, we define a cubic polynomial $`𝒱(\xi )=C_{IJK}\xi ^I\xi ^J\xi ^K`$, where $`C_{IJK}`$ is a constant rank-3 symmetric tensor. This polynomial determines a symmetric rank-2 tensor
$$a_{IJ}(\xi )=\frac{1}{3}\frac{}{\xi ^I}\frac{}{\xi ^J}\mathrm{ln}𝒱(\xi ).$$
The parameters $`\xi ^I`$ can be interpreted as coordinate functions for an $`(n_V+1)`$-manifold, called the ambient space. The tensor $`a_{IJ}`$, which may have indefinite signature, defines a metric on this space. However, the coordinates are restricted via $`𝒱(\xi )>0`$ so that the metric is positive definite, which means that the manifold is Riemannian. The equation $`𝒱(\xi )=k`$ ($`k^+`$) defines a family of real hypersurfaces, and in particular
$$𝒱(\xi )=1$$
(4)
defines a real $`n_V`$-manifold corresponding to the scalar manifold $`_R`$. The functions $`h^I`$ and $`h_{\stackrel{~}{x}}^I`$ that appear in fermionic terms of the Lagrangian and susy transformations are directly proportional to $`\xi ^I|_{𝒱=1}`$ and $`\xi _{,\stackrel{~}{x}}^I|_{𝒱=1}`$, respectively; the $`h^I`$ are essentially embedding coordinates of $`_R`$ in the ambient space. In MESGT , it was shown that, when (4) holds, the $`C_{IJK}`$ may be put in the “canonical form”
$$\begin{array}{cc}& C_{000}=1,C_{0ij}=\frac{1}{2}\delta _{ij},\hfill \\ & C_{00i}=0,C_{ijk}=\text{arbitrary},\hfill \end{array}$$
(5)
where $`I=(0,i)`$, $`i=1,\mathrm{},n_V`$. We denote the restriction of the ambient space metric to $`_R`$ as
$$\underset{IJ}{\overset{}{a}}=a_{IJ}|_{𝒱=1}.$$
The metric of the scalar manifold is then the pullback of the restricted ambient space metric to $`_R`$:
$$g_{\stackrel{~}{x}\stackrel{~}{y}}=\frac{3}{2}\underset{IJ}{\overset{}{a}}h_{,\stackrel{~}{x}}^Ih_{,\stackrel{~}{y}}^J.$$
Both of these metrics are positive definite due to the constraint $`𝒱>0`$.
The $`C_{IJK}`$ tensor completely determines the MESGT Lagrangian MESGT . Therefore, the global symmetry group of the Lagrangian is given by the symmetry group, $`G`$, of this tensor, along with automorphisms of the $`𝒩=2`$ superalgebra; that is, the symmetry group of the Lagrangian is $`G\times SU(2)_R`$. Since $`G`$ consists of symmetries of the full Lagrangian, they are symmetries of the scalar sector in particular, and therefore isometries of the scalar manifold $`_R`$: $`GIso(_R)`$ (the $`SU(2)_R`$ action is trivial on the scalars). The full Lagrangian, however, is not necessarily invariant under the full group $`Iso(_R)`$.
A subgroup $`KG`$ may then be gauged if $`n_V+1\text{dim}[K]`$. The vector fields decompose into
$$𝐧_𝐕+\mathrm{𝟏}=\text{adj}(K)\text{non-singlets}(K)\text{singlets}(K).$$
Under infinitesimal $`K`$-transformations parametrized by $`\alpha ^I`$, the bosonic fields transform as:
$$\begin{array}{cc}\hfill \delta _\alpha A_{\widehat{\mu }}^I& =\frac{1}{g}_{\widehat{\mu }}\alpha ^I+\alpha ^Jf_{JK}^IA_{\widehat{\mu }}^K\hfill \\ \hfill \delta _\alpha \varphi ^{\stackrel{~}{x}}& =\alpha ^IK_I^{\stackrel{~}{x}},\hfill \end{array}$$
(6)
where $`K_I^{\stackrel{~}{x}}`$ are a set of $`n_V`$ Killing vectors on the scalar manifold parametrized by the $`\varphi ^{\stackrel{~}{x}}`$; and $`f_{JK}^I`$ and $`\alpha ^I`$ vanish if any index corresponds to a spectator vector field <sup>2</sup><sup>2</sup>2If we gauge a compact group, there is always at least one singlet (spectator) vector field, which can be identified as the graviphoton..
Now the $`C_{IJK}`$ must be a rank-3 symmetric invariant of $`K`$. If $`K`$ is compact, then only $`C_{ijk}`$ (see (5)) must be a rank-3 symmetric invariant. The 5D bosonic YMESGT Lagrangian is MESGT
$$\begin{array}{cc}& \widehat{e}^1_{bos}=\hfill \\ & \frac{1}{2}\widehat{R}\frac{1}{4}\underset{IJ}{\overset{}{a}}_{\widehat{\mu }\widehat{\nu }}^I^{J\widehat{\mu }\widehat{\nu }}\frac{1}{2}g_{\stackrel{~}{x}\stackrel{~}{y}}D_{\widehat{\mu }}\varphi ^{\stackrel{~}{x}}D^{\widehat{\mu }}\varphi ^{\stackrel{~}{y}}\hfill \\ & +\frac{\widehat{e}^1}{6\sqrt{6}}C_{IJK}ϵ^{\widehat{\mu }\widehat{\nu }\widehat{\rho }\widehat{\sigma }\widehat{\lambda }}\{F_{\widehat{\mu }\widehat{\nu }}^IF_{\widehat{\rho }\widehat{\sigma }}^JA_{\widehat{\lambda }}^K\hfill \\ & +\frac{3}{2}gF_{\widehat{\mu }\widehat{\nu }}^IA_{\widehat{\rho }}^J(f_{LM}^KA_{\widehat{\sigma }}^LA_{\widehat{\lambda }}^M)\hfill \\ & +\frac{3}{5}g^2(f_{GH}^JA_{\widehat{\nu }}^GA_{\widehat{\rho }}^H)(f_{LF}^KA_{\widehat{\sigma }}^LA_{\widehat{\lambda }}^F)A_{\widehat{\mu }}^I\},\hfill \end{array}$$
(7)
where hats indicate five-dimensional quantities and $`\widehat{e}`$ is the determinant of the fünfbein. The $`A_\mu ^I`$ are abelian vector fields and <sup>3</sup><sup>3</sup>3$`A_{[\alpha }B_{\beta ]}\frac{1}{2}(A_\alpha B_\beta A_\beta B_\alpha )`$
$$\begin{array}{cc}\hfill F_{\widehat{\mu }\widehat{\nu }}^I& 2_{[\widehat{\mu }}A_{\widehat{\nu }]}^I\hfill \\ \hfill _{\widehat{\mu }\widehat{\nu }}^I& F_{\widehat{\mu }\widehat{\nu }}^I+2gA_{[\widehat{\mu }}^IA_{\widehat{\nu }]}^J\hfill \\ \hfill D_{\widehat{\mu }}\varphi ^{\stackrel{~}{x}}& _{\widehat{\mu }}\varphi ^{\stackrel{~}{x}}+gK_I^{\stackrel{~}{x}}A_{\widehat{\mu }}^I.\hfill \end{array}$$
The supersymmetry transformations of the bosonic fields are
$$\begin{array}{cc}\hfill \delta \widehat{e}_{\widehat{\mu }}^{\widehat{m}}=& \frac{1}{2}\overline{ϵ}^i\mathrm{\Gamma }^m\mathrm{\Psi }_{\mu i}\hfill \\ \hfill \delta A_{\widehat{\mu }}^I=& \frac{1}{2}h_{\stackrel{~}{p}}^I\overline{ϵ}^i\mathrm{\Gamma }_\mu \lambda _i^{\stackrel{~}{p}}+\frac{i\sqrt{6}}{4}h^I\widehat{\overline{\mathrm{\Psi }}}_{\widehat{\mu }}^iϵ_i\hfill \\ \hfill \delta \varphi ^{\stackrel{~}{x}}=& \frac{i}{2}f_{\stackrel{~}{p}}^{\stackrel{~}{x}}\overline{ϵ}^i\lambda _i^{\stackrel{~}{p}},\hfill \end{array}$$
(8)
where the $`f_{\stackrel{~}{p}}^{\stackrel{~}{x}}`$ are $`n_V`$-bein of the real scalar manifold $`_V`$.
Dimensional reduction of 5D $`𝒩=2`$ YMESGT
In the “upstairs” orbifold construction, one starts with a 5D theory, and compactifies on $`S^1`$. It is sufficient for our purposes to use the dimensionally reduced theory, consisting of those fields satisfying $`_5\mathrm{\Phi }=0`$. The dimensional reduction breaks the 5D local Lorentz invariance to a 4D local Lorentz invariance. The four local symmetries that are broken can be used to fix four degrees of freedom in the fünfbein. Splitting $`\widehat{\mu }=(\mu ,5)`$, we choose the parametrization for the fünfbein to be MESGT
$$\widehat{e}_{\widehat{\mu }}^{\widehat{m}}=\left(\begin{array}{ccc}e^{\frac{\sigma }{2}}e_\mu ^m& & 2e^\sigma C_\mu \\ 0& & e^\sigma \end{array}\right).$$
Since $`\widehat{g}_{\widehat{\mu }\widehat{\nu }}=\widehat{e}_{\widehat{\mu }}^{\widehat{m}}\widehat{e}_{\widehat{\nu }}^{\widehat{n}}\eta _{\widehat{m}\widehat{n}}`$,
$$\begin{array}{cc}\hfill \widehat{g}_{\mu \nu }=& e^\sigma g_{\mu \nu }+4e^{2\sigma }C_\mu C_\nu \hfill \\ \hfill \widehat{g}_{\mathrm{5\hspace{0.17em}5}}=& e^{2\sigma }\hfill \\ \hfill \widehat{g}_{\mu \mathrm{\hspace{0.17em}5}}=& 2e^{2\sigma }C_\mu .\hfill \end{array}$$
(9)
Furthermore, let $`A_{\widehat{\mu }}^I=(A_\mu ^I,A^I).`$ Under infinitesimal local coordinate transformations of the compact coordinate parameterized by $`\xi ^5(x^\mu )`$, the 4D fields $`A_\mu ^I`$ and $`C_\mu `$ transform as
$$\begin{array}{cc}\hfill \delta _{\xi ^5}A_\mu ^I& =_\mu \xi ^5A^I\hfill \\ \hfill \delta _{\xi ^5}C_\mu & =2_\mu \xi ^5,\hfill \end{array}$$
(10)
with the remaining 4D bosonic fields being invariant. One can interpret $`\xi ^5(x^\mu )`$ as a parameter for local $`U(1)`$ transformations, for which $`C_\mu `$ is a gauge field. Note that the vector fields $`A_\mu ^I`$ transform non-trivially under these $`U(1)`$ transformations. To properly dimensionally reduce the theory, one must make field redefinitions in order to obtain $`U(1)`$ (KK)-invariant fields
$$A_\mu ^IA_\mu ^I+2C_\mu A^I.$$
The dimensionally reduced bosonic Lagrangian becomes GMZ05a
$$\begin{array}{cc}\hfill e^1_{DR}=& \frac{1}{2}R\frac{3}{4}\underset{IJ}{\overset{}{a}}D_\mu \stackrel{~}{h}^ID^\mu \stackrel{~}{h}^J\frac{1}{2}e^{2\sigma }\underset{IJ}{\overset{}{a}}D_\mu A^ID^\mu A^J\hfill \\ & (\frac{1}{2}e^{3\sigma }+e^\sigma \underset{IJ}{\overset{}{a}}A^IA^J)C_{\mu \nu }C^{\mu \nu }\frac{1}{4}e^\sigma \underset{IJ}{\overset{}{a}}_{\mu \nu }^I^{\mu \nu J}e^\sigma \underset{IJ}{\overset{}{a}}A^I_{\mu \nu }^JC^{\mu \nu }\hfill \\ & +\frac{e^1}{2\sqrt{6}}C_{IJK}ϵ^{\mu \nu \rho \sigma }(A^I_{\mu \nu }^J_{\rho \sigma }^K+2A^IA^J_{\mu \nu }^KC_{\rho \sigma }+\frac{4}{3}A^IA^JA^KC_{\mu \nu }C_{\rho \sigma })\hfill \\ & g^2P,\hfill \end{array}$$
(11)
where $`\stackrel{~}{h}^Ie^\sigma h^I`$, and
$$P\frac{3}{4}e^{3\sigma }\underset{IJ}{\overset{}{a}}(A^Kf_{KL}^Ih^L)(A^Mf_{MN}^Jh^N),$$
(12)
and
$`D_\mu A^I`$ $``$ $`_\mu A^I+gA_\mu ^Jf_{JK}^IA^K`$ (13)
$`D_\mu \stackrel{~}{h}^I`$ $``$ $`_\mu \stackrel{~}{h}^I+gA_\mu ^Jf_{JK}^I\stackrel{~}{h}^K`$ (14)
$`C_{\mu \nu }`$ $``$ $`2_{[\mu }C_{\nu ]}.`$ (15)
Just as $`A_\mu ^I`$ was redefined to be KK-invariant, we make the further redefinitions
$`\mathrm{\Psi }_\mu ^i`$ $``$ $`\mathrm{\Psi }_\mu ^i+\mathrm{\Psi }_{\dot{5}}^iC_\mu `$ (16)
$`\mathrm{\Gamma }_\mu `$ $``$ $`\mathrm{\Gamma }_\mu +\mathrm{\Gamma }_{\dot{5}}C_\mu ,`$ (17)
so that $`\mathrm{\Psi }_\mu ^i`$ and $`\mathrm{\Gamma }_\mu `$ are now KK-invariant. The dimensionally reduced susy transformations of the bosonic fields are
$$\begin{array}{cc}\hfill \delta ^{}e_\mu ^m=& \frac{1}{2}\overline{ϵ}^i\mathrm{\Gamma }^m\mathrm{\Psi }_{\mu i}^{(4)}\hfill \\ \hfill \delta A_\mu ^I=& \frac{1}{2}h_{\stackrel{~}{p}}^I\overline{ϵ}^i\mathrm{\Gamma }_\mu \lambda _i^{\stackrel{~}{p}}+\frac{1}{2}ie^\sigma \overline{\mathrm{\Psi }}_\mu ^iϵ_i(\frac{\sqrt{6}}{2}\stackrel{~}{h}^I+A^I)\hfill \\ \hfill \delta A^I=& \frac{1}{2}h_{\stackrel{~}{p}}^I\overline{ϵ}^i\mathrm{\Gamma }_{\dot{5}}\lambda _i^{\stackrel{~}{p}}+\frac{\sqrt{6}}{4}i\overline{\psi }^iϵ_ih^I\hfill \\ \hfill \delta e^\sigma =& \frac{1}{2}\overline{ϵ}^i\mathrm{\Gamma }^5\psi _i\hfill \\ \hfill \delta C_\mu =& \frac{1}{2}e^\sigma \overline{ϵ}^i\mathrm{\Gamma }^5\mathrm{\Psi }_{\mu i}\hfill \\ \hfill \delta \varphi ^x=& \frac{i}{2}f_a^x\overline{ϵ}^i\lambda _i^a,\hfill \end{array}$$
(18)
where $`\mathrm{\Psi }_{\widehat{\mu }}^i=(\mathrm{\Psi }_\mu ^i,\psi ^i)`$; $`\delta ^{}`$ denotes the “bare” susy transformation from five dimensions plus a local Lorentz transformation to maintain the condition $`\widehat{e}_{\dot{5}}^m=0`$; and we have identified the four-dimensional gravitini to be
$$\mathrm{\Psi }_{\mu i}^{(4)}e_\mu ^n\{\mathrm{\Psi }_{ni}+\frac{1}{2}(\mathrm{\Gamma }^n)^1\mathrm{\Gamma }^5\mathrm{\Psi }_{5i}\}.$$
(19)
The dimensionally reduced Lagrangian can be written in terms of a Kähler scalar manifold with complex scalars
$$z^I\frac{1}{3}A^I+\frac{i}{\sqrt{2}}\stackrel{~}{h}^I.$$
(20)
More details of the dimensional reduction to 4D $`𝒩=2`$ supergravity can be found in GMZ05a .
Boundary propagating multiplets
It’s clear from appendix B that the action of $`_2`$ on the supersymmetry spinors $`ϵ^i`$ necessarily requires half of the components to be odd, so that the original eight supersymmetry currents will be broken to four on the boundaries, so that there is at most $`𝒩=1`$ susy there. In terms of symplectic-Majorana spinors $`ϵ^i`$, the $`_2`$ action is represented as
$$i\mathrm{\Gamma }^5ϵ^1\text{and}i\mathrm{\Gamma }^5ϵ^2.$$
(The 4-component eigenspinors of the $`_2`$ action are linear combinations of the two symplectic-Majorana spinors.) The zero mode susy parameters can be written in the upstairs picture as (see table LABEL:tab:Table17 in appendix B)
$$ϵ^1=\left(\begin{array}{c}\eta \\ ϵ(x^5)e\zeta ^{}\end{array}\right),ϵ^2=\left(\begin{array}{c}ϵ(x^5)\zeta \\ e\eta ^{}\end{array}\right),$$
and so don’t have a well-defined limit on the fixed-planes. In the downstairs picture, on the other hand, fermions will have a well-defined limit at the boundaries (see JB e.g.); the fields in (3) can be written at the boundaries either as left-chiral fermions with their right-chiral conjugates:
$$\lambda ^{\stackrel{~}{p}\mathrm{\hspace{0.33em}1}}=\left(\begin{array}{c}\delta ^{\stackrel{~}{p}}\\ 0\end{array}\right),\lambda ^{\stackrel{~}{p}\mathrm{\hspace{0.33em}2}}=\left(\begin{array}{c}0\\ e\delta ^{\stackrel{~}{p}}\end{array}\right),$$
or right-chiral fermions with their left-chiral conjugates, which we denote with a bar:
$$\overline{\lambda }^{\stackrel{~}{p}\mathrm{\hspace{0.33em}1}}=\left(\begin{array}{c}0\\ e\gamma ^{\stackrel{~}{p}}\end{array}\right),\overline{\lambda }^{\stackrel{~}{p}\mathrm{\hspace{0.33em}2}}=\left(\begin{array}{c}\gamma ^{\stackrel{~}{p}}\\ 0\end{array}\right).$$
Assuming for simplicity that the only boundary-propagating vector fields are gauge fields for the bulk gauge group $`K`$, consistent parity assignments in the upstairs picture allow the following boundary multiplets in the downstairs picture:
| Multiplet | Representation | Type |
| --- | --- | --- |
| $`\{g_{\mu \nu },\mathrm{\Psi }_\mu \}`$ | $`K_\alpha `$ singlet | |
| $`\{A_\mu ^\alpha ,\lambda ^{\rho i}\}`$ | $`\text{adj}[K_\alpha ]`$ | $``$ |
| $`\{\overline{\lambda }^{pi},z^a\}`$ | $`R_V[K/K_\alpha ]`$ | $``$ |
| $`\{\mathrm{\Psi }_{\dot{5}},z^0\}`$ | $`K_\alpha `$ singlet | |
where we’ve split the index $`I=(0,\alpha ,a)`$; $`\stackrel{~}{x}=(x,\chi )`$; and $`\stackrel{~}{p}=(p,\rho )`$. We have denoted the surviving gauge group on the boundaries as $`K_\alpha `$, and $`R_V[K/K_\alpha ]`$ is the representation of the coset space elements. The value of $`n^{}`$ in $`\alpha =1,\mathrm{},n^{}`$ and $`a=(n^{}+1),\mathrm{},(n_V+1)`$ is arbitrary for now, so there is complete freedom in choosing parities for vector multiplets; we will consider this more carefully in section V. The second to last multiplet consists of a chiral multiplet in a real representation and its CPT conjugate. The case in which there are $`K`$-singlet vector fields propagating on the boundaries is straightforward.
What happens when a non-compact group is gauged in five dimensions? If the non-compact gauge fields were assigned even parity, then a non-compact gauge group would appear in the 4D theory. However, there would not be the proper degrees of freedom to give a ground state with compact gauge symmetry since the scalar degrees of freedom $`A^I`$ needed to form massive $`𝒩=1`$ vector multiplets must have odd parity. Therefore, the non-compact gauge fields must be assigned odd parity. We will then get $`𝒩=1`$ chiral multiplets in the coset $`K/H`$, with $`H`$ the maximal compact subgroup of $`K`$. This is a novel way of obtaining a 4D Higgs sector, along the lines of previous Higgs-gauge unifications in higher dimensions gaugehiggs , since there will be a distinct scalar potential for these scalars.
### III.2 Hypermultiplet couplings general coupling
A colection of $`n_H`$ hypermultiplets in five dimensions consist of $`2n_H`$ fermions and $`4n_H`$ real scalars, the latter parametrizing a quaternionic $`n_H`$-manifold $`_Q`$ with tangent space group $`USp(2n_H)\times SU(2)_R`$. We write the hypermultiplets as
$$\{\zeta ^A,q^{\stackrel{~}{X}}\},$$
where $`\stackrel{~}{X}=1,\mathrm{},4n_H`$ are the curved indices of $`_Q`$; and $`A=1,\mathrm{},2n_H`$ are flat, $`USp(2n_H)`$ indices. The $`4n_H`$-bein $`f_{iA}^{\stackrel{~}{X}}`$ relate scalar manifold curved and flat space metrics
$$g_{\stackrel{~}{X}\stackrel{~}{Y}}f_{iA}^{\stackrel{~}{X}}f_{jB}^{\stackrel{~}{Y}}=ϵ_{ij}C_{AB},$$
where $`i,j=1,2`$ are $`SU(2)_R`$ indices. Note that, in contrast to the case of vector multiplets, the scalars form $`2n_H`$ $`SU(2)_R`$-doublets, while the $`2n_H`$ fermions are $`SU(2)_R`$-singlets.<sup>4</sup><sup>4</sup>4In this paper, we are not considering gaugings of $`SU(2)_R`$ or its subgroups. In 2-component spinor notation, we write the fermions as
$$\zeta ^A=\left(\begin{array}{c}\zeta _1^A\\ \zeta _2^A\end{array}\right).$$
(21)
If non-trivial isometries of $`_Q`$ are gauged, they act on the scalars as
$$\delta _\alpha q^{\stackrel{~}{X}}=\alpha ^IK_I^{\stackrel{~}{X}}(q),$$
(22)
where the $`K_I^{\stackrel{~}{X}}`$ are the Killing fields on the quaternionic scalar manifold; and $`\alpha ^I`$ are the same local transformation parameters as in the pure YMESGT sector. The susy transformations for the scalars are
$$\delta q^{\stackrel{~}{X}}=i\overline{ϵ}^i\zeta ^Af_{iA}^{\stackrel{~}{X}}.$$
(23)
The 5D bosonic hypermultiplet Lagrangian (coupled to a YMESGT) is
$$\widehat{e}^1_{hyper}=\frac{1}{2}g_{\stackrel{~}{X}\stackrel{~}{Y}}𝒟_{\widehat{\mu }}q^{\stackrel{~}{X}}𝒟^\mu q^{\stackrel{~}{Y}}2g^2V_{iA}V^{iA},$$
(24)
where
$$V^{iA}\frac{\sqrt{6}}{4}h^IK_I^{\stackrel{~}{X}}f_{\stackrel{~}{X}}^{Ai}$$
and the $`K`$-covariant derivative is
$$D_\mu q^{\stackrel{~}{X}}𝒟_\mu q^{\stackrel{~}{X}}+gA_{\widehat{\mu }}^IK_I^{\stackrel{~}{X}}(q),$$
where $`𝒟_\mu `$ is the covariant derivative associated with the tangent space Lorentz and $`USp(2n_H)\times SU(2)`$ connections. The dimensional reduction of the bosonic Lagrangian is straightforward and will not be quoted.
Boundary propagating multiplets
Half of the hypermultiplet field content is required to have odd parity; therefore, let’s split the index $`\stackrel{~}{X}=(X,\chi )`$, with $`X=1,\mathrm{},2n_H`$ and $`\chi =2n_H+1,\mathrm{},4n_H`$. We let $`q^X`$ be the even parity fields, and $`q^\chi `$ the odd fields. Similarly, we spit the index $`A=(n,\stackrel{~}{n})`$ with $`n=1,\mathrm{},n_H`$ and $`\stackrel{~}{n}=n_H+1,\mathrm{},2n_H`$. If we couple hypermultiplets in the quaternionic-dimensional representation $`R_H[K]`$ of the gauge group to a 5D YMESGT, the multiplets with boundary propagating modes will be
| Multiplet | Representation | Type |
| --- | --- | --- |
| $`\{\zeta _1^n,q^{X_1}\}`$ | $`R_H[K_\alpha ]`$ | $``$ or $``$ |
| $`\{\zeta _2^{\stackrel{~}{n}},q^{X_2}\}`$ | $`\overline{R}_H[K_\alpha ]`$ | $``$ or $``$ |
where we have further split $`X=(X_1,X_2)`$ with $`X_1=1,\mathrm{},n_H`$ and $`X_2=n_H+1,\mathrm{},2n_H`$. We get a left-chiral multiplet and its CPT conjugate. Here $`R_H[K_\alpha ]`$ is the decomposition of $`R_H[K]`$ under the group $`K_\alpha K`$, and is the real-dimensional representation (see appendix A for conventions).
Example
Consider the “unified” MESGT with $`SU(5,1)`$ global symmetry group GZ03 ; M05a , whose vector fields are in 1-1 correspondence with the traceless elements of the Lorentzian Jordan algebra $`J_{(1,5)}^{}`$ GZ03 . The theory is “unified” in the sense that all of the vector fields of the 5D theory (including the bare graviphoton) furnish the $`\text{adj}[SU(5,1)]`$. The $`C_{IJK}`$ tensor is a rank-3 symmetric invariant of the global symmetry group, so its components are proportional to the $`d`$-symbols of $`SU(5,1)`$.
As in M05a , we can now couple hypermultiplets whose scalars parametrize the quaternionic manifold
$$_Q=\frac{E_7}{SO(12)\times SU(2)}$$
(25)
to the MESGT based on $`J_{(1,5)}^{}`$, gauging the common $`SU(5,1)`$ subgroup. Then the five-dimensional ground state would have at most an $`SU(5)\times U(1)`$ gauge group coupled to hypermultiplets in the $`\mathrm{𝟏}\mathrm{𝟓}\mathrm{𝟏𝟎}`$. We may then make parity assignments by associating each type of index with an $`SU(5)`$ representation:
| $`\alpha `$ | $`a`$ | $`0`$ | $`n`$ | $`\stackrel{~}{n}`$ |
| --- | --- | --- | --- | --- |
| $`\text{adj}[SM]`$ | $`SU(5)/SM\mathrm{𝟓}\overline{\mathrm{𝟓}}`$ | $`\mathbf{\hspace{0.17em}1}`$ | $`\mathbf{\hspace{0.17em}1}\overline{\mathrm{𝟓}}\mathrm{𝟏𝟎}`$ | $`c.c.`$ |
The 4D low energy effective theory of each boundary will have an $`𝒩=1`$ supergravity multiplet; SM gauge multiplets; weak doublet and color triplet chiral multiplets both with a scalar potential; and left-chiral matter multiplets (including a sterile fermion multiplet) along with their right-chiral conjugates. There are also the generic singlet left- and right-chiral multiplets coming from the 5D supergravity multiplet, as well as chiral multiplets in the symmetric space $`SU(5)/SM`$.
Remark:
At least within the category of homogeneous quaternionic manifolds, all occurences of hypermultiplets in the $`\mathrm{𝟓}`$ and $`\mathrm{𝟏𝟎}`$ of $`SU(5)`$ come from spaces admitting an $`SO(10)`$ isotropy subgroup under which these $`𝒩=2`$ hypermultiplets join a singlet to form the $`\mathrm{𝟏𝟔}`$ (see M05a ). In the previous example, we made a choice to truncate the $`𝒩=1`$ left-chiral multiplets in the $`\overline{\mathrm{𝟏𝟔}}`$ of $`SO(10)`$ (and their right-chiral conjugates in the $`\mathrm{𝟏𝟔}`$).
### III.3 Tensor multiplet couplings
When a MESGT with $`n_V`$ abelian vector multiplets is gauged, the symmetry group of the Lagrangian is broken to the gauge group $`KG`$. The $`n_V+1`$ vector fields decompose into $`K`$-reps
$$𝐧_𝐕+\mathrm{𝟏}=\text{adj}(K)\text{non-singlets}(K)\text{singlets}(K).$$
Such a gauging requires the non-singlet vector fields to be dualized to anti-symmetric tensor fields GZ:99dec satisfying a field equation that serves as a “self-duality” constraint selfdual (thus keeping the number of degrees of freedom the same):
$$B_{\mu \nu }^M=c_N^Mϵ_{\mu \nu }^{\rho \sigma \lambda }_{[\rho }B_{\sigma \lambda ]}^N+\mathrm{},$$
(26)
where $`c_N^M`$ has dimensions of inverse mass; square brackets denote anti-symmetric permutations; and ellipses denote terms involving other fields. A tensor field does not require an abelian invariance to remain massless.
We have already discussed the scalar sector of a pure 5D YMESGT. When tensor multiplets are coupled the scalar manifold is again a real Riemannian space, but which cannot be decomposed globally as a product of “vector” and “tensor” parts. We can, of course, identify an orthogonal frame of scalars at each point of the manifold: the vector multiplets are associated with the combination $`h_{\stackrel{~}{x}}^I\varphi ^{\stackrel{~}{x}}`$ at a given point, while the tensor multiplets are associated with the independent combination $`h_{\stackrel{~}{x}}^M\varphi ^{\stackrel{~}{x}}`$. Similarly, the combination of fermions $`h_{\stackrel{~}{p}}^I\lambda ^{\stackrel{~}{p}i}`$ are associated with vector multiplets, while $`h_{\stackrel{~}{p}}^M\lambda ^{\stackrel{~}{p}i}`$ with tensor multiplets. We will write $`\varphi ^{\stackrel{~}{x}}`$ and $`\varphi ^{\stackrel{~}{m}}`$ to denote the scalar partners of the vector and tensors, respectively, at any given point of the scalar manifold. Similarly, we write $`\lambda ^{\stackrel{~}{p}i}`$ and $`\lambda ^{\stackrel{~}{\mathrm{}}i}`$ as the fermionic partners of the vector and tensor fields, respectively. It is then implicitly understood that the meaning of this notation is given by the above discussion.
When tensors are present, we will use indices $`I,J,K`$ for 5D vector fields and $`M,N,P`$ for 5D tensor fields. Then $`n_T`$ tensor multiplets are
$$\{B_{\widehat{\mu }\widehat{\nu }}^M,\lambda ^{\stackrel{~}{\mathrm{}}i},\varphi ^{\stackrel{~}{m}}\}.$$
To be consistent with the gauge symmetry, the components of the $`C`$-tensor are constrained to be GZ:99dec :
$$\begin{array}{c}\hfill C_{IMN}=\frac{\sqrt{6}}{2}\mathrm{\Omega }_{NP}\mathrm{\Lambda }_{IM}^P\\ \hfill C_{MNP}=0,C_{MIJ}=0,\end{array}$$
(27)
where $`\mathrm{\Omega }_{NP}`$ is antisymmetric and $`\mathrm{\Lambda }_{IM}^P`$ are symplectic $`K`$-representation matrices appearing in the $`K`$-transformation of the tensor fields: $`\delta _\alpha B_{\mu \nu }^M=\alpha ^I\mathrm{\Lambda }_{IN}^MB_{\mu \nu }^N.`$ Furthermore, $`C_{IJK}`$ must be a rank-three symmetric $`K`$-invariant tensor. Note: We are assuming a compact or non-semisimple gauge group $`K`$; see Bergshoeff for more general couplings where $`C_{MIJ}0`$.
The new or modified terms in the bosonic 5D Lagrangian involving are GZ:99dec
$$\begin{array}{cc}& _T=\frac{\widehat{e}}{4}\underset{MN}{\overset{}{a}}B_{\widehat{\mu }\widehat{\nu }}^MB^{N\widehat{\mu }\widehat{\nu }}\frac{\widehat{e}}{2}\underset{IM}{\overset{}{a}}_{\widehat{\mu }\widehat{\nu }}^IB^{M\widehat{\mu }\widehat{\nu }}\hfill \\ & +\frac{1}{4g}ϵ^{\widehat{\mu }\widehat{\nu }\widehat{\rho }\widehat{\sigma }\widehat{\lambda }}\mathrm{\Omega }_{MN}B_{\widehat{\mu }\widehat{\nu }}^M_{\widehat{\rho }}B_{\widehat{\sigma }\widehat{\lambda }}^N\hfill \\ & \frac{e}{2}g_{\stackrel{~}{x}\stackrel{~}{y}}𝒟_{\widehat{\mu }}\varphi ^{\stackrel{~}{x}}𝒟^{\widehat{\mu }}\varphi ^{\stackrel{~}{y}}+\frac{1}{2\sqrt{6}}C_{MNI}ϵ^{\widehat{\mu }\widehat{\nu }\widehat{\rho }\widehat{\sigma }\widehat{\lambda }}B_{\widehat{\mu }\widehat{\nu }}^MB_{\widehat{\rho }\widehat{\sigma }}^NA_{\widehat{\lambda }}^I,\hfill \end{array}$$
where $`\varphi ^{\stackrel{~}{x}}`$ are scalars in both vector and tensor multiplets, and
$$𝒟_{\widehat{\mu }}\varphi ^{\stackrel{~}{x}}=_{\widehat{\mu }}\varphi ^{\stackrel{~}{x}}+K_I^{\stackrel{~}{x}}A_{\widehat{\mu }}^I,$$
with $`K_I^{\stackrel{~}{x}}`$ the Killing vectors on the vector/tensor scalar manifold. The 5D field equations for the $`B_{\widehat{\mu }\widehat{\nu }}^M`$ are (in terms of forms)
$$\begin{array}{cc}& {}_{}{}^{}DB^M=g\mathrm{\Omega }^{MN}\underset{M\stackrel{~}{I}}{\overset{}{a}}^{\stackrel{~}{I}},\hfill \\ & \text{where}^{\stackrel{~}{I}}=\left(\begin{array}{c}^I\\ B^M\end{array}\right).\hfill \end{array}$$
(28)
The presence of non-trivially charged tensors also introduces a scalar potential $`P^{(T)}`$ that was not present in the case of pure YMESGTs. The term in the Lagrangian is GZ:99dec
$$\begin{array}{cc}& _{P^{(T)}}=2g^2\widehat{e}W^{\stackrel{~}{p}}W^{\stackrel{~}{p}},\hfill \\ & \text{where}W^{\stackrel{~}{p}}=\frac{\sqrt{6}}{8}h_M^{\stackrel{~}{p}}\mathrm{\Omega }^{MN}h_N.\hfill \end{array}$$
(29)
Dimensional reduction
In the dimensional reduction, we parametrize the tensor field as
$$B_{\widehat{\mu }\widehat{\nu }}^M=\left(\begin{array}{ccc}0& & \stackrel{~}{A}_\nu ^M\\ \stackrel{~}{A}_\mu ^M& & B_{\mu \nu }^M\end{array}\right),$$
where tildes have been used to help distinguish from 4D vectors arising as components of 5D vectors. The resulting Lagrangian containing both $`\stackrel{~}{A}_\mu ^M`$ and $`B_{\mu \nu }^M`$ is analogous to the 1st order formulation of the Freedman-Townsend model Freedman . One can obtain a 2nd order formulation in terms of the $`\stackrel{~}{A}_\mu ^M`$ by using the Euler-Lagrange equations, which appear as constraints relating these fields in the dimensionally reduced theory. However, in the case of dimensional reduction, there is an obstruction GMZ05b to obtaining a local 2nd order Lagrangian in terms of the $`B_{\mu \nu }^M`$. This will not be present in the case of the orbifold.
The $`\xi ^5`$ transformations of the dimensionally reduced fields $`B_{\mu \nu }^M`$ and $`\stackrel{~}{A}_\mu ^M`$ are
$$\begin{array}{cc}\hfill \delta _{\xi ^5}B_{\mu \nu }^M& =_\mu \xi ^5\stackrel{~}{A}_\nu ^M_\nu \xi ^5\stackrel{~}{A}_\mu ^M\hfill \\ \hfill \delta _{\xi ^5}\stackrel{~}{A}_\mu ^M& =0.\hfill \end{array}$$
(30)
We must therefore make a field redefinition
$$B_{\mu \nu }^MB_{\mu \nu }^M4C_{[\mu }\stackrel{~}{A}_{\nu ]}^M$$
(31)
so that the $`B_{\mu \nu }^M`$ are now KK-invariant. The full dimensional reduction of tensor-coupled YMESTs is given in GMZ05b . The new or modified terms in the 5D bosonic Lagrangian are
$$\begin{array}{cc}& e^1^{(4)}=e^{2\sigma }\stackrel{}{a}_{IM}(D_\mu A^I)\stackrel{~}{A}^{\mu M}\frac{1}{2}e^{2\sigma }\stackrel{}{a}_{MN}\stackrel{~}{A}_\mu ^M\stackrel{~}{A}^{\mu N}\hfill \\ & \frac{3}{4}\stackrel{}{a}_{\stackrel{~}{I}\stackrel{~}{J}}(𝒟_\mu \stackrel{~}{h}^{\stackrel{~}{I}})(𝒟^\mu \stackrel{~}{h}^{\stackrel{~}{J}})+\frac{e^1}{2\sqrt{6}}C_{MNI}ϵ^{\mu \nu \rho \sigma }B_{\mu \nu }^MB_{\rho \sigma }^NA^I\hfill \\ & +\frac{e^1}{g}ϵ^{\mu \nu \rho \sigma }\mathrm{\Omega }_{MN}C_{\mu \nu }\stackrel{~}{A}_\rho ^M\stackrel{~}{A}_\sigma ^N+\frac{e^1}{g}ϵ^{\mu \nu \rho \sigma }\mathrm{\Omega }_{MN}B_{\mu \nu }^MD_\rho \stackrel{~}{A}_\sigma ^N\hfill \\ & \frac{1}{4}e^\sigma \stackrel{}{a}_{MN}B_{\mu \nu }^MB^{N\mu \nu }\frac{1}{2}e^\sigma \stackrel{}{a}_{IM}(_{\mu \nu }^I+2C_{\mu \nu }A^I)B^{M\mu \nu }\hfill \\ & g^2P,\hfill \end{array}$$
(32)
where
$`D_\mu \stackrel{~}{A}_\nu ^M`$ $``$ $`_\mu \stackrel{~}{A}_\nu ^N+gA_\mu ^I\mathrm{\Lambda }_{IP}^N\stackrel{~}{A}_\nu ^P`$ (33)
$`𝒟_\mu \stackrel{~}{h}^{\stackrel{~}{I}}`$ $``$ $`_\mu \stackrel{~}{h}^{\stackrel{~}{I}}+gA_\mu ^IM_{I\stackrel{~}{K}}^{\stackrel{~}{I}}\stackrel{~}{h}^{\stackrel{~}{K}}.`$ (34)
The total scalar potential, $`P`$, is now
$$\begin{array}{cc}\hfill P=& \mathrm{\hspace{0.33em}2}e^\sigma W^{\stackrel{~}{p}}W^{\stackrel{~}{p}}\hfill \\ & +\frac{3}{4}e^{3\sigma }\stackrel{}{a}_{\stackrel{~}{I}\stackrel{~}{J}}(A^IM_{I\stackrel{~}{K}}^{\stackrel{~}{I}}h^{\stackrel{~}{K}})(A^JM_{J\stackrel{~}{L}}^{\stackrel{~}{J}}h^{\stackrel{~}{L}}),\hfill \end{array}$$
(35)
where $`W^{\stackrel{~}{p}}`$ is defined in eqn (29), and
$$M_{(I)\stackrel{~}{K}}^{\stackrel{~}{J}}=\left(\begin{array}{cc}f_{IK}^J& 0\\ 0& \mathrm{\Lambda }_{IM}^N\end{array}\right),$$
(36)
and $`\stackrel{~}{I}=(I,M)`$.
Parity assignments of tensor sector
Since (30) is only true for KK-transformations connected to the identity, the $`\stackrel{~}{A}_\mu ^M`$ are not necessarily even under $`_2`$ action. However, the above do lead to the constraint
$$P(\stackrel{~}{A}_\mu ^M)=P(B_{\mu \nu }^M),$$
componentwise. These two fields do not describe independent propagating degrees of freedom since they are related by a constraint equation (coming from the fact that the 5D tensors satisfied a “self-duality” field equation (26) reducing the number of propagating modes):
$$B_{\mu \nu }^M=c_N^M(^{}D\stackrel{~}{A}^N)_{\mu \nu }+\mathrm{},$$
(37)
where $`c_N^M`$ is proportional to $`\mathrm{\Omega }^{MP}\underset{PN}{\overset{}{a}}`$; $``$ is the Hodge operator; and the dots indicate terms involving other fields. There are two classes of assignments we can make, characterized by the parity of the symplectic form $`\mathrm{\Omega }_{MN}`$ on the vector space spanned by the 5D tensors.
#### III.3.1 Odd Parity $`\mathrm{\Omega }_{MN}`$
If the self-duality relation is used to express all tensor fields in terms of the vectors $`\stackrel{~}{A}_\mu ^M`$, the mass of the $`\stackrel{~}{A}_\mu ^M`$ is non-vanishing at the orbifold fixed points. However, there are insufficient fermionic degrees of freedom to form massive $`𝒩=1`$ vector multiplets. Therefore, we must use the Euler-Lagrange equations to write $`\stackrel{~}{A}_\mu ^MB_{\mu \nu }^M`$. There is an obstruction to doing this with the dimensionally reduced 1st order Lagrangian obtained by using KK-invariant field redefinitions GMZ05b , in particular due to the topological term of the form $`\mathrm{\Omega }_{MN}C_{\mu \nu }\stackrel{~}{A}_\mu ^M\stackrel{~}{A}_\nu ^N`$, but the problem is avoided in the orbifold theory since the obstruction vanishes on the 4D fixed planes.
Once the 2nd order Lagrangian with tensors is obtained, the $`B_{\mu \nu }^M`$ can then be Hodge dualized to scalars $`B^M`$ by adding a term of the form
$$e^1\mathrm{\Delta }=ϵ^{\mu \nu \rho \sigma }\mathrm{\Omega }_{MN}_{\mu \nu \rho }^MD_\sigma B^N$$
(38)
to the Lagrangian (if this step is performed before integrating over the fifth dimension, one needs a $`\delta (x^5)`$ factor), where $`D_\rho `$ is the gauge covariant derivative acting on the scalars, and $`_{\mu \nu \rho }^M=3!_{[\mu }B_{\nu \rho ]}^M`$.
The multiplets that will propagate on the fixed planes are
| Multiplet | Representation | Type |
| --- | --- | --- |
| $`\{\overline{\lambda }^\mathrm{}i,z^M\}`$ | $`𝐍\overline{𝐍}`$ | $``$ |
That is, there are left-chiral multiplets in a real representation along with the CPT conjugate right-chiral multiplets.
Example
Consider the “unified” 5D MESGT based on the Lorentzian Jordan algebra $`J_{(1,5)}^{}`$, whose global symmetry group is $`SU(5,1)`$ GZ03 . We can couple this theory to hypermultiplets whose scalars parametrize the particular scalar manifold (25). If we gauge the common $`SU(5)\times U(1)SU(5,1)`$ subgroup, we will get $`SU(5)\times U(1)`$ gauge multiplets, along with tensor multiplets in the $`\mathrm{𝟓}\overline{\mathrm{𝟓}}`$ and hypermultiplets in the $`\mathrm{𝟏}\mathrm{𝟓}\mathrm{𝟏𝟎}`$. This is then similar to the ground state theory in the $`SU(5,1)`$ gauging example before, but with some important differences; one of which being that the 5D bare graviphoton $`A_{\widehat{\mu }}^0`$ does not take part in gauging the isometries of the scalar manifold. We can make the assignments
| $`\alpha `$ | $`a`$ | $`0`$ | $`M`$ | $`n`$ | $`\stackrel{~}{n}`$ |
| --- | --- | --- | --- | --- | --- |
| $`\text{adj}[SM\times U(1)]`$ | $`SU(5)/SM`$ | $`\mathrm{𝟏}`$ | $`\mathrm{𝟓}\overline{\mathrm{𝟓}}`$ | $`\mathrm{𝟏}\overline{\mathrm{𝟓}}\mathrm{𝟏𝟎}`$ | $`c.c.`$ |
The propagating modes along the fixed planes will be $`SU(3)\times SU(2)\times U(1)^2`$ gauge fields; weak doublet (Higgs) chiral multiplets; color triplet chiral multiplets; and left-chiral matter multiplets (including a sterile fermion multiplet) with their CPT conjugates. Again, there is also the generic singlet spin-1/2 multiplet coming from the 5D supergravity multiplet, and chiral multiplets in the symmetric space $`SU(5)/SM`$. All of these multiplets are tree-level massless, while the scalars in the $`\mathrm{𝟓}\overline{\mathrm{𝟓}}`$ have a potential term.
#### III.3.2 Even Parity for $`\mathrm{\Omega }_{MN}`$
If $`\mathrm{\Omega }_{MN}`$ has even parity, then the $`_2`$ action acts reducibly on the symplectic vector space and projects out half of the bosonic fields. Consistency with supersymmetry requires that we use the field equation (37) to eliminate $`B_{\mu \nu }^M`$ from the Lagrangian. Splitting $`M=(,\overline{})`$ with $`=1,\mathrm{},n_T/2`$ and $`\overline{}=n_T/2,\mathrm{},n_T`$, the $`\stackrel{~}{A}_\mu ^{}`$ then have a mass matrix
$$M_𝒩e^{2\sigma }\underset{𝒩}{\overset{}{a}}.$$
Therefore, the $`n_T/2`$ vectors, $`n_T`$ spin-1/2 fields, and $`n_T/2`$ scalars form $`n_T/2`$ massive $`𝒩=1`$ vector multiplets.
The multiplets with propagating modes on the boundaries are
| Multiplet | Representation | Type |
| --- | --- | --- |
| $`\{A^{},\lambda ^{\overline{\mathrm{}}i},h^{}\}`$ | $`𝐍_{}`$ | $``$ |
The notation for the representation means that the gauge group at the fixed points must support a real $`𝐍`$ (whereas the 5D gauge group had to support a complex $`𝐍`$). Let’s illustrate this with an example.
Example
The minimal example in which one is left with a group containing SM is where the 5D gauge group is $`SU(10)\times U(1)`$. Starting with the “unified” MESGT defined by the Lorentzian Jordan algebra $`J_{(1,10)}^{}`$ and with $`SU(N,1)`$ global symmetry of the Lagrangian, we can gauge the $`SU(10)\times U(1)`$ subgroup, yielding tensors in the $`\mathrm{𝟏𝟎}\overline{\mathrm{𝟏𝟎}}`$. If the symplectic form has even parity, then the orbifold conditions require the group to be broken to at least $`SO(10)\times U(1)`$, under which we have massive vector multiplets in the (real) $`\mathrm{𝟏𝟎}`$. This theory will appear as a ground state of the 4D $`𝒩=1`$ theory with gauge group $`(SO(10)\times U(1))T^{10}`$, analogous to the $`𝒩=2`$ case discussed in GMZ05b . There are also chiral multiplets from the broken gauge multiplets forming the $`\mathrm{𝟓𝟒}`$, along with their CPT conjugates.
## IV Objects other than fields
There are field-dependent and independent objects that appear in the Lagrangian and supersymmetry transformations that carry $`_2`$ parity. In particular, the field independent objects are the $`C_{IJK}`$ tensor defining the MESGT that exists prior to gauging; the structure constants $`f_{JK}^I`$ and transformation parameters $`\alpha ^I(x)`$ of the 5D gauge group; and the symplectic tensor $`\mathrm{\Omega }_{MN}`$ and transformation matrices $`\mathrm{\Lambda }_{IN}^M`$ in the tensor coupled theory. These contain an implicit $`ϵ(x^5)`$ or $`\kappa (x^5)`$ factor when assigned odd parity. The field dependent objects are the restricted ambient space metric $`\underset{IJ}{\overset{}{a}}(\varphi )`$ and scalar manifold metrics $`g_{xy}(\varphi )`$ and $`g_{XY}(q)`$; the $`h_p^I(\varphi )`$; the scalar vielbein $`f_x^p(\varphi )`$ and $`f_{iA}^X(q)`$; the Killing vectors on the scalar manifold $`K_x^I(\varphi )`$ and $`K_X^I(q)`$. These vanish when assigned odd parity due to the form of the bosonic field expansion in eqn (1).
Pure YMESGT
In MESGT , it was shown that the $`C_{IJK}`$ defining a MESGT may be put in a “canonical” basis satisfying the positivity of $`𝒱=C_{IJK}\zeta ^I\zeta ^J\zeta ^K`$ (see (5)). The parity assignments of the components are determined by requiring the polynomial $`𝒱`$ to be invariant under $`_2`$ action. Splitting $`i=(\alpha ,a)`$, we have
| Even | Odd |
| --- | --- |
| $`C_{000}C_{0ab}C_{0\alpha \beta }C_{abc}C_{a\alpha \beta }`$ | $`C_{0a\alpha }C_{\alpha \beta \gamma }C_{\alpha ab}`$ |
There is freedom in choosing $`ϵ(x^5)`$ or $`\kappa (x^5)`$ as the jumping function for odd components.
Consistency of the infinitesimal gauge transformations (6) require parity assignments for the $`f_{JK}^I`$ and $`\alpha ^I(x)`$ to be
| Even | Odd |
| --- | --- |
| $`f_{\beta \gamma }^\alpha f_{ab}^\alpha `$ | $`f_{bc}^af_{\alpha \beta }^a`$ |
| $`f_{a\beta }^0f_{\alpha 0}^0`$ | $`f_{\alpha \beta }^0`$ |
| $`\alpha ^\beta `$ | $`\alpha ^0\alpha ^b`$ |
where $`f_{JK}^I`$ vanishes if any of the indices correspond to 5D spectator vector fields <sup>5</sup><sup>5</sup>5This will be true, e.g., for the “bare graviphoton” $`A_\mu ^0`$ if the 5D gauge group is compact.; and permutations of the indices have the same parity. The (upstairs picture) gauge transformation parameters are subject to an expansion on $`S^1/_2`$ just as in (1). Consistency of the 5D gauge algebra requires that odd $`f_{JK}^I`$ be redefined by $`ϵ(x^5)f_{JK}^I`$, where the $`f_{JK}^I`$ are now even <sup>6</sup><sup>6</sup>6Note that the orbifold fixed planes are not oriented surfaces..
The components of the restricted ambient space metric and scalar manifold metric have parities determined by the requirement that the line elements of those spaces be preserved:
| Even | Odd |
| --- | --- |
| $`\underset{\alpha \beta }{\overset{}{a}}\underset{ab}{\overset{}{a}}`$ | $`\underset{a\beta }{\overset{}{a}}\underset{0\alpha }{\overset{}{a}}`$ |
| $`\underset{00}{\overset{}{a}}\underset{0a}{\overset{}{a}}`$ | |
| $`g_{xy}g_{\chi \psi }`$ | $`g_{\chi y}`$ |
Consistency of the gauge transformations (6) determine the parities of the scalar manifold Killing fields:
| Even | Odd |
| --- | --- |
| $`K_\alpha ^x`$ | $`K_0^xK_a^x`$ |
| $`K_0^\chi K_a^\chi `$ | $`K_\alpha ^\chi `$ |
The non-zero components $`K_a^\chi `$, which are a set of sections of the normal bundle over the 4D scalar manifold, determine the form of the scalar potential involving the $`\varphi ^x`$ and $`A^a`$ at the fixed points:
$$g_{\chi \psi }K_a^\chi K_b^\psi A^aA^b,$$
(39)
where $`g_{\chi \psi }`$ is the metric determined by the normal bundle connection.
Finally, the functions $`h^I`$ and $`h_p^I`$; the scalar vielbein $`f_x^p`$; and the functions $`h_x^I=h_p^If_x^p`$ are required to satisfy
| Even | Odd |
| --- | --- |
| $`h^0h^a`$ | $`h^\alpha `$ |
| $`h_p^0h_p^ah_\rho ^\alpha `$ | $`h_\rho ^ah_p^\alpha `$ |
| $`f_x^pf_\chi ^\rho `$ | $`f_\chi ^pf_x^\rho `$ |
| $`h_x^0h_x^ah_\chi ^\alpha `$ | $`h_\chi ^ah_x^\alpha `$ |
Tensor couplings
In the tensor-coupled theory, parity assignments depend on the parity of the symplectic form $`\mathrm{\Omega }_{MN}`$. The parities of the additional $`C`$-tensor components must be
| Even | Odd |
| --- | --- |
| $`C_{\overline{𝒩}\alpha }`$ | $`C_{𝒩\alpha }C_{\overline{}\overline{𝒩}\alpha }`$ |
| $`C_{𝒩a}C_{\overline{}\overline{𝒩}a}`$ | $`C_{\overline{𝒩}a}`$ |
| $`P(\mathrm{\Omega }_{MN})=+\mathrm{\Omega }_{MN}`$ | |
| Even | Odd |
| --- | --- |
| $`C_{MNa}`$ | $`C_{MN\alpha }`$ |
| $`P(\mathrm{\Omega }_{MN})=\mathrm{\Omega }_{MN}`$ | |
Consistency of the gauge transformations require the representation matrices to satisfy
| Even | Odd |
| --- | --- |
| $`\mathrm{\Lambda }_{\alpha 𝒩}^{}\mathrm{\Lambda }_{\alpha \overline{𝒩}}^\overline{}`$ | $`\mathrm{\Lambda }_{\alpha 𝒩}^\overline{}\mathrm{\Lambda }_{\alpha \overline{𝒩}}^{}`$ |
| $`\mathrm{\Lambda }_{0𝒩}^\overline{}\mathrm{\Lambda }_{0\overline{𝒩}}^{}`$ | $`\mathrm{\Lambda }_{0\overline{𝒩}}^\overline{}\mathrm{\Lambda }_{0𝒩}^{}`$ |
| $`\mathrm{\Lambda }_{a𝒩}^\overline{}\mathrm{\Lambda }_{a\overline{𝒩}}^{}`$ | $`\mathrm{\Lambda }_{a\overline{𝒩}}^\overline{}\mathrm{\Lambda }_{a𝒩}^{}`$ |
| $`P(\mathrm{\Omega }_{MN})=+\mathrm{\Omega }_{MN}`$ | |
| Even | Odd |
| --- | --- |
| $`\mathrm{\Lambda }_{\alpha N}^M`$ | $`\mathrm{\Lambda }_{0N}^M\mathrm{\Lambda }_{aN}^M`$ |
| $`P(\mathrm{\Omega }_{MN})=\mathrm{\Omega }_{MN}`$ | |
If odd gauge transformation parameters $`\alpha ^a`$ are expanded as in eqn (1), odd $`\mathrm{\Lambda }_{IN}^Mϵ(x^5)\mathrm{\Lambda }_{IN}^M`$ for consistency of the gauge algebra; the relation $`C_{MNI}\mathrm{\Omega }_{MP}\mathrm{\Lambda }_{IN}^P`$ from eqn (27) then requires odd $`C_{IMN}ϵ(x^5)C_{IMN}`$. Note that when $`\mathrm{\Omega }_{MN}`$ has odd parity, the set of coefficients $`c_N^M\mathrm{\Omega }^{MP}\underset{PN}{\overset{}{a}}`$ are odd; a choice that was made in YL:03dec .
As in the pure YMESGT case, the ambient space and scalar manifold line elements should be preserved under the $`_2`$ action so that
| Even | Odd |
| --- | --- |
| $`\underset{𝒩}{\overset{}{a}}\underset{\overline{}\overline{𝒩}}{\overset{}{a}}`$ | $`\underset{\overline{𝒩}}{\overset{}{a}}`$ |
| $`\underset{\overline{}\alpha }{\overset{}{a}}\underset{a}{\overset{}{a}}`$ | $`\underset{\alpha }{\overset{}{a}}\underset{\overline{}a}{\overset{}{a}}`$ |
| $`g_{\stackrel{~}{x}\stackrel{~}{y}}g_{mn}`$ | $`g_{\stackrel{~}{x}m}`$ |
| $`P(\mathrm{\Omega }_{MN})=+\mathrm{\Omega }_{MN}`$ | |
| Even | Odd |
| --- | --- |
| $`\underset{aM}{\overset{}{a}}`$ | $`\underset{\alpha M}{\overset{}{a}}`$ |
| $`g_{\stackrel{~}{x}\stackrel{~}{y}}g_{mn}`$ | $`g_{\stackrel{~}{x}m}`$ |
| $`P(\mathrm{\Omega }_{MN})=\mathrm{\Omega }_{MN}`$ | |
Finally, the functions $`h^M(\varphi )`$ and $`h_{\mathrm{}}^M`$; the vielbein $`f_m^{\mathrm{}}`$; and the functions $`h_x^I=f_x^{\mathrm{}}h_{\mathrm{}}^I`$ satisfy
| Even | Odd |
| --- | --- |
| $`h^{}(\varphi )`$ | $`h^\overline{}(\varphi )`$ |
| $`h_m^{}h_{\overline{m}}^\overline{}`$ | $`h_{\overline{m}}^{}h_m^\overline{}`$ |
| $`h_{\mathrm{}}^{}h_\overline{\mathrm{}}^\stackrel{~}{}`$ | $`h_{\mathrm{}}^\stackrel{~}{}h_\overline{\mathrm{}}^{}`$ |
| $`P(\mathrm{\Omega }_{MN})=+\mathrm{\Omega }_{MN}`$ | |
| Even | Odd |
| --- | --- |
| $`h^M(\varphi )`$ | |
| $`h_m^M`$ | |
| $`h_{\mathrm{}}^M`$ | $`h_\overline{\mathrm{}}^M`$ |
| $`P(\mathrm{\Omega }_{MN})=\mathrm{\Omega }_{MN}`$ | |
where we have split $`\stackrel{~}{\mathrm{}}=(\mathrm{},\overline{\mathrm{}})`$ and $`\stackrel{~}{m}=(m,\overline{m})`$.
Hypermultiplet couplings
The parity assigments for the Killing vectors and vielbein of the quaternionic scalar manifold are required to be
| Even | Odd |
| --- | --- |
| $`K_\alpha ^{X_1}K_a^{X_2}`$ | $`K_a^{X_1}K_\alpha ^{X_2}`$ |
| $`f_{1A_1}^{X_1}f_{2A_2}^{X_1}`$ | $`f_{2A_1}^{X_1}f_{1A_2}^{X_1}`$ |
| $`f_{2A_1}^{X_2}f_{1A_2}^{X_2}`$ | $`f_{1A_1}^{X_2}f_{2A_2}^{X_2}`$ |
## V Parity assignments of the scalar sector
In the previous sections, we listed the boundary multiplets by “orbifolding” 5D YMESGTs coupled to tensor multiplets. It’s clear, for example, that 5D vector multiplets should yield 4D boundary vector multiplets when the associated scalar manifold embedding functions $`h^I(\varphi )`$ are assigned odd parity, so that the parities of the $`h^I`$ induce parity assignments for the physical scalars $`\varphi ^{\stackrel{~}{x}}`$. However, it is not immediately clear how the assignments are transmitted from one to the other. That is, how are the physical scalars truncated in light of the fact that the ‘odd’ scalar partners of the ‘even’ vectors $`A_{\widehat{\mu }}^\alpha `$ are the combinations $`h_{\stackrel{~}{x}}^\alpha \varphi ^{\stackrel{~}{x}}`$, where $`h_{\stackrel{~}{x}}^\alpha h_{,\stackrel{~}{x}}^\alpha `$. We will consider this using two examples based on symmetric “very special” real scalar manifolds.
The first example is of the “generic Jordan” family of Maxwell-Einstein supergravity theories (MESGTs) with symmetric scalar manifolds of the form MESGT
$$_R=SO(1,1)\times \frac{SO(n_V1,1)}{SO(n_V1)},$$
where $`n_V`$ is the number of vector multiplets. The cubic polynomial for the theory in the absence of an orbifold is $`𝒱=C_{IJK}\xi ^I\xi ^J\xi ^K`$, where
$$C_{000}=1,C_{00i}=0,C_{0ij}=\frac{1}{2}\delta _{ij},$$
$$C_{111}=\frac{1}{\sqrt{2}},C_{1ab}=+\frac{1}{\sqrt{2}}\delta _{ab},$$
with $`a,b=2,\mathrm{},n_V`$. On $`_4\times S^1/_2`$, the $`C_{IJK}`$ can have odd components satisfying jumping conditions, though we take $`h^0`$ to always have even parity.
For example, assigning odd parity to $`h^1`$ and $`C_{1ij}`$ (and redefining $`C_{1ij}ϵ(x^5)C_{1ij}`$), the polynomial is
$$\begin{array}{cc}\hfill 𝒱=& \left(\frac{2}{3}\right)^{\frac{3}{2}}[(\xi ^0)^3\frac{3}{2}\xi ^0\delta _{ij}\xi ^i\xi ^j\frac{ϵ(x^5)}{\sqrt{2}}(\xi ^1)^3\hfill \\ & +\frac{3ϵ(x^5)}{\sqrt{2}}\xi ^1[(\xi ^2)^2+\mathrm{}+(\xi ^{n_V})^2]].\hfill \end{array}$$
Restricting to the scalar hypersurface $`𝒱=1`$, the terms with $`h^1\xi ^1|_{𝒱=1}`$ vanish (since $`h^1`$ is a parity-odd function of scalars) so that, at the orbifold fixed points,
$$𝒱|_{fp}=h^0\left\{(h^0)^2\frac{3}{2}\delta _{ab}h^ah^b\right\}.$$
Remark: In general, 4D $`𝒩=1`$ supergravity theories are in 1-1 correspondence with Hodge manifolds. The 4D $`𝒩=1`$ supergravity theory we obtain from orbifolding is of a special class based on a (not necessarily irreducible) cubic polynomial satisfying $`𝒱_{fp}(\text{Im}(z))=e^{3\sigma }>0`$.
In the family of ‘generic Jordan’ MESGTs, however, the solution to the condition $`𝒱=1`$ is simplest by making linear transformations $`\stackrel{ˇ}{\xi }^I=M_J^I\xi ^J`$ to a ‘non-canonical’ basis scalarmfds (before orbifolding the theory):
$$\begin{array}{cc}& M_0^0=\frac{1}{\sqrt{3}},M_1^0=\sqrt{\frac{2}{3}}\hfill \\ & M_0^1=\sqrt{\frac{2}{3}},M_1^1=\frac{1}{\sqrt{3}}\hfill \\ & M_ı^ȷ=\delta _ı^ȷ,\hfill \end{array}$$
where $`ı,ȷ=2,\mathrm{},n_V`$; and all other $`M_J^I`$ are zero. The solution is then
$$\begin{array}{cc}& 𝒱=\frac{3^{3/2}}{2}\stackrel{ˇ}{h}^0\stackrel{ˇ}{h}^2=1\hfill \\ & \stackrel{ˇ}{h}^0=\frac{1}{\sqrt{3}\varphi ^2},\stackrel{ˇ}{h}^i=\sqrt{\frac{2}{3}}\varphi ^i,\hfill \end{array}$$
where $`i=1,\mathrm{},n_V`$ and $`A`$ denotes the Minkowski norm with signature $`(+\mathrm{})`$. In this basis, assigning odd parity to the $`h^\alpha `$ is equivalent to assigning odd parity to the $`\varphi ^\alpha `$, and truncation of the scalar $`h_{\stackrel{~}{x}}^\alpha \varphi ^{\stackrel{~}{x}}`$ is obtained by $`\varphi ^\alpha =0`$.
Let’s contrast this with another example: the “generic non-Jordan” family MESGT of MESGTs with cubic polynomial
$$𝒱=\sqrt{2}\stackrel{ˇ}{\xi }^0(\stackrel{ˇ}{\xi }^1)^2\stackrel{ˇ}{\xi }^1\underset{ȷ}{}(\stackrel{ˇ}{\xi }^ȷ)^2,$$
with solution to $`𝒱=1`$ scalarmfds :
$$\begin{array}{cc}\hfill \stackrel{ˇ}{h}^0=& \frac{1}{\sqrt{3}(\varphi ^1)^2}+\frac{1}{\sqrt{3}}\varphi ^1\underset{ȷ}{}(\varphi ^ȷ)^2\hfill \\ \hfill \stackrel{ˇ}{h}^1=& \left(2/3\right)^{\frac{3}{2}}\varphi ^1\hfill \\ \hfill \stackrel{ˇ}{h}^ȷ=& \left(2/3\right)^{\frac{3}{2}}\varphi ^1\varphi ^ȷ\hfill \end{array}$$
where $`ȷ=2,\mathrm{},n_V`$ and $`\stackrel{ˇ}{h}^I`$ is not in the canonical basis.
Remarks
(i) The above solution requires that $`\stackrel{ˇ}{h}^0`$ and $`\stackrel{ˇ}{h}^1`$ must have the same parity.
(ii) The truncation of scalars $`h_{\stackrel{~}{x}}^\alpha \varphi ^{\stackrel{~}{x}}`$ still corresponds to the truncation $`\varphi ^\alpha =0.`$
The first remark tells us that assigning parities consistently in a given YMESGT requires a choice of basis $`\xi ^I`$, and the solution to the condition $`𝒱=1`$. As for the second remark, the requirement that the scalar combination $`\stackrel{ˇ}{h}_{\stackrel{~}{x}}^\alpha \varphi ^{\stackrel{~}{x}}`$ be odd, and so vanish at the orbifold fixed points, allows for a continuous family of vacua parametrized by the scalar vevs. For example, if $`A_\mu ^2`$ has even parity, then $`\stackrel{ˇ}{h}_{\stackrel{~}{x}}^2\varphi ^{\stackrel{~}{x}}`$ must have odd parity so that there is a family of vacua with the direction normal to the flows being $`\varphi ^1_{\varphi ^2}+\varphi ^2_{\varphi ^1}`$ (this is the direction in which the propagating scalar is truncated). However, as $`\stackrel{ˇ}{h}^2`$ must also be odd, this requires $`\varphi ^2`$ to be odd so that it vanishes at the orbifold fixed points. The point is that the theory is simply restricted to $`\varphi ^2=0`$ vacua, connected to the basepoint $`\stackrel{ˇ}{h}^I=(1,0,\mathrm{},0)`$ of the original 5D theory.
We now turn briefly to the implications of odd parity assignments for the $`h^I`$. Once we assign odd parities to a set $`\xi ^\alpha `$, these must vanish on the fixed planes in a basis independent way (as well as $`h^\alpha \xi ^\alpha |_{𝒱=1}`$). For this to be true, the theory must lie on the boundary of the classical Kähler cone, so that a consistent formulation of supergravity requires additional massless states located at the fixed points. These, in turn, admit a higher dimensional interpretation as 11D supergravity CYcomp membranes wrapping collapsing Calabi-Yau 2-cycles W96b . The construction of consistent minimal ($`𝒩=2`$) five-dimensional supergravity on the extended Kähler cone was studied in mohaupt ; but in the present case, the theory on the fixed planes is 4D $`𝒩=1`$, which is a less restrictive framework to work with. In addition, the quantum Kähler cone is generally different from the classical one AGM . Therefore, the determination of these additional massless states is outside the scope of the present paper. We end by noting that, due to the collapsing 2-cycles, the dual CY 4-cycles may collapse to 2- or 0-cycles; this is encoded in the functions $`h_I=C_{IJK}h^Jh^K`$, which are rescaled 4-cycle volumes. Therefore, the choice of jumping function for odd $`C_{IJK}`$ determines whether the 4-cycle collapses completely, and one must choose appropriately.
## VI Extension to $`\mathrm{\Gamma }=_2\times _2`$
There are a couple of phenomenological issues that make the $`S^1/_2`$ orbifold models too simplistic. First, there are always massless chiral multiplets in real representations when a gauge group is broken at the orbifold fixed planes (though these may contain MSSM Higgs fields). Second, all chiral multiplets come in complete representations of the 5D gauge group, which can lead to unwanted fields charged under the Standard Model gauge group. The boundary conditions described by an $`S^1/(_2\times _2)`$ ph/0011311 construction are for the most part capable of resolving these issues.
An exception is the tensor sector: although there is a choice in assignment of parity for the symplectic form $`\mathrm{\Omega }_{MN}`$, we cannot assign $`(+)`$ parity under $`_2\times _2`$ action (it leads to inconsistencies in assignments for the fields). Furthermore, given a choice of $`\mathrm{\Omega }_{MN}`$ parity, there wasn’t a choice of parity assignments in the $`\mathrm{\Gamma }=_2`$ case since supersymmetry dictated the results. Therefore, the situation with tensors is no different in the $`S^1/(_2\times _2)`$ construction. This means that, e.g., tensor multiplets do not allow a doublet-triplet resolution via parity assignments (see the example in section III).
An expansion of a field $`\mathrm{\Phi }(x,x^5)`$ on $`S^1/(_2\times _2)`$ will be of the form
$$\begin{array}{cc}\hfill \mathrm{\Phi }^{(++)}(x,x^5)& =\underset{n}{}\mathrm{\Phi }^{(n)}(x)\mathrm{cos}\left[2nx^5/R\right]\hfill \\ \hfill \mathrm{\Phi }^{(+)}(x,x^5)& =\underset{n}{}\mathrm{\Phi }^{(n)}(x)(A_{(n)}\mathrm{cos}[(2n+1)x^5/R]+\mathrm{}\hfill \\ \hfill \mathrm{\Phi }^{(+)}(x,x^5)& =\underset{n}{}\mathrm{\Phi }^{(n)}(x)(C_{(n)}\mathrm{sin}[(2n+1)x^5/R]+\mathrm{}\hfill \\ \hfill \mathrm{\Phi }^{()}(x,x^5)& =\underset{n}{}\mathrm{\Phi }^{(n)}(x)(E_{(n)}\mathrm{sin}[2nx^5/R]+\mathrm{},\hfill \end{array}$$
where ellipses denote terms with $`ϵ(x^5)`$ factors as in eqn (1), with expansion constants $`B_{(n)},D_{(n)},F_{(n)}`$, respectively. Once again, bosonic fields cannot have $`ϵ`$ factors since the upstairs picture Lagrangian and equations of motion would involve $`\delta ^2`$, where $`\delta `$ is the Dirac distribution. For those, we must set $`B_{(n)}=D_{(n)}=F_{(n)}=0`$. Therefore, bosonic $`\mathrm{\Phi }^{(+)}(x,x^5)`$ vanish at $`x^5=0`$ and bosonic $`\mathrm{\Phi }^{(+)}(x,x^5)`$ vanish at $`x^5=\pi R/2`$; fermionic fields are not generally well-defined on the orbifold fixed planes. Let $`P(\mathrm{\Phi })`$ be the parity of $`\mathrm{\Phi }`$ under the first $`_2`$ factor, and $`P^{}(\mathrm{\Phi })`$ denote the parity under the second factor. Taking the covering space to be $`[\pi R,\pi R]`$ (with $`\{\pi R\}\{\pi R\}`$) as before, the orbifold now has fixed points at $`\{0\},\{\pi R/2\}`$.
### VI.1 Vector sector
In the previous sections, we made an index split for quantities with $`\pm 1`$ parity under the single $`_2`$. We will make a further index splitting for quantities with the four possible values $`\{\pm 1,\pm 1\}`$ for the parity $`\{P(\mathrm{\Phi }),P^{}(\mathrm{\Phi })\}`$:
$$i=(\alpha ,\alpha ^{},a,a^{})\stackrel{~}{p}=(\rho ,\rho ^{},p,p^{}).$$
A given assignment of $`_2\times _2`$ parity to an object will consist of the union of two assignments in the $`S^1/_2`$ construction. Fields from the 5D vector multiplets will have the following assignments:
| $`++`$ | $`+`$ | $`+`$ | $``$ |
| --- | --- | --- | --- |
| $`A_\mu ^\alpha `$ | $`A_\mu ^\alpha ^{}`$ | $`A_\mu ^a`$ | $`A_\mu ^a^{}`$ |
| $`A^a^{}`$ | $`A^a`$ | $`A^\alpha ^{}`$ | $`A^\alpha `$ |
| $`h^a^{}`$ | $`h^a`$ | $`h^\alpha ^{}`$ | $`h^\alpha `$ |
| $`\delta ^\rho `$ | $`\delta ^\rho ^{}`$ | $`\delta ^p`$ | $`\delta ^p^{}`$ |
| $`\gamma ^p^{}`$ | $`\gamma ^p`$ | $`\gamma ^\rho ^{}`$ | $`\gamma ^\rho `$ |
Note: the bare graviphoton $`A_\mu ^0`$ always has $`()`$ parity (so $`A^0`$ has (++) parity). The range of $`\mathrm{}_1`$, $`\mathrm{}_2`$, and $`\mathrm{}_3`$ in $`\alpha =1,\mathrm{},\mathrm{}_1`$; $`\alpha ^{}=\mathrm{}_1+1,\mathrm{},\mathrm{}_2`$; $`a=\mathrm{}_2+1,\mathrm{},\mathrm{}_3`$; and $`a^{}=\mathrm{}_3+1,\mathrm{},n_V`$, are arbitrary.
The fields with $`(+)`$ or $`(+)`$ eigenvalues have massive $`n=0`$ modes on the fixed planes for the same reason that any Kaluza-Klein field does: there is excitation in the $`x^5`$ direction. In the low energy effective theory, such fields will fall into massive $`𝒩=1`$ multiplets in four dimensions due to terms in the Lagrangian with $`_5\mathrm{\Phi }^+`$ or $`_5\mathrm{\Phi }^+`$.
In contrast to the $`S^1/_2`$ construction, we can now choose to make massive multiplets out of unwanted light chiral multiplets in real representations by choosing there to be no $`a^{},p^{}`$ indices. Alternatively, we can keep a subset of those massless chiral multiplets (in a real representation) such that they no longer furnish complete $`K`$-representations. The multiplets with propagating modes on the boundaries and their properties are listed below (we have decomposed the representation $`R_V[K]=\text{adj}[K_\alpha ]R_V^1[K_\alpha ]R_V^2[K_\alpha ]R_V^3[K_\alpha ]`$):
| Multiplet | Representation | Type | Boundary | Mass |
| --- | --- | --- | --- | --- |
| $`\{A_\mu ^\alpha ,\lambda ^{\rho i}\}`$ | $`\text{adj}[K_\alpha ]`$ | Real | Both | Massless |
| $`\{\overline{\lambda }^{p^{}i},z^a^{}\}`$ | $`R_V^1[K_\alpha ]`$ | Real | Both | Potential |
| $`\{\mathrm{\Psi }_5^i,z^0\}`$ | $`K_\alpha `$-singlet | Real | Both | Massless |
| $`\{A_\mu ^\alpha ^{},\lambda ^{\rho ^{}i}\}`$ | $`R_V^2[K_\alpha ]`$ | Real | $`y=0`$ | $`𝒪(1/R)`$ |
| $`\{\overline{\lambda }^{pi},z^a\}`$ | $`R_V^3[K_\alpha ]`$ | Real | $`y=0`$ | $`𝒪(1/R)`$ |
| $`\{A_\mu ^a,\lambda ^{pi}\}`$ | $`R_V^3[K_\alpha ]`$ | Real | $`y=\pi R`$ | $`𝒪(1/R)`$ |
| $`\{\overline{\lambda }^{\rho ^{}i},z^\alpha ^{}\}`$ | $`R_V^2[K_\alpha ]`$ | Real | $`y=\pi R`$ | $`𝒪(1/R)`$ |
Example
Let’s revisit the $`SU(5,1)`$ example based on the Lorentzian Jordan algebra $`J_{(1,5)}^{}`$. We can obtain chiral multiplets (with a scalar potential) in the $`(\mathrm{𝟏},\mathrm{𝟐})(\mathrm{𝟏},\overline{\mathrm{𝟐}})`$ of $`SU(3)\times SU(2)\times U(1)`$ (along with a spin-1/2 gauge singlet multiplet). Let the indices correspond to:
| $`I`$ | $`\alpha `$ | $`a^{}`$ | $`\alpha ^{}`$ | $`0`$ |
| --- | --- | --- | --- | --- |
| $`SU(5,1)`$ | $`SM`$ | $`(1,2)(1,\overline{2})`$ | $`(3,1)(\overline{3},1)`$ | (1,1) |
| | Gauge | | $`(3,2)(\overline{3},2)`$ | |
The $`A_\mu ^\alpha `$ correspond to Standard Model gauge fields propagating on both fixed planes; the remaining vector fields either sit in massive multiplets, or are simply projected out. In particular, we take the $`A_\mu ^\alpha ^{}`$ to be the $`(3,2)(\overline{3},2)`$ vectors (X, Y bosons) and color triplet vectors $`(3,1)(\overline{3},1)`$, which will propagate in massive supermultiplets in the effective theory of the $`y=0`$ plane. This implies that massive spin-1/2 multiplets in the $`[(3,2)(3,1)][c.c.]`$ will propagate in the effective theory of the $`y=\pi R`$ plane. Next, let the $`A_\mu ^a^{}`$ denote the vectors in the $`(1,2)(1,\overline{2})`$, which means there will be chiral multiplets in this representation at both fixed planes (with scalar potential terms). Finally, we get conjugate pairs of massless chiral gauge singlet multiplets from the 5D supergravity multiplet. There are no fields with index $`a`$ in this example.
### VI.2 Hypermultiplet sector
So far, we have not been able to obtain massless chiral multiplets in complex representations of the boundary gauge group. Once again, the only way to do this (starting only from a 5D bulk theory) is to couple 5D hypermultiplets. We can make an index split as in the previous cases:
$$\stackrel{~}{X}=(X,X^{},\mathrm{\Omega },\mathrm{\Omega }^{})A=(n,n^{},\stackrel{~}{n},\stackrel{~}{n}^{}),$$
where the fields have the following parity assignments under $`_2\times _2`$:
| $`++`$ | $`+`$ | $`+`$ | $``$ |
| --- | --- | --- | --- |
| $`q^X`$ | $`q^X^{}`$ | $`q^\mathrm{\Omega }`$ | $`q^\mathrm{\Omega }^{}`$ |
| $`\xi _1^n`$ | $`\xi _1^n^{}`$ | $`\xi _1^{\stackrel{~}{n}}`$ | $`\xi _1^{\stackrel{~}{n}^{}}`$ |
| $`\xi _2^{\stackrel{~}{n}}`$ | $`\xi _2^{\stackrel{~}{n}^{}}`$ | $`\xi _2^n`$ | $`\xi _2^n^{}`$ |
The fields with $`(+)`$ and $`(+)`$ eigenvalues have massive $`n=0`$ modes, and so should fall into massive spin-1/2 multiplets. Therefore, the indices $`n`$ and $`\stackrel{~}{n}`$ are required to be in 1-1 correspondence as are the indices $`n^{}`$ and $`\stackrel{~}{n}^{}`$. However, there is no constraint between unprimed and primed indices, and each pair has an arbitrary range. If $`K`$ is the 5D gauge group, and $`K_\alpha `$ is the boundary gauge group, the $`K_\alpha `$-representations of the massless chiral multiplets at the boundaries no longer need to form complete $`K`$-representations.
The multiplets with boundary propagating modes are listed below, along with some of their properties. Start with $`n_H`$ 5D hypermultiplets in the $`R_H[K]`$ of the 5D gauge group $`K`$. Let the gauge group at the orbifold fixed points be $`K_\alpha `$ so that under this group, the $`R_H[K]`$ decomposes into the representation
$$R_H[K_\alpha ]=R_H^1[K_\alpha ]R_H^2[K_\alpha ],$$
where the indices $`1`$ and $`2`$ denote the splitting of $`\stackrel{~}{X}`$ into $`(X,X^{})`$. At the fixed points, the hypermultiplets split into chiral multiplets with indices split into $`(X,\mathrm{\Omega }^{};X^{},\mathrm{\Omega })`$, and are in the representations $`R_H^i[K_\alpha ]`$ or $`R_H^i[K_\alpha ]\overline{R}_H^i[K_\alpha ]`$:
| Multiplet | Representation | Type | Boundary | Mass |
| --- | --- | --- | --- | --- |
| $`\{\xi ^n,q^{X_1}\}`$ | $`R_H^1[K_\alpha ]`$ | $``$ or $``$ | Both | Potential |
| $`\{\xi ^{\stackrel{~}{n}},q^{X_2}\}`$ | $`\overline{R}_H^1[K_\alpha ]`$ | $``$ or $``$ | Both | Potential |
| $`\{\xi ^A^{},q^X^{}\}`$ | $`R_H^2[K_\alpha ]c.c.`$ | $``$ | $`y=0`$ | $`𝒪(1/R)`$ |
| $`\{\xi ^A^{},q^\mathrm{\Omega }\}`$ | $`R_H^2[K_\alpha ]c.c.`$ | $``$ | $`y=\pi R`$ | $`𝒪(1/R)`$ |
We have split $`X=(X_1,X_2)`$ such that $`X_1=1,\mathrm{},n_H`$ and $`X_2=n_H+1,\mathrm{},2n_H`$. Also, $`A^{}=(n^{},\stackrel{~}{n}^{})`$ is a $`USp(2m)`$ index.
Example
Consider the $`SU(5)`$ YMESGT with $`C_{IJK}`$ as in (5) (where $`C_{ijk}`$ are the $`d`$-symbols of $`SU(5)`$), coupled to the minimal amount of Higgs and matter content in the bulk. This can be realized by coupling the YMESGT to hypermultiplets whose scalars parametrize the quaternionic manifold <sup>7</sup><sup>7</sup>7By allowing an additional singlet hypermultiplet, we can instead couple the exceptional scalar manifold $`\frac{E_8}{SU(6)\times SU(2)}`$.
$$_Q=\frac{SU(27n,2)}{SU(27n)\times SU(2)\times U(1)},$$
resulting in a coupling of $`n`$ sets of hypermultiplets in the $`\mathrm{𝟏}3(\mathrm{𝟓})\mathrm{𝟏𝟎}`$ of $`SU(5)`$ M05a .
Suppose we are going to break $`SU(5)SU(3)\times SU(2)\times U(1)`$; focussing on the hypermultiplet sector, we can make the following assignments
| $`n`$ | $`\stackrel{~}{n}`$ | $`n^{}`$ | $`\stackrel{~}{n}^{}`$ |
| --- | --- | --- | --- |
| $`Matter(1,2)(1,\overline{2})`$ | $`c.c.`$ | $`(3,1)(\overline{3},1)`$ | $`c.c.`$ |
This will result in a low energy theory at both boundaries with Standard Model chiral matter multiplets; a pair of left-chiral Higgs doublets and their CPT conjugates; and a pair of massive spin-1/2 color triplet multiplets, all at both boundaries.
## VII Conclusion
We have found parity assignments for fields and other objects in five-dimensional Yang-Mills-Einstein supergravity coupled to tensor and hypermultiplets on $`M_4\times S^1/\mathrm{\Gamma }`$, allowing for general gauge symmetry breaking. We have used the dimensionally reduced Lagrangians, truncating our attention to the zero mode sector of the theory, though the parity assignments are true for the full supergravity theory.
For $`\mathrm{\Gamma }=_2`$, the generic result is that bulk gauge multiplets for a group $`K`$ in the bulk yield boundary gauge multiplets for the remaining group $`K_\alpha `$, and chiral multiplets in $`K/K_\alpha `$ forming a real representation of $`K_\alpha `$; bulk hypermultiplets yield chiral multiplets on the boundary forming either real or complex $`K_\alpha `$-reps; and bulk tensor multiplets may yield either chiral multiplets in real $`K_\alpha `$-reps or massive charged vector multiplets, depending on the parity for the components of the associated symplectic metric. The boundary theories involving massive vector multiplets charged under $`K_\alpha `$ are the analogues of the $`𝒩=2`$ dimensionally reduced theories in GMZ05b .
The extension to the case $`\mathrm{\Gamma }=_2\times _2`$ allows for light boundary chiral multiplets in incomplete $`K`$-representations except those coming from 5D tensor multiplets. This is the usual mechanism for ‘doublet-triplet’ splitting in orbifold-GUTs.
A novel feature of supergravity theories is that one can gauge a non-compact group. The vectors representing the non-compact generators must be given odd parity, yielding a 4D compact gauge group and chiral multiplets in real representations.
While not examined here, the classical scalar potential on the boundaries can be determined from the form of the parity assignments and the dimensionally reduced Lagrangians, which are available in the literature quoted. There is a positive-definite contribution from the 5D pure YMESGT and tensor sector, while hypermultiplet couplings can contribute negative scalar potential terms. The potentials from compact and non-compact gauge multiplets are distinct. The overall potential can then be analyzed for groundstates spontaneously breaking supersymmetry or electroweak symmetry.
Although boundary-localized fields are not presumed, assigning even parity to 4D $`𝒩=1`$ vector multiplets requires the moduli space of the boundary theory to lie on the boundary of the clasical Kähler cone. Therefore, a consistent supergravity description requires boundary-localized fields to be added, which in a compactified 11D supergravity description correspond to membrane states that become massless when wrapping collapsed Calabi-Yau 2-cycles (the collapse occuring over the fixed points of the effective 5D theory). In particular, gauge symmetry enhancement may occur at these points.
## VIII Future Directions
We will examine some issues regarding symmetries, anomalies, and the presence of the singlet scalars appearing in the orbifold theory. There is always a QCD-type axion present in 5D supergravity on $`S^1/\mathrm{\Gamma }`$, and we will describe the form of the couplings.
Now that parities of various objects appearing in supergravity have been listed, one can examine the terms in the Lagrangian, including $`_5\mathrm{\Phi }`$ terms and some of the fermion interactions that were mostly neglected here. In the upstairs picture, as in the case of simple supergravity JB , this will result in a bulk Lagrangian along with terms supported only at the orbifold fixed points. In the downstairs picture, this simply corresponds to a 5D Lagrangian with boundary conditions on the fields. This will allow us to examine the supersymmetry transformation laws for the fields in the presence of orbifold fixed points, which will lead to a study of global solutions to the Killing spinor equations. Much work has been done using such an analysis in the case of gauged simple supergravity (see e.g. JB and references therein), mostly in the context of braneworld scenarios.
Acknowledgements
The author would like to thank Jonathan Bagger for enlightening conversation, and Murat Günaydin for helpful comments regarding the writing of the manuscript.
Appendices
## Appendix A Notation and Conventions
We use the following notation for representations of Lie groups. Consider a set of $`m`$ 5D $`𝒩=2`$ hypermultiplets, which contain $`4m`$ real scalars that form the $`(𝐦\overline{𝐦})(𝐦\overline{𝐦})`$ of a group $`G`$ (assuming that we are not dealing with pseudoreal representations), where the dimension of the representation is real. Since the scalars form $`m`$ quaternions, the representation can be written as $`𝐦_{}`$. We will simply say that the $`m`$ hypermultiplets are in the $`𝐦`$ of $`G`$, dropping the subscript. On the other hand, consider $`n`$ 4D $`𝒩=1`$ left-chiral multiplets in the $`𝐧`$ of $`G`$ and their $`n`$ right-chiral conjugate partners in the complex conjugate $`\overline{𝐧}`$ of $`G`$, where again we use the real dimension of the representation. If the chiral multiplets all combine with their conjugates to form massive spin-1/2 multiplets, we will use the complex dimension $`n_{}`$ and simply say the spin-1/2 multiplets are in the $`𝐧`$ of $`G`$, dropping the subscript.
We use the signature $`\eta _{mn}=diag(,+,+,+,+)`$. For gamma matrices, we take
$$\mathrm{\Gamma }^\mu =\left(\begin{array}{ccc}0& & \sigma ^m\\ \sigma ^m& & 0\end{array}\right)\mathrm{\Gamma }^5=\left(\begin{array}{ccc}i& & 0\\ 0& & i\end{array}\right)$$
where $`\sigma ^m`$ are the spacetime Pauli matrices, and $`m`$ is a flat spacetime index. The charge conjugation matrix is taken to be
$$C=\left(\begin{array}{ccc}e& & 0\\ 0& & e\end{array}\right)\text{where}e=\left(\begin{array}{ccc}0& & 1\\ 1& & 0\end{array}\right).$$
The charge conjugation matrix therefore satisfies
$$C^T=C=C^1\text{and}C\mathrm{\Gamma }^mC^1=(\mathrm{\Gamma }^m)^T.$$
In five spacetime dimensions, there are a minimum of eight supercharges; since there is a global $`SU(2)_R`$ symmetry, it is convenient to use symplectic-Majorana spinors, which form an explicit $`SU(2)_R`$ doublet. Given a 4-component spinor $`\lambda `$, the Dirac conjugate is defined by
$$\overline{\lambda }^i=(\lambda _i)^{}\mathrm{\Gamma }^0,$$
where $`i`$ is an $`SU(2)_R`$ index, which is raised and lowered according to
$$\lambda ^i=ϵ^{ij}\lambda _j,\lambda _j=\lambda ^iϵ_{ij},$$
with $`ϵ_{12}=ϵ^{12}=1`$. Then a symplectic-Majorana spinor is one that satisfies
$$\overline{\lambda }^i=\lambda ^{iT}C.$$
We will take the following form for our Majorana spinors showing the 2-component spinor content:
$$\lambda ^1=\left(\begin{array}{c}\xi \\ e\zeta ^{}\end{array}\right)\lambda ^2=\left(\begin{array}{c}\zeta \\ e\xi ^{}\end{array}\right).$$
## Appendix B Parity assignments for fermionic fields
The parity assignments for the components of the symplectic Majorana spinors in eqn (3) and (21) are listed here. Some symplectic pairs $`\lambda ^1,\lambda ^2`$ become left/right chiral spinors on the spacetime boundaries, while others denoted in the paper as $`\overline{\lambda }^1,\overline{\lambda }^2`$ become right/left chiral spinors.
| Even | Odd |
| --- | --- |
| $`\alpha _\mu \alpha _\mu ^{}`$ | $`\beta _\mu \beta _\mu ^{}`$ |
| $`\beta _{\dot{5}}\beta _{\dot{5}}^{}`$ | $`\alpha _{\dot{5}}\alpha _{\dot{5}}^{}`$ |
| $`\gamma ^p\gamma ^p`$ | $`\delta ^p\delta ^p`$ |
| $`\delta ^\rho \delta ^\rho `$ | $`\gamma ^\rho \gamma ^\rho `$ |
| $`\eta \eta ^{}`$ | $`\zeta \zeta ^{}`$ |
| $`\zeta _1^n\zeta _2^{\stackrel{~}{n}}`$ | $`\zeta _2^n\zeta _1^{\stackrel{~}{n}}`$ | |
warning/0506/cond-mat0506323.html | ar5iv | text | # Dynamical projection of atoms to Feshbach molecules at strong coupling
## Abstract
The dynamical atom/molecule projection, recently used to probe fermion pairing, is fast compared to collective fermion times, but slow on the Feshbach resonance width scale. Theory of detuning-induced dynamics of molecules coupled to resonantly associating atom pairs, employing a time-dependent many-body Green’s function approach, is presented. An exact solution is found, predicting a $`1/3`$ power law for molecule production efficiency at fast sweep. The results for $`s`$\- and $`p`$-wave resonances are obtained and compared. The predicted production efficiency agrees with experimental observations for both condensed and incoherent molecules away from saturation.
Cold atomic Fermi gases, magnetically tuned to a Feshbach resonance region fermion\_expts , host an intriguing strongly interacting many-body system Timmermans01 ; Holland01 ; Timmermans99 ; Griffin02 . Recently, pairing of fermions near the resonance was probed with the help of dynamical projection of atomic state on the molecular state Regal04 ; Zwierlein04 , achieved by a sweep through the resonance, followed by the detection of molecular Bose-Einstein condensate. The sweep could be made very fast compared to typical fermion time scales, such as the collision frequency or inverse Fermi bandwidth and pairing energy gap, making the process a “snapshot probe” with regard to the collective fermion processes.
On a single particle level, however, broad Feshbach resonances studied in Refs. Regal04 ; Zwierlein04 , exhibit strong atom-molecule coupling in a relatively wide detuning range. In this sense, the sweep speed Regal04 ; Zwierlein04 corresponds to essentially adiabatic atom/molecule conversion, slow on the scale of the resonance width. For example, in the JILA experiment Regal04 , the Feshbach resonance width $`\mathrm{\Delta }B10\mathrm{G}`$ translates into $`\mathrm{\Delta }\nu 180\mathrm{MHz}`$ in detuning frequency, while the characteristic time of the magnetic field sweep, $`\tau _{}=(dt/d\nu )^{1/2}1\mu \mathrm{s}`$, is about $`10^2`$ times longer than $`\mathrm{\Delta }\nu ^1`$. A similar estimate applies to the MIT fast projection experiment Zwierlein04 . Somewhat paradoxically, the fermions involved in this “slow” molecule formation are the same whose many-body state is being analyzed by the dynamical “snapshot” projection. Thus a correct physical picture of the molecular state swept through the resonance must combine the adiabatic single-particle and the “snapshot” many-body aspects in a seamless way.
Continuing efforts to use atom/molecule projection as investigative tool call for better understanding of the driven molecular state. The Landau-Zener model Mies00 ; Goral04 , which fits the data well near saturation Hodby04 , focuses on the adiabatic aspects, ignoring molecule dissociation into continuous spectrum of atom pairs. The dynamical mean field approach Kokkelmans02 ; Kohler03 ; Yurovsky04 ; Duine04 ; Haque04 ; Javanainen04 ; Dicke\_FR ; Goral05 , which can be justified for bosons in the atomic BEC regime, lacks firm foundation in the fermion case. Recently, the many-body state overlap models Diener04 ; Avdeenkov04 ; Perali05 ; Altman05 were put forward. While providing some guidance, these approaches do not account for the experimentally relevant situation of broad resonance Regal04 ; Zwierlein04 when the “snapshot” many-body projection is slow on the scale of individual molecule formation.
Our objective is to describe molecules at a sweep fast compared to the elastic collisions, when only the quantum-mechanical processes are relevant. We develop a theoretical framework which accounts for resonance dissociation/association in the presence of time-dependent detuning as well as for fermion pairing correlations. We describe the molecules swept through the resonance using a time-dependent Green’s function which fully accounts free relative motion of the atoms associating to form molecules. While our method is quite general and applicable to Feshbach resonances with any angular momentum, here we focus, for the sake of concreteness, on the $`s`$-wave case. We consider the evolution from equilibrium at $`\nu =\nu _0`$, followed by an abrupt linear sweep:
$$\nu (t)=\{\begin{array}{cc}\nu _0,\hfill & t<0\hfill \\ \nu _0\alpha t,\hfill & t>0\hfill \end{array}$$
(1)
with $`\alpha `$ the sweep rate. The generalization to the $`p`$-wave resonances Ticknor04 ; Zhang04 ; Schunck05 ; Gurarie04 ; Chevy04 will be straightforward (see below).
Finding the time-dependent Green’s function is a nontrivial mathematical problem, here solved exactly using an idea similar to that of the Wiener-Hopf method. The important time scale, characterizing the adiabaticity of the sweep (see Fig. 1 inset), is found to be
$$\tau _0=\left(\mathrm{}\lambda ^2/\alpha ^2\right)^{1/3},\lambda =g^2m^{3/2}/4\pi \mathrm{}^3,$$
(2)
with $`g`$ the atom-molecule coupling (see Eq.(3)), and $`m`$ the atom mass. The time scale $`\tau _0`$ can also be inferred, as noted by Altman and Vishwanath Altman05 , from the adiabaticity condition $`\dot{\omega }\omega ^2`$ for the time-dependent molecule energy $`\mathrm{}\omega `$. Different regimes arise depending on the relation between $`\tau _0`$ and $`\nu _0/\alpha `$, the time it takes the sweep to reach the resonance (Fig. 1). The atom-to-molecule transformation takes place at times less than $`\tau _0`$ after crossing the resonance, where the evolution is nonadiabatic. At later times, the molecules, dressed by atom pairs, evolve adiabatically. For a fast sweep, $`\alpha \tau _0\nu _0`$, the number of produced molecules scales with the sweep rate as $`\alpha ^{1/3}`$, while for slower sweep, $`\alpha \tau _0\nu _0`$, the number of molecules scales as $`\alpha ^1`$.
These results agree with the molecular number and condensate production efficiency reported by JILA group Regal04 . The sweep speeds $`|dt/dB|1080\mu \mathrm{s}/\mathrm{G}`$ Regal04 correspond to $`\nu _0/\alpha 1100\mu \mathrm{s}`$ with $`\nu _0=0.11\mathrm{G}`$ in the magnetic field units. The characteristic atom-molecule coupling $`\lambda ^21\mathrm{GHz}`$ gives the adiabaticity time $`\tau _01020\mu \mathrm{s}`$ depending on the sweep speed. Thus with $`0.2<\alpha \tau _0/\nu _0<10`$ both the fast and the slow regimes are realized. Indeed, the molecule number obtained for different sweep speeds below saturation (see Fig. 5 in Ref. Regal04 displaying the data for $`\nu _0=0.12\mathrm{G}`$) can be fitted quite accurately with the $`1/3`$ power law dependence, $`N_m|dt/dB|^{1/3}`$, in agreement with our results. Also reasonable, by the order of magnitude, is the predicted total number of produced molecules. Our conclusions regarding the incoherent molecule production channel are consistent with the observed independence of the condensate fraction Regal04 of the sweep speed. We obtain the same production efficiency for condensed and incoherent molecules (Eq.(22)), except near saturation.
Let us recall the form of the two-channel Hamiltonian Timmermans99 , describing pairs of fermions binding to form molecules at the resonance:
$$=_a+_m+\underset{𝐩,𝐩^{}}{}(gb_{𝐩+𝐩^{}}^+a_𝐩a_𝐩^{}+\mathrm{h}.\mathrm{c}.)$$
(3)
with $`_a=_{𝐩\sigma }\frac{p^2}{2m}a_{𝐩\sigma }^+a_{𝐩\sigma }`$, $`_m=_𝐤(\nu +\frac{k^2}{4m})b_𝐤^+b_𝐤`$, and $`a_{𝐩\sigma }`$, $`b_𝐤`$ the atom and molecule operators, $`\sigma `$ the spin ($`\mathrm{}=1`$). The detuning $`\nu `$ is determined by the molecule and two-atom Zeeman energy difference, $`\nu =\mathrm{\Delta }\mu \left(BB_0\right)`$.
The single molecule Green’s function, obtained from Dyson equation Bruun04 , has the form
$$G(\omega ,k)=\frac{1}{\stackrel{~}{\omega }\nu \mathrm{\Sigma }(\stackrel{~}{\omega })},\stackrel{~}{\omega }=\omega \frac{k^2}{4m}+i0,$$
(4)
where $`\mathrm{\Sigma }(\omega )=\lambda (\omega )^{1/2}`$ is the self-energy describing molecule dissociation ($`s`$-wave), which arises after integrating over the 3d density of atom pair states $`N(ϵ)ϵ^{1/2}`$ along with a suitable ultraviolet regularization Griffin02 .
For time-independent $`\nu `$, the molecular state dressed by atom pairs, is described by the Green’s function pole:
$$G_0(\omega )=\frac{Z(\omega )}{\omega \omega (k)+i0}$$
(5)
with $`\omega (k)`$ given by $`\stackrel{~}{\omega }\mathrm{\Sigma }(\stackrel{~}{\omega })=\nu `$. Near the resonance, at $`|\nu |\mathrm{\Delta }E_{}=\lambda ^2,`$ neglecting $`\omega `$ compared to $`\mathrm{\Sigma }(\omega )`$, one obtains molecular energy quadratic in detuning:
$$\omega (k)=(\nu /\lambda )^2+k^2/4m.$$
(6)
At $`\nu <0`$, Eq.(6) gives the energy of molecules, while at $`\nu >0`$ it describes a resonance in the two-fermion scattering mediated by virtual molecules virtual\_resonance . The residue $`Z`$ defines the bare molecule weight in the physical molecule state, $`Z^1(\omega )=dG^1/d\omega =1+\frac{\lambda }{2}(\stackrel{~}{\omega })^{1/2}`$, which varies from zero to one across the resonance, at $`\mathrm{\Delta }E_{}\nu <0`$. At relatively small detuning, $`|\nu /\mathrm{\Delta }E_{}|1`$, $`Z`$ increases linearly: $`Z(\omega )2|\nu |/\lambda ^2`$.
To investigate molecule formation at the resonance, we consider the Green’s function for the problem with time-dependent detuning $`\nu (t)`$. In this case, due to nonlocal character of $`\mathrm{\Sigma }`$ in the time domain, the molecule evolution is described by an integral-differential equation Duine04 ; Haque04
$$\left(i_t\nu (t)\frac{k^2}{4m}\right)b_k(t)\mathrm{\Sigma }_k(t,t^{})b_k(t^{})𝑑t^{}=\eta _k(t)$$
(7)
with $`\eta _k(t)=ge^{ikx}\psi _{}(x,t)\psi _{}(x,t)d^3x`$ the pairing field, and $`\psi _\sigma (x,t)=_pa_{p,\sigma }e^{ipxiϵ_pt}`$. Here the self-energy is
$$\mathrm{\Sigma }_k(t>t^{})=\mathrm{\Sigma }_k(\stackrel{~}{\omega })e^{i\omega (tt^{})}\frac{d\omega }{2\pi }=\frac{ae^{i\frac{k^2}{4m}(tt^{})}}{(tt^{})^{3/2}},$$
(8)
$`a=\frac{\lambda }{2\sqrt{i\pi }}`$, and $`\mathrm{\Sigma }(t<t^{})`$ vanishes due to the causality.
The pairing field $`\eta `$, which acts as a source in Eq.(7), should be taken as a c-number for the condensed molecules (with $`k=0`$), and as an operator for the incoherent molecules. Generally, its correlation function includes both the coherent and incoherent parts:
$$\overline{\eta }_{k,\omega }\eta _{k,\omega }=(2\pi )^4|\eta _0|^2\delta (\omega \mu )\delta (k)+K(\omega ,k),$$
(9)
where $`\eta _0`$ describes a finite amplitude for two fermions to have opposite momenta in the paired state, with $`\mu 2E_F`$ the chemical potential of a pair, and $`K(\omega ,k)=\overline{\eta }_{k,\omega }\eta _{k,\omega }`$ the dynamical pair correlator which is nonzero even for ideal Fermi gas. We first consider the coherent molecule production, treating both $`\eta =\eta _0e^{i\mu t}`$ and $`b(t)`$ as c-numbers. The incoherent pair source $`K(\omega ,k)`$ will be discussed subsequently below.
The evolution problem (7) is non-elementary due to nonlocality of $`\mathrm{\Sigma }(t,t^{})`$. Our approach employs an idea similar to that used in the Wiener-Hopf method. We first handle an auxiliary problem in which the linear sweep $`\nu (t)=\nu _0\alpha t`$ extends from $`\mathrm{}`$ to $`\mathrm{}`$, and then modify the solution to describe the situation of interest (1).
The auxiliary problem in question is to find $`b(t)`$, $`\mathrm{}<t<\mathrm{}`$, which obeys a linear integral-differential equation
$$\left(\widehat{\omega }\nu _0\mathrm{\Sigma }(\widehat{\omega })+\alpha t\right)b(t)=\eta (t),\widehat{\omega }=i_t$$
(10)
with a source term $`\eta (t)`$ of a generic form. It is convenient to go to Fourier representation, in which $`t=i_\omega `$ and the problem is reduced to an ordinary differential equation $`\left(\omega \mathrm{\Sigma }(\omega )\nu _0i\alpha _\omega \right)b(\omega )=\eta (\omega )`$ for $`b(\omega )=e^{i\omega t}b(t)𝑑t`$. Solution of this equation, first order in $`_\omega `$, is found using the gauge transformation $`b(\omega )e^{i\phi (\omega )}b(\omega )`$ with the phase $`\phi `$ satisfying
$$\alpha \phi ^{}(\omega )=D_0(\omega )\omega \mathrm{\Sigma }(\omega )\nu _0.$$
(11)
This problem is solved by the function
$$b(\omega )=i\alpha ^1e^{i\phi (\omega )}_\omega ^+\mathrm{}e^{i\phi (\omega ^{})}\eta (\omega ^{})𝑑\omega ^{}.$$
(12)
To verify (12), one can compare its behavior to that of $`b(\omega )=Ce^{i\phi (\omega )}`$, the solution to the homogeneous problem (10). For $`\omega `$ large and positive, since $`(\omega i0)^{1/2}=i\sqrt{\omega }`$, we obtain the asymptotic behavior $`D_0i\lambda \omega ^{1/2}`$, $`\phi i(2\lambda /3\alpha )\omega ^{3/2}`$, $`e^{\pm i\phi }e^{\pm a\omega ^{3/2}}`$ ($`a=2\lambda /3\alpha `$). Thus $`C=0`$ is required to eliminate exponential growth. Indeed, the asymptotic behavior of the integral in (12) at large positive $`\omega `$ is non-exponential: $`_\omega ^+\mathrm{}e^{i\phi (\omega ^{})}\eta (\omega ^{})𝑑\omega ^{}e^{i\phi (\omega )}\eta (\omega )/\phi ^{}(\omega )`$. (For any physical source, $`\eta (\omega \mathrm{})`$ is algebraic.) At the same time, the behavior at large negative $`\omega `$ does not require special attention: $`\phi `$ is real at $`\omega <0`$, and so the exponentials $`e^{\pm i\phi }`$ oscillate as $`e^{\pm ia(\omega )^{3/2}}`$ without giving rise to “dangerous” asymptotic behavior.
Now, having found the solution for the sweep spanning the entire range $`\mathrm{}<t<+\mathrm{}`$, let us consider the sweep trajectory (1). In this case, it is convenient to represent the function $`b(t)`$ as a sum $`b_<(t)+b_>(t)`$, with $`b_{>,<}(t)`$ nonzero only at $`t0`$ ($`t0`$), respectively, obtained by restricting $`b(t)`$ on the half-line $`t0`$ ($`t0`$). Then the evolution equation, in operator form written as $`\left(\widehat{\omega }\nu (t)\mathrm{\Sigma }(\widehat{\omega })\right)b(t)=\eta (t)`$, can be represented as
$$\widehat{D}_0b_<+(\widehat{D}_0+\alpha \widehat{t})b_>=\eta _<+\eta _>$$
(13)
with $`\eta _{>,<}=\theta (\pm t)\eta (t)`$ having the same meaning as $`b_{>,<}(t)`$, and $`\widehat{D}_0D_0(\widehat{\omega })`$ defined in Eq.(11).
Let us project the terms on the left hand side on the regions $`t0`$, $`t0`$, taking into account the constraints due to causality. The integral operator $`\mathrm{\Sigma }`$ acts only forward, not backward in time, the property explicit in Eq.(8). Thus the function $`\left(D_0(\widehat{\omega })+\alpha \widehat{t}\right)b_>`$ is nonzero only at $`t>0`$, while the function $`D_0(\widehat{\omega })b_<`$ has both the $`t>0`$ and the $`t<0`$ parts. This observation allows to write the problem as two separate problems for $`b_{>,<}(t)`$ as follows:
$$\left[\widehat{D}_0b_<\right]_<=\eta _<,(\widehat{D}_0+\alpha \widehat{t})b_>+\left[\widehat{D}_0b_<\right]_>=\eta _>,$$
(14)
where $`[\mathrm{}]_{<,>}`$ denotes the part of the function at $`t>0`$ ($`t<0`$), with zero value on the opposite half-line. Now, we can solve the first equation for $`b_<`$ and substitute the result in the second equation, which (after some algebra) can be brought to the form
$$(\widehat{D}_0+\alpha \widehat{t})b_>=\widehat{D}_0\left[\widehat{D}_0^1\eta \right]_>.$$
(15)
We note that $`b_>`$ and the function on the right-hand side are both nonzero only at $`t>0`$. This allows to treat this equation as Eq.(10), formally extending the linear time dependence $`\alpha t`$ to negative $`t`$. Using the above result, we obtain the answer in Fourier representation of the form (12) with $`\eta `$ replaced by
$$\stackrel{~}{\eta }(\omega )=\widehat{D}_0\left[\widehat{D}_0^1\eta \right]_>=D_0(\omega )\frac{D_0^1(\omega ^{})\eta (\omega ^{})}{\delta i(\omega \omega ^{})}\frac{d\omega ^{}}{2\pi }.$$
Now, let us consider the source $`\eta (t)=\eta _0e^{i\mu t}`$, describing coherent fermion pairs with the chemical potential $`\mu /2`$ per particle. In this case, $`\eta (\omega )=2\pi \eta _0\delta (\omega \mu )`$ and $`\stackrel{~}{\eta }(\omega )=D_0^1(\mu )\eta _0D_0(\omega )/(\delta i(\omega \mu ))`$. Inserting $`\stackrel{~}{\eta }`$ in Eq.(12), and using the identity $`(\delta i(\omega \mu ))^1=_0^{\mathrm{}}e^{i(\omega \mu )\tau }𝑑\tau `$, we find a closed form representation
$`b(\omega )={\displaystyle \frac{A\eta _0}{\delta i(\omega \mu )}}+\mathrm{\Delta }b(\omega ),`$ (16)
$`\mathrm{\Delta }b(\omega )=iA\eta _0e^{i\phi (\omega )}{\displaystyle _0^{\mathrm{}}}e^{i\mu \tau }\tau {\displaystyle _\omega ^+\mathrm{}}e^{i\omega ^{}\tau i\phi (\omega ^{})}𝑑\omega ^{}𝑑\tau `$
with $`A=D_0^1(\mu )`$. (To obtain (16), we transformed the integral over $`\omega ^{}`$ by writing $`D_0(\omega ^{})=\alpha d(\omega ^{}\tau \phi (\omega ^{}))/d\omega ^{}\alpha \tau `$ and integrating by parts.) Since the first term of (16) gives the would be $`b(\omega )`$ in the absence of the sweep, $`\mathrm{\Delta }b(\omega )`$ describes the effect of the sweep.
Now, let us analyze the asymptotic behavior of $`b(t)=e^{i\omega t}b(\omega )𝑑\omega /2\pi `$ at large positive $`t\tau _0,\nu _0/\alpha `$. In this case, the integral over $`\omega `$ is controlled by large negative $`\omega `$, which can be seen with the help of the stationary phase approximation. Indeed, the saddle point $`\omega _{}`$ of $`\omega t+\phi (\omega )`$, obtained from $`\phi ^{}=t`$, at $`t+\mathrm{}`$ implies $`\omega \mathrm{}`$. With that in mind, we obtain the asymptotic for $`b(t)`$ by setting the lower integration limit in Eq.(16) at $`\omega =\mathrm{}`$, leading to the central result of this work:
$`\mathrm{\Delta }b(t)={\displaystyle \frac{A\eta _0}{2\pi i}}F^{}(t){\displaystyle _0^{\mathrm{}}}e^{i\mu \tau }\tau F(\tau )𝑑\tau ,`$ (17)
$`F(t)={\displaystyle _{\mathrm{}}^+\mathrm{}}e^{i\omega ^{}ti\phi (\omega ^{})}𝑑\omega ^{}.`$ (18)
The qualitative behavior of $`F(t)`$ can be analyzed using the stationary phase approximation. We obtain $`F(t)=\sqrt{2\pi i/\phi ^{\prime \prime }(\omega _{})}e^{i\omega _{}\tau i\phi (\omega _{})}`$, where the stationary phase equation for $`\omega _{}`$, given by $`D_0(\omega _{})+\alpha t=0`$, has a real solution $`\omega _{}=\alpha ^2(tt_0)^2/\lambda ^2`$ only for $`t>t_0=\nu _0/\alpha `$. Relating the curvature $`\phi ^{\prime \prime }`$ and the Green’s function residue, $`\alpha \phi ^{\prime \prime }=D_0^{}=Z^1`$, yields the asymptotic form
$$F(t>t_0)=\left(2\pi i\alpha Z(\omega _{})\right)^{1/2}e^{i\frac{\alpha ^2}{3\lambda ^2}(tt_0)^3}$$
(19)
with $`Z(\omega _{})=2\alpha (tt_0)/\lambda ^2`$. (The self-energy-dominated $`D_0(\omega )=\nu _0\mathrm{\Sigma }(\omega )`$, appropriate for broad $`s`$-wave Feshbach resonance, was used in the above estimates.) Thus $`F(t)`$ grows as $`(tt_0)^{1/2}`$ and oscillates at $`(tt_0)/\tau _01`$, decreasing exponentially at $`tt_0<0`$.
To apply Eq.(17) to the experimental situation we take into account that $`\mu \mathrm{}/\tau _0,\mathrm{}/t_0`$. (Indeed, $`\mu 2E_F`$, with $`E_F=0.35\mu \mathrm{K}=50\mathrm{KHz}`$ in Ref. Regal04 .) We evaluate the integral in (17) using the stationary phase form (19):
$$\mathrm{\Delta }b(t)=(2\alpha Z(t))^{1/2}A\eta _0e^{i\frac{\alpha ^2}{3\lambda ^2}(tt_0)^3}\left(C_1t_0+C_2\tau _0\right),$$
(20)
$`C_1=(i/3)^{1/2}\mathrm{\Gamma }(1/2)`$, $`C_2=(i/3)^{1/6}\mathrm{\Gamma }(5/6)`$. The asymptotic number of molecules at time $`t`$ is evaluated as $`N_m^{(0)}=|\mathrm{\Delta }b(t)|^2`$ (see Fig.1). The fast and slow sweep regimes can be identified: at $`t_0\tau _0`$ we obtain $`N_m\alpha ^{1/3}`$, while at $`t_0\tau _0`$ we have $`N_m\alpha ^1`$.
The incoherent molecule production can be studied in a similar manner. Using the operator $`\eta _k`$ as a source in (8), and averaging over its dynamical correlations (9) we obtain the molecule momentum distribution
$$N_m(k)=\underset{\omega }{}\left|\frac{A}{2\pi }F^{}(t)\underset{0}{\overset{\mathrm{}}{}}e^{i\omega \tau }\tau F(\tau )𝑑\tau \right|^2K(\omega ,k).$$
(21)
with $`A=D_0^1(\omega )`$. As a function of frequency, $`K`$ is nonzero at $`\omega 2E_F`$. At $`E_F\mathrm{}/\tau _0,\mathrm{}/t_0`$, the expression $`|\mathrm{}|^2`$ is $`\omega `$-independent, as above. Factoring it out, we conclude that the condensed and incoherent molecule production efficiencies are identical. The molecule condensate fraction is then expressed through the fermion pair fraction:
$$\frac{N_m^{(0)}}{N_m^{(0)}+_kN_m(k)}=\frac{|\psi _{}(x)\psi _{}(x)|^2}{\widehat{n}_{}(x)\widehat{n}_{}(x)}$$
(22)
$`\widehat{n}_\sigma (x)=\overline{\psi }_\sigma (x)\psi _\sigma (x)`$. We note that the incoherent contribution exists even in the absence of pairing. For ideal fermions at density $`n`$, we have $`_\omega K(\omega ,k)=\frac{1}{2}g^2n(1u)^2(1+u/2)\theta (1u)`$, $`u=k/2p_F`$, which corresponds to a broad molecule momentum distribution with $`k2p_F`$.
The approach presented above yields accurate results for the atom/molecule projection in a wide range of sweep rates, fast and slow, as long as the times $`\tau _0`$, $`t_0`$ are short on the scale of $`E_F`$. The only limitation stems from the assumption of a steady source, which describes the situation when the fraction of atoms converted into molecules is small. The depletion effects, which are different for the condensed and incoherent molecules, can be incorporated in the framework of a quantum kinetic equation.
The above method is applicable to the $`p`$-wave Feshbach resonance case, with the essential modification in the self-energy form $`\mathrm{\Sigma }(\omega )(\omega )^{3/2}`$ Gurarie04 ; Chevy04 . This self-energy is an irrelevant perturbation near the resonance, so we have the atom/molecule conversion time $`\tau _0=\alpha ^{1/2}`$ and $`Z=1`$, as for weak coupling. Thus one can use $`D_0(\omega )=\omega \nu _0`$ in (11), yielding the result identical to (17) with
$$F(t)=e^{i\omega (tt_0)+i\omega ^2/2\alpha }𝑑\omega =(2\pi i\alpha )^{1/2}e^{i\alpha (tt_0)^2/2}.$$
The number of produced molecules, both condensed and incoherent, scales as inverse sweep rate $`\alpha ^1`$. The production is less efficient than in the $`s`$-channel case due to weaker coupling at resonance.
In summary, molecule production at Feshbach resonance is considered as a many-body problem for which the exact Green’s function is obtained using Wiener-Hopf method. The theory is applied to the $`s`$-wave and $`p`$-wave resonances. The slow and fast sweep regimes are identified in the $`s`$-wave case, controled by the adiabaticity time scale (2). The predicted power law $`1/3`$ for the molecule production, as well as the total molecule number, are found to be in agreement with observations away from saturation Regal04 . The independence of the produced condensate fraction on the sweep rate observed at fast sweep Regal04 is also explained by this theory.
We are grateful to Dmitry Petrov for useful comments. |
warning/0506/quant-ph0506251.html | ar5iv | text | # Superbroadcasting of mixed states
## Abstract
We derive the optimal universal broadcasting for mixed states of qubits. We show that the no-broadcasting theorem cannot be generalized to more than a single input copy. Moreover, for four or more input copies it is even possible to purify the input states while broadcasting. We name such purifying broadcasting superbroadcasting.
Broadcasting—namely distributing information over many users—suffers in-principle limitations when the information is quantum, and this poses a critical issue in quantum information theory, for distributed processing and networked communications. For pure states an ideal broadcasting coincides with the so-called quantum cloning, corresponding to an ideal device capable of producing from a finite number $`N`$ of copies of the same state $`|\psi `$ a larger number $`M>N`$ of output copies of the same state, for a given set of input states. Since such a transformation is not isometric, it cannot be achieved by any physical machine on a generally nonorthogonal set of states (this is essentially the content of the no-cloning theorem Wootters82 ; Dieks82 ; Yuen ). The situation is more involved when the states are mixed, since from the point of view of each single user the local mixed state is indistinguishable from the partial trace of an entangled state, and there are infinitely many joint states corresponding to ideal broadcasting. For this reason in the literature fuchs the word broadcasting is used technically to denote a map whose output has identical local states, versus the word cloning used for the case of tensor product of identical states.
Since ideal cloning is not possible, the quantum information encoded on pure states can be broadcast only approximately, and this posed the problem of optimizing the broadcasting e. g. by maximizing an input-output fidelity equally well on all pure states. In the literature this kind of optimized broadcasting is called optimal universal cloning Buzek ; Gisin ; sdc ; Werner . For mixed states the no-cloning theorem is not logically sufficient to forbid ideal broadcasting. In Ref. fuchs the impossibility of broadcasting has been proved in the case of one input copy and two output copies for a set of density operators generally non mutually commuting. Later, in the literature (see, for example, Ref. clifton ) this result has been often implicitly considered as the generalization of the no-cloning theorem to the case of mixed input states. In the present paper we will show that this assertion cannot be generalized to more than a single input copy. In particular, for numbers of input copies $`N4`$ the no-broadcasting theorem does not hold, and it is even possible to purify while broadcasting. We named such a procedure superbroadcasting.
We now present the theoretical derivation of our result.
Let us consider a general broadcasting channel from $`N`$ to $`M`$ copies, namely a completely positive (CP) trace-preserving map from states on $`_{\mathrm{in}}^N`$ to states on $`_{\mathrm{out}}^M`$ that is invariant under permutations of input copies and of output copies. Moreover, we take the broadcasting to be universal, namely the broadcasting map $`\mathrm{B}`$ is covariant under the group of unitary transformations of $``$, more precisely
$$\mathrm{B}(U^N\rho ^NU^{}{}_{}{}^{N})=U^M\mathrm{B}(\rho ^N)U^{}{}_{}{}^{M}.$$
(1)
We will restrict attention to qubits, namely $`^2`$. Upon using the Choi-Jamiolkowsky representation dalop
$$\begin{array}{cc}\hfill R_\mathrm{B}=& \mathrm{B}\mathrm{I}(|II|),\hfill \\ \hfill \mathrm{B}(Q)=& \mathrm{Tr}_{\mathrm{in}}[(I_{\mathrm{out}}Q^\tau )R_\mathrm{B}]\hfill \end{array}$$
(2)
where $`Q`$ denotes a state on $`_{in}`$, and $`R_\mathrm{B}`$ is a positive operator on $`_{\mathrm{out}}_{\mathrm{in}}`$, the covariance condition (1) is equivalent to invariance of $`R_\mathrm{B}`$ under the group representation $`U_g^MU_g^{}^N`$, $`U_g`$ denoting the $`j=\frac{1}{2}`$ representation, for $`g𝕊𝕌(2)`$ \[the symbol $`|I`$ denotes the maximally entangled vector $`|I=_n|n|n`$, and <sup>τ</sup> denotes transposition with respect to the orthonormal basis $`\{|n\}`$\]. In the Choi-Jamiolkowsky representation the trace-preserving condition on the CP map reads
$$\mathrm{Tr}_{\mathrm{out}}[R_\mathrm{B}]=I_{\mathrm{in}},$$
(3)
where $`I_{\mathrm{in}}`$ denotes the identity on $`_{\mathrm{in}}`$. For the unitary group $`𝕊𝕌(2)`$ the complex conjugate representation of any unitary representation, say $`V_g`$, is unitarily equivalent to the direct representation, i. e. $`V_g^{}=CV_gC^{}`$, under the $`\pi `$-rotation $`C`$ around the $`y`$ axis. The explicit form of $`C`$ actually depends on the particular representation $`V_g`$: for the tensor representation $`U_g^N`$ one has $`Ci\sigma _y^N`$. It is then convenient to rewrite the map as follows
$$\mathrm{B}(Q)=\mathrm{Tr}_{\mathrm{in}}[(I_{\mathrm{out}}\stackrel{~}{Q})S_\mathrm{B}]$$
(4)
with
$$\stackrel{~}{Q}CQ^\tau C^{},S_\mathrm{B}(I_{\mathrm{out}}C)R_\mathrm{B}(I_{\mathrm{out}}C^{}),$$
(5)
and now covariance of the CP map $`\mathrm{B}`$ corresponds to invariance of $`S_\mathrm{B}`$ under the representation $`U_g^{(N+M)}`$. A tensor product representation $`U_g^L`$ decomposes into irreducible components according to the Wedderburn decomposition of spaces
$$^L=\underset{j=L/2}{\overset{L/2}{}}_j^{d_j},$$
(6)
where $`x`$ denotes the fractional part of $`x`$ (i. e. $`L/2=0`$ for $`L`$ even and $`L/2=1/2`$ for $`L`$ odd), and the multiplicity $`d_j`$ can be evaluated by recurrence on $`L`$ by adding a qubit at a time, giving $`d_j=\frac{2j+1}{L/2+j+1}\left(\genfrac{}{}{0pt}{}{L}{L/2+j}\right)`$ cirekma . Eq. (6) is also called Clebsch-Gordan series. The spaces $`_j`$ and $`^{d_j}`$ are called representation and multiplicity spaces, respectively. With the above decomposition the group representation writes $`U_g^L=_{j=L/2}^{L/2}U_g^{(j)}I_{d_j}`$, whereas an operator invariant under $`U_g^L`$ has the form $`_{j=L/2}^{L/2}I_jW^{(j)}`$, $`I_j`$ denoting the identity over the representation space $`_j`$, and $`W^{(j)}`$ an operator on the multiplicity space $`^{d_j}`$. On the other hand, an operator invariant under the permutation group $`_L`$ of the $`L`$ copies of the representation has the form $`_{j=L/2}^{L/2}Z_jI_{d_j}`$, where $`Z_j`$ is any operator on the representation space $`_j`$ (this is the so-called Schur-Weyl duality) ford . Since the operator $`S_\mathrm{B}`$ is invariant under $`_M\times _N`$ it must be of the form $`S_\mathrm{B}=_{j=M/2}^{M/2}_{l=N/2}^{N/2}S_{jl}I_{d_j}I_{d_l}`$, where $`S_{jl}`$ is a positive operator over $`_j_l`$. By further decomposing $`_j_l=_{J=|jl|}^{j+l}_J`$ into invariant subspaces and imposing invariance of $`S_\mathrm{B}`$ under $`U_g^{(M+N)}`$, one obtains the general form
$$S_\mathrm{M}=\underset{j=M/2}{\overset{M/2}{}}\underset{l=N/2}{\overset{N/2}{}}\underset{J=|jl|}{\overset{j+l}{}}s_{j,l,J}P_J^{(j,l)}I_{d_j}I_{d_l},$$
(7)
for positive coefficients $`s_{j,l,J}`$, $`P_J^{(j,l)}`$ denoting the orthogonal projector over the irreducible representation $`J`$ coming from the couple $`j,l`$.
The trace preservation condition is now equivalent to
$`\mathrm{Tr}_{\mathrm{out}}[S_\mathrm{M}]=`$ (8)
$`{\displaystyle \underset{j=M/2}{\overset{M/2}{}}}{\displaystyle \underset{l=N/2}{\overset{\frac{N}{2}}{}}}\mathrm{Tr}_j\left[{\displaystyle \underset{J=|jl|}{\overset{j+l}{}}}d_js_{j,l,J}P_J^{(j,l)}\right]I_{d_l}=I_{\mathrm{in}}.`$
Since $`\mathrm{Tr}_j[P_J^{(j,l)}]`$ is invariant under $`U_g^{(l)}`$, one can easily see that $`\mathrm{Tr}_j[P_J^{(j,l)}]=\frac{2J+1}{2l+1}I_l`$, whence the latter condition becomes
$$\underset{l=N/2}{\overset{N/2}{}}\underset{j=M/2}{\overset{M/2}{}}\underset{J=|jl|}{\overset{j+l}{}}d_js_{j,l,J}\frac{2J+1}{2l+1}I_lI_{d_l}=I_{\mathrm{in}},$$
(9)
namely
$$\underset{j=M/2}{\overset{M/2}{}}\underset{J=|jl|}{\overset{j+l}{}}d_js_{j,l,J}\frac{2J+1}{2l+1}=1,N/2l\frac{N}{2},$$
(10)
with positive coefficients $`s_{j,l,J}`$.
Upon writing the input state $`\stackrel{~}{Q}=\stackrel{~}{\rho }^N`$ in the Bloch vector form, we have the decomposition
$$\begin{array}{cc}& \stackrel{~}{\rho }^N=\left[\frac{1}{2}(Ir\stackrel{}{k}\stackrel{}{\sigma })\right]^N\hfill \\ & =(r_+r_{})^{N/2}\underset{l=N/2}{\overset{N/2}{}}\underset{n=l}{\overset{l}{}}\left(\frac{r_{}}{r_+}\right)^n|lnln|I_{d_l},\hfill \end{array}$$
(11)
where $`0r1`$, and $`r_\pm \frac{1}{2}(1\pm r)`$, and $`|ln`$ denotes the eigenstate of the angular momentum component $`\stackrel{}{k}\stackrel{}{J}^{(l)}`$ with eigenvalue $`n`$. From Eq. (10) we see that the broadcasting channels from $`N`$ to $`M`$ make a convex set, with the extreme points classified by functions $`\phi `$ and $`\mathrm{\Phi }`$ corresponding to a given choice $`j=\phi (l)`$, $`J=\mathrm{\Phi }(l)`$, namely to the choice of coefficients
$$s_{j,l,J}^{(\phi ,\mathrm{\Phi })}=\frac{2l+1}{2J+1}\frac{1}{d_j}\delta _{j,\phi (l)}\delta _{J,\mathrm{\Phi }(l)},$$
(12)
or to the Choi-Jamiolkowsky operator
$$S_\mathrm{M}^{(\phi ,\mathrm{\Phi })}=\underset{l=N/2}{\overset{N/2}{}}\frac{2l+1}{2\mathrm{\Phi }(l)+1}\frac{1}{d_{\phi (l)}}P_{\mathrm{\Phi }(l)}^{(\phi (l),l)}I_{d_{\phi (l)}}I_{d_l}.$$
(13)
Using the expression (13) for extremal broadcasting channels and Eq. (11) for the input state we can evaluate the output state
$$\begin{array}{cc}& \mathrm{M}_{(\phi ,\mathrm{\Phi })}(\rho ^N)=(r_+r_{})^{N/2}\underset{l=N/2}{\overset{N/2}{}}\frac{2l+1}{2\mathrm{\Phi }(l)+1}\frac{1}{d_{\phi (l)}}\hfill \\ & \times \underset{n=l}{\overset{l}{}}\left(\frac{r_{}}{r_+}\right)^n\mathrm{Tr}_l[(I_{\phi (l)}|lnln|)P_{\mathrm{\Phi }(l)}^{(\phi (l),l)}]I_{d_{\phi (l)}}.\hfill \end{array}$$
(14)
In terms of Clebsch-Gordan coefficients, this can be rewritten as
$$\begin{array}{cc}& \mathrm{M}_{(\phi ,\mathrm{\Phi })}(\rho ^N)=(r_+r_{})^{N/2}\hfill \\ & \times \underset{l=N/2}{\overset{N/2}{}}\frac{2l+1}{2\mathrm{\Phi }(l)+1}\frac{d_l}{d_{\phi (l)}}\underset{n=l}{\overset{l}{}}\left(\frac{r_{}}{r_+}\right)^n\hfill \\ & \times \underset{m=\phi (l)}{\overset{\phi (l)}{}}\mathrm{\Phi }(l)m+n|\phi (l)m,ln^2|\phi (l)m\phi (l)m|I_{d_{\phi (l)}}.\hfill \end{array}$$
(15)
Now, we are interested in the single output copy, which is the broadcast state. This is given by the partial trace of Eq. (15) over $`M1`$ copies. The evaluation of the partial trace needs the matching between the Wedderburn decomposition and the qubit tensor product representation. According to the Schur-Weyl duality the multiplicity space of the Wedderburn decomposition supports a unitary irreducible representation of the permutation group $`_M`$ of the $`M`$ qubits. Therefore, one has the identity for any operator $`X_j`$ on $`_j^{d_j}`$
$$\underset{l_M}{}\pi _lX_j\pi _l^{}=\frac{M!}{d_j}\mathrm{Tr}_{^{d_j}}[X_j]I_{d_j}$$
(16)
where $`\pi _l`$ denotes the generic permutation. In particular, for $`X_j=|jmjm||11|`$, $`|1`$ denoting any fixed vector of $`^{d_j}`$, one has
$$|jmjm|I_{d_j}=\frac{d_j}{M!}\underset{l_M}{}\pi _lX_j\pi _l^{}$$
(17)
Clearly, one can always choose the given vector of the irreducible representation as cirekma
$$|jm|1=|jm|\mathrm{\Psi }_{}^{\frac{M}{2}j},$$
(18)
where $`|\mathrm{\Psi }_{}`$ denotes the singlet. We can then take the partial trace of both sides of Eq. (17). For each permutation, say $`\pi _s`$, which exchanges the last qubit with one belonging to a singlet, one has $`\mathrm{Tr}_{M1}[\pi _sX_j\pi _s^{}]=\frac{I}{2}`$, and we have $`(M2j)(M1)!`$ permutations of this kind. On the other hand, for each permutation, say $`\pi _m`$, which exchanges the last qubit with one belonging to the $`j`$-multiplet, one has $`\mathrm{Tr}_{M1}[\pi _mX_j\pi _m^{}]=\mathrm{Tr}_{j\frac{1}{2}}[|jmjm|]`$ and there are $`2j(M1)!`$ permutations of this kind. Using the explicit form of the Clebsch-Gordan coefficients one can derive the following identity
$$\mathrm{Tr}_{j\frac{1}{2}}[|jmjm|]=\frac{1}{2}I+\frac{m}{2j}\stackrel{}{k}\stackrel{}{\sigma }.$$
(19)
Substituting the above formula when performing the partial trace of both sides of Eq. (17), one obtains the following expression for the single copy output density operator
$$\begin{array}{cc}& \rho _{(\phi ,\mathrm{\Phi })}^{}(r)=(r_+r_{})^{N/2}\underset{l=N/2}{\overset{N/2}{}}\frac{2l+1}{2\mathrm{\Phi }(l)+1}d_l\underset{m=\phi (l)}{\overset{\phi (l)}{}}\hfill \\ & \times \underset{n=l}{\overset{l}{}}\left(\frac{r_{}}{r_+}\right)^n\mathrm{\Phi }(l)m+n|\phi (l)m,ln^2\frac{1}{2}(I+\frac{2m}{M}\stackrel{}{k}\stackrel{}{\sigma }).\hfill \end{array}$$
(20)
We are now in position to analyse the broadcast state, in particular its Bloch vector. In Eq. (20) we see that the input and the output Bloch vectors are parallel, and clearly $`[\rho ^{},\rho ]=0`$. On the other hand, the length of the output Bloch vector is given by
$$\begin{array}{cc}& r_{(\phi ,\mathrm{\Phi })}^{}(r)=(r_+r_{})^{N/2}\underset{l=N/2}{\overset{N/2}{}}\frac{2l+1}{2\mathrm{\Phi }(l)+1}d_l\hfill \\ & \times \underset{m=\phi (l)}{\overset{\phi (l)}{}}\underset{n=l}{\overset{l}{}}\left(\frac{r_{}}{r_+}\right)^n\mathrm{\Phi }(l)m+n|\phi (l)m,ln^2\frac{2m}{M}\hfill \end{array}$$
(21)
We are now interested in maximizing the length of the output Bloch vector. Since $`r^{}`$ is linear on the convex set of broadcasting channels, we just need to consider extremal maps, and look for the maximum $`r_{opt}^{}(r)=\mathrm{max}_{(\phi ,\mathrm{\Phi })}\{r_{(\phi ,\mathrm{\Phi })}^{}(r)\}`$. It is possible to provenew that the maximal $`r_{(\phi ,\mathrm{\Phi })}^{}(r)`$ is achieved for $`\phi (l)=M/2`$ and for $`\mathrm{\Phi }(l)=\left|l\frac{M}{2}\right|`$, independently on $`r`$. For pure states these optimal maps coincide with those of optimal universal cloning transformationsBuzek ; Gisin ; sdc ; Werner . Also, it can be shownnew that our optimal map gives the same results achievable using the procedure of Ref. cirekma .
As an example, in Fig. 2 we plot the scaling factor $`p(r)=r_{opt}^{}(r)/r`$ for the maps maximizing $`r^{}`$ for $`N=5`$ and several values of $`M`$. One can see that for a wide range of values of $`r`$, one has $`p(r)1`$. This corresponds to a purification of the local states, and since one also has a number of copies at the output $`M>N`$ greater than the number of inputs, it is actually a broadcasting with simultaneous purification, what we call superbroadcasting. Clearly, for $`MN`$ one has more purification than for $`M>N`$, corresponding to the purification protocol cirekma . The superbroadcasting occurs for $`N4`$ input copies. As a rule, one has purification below some value $`r_{}(N,M)`$ of the input purity, for a bounded number $`MM_{}(N)`$ of the output copies.
In Fig. 3 we plot $`r_{}(N,N+1)`$ and $`r_{}(N,M_{}(N))`$ versus the number of input copies $`N`$. After the threshold at $`N=4`$ corresponding to $`r_{}(4,5)=0.787`$, one has a monotonic increase of $`r_{}(N,N+1)`$ and $`r_{}(N,M_{}(N))`$ toward asymptotic purity, with power laws $`2N^2`$ and $`N^1`$, respectively. For larger $`M`$ one has a generally higher threshold for $`N`$, and smaller values of $`r_{}(N,M)`$. For $`N=4`$ one has superbroadcasting for up to $`M=7`$, for $`N=5`$ up to $`M=21`$, and for $`N=6`$ up to $`M=\mathrm{}`$. Notice that perfect broadcasting (corresponding to $`p(r)=1`$) can be achieved under the same conditions of superbroadcasting, (clearly generally by a different map). We remind that we have considered boradcasting of universally covariant sets of mixed states. Indeed, for smaller sets of input states it can be shown that superbroadcasting is possible also for $`N=3`$ input copies (as for equatorial phase-covariant mixed statesnew ), and, for even smaller sets one cannot exclude superbroadcasting also for $`N=2`$.
In conclusion, we have derived the optimal universal broadcasting for mixed states of qubits, optimal in the sense that it maximizes the purity of local states. For pure states and $`M>N`$ the map coincides with the optimal universal cloning transformationBuzek ; Gisin ; sdc ; Werner , whereas for $`NM`$ it is equivalent to the optimal purification map of Ref. cirekma . Thus our optimal broadcasting map generalizes/interpolates between optimal cloning and optimal purification. We have shown that the no-broadcasting theoremfuchs for noncommuting mixed states cannot be generalized to more than a single input copy, and for $`N4`$ input copies one can even purify the state while broadcasting, below some maximum value of the purity. We named such phenomenon superbroadcasting. The possibility of superbroadcasting does not correspond to an increase of the available information about the original input state $`\rho `$, due to detrimental correlations between the local broadcast copies, which does not allow to exploit their statistics. This phenomenon was already noticed in Ref. keylwer , in an asymptotic analysis of the rate of optimal purification procedures. Notice that the correlations alone among qubits cannot be erased by any physical process, since the de-correlating map which sends a state to the tensor product of its partial traces is non linear. From the point of view of single users our broadcasting protocol is actually a purification (for states sufficiently mixed), and the same broadcasting process transfers some noise from the local states to the correlations between them. We think that the present result opens new interesting perspectives in the ability of distributing quantum information in a noisy environment.
This work has been co-founded by the EC under the program ATESIT (Contract No. IST-2000-29681), and QUPRODIS (Contract No. IST-2002-38877). P.P. acknowledges support from the INFM under project PRA-2002-CLON. G.M.D. acknowledges partial support by the MURI program administered by the U.S. Army Research Office under Grant No. DAAD19-00-1-0177. |
warning/0506/astro-ph0506037.html | ar5iv | text | # ISO observations of the Wolf-Rayet galaxies NGC 5430, NGC 6764, Mrk 309 and VII Zw 1911footnote 1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) with the participation of ISAS and NASA
## 1 Introduction
### 1.1 WR galaxies
Wolf-Rayet (WR) galaxies are defined as those galaxies whose integrated spectra contain a broad emission feature at HeII $`\lambda `$4686 Conti (1991). This feature has a full-width half-max (FWHM) of about 10–20Å and is a typical signature of WR stars. Though Seyfert galaxies and active galactic nuclei (AGN) often show a HeII $`\lambda `$4686 emission line, WR galaxies can be distinguished from them by their relatively narrow nebular emission lines. WR galaxies are found exclusively among emission line (EL) galaxies, where the photoionization of the nebular line is stellar in origin, and also possess a very blue continuum which is indicative of a large population of young hot massive stars. The broad HeII $`\lambda `$4686 emission feature is very prominent in the spectra of Galactic and LMC WR stars. By comparing the luminosity and width of this feature as it appears in the spectrum of a WR galaxy, with the corresponding emission lines from Galactic WR stars, an estimate of the number of WR stars in a WR galaxy can be made.
The first galaxy in which this WR feature was discovered Allen et al. (1976) was the blue compact dwarf He 2-10, though NGC 6764 and Mrk 309 were the first objects to be actually referred to as WR galaxies Osterbrock and Cohen (1982). The first comprehensive catalogue was compiled by Conti (1991), and included approximately 40 galaxies. Since then, the number of known WR galaxies has grown rapidly to more than 130. These have been most recently catalogued by Schaerer et al. (1999). Many of the new members of this catalogue show additional features from WR stars in their spectra. For example, the broad emission lines of NIII $`\lambda `$4640 and/or CIII $`\lambda `$4650 as well as CIV $`\lambda `$5808 are among the strongest optical lines in WN and WC stars and are increasingly being detected.
WR galaxies are found among a large variety of morphological types, from low-mass blue compact dwarfs and irregular galaxies, to massive spirals and luminous merging IRAS galaxies. There are systems where WR stars are found in singular giant HII regions (e.g. Tol 89 within the galaxy NGC 5398), in the nucleus or core (e.g. in the barred spiral LINER NGC 6764), in knots (e.g. at the end of the bar in the barred spiral NGC 5430) and in interacting members of compact groups (e.g. HCG 31A and C) O’Halloran et al (2002). Given the wide range of morphological types, an age $`10`$ Myr and high initial masses of $`\mathrm{M}_{ini}`$ $`35`$ $`\mathrm{M}_{}`$ Maeder & Conti (1994), WR galaxies are therefore ideal objects to study the early phases of starbursts, determine burst properties (age, duration, SFR) and to constrain parameters (i.e. slope and upper mass cut-off) of the upper part of the initial mass function. Conversely studies of the stellar populations in super star clusters frequently formed in starbursts and WR galaxies Conti & Vacca (1994); Meurer et al. (1995) can also place constraints on stellar evolution models for massive stars (e.g at extremely low metallicities) which are inaccessible in the Local Group.
### 1.2 Why observe WR galaxies with ISO?
Observations of four out of six WR galaxies which were part of the ISO WRHIIGAL program O’Halloran et al (2000, 2002) are presented in this paper, along with supplemental archival observations from the IRGAL\_1 (PI: Tonaka), GALXISM (PI: Smith), MPEXGAL1 (PI: Genzel) and WMFP15\_A (PI: Wozniak) programs. The main goal of the original WRHIIGAL program was to systematically investigate a sample of WR galaxies in order to try to understand what induces such a massive burst of star formation in these systems, and to investigate the mid and far-IR characteristics of such objects. The fact that WR galaxies are roughly a coeval sample makes comparisons with other WR and starburst galaxies important. Given the importance of the mid-infrared region of the spectrum for the physics of star formation and for exploring the links between massive star formation in galaxies and AGN, acquisition of data in this wavelength range was crucial. However, no other instruments prior to the launch of ISO could provide spatially and spectrally resolved mid-infrared data on a sample of WR and starburst galaxies - ISO data was crucial in opening up this important area of research, such as those obtained using mid-IR data from surveys of starburst and blue compact galaxies in the ISO WRHIIGAL and HAROA programs Steel et al. (1996); O’Halloran et al (2000, 2002); Metcalfe et al. (2005); O’Halloran et al (2005). The ISOCAM, ISOPHOT and LWS observations provided maps and spectra of the galaxies in one, or in combinations of, the Unidentified Infrared Bands (UIB) and dust or nebular line emission. They enabled a determination of the spatial relationship between the optical star forming regions and the infrared emission from the young WR and starburst regions, allowing a discrimination of the relative contribution to the infrared flux from dust, nebular emission and polycyclic aromatic hydrocarbon (PAH) molecules. The advent of the Spitzer Space Telescope, optimized to conduct research in a similar wavelength range but with greater spatial and spectral resolution, will further open up this research field over the forthcoming years.
### 1.3 The sample of WR galaxies
Four galaxies displaying WR features were observed, as part of a wider survey of starburst and WR galaxies (the WRHIIGAL and HAROA programs) Steel et al. (1996); O’Halloran et al (2000, 2002); Metcalfe et al. (2005); O’Halloran et al (2005) by the Infrared Space Observatory, are discussed in this section. Table 1 lists the basic properties of each galaxy.
#### 1.3.1 NGC 5430
NGC 5430 (Mrk 799) is a nearby S-shaped barred spiral galaxy (SBb) located at a distance of 39.8 Mpc. Keel (1982) detected the broad 4650Å emission feature in the knot southeast of the nucleus, at the end of the bar. This knot is an extremely luminous HII region Keel (1982, 1987) though of unusually low luminosity at 4.885 GHz Condon et al. (1982), possibly indicating an anomalously low supernova rate.
#### 1.3.2 NGC 6764
At a distance of 32.4 Mpc, NGC 6764 is a nearby S-shaped barred spiral galaxy (SBbc), somewhat similar to NGC 5430. Originally considered to be a Seyfert 2 galaxy by Rubin et al. (1975), Heckman (1980) later classified NGC 6764 as a LINER on account of its large OI $`\lambda `$6300/OIII $`\lambda `$5007 ratio of $`0.3`$. This galaxy contains a nuclear stellar optical continuum source which extends over $`1.6^{\prime \prime }`$ Rubin et al. (1975).
#### 1.3.3 Mrk 309
At a distance of 174.1 Mpc, Mrk 309 is the most distant object in the sample. It was identified as a galaxy with a bright UV continuum, noticeable H$`\alpha `$ emission and possible Seyfert characteristics Markarian & Lipovetskii (1971).
Mrk 309 appears in the sample of 212 emission-line galaxies extracted from the Universidad Complutense de Madrid (UCM) lists Zamorano et al. (1994) and was described as a Wolf-Rayet dominated nucleus galaxy with various HII regions outside the nucleus. The R-band image of Mrk 309 presented by Vitores et al. (1996) shows that the southern arm is more prominent than the northern and ends in a knot approximately $`12^{\prime \prime }`$ from the nucleus.
#### 1.3.4 VII Zw 19
VII Zw 19 appears in the list of WR galaxies discovered by Kunth and Joubert (1985) from a search among a sample of 45 blue EL galaxies. A broad emission band between 4600Å and 4711Å was the criteria used to decide that this EL galaxy is also a WR galaxy. Beck (2000) notes that the radio and H$`\alpha `$ emission from this galaxy have the same basic structure of a very strong central source embedded in a weak envelope. The central non-thermal source appears to be unresolved and is about 800 times more luminous than Cas A. Recent MERLIN and VLA observations Beck et al. (2004) note that VII Zw 19 resembles M82 in its radio and infrared spectrum, however the starburst region of VII Zw 19 is twice the size and twice as luminous as that of M82. VII Zw 19 is situated at a distance of 65.2 Mpc.
We present ISOCAM, ISOPHOT and LWS observations of these galaxies as part of an observing program consisting of several galaxies exhibiting Wolf-Rayet signatures. The observations and data reduction are presented in Sect. 2. The results are contained in Sect. 3 and discussed in Sect. 4. The conclusions are summarized in Sect. 5.
## 2 Observations and Data Reduction
The ISO Kessler et al. (2003) observations were obtained using the mid-infrared camera ISOCAM Blommaert et al. (2003), the spectrometric mode of the ISO photopolarimeter ISOPHOT Laureijs et al (2003a) and the medium-resolution grating mode of the long wavelength spectrometer LWS Gry et al (2003). The astronomical observing template (AOT) used numbers and the observing log for the ISO observations is presented in Table 2.
### 2.1 ISOCAM
The observations had the following configuration: 1.5, 3 and 6<sup>′′</sup> PFOV (for Mrk 309 LW2/LW3, NGC 5430 & NGC 6764 LW2/LW3 and LW10 observations respectively), integration time of 2.1 s and $``$100 readouts (excluding discarded stabilisation readouts). The diameter of the point spread function (PSF) central maximum at the first Airy minimum is 0.84 x $`\lambda `$($`\mu `$m) arcseconds. The FWHM is about half that amount. A discussion of the data reduction steps for the CAM data can be found in O’Halloran et al. (2002).
### 2.2 PHT-S
PHT-S consists of a dual grating spectrometer with a resolving power of 90 Laureijs et al (2003b). Band SS covers the range 2.5 - 4.8 $`\mu `$m, while band SL covers the range 5.8 - 11.6 $`\mu `$m. Two different chopping modes were used. The NGC 5430 and Mrk 309 PHT-S spectra were obtained by pointing the $`24^{\prime \prime }\times 24^{\prime \prime }`$ aperture of PHT-S alternatively towards the peak of the emission (for 512 seconds) and then towards two background positions off the galaxy (256 seconds each), using the ISOPHOT focal plane chopper. The NGC 6764 and VII Zw 19 PHT-S spectra were obtained on the other hand by operating the PHT-S aperture in rectangular chopping mode. The satellite pointed to a position between the source and an off-source position, and the chopper moved then alternatively between these two positions. The source was always in the positive beam in the spacecraft Y-direction. The calibration of the spectra was performed by using a spectral response function derived from several calibration stars of different brightness observed in chopper mode Acosta-Pulido et al. (2000). The relative spectrometric uncertainty of the PHT-S spectrum is about 20% when comparing different parts of the spectrum that are more than a few microns apart. The absolute photometric uncertainty is $``$ 10% for bright calibration sources. A discussion of the data reduction steps for the PHOT data can be found in O’Halloran et al. (2002).
### 2.3 LWS
The LWS spectra were obtained with the aperture centered on the 6.7 $`\mu `$m ISOCAM map for NGC 5430 and on the nuclear region of the galaxy for NGC 6764. It was assumed that for both objects the source was completely included in the beam of the LWS instrument, so that no extended-source correction was necessary. The beam of LWS was slightly elliptical and its FWHM varied between 65<sup>′′</sup> and 85<sup>′′</sup>, depending on wavelength and direction (Swinyard et al. 1998). The grating was fully scanned 6 times over the entire wavelength range. A spectral sampling interval of 4 was employed to give 4 spectral points per resolution element in each of the scans. All the datasets were first processed with version 10 of the LWS pipeline processing software, OLP V10.1 Gry et al (2003). The LWS Interactive Analysis package LIA Sidher et al (1997) was then used to further process the output of the standard pipeline. The nominal, fixed, dark current values, as determined from the dedicated measurements in revolution 650 Swinyard et al (2000) were subtracted from the data. The data was rebinned to one point per LWS detector, employing a scan-averaging method described by Sidher et al. (2000), yielding 10 bandpass continuum estimates spanning the LWS spectral range, and are plotted in Fig 3.
## 3 Results
### 3.1 ISOCAM maps
#### 3.1.1 NGC 5430
Deconvolved 6.7, 12 and 14.3 $`\mu `$m maps of NGC 5430 overlaid on a DSS image are presented in Figs. 1(I), (II) and (III) respectively. It can clearly be seen that the ISOCAM contours follow the general form of the galaxy, especially on the 12 $`\mu `$m map. In addition three sources were clearly detected in the higher resolution 12 $`\mu `$m map and are labelled A, B and N, though source B was also detected at 6.7 and 14.3 $`\mu `$m. The source B coincides with the highly luminous southeast HII region and contains the majority of the WR population. The weaker source A is probably associated with a less luminous HII region in NGC 5430. The nucleus of NGC 5430 is denoted N in Fig. 1(II). The derived fluxes for NGC 5430 are given in Table 3.
#### 3.1.2 NGC 6764
Deconvolved 6.7, 12 and 14.3 $`\mu `$m maps of NGC 6764 overlaid on a DSS image are presented in Figs. 1(IV), (V) and (VI). Four infrared sources were detected on the higher resolution 12 $`\mu `$m map and are labelled A, B, C and N. The weaker sources A and B are probably associated with star formation regions in the arms of the galaxy, while C may also be a star formation region or an external galaxy. The derived ISOCAM fluxes are given in Table 3.
#### 3.1.3 Mrk 309
The 7.7 and 14.3 $`\mu `$m maps of Mrk 309 are presented in Figs. 1(VII) and 1(VIII). Due to the low signal to noise ratio of the 7.7 $`\mu `$m image, deconvolution was not reliable. The nucleus of Mrk 309 was detected in both maps, though more weakly at 7.7$`\mu `$m than at 14.3 $`\mu `$m. At low flux levels the 7.7 $`\mu `$m contours follow the outline of the galaxy, while the 14.3 $`\mu `$m emission is more compact. The ISOCAM infrared fluxes are given in Table 3.
#### 3.1.4 VII Zw 19
No observations of VII Zw 19 were performed using ISOCAM.
### 3.2 PHT-S spectra
#### 3.2.1 NGC 5430
The PHT-SL spectrum of the nuclear region of NGC 5430 is presented in Fig. 2a. The main PAH bands at 6.2, 7.7, 8.6 and 11.3$`\mu `$m were easily detected, along with the \[ArII\] 6.99 $`\mu `$m feature. Two additional PHT-SL spectra of the regions labelled A and B were also obtained and are presented in Figs. 2b and 2c. In both spectra the PAH bands plus \[ArII\] at 6.99 $`\mu `$m were well detected, along with a blend of features at 10.6 $`\mu `$m in region A. The fluxes for the identified features are given in Table 4.
#### 3.2.2 NGC 6764
Three PHT-SL spectra of the nucleus and regions A and B were obtained and are presented in Figs. 2d-2f. The main PAH bands at 6.2, 7.7, 8.6 and 11.3 $`\mu `$m were easily detected in the nuclear region, though only weakly detected at star formation region A. The spectrum of region B was very noisy and no features were reliably identified. The fluxes for the identified features are given in Table 5.
#### 3.2.3 Mrk 309
The PHT-SL spectrum of Mrk 309, shown in Fig. 2g was very noisy with no reliable detections of the PAH emission bands, and thus flux determinations were not made for features in this spectrum.
#### 3.2.4 VII Zw 19
The PHT-SL spectrum of the nucleus of the VII Zw 19 is presented in Fig. 4h. The main PAH bands at 6.2, 7.7, 8.6 and 11.3 $`\mu `$m are easily detected, along with \[ArII\] at 6.99 $`\mu `$m and a possible S(3) pure rotational line, $`\upsilon `$ = 0-0, of molecular hydrogen at 9.7 $`\mu `$m. The feature at 10.65 $`\mu `$m may be a blend of features. The identified line features and line fluxes are given in Table 5.
### 3.3 LWS spectra
The LWS spectra for the far-infrared lines \[OI\] at 63 $`\mu `$m and \[CII\] at 158 $`\mu `$m for both NGC 5430 and NGC 6764 are presented in Fig. 3, while the line fluxes are listed in Table 6. The continuum source strength, or applicable upper-limits, recorded in the ten LWS detectors by binning data across the spectral range of each detector are listed in Table 7.
## 4 Discussion
### 4.1 Previous surveys
#### 4.1.1 NGC 5430
Keel (1982) estimated that $`5\times 10^3`$ WN7 stars and $`10^4`$ WC8 stars reside in this knot as well as a still larger number of O stars. With improved spectra Keel (1987) later claims that though HeII $`\lambda `$4686 and \[NIII\] $`\lambda `$4640 are detected, \[CIV\] is not and therefore the WR stars in the knot might be regarded as some species of WN stars. Keel (1987) also identifies 26 other “normal” HII regions present in NGC 5430 other than the southeast knot and the nucleus.
According to Keel (1987) the observed star formation in the SE knot is both sudden and transient. The nebula is expanding, probably driven by stellar winds, and dissipating on a timescale of $`10^7`$ years. Star formation elsewhere in the galaxy appears to be proceeding at a normal rate and with a normal HII region luminosity function Keel (1987). The nebular abundances in the SE knot are rather high for a position so far out in a disk or bar and are more like those seen in nuclei. This suggests that this knot might be a separate object resulting from the collision of a dwarf irregular galaxy with the disk and bar of NGC 5430. Furthermore, the H$`\alpha `$ luminosity and emission-line ratios are quite comparable to those seen in such systems as Mrk 108 Keel et al. (1985). Using starburst models by Cervino and Mas-Hesse (1994) and Leitherer and Heckman (1995), Contini et al. (1996) estimated the age of the starburst in the SE knot to be 3 Myr and 4 Myr respectively. For the nucleus the estimates are 8 Myr and 9 Myr respectively. These estimates are certainly compatible with the scenario of separate star formation mechanisms in the nucleus and SE knot of NGC 5430.
#### 4.1.2 NGC 6764
Osterbrock and Cohen (1982) detected WR emission features in the spectrum of the nucleus of this narrow emission line galaxy. From the equivalent widths of the WR emission features at 4650Å, Osterbrock and Cohen (1982) estimate that the number of WN and WC stars is $`5.0\times 10^4`$ and $`9.0\times 10^4`$ respectively. Confirmation of the existence of the WR features detected by Osterbrock and Cohen (1982) was made by Eckart et al. (1996), who claim that the number of WO stars is negligible and WN stars dominate the WR star population. Detection of the broad \[CIII\] $`\lambda `$5696 and \[CIV\] $`\lambda `$5808 lines by Kunth and Contini (1999) indicate the existence of some WC stars. Eckart et al. (1996) estimate that there are a total of 3600 WR stars contributing to the ionizing flux in the nucleus of NGC 6764. These authors also note that the WR feature is spatially extended and most of it originates in the $`1.6^{\prime \prime }`$ diameter nuclear optical continuum source. Using the same method as for NGC 5430, Contini et al. (1996) estimates the age of the starburst in the nucleus of NGC 6764 to be 5.0 Myr and 6.5 Myr depending on which models are used Cervino & Mas-Hesse (1994); Leitherer & Heckman (1995). The presence of WR stars Osterbrock and Cohen (1982) indicates that the starburst is very young ($`6`$ Myr) and that many massive stars were born during the burst. Eckart et al. (1996) conclude that the nucleus of NGC 6764 has recently undergone or is currently undergoing an intense starburst with a characteristic timescale of a few times $`10^7`$ years.
#### 4.1.3 Mrk 309
Arkelian et al. (1972) noted the presence of diffuse H$`\alpha `$ emission, as well as the lines NII $`\lambda `$$`\lambda `$6548, 6583. The presence of these lines were later confirmed by Afanasev et al. (1980). Osterbrock and Cohen (1982) detected a blend of three broad emission features at 4650Å in the nucleus of this galaxy: NIII $`\lambda `$4640, CIV $`\lambda `$4660 and HeII $`\lambda `$4686. These lines are slightly narrower than in NGC 6764 and hence a slightly lower proportion of the continuum radiation at 4650Å comes from WR stars. However, Conti (1991) argues that the emission feature at 4660Å is attributable to FeIII rather than CIV, and is consistent with the deficiency in the strength of the CIII $`\lambda `$4650 emission feature as noted by Osterbrock and Cohen (1982). CIV $`\lambda `$5808 and CIII $`\lambda `$5696 emission, attributed to WC stars, has also been tentatively detected Osterbrock and Cohen (1982). From the equivalent widths of the WR emission features at 4650 Å these authors deduced that 9$`\%`$ of the continuum radiation at this wavelength comes from WR stars and hence calculated the number of WN and WC stars to be approximately $`0.9\times 10^4`$ and $`1.5\times 10^4`$ respectively.
With the exception of the WR features at 4650Å, Osterbrock and Cohen (1982) point out that the emission line spectra of Mrk 309 is typical of low-ionization galaxies with gas in their nuclei apparently photoionized by early-type stars, as in HII regions. Indeed Mazzarella and Balzano (1986) places Mrk 309 in the HII spectral class. Osterbrock and Cohen (1982) concluded that the number of O stars in Mrk 309 is comparable to the number of WR stars and hence massive star formation must have occurred recently.
### 4.2 Dust components and luminosities
There is growing evidence for the existence of several components within the dust distribution of galaxies Klaas et al. (2001). Specifically this can be divided into warm dust components associated with star formation regions and a spatially extended distribution of cold dust. Warm dust emission is generally associated with very small dust grains which are heated by the single-photon absorption process to temperatures up to several hundred Kelvin Desert et al. (1990), and are not in thermal equilibrium with their environment Calzetti et al. (2000). Emission from very small dust grains account for the characteristics of the mid-IR at wavelengths $``$ 40 $`\mu `$m. The cold dust emission is associated with classical large grains, emitting at wavelengths in excess of 80 $`\mu `$m in thermal equilibrium with their environment Calzetti et al. (2000). The large grains are heated by the normal interstellar radiation field. These large dust grains account for practically all the emission longward of 80 $`\mu `$m in galaxies.
To model the dust continuum at mid and far-infrared wavelengths, modified blackbody functions comprised of a greybody component were fitted to the IRAS and ISO fluxes for the galaxies in the WR sample. The spectral energy distributions of NGC 5430, NGC 6764, Mrk 309 and VII Zw 19 using IRAS, ISOCAM, PHT-SL and LWS fluxes are presented in Figs. 4a-d. The dust model, denoted by the solid curve was fitted to the data and contains two separate dust populations along with the PAH bands Boulanger et al. (1998): a warm dust component at 135-179 K and a cooler dust component at 39-64 K and each component is indicated in Figs. 4a-d.
For VII Zw 19 at $`\lambda `$ $``$ 300$`\mu `$m, it is noticeable that the emission is flattening, possibly due to the presence of a very cold dust component. However overall, the results are consistent with similar modelling by Siebenmorgen et al. (1999).
To determine the mass of each dust component, several parameters need to be determined. The luminosity of each component was determined from the integrated fluxes of each greybody, along with the dust temperature and then a dust mass for each component was determined Klaas et al. (2001). Using this method, the IR luminosities and dust component masses in similar starburst galaxies such as the Antennae and NGC 6240 were determined in order to check the validity of the derived values Klaas et al (1997).
For each target, the luminosity for the galaxy, the dust component luminosity and the derived dust masses are given in Table 8. For the warm and cool dust components, the derived values are quite similar to starbursts such as the Antennae, NGC 6240 Klaas et al (1997), NGC 1741 O’Halloran et al (2002) and Mrk 297 Metcalfe et al. (2005).
### 4.3 Current star formation
#### 4.3.1 Diagnostics using ISOCAM data
In order to determine the current state of star formation within our sample, diagnostics such as the 14.3/6.75 $`\mu `$m flux ratio can be used as they act as standard ratio diagnostics for this type of investigation Vigroux et al (1999); Helou (1999). Since the 14.3 $`\mu `$m flux is dominated primarily by emission from very small dust grains and the 6.75 and 7.7 $`\mu `$m fluxes are dominated by PAHs, the 14.3/7.7 $`\mu `$m ratio provides a diagnostic similar to the 14.3/6.75 $`\mu `$m ratio, allowing a determination of the current state of the starburst. These diagnostic ratios generally decrease as interactions develop and starbursts age Vigroux et al (1999); Helou (1999) \- the highly ionizing O stars in the burst die off, thus no longer destroying nearby PAHs, plus emission from nearby dust heated by massive stars also decreases. For the three galaxies with available ISOCAM data, the derived ratios are given in Table 9. The 6.7 $`\mu `$m flux for Mrk 309 was synthesized from the PHT-SL data, a technique used previously when faced with the lack of 6.75 $`\mu `$m ISOCAM data O’Halloran et al (2000, 2002). The obtained 14.3/7.7 14.3/6.75 and 14.3/7.7 $`\mu `$m ratios were indicative of strong ongoing star formation, with ratios generally well above 2, similar to those found in earlier works for strong starbursts Vigroux et al (1999).
#### 4.3.2 Star formation rate
Using the calibrations in Kennicutt (1998) and Wilke et al. (2004), we can calculate the global star formation rate for the sample from the derived IR luminosities. Global SFRs of 9.6, 5.0, 24.1 and 12.6 M yr<sup>-1</sup> were calculated for NGC 5430, NGC 6764, Mrk 309 and VII Zw 19 respectively. These strong levels of ongoing star formation are consistent with the values from the infrared diagnostics.
### 4.4 Do any galaxies within the sample harbour an AGN?
#### 4.4.1 Mid-IR diagnostics
Historically, 2 galaxies of the sample have been suspected of harbouring AGN. Mrk 309 was identified as a galaxy with a bright UV continuum and noticeable H$`\alpha `$ emission by Markarian (1971) who noted the possibility of a Seyfert nucleus. Afanasev (1980) reported H$`\alpha `$ and \[NII\] $`\lambda `$6548, $`\lambda `$6583 and stated that Mrk 309 may be a Seyfert 2 galaxy. NGC 6764 has been classified as a LINER galaxy based on optical spectroscopy Osterbrock and Cohen (1982). NIR spectroscopy of the nuclear region Schinnerer et al. (2000) reveals that star formation is confined to two regions, one with a radius less than 100 pc containing the WR stars. NGC 6764 also shows strong variations in ROSAT X-ray flux density by at least a factor of 2 on timescales of 7 days Schinnerer et al. (2000), suggesting the presence of a compact AGN with R $``$ 10<sup>3</sup> AU. No previous survey has suggested that NGC 5430 or VII Zw 19 harbours an AGN.
In order to test the possibility that the central region harbours an AGN, diagnostic tools using ISO data are available to probe the nature of the activity within the central starburst region. The ratio of the integrated PAH luminosity and the 40 to 120 $`\mu `$m IR luminosity Lu et al. (1999) discriminates between starbursts, AGN and normal galaxies. The lower the ratio, the more likely the galaxy harbours an AGN, due to high very small dust grain emission powered by AGN longward of 10 and shorter than 50 $`\mu `$m, plus PAH destruction near the AGN Vigroux et al (1999). Similarly, the ratio of the 7.7 $`\mu `$m PAH flux to the continuum level at this wavelength can provide a measure of the level of activity within the nucleus Genzel et al. (1998); Clavel et al. (2000); Laureijs et al (2000), as a high ratio indicates the lack of an AGN excited dust component. The derived ratio values are given in Table 9 for each galaxy. Given the high L(PAH)/L(40-120 $`\mu `$m) and F(PAH 7.7 $`\mu `$m)/F(7.7 $`\mu `$m continuum) values, it can be stated that most of these galaxies are home to only a compact burst of star formation and that an AGN is not required within the central bursts. - values of $``$ 1.5 and 0.06 would be indicative of a strongly dominant AGN component. The exception is NGC 6764, whose F(PAH 7.7 $`\mu `$m)/F(7.7 $`\mu `$m continuum) value of 1.22 is consistent with the presence of an AGN, yet the L(PAH)/L(40-120 $`\mu `$m) is more in line with a starburst, a finding which suggests the presence of a compact AGN dominated by a strong starburst component.
#### 4.4.2 Far-IR diagnostics
Atomic oxygen and ionized carbon are the principal coolants of the gaseous interstellar medium via their fine structure lines in the far infrared (FIR). In particular, \[OI\] at 63 $`\mu `$m and \[CII\] at 158 $`\mu `$m dominate the cooling in the photodissociation regions associated with massive stars such as Wolf Rayets, along with \[OI\] at 146 $`\mu `$m and \[OIII\] at 88 $`\mu `$m Malhotra et al. (1997). The \[OI\] and \[CII\] features are also produced in the warm atomic gas behind dissociative shocks, in HII regions or in photodissociation regions (PDRs), while \[OIII\] is more associated with denser environments within HII regions Malhotra et al. (1997); Braine & Hughes (1999).
The \[OI\] 63 $`\mu `$m and \[CII\] 158 $`\mu `$m lines were well detected in both NGC 5430 and NGC 6764 and allowed a determination of the L<sub>CII</sub>/L<sub>FIR</sub> and (L<sub>CII</sub>+L<sub>OI</sub>)/L<sub>FIR</sub> ratios, which could be used to probe the nature of the environment within the galaxy. For NGC 5430, values of L<sub>CII</sub>/L<sub>FIR</sub> = $`1.9\times 10^3`$ and (L<sub>CII</sub>+L<sub>OI</sub>)/L<sub>FIR</sub> = $`1.1\times 10^2`$ were determined, given an IRAS flux ratio F(60 $`\mu `$m)/F(100 $`\mu `$m) = 0.55. For NGC 6764, values of L<sub>CII</sub>/L<sub>FIR</sub> = $`2.9\times 10^3`$ and (L<sub>CII</sub>+L<sub>OI</sub>)/L<sub>FIR</sub> = $`1.9\times 10^2`$ were determined, given an IRAS flux ratio F(60 $`\mu `$m)/F(100 $`\mu `$m) = 0.53. The L<sub>CII</sub>/L<sub>FIR</sub> and (L<sub>CII</sub>+L<sub>OI</sub>)/L<sub>FIR</sub> ratios are consistent with those from other starburst galaxies, given an IRAS flux ratio F(60 $`\mu `$m)/F(100 $`\mu `$m) = 0.66 Malhotra et al. (1997); Braine & Hughes (1999), though for higher dust temperatures and star formation rates these ratios decrease Malhotra et al. (1997).
## 5 Conclusions
Observations of four of a sample of WR galaxies using the Infrared Space Observatory were presented in this paper. ISOCAM maps of NGC 5430, Mrk 309 and NGC 6764 revealed the location of star formation regions in each galaxy, while ISOPHOT spectral observations from 4 to 12 $`\mu `$m detected the ubiquitous PAH bands in the nuclei of the targets, and several of the disk star forming regions. Strong \[OI\] and \[CII\] lines were detected in LWS spectra of NGC 5430 and NGC 6764.
Using a combination of ISO and IRAS flux densities, a dust model based on the sum of two modified blackbody components were successfully fitted to the available data. The modified blackbody functions were comprised of a $`1/\lambda `$ component with temperatures of $``$40–64 K and $``$140–180 K depending on the galaxy. These models accounted for the far-infrared emission and were then used to calculate new values for the total IR luminosities for each galaxy and the size of the various dust populations. For the dust components, the derived values are quite similar to starbursts such as the Antennae and Mrk 297.
The high ISOCAM flux ratios are indicative of strong, ongoing star formation in these galaxies. The high L(PAH)/L(40-120 $`\mu `$m) and F(PAH 7.7 $`\mu `$m)/F(7.7 $`\mu `$m continuum) values also state that most of these galaxies are home to only a compact burst of star formation and that an AGN is not required within the central bursts. The exception however is NGC 6764, whose F(PAH 7.7 $`\mu `$m)/F(7.7 $`\mu `$m continuum) value of 1.22 is consistent with the presence of an AGN, yet the L(PAH)/L(40-120 $`\mu `$m) is more in line with a starburst, a finding in line with a compact AGN dominated by the starburst.
###### Acknowledgements.
The ISOCAM data presented in this paper was analyzed using ‘CIA’, a joint development by the ESA Astrophysics Division and the ISOCAM Consortium. The ISOCAM Consortium is led by the ISOCAM PI, C. Cesarsky, Direction des Sciences de la Matiere, C.E.A., France. The ISOPHOT data presented in this paper was reduced using PIA, which is a joint development by the ESA Astrophysics Division and the ISOPHOT consortium. The ISO Spectral Analysis Package (ISAP) is a joint development by the LWS and SWS Instrument Teams and Data Centers. Contributing institutes are CESR, IAS, IPAC, MPE, RAL and SRON. |
warning/0506/hep-ph0506237.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Future neutrino oscillation experiments with beams will ultimately lead to precise measurements of the leading solar and atmospheric oscillation parameters. This will also lead to precise measurements of $`\mathrm{sin}^22\theta _{13}`$, to a determination of the mass hierarchy, and measurements of leptonic CP violation. Long baseline neutrino oscillation experiments constitute therefore a very valuable physics program, since they will provide the most precise information on flavour which will be crucial for progress in our understanding of the origin and the structure of flavour. There exist different experimental paths how precision oscillation measurements can be achieved. Beyond the planned superbeam neutrino oscillation experiments T2K and NO$`\nu `$A ), one could build upgrades of superbeams (T2HK ), neutrino factories , or $`\beta `$-beams . The optimization among these alternatives depends on the outcome of on-going experiments which determine or limit, for example, the actual value of $`\mathrm{sin}^22\theta _{13}`$, which determines the CP violation and mass hierarchy measurement reaches. More important are, however, the generic differences between the different options. Since superbeam upgrades are limited by the intrinsic beam background, their performance becomes rapidly suppressed below $`\mathrm{sin}^22\theta _{13}10^3`$. Contrary to that, neutrino factories have a clean beam which contains 50% neutrinos and 50% anti-neutrinos of different flavours, which means that they are not limited by an intrinsic beam background. However, neutrino factory detectors must be able to discriminate extremely well between right-sign and wrong-sign muons to discriminate neutrinos and anti-neutrinos . This can be done with magnetized detectors, but it heavily affects the efficiencies especially at low energies, limiting the accessible L/E range. In addition, for a $`\mathrm{sin}^22\theta _{13}`$ or CP violation search experiment, high event rates turn out to be very important, which means that combinations of high neutrino energies and quite long baselines are favored (see, e.g. ). $`\beta `$-beams, on the other side, face neither of these principle disadvantages: They use a pure electron neutrino (anti-neutrino) beam and they can be easily operated at the oscillation maximum at a rather reasonable baseline. Besides there also is a proposal for a very low $`\gamma 10`$ $`\beta `$-beam which could be used to study neutrino nucleon interactions . Although the insights gained at such a facility would be invaluable to improve our theoretical understanding of neutrino cross sections, we will not discuss this option any further as it would be beyond the scope of this work.
All these novel types of experiments have certainly a number of technological issues which must be further investigated before such facilities can be built. However, the physics questions, such as the optimal $`\gamma `$ and baseline combinations for $`\beta `$-beams, should be understood first in order to identify the best options which we should be aiming for.
The answers to such questions depends obviously on the chosen detector technology, external constraints to the setup, and the physical observable for which the optimization of the $`\beta `$-beam is performed. Some of these issues have been previously discussed in detail with a somewhat different focus in .
This study addresses the most relevant of these questions. We analyze in detail the optimization of $`\beta `$-beams and we compare the physics potential of $`\beta `$-beams to the one of superbeam upgrades and neutrino factories, where a special emphasis is put on an “equal footing” comparison, i.e., we choose comparable detector sizes and running times.
This study is organized as follows: In chapter 2, we describe the simulation techniques and the parameters used for the beam and for the detectors of this study. The optimization of the $`\mathrm{sin}^22\theta _{13}`$ sensitivity is then discussed in chapter 3. Special emphasis is here put on the dependence on $`\gamma `$, the baseline and the $`\gamma `$ scaling of the number of decays per year. Chapter 4 contains then a discussion of the optimization with respect to the mass hierarchy and CP violation. A summary and conclusions will be given finally in section 5.
## 2 Experiment simulation
In this section, we first describe the properties of the $`\beta `$-beam and the techniques used for the simulation of the experiment. One of the main objectives of this study is to investigate the performance of the $`\beta `$-beam in a wide range of $`\gamma `$. Different choices of $`\gamma `$ imply different optimal detector technologies and the description of the detector parameterizations is therefore a key element of this section.
### 2.1 Simulation techniques
All experiment simulations in this study are performed with the GLoBES software . The assumed “true” solar and atmospheric oscillation parameters, which are used as input for the calculation of simulated event rates with GLoBES, are, unless stated otherwise :
$$\begin{array}{ccc}\mathrm{\Delta }m_{31}^2=2.510^3\mathrm{eV}^2,& & \mathrm{sin}^22\theta _{23}=1\\ & & \\ \mathrm{\Delta }m_{21}^2=8.210^5\mathrm{eV}^2,& & \mathrm{sin}^22\theta _{12}=0.83.\end{array}$$
(1)
The errors of these parameters are given in Refs. will also be included in our analysis. For $`\mathrm{sin}^22\theta _{13}`$, we only allow values below the CHOOZ bound , i.e., $`\mathrm{sin}^22\theta _{13}0.1`$. In some cases, we additionally demonstrate the effects for the choice of a smaller true value of $`\mathrm{\Delta }m_{31}^2=1.510^3\mathrm{eV}^2`$, representing the lower end of the currently allowed 90% CL region and, how this affects the performance of the simulated experiments.
The solar oscillation parameters and their errors are included as external input which affect the performance of the discussed setups via correlations. We assume a precision of 5% for $`\mathrm{\Delta }m_{21}^2`$ and 10% for $`\theta _{12}`$. This corresponds approximately to the current precisions such that this should be a conservative assumption at the time of the discussed experiments . In addition, $`\beta `$-beams cannot determine the leading atmospheric parameters with high precision. Therefore, we combine all of our $`\beta `$-beam simulations, which include effects due to correlations with the atmospheric parameters, with an equivalent of 10 years of T2K running<sup>1</sup><sup>1</sup>1See Refs. for details of the T2K description within GLoBES.. Data corresponding to this scenario in their combined statistics should be available from T2K, NO$`\nu `$A, and atmospheric experiments at the time when a $`\beta `$-beam is analyzed. However, we only use the T2K disappearance channels (and leave out the appearance channels) in order to avoid confusion in the interpretation of $`\beta `$-beam appearance data. This assures that we only include external information on the leading atmospheric parameters. This approach may seem to be more complicated than assuming external precisions for the leading atmospheric parameters. However, if one discusses the effects of parameter degeneracies, it is always a difficult issue where (at which point in parameter space) these external precisions for the degenerate solutions should be centered, i.e. where the external measurement gives the degenerate solution. Thus, simply assuming external precisions added at some points in parameter space would clearly over-estimate the performance of the $`\beta `$-beams by adding unwanted prior contributions. Including in the analysis an equivalent of 10 years of T2K disappearance running is therefore a simple and quite realistic approach to deal with this problem without creating “artificial” topologies. For the neutrino factory and T2HK simulations, however, it is reasonable to assume that $`\mathrm{\Delta }m_{31}^2`$ and $`\theta _{23}`$ are measured by far best with their own disappearance channels, i.e., we do not impose any external precision on the leading atmospheric parameters there. Eventually, we assume a constant matter density profile with a 5 % uncertainty on the value of the baseline-averaged matter density, where the uncertainty takes into account matter density uncertainties as well as matter density profile effects . If the baseline of an experiment is changed, we also re-compute and use the average matter density for this baseline.
### 2.2 Beam characteristics
The neutrino beam discussed in this study is assumed to originate from the beta decay of <sup>6</sup>He and <sup>18</sup>Ne isotopes in straight sections of a storage ring. It is assumed that the energy or equivalently the relativistic $`\gamma `$ value can be chosen. The corresponding decay channels are:
$$\begin{array}{ccc}{}_{}{}^{18}Ne& & {}_{}{}^{18}F+e^++\nu _e\\ & & \\ {}_{}{}^{6}He& & {}_{}{}^{6}Li+e^{}+\overline{\nu }_e\end{array}$$
(2)
which leads correspondingly to pure electron neutrino or electron anti-neutrino beams without any intrinsic beam contamination. In the rest frame of the particular decay the neutrinos are emitted isotropically and the energy spectrum is given by the well-known beta decay with a certain endpoint energy for a given isotope. In the storage ring, the spectrum is boosted and it becomes in the laboratory system :
$$\frac{d\varphi }{dE_\nu }\frac{E_\nu ^2}{\gamma }\left(1\frac{E_\nu }{2\gamma (E_0+m_e)}\right)\sqrt{\left(1\frac{E_\nu }{2\gamma (E_0+m_e)}\right)^2\left(\frac{m_e}{E_0+m_e}\right)^2}.$$
(3)
<sup>18</sup>Ne and <sup>6</sup>He are in principle not the only possible choices for the isotopes, but we will use here the same assumption as most of the existing literature . The endpoint energies $`E_0`$ of these two isotopes are very similar. This leads to the nice feature, that they give approximately the same mean neutrino energy at the same $`\gamma `$, i.e.we have $`\gamma (^6He)=\gamma (^{18}Ne)`$ and there is no obvious gain in increasing one of the two $`\gamma `$’s (cf., Ref. )<sup>2</sup><sup>2</sup>2Basically it is fortunate to have a $`\gamma `$ as high as possible for each neutrinos and anti-neutrinos independently.. The endpoint energies are $`E_0=3.4\mathrm{MeV}`$ for <sup>18</sup>Ne and $`E_0=3.5\mathrm{MeV}`$ for <sup>6</sup>He. Note that we neglect the fact that there are different exited states in the daughter nuclei of the decay, which additionally lead to negligible small contributions to the spectra with different endpoint energies.
In Ref. it is assumed that $`2.910^{18}`$ <sup>6</sup>He and $`1.110^{18}`$ <sup>18</sup>Ne decays per year can be achieved at a $`\beta `$-beam scenario with the acceleration $`\gamma (^{18}Ne)=100`$ and $`\gamma (^6Ne)=60`$ at the same time. If not stated differently, we assume that the number of decays per year stays constant with varying $`\gamma `$. However, since this assumption is most likely not justified from technological considerations, we also describe effects where the number of decays varies with $`\gamma `$ with the above values as normalization at $`\gamma (^{18}Ne)=100`$ and $`\gamma (^6Ne)=60`$.
For all discussed $`\beta `$-beam setups, we choose a total running time of 8 years. This implies that the number of ion decays above are assumed to be reached by either the simultaneous operation of the neutrino and anti-neutrino beams, or their double numbers by the successive operation of four years neutrinos and four years anti-neutrinos. We will also discuss deviations from equal neutrino and anti-neutrino running which could be achieved by successive operation option or by changing the ratio of stored ions. However, as default we use equal running times for neutrinos and anti-neutrinos, since the performance for quantities highly correlated with $`\delta _{\mathrm{CP}}`$ is usually best for equal statistics in the neutrino and anti-neutrino channels, and the higher anti-neutrino flux turns out to compensate for the lower anti-neutrino cross section very well.
### 2.3 Detector technologies
A general requirement for any $`\beta `$-beam detector is to have good muon-electron separation capabilities and to have an efficient neutral current rejection. At the same time, the technology must be available and cost effective to allow in time a scaling to large detectors. For lower values of $`\gamma `$, certainly Water Cherenkov detectors (WC) fulfill these criteria . However, at higher $`\gamma `$ values, the lack of background discrimination in WCs becomes a huge problem and other detector types, such as calorimeters or TPCs (Time Projection Chambers) are more suitable . The precise value of $`\gamma `$ where this turnover happens seems to be an unresolved question and quite different views can be found in the literature . We will describe our own approach to this problem and we will discuss our findings in relationship to the existing literature. Our choice for large values of $`\gamma `$ is the so called “Totally Active Scintillator Detector” (TASD) for reasons which will be discussed below. The two detector technologies used in this study, namely WC and TASD, are described in more detail in the following subsections.
#### 2.3.1 Water Cherenkov detector
Water Cherenkov detectors are well suited to distinguish muon neutrinos from electron neutrinos. However, background rejection can be a problem in using a WC detector in combination with a $`\beta `$-beam. The main source of background to the muon neutrino appearance search will be the flavour-blind neutral current events which are mistaken for muon neutrino charged current events. The most critical neutral current events are those where one or several energetic pions are involved, which implies that there is basically no background below the pion production threshold around $`200\mathrm{MeV}`$. Therefore, one solution would be to tune $`\gamma `$ to a low value where most of the neutrinos in the beam are below this threshold . In that case there would be no energy information, since the Fermi-motion of the nucleons would induce an energy smearing of about $`100\mathrm{MeV}`$. This would reduce the $`\beta `$-beam to a mere counting experiment, which would have only a very limited physics reach . Above the pion threshold, the feasibility of using a WC detector depends on the ability to correctly identify pions and to reject neutral current events. The pion identification works, in principle, by identifying its decay process and it seems to be possible up to some level. There are very different statements in the literature how well this can be done . The different results can to a large extent be attributed to the different level of detail used in the detector simulation. Nonetheless we will show a direct comparison of the two simulations and our parameterization at a reference scenario with $`\gamma =150`$ at a baseline of $`L=300\mathrm{km}`$ and an overall exposure equivalent to $`5000\mathrm{kt}\mathrm{y}`$, where we call the one from Ref. “case A” and the one from Ref. “case B”. As one can see from Table 1, the number of signal events is very similar in both cases, but the number of background events is very different. Moreover, the shape of the background spectra is very different and the background events are much more concentrated at low energies for case B.
The simulation of case B is based on the Super-Kamiokande Monte-Carlo and it seems to be more detailed in its treatment of detector effects. For this reason we use a parameterization which is, in total rates very close to case B, as can also be seen from Table 1. Note, however, that even though the Monte-Carlo used for case B has been well-tested in Super-Kamiokande, it is important to keep in mind that such simulations rely on physical input such as cross sections, which are not very well known. Moreover, the assumption that the response of a 20 times larger detector is the same as that of Super-Kamiokande is implied there.
In order to describe the energy response of the detector in our study, we divide the signal events into samples of quasi-elastic events (QE) and inelastic events (IE). Only for the QE sample, it is possible to accurately reconstruct the neutrino energy from the charged lepton. For IE events, the reconstructed energy will always lie below the true (incident) neutrino energy because the hadronic component of the interaction cannot be seen by a WC detector. Since the separation of those two event samples is fraud with a large error, we will use the same technique as described in Ref. . This means that the total rates number of all $`\mathrm{IE}+\mathrm{QE}`$ events is taken and in addition the spectrum of the QE event sample is used to obtain spectral information. In order to avoid double counting of events, the QE event sample is taken only with a free normalization. In this approach, no particular assumption about the event by event separation has to be made, because this approach is purely statistical. In fact, the real experiment might even perform better since there actually could be some event by event separation. For the spectral analysis, we assume that the Fermi-motion is the main component of the resulting energy resolution function. Therefore, a constant width of $`85\mathrm{MeV}`$ in a Gaußian energy resolution function is taken, in order to describe the energy reconstruction of the QE sample. For the background distribution, we make the assumption that every neutrino which interacts via neutral currents is reconstructed with an energy distribution which is flat from zero up to the true neutrino energy. In this way we obtain a background which is peaked at low energies very similar to the one in case B. We do not take into account any other background source, since in Ref. it was shown that atmospheric neutrinos only give a very small contribution.
Figure 1 shows a comparison of the physics output of the cases A, B and our parameterization. The allowed regions in the $`\theta _{13}`$-$`\delta `$ plane at a confidence level of $`2\sigma `$ for 2 d.o.f. are shown assuming that all the other parameters are exactly known and have the values as in Eq. (1).
Case B (dark gray/red curves) and our parameterization (black/blue curves) agree very well throughout the parameter space. The numbers used for our parameterization can be found in Table 2.
In order to be able to use $`\gamma `$ values in the range from 50 to 500, the energy range is chosen as $`0.23.0\mathrm{GeV}`$ and is divided into bins of $`100\mathrm{MeV}`$ width, such that the total number of bins is 28. All efficiencies are constant with exception of the first bin where they are only $`1/2`$ of the value given in Table 2 to account for threshold effects. The numbers for the disappearance channel could be quite different (in fact, the efficiencies might be much higher), but we checked that the final results do not depend on this assumption. Therefore, because of simplicity, we take the same numbers as for the appearance channels. We also include systematic uncertainties on the normalization of signal and background events as given in Table 2, where all errors are assumed to be fully uncorrelated. This is a conservative approximation and does not affect the result from the appearance measurement. The disappearance measurement, on the other hand, would require a total absolute error of less than 1% to yield any information on $`\theta _{13}`$ . This number, however, seems very difficult to be reached, if at all, for any experiment which involves neutrino-nucleus cross-sections at low energies. Our parameterization has been calibrated against case B for $`\gamma =150`$. Therefore using different values of $`\gamma `$ involves some extrapolation. Our parameterization should nevertheless be reliable from $`\gamma `$ 100 up to 350, and it should additionally reproduce the qualitative features of the $`\gamma `$-scaling within a range from 50 to 500.
#### 2.3.2 Totally active scintillator detector
For even larger values of $`\gamma `$, a detector which can also measures the hadronic energy deposition is required. The reason is that the fraction of inelastic events in the whole event sample increases with increasing $`\gamma `$, because the neutrino spectra are extended to higher energies. The techniques which have traditionally been used for that purpose are tracking calorimeters and TPCs (such as a large liquid Argon TPC as described in Ref. ). The latter technology has certainly a great potential in neutrino physics, but given the fact that background issues are not the primary concern, we will discuss the more traditional and better understood option of a tracking calorimeter. Basically, there exist three different approaches:
* magnetized iron plates, interleaved with scintillator bars
* low-Z material (such as particle board), interleaved with scintillator bars
* all active detector made of liquid scintillator and plastic tubes
The big advantage of a (magnetized) iron calorimeter is usually the ability to determine the charge of muons, but this is pointless for a $`\beta `$-beam since there is no appearance of wrong sign muons like at a neutrino factory. For the other options, the advantages and disadvantages have been very carefully addressed in the preparation of the NO$`\nu `$A proposal . We decided to use for our study the same technology as in the NO$`\nu `$A proposal, which is the so called “Totally Active Scintillator Detector” (TASD). The totally active design provides a superior energy resolution and background rejection at reasonable efficiencies. For our parameterization, we follow closely the work done by the NO$`\nu `$A collaboration, the only problem being that all studies have been done for $`\nu _e`$ appearance, whereas we look for $`\nu _\mu `$ appearance. The latter should be much easier because the muon track is much more difficult to be confused with a neutral current event. Therefore, our parameterization is on the conservative side, which does not affect our conclusions since the TASD very effectively rejects backgrounds. The numbers we use for efficiencies and systematical errors are given in Table 3 and are taken from Refs. .
The energy window reaches from $`0.5\mathrm{GeV}`$ up to the endpoint of the neutrino spectrum and is divided into 20 bins. The energy resolution is given by a Gaußian with a width of 3% $`\sqrt{E}`$ for muon neutrinos and 6% $`\sqrt{E}`$ for electron neutrinos. The background is assumed to have the same shape as the signal. But note that the shape of the background is not much of an issue in the case of a TASD detector since the background is very small. We checked that a background of the same total magnitude which is distributed like $`E_\nu ^1`$ gives basically the same results.
### 2.4 Experiment configurations and event rates
In order to compare the different detector options we still need to define the detector size and the luminosity of the beam. The fiducial volume of a WC detector seems to lie naturally of the order of $`1\mathrm{Mt}`$, since this suits also other applications as proton decay searches and there exists various proposals of this type, see e.g. Ref. . For definiteness, we assume our WC detector of this type, but with a somewhat more affordable fiducial volume of $`500\mathrm{kt}`$. For the TASD there exists currently only the NO$`\nu `$A proposal for a $`30\mathrm{kt}`$ detector, and $`50\mathrm{kt}`$’s seem to be feasible. We will use therefore the latter (larger) value as the detector mass of the TASD within this study. Within the usual uncertainties, these two detectors should also have a comparable price. These two choices lead to the typical signal and background event rates shown in Table 4 for the true parameters of Eq. (1) and $`\mathrm{sin}^22\theta _{13}=0.1`$. Note, however, that these numbers are calculated under the assumption of a constant number of isotope decays. This will of course change when we include the $`\gamma `$-scaling in subsequent sections.
## 3 Optimization of $`\mathrm{sin}^22\theta _{13}`$ sensitivity
In this section, we focus on the optimization of detecting a finite value and/or measuring $`\mathrm{sin}^22\theta _{13}`$. For that, we introduce performance indicators for the sensitivity to $`\mathrm{sin}^22\theta _{13}`$ and discuss the principle degrees of freedom for the optimization. Furthermore, we illustrate the optimization in the most relevant directions of the parameter space and choose specific setups. Eventually, we compare the performance to other established techniques, such as neutrino factories or superbeam upgrades.
### 3.1 Degrees of freedom and performance indicators
The optimization of a $`\beta `$-beam experiment involves a number of issues:
* What is the optimal $`\gamma `$? Obviously, the detector technology is a major issue in this optimization. In addition, external constraints may cause the number of decays per year not to stay constant with increasing $`\gamma `$.
* What is the optimal baseline for a given $`\gamma `$?
* How long should one run in the $`\nu `$ and $`\overline{\nu }`$ running mode? At the same or different $`\gamma `$’s? Can one run these modes simultaneously?
* What isotopes should one use? How many decays per year are realistic?
We will discuss some of these issues in greater detail below, while we will make reasonable assumptions in other cases. Let us first repeat the main assumptions which have already been discussed: We assume neutrino and anti-neutrino beams from the decay of <sup>6</sup>He and <sup>18</sup>Ne isotopes with the reference numbers for the decays per year at $`\gamma (^6He)=60`$ and $`\gamma (^{18}Ne)=100`$ as in Ref. . In general, we assume a successive operation with neutrinos and anti-neutrino running, since we will allow a variation of the neutrino versus antineutrino running time in some cases. Furthermore, we fix $`\gamma (^6He)=\gamma (^{18}Ne)`$, since there is no obvious gain in increasing one of the two $`\gamma `$’s (cf., Ref. ). We also use a constant, $`\gamma `$ independent detector mass for an assumed Water Cherenkov (WC) or a Totally Active Scintillator Detector (TASD). The WC detector will be used only below $`\gamma 500`$, since for higher energies the WC detector will be dominated by non-QE events and the background parameterization is rather unclear.
Somewhat more complicated is the issue of the $`\gamma `$-dependence of the beam. Initially we will assume that the number of decays per year does not depend on $`\gamma `$, but this is definitively not realistic. We will therefore discuss in detail the impact of a scaling with $`\gamma `$ later.
We use two different performance indicators for $`\mathrm{sin}^22\theta _{13}`$. We define the $`\mathrm{sin}^22\theta _{13}`$ sensitivity as the largest value of $`\mathrm{sin}^22\theta _{13}`$ which fits $`\mathrm{sin}^22\theta _{13}=0`$, i.e., the $`\mathrm{sin}^22\theta _{13}`$ sensitivity tests the hypothesis $`\mathrm{sin}^22\theta _{13}=0`$. It corresponds to the new exclusion limit if an experiment does not observe a signal. Since the simulated rate vector is computed for $`\mathrm{sin}^22\theta _{13}=0`$, this sensitivity does not depend on the true (simulated) values of $`\mathrm{sin}^22\theta _{13}`$, $`\delta _{\mathrm{CP}}`$, or the mass hierarchy (cf., Appendix C of Ref. ). However, there are strong correlations and degeneracies in the fit rate vector because any combination of parameters which fits $`\mathrm{sin}^22\theta _{13}=0`$ destroys the $`\mathrm{sin}^22\theta _{13}`$ sensitivity. In particular, we compute the statistics and systematics $`\mathrm{sin}^22\theta _{13}`$ sensitivity for fixed $`\delta _{\mathrm{CP}}=0`$ and the other oscillation parameters fixed to their simulated values. The correlation with $`\delta _{\mathrm{CP}}`$ and the other oscillation parameters will then be included by the projection of the fit manifold onto the $`\mathrm{sin}^22\theta _{13}`$-axis as the correlation bar of our figures, where only the best-fit solution is used. Any disconnected solution at the chosen confidence level which fits $`\mathrm{sin}^22\theta _{13}=0`$ is treated as degeneracy, such as a $`(\delta _{\mathrm{CP}},\theta _{13})`$ or mass hierarchy degeneracy. Thus, we treat connected degenerate solutions (with the best-fit manifold) as correlations, and disconnected degenerate solutions as degeneracies.
In addition to the $`\mathrm{sin}^22\theta _{13}`$ sensitivity, we show the $`\mathrm{sin}^22\theta _{13}`$ discovery reach in some cases, which tests the hypothesis of observing a signal for a given set of true values. Thus, the $`\mathrm{sin}^22\theta _{13}`$ discovery reach strongly depends on the true values of $`\mathrm{sin}^22\theta _{13}`$, $`\delta _{\mathrm{CP}}`$, and the mass hierarchy. However, there are almost no correlations or degeneracies, since the fit rate vector is computed for $`\mathrm{sin}^22\theta _{13}=0`$. In principle, the risk-minimized (with respect to the possible true values) $`\mathrm{sin}^22\theta _{13}`$ discovery reach is comparable to the $`\mathrm{sin}^22\theta _{13}`$ sensitivity – though the problem is not exactly symmetric.
### 3.2 Performance as function of $`\gamma `$ and baseline optimization
A first important optimization question concerns the optimal value of $`\gamma `$. Naively this question seems to be trivial for a fixed number of decays per year - the higher the $`\gamma `$, the better. However, there is a strong dependence on detector technology, since non-QE events will start to dominate a WC detector for higher energies, and the efficiency of a TASD (or iron calorimeter) is very low at low energies because of too short tracks compared to the positional information of the detector. We show therefore in Figure 2 the $`\mathrm{sin}^22\theta _{13}`$ sensitivity for the WC detector (left) and the TASD (right) as function of $`\gamma `$, where the parameterization of the WC detector is most reliable in the unshaded region. For the $`L/\gamma `$, we assume 1.3 for both detectors. Note, however, that we will find that $`L/\gamma =1.3`$ is not optimal for the Water Cherenkov detector if one includes correlations and degeneracies. For the TASD, the rule “the higher the $`\gamma `$, the better” obviously applies if the stored ion decays are constant in $`\gamma `$, since the muon tracks are in average fully contained in the shown range. For the WC detector, this rule also applies in the unshaded region, but it is likely that background domination and non-QE events will affect the performance for higher $`\gamma `$. In addition, the dependence on $`\gamma `$ is much more shallow in the considered range, and the impact of systematics increases with increasing $`\gamma `$. Note that though the Water Cherenkov detector has the better systematics only performance at the upper end of the unshaded region (left panel), both detectors perform very similar after including correlations and degeneracies in this range.
The baseline dependence of the $`\mathrm{sin}^22\theta _{13}`$ sensitivity on $`L/\gamma `$ is shown in Figure 3 for different choices of (fixed) $`\gamma `$ as given in the plot labels. The optimal $`\mathrm{sin}^22\theta _{13}`$ sensitivity can in all cases be achieved from a statistical and systematical point of view for $`L/\gamma 0.81.3`$, where the appearance events only come from the first oscillation maximum. For higher values of $`L/\gamma `$ the statistics $`\mathrm{sin}^22\theta _{13}`$ sensitivity decreases due to the $`1/L^2`$ dependence of the flux although the actual position of the first oscillation maximum, taking into account the average neutrino energy, would be at $`L/\gamma 2.1`$. However, longer baselines, where appearance events from the second oscillation maximum enter the energy window of the analysis, have altogether a better potential to resolve correlations and degeneracies. The combined effect of the first and second oscillation maximum together leads to a better determination of the oscillation pattern and larger matter effects, though the statistic limit becomes worse. For the WC detector, we therefore choose $`L/\gamma =2.6`$, where the impact of correlations and degeneracies is marginal and the overall performance is significantly improved. We demonstrate in Appendix A in detail where this better performance comes from and why it is only visible after the inclusion of correlations and degeneracies. Furthermore, we choose $`L/\gamma =1.3`$ for the TASD, since a larger $`L/\gamma `$ hardly improves the performance. In addition, we have tested that if the number of decays per year significantly drops with increasing $`\gamma `$, the minimum of Setup 3 is shifted towards longer baselines. Note that in all cases the flatness in $`L/\gamma `$ implies that the precise baseline is not so important for the overall $`\mathrm{sin}^22\theta _{13}`$ sensitivity. In addition, the choice of these values for $`L/\gamma `$ cannot be entirely based upon this figure and will be later justified in the context of different performance indicators.
Let us now directly compare the two detector technologies in Figure 4, where the gray curves correspond to the WC detector and the black curves to the TASD. For small $`\gamma `$, the WC detector obviously has the best performance, whereas for higher $`\gamma `$, the TASD is the way to go for. The crossing point between the two technologies depends on the confidence level and lies somewhere in the interval $`250\gamma 500`$. Note, however, that the TASD performance is dominated by correlations and degeneracies (cf., Figure 2), which means that its discovery potential will certainly be better for $`\gamma =500`$ at all confidence levels. In addition, the parameterization of the WC detector is not very reliable in this region anymore.
In order to discuss some effects in greater detail we use well defined setups/representatives in order to evaluate the requirements for a $`\beta `$-beam. Comparing with Figure 4, there are three interesting (approximate) ranges:
1. Low $`\gamma 300`$: This range can be probed with relatively “small” accelerators, such as of SPS size, and WC detectors. The physics potential could compete with superbeam upgrades, such as an upgrade of T2K to a Hyperkamiokande detector, or an intermediate step towards a neutrino factory.
2. Medium $`300\gamma 800`$: In this case, larger accelerator rings are required, for example of Tevatron size. The detector technology could be WC or TASD. The physics goal could be to compete with superbeam upgrades or even small neutrino factories.
3. Large $`\gamma 800`$: Very large accelerators of the size of the LHC are required<sup>3</sup><sup>3</sup>3Note, however, that even though the LHC would have enough energy, it certainly does not have enough RF acceleration power, i.e., a new huge accelerator would be required in this case.. TASD or iron calorimeters are possible choices for the detector technologies. In this case, the goal is clearly competition with neutrino factories for all of the relevant measurements.
These $`\gamma `$-ranges are shown in many figures of this and the next chapter. Obviously, the requirements for all of these ranges are somewhat different, which means that it will not only be sufficient to compare $`\beta `$-beam with $`\beta `$-beam, but it will also be necessary to compare individual $`\beta `$-beams from each range with its competitors. For this purpose, we define three setups from these three ranges (cf., Figure 4):
* Setup 1 ($`\gamma =200`$, WC): In this case, the WC detector representation is well-established and $`\gamma `$ should not be too high for SPL-sized accelerators.
* Setup 2 ($`\gamma =500`$, TASD): This setup represents the lowest $`\gamma `$ where a TASD will likely perform better that a WC detector at all confidence levels.
* Setup 3 ($`\gamma =1000`$, TASD): This representative corresponds to a very sophisticated option close the upper limit of what is doable.
### 3.3 Isotope decay scaling
We have already mentioned that the number of decays per year is most likely not constant in $`\gamma `$. In order to include this effect, we use the following power law parameterization, which should be justified for a certain $`\gamma `$-range, to describe this scaling with $`\gamma `$ ($`i=1`$ for $`{}_{}{}^{18}Ne`$: neutrinos, $`i=2`$ for $`{}_{}{}^{6}He`$: anti-neutrinos):
$$N^i=N_0^i\left(\frac{\gamma _0^i}{\gamma }\right)^n$$
(4)
Here $`N_0^i`$ is determined by our reference point at $`\gamma _0^1=100`$, $`\gamma _0^2=60`$. We can now discuss different cases for $`n`$, which leads to different optimization strategies:
* $`n=0`$: The number of decays per year is fixed. This implies that the accelerator and storage ring has to scale appropriately with $`\gamma `$ in a non-trivial manner.
* $`0<n<1`$: This seems to be the most likely range of realistic cases. The number of decays per year becomes constrained with increasing $`\gamma `$ by the geometry of the accelerator and decay ring and $`\gamma `$ increased lifetime of the isotopes in the laboratory system.
* $`n1`$: This case corresponds to a fixed setup constraining the performance. Given the scaling of the baseline, the number of events stays approximately constant as function of $`\gamma `$. A realistic constraint for the SPS would be, for example, $`n1`$ from the number of merges in the decay ring and the number of ions per bunch .
* $`n>1`$: In this case, it clearly does not make sense to go to higher $`\gamma `$’s, since the event rate decreases with $`\gamma `$ if we stay in the oscillation maximum
* $`n<0`$: The number of decays per year increases with $`\gamma `$. This hypothetical (but technologically unlikely) possibility requires that the accelerator and decay ring over-proportionally scale with $`\gamma `$.
We further on consider the range $`0n1`$ to be realistic. However, it is conceivable that, for a given setup, the performance will scale with $`n0`$ in the beginning, and change into an $`n1`$ scaling in the saturation regime.
We show in Figure 5 the $`\gamma `$-scaling of the final $`\mathrm{sin}^22\theta _{13}`$ sensitivity (including correlations and degeneracies) for the different detector technologies and $`0n1`$ (bands), where the solid lines correspond to $`n=0.5`$. For $`n1`$ (upper ends of bands), the performance is rather flat in $`\gamma `$. This simply means that the increase in cross section is compensated by the $`\gamma `$-scaling. For $`n0`$ (lower end of bands), the $`\mathrm{sin}^22\theta _{13}`$ sensitivity scales as discussed above. It is interesting to observe that Setups 1 and 2 are much less affected by a decreasing number of decays than Setup 3, which over-proportionally looses sensitivity for $`n>0`$. Thus, for Setup 3, it is crucial that the accelerator and storage design allow enough decays per year, whereas for Setups 1 and 2 a certain loss in the number of decays results only in about a factor of two weaker sensitivity limit. This means that, given a specific accelerator (such as the LHC), it makes only sense to discuss very high $`\gamma `$ setups if it is guaranteed that enough ions can be stored. Note that the dependence on $`n`$ is therefore an important constraints for the $`\gamma `$ optimization of $`\beta `$-beams, since, for the TASD, it constrains the rule “the larger $`\gamma `$, the better”.
### 3.4 Comparison with other technologies
Next we compare the $`\mathrm{sin}^22\theta _{13}`$ sensitivity of $`\beta `$-beams to neutrino factories and superbeam upgrades. For the neutrino factory, we use the representative NuFact-II from Ref. which uses only $`\nu _\mu `$ appearance and disappearance in two different baseline configurations $`L=\mathrm{3\hspace{0.17em}000}\mathrm{km}`$ and $`L=\mathrm{7\hspace{0.17em}500}\mathrm{km}`$ (NF@3000km and NF@7500km). This assumes $`4\mathrm{MW}`$ target power (corresponding to about $`1.0610^{21}`$ useful muon decays per year), a $`50\mathrm{kt}`$ fiducial volume magnetized iron calorimeter, and $`50\mathrm{GeV}`$ neutrinos. The longer baseline corresponds to the “magic baseline”, where correlations and degeneracies are resolved, but no $`\delta _{\mathrm{CP}}`$ measurement is possible . The $`\mathrm{sin}^22\theta _{13}`$ sensitivity and mass hierarchy sensitivity of the neutrino factory at the magic baseline is therefore very competitive. We include this option because we want to compare optimized setups (for particular purposes) with optimized setups later, i.e., the optimal $`\beta `$-beam for the mass hierarchy measurement may also have a different baseline that the one for the CP violation measurement. For the superbeam upgrade, we choose T2HK from Refs. simulated in Ref. , but we reduce the fiducial mass to $`500\mathrm{kt}`$ and use the same detector as in this study in order to be comparable to our WC detector for Setup 1. This superbeam upgrade also assumes a target power of $`4\mathrm{MW}`$ and we call it therefore T2HK. The different setups are summarized in Table 5.
Figure 6 shows the $`\mathrm{sin}^22\theta _{13}`$ sensitivity for these setups. First, it is interesting to observe that all of the $`\beta `$-beam representatives are competitive to the T2HK setup in all cases. The neutrino factory at $`L=\mathrm{3\hspace{0.17em}000}\mathrm{km}`$ can, for the chosen parameter values, not resolve the $`(\delta ,\theta _{13})`$-degeneracy at the $`3\sigma `$ confidence level<sup>4</sup><sup>4</sup>4This behaviour could in principle change once additional information like $`\nu _\tau `$ appearance is available, and therefore the final $`\mathrm{sin}^22\theta _{13}`$ sensitivity is much worse than that of the $`\beta `$-beams. The neutrino factory at $`L=\mathrm{7\hspace{0.17em}500}\mathrm{km}`$ has the best overall performance because it is hardly affected by the correlation with $`\delta _{\mathrm{CP}}`$. Note that Setup 1 has a better final $`\mathrm{sin}^22\theta _{13}`$-sensitivity than Setup 2 because of the choice of $`L/\gamma =2.6`$. Though the statistics limit is worse for this choice, the correlation and degeneracy bars become very small and lead to a better final sensitivity than for $`L/\gamma =1.3`$.
In order to complete the picture, we show in Figure 7 the $`\mathrm{sin}^22\theta _{13}`$ discovery reach as function of the true values of $`\mathrm{sin}^22\theta _{13}`$, $`\delta _{\mathrm{CP}}`$ (stacked to the “CP fraction”), and the mass hierarchy. Though the $`\mathrm{sin}^22\theta _{13}`$ sensitivity of Setup 1 is slightly better than one of Setup 2, one can clearly see that there is a hierarchy in the discovery reach: The choice of Setups 1, 2 and 3 implies that their differences correspond to approximately equal improvements in terms of fraction of the parameter space. The large impact of correlations on the $`\mathrm{sin}^22\theta _{13}`$ sensitivity of Setup 2 (cf., Figure 6), which mainly comes from the correlation with $`\delta _{\mathrm{CP}}`$, implies that the discovery reach is in many cases of $`\delta _{\mathrm{CP}}`$ better than the one of Setup 1. Note that the relative position of the Setup 1 curve would not significantly change with a different choice of $`L/\gamma =1.3`$ because in this case the shape of the curves would be very similar, but the systematics limit of Setup 2 is slighlty better (cf., Figure 2). For a normal hierarchy, the discovery reach of Setup 3 covers considerably less parameter space than that one of a neutrino factory at $`L=\mathrm{3\hspace{0.17em}000}\mathrm{km}`$. For the inverted hierarchy, the matter effect enhancement of the lower neutrino factory anti-neutrino rate (instead of the higher neutrino rate) leads to a relatively degraded reach for the neutrino factory baselines, whereas the event numbers of the $`\beta `$-beams are rather similar for neutrinos and anti-neutrinos. However, the neutrino factory at $`\mathrm{3\hspace{0.17em}000}\mathrm{km}`$ covers the most parameter space in both cases, and the superbeam upgrade T2HK by far the least. Note that the neutrino factory behavior for the $`\mathrm{sin}^22\theta _{13}`$ sensitivity and $`\mathrm{sin}^22\theta _{13}`$ discovery reach is completely different. For the $`\mathrm{sin}^22\theta _{13}`$ sensitivity, very few particular combinations of parameters prevent a strong $`\mathrm{sin}^22\theta _{13}`$ limit, whereas for the $`\mathrm{sin}^22\theta _{13}`$ discovery reach, correlations and degeneracies are of secondary importance. A neutrino factory is therefore clearly a $`\mathrm{sin}^22\theta _{13}`$ discovery machine. Note that the discovery reach fit rate vector is computed for (fixed) $`\mathrm{sin}^22\theta _{13}=0`$, which means there is only a substantial impact of correlations if the solar appearance term contributes significantly to the appearance rate.<sup>5</sup><sup>5</sup>5It turns out that these correlations have some impact on the discovery reach of Setup 1 and the neutrino factory baselines, especially $`L=\mathrm{3\hspace{0.17em}000}\mathrm{km}`$, since both, having a long enough baseline and being far off the matter density resonance (in energy) increase the relative importance of the solar appearance term. Since the solar appearance term does not depend on the mass hierarchy, there is no contribution of the $`\mathrm{sgn}(\mathrm{\Delta }m_{31}^2)`$-degeneracy.
In summary, all of the discussed $`\beta `$-beam options could be interesting alternatives to a superbeam upgrade or intermediate options towards a neutrino factory. Especially Setup 1, which uses the identical detector needed for other applications, might be an interesting transition candidate. Since the $`\beta `$-beams have a clean composition of electron neutrinos, they are not affected by an intrinsic contamination of muon events and are therefore not systematics limited close to $`\mathrm{sin}^22\theta _{13}10^3`$ such as superbeams are. We will discuss the sensitivity to the mass hierarchy and CP violation in the next section in order to get a broader perspective of the problem, and to evaluate if a $`\beta `$-beam could, in principle, replace a neutrino factory.
## 4 Sensitivity to mass hierarchy and CP violation
Beyond detecting a finite value of $`\mathrm{sin}^22\theta _{13}`$, the mass hierarchy and CP violation sensitivities are the two most interesting quantities to be measured by the discussed experiments. We will first introduce performance indicators for these quantities. Then we will discuss optimization aspects for the mass hierarchy and CP violation determination. Finally, we will compare the performances of the $`\beta `$-beam setups with other experiments.
### 4.1 Performance indicators
We define sensitivity to a particular mass hierarchy (normal or inverted) if the wrong-sign solution can be excluded at a certain confidence level. This implies that a $`\mathrm{sgn}(\mathrm{\Delta }m_{31}^2)`$-degeneracy or mixed degeneracy will destroy the mass hierarchy sensitivity. The mass hierarchy sensitivity does not only depend on the simulated hierarchy, but also on the true values of $`\mathrm{sin}^22\theta _{13}`$ and $`\delta _{\mathrm{CP}}`$. Since it is not possible to show the full parameter space for mass hierarchy sensitivity simultaneously, we use in many cases the sensitivity for the true $`\delta _{\mathrm{CP}}=0`$. In other cases, we show the mass hierarchy sensitivity as function of $`\mathrm{sin}^22\theta _{13}`$ and $`\delta _{\mathrm{CP}}`$, where we stack the true values of $`\delta _{\mathrm{CP}}`$ to the “Fraction of $`\delta _{\mathrm{CP}}`$” (CP fraction) similar to Figure 7. A mass hierarchy sensitivity for a CP fraction $`1`$ corresponds to mass hierarchy sensitivity for any values of $`\delta _{\mathrm{CP}}`$ (worst case in $`\delta _{\mathrm{CP}}`$), and a mass hierarchy sensitivity for a CP fraction $`0`$ to mass hierarchy sensitivity for the best case $`\delta _{\mathrm{CP}}`$. A CP fraction $`0.5`$ refers to the “typical” value of $`\delta _{\mathrm{CP}}`$, i.e., 50% of all cases of $`\delta _{\mathrm{CP}}`$. Note that the CP fraction is a probabilistic measure in the sense that only one of these values can be realized by nature. Assuming a uniform distribution of $`\delta _{\mathrm{CP}}`$, it directly corresponds to the probability to discover the mass hierarchy for a chosen $`\mathrm{sin}^22\theta _{13}`$.
We define sensitivity to CP violation if the CP conserving solutions $`\delta _{\mathrm{CP}}=0`$ and $`\delta _{\mathrm{CP}}=\pi `$ can be excluded at a certain confidence level. This implies that any degenerate solution overlapping a CP conserving solution destroys the sensitivity to CP violation. Note that this sensitivity clearly differs from the parameter sensitivity to a specific parameter value of $`\delta _{\mathrm{CP}}`$, which includes the sensitivity to the special value $`\delta _{\mathrm{CP}}=0`$. In some cases, we only show the parameter sensitivity to maximal CP violation $`\delta _{\mathrm{CP}}=\pi /2`$ or $`\delta _{\mathrm{CP}}=3\pi /2`$. However, since any value of $`\delta _{\mathrm{CP}}\{0,\pi \}`$ violates CP, we also show the parameter sensitivity to any CP violation in other cases, i.e., the sensitivity to CP violation as function of the true values of $`\mathrm{sin}^22\theta _{13}`$ and $`\delta _{\mathrm{CP}}`$ (which, in principle, also depends on the simulated mass hierarchy). Similar to the mass hierarchy sensitivity, we then stack the values of $`\delta _{\mathrm{CP}}`$ to the “Fraction of $`\delta _{\mathrm{CP}}`$”. Note, however, that no experiment can have a CP fraction $`1.0`$ for the CP violation sensitivity at any point, since there will be no CP violation sensitivity for values of $`\delta _{\mathrm{CP}}`$ close to $`0`$ and $`\pi `$.
### 4.2 Scaling with $`\gamma `$ and optimization
Similar to the $`\mathrm{sin}^22\theta _{13}`$ sensitivity, one can discuss the mass hierarchy and CP violation sensitivities as function of $`\gamma `$ for the different detector technologies. This comparison is shown in Figure 8 for the chosen range of $`L/\gamma `$. Since higher $`\gamma `$ implies a longer baseline, it also implies stronger matter effects, where we use the average density along a specific baseline. Therefore the mass hierarchy sensitivity also improves with higher values of $`\gamma `$ (cf., left plot). The different choice of $`L/\gamma `$ for the WC detector implies that the mass hierarchy sensitivity is already present at about half of the $`\gamma `$ for the TASD. There is no substantial difference between the normal and inverted mass hierarchies, because all setups use approximately equal neutrino and anti-neutrino rates. For the CP violation sensitivity, higher $`\gamma `$’s are, in principle, favorable, since they imply larger event rates. However, for very high $`\gamma `$, the missing energy resolution of the non-QE events (WC) and the matter density uncertainties (TASD) act counter-productive. For the $`\beta `$-beams, there seem to exist only very little problems with degeneracies for the CP violation sensitivity, because the measurement at the oscillation maximum helps to resolve the degeneracies.
It is interesting to discuss what the choices of $`L/\gamma `$ for Setups 1, 2, and 3 are really optimized for. So far, we have demonstrated that the choice of $`L/\gamma =2.6`$ for the WC detector and $`L/\gamma =1.3`$ for the TASD are quasi-optimal for the $`\mathrm{sin}^22\theta _{13}`$ sensitivity and lead to a clear hierarchy of these setups for the $`\mathrm{sin}^22\theta _{13}`$ discovery reach. We show in Figure 9 the $`L/\gamma `$-dependence for all three setups and the $`\mathrm{sin}^22\theta _{13}`$ sensitivity (black solid curves), maximal CP violation sensitivity (dashed curves), and normal mass hierarchy sensitivity for $`\delta _{\mathrm{CP}}=0`$ (gray curves). The thick vertical lines correspond to our choices of $`L/\gamma `$, whereas the thin lines represent alternative optimization strategies. Setup 1 is apparently optimized for the $`\mathrm{sin}^22\theta _{13}`$ sensitivity, where the large $`L/\gamma `$ clearly favors the mass hierarchy sensitivity and hardly affects the CP violation sensitivity. Thus, it represents a good compromise for all quantities. Setups 2 and 3 are optimized for CP violation, where the $`\mathrm{sin}^22\theta _{13}`$ sensitivity is in both cases very close to the optimum. For Setups 2 and 3, however, the mass hierarchy sensitivity would be considerably better close to the second oscillation maximum $`L/\gamma 2.6`$, while the $`\mathrm{sin}^22\theta _{13}`$ sensitivity would be hardly degraded. The choice of the baseline depends therefore for Setups 2 and 3 on the priorities, i.e. if one optimizes for mass hierarchy or CP violation measurements. Since we also use the setup NF@7500km for the neutrino factory, we will show Setups 2 and 3 at the second oscillation maximum in some cases for a fair comparison of the mass hierarchy sensitivity. Note that Setup 3 could, in principle, also be operated at the “magic baseline”, where the mass hierarchy and $`\mathrm{sin}^22\theta _{13}`$ sensitivities are only somewhat worse than optimal, but no CP violation measurement is possible. The only quantity which is not shown here is the $`\mathrm{sin}^22\theta _{13}`$ discovery reach. One can show that it is substantially degraded for the second oscillation maximum at Setups 2 or 3 ($`L=2.6\gamma `$). In particular, Setup 2 would perform much worse than Setup 1. However, the choice of the first oscillation maximum instead of the second would hardly change the parameter space coverage of Setup 2. Thus, our choices of $`L/\gamma `$ are consistent with the primary objective to discover $`\mathrm{sin}^22\theta _{13}`$.
Recently the important issue has been raised that the neutrino event rates might be substantially suppressed compared to the anti-neutrino event rates. We show therefore in Figure 10 the dependence on the neutrino running fraction for all three setups and the $`\mathrm{sin}^22\theta _{13}`$ sensitivity (black solid curves), maximal CP violation sensitivity (dashed curves), and normal mass hierarchy sensitivity for $`\delta _{\mathrm{CP}}=0`$ (gray curves). The vertical thick lines correspond to our choice of 50% neutrino and 50% anti-neutrino running. The neutrino running fraction $`f`$ is the fraction of neutrino running divided by the total running time of eight years, i.e. the experiment runs $`f\times 8`$ years with neutrinos and $`(1f)\times 8`$ years with anti-neutrinos. From Figure 10, we find that all setups are at the optimal performance for the CP violation measurements (vertical lines). Setups 1 and 2 also have optimal $`\mathrm{sin}^22\theta _{13}`$ sensitivity, whereas Setup 3 has optimal mass hierarchy sensitivity. As far as the symmetry of the plots is concerned, running with only anti-neutrinos is clearly favored compared to running with only neutrinos (extreme cases), because we have somewhat higher anti-neutrino event rates and lower backgrounds, i.e. the absolute rate is higher for the anti-neutrino case. For the inverted hierarchy (Figure 10 is shown for the normal hierarchy), only running with neutrinos is even slightly more disfavored because of the matter suppression of the neutrino rate. It is interesting to note that even rather substantial deviations from a symmetric neutrino and anti-neutrino operation does not have extreme effects on the measurements. Setup 3 is most affected by such deviations, where a lower neutrino fraction means better statistics and thus a better $`\mathrm{sin}^22\theta _{13}`$ sensitivity, but it creates an imbalance between neutrinos and anti-neutrinos affecting the CP violation sensitivity. Nevertheless, it does not make sense to run with neutrinos or anti-neutrinos only, since this ratio would lead to degrading the sensitivities by an order of magnitude. Setup 1, for example, then looses its competitiveness compared to superbeam upgrades. In all cases, at least 10%-20% of neutrino running is necessary, which corresponds rescaled to at least a total number of $`110^{18}`$ useful <sup>18</sup>Ne decays plus $`2610^{18}`$ useful <sup>6</sup>He decays.
### 4.3 Comparison with other technologies
In Figure 11 the sensitivity to CP violation is shown for the normal (left) and inverted (right) mass hierarchy for different experiments as function of the true values of $`\mathrm{sin}^22\theta _{13}`$ and $`\delta _{\mathrm{CP}}`$ at the $`3\sigma `$ confidence level. This figure clearly demonstrates that for a normal mass hierarchy all of the discussed $`\beta `$-beam options have an impressing CP violation sensitivity very competitive to the neutrino factory, because $`(\delta _{\mathrm{CP}},\mathrm{sin}^22\theta _{13})`$-degeneracy and “$`\pi `$-transit” of the $`\mathrm{sgn}(\mathrm{\Delta }m_{31}^2)`$-degeneracy destroy the CP violation sensitivity of the neutrino factory at many places. Note that these degeneracy problems could, in principle, be resolved by a combination with the magic baseline, but a much better sensitivity than that of Setup 3 is unlikely to be achieved. As far as the $`\mathrm{sin}^22\theta _{13}`$ reach is concerned (in horizontal direction), there is again a clear hierarchy among Setups 1, 2, and 3. For large values of $`\mathrm{sin}^22\theta _{13}`$, however, matter density uncertainties affect the longer baselines, and Setup 2 has to deal with some problems due to degeneracies (left plot). For optimal $`\mathrm{sin}^22\theta _{13}`$, Setup 3 can establish CP violation for more than 90% of all values of $`\delta _{\mathrm{CP}}`$, whereas the neutrino factory is limited to about 80%. For the inverted hierarchy, the matter effects enhance the anti-neutrino rate, which means that the neutrino and (lower) anti-neutrino rates at the neutrino factory are getting more equal statistical weight and the correlations can easier be resolved. Balanced event rates of the $`\beta `$-beams lead, on the other hand, to very little impact of the mass hierarchy. For T2HK, the somewhat lower anti-neutrino rate implies a similar behavior to the neutrino factory.
In order to compare the mass hierarchy sensitivity of all options, we show in Figure 12 the mass hierarchy sensitivity as function of $`\mathrm{sin}^22\theta _{13}`$ and $`\delta _{\mathrm{CP}}`$ for the normal (left) and inverted (right) mass hierarchy for the Setups defined in the last section and the neutrino factory and superbeam representatives. Given the choice of $`L/\gamma `$, Setups 1 and 2 have a very similar mass hierarchy sensitivity because of the very similar baselines, whereas Setup 3 is substantially better. In all cases, the neutrino factory at the magic baseline covers by far the most parameter space, whereas the performance of NF@3000km is very close to the one of Setup 3. The superbeam setup T2HK can only establish the mass hierarchy for a very small fraction of $`\delta _{\mathrm{CP}}`$ because of its short baseline. Note, however, that other superbeam upgrades (such as FNAL-Homestake or BNL-Homestake) could have a much better mass hierarchy sensitivity . The relative performance of the neutrino factory baselines is degraded for the inverted versus normal hierarchy, because the event rates are not evenly distributed between neutrinos and anti-neutrinos. As we have discussed in the last section, a longer baseline for Setups 2 and 3 improves the mass hierarchy sensitivity drastically. It turns out that for $`L/\gamma =2.6`$, Setup 2 is comparable to the neutrino factory at 3000 km, whereas Setup 3 is almost as good as the neutrino factory at the magic baseline.
A very interesting (though not very likely) case for the $`\beta `$-beams could occur if $`\mathrm{\Delta }m_{31}^2`$ turns out to be at the lower end of the currently allowed 90% CL region , i.e., $`\mathrm{\Delta }m_{31}^20.0015\mathrm{eV}^2`$. Since $`\mathrm{\Delta }m_{31}^2`$ will be measured to a high precision soon by MINOS, T2K, and NO$`\nu `$A, it is straightforward to re-optimize the $`\beta `$-beams by just moving the detector back into the oscillation maximum. For the neutrino factory, however, the oscillation maximum for the mean energy is at a very long baselines $`L\mathrm{7\hspace{0.17em}500}\mathrm{km}`$ anyway, which means that moving the detector to an even longer baseline should have an effect similar to choosing the magic baseline scenario directly. In addition, other constraints may prevent the selection of longer baselines. We illustrate the effect of a smaller $`\mathrm{\Delta }m_{31}^2`$ in Figure 13, where the arrows indicate the shift. In this figure, the $`\beta `$-beam baselines are re-scaled according to $`LL\times 0.0025/0.0015`$ in order to stay in the oscillation maximum, whereas the other baselines are fixed. In almost all cases the experiments loose sensitivity. However, the relative shift for the neutrino factories is in some cases much larger because the smaller $`\mathrm{\Delta }m_{31}^2`$ means that the oscillation peak is shifted to lower energies where the charge identification requirement leads to lower efficiencies. In particular, the CP violation sensitivity of the neutrino factory is highly affected. For the mass hierarchy sensitivity, the neutrino factory at the magic baseline is still the best experiment. Note that Setups 1 and 2 are hardly affected by the different value of $`\mathrm{\Delta }m_{31}^2`$, since a smaller value of $`\mathrm{\Delta }m_{31}^2`$ means a longer baseline and the stronger matter effects partially (Setup 2) or fully (Setup 1) compensate for the drop in $`1/L^2`$. To be fair, this comparison is only shown for selected parameter choices and for the assumption of unflexible neutrino factory baselines. Indeed, for a smaller value of $`\mathrm{\Delta }m_{31}^2`$, a dedicated study is required which re-optimizes all potential experiments. However, one can easily see from this figure that the charge identification cuts at low energies for the neutrino factory imply that one quickly ends up at inconveniently long baselines for such an optimization. This behavior is not expected for the $`\beta `$-beams.
## 5 Summary and conclusions
In this study, we have discussed various optimization aspects of $`\beta `$-beams and we compared the physics potential to neutrino factories and superbeam upgrades. Two central parameters are the gamma factor and the baseline of the $`\beta `$-beam which were taken to be free parameters. We considered two different detector technologies in connection with the $`\beta `$-beam, namely a Water Cherenkov detector and a Totally Active Scintillator Detector with $`500\mathrm{kt}`$ and $`50\mathrm{kt}`$ fiducial mass, respectively. The Water Cherenkov detector was also considered as a target for a superbeam upgrade which we called T2HK. For the comparison with the neutrino factory we used a $`50\mathrm{kt}`$ magnetized iron detector at baselines of 3000 km and 7500 km. An important aspect concerns the assumptions about the number of ion decays per year. One scenario was to assume that the number of decays does not depend on $`\gamma `$. However, this is for a number of reasons technologically not realistic and we studied therefore scenarios where the number of decays per year scales with $`\gamma `$ like a power law. For the superbeam upgrade and the neutrino factory “standard” beam luminosities were assumed (see section 3.4 for details). As performance indicators, we have considered the $`\mathrm{sin}^22\theta _{13}`$ sensitivity (corresponding to the exclusion limit which can be achieved by an experiment), the $`\mathrm{sin}^22\theta _{13}`$ discovery reach, the (normal and inverted) mass hierarchy discovery reach, and the CP violation discovery reach for both hierarchies. Specific experimental setups were defined to allow a more detailed comparison.
Our main results for the discovery reaches are summarized in Figure 14, where the bands reflect the impact of the true value of $`\delta _{\mathrm{CP}}`$. We find that the choice of the optimal $`\gamma `$ clearly depends on the objectives of the $`\beta `$-beam experiment and external constraints. We in general find good agreement with existing studies showing that $`\gamma `$ should be at least high enough to avoid the Fermi-motion dominated regime in order to have sufficient energy information . Low $`\gamma 300`$ can be achieved with relatively “small” accelerators, such as of SPS size in combination with Water Cherenkov detectors. All of the discussed performance indicators imply that a $`\beta `$-beam in this range clearly outperforms the T2HK superbeam upgrade using the same detector since it is not limited by the intrinsic beam background. In fact, the CP violation discovery reach is already quite close to the optimum even compared with higher gamma or neutrino factory setups. Note that this excellent simultaneous sensitivity to $`\mathrm{sin}^22\theta _{13}`$, mass hierarchy, and CP violation is achieved by including the second oscillation maximum to disentangle correlations and degeneracies. An operation at a shorter baseline would significantly affect the final mass hierarchy and $`\mathrm{sin}^22\theta _{13}`$ sensitivities, and it would hardly help the CP violation sensitivity. Note that we have chosen the ion decay rates such that there are approximately equal neutrino and anti-neutrino event rates for all setups. This balanced concept implies that the $`\beta `$-beam performance hardly depends on the mass hierarchy and the performance for CP violation is excellent. However, if the neutrino rate were significantly lower than the anti-neutrino rate, then the CP violation sensitivity would be affected.
The range $`300\gamma 800`$, requires already relatively large accelerator rings, such as of Tevatron size. As detector technology, we have chosen the Totally Active Scintillator Detector (TASD) in this case, because it is rather predictable in this gamma-range. However, it is not excluded that the extrapolation to a ten times larger Water Cherenkov could result in a better performance. For the setup in this and the larger gamma range, we find that the choice of $`L/\gamma =1.3`$ or $`2.6`$ clearly depends on the optimization goals. For $`\mathrm{sin}^22\theta _{13}`$ discovery and CP violation potential, shorter baselines are favored, whereas sensitivity to the mass hierarchy favours longer baselines. The $`\mathrm{sin}^22\theta _{13}`$ is rather indifferent with respect to the baseline choice, since statistics becomes worse for longer baselines, while the matter effects become stronger and the second oscillation maximum helps to resolve correlations and degeneracies. Except for CP violation, even the optimized medium $`\gamma `$ setup is not competitive to the neutrino factory, but it could be an interesting step towards it.
Large $`\gamma 800`$ would require very large and powerful accelerators of LHC size, where we use TASD as detector technology. We find $`\mathrm{sin}^22\theta _{13}`$ (sensitivity and discovery) and mass hierarchy discovery reaches close to the neutrino factory setups, and a CP violation discovery reach better than the one of the neutrino factory setups (cf., Figure 6 and Figure 14). In this case, a $`\beta `$-beam could clearly be a competitor to a neutrino factory, if technically feasible for such a high $`\gamma `$.
We have assumed initially that the number of ion decays per year is constant in $`\gamma `$ and we investigated the impact of external constraints to this assumption. We found that especially the high gamma setup suffers from modifications in the scaling of ion decays. This implies that the ion luminosity is a critical factor for this setup. In addition, we compare in Figure 14 setups optimal for the individual purposes, i.e., differently optimized $`\beta `$-beams with correspondingly optimized neutrino factories. If we required optimal sensitivity to all quantities, we would need two neutrino factory or two $`\beta `$-beam baselines. Though a storage ring with two decay sections could, in principle, be possible for a $`\beta `$-beam, it might be more compact in the neutrino factory case. If, however, one of the measurements becomes obsolete, such as the mass hierarchy which might be determined by a supernova explosion, the $`\beta `$-beam baseline optimization seems to be quite straightforward, whereas the neutrino factory may still require another baseline to resolve degeneracies. For a precision measurement of $`\delta _{\mathrm{CP}}`$, for example, the “magic baseline” could be required for unfortunate values of $`\delta _{\mathrm{CP}}`$ .
The $`\beta `$-beam has in all cases the advantage that the baseline can be freely chosen such that one is measuring at the oscillation maximum. For the case of a much smaller $`\mathrm{\Delta }m_{31}^2`$ than the current best-fit value, first tests suggest that the performance improves compared to neutrino factories. However, this less likely case requires further study because one should also optimize the neutrino factory for such a different value of $`\mathrm{\Delta }m_{31}^2`$.
We conclude that a lower gamma $`\beta `$-beam is certainly an interesting physics alternative to a large superbeam upgrade and a higher gamma $`\beta `$-beam could be an competitive alternative to a neutrino factory. In all cases, the attractiveness of the $`\beta `$-beams depends clearly on the ability to produce enough isotope decays for both neutrinos and anti-neutrinos. Especially for $`\delta _{\mathrm{CP}}`$ measurements, the $`\beta `$-beams might then outperform all existing techniques, whereas for $`\mathrm{sin}^22\theta _{13}`$ discovery and mass hierarchy sensitivity the neutrino factory is ultimately a better choice. Except from $`\mathrm{sin}^22\theta _{13}`$, mass hierarchy, and CP violation measurements, there is more physics to be done with a neutrino factory. In the neutrino oscillation sector, for example, there is sensitivity to the leading atmospheric parameters. The case $`\mathrm{sin}^22\theta _{13}=0`$ would suggest to use the $`\nu _\mu `$ disappearance channel at a very long baseline for mass hierarchy and MSW effect measurements . This role is unlikely to be replaced by the electron neutrino disappearance channel of a $`\beta `$-beam. Nevertheless, we conclude that $`\beta `$ beams constitute a very interesting option for future precision neutrino oscillation experiments. Further technological feasibility studies are clearly well motivated to explore if a $`\beta `$ beam can be realized. The technical feasibility, the financial effort, and the physics potential of a $`\beta `$-beam and a neutrino factory have to be compared then again before an ultimate decision is made.
#### Acknowledgments
We would like to thank M. Lindroos, P. Litchfield, M. Mezzetto, and L. Mualem for useful information. This work has been supported by SFB 375 and the Graduiertenkolleg 1054 of Deutsche Forschungsgemeinschaft. WW would like to acknowledge support from the W. M. Keck Foundation and NSF grant PHY-0070928. In addition, he would like to thank the theory groups at TUM and Wisconsin for their warm hospitality during his visits, where parts of this work have been carried out.
## Appendix A: Water Cherenkov detector and 2nd oscillation maximum
The better performance of Setup 1 at the longer baseline with $`L/\gamma =2.6`$ (although statistics drops with $`1/L^2`$) can be understood with Figure 15 where the comparison of $`L/\gamma =1.3`$ (left column) and $`L/\gamma =2.6`$ (right column) is shown. The appearance spectra for $`\mathrm{sin}^22\theta _{13}=0.01`$ and $`\delta _{CP}=\pi /4`$ is shown in the first row of Figure 15 for neutrinos (black curve) and anti-neutrinos (gray curve). For Setup 1 at $`L/\gamma =1.3`$ one can clearly see, that only the first oscillation maximum contributes to the whole appearance spectra, while for the one at $`L/\gamma =2.6`$ the first oscillation appearance maximum is shifted to higher neutrino energies and appearance events from the second oscillation maximum enter the energy window of the analysis from lower energies but the overall event rates are decreased. In the second row, we show the allowed regions in the $`\theta _{13}`$-$`\delta `$ plane at 1, 2, and 3 $`\sigma `$ for the same true oscillation parameters (indicated by the black dot) where only systematical errors are taken into account. The allowed regions for the $`L/\gamma =2.6`$ scenario are somewhat larger due to the lower statistics. But only if also correlations and degeneracies are included (third row) one can see the impact of the second oscillation maximum. The degenerate solution fitted with $`\mathrm{\Delta }m_{31}^2<0`$ is smaller for the $`L/\gamma =2.6`$ scenario and does not reach to higher values of $`\theta _{13}`$ than the region that contains the best-fit value. Additionally the $`\mathrm{\Delta }\chi ^2`$ value at the local minimum of the degenerate solution is much higher with $`L/\gamma =2.6`$ than the one with $`L/\gamma =1.3`$ and does not even appear at 1 $`\sigma `$.
In order to understand better the impact of the appearance events from the second oscillation maximum in the $`L/\gamma =2.6`$ scenario, we divided the whole data set of Setup 1 in two separate data sets which only contain the appearance events from the first or second oscillation maximum. Data set I reaches from 0.2 to 0.7 GeV and data set II reaches from 0.7 to 1.6 GeV as can be seen in the upper right picture of Figure 15. In Figure 16 we compare again the allowed regions in the $`\theta _{13}`$-$`\delta `$ plane at 1, 2, and 3 $`\sigma `$, now for data set I (left column), data set II (middle column) and both data sets combined (right column). In the first row we only consider systematics and in the second row correlations and degeneracies are switched on. One can clearly see that due to extremely low statistics with only the appearance events from data set 1 the allowed regions are strongly expanded. The allowed regions from data set II are highly improved and in the case ”systematics only” even somewhat smaller than for the combination of both data sets. This effect comes from the fact that most of the background events reconstruct at smaller energies (i.e. within data set I) and therefor the S/B ratio is smaller for data set I only. But for the degenerate solution fitted with $`\mathrm{\Delta }m_{31}^2<0`$ the combination of both data sets results in an improvement since the local degenerate minimum for data set II lies at a different point in the parameter plane than the one for data set I. |
warning/0506/physics0506021.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Economic development generates significant structural transformations, such as changes in the demographic condition and the production structure. As mentioned in chapter 10 of ref. the most important structural feature of developing economies is the distinction between rural and urban sectors. The agriculture plays a key role in the development of the urban sector of an economy: it provides both the food surplus that enables the urban sector to survive, and the supply of labor to the expanding industrial urban sector. As suggested in chapter 5 of ref. , the fundamental part of the transformation from mostly dispersed rural and agrarian countries in a more concentrated urban and industrial is the flux of a large number of individuals through migration from rural areas to urban areas.
In this paper we examine the rural-urban migration phenomena which takes place during the industrialization process. The analysis is carried out by using an agent-based computational model which aims to describe some of the main structural features of a developing economy. We look at the rural-urban migration as a discrete choice problem, which allows us to formalize the migration process by using an Ising like model. We modelled the migratory decision in the usual manner considering the force exerted by the difference of earnings between the sectors. Moreover, it is also included a new factor, the influence from neighbors like in the Ising model.
The paper is organized as follows. Section 2 describes the economic setting, i.e., the typical dual economic structure (rural versus urban sector) in industrializing countries. In section 3 the migration process is modelled within a statistical mechanics framework. Section 4 presents the simulations and the main results. Finally, section 5 shows the concluding remarks.
## 2 The Economic Setting
Let us consider an economic system formed by two sectors, one rural and the other urban. The main differences between these sectors are the sort of goods they produce, the technologies used by firms and the framework of wage determination. Such a dual structure is typically used by the economic literature which investigates the rural-urban migration phenomena . The basic features of the dual economy will be described in subsection 2.1 and 2.2. Subsection 2.3 shows how the equilibrium macrostate of the economic system is determined.
### 2.1 The urban sector
The urban productive sector is formed by firms specialized in the production of manufacturated goods. The output of the $`ith`$ firm $`Y_{mi}`$ depends positively on both the amount of employed workers $`N_{mi}`$ and the effort $`\epsilon `$, spent by each worker to perform his job. Based on the classical rural-urban migration theory , we assume that the stock of capital during the analysis period is given. Supposing a standard geometrical functional form , the production function of the manufacturing firm is the well-known Cobb-Douglas<sup>4</sup><sup>4</sup>4The modelling of employment and wage determination of the urban sector is based on the efficiency-wage approach. See chapter 10 of ref. and section V of ref. for further details.
$$Y_{mi}=A_m\left(\epsilon N_{mi}\right)^\alpha ,$$
(1)
where $`0<\alpha <1`$ and $`A_m>0`$ are parametric constants.
By using the functional form originally suggested by Summers and slightly modified by Romer , the urban worker’s effort can be defined as a function of the real wage paid by the manufacturing firm, the urban unemployment rate $`u`$ and the alternative wage $`w_m`$, which is paid by other firms of the same sector. Then the effort function is given by
$$\epsilon =\{\begin{array}{c}\left[\frac{w_{mi}(1bu)w_m}{(1bu)w_m}\right]^\eta ,\text{if}w_{mi}>(1u)w_m,\hfill \\ 0,\text{otherwise,}\hfill \end{array}$$
(2)
where $`0<\eta <1`$ and $`b>0`$ are parametric constants.
The $`Z_m`$ manufacturing firms which form the manufacturing sector seek to maximize their real profits, measured in units of the manufacturated good, by choosing wages and employment freely. Given eq. (1) and eq. (2) the real profit of $`i`$th manufacturing firm is
$$A_m\left[\left(\frac{w_{mi}(1bu)w_m}{(1bu)w_m}\right)^\eta N_{mi}\right]^\alpha w_{mi}N_{mi}.$$
(3)
The maximization condition of eq. (3) can be found using the first-order condition for a maximum, which result is<sup>5</sup><sup>5</sup>5The second-order condition for a maximum is also satisfied.
$$w_{mi}=\frac{(1bu)w_m}{1\eta }$$
(4)
and
$$N_{mi}=\left[\frac{\alpha A_m\eta ^{\alpha \eta }(1\eta )^{1\alpha \eta }}{(1bu)w_m}\right]^{\frac{1}{1\alpha }}.$$
(5)
In equilibrium, all these firms choose the same wage , i.e., $`w_{mi}=w_m`$ $`(i=1,2,3,\mathrm{},Z_m)`$. Then, from equation (4) the equilibrium urban unemployment rate is
$$u=\frac{\eta }{b}.$$
(6)
By definition, the urban unemployment rate is the ratio between the number of unemployed workers and the urban population $`(N_uZ_mN_{mi})/N_u,`$ where $`N_u`$ is the amount of workers localized in the urban sector. The previous definition must be consistent in each period to the equilibrium value of (6). The employment level of the manufacturing firm which obeys this consistency condition is obtained equaling the equilibrium in eq. (6) to the previous definition:
$$N_{mi}=\left(1\frac{\eta }{b}\right)\frac{N_u}{Z_m}.$$
(7)
Taking eq. (2), evaluated in the equilibrium, and eq. (7) and replace both in eq. (1), the aggregated production of the manufacturing sector, $`Z_mY_{mi}`$, is given by
$$Y_m=\xi _1N_{u}^{}{}_{}{}^{\alpha },$$
(8)
where $`\xi _1=A_mZ_m^{1\alpha }\left[\left(\frac{\eta }{1\eta }\right)^\eta \left(1\frac{\eta }{b}\right)\right]^\alpha `$.
By using eqs. (5), (6) and eq. (7) one can obtain the equilibrium wage of the manufacturing sector:
$$w_m=\xi _2N_{u}^{}{}_{}{}^{\alpha 1},$$
(9)
where $`\xi _2=\alpha A_m\left(\frac{\eta }{1\eta }\right)^{\alpha \eta }\left[\left(1\frac{\eta }{b}\right)\frac{1}{Z_m}\right]^{\alpha 1}`$.
Given the parametric constants that specify the technology and the effort sensitivity to wage, as well as the size of the manufacturing productive sector, it is possible to see that the equilibrium of the urban sector depends directly upon the urban population $`N_u`$.
### 2.2 The rural sector
In the rural sector the farm $`i`$ produces an agricultural output $`Y_{ai}`$ by employing an amount of workers $`N_{ai}`$. The output is obtained by using a Cobb-Douglas production function
$$Y_{ai}=A_a\left(N_{ai}\right)^\varphi ,$$
(10)
where $`0<\varphi <1`$ and $`A_a>0`$ are parametric constants. We suppose that both the land endowment and the stock of capital of the farm are given during the period of analysis as assumed by refs. and .
Differently from the urban sector, farms are price-takers and the real wage is adjusted up to the point in which there is no unemployment in this sector . This implies that the rural population will match the aggregated employment in the rural sector. Therefore, the equilibrium employment level of the farm $`i`$ is
$$N_{ai}=\frac{NN_u}{Z_a},$$
(11)
where $`Z_a`$ is the amount of farms which constitute the agricultural sector and $`N`$ is the total population of the economic system.
From eq. (10) and eq. (11) the aggregated production of the rural sector, $`Z_aY_{ai}`$, is
$$Y_a=\xi _3\left(NN_u\right)^\varphi ,$$
(12)
where $`\xi _3=A_aZ_a^{1\varphi }`$.
Thus, the profit maximizing of the farms imply that the rural real wage expressed in units of the manufactured good becomes equal to the marginal product of agricultural labor in units of manufacturing good<sup>6</sup><sup>6</sup>6This marginal product is the derivate of the production function, eq. (10), with respect to $`N_{ai}`$ multiplied by $`p`$.:
$$w_a=\xi _4p\left(NN_u\right)^\varphi ,$$
(13)
where $`\xi _4=\left(A_a\varphi /Z_a^{\varphi 1}\right)`$ and $`p`$ is the price of the agricultural good expressed in units of the manufactured good.
Like in the urban sector, the equilibrium state of the rural sector depends on the urban population, as the size of total population of the economy is fixed.
### 2.3 The macrostate of economic system
As proposed by Harris and Todaro , the terms of trade between the rural and urban sectors, measured by the price $`p`$, depend on the relative scarcity of agricultural and manufacturated goods. This can be measured by the ratio $`Y_m/Y_a`$. The greater this ratio the greater will be the scarcity of agricultural good, which implies an increase of the agricultural good price in units of manufacturated good. Formally, given the urban population, the equilibrium relative price of the agricultural good is<sup>7</sup><sup>7</sup>7In the literature on rural-urban migration it is usual, because of analytical simplicity, to consider $`p`$ constant . This is true in the special case when $`\gamma =0`$ in eq. (14).
$$p=\rho \left(\frac{Y_m}{Y_a}\right)^\gamma $$
(14)
where $`\rho >0`$ and $`\gamma >0`$ are parametric constants.
Therefore, given the size of urban population, by using equations (6-9) one can calculate the state of urban sector. Likewise, the rural sector state is determined by means of equations (11-14). The equilibrium state of both sectors will be modified if a migratory flux changes the population distribution of the economic system.
## 3 Migratory process: a statistical mechanics approach
As argued by Harris and Todaro , individuals take their decisions of migrating or not by considering the differential of expected wages between their present sector and the sector they intend to go. However, other authors have taken into account additional reasons. Based on the formalization from statistical mechanics applied to socioeconomic phenomena , in this section we propose an agent-based computational model to describe the rural-urban migratory process. This model is focused on the influence that individuals suffer in the reference group that they are included. The emergent properties will be analyzed taking into account the standard effect of labor allocation Harris-Todaro mechanism, which is based on the expected differential wages between sectors. This analysis will also be concerned on the effect of social neighborhood, often mentioned by other authors but not yet formalized.
The main feature of the decision process is that each worker reviews his sectorial location after a period of time spent in that sector. We exclude, by assumption, the possibility that the worker may simultaneously supply his labor force to both sectors. Thus, only two choices are admitted: stay in the sector in which he was during previous periods or migrate.
In order to model the migration process by allowing only discrete choices, each worker has its state defined by $`\sigma _i\{1,+1\}`$, where $`\sigma _i=1`$ means that the worker is at the rural sector; otherwise, $`\sigma _i=+1`$, representing the urban sector.
In our model, during the decision process, explicit and observable incentives are taken into account by each potential migrant. This is called a deterministic private utility , given by
$$U_i=H(t)\sigma _i,$$
(15)
where $`H(t)=k\omega _e`$, $`k>0`$ is a parametric constant and $`\omega _e`$ is the expected urban-rural differential wage. The expected urban-rural wages in function of $`H(t)`$ are specified as follows.
Jobs are allocated at random when manufacturing firms are faced with more applicants than jobs avaliable . It means that in each time step all urban workers have the same probability to find an urban job. Under such a hypothesis, the term $`(1u)`$ is the probability of an urban worker to obtain a job. Hence, $`(1u)w_m`$ is the expected urban wage. Assuming that the rural wage is perfectly flexible there is no unemployment in the rural sector. Then, the probability to find a job in the rural sector is $`1`$. Therefore, the rural wage $`w_a`$ is the same as the expected wage in this sector. In sum, the expected differential of wage between urban and rural sectors is
$$\omega _e=(1u)w_mw_a.$$
(16)
Besides, the worker $`i`$ is also under the influence of other workers, his social neighborhood , denoted by $`n_i`$. The measure of such influence, that is, the deterministic social utility , is given by
$$S_i=J\underset{jn_i}{}\sigma _i\sigma _j,$$
(17)
where $`J>0`$ is a parametric constant. The term $`J`$ represents the interaction weight which relates the worker $`i`$’s choice to the neighbor $`j`$’s choice. This is assumed to be nonnegative, by representing the hypothesis that the worker seeks to conform to the behavior of his neighbors . The interactions among neighbors are assessed in the workers’ nearest neighbors or in the next nearest neighbors.
Then, following references and , we assume that payoff of worker $`i`$, which is his deterministic total (private and social) utility can be obtained replacing eq. (16) in eq. (15) and summing with eq. (17):
$$_i=k\left[(1u)w_mw_a\right]\sigma _i+J\underset{jn_i}{}\sigma _i\sigma _j.$$
(18)
Therefore, this system can be described by the well-known ferromagnetic Ising model, in the presence of an external time-dependent magnetic field:
$$=H(t)\underset{i=1}{\overset{N}{}}\sigma _iJ\underset{<ij>}{}\sigma _i\sigma _j.$$
(19)
In each time step, each worker reviews his decision about the sectorial location with probability $`a`$, called activity . Then, there is a part of the population that reviews their decisions and becomes potential migrants.
The potential migrant $`i`$ becomes an actual migrant depending on the comparison between his deterministic total utility $`_i`$ and his non observable and idiosyncratic motivations $`\mu _i`$, called random private utility . The term $`\mu _i`$ represents the net difference between the random private utilities that the potential migrant assigns to the sector he intends to move and his present sector.
In each period, if $`\mu _i>_i`$, the potential migrant $`i`$ becomes an actual migrant; otherwise, this does not happen. Supposing that $`\mu _i`$ is a random variable logistically distributed , the probability that the potential migrant effectively migrates is given by a cumulative distribution:
$$Pr_i=\frac{1}{1+e^{\beta _i}},$$
(20)
where $`\beta >0`$ is a parametric constant that in this context measures the heterogeneity of workers concerning to the migration propensity. Equation (20) is a measure of the probability that a worker $`i`$, who is reviewing his location strategy, stays in the sector that he is localized at that time. The higher his deterministic total utility, eq. (18), the higher the probability that no change will take place.
## 4 Simulation
To carry out the simulation of the economic system described in the previous sections, each worker is placed in one of the sites of a square lattice. The state of each site (worker) is set as mentioned before: $`\sigma _i=+1`$ for urban workers and $`\sigma _i=1`$ for rural ones. It is important to emphasize that the state of these sites represent the sectorial allocation of each worker, i.e., whether an individual is suppling his labor force in the urban or rural sector. It means that the coordinates of the lattice sites are not related to spatial distribution of workers.
To set up the initial state of the system, all workers are randomly distributed in the lattice. At time $`t=0`$ there is the initial condition that $`20\%`$ of the population is urban. In other words, initially, $`20\%`$ of the sites will be assigned with $`\sigma _i=+1`$ and the remaining $`80\%`$, $`\sigma _i=1`$. The reason for this initial distribution is because these are the values which have usually been observed in developing countries before the urbanization process initiates.
The next step in the simulation is to calculate the equilibrium state variables of the urban sector, by using eqs. (6-9), and of the rural sector by using eqs. (11-14). Since the state variables of both sectors are known, it is necessary to define the amount of workers that will review their sectorial location, i.e. those one who will become potential migrants. To do this, it is assumed that the probability that a worker will become a potential migrant is given by the activity $`a`$, as defined by Stauffer and Penna . All those selected as potential migrants will have their private utility calculated by eq. (18).
In order to conclude the reviewing process, the probability defined in eq. (20) is assessed. Then, a random number $`rn[0,1]`$ is generated from an uniform distribution. If $`rn>Pr`$, then the potential migrant becomes an actual migrant; otherwise, no change takes place.
As soon as the potential migrants end their reviewing process, a new sectorial distribution is obtained. Knowledge of the new urban population allows the macrostate of the economic system to be reset. Therefore, the state variables of both sectors have to be calculated again. The whole procedure described above will be repeated as many times as we set in the simulation. The stopping criteria used by us is halting the simulation some steps after the moment when the system reaches equilibrium.
Figure 1 shows the proportion of workers in the urban sector $`n_u\frac{N_u}{N}`$, from now on called urban share, plotted in three different combination of the parameters $`J`$ and $`k`$. It is necessary to remind that the parameters $`J`$ and $`k`$ adjust the instensity of the deterministic private utility, eq. (15), and deterministic social utility, eq. (17), respectively. From top to bottom the set of parameters used in the plotting are $`(J>0`$, $`k>0)`$, $`(J=0`$, $`k>0)`$ and $`(J>0`$, $`k=0)`$.
Firstly, consider the case $`(J=0`$, $`k>0)`$ plotted in Fig. 1. In this case, the review conducted by the agents is guided only by the deterministic private utility, which in turn depends on the expected urban-rural difference of wages. As in models of classical theory of migration , when the expected urban wage is higher than the rural wage, it implies in a continuous growth of the urban share, as well as a relatively fast convergence towards the equilibrium.
Secondly, consider the case where both effects are taken into account, $`(J>0`$, $`k>0)`$. Like the previous case, the rural-urban migratory process occurs again, however, the system reaches a higher value of the equilibrium urban share, though it takes more time for such outcome. This difference is caused by the parameter $`J>0`$, what means that the influence of the social neighborhood is considered. To better understand this behavior, it should be reminded that the process of sectorial position revision depends on the deterministic private utility and the social private utility. Then, when $`J>0`$ the influence of social neighborhood is being exerted, i.e., each worker attempts to adjust his choice according to the sectorial position of his neighbors. The existence of such an influence causes two different effects during the process of convergence towards equilibrium. In the first moment, when the neighborhood are mainly rural, the influence from neighbors slows the rural-urban migratory flux, increasing the time necessary to reach equilibrium. In the second moment, when the neighborhood become mainly urban, the influence reinforces the attraction from the high expected urban wage, leading to higher equilibrium urban share.
Finally, we consider the case $`(J>0`$, $`k=0)`$, with only neighborhood effects shown. In this case, the potential migrants consider only the sectorial position of the neighborhood and do not take into account the expected differential of wages. The pure effect due from neighborhood leads to the extinction of the urban sector. This is not an empirically important case, as it has not been observed in developing economies.
In Figure 2, another important feature caused by the migratory dynamics is the expansion of per capita income $`(Y_m+pY_a)/N`$. This result matches to the economic data in which in countries with high per capita income the proportion of the population living in rural area is low .
In the initial state of the system the configuration was randomly set with 20% of the sites assigned $`\sigma _i=+1`$, urban workers, and the rest $`\sigma _i=1`$, rural workers as shown in Fig. 3a. The final state of the dynamics by using $`(J>0`$, $`k>0)`$ can be visualized in Figure 3b. Now the infinite cluster is formed by sites $`\sigma _i=+1`$ representing the urban concentration caused by the migratory process. Several others clusters are formed by sites $`\sigma _i=1`$.
Figures 4 and 5 show the average magnetization $`m=\sigma _i/N`$ and the expected wages ratio $`r_e(1u)w_m/w_a`$, respectively. Both figures are plotted as function of the ratio $`J/k`$ ($`k`$ kept constant) measuring the relative intensity between these parameters.
Figure 4 has plotted in its vertical axis the average magnetization calculated during a period after the system have reached equilibrium. To values of $`J/k`$ less than the critical threshold the net magnetization is $`m0.4`$ representing an urban share about $`n_u=0.70`$. By increasing the ratio $`J/k`$ after this critical threshold the system goes to a new regime, changing completely its net magnetization.
Figure 5 is a plotting of expected wage ratio as function of $`J/k`$. To values $`J/k9.0`$ the ration is $`r_e1.0`$, what indicates that the expected urban wage and the rural wage converge to the same value. This property is known as Harris-Todaro equilibrium condition . Hence, in a economic system where internal migration occurs freely the absolute difference between the rural and urban wages can persist if workers consider the possibility of unemployment. After the threshold $`J/k>9.0`$, $`r_e`$ has its maximum value around $`2.8`$ which shows that the urban expected wage is $`2.8`$ times greater than the rural wage. Even having this ratio increasing the value of the worker private utility, eq. (15), the equilibrium of the system is $`m0.29`$, i.e., a rural concentration of $`64.9\%`$. The explanation of this outcome is that after a given threshold the values of $`J`$ are in such a range that the social utility, eq. (17), is many times higher than the private utility. In other words, in such range, it does not matter if the expected wage is attractive in the urban sector because the strongest factor in the migration decision is the influence of the neighborhood, i.e., agents tend to mimic the behavior of other agents.
Simulations plotted in Figure 6 indicate that when the size $`N`$ of the lattice increases the equilibrium urban share $`n_u`$ will change. For a given heterogeneity of the agents $`\beta `$, there is a power law relating equilibrium urban share and the inverse of lattice size. This can be formalized in the expression below
$$n_u=A\left(\frac{1}{N}\right)^\theta ,$$
(21)
where $`A`$ and $`\theta `$ are constants which have to be estimated. To carry out the estimation of these constants we evaluated a linear regression of the log-linear version of eq. (21). In Table 1 one can find the estimation of the constants $`A`$ and $`\theta `$ based on data generated for five different values of $`\beta `$.
The estimation of the constants are approximately the same when using slightly different values of agent heterogeneity $`\beta `$. For example, by using any pair of constants $`A`$ and $`\theta `$ from Table 1, the estimation of equilibrium urban share by eq. (21) is $`n_u=0.61`$ for an economy with 50 million of workers.
| $`\beta `$ | | $`A`$ | | $`\theta `$ |
| --- | --- | --- | --- | --- |
| 1.5 | | 1.064 | | 0.032 |
| 2.0 | | 1.061 | | 0.031 |
| 2.5 | | 1.061 | | 0.031 |
| 3.0 | | 1.064 | | 0.032 |
| 3.5 | | 1.066 | | 0.032 |
Table 1. Estimates of parameters $`A`$ and $`\theta `$ for different values of $`\beta `$.
The effects of the ratio $`J/k`$, together with size of population $`N=L^2`$, are shown in Fig. 7. The different values of equilibrium urban share are plotted in a grey scale. The first property observed in this figure is the existence of several phase states which depend on the values of $`J/k`$ and $`N`$. Each phase state is characterized by a constant equilibrium urban share. The topology of Figure 7 is in agreement with the results shown in Figs. 4 and 6, demonstrating that the properties of equilibrium macrostate depends on the combination of these parameters.
In Figure 8 is plotted the equilibrium urban share as function of the parameter $`\beta `$. For values of $`\beta `$ tending to zero the equilibrium urban share tends to 0.5 (or $`m=0`$), which implies in a null urban concentration (null average magnetization), even though there is an expected urban wage higher than the rural wage. In fact, eq. (20) shows that the smaller $`\beta `$ the higher the idiosyncratic and non-observed proportion of the worker’s behavior related to the migration propensity. If $`\beta =0`$, the choices $`\sigma _i=+1`$ and $`\sigma _i=1`$ have the same probability to occur being independent of the expected differential of wages. In sum, when the heterogeneity of the workers related to the decision of migration increases, the urban concentration will decline in the long run. On the other hand, when the heterogeneity of the agents decreases, i.e., $`\beta `$ increases, the equilibrium urban share is invariable after a threshold.
The decision of migration is not taken simultaneously by all individuals. In order to simulate this behavior, the parameter called activity $`a`$ is used. It gives the probability that a worker will review the decision about his sectorial location. More specifically, $`a`$ represents the fraction of the population which will go through the reviewing process. This fraction of individuals is randomly selected and changes in each time step. In Fig. 9 variation of $`a`$ in different simulations shows that the time needed for the system to reach equilibrium is proportionally inverse to the value of the activity. Therefore, the time needed to reach the equilibrium state is strongly related to the amount of individuals which review their sectorial decision.
## 5 Conclusion
This paper has developed an agent-based computational model to analyse the rural-urban migration phenomena. The basic scenario was made of an economic system formed by two sectors, rural and urban, which differed in term of the goods produced, the production technology and the mechanism of wage determination.
By assuming the sectorial migration decision as discrete choice in a milieu of decentralized and non-coordinated decision-making, the rural-urban migration process was formalized as an Ising like model. The simulations showed aggregate regularities which indicates that decentralized migration decisions can lead to the emergence of equilibrium macrostates with features observed in developing economies. First, the simulation having an initial macrostate with population predominantly rural and expected urban wage higher than rural wage provoked a transitional rural-urban migratory dynamics, with continuous growth of the urban share. This is a key feature of the phenomena called in ref. as urban transition.
Second, simulations also showed that, during the rural-urban migration process, the reduction of the rural share takes place together with the increasing of per capita income of the economy. Such an inverse relation between rural share and per capita income is one of the most robust facts detected in economic statistics .
Third, the transitional rural-urban migratory dynamics converged towards an equilibrium macrostate. The features of this transitional dynamics and equilibrium are sensitive to the relative weight between private and social effects (utilities) as well as the degree of heterogeneity of agents concerning the migration propensity. When the social interaction component is relatively stronger and below a critical threshold the transitional dynamics towards equilibrium is delayed and reaches a higher equilibrium urban share. With a high heterogeneity of agents, $`\beta 0`$, this generates the end of the pulling force due the high expected urban wage what makes the system to reach an equilibrium macrostate with an urban share $`n_u=0.5`$. On the other hand, with a moderate heterogeneity of agents, $`\beta >1`$, the equilibrium urban shares will be set in a empirically reasonable range ($`n_u0.6`$).
The analysis shown in this paper suggests that a deeper investigation can still be carried out, which adopt alternative hypothesis mainly regarding the private and social utilities as well as other assumptions employed in our model.
## Acknowledgments
We would like to thank Dietrich Stauffer, D. N. Dias, T. Lobo for their contributions and Dr. Renato P. Colistete for his comments. Jaylson J. Silveira acknowledges research grants from CNPq. Aquino L. Espíndola thanks CAPES for the financial support. T. J. P. Penna thanks CNPq for the fellowship. |
warning/0506/math0506276.html | ar5iv | text | # Hilbert-Schmidt groups as infinite-dimensional Lie groups and their Riemannian geometry
## 1. Introduction
The results of this article are inspired by our previous study of heat kernel measures on infinite-dimensional Lie groups in , , , . The main tool in these papers was the theory of stochastic differential equations in infinite dimensions. The present paper, however, is entirely non probabilistic. It is organized as follows: in Sections 2 and 3 we discuss the exponential map for a certain class of infinite-dimensional groups, and in Sections 5 and 6 we introduce notions of Riemannian geometry for Hilbert-Schmidt groups and compute the Ricci curvature for several examples.
### 1.1. Motivation: Wiener measures and geometry
In our previous papers we were concerned with a pair of infinite dimensional Lie groups, $`G_{CM}G_W`$, related to each other in much the same way that the Cameron-Martin Hilbert space, $`H_1([0,1])`$, is related to Wiener space, $`C_{}([0,1])`$: it is well understood that the geometry of the Hilbert space $`H_1([0,1])C_{}([0,1])`$ controls Wiener measure on $`C_{}([0,1])`$, even though $`H_1([0,1])`$ is a subspace of Wiener measure zero. In the papers , , , we constructed an analog of Wiener measure on an infinite dimensional group $`G_W`$ as the “heat kernel” (evaluated at the identity of $`G_W`$) associated to the Laplacian on the dense subgroup $`G_{CM}`$. To this end one must choose an inner product on the Lie algebra, $`𝔤_{CM}`$, of $`G_{CM}`$ in order to introduce a left invariant Riemannian metric on $`G_{CM}`$. The Lie algebra $`𝔤_{CM}`$ determines a Laplacian, whose heat kernel measure actually lives on the larger group $`G_W`$. It can be shown that in general $`G_{CM}`$ is a subgroup of measure zero. More features of this group have been discussed in , , . In these papers we constructed the heat kernel measure by probabilistic techniques. We used a stochastic differential equation based on an infinite dimensional Brownian motion in the tangent space at the identity of $`G_W`$ whose covariance is determined by the inner product on $`𝔤_{CM}`$.
Just as in the case of the classical Cameron-Martin space, where the Sobolev norm on $`H_1([0,1])`$ is much stronger than the supremum norm on $`C([0,1])`$, so also the norm on $`𝔤_{CM}`$ must be much stronger than the natural norm on the tangent space at the identity of $`G_W`$ in order for the heat kernel measure to live on $`G_W`$. If $`G_W`$ is simply the additive group, $`C_{}([0,1])`$, and $`G_{CM}`$ is the additive group $`H_1([0,1])`$, then this heat kernel construction reproduces the classical Wiener measure.
In this paper we address questions relating to the geometry of $`G_{CM}`$, with a view toward eventual application to further understanding of the heat kernel measure on $`G_W`$. On a finite dimensional Riemannian manifold properties of the heat kernel measure are intimately related to the Ricci curvature of the manifold, and in particular to lower bounds on the Ricci curvature . We are going to compute the Ricci curvature for several classes of infinite dimensional groups, in particular, for those groups $`G_{CM}`$ whose heat kernel measure on $`G_W`$ we have already proven the existence of in , , , . Our results show that the Ricci curvature is generally not bounded below, even in the cases when we were able to construct the heat kernel measure on $`G_W`$ (e.g. for the group $`SO_{HS}`$). One of the implications of our results is that the methods used to prove quasi-invariance of the heat kernel measure in the finite-dimensional case are not applicable for the settings described in our earlier papers. We also compute the Ricci curvature of groups $`G_W`$. All these groups are Hilbert-Schmidt groups which are described below.
### 1.2. Hilbert-Schmidt groups as Lie groups and their Riemannian geometry
Denote by $`HS`$ the space of Hilbert-Schmidt operators on a real separable Hilbert space $`H`$. Let $`B(H)`$ be the space of bounded operators on $`H`$, and let $`I`$ be the identity operator. Denote by $`GL(H)`$ the group of invertible elements of $`B(H)`$. The general Hilbert-Schmidt group is $`GL_{HS}=GL(H)(HS+I)`$. In the papers , , , we proved the existence and some basic properties of the heat kernel measures on certain classical subgroups of $`GL_{HS}`$, namely, $`SO_{HS}`$ and $`Sp_{HS}`$. The Lie algebras of these groups are closed subspaces of $`HS`$ in the Hilbert-Schmidt norm. In the setting described above $`GL_{HS}`$, $`SO_{HS}`$ and $`Sp_{HS}`$ are examples of the group $`G_W`$.
But the corresponding Cameron-Martin subgroups are, of necessity, only dense subgroups of $`G_W`$. They are determined by their tangent space, $`𝔤_{CM}`$, at the identity. In order to get the corresponding heat kernel measure to live on $`GL_{HS}`$, $`SO_{HS}`$ or $`Sp_{HS}`$, respectively, the tangent space $`𝔤_{CM}`$ must be given a Hilbert norm $`||`$ which is much stronger than the Hilbert-Schmidt norm. The result is that the commutator bracket of operators may not be continuous in this norm. That is, $`|ABBA|C|A||B|`$ may fail for any constant $`C`$ as $`A`$ and $`B`$ run over $`𝔤_{CM}`$. Consequently $`𝔤_{CM}`$ may not really be a Lie algebra. Rather, the commutator bracket may be only densely defined as a bilinear map into $`𝔤_{CM}`$. In Section 4 we will give a class of examples of groups contained in $`GL_{HS}`$ such that $`𝔤_{CM}`$ is not closed under the commutator bracket and in this sense is not a Lie algebra. But in most of the examples we consider this is not the case.
In this paper we are going first to address the problem of the relation of the tangent space, $`T_I(G_{CM})`$, to the Lie algebra structure of some dense subspace of $`T_I(G_{CM})`$. As in the case of classical Wiener space, where polygonal paths play a central technical role in the work of Cameron and Martin, so also it seems to be unavoidable, for the purposes of , , , , to make use of a group $`GG_{CM}`$ which is in some sense dense in $`G_{CM}`$ and which is itself a union of an increasing sequence of finite dimensional Lie groups: $`G=\underset{n=1}{\overset{\mathrm{}}{}}G_n`$, $`G_nG_{n+1}`$. In Sections 2 and 3 all of our groups will be taken to be subgroups of $`GL(H)`$ rather than $`GL_{HS}`$. Then the Lie algebra of $`G_n`$, $`𝔤_n`$, is a finite dimensional subspace of $`B(H)`$, and is closed under the commutator bracket. Let $`𝔤=\underset{n=1}{\overset{\mathrm{}}{}}𝔤_n`$. Then $`𝔤`$ is a also a Lie algebra under the commutator bracket. But $`𝔤`$ cannot be complete in any norm in the infinite dimensional case.
Our goal in Sections 2 and 3 is to study the completion, $`𝔤_{\mathrm{}}`$, of $`𝔤`$ in some (strong) Hilbert norm. In these sections we will assume that the completion $`𝔤_{\mathrm{}}`$ actually embeds into $`B(H)`$ and that the commutator bracket is continuous in this norm. We will characterize the group, $`G_{CM}`$, generated by $`\mathrm{exp}(𝔤_{\mathrm{}})`$ in this case, and show that the exponential map covers a neighborhood of the identity. These groups are examples of so called Baker-Campbell-Hausdorff Lie groups (e.g. , , , , , , , ). Let us mention here that the question of whether the exponential map is a local diffeomorphism into an infinite-dimensional Lie group has a long history. Our treatment is different in two major aspects. The first one is the choice of an inner product on $`𝔤`$ and corresponding norm on $`𝔤`$. As we mentioned earlier the heat kernel analysis on $`GL_{HS}`$ forces us to choose an inner product on $`𝔤`$ which is different from the Hilbert-Schmidt inner product. We will assume that the commutator bracket is continuous on $`𝔤_{\mathrm{}}`$ in the extended norm $`||`$, namely, $`|[x,y]|C|x||y|`$, where the constant $`C`$ is not necessarily $`2`$. In most results on Banach-Lie groups this constant is assumed to be $`2`$ (e.g. , ). Quite often the underlying assumption is that $`𝔤`$ is a Banach algebra, and that $`𝔤`$ is complete. None of these is assumed in our case since we wish to deal with examples without these restrictions.
In Sections 5 and 6 we will compute the Ricci curvature in two major cases: when the norm on $`𝔤`$ is the Hilbert-Schmidt norm, and when it strongly dominates the Hilbert-Schmidt norm. Then we will examine how the lower bound of the Ricci curvature depends on the choice of the strong Hilbert norm. We will extend Milnor’s definitions of curvature on Lie groups to our infinite dimensional context for this purpose. Our results show that the Ricci curvature is generally not bounded below, and in some cases is identically minus infinity.
### 1.3. Historical comments
We give references to the following mathematical literature addressing different features of infinite-dimensional Lie groups and exponential maps. Our list is certainly not complete, since the subject has been studied for many years. There are several reviews on the subject (e.g. , , , , ). Possible non-existence of an exponential map is addressed in , , . Super Lie algebras of super Lie groups have been studied in , . Direct and inductive limits of finite-dimensional groups and their Lie group structures have been discussed in , , , , , , , . Our main tool in proving that the exponential map is a local diffeomorphism is the Baker-Campbell-Dynkin-Hausdorff formula. Note that in the terminology of Banach-Lie groups it means that we prove that the groups we consider are Baker-Campbell-Hausdorff Lie groups (e.g. , , , , , , , ). Some of these articles also address the issue of completeness of the space over which a Lie group is modeled. We show that under Completeness Assumption 2.1 and Assumption 3.1 the Cameron-Martin group $`G_{CM}`$ is complete in the metric induced by the inner product on $`𝔤`$. In addition, in Section 4 we show that a natural completion of the infinite-dimensional Lie algebra $`𝔤`$ is not a Lie algebra in general without these assumptions.
One of the main contributors to the field of connections between differential geometry and stochastic analysis is P. Malliavin, who wrote a survey on the subject in . A book on stochastic analysis on manifolds has been written by E.Hsu . In conclusion we refer to works of B. Driver, S. Fang, D. Freed, E. Hsu, D. Stroock, I. Shigekawa, T. Wurzbacher et al dealing with infinite-dimensional Riemannian geometry and its applications to stochastic analysis (, , , , , , , , , , , , , , , , , , , , , ). They mostly concern loop groups, path spaces, their central extensions etc. In these cases the Riemannian geometry on the infinite-dimensional manifold is induced by the geometry of the space in which the loops or paths lie. The situation we consider is quite different.
Acknowledgment. I thank B. Driver and P. Malliavin who asked me about the Ricci curvature of the Hilbert-Schmidt groups, and especially B. Driver who suggested to use J. Milnor’s results for finite-dimensional Lie groups. Section 3 benefited greatly from my discussions with A.Teplyaev about the exponential map. Finally, I am very grateful to L.Gross for carefully reading the manuscript and suggesting significant improvements to the text.
## 2. Completeness Assumption and the definition of the Cameron-Martin group
In this section we study a Lie group associated with an infinite-dimensional Lie algebra. The results of this section are not restricted to the Hilbert-Schmidt operators. We begin with an informal description of the setting. Let $`𝔤`$ be a Lie subalgebra of $`B(H)`$, the space of bounded linear operators on a separable Hilbert space $`H`$. The group under consideration is a subgroup of $`GL(H)`$, the group of invertible elements of $`B(H)`$. The space $`B(H)`$ is the natural (infinite-dimensional) Lie algebra of $`GL(H)`$ with the operator commutator as the Lie bracket.
We assume that $`𝔤`$ is equipped with a Hermitian inner product $`(,)`$, and the corresponding norm is denoted by $`||`$. In an infinite-dimensional setting $`𝔤`$ might not be complete. We will always work with the situation when $`𝔤`$ has a completion which is a subspace of $`B(H)`$. But as we will see in Section 4, the most natural candidate for such a completion of $`𝔤`$ might not be closed under taking the Lie bracket. Similarly, when we look at the Lie group corresponding to the infinite-dimensional Lie algebra $`𝔤`$, it might not be complete in the metric induced by the inner product on the Lie algebra.
The infinite-dimensional Lie algebra $`𝔤`$ is described by finite-dimensional approximations. Let $`G_1G_2\mathrm{}G_n\mathrm{}B(H)`$ be a sequence of connected finite-dimensional Lie subgroups of $`GL(H)`$. Denote by $`𝔤_nB(H)`$ their Lie algebras. We will consider the Lie algebra $`𝔤=\underset{n=1}{\overset{\mathrm{}}{}}𝔤_n`$.
###### Assumption 2.1 (Completeness Assumption).
Throughout this section we assume that there is a subspace $`𝔤_{\mathrm{}}`$ of $`B(H)`$ such that the Lie algebra $`𝔤`$ is contained in $`𝔤_{\mathrm{}}`$ and the given inner product $`(,)`$ on $`𝔤`$ extends to $`𝔤_{\mathrm{}}`$, which is complete with respect to this inner product. We will abuse notation by using $`(,)`$ to denote the extended inner product on $`𝔤_{\mathrm{}}`$ and by $`||`$ the corresponding norm. We also assume that $`𝔤`$ is dense in $`𝔤_{\mathrm{}}`$ in the norm $`||`$.
We will discuss this assumption in more detail in Section 4 in the case of the Hilbert-Schmidt groups. In particular, we will show that $`𝔤_{\mathrm{}}`$ is not a Lie algebra in some cases of particular interest.
###### Notation 2.2.
Let $`C_{CM}^1`$ denote the space of paths $`g:[0,1]GL(H)`$ such that
1. $`g(s)`$ is continuous in the operator norm,
2. $`\dot{g}=\frac{dg}{ds}`$ exists in $`B(H)`$ equipped with the operator norm,
3. $`\dot{g}`$ is piecewise continuous in the operator norm,
4. $`g^{}=g^1\dot{g}`$ is in $`𝔤_{\mathrm{}}`$, and $`g^{}`$ is piecewise continuous in the norm $`||`$.
Let
$$d(y,z)=\underset{g}{inf}\{\underset{0}{\overset{1}{}}|g^1\dot{g}|ds\},$$
where $`g`$ runs over $`C_{CM}^1`$ with $`g(0)=y,g(1)=z`$. We set $`d(y,z)=\mathrm{}`$ if there is no such path $`g`$. Note that $`d`$ depends on the norm $`||`$ on $`𝔤`$.
###### Notation 2.3.
$`G_{CM}=\{xB(H):d(x,I)<\mathrm{}\}`$.
###### Proposition 2.4.
$`G_{CM}`$ is a group, and $`d`$ is a left-invariant metric on $`G_{CM}`$.
###### Proof.
The proof for the first part is the same as for any finite-dimensional Lie group. In particular, if $`f:[0,1]G_{CM},f(0)=x,f(1)=y`$, $`x,yG_{CM}`$, then $`h(s)=y^1f(s)`$ is a curve connecting $`y^1x`$ and $`I`$, and $`\left|h^1\dot{h}\right|=\left|f^1\dot{f}\right|`$. Therefore $`d(y^1x,I)=d(x,y)`$, thus $`y^1xG_{CM}`$. ∎
###### Definition 2.5.
$`G_{CM}`$ is called *the Cameron-Martin group*.
## 3. The Cameron-Martin group and the exponential map
###### Assumption 3.1 (Continuity Assumption on the Lie bracket).
Throughout this section we assume that $`𝔤_{\mathrm{}}`$ is closed under taking the commutator bracket, and that the commutator bracket is continuous on $`𝔤_{\mathrm{}}`$, that is, there is $`C>0`$ such that
$$\left|[x,y]\right|C\left|x\right|\left|y\right|$$
for all $`x,y𝔤_{\mathrm{}}`$.
###### Remark 3.2.
Assumption 3.1 is satisfied for any Banach algebra with $`C=2`$. In particular, it holds for the operator norm $``$ and the Hilbert-Schmidt norm $`||_{HS}`$, but these are not the norms we are going to consider. The necessity to use the space $`𝔤_{\mathrm{}}`$ with the norm $`||`$ is dictated by the needs of the heat kernel measure construction carried out in , , , .
###### Remark 3.3.
Assumption 3.1 implies that the operator $`adh`$ is bounded on $`𝔤_{\mathrm{}}`$, namely, $`adhC|h|`$ where the constant $`C`$ is as in Assumption 3.1.
###### Theorem 3.4.
If Assumption 3.1 is satisfied, then the exponential map is a diffeomorphism from a neighborhood of $`0`$ in $`𝔤_{\mathrm{}}`$ onto a neighborhood of $`I`$ in $`G_{CM}`$.
As before let $`G_1G_2\mathrm{}G_n\mathrm{}B(H)`$ be a sequence of connected finite-dimensional Lie subgroups of $`GL(H)`$.
###### Theorem 3.5.
Suppose Assumption 3.1 holds. Then the group $`G_{\mathrm{}}=\underset{n=1}{\overset{\mathrm{}}{}}G_n`$ is dense in $`G_{CM}`$ in the metric $`d`$, and the Cameron-Martin group $`G_{CM}`$ is complete in the metric $`d`$.
In order to prove the surjectivity of the exponential map onto a neighborhood of the identity in the Cameron-Martin group $`G_{CM}`$, it is necessary to prove that, for example, if $`x`$ and $`y`$ are small (in the norm $`||`$) elements of $`𝔤_{\mathrm{}}`$, then $`e^xe^y=e^z`$ for some element $`z`$ in $`𝔤_{\mathrm{}}`$. This makes it unavoidable to use some version of the Baker-Campbell-Dynkin-Hausdorff (BCDH) formula (e.g. , ).
We begin by proving several preliminary results. Most of these lemmas were first proven in a somewhat different form in . We give brief proofs here to make the exposition complete. Lemma 3.8 and Proposition 3.7 are interesting in themselves. They show, with the help of the BCDH formula, that $`\mathrm{log}`$ has a derivative with values in $`𝔤_{\mathrm{}}`$. The standard integral formula for the logarithm $`\mathrm{log}(I+x)=_0^1x(I+sx)^1𝑑s`$, although more useful in many contexts, does not easily yield values in $`𝔤_{\mathrm{}}`$.
###### Notation 3.6.
Let $`gC_{CM}^1`$ be such that $`g(t)I<1`$ for any $`t[0,1]`$. Define the logarithm of $`g`$ by
$$h(t)=\mathrm{log}g(t)=\underset{n=1}{\overset{\mathrm{}}{}}\frac{(1)^{n1}}{n}(g(t)I)^n.$$
###### Proposition 3.7.
\[The derivative of $`log`$ as a series\] Let $`A(t)=g(t)^1\dot{g}(t)`$, and for any $`x𝔤_{\mathrm{}}`$ define
(3.1)
$$F(x,t)=A(t)+\frac{1}{2}[x,A(t)]\underset{p=1}{\overset{\mathrm{}}{}}\frac{1}{2(p+2)p!}[\mathrm{}[[x,A(t)],\stackrel{p}{\stackrel{}{x],\mathrm{},x}}].$$
Then the series converges in $`𝔤_{\mathrm{}}`$.
###### Proof.
Indeed, by Assumption 3.1
$$\begin{array}{c}|F(x,t)||A(t)|+\frac{C|x||A(t)|}{2}+\underset{p=1}{\overset{\mathrm{}}{}}\frac{1}{2(p+2)p!}|[\mathrm{}[[x,A(t)],\stackrel{p}{\stackrel{}{x],\mathrm{},x}}]|\hfill \\ \hfill |A(t)|+\frac{C|x||A(t)|}{2}+|A(t)|\underset{p=1}{\overset{\mathrm{}}{}}\frac{C^{p+1}|x|^{p+1}}{2(p+2)p!}<\mathrm{},\end{array}$$
where $`C`$ is the same constant as in Assumption 3.1. In particular, this means that $`F(x,t)𝔤_{\mathrm{}}`$ for any $`x𝔤_{\mathrm{}}`$ and any $`0t1`$. ∎
###### Lemma 3.8.
\[The derivative of $`log`$\] Let $`gC_{CM}^1`$, $`h=\mathrm{log}g`$, then
(3.2)
$$\dot{h}(t)=F(h,t),$$
where $`F(x,t)`$ is defined by Equation (3.1).
###### Proof.
Indeed, $`h(t+s)=\mathrm{log}g(t+s)=\mathrm{log}(g(t)g(t)^1g(t+s))`$. Let
$$f(t,s)=\mathrm{log}(g(t)^1g(t+s)).$$
Then $`h(t+s)=\mathrm{log}(e^{h(t)}e^{f(t,s)})=BCDH(h(t),f(t,s))`$, where $`BCDH(x,y)`$ is given by the Baker-Campbell-Dynkin-Hausdorff formula for $`x,y𝔤`$
(3.3)
$$\begin{array}{c}BCDH(x,y)=\mathrm{log}(\mathrm{exp}x\mathrm{exp}y)=\hfill \\ \hfill \underset{m}{}\underset{p_i,q_i}{}\frac{(1)^{m1}}{m(\underset{i}{}p_i+q_i)}\frac{[\mathrm{}\stackrel{p_1}{\stackrel{}{[x,x],\mathrm{},x}}],\stackrel{q_1}{\stackrel{}{y],\mathrm{},y}}],\stackrel{q_m}{\stackrel{}{\mathrm{},y],\mathrm{}y}}]}{p_1!q_1!\mathrm{}p_m!q_m!}.\end{array}$$
Note that
$$\begin{array}{c}\frac{df(t,s)}{ds}=\underset{\epsilon 0}{lim}\frac{f(t,s+\epsilon )f(t,s)}{\epsilon }=\hfill \\ \hfill \underset{n=1}{\overset{\mathrm{}}{}}\frac{(1)^{n1}}{n}\underset{k=1}{\overset{n}{}}(g(t)^1g(t+s)I)^{k1}g(t)^1\dot{g}(t+s)\left(g(t)^1g(t+s)I\right)^{nk},\end{array}$$
and therefore
$$\frac{df(t,s)}{ds}|{}_{s=0}{}^{}=g(t)^1\dot{g}=A(t).$$
Note that $`f(t,0)=0`$, and therefore
$$\begin{array}{c}\dot{h}(t)=\underset{s0}{lim}\frac{h(t+s)h(t)}{s}=\underset{s0}{lim}\frac{BCDH(h(t),f(t,s))h(t)}{s}=\hfill \\ \hfill A(t)+\frac{1}{2}[h(t),A(t)]+\frac{d}{ds}|{}_{s=0}{}^{}\underset{p}{}\frac{(1)^1}{2(p+2)}\frac{[\mathrm{}[[h(t)f(s,t)],\stackrel{p}{\stackrel{}{h(t)],\mathrm{},h(t)}}]}{1!1!p!0!}=\\ \hfill A(t)+\frac{1}{2}[h(t),A(t)]\underset{p=1}{\overset{\mathrm{}}{}}\frac{1}{2(p+2)p!}[\mathrm{}[[h(t),A(t)],\stackrel{p}{\stackrel{}{h(t)],\mathrm{},h(t)}}]\end{array}$$
Denote by $`d_n`$ the distance metric on $`G_n`$ corresponding to the norm $`||`$ restricted to the Lie algebra $`𝔤_n`$, and define the metric $`d_{\mathrm{}}=inf_nd_n`$. As before let $`G_{\mathrm{}}=\underset{n=1}{\overset{\mathrm{}}{}}G_n`$. Then $`G_{\mathrm{}}`$ is a group contained in the Cameron-Martin group $`G_{CM}`$. Moreover, for any $`x,yG_{\mathrm{}}`$ we have $`d(x,y)d_{\mathrm{}}(x,y)`$.
###### Lemma 3.9.
Let $`0<L<\mathrm{ln}2/2C`$, where $`C`$ is the same constant as in Assumption 3.1. Then there is a positive constant $`M`$ such that
$`\left|d(I,e^x)|x|\right|M|x|^2,`$
$`\left|d_{\mathrm{}}(I,e^x)|x|\right|M|x|^2.`$
for any $`x𝔤`$ provided $`|x|<L`$.
###### Proof.
First of all, the proof is the same for the metrics $`d`$ and $`d_n`$ with the same constants. We will show how to prove the estimate for the metric $`d`$. Joining $`I`$ to $`e^x`$ by the path $`se^{sx}`$, $`0s1`$, we see that
$$d(I,e^x)|x|_0^1|e^{sx}\dot{e^{sx}}|𝑑s|x|=0.$$
Now let $`g(s)`$ and $`h(s)=\mathrm{log}g(s)`$ be as in Notation 3.6. As before, let $`A(s)=g^1(s)\dot{g}(s)`$, and therefore
$$A(s)=g^1(s)\dot{g}(s)=\underset{k=0}{\overset{\mathrm{}}{}}\frac{(adh)^k(\dot{h})}{(k+1)!}=\dot{h}[h,\dot{h}]+\frac{[h,[h,\dot{h}]]}{2!}\frac{[h,[h,[h,\dot{h}]]]}{3!}+\mathrm{}$$
Then by Remark 3.3 for any smooth path $`g(s)`$ from $`I`$ to $`e^x`$
$$\begin{array}{c}0|x|d(I,e^x)=\left|_0^1\dot{h}(s)𝑑s\right|_0^1\left|g^1(s)\dot{g}(s)\right|𝑑s\hfill \\ \hfill \left|_0^1\dot{h}(s)𝑑s\right|\left|_0^1g^1(s)\dot{g}(s)𝑑s\right|\left|_0^1\dot{h}(s)g^1(s)\dot{g}(s)ds\right|=\\ \hfill \left|_0^1\underset{k=1}{\overset{\mathrm{}}{}}\frac{(adh)^k(\dot{h})}{(k+1)!}ds\right|\underset{k=1}{\overset{\mathrm{}}{}}\frac{C^k}{(k+1)!}_0^1|h(s)|^k|\dot{h}(s)|𝑑s,\end{array}$$
where the constant $`C`$ is the same as in Assumption 3.1. In particular, if $`h(s)=sx,0s1`$, we have
$$\left|d(I,e^x)|x|\right|\underset{k=1}{\overset{\mathrm{}}{}}\frac{C^k|x|^{k+1}}{(k+1)!}_0^1s^k𝑑s|x|^2\underset{k=1}{\overset{\mathrm{}}{}}\frac{C^kL^{k1}}{(k+1)!(k+1)}.$$
Thus we can take
$$M=\underset{k=1}{\overset{\mathrm{}}{}}\frac{C^kL^{k1}}{(k+1)!(k+1)}.$$
Note that the above proof would go through for $`d_n`$, where $`n`$ is such that $`x𝔤_n`$. Thus the statement of Lemma 3.10 holds for $`d`$ replaced by $`d_{\mathrm{}}`$ with the same constant $`M`$. ∎
###### Corollary 3.10.
By choosing $`h(s)=sy+(1s)x`$ for $`0s1`$ in the above proof we have that
$`\left|d(e^x,e^y)|xy|\right|M|xy|^2,`$
$`\left|d_{\mathrm{}}(e^x,e^y)|xy|\right|M|xy|^2.`$
for any $`x,y𝔤`$ provided $`|x|<L`$ and $`|y|<L`$.
###### Lemma 3.11.
Let $`gC_{CM}^1`$. Suppose that
$$g(0)=I,g(t)I<1\text{ for any }t[0,1]\text{ and }|\mathrm{log}g(1)|<L,$$
where $`L`$ is the same as in Corollary 3.10. Then there is a sequence $`g_nG_n`$ such that $`\underset{n\mathrm{}}{lim}d(g_n,g(1))=0`$.
###### Proof.
As before let $`h(s)=\mathrm{log}g(s)`$. Note that by Proposition 3.7 $`h(t)𝔤_{\mathrm{}}`$. Therefore there are $`h_n𝔤_n`$ such that $`|h(1)h_n|\underset{n\mathrm{}}{\overset{}{}}0`$. Let $`g_n=e^{h_n}`$. By Corollary 3.10 we have
$$d(e^{h(1)},g_n)C_2|h(1)h_n|\underset{n\mathrm{}}{\overset{}{}}0.$$
By a direct calculation we have $`e^{h(t)}=g(t)`$. ∎
###### Proposition 3.12.
Suppose Assumption 3.1 holds, then
$$d(x,y)=d_{\mathrm{}}(x,y)$$
for any $`x,yG_{\mathrm{}}`$.
###### Proof.
Both metrics $`d`$ and $`d_{\mathrm{}}`$ are left-invariant on $`G_{\mathrm{}}`$, therefore we can assume that $`x=I`$. As we pointed out earlier, for any $`yG_{\mathrm{}}`$ we have $`d(I,y)d_{\mathrm{}}(I,y)`$. Therefore we only need to prove that $`d_{\mathrm{}}(I,y)d(I,y)`$ for any $`yG_{\mathrm{}}`$.
Denote $`d(I,y)=D`$. For any $`\epsilon ^{}>0`$ there is a path $`g(s)C_{CM}^1`$ such that $`g(0)=I,g(1)=y`$ and
$$D_0^1|g^1(s)\dot{g}(s)|𝑑sD+\epsilon ^{}.$$
Denote $`\epsilon =_0^1|g^1(s)\dot{g}(s)|𝑑sD`$. Let $`m>1,m`$. There exist $`t_k,0=t_0<t_1<\mathrm{}t_{m1}<t_m=1`$, $`k=0,1,\mathrm{},m`$ such that
$$_{t_k}^{t_{k+1}}|g^1(s)\dot{g}(s)|𝑑s=\frac{D+\epsilon }{m}\text{ and }|\mathrm{log}y_{k+1}^1y_k|L,$$
where $`y_k=g(t_k)`$ and $`L<\mathrm{min}\{\mathrm{ln}2/2C,1/M\}`$. Note that in general $`y_k`$ is in the Cameron-Martin group $`G_{CM}`$, not in the group $`G_{\mathrm{}}`$. By Lemma 3.11 there exist $`x_kG_{\mathrm{}}`$ such that $`d(y_k,x_k)<\epsilon /m`$. By applying Corollary 3.10 twice we have that
$$\begin{array}{c}d_{\mathrm{}}(x_k,x_{k+1})|\mathrm{log}x_{k+1}^1x_k|+M|\mathrm{log}x_{k+1}^1x_k|^2\hfill \\ \hfill d(x_k,x_{k+1})+2M|\mathrm{log}x_{k+1}^1x_k|^2d(x_k,x_{k+1})+2M|\mathrm{log}y_{k+1}^1y_k|^2.\end{array}$$
Corollary 3.10 also implies that
$$|\mathrm{log}y_{k+1}^1y_k|M|\mathrm{log}y_{k+1}^1y_k|^2d(y_k,y_{k+1})\frac{D+\epsilon }{m},$$
and since $`L<1/M`$
$$|\mathrm{log}y_{k+1}^1y_k|\frac{D+\epsilon }{m(1ML)}.$$
Thus
$$\begin{array}{c}d_{\mathrm{}}(x_k,x_{k+1})d(x_k,x_{k+1})+2M\left(\frac{D+\epsilon }{m}\right)^2\hfill \\ \hfill d(x_k,y_k)+d(y_k,y_{k+1})+d(y_{k+1},x_{k+1})+2M\left(\frac{D+\epsilon }{m}\right)^2\\ \hfill d(y_k,y_{k+1})+\frac{2\epsilon }{m}+2M\left(\frac{D+\epsilon }{m}\right)^2\frac{D+\epsilon }{m}+\frac{2\epsilon }{m}+2M\left(\frac{D+\epsilon }{m}\right)^2.\end{array}$$
Finally,
$$d_{\mathrm{}}(I,y)m\left(\frac{D+\epsilon }{m}+\frac{2\epsilon }{m}+2M\left(\frac{D+\epsilon }{m}\right)^2\right)=D+3\epsilon +\frac{2M(D+\epsilon )^2}{m}.$$
Recall that $`D=d(I,y)`$, and $`\epsilon `$ and $`m`$ are arbitrary. ∎
Theorem 3.4 now is a direct consequence of Lemma 3.8 and Lemma 3.10 which imply that the exponential and logarithmic functions are well-defined and differentiable in neighborhoods of the identity and zero respectively.
###### Proof of Theorem 3.5.
First of all, $`G_{\mathrm{}}G_{CM}`$ since $`G_nG_{CM}`$ for all $`n`$. Let $`gG_{\mathrm{}}`$, $`kG_{CM}`$, and suppose we have a path $`g(t):[0,1]G_{CM},g(t)C_{CM}^1,g(0)=g,g(1)=k`$. Without a loss of generality we can assume that $`g(t)g<1`$, $`d(\mathrm{log}k,\mathrm{log}g)<L`$, where $`L`$ is the same as in Lemma 3.10. Otherwise the path $`\mathrm{log}g(t)`$ can be divided into a finite number of subpaths satisfying the condition. Lemma 3.11 implies that for any $`\epsilon >0`$ there is $`mG_{\mathrm{}}`$ such that $`d(m,g^1k)<\epsilon `$. Then
$$d(m,g^1k)=d(gm,k)<\epsilon .$$
Note that $`gmG_{\mathrm{}}`$, and therefore we have shown that elements of $`G_{CM}`$ can be approximated by elements of $`G_{\mathrm{}}`$. ∎
## 4. Is $`𝔤_{\mathrm{}}`$ a Lie algebra?
As in the previous section, let $`H`$ be a separable Hilbert space, $`B(H)`$ be the space of bounded operators on $`H`$, and $`I`$ be the identity operator. Here we restrict ourselves to the case of the Hilbert-Schmidt groups. By $`HS`$ we denote the space of Hilbert-Schmidt operators on $`H`$. The space $`HS`$ is equipped with the Hilbert-Schmidt inner product $`(,)_{HS}`$. Let $`GL(H)`$ be the group of invertible elements of $`B(H)`$. Then a Hilbert-Schmidt group is a closed subgroup of $`GL(H)`$ such that $`AIHS`$ for any $`AG`$. Note that the set $`\{AGL(H):AIHS\}`$ is a group.
In Section 3 we considered a sequence of finite-dimensional groups $`G_n`$. Suppose now that in addition $`G_nI+HS`$. Then their Lie algebras satisfy $`𝔤_nHS`$. As before we assume that the infinite-dimensional Lie algebra $`𝔤=\underset{n=1}{\overset{\mathrm{}}{}}𝔤_n`$ is equipped with a Hermitian inner product $`(,)`$, and the corresponding norm is denoted by $`||`$. Throughout this section we assume the following modified version of Assumption 2.1.
###### Assumption 4.1 (Hilbert-Schmidt Completeness Assumption).
There is a subspace $`𝔤_{\mathrm{}}`$ of $`HS`$ such that the Lie algebra $`𝔤`$ is contained in $`𝔤_{\mathrm{}}`$ and the given inner product $`(,)`$ on $`𝔤`$ extends to $`𝔤_{\mathrm{}}`$, which is complete with respect to this inner product. As before we will abuse notation by using $`(,)`$ to denote the extended inner product on $`𝔤_{\mathrm{}}`$ and by $`||`$ the corresponding norm. We assume that $`𝔤`$ is dense in $`𝔤_{\mathrm{}}`$ in the norm $`||`$.
###### Remark 4.2.
In our earlier papers we assumed that $`(,)_𝔤_{\mathrm{}}`$ is given by $`(x,y)_𝔤_{\mathrm{}}=(x,Q^1y)_{HS}`$, where $`Q`$ is a one-to-one nonnegative trace class operator on $`HS`$ for which each $`𝔤_n`$ is an invariant subspace. The assumption that $`Q`$ is trace-class assures that the heat kernel measure constructed in our previous work (, , , ) actually lives in $`HS+I`$. In the present paper we do not assume that $`Q`$ is trace-class unless it is stated explicitly. Moreover, we do not use the operator $`Q`$, but rather describe the assumptions on $`Q`$ in terms of an orthonormal basis of $`𝔤_{\mathrm{}}`$.
In the next statement we use the fact that we can view an element of $`HS`$ as an infinite matrix $`A=\{a_{ij}\}_{i,j=1}^{\mathrm{}}`$ such that the sum $`_{i,j}|a_{ij}|^2`$ is finite. Then $`e_{ij}`$, the matrices with $`1`$ at the $`ij`$th place and $`0`$ at all other places, form an orthonormal basis of $`HS`$ with the inner product $`(,)_{HS}`$. Let us describe an example of the setting introduced above. Namely, let $`𝔤_{\mathrm{}}`$ be the vector space generated by the orthonormal basis $`\xi _{ij}=\lambda _{ij}e_{ij},(i,j)A\times `$ for some $`\lambda _{ij}>0`$. Then the inner product on $`𝔤_{\mathrm{}}`$ is determined by $`(\xi _{ij},\xi _{km})_𝔤_{\mathrm{}}=(e_{ij},e_{km})_{HS}`$. It turns out that $`𝔤_{\mathrm{}}`$ might not be a Lie algebra.
###### Proposition 4.3.
There exists a sequence of positive numbers $`\lambda _{ij}`$ such that $`𝔤_{\mathrm{}}`$ is not a Lie algebra.
###### Proof.
Let $`\lambda _{ij}=\lambda _{ji}`$ for any $`i,j`$, and $`x,y𝔤_{\mathrm{}}`$ be such that
$$x=\underset{l=0}{\overset{\mathrm{}}{}}x_l(\xi _{3l+1,3l+2}\xi _{3l+2,3l+1}),y=\underset{l=0}{\overset{\mathrm{}}{}}y_l(\xi _{3l+1,3l+3}\xi _{3l+3,3l+1}).$$
Then
$`[x,y]`$ $`={\displaystyle \underset{l,n}{}}x_ly_n[\xi _{3l+1,3l+2}\xi _{3l+2,3l+1},\xi _{3n+1,3n+3}\xi _{3n+3,3n+1}]`$
$`={\displaystyle \underset{l}{}}x_ly_l{\displaystyle \frac{\lambda _{3l+1,3l+2}\lambda _{3l+1,3l+3}}{\lambda _{3l+2,3l+3}}}(\xi _{3l+2,3l+3}\xi _{3l+3,3l+2}).`$
Thus
$$|[x,y]|^2=\underset{l}{}x_l^2y_l^2\frac{\lambda _{3l+1,3l+2}^2\lambda _{3l+1,3l+3}^2}{\lambda _{3l+2,3l+3}^2}.$$
Let $`x_l=y_l=\lambda _{3l+1,3l+2}=\lambda _{3l+1,3l+3}=a_l`$ and $`\lambda _{3l+2,3l+3}=b_l`$, where $`a_l`$ and $`b_l`$ are $`\mathrm{}^2`$-sequences. Then $`|[x,y]|^2=_la_l^8b_l^2=\mathrm{}`$ if, for example, $`a_l=1/l`$ and $`b_l=1/l^4`$.
The next result shows that there exist inner products such that Continuity Assumption 3.1 on the Lie bracket is satisfied.
###### Theorem 4.4.
Suppose $`\lambda _{i,j}=\lambda _i\lambda _j`$, $`i,j`$, then for any $`x,y𝔤_{\mathrm{}}`$
$$|[x,y]|2\underset{i}{sup}\lambda _i^2|x||y|.$$
###### Proof.
Let $`x=_{i,j=1}^{\mathrm{}}x_{i,j}\xi _{ij}`$ and $`y=_{k,m=1}^{\mathrm{}}y_{k,m}\xi _{km}`$. Then
$$xy=\underset{i,j,m}{}x_{i,j}y_{j,m}\frac{\lambda _{i,j}\lambda _{j,m}}{\lambda _{i,m}}\xi _{im}$$
and therefore
$`|xy|^2`$ $`={\displaystyle \underset{i,m}{}}\left({\displaystyle \underset{j}{}}x_{i,j}y_{j,m}{\displaystyle \frac{\lambda _{i,j}\lambda _{j,m}}{\lambda _{i,m}}}\right)^2`$
$`={\displaystyle \underset{i,m}{}}\left({\displaystyle \underset{j}{}}x_{i,j}y_{j,m}\lambda _j^2\right)^2{\displaystyle \underset{i,j}{}}x_{i,j}^2\lambda _j^2{\displaystyle \underset{k,m}{}}y_{k,m}^2\lambda _k^2.`$
###### Corollary 4.5.
If $`\lambda _{i,j}=\lambda _i\lambda _j`$, $`i,j`$, and $`\{\lambda _i\}_{i=1}^{\mathrm{}}\mathrm{}^p`$ for any $`p>0`$, then Assumption 3.1 on the Lie bracket is satisfied. In particular, if $`p=2`$, then the operator $`Q`$ mentioned in Remark 4.2 is trace-class.
## 5. Riemannian geometry of the Hilbert-Schmidt groups: definitions and preliminaries
The goal of the next two sections is to see if there exists a natural Lie algebra for $`G_{CM}`$ and an inner product on it such that the Ricci curvature is bounded from below. The first obstacle in answering such a question is the absence of geometric definitions. We chose to follow the work of J. Milnor for finite-dimensional Lie groups in . There he described the Riemannian geometry of a Lie group with a Riemannian metric invariant under left translation. One of his aims was to see how the choice of an orthonormal basis of the (finite-dimensional) Lie algebra determines the curvature properties of the corresponding Lie group. This is exactly the question we study, but in infinite dimensions: how the choice of the inner product on $`𝔤`$ changes the Riemannian geometry of the group $`G_{CM}`$. We consider general norms on $`𝔤`$ which are diagonal in a certain sense. This allows us to compute the Ricci curvature in two important cases: the first case is when the norm on $`𝔤`$ is the Hilbert-Schmidt norm, and the second one is when the norm on $`𝔤`$ is determined by a nonnegative trace class operator on $`HS`$. The latter assumption assures that the corresponding heat kernel measure constructed in our previous work (, , , ) actually lives in $`HS+I`$. We use finite-dimensional approximations to $`𝔤`$ to define the sectional and Ricci curvatures. Our results show that for the general, orthogonal and upper triangular Hilbert-Schmidt algebras the Ricci curvature generally is not bounded from below. Moreover, for the upper triangular Hilbert-Schmidt algebra the Ricci curvature is identically minus infinity.
Let $`𝔤`$ be a infinite-dimensional Lie algebra equipped with an inner product $`(,)`$. We assume that $`𝔤`$ is complete. By Theorem 4.4 this is the case for all examples we consider later in this Section.
###### Definition 5.1.
The Levi-Civita connection $`_x`$ is defined by
$$(_xy,z)=\frac{1}{2}(([x,y],z)([y,z],x)+([z,x],y))$$
for any $`x,y,z𝔤`$.
###### Definition 5.2.
1. The Riemannian curvature tensor $`R`$ is defined by
$$R_{xy}=_{[x,y]}_x_y+_y_x,x,y𝔤.$$
2. For any orthogonal $`x,y`$ in $`𝔤`$
$$K(x,y)=(R_{xy}(x),y)$$
is called the sectional curvature.
3. Let $`\{\xi _i\}_{i=1}^{\mathrm{}}`$ be an orthonormal basis of $`𝔤`$, $`N`$ be finite, then
$$R^N(x)=\underset{i=1}{\overset{N}{}}K(x,\xi _i)=\underset{i=1}{\overset{N}{}}(R_{x\xi _i}(x),\xi _i)$$
is the truncated Ricci curvature.
4. Let $`N`$ be finite, then
$$\widehat{R}^N(x)=\underset{i=1}{\overset{N}{}}R_{\xi _ix}(\xi _i)$$
is the truncated self-adjoint Ricci curvature or transformation.
The self-adjoint Ricci transformation is a convenient computational tool. First of all,
$$R^N(x)=(\widehat{R}^N(x),x).$$
Then if $`\{\xi _i\}_{i=1}^{dim𝔤}`$ is an orthonormal basis which diagonalizes $`\widehat{R}`$, that is, $`\widehat{R}^N(\xi _i)=a_i\xi _i`$, then
$$R^N(x)=\underset{i=1}{\overset{N}{}}a_ix_i^2,x=\underset{i=1}{\overset{dim𝔤}{}}x_i\xi _i.$$
The numbers $`a_i`$ are called the principal Ricci curvatures.
As in Section 4, let $`\{e_{ij}\}_{i,j=1}^{\mathrm{}}`$ be the standard basis of the space of Hilbert-Schmidt operators $`HS`$. The Lie bracket for these basis elements can be written as
$$[e_{ij},e_{km}]=\delta _{j,k}e_{im}\delta _{i,m}e_{kj}$$
where $`\delta _{ln}`$ is Kronecker’s symbol. As before, we study a subspace $`𝔤_{\mathrm{}}`$ of $`HS`$ generated by an orthonormal basis $`\xi _{ij}=\lambda _i\lambda _je_{ij}`$ for some $`\lambda _i>0`$, $`(i,j)A\times `$.
## 6. Ricci curvature
The main results of this section are Theorems 6.1, 6.8 and 6.10. For the skew-symmetric and triangular infinite matrices, the orthonormal basis we use actually diagonalizes the truncated self-adjoint Ricci curvature. Then the results of this section show that if we define the Ricci curvature as the limit of the truncated Ricci curvature as the dimension goes to $`\mathrm{}`$, it is not bounded from below. Moreover, for the upper triangular matrices the Ricci curvature is identically negative infinity.
### 6.1. General Hilbert-Schmidt algebra
Let $`\{\lambda _i\}_{i=1}^{\mathrm{}}`$ be a bounded sequence of strictly positive numbers. In this section we consider the infinite dimensional Lie algebra $`𝔤_{\mathrm{}}`$ generated by the orthonormal basis $`\xi _{ij}=\lambda _i\lambda _je_{ij}`$. Then $`𝔤_{\mathrm{}}`$ is a Lie subalgebra of $`𝔤𝔩_{HS}`$. Recall that if the sequence $`\{\lambda _i\}_{i=1}^{\mathrm{}}`$ is bounded, then by Theorem 4.4, Continuity Assumption 3.1 is satisfied for the corresponding norm on $`𝔤_{\mathrm{}}`$.
###### Theorem 6.1.
1. Let $`N>\mathrm{max}i,j`$, then the truncated Ricci curvature is
$$R_{ij}^N=R^N(\xi _{ij})=\frac{1}{4}(6\delta _{i,j}\lambda _i^44\delta _{i,j}\lambda _i^4N2\lambda _i^4N2\lambda _j^4N+2\underset{m=1}{\overset{N}{}}\lambda _m^4).$$
2. Suppose that $`\{\lambda _i\}_{i=1}^{\mathrm{}}\mathrm{}^2`$. Then
$$\underset{N\mathrm{}}{lim}R^N(\xi _{ij})=\frac{1}{2}\underset{N\mathrm{}}{lim}(\lambda _i^4N\lambda _j^4N+\underset{m=1}{\overset{N}{}}\lambda _m^4)=\mathrm{},$$
if $`ij`$.
$$\underset{N\mathrm{}}{lim}R^N(\xi _{ii})=\frac{1}{2}\underset{N\mathrm{}}{lim}(3\lambda _i^44\lambda _i^4N+\underset{m=1}{\overset{N}{}}\lambda _m^4)=\mathrm{}.$$
3. For the Hilbert-Schmidt inner product, the truncated Ricci curvature is
$$R_{ij}^N=\frac{1}{2}(3\delta _{i,j}2\delta _{i,j}NN).$$
This theorem is a direct consequence of the following results. First of all, note that the Lie bracket can be written
$$[\xi _{ij},\xi _{km}]=\delta _{j,k}\lambda _j^2\xi _{im}\delta _{i,m}\lambda _i^2\xi _{kj},$$
where $`\delta _{l,n}`$ is Kronecker’s symbol. Denote $`_{ab}=_{\xi _{ab}}`$, then
###### Lemma 6.2.
$$_{ab}\xi _{cd}=\frac{1}{2}(\delta _{b,c}\lambda _b^2\xi _{ad}\delta _{a,d}\lambda _a^2\xi _{cb}\delta _{a,c}\lambda _d^2\xi _{db}+\delta _{b,d}\lambda _c^2\xi _{ac}+\delta _{b,d}\lambda _a^2\xi _{ca}\delta _{a,c}\lambda _b^2\xi _{bd})$$
###### Proof.
$$\begin{array}{c}(_{ab}\xi _{cd},\xi _{ef})=\frac{1}{2}\left(([\xi _{ab},\xi _{cd}],\xi _{ef})([\xi _{cd},\xi _{ef}],\xi _{ab})+([\xi _{ef},\xi _{ab}],\xi _{cd})\right)=\hfill \\ \hfill \frac{1}{2}(\delta _{b,c}\lambda _b^2(\xi _{ad},\xi _{ef})\delta _{a,d}\lambda _a^2(\xi _{cb},\xi _{ef})\delta _{d,e}\lambda _d^2(\xi _{cf},\xi _{ab})\\ \hfill +\delta _{c,f}\lambda _c^2(\xi _{ed},\xi _{ab})+\delta _{f,a}\lambda _f^2(\xi _{eb},\xi _{cd})\delta _{e,b}\lambda _e^2(\xi _{af},\xi _{cd}))=\\ \hfill \frac{1}{2}(\delta _{b,c}\delta _{e,a}\delta _{f,d}\lambda _b^2\delta _{a,d}\delta _{e,c}\delta _{f,b}\lambda _a^2\delta _{a,c}\delta _{e,d}\delta _{f,b}\lambda _d^2\\ \hfill +\delta _{b,d}\delta _{e,a}\delta _{f,c}\lambda _c^2+\delta _{b,d}\delta _{e,c}\delta _{f,a}\lambda _a^2\delta _{a,c}\delta _{f,d}\delta _{e,b}\lambda _b^2).\end{array}$$
###### Lemma 6.3.
$$\begin{array}{c}R_{ij,km}\xi _{ij}=\hfill \\ \hfill \frac{1}{4}(3\delta _{j,k}\delta _{j,m}\lambda _i^2\lambda _j^2\xi _{ii}+2\delta _{i,k}\delta _{j,m}(\lambda _i^4+\lambda _j^4)\xi _{ij}+2\delta _{i,m}\delta _{j,k}\lambda _i^2\lambda _j^2\xi _{ij}\\ \hfill +2\delta _{i,k}\delta _{j,m}\lambda _i^2\lambda _j^2\xi _{ji}+3\delta _{i,k}\delta _{i,m}\lambda _i^2\lambda _j^2\xi _{jj}\\ \hfill \delta _{i,m}\lambda _i^2\lambda _k^2\xi _{ik}2\delta _{i,j}\delta _{i,m}\lambda _i^2\lambda _k^2\xi _{ik}+\delta _{i,m}\lambda _i^4\xi _{ki}2\delta _{i,j}\delta _{i,m}\lambda _i^4\xi _{ki}\\ \hfill 4\delta _{i,k}\lambda _i^4\xi _{im}+\delta _{i,k}\lambda _j^4\xi _{im}+\delta _{i,k}\lambda _m^4\xi _{im}2\delta _{i,j}\delta _{i,k}\lambda _i^4\xi _{im}\delta _{i,k}\lambda _i^2\lambda _m^2\xi _{mi}2\delta _{i,j}\delta _{i,k}\lambda _i^2\lambda _m^2\xi _{mi}\\ \hfill \delta _{j,m}\lambda _j^2\lambda _k^2\xi _{jk}+\delta _{j,m}\lambda _i^4\xi _{kj}+\delta _{j,m}\lambda _k^4\xi _{kj}4\delta _{j,m}\lambda _j^4\xi _{kj}+\delta _{j,k}\lambda _j^4\xi _{jm}\delta _{j,k}\lambda _j^2\lambda _m^2\xi _{mj})\end{array}$$
###### Proof.
The Riemannian curvature tensor applied to $`\xi _{ij}`$ is
$$\begin{array}{c}R_{\xi _{ij}\xi _{km}}\xi _{ij}=R_{ij,km}\xi _{ij}=_{[\xi _{ij},\xi _{km}]}\xi _{ij}_{ij}_{km}\xi _{ij}+_{km}_{ij}\xi _{ij}=\hfill \\ \hfill \delta _{j,k}\lambda _j^2_{im}\xi _{ij}\delta _{i,m}\lambda _i^2_{kj}\xi _{ij}_{ij}_{km}\xi _{ij}+_{km}_{ij}\xi _{ij}=\\ \hfill \frac{1}{2}\delta _{j,k}\lambda _j^2(\delta _{i,m}\lambda _i^2\xi _{ij}\delta _{i,j}\lambda _i^2\xi _{im}\lambda _j^2\xi _{jm}\lambda _m^2\xi _{mj}+2\delta _{m,j}\lambda _i^2\xi _{ii})\\ \hfill \frac{1}{2}\delta _{i,m}\lambda _i^2(\lambda _i^2\xi _{ki}+\lambda _k^2\xi _{ik}+\delta _{j,i}\lambda _j^2\xi _{kj}\delta _{k,j}\lambda _k^2\xi _{ij}2\delta _{i,k}\lambda _j^2\xi _{jj})\\ \hfill \frac{1}{2}_{ij}(\delta _{i,m}\lambda _m^2\xi _{kj}\delta _{k,j}\lambda _k^2\xi _{im}\delta _{i,k}\lambda _j^2\xi _{jm}+\delta _{m,j}\lambda _i^2\xi _{ki}+\delta _{m,j}\lambda _k^2\xi _{ik}\delta _{i,k}\lambda _m^2\xi _{mj})\\ \hfill +\frac{1}{2}_{km}(\lambda _i^2\xi _{ii}+\lambda _i^2\xi _{ii}+\delta _{i,j}\lambda _j^2\xi _{ij}\delta _{i,j}\lambda _i^2\xi _{ij}2\lambda _j^2\xi _{jj})\end{array}$$
###### Lemma 6.4.
The sectional curvature is
$$\begin{array}{c}K(\xi _{ij},\xi _{km})=(R_{ij,km}\xi _{ij},\xi _{km})=\hfill \\ \hfill \frac{1}{4}(6\delta _{i,k}\delta _{i,m}\delta _{j,k}\delta _{j,m}\lambda _i^4+2\delta _{i,k}\delta _{j,m}\lambda _i^4+2\delta _{i,k}\delta _{j,m}\lambda _j^42\delta _{i,k}\delta _{i,m}\delta _{k,m}\lambda _i^4\\ \hfill +\delta _{i,m}\lambda _i^42\delta _{i,j}\delta _{i,m}\lambda _i^44\delta _{i,k}\lambda _i^4+\delta _{i,k}\lambda _j^4+\delta _{i,k}\lambda _m^42\delta _{i,j}\delta _{i,k}\lambda _i^42\delta _{j,k}\delta _{j,m}\delta _{m,k}\lambda _j^4\\ \hfill +\delta _{j,m}\lambda _i^4+\delta _{j,m}\lambda _k^44\delta _{j,m}\lambda _j^4+\delta _{j,k}\lambda _j^4).\end{array}$$
###### Proof.
$$\begin{array}{c}(R_{ij,km}\xi _{ij},\xi _{km})=\hfill \\ \hfill \frac{1}{4}(3\delta _{i,k}\delta _{i,m}\delta _{j,k}\delta _{j,m}\lambda _i^2\lambda _j^2+2\delta _{i,k}\delta _{j,m}\lambda _i^4+2\delta _{i,k}\delta _{j,m}\lambda _j^4\\ \hfill +2\delta _{i,k}\delta _{i,m}\delta _{j,k}\delta _{j,m}\lambda _i^2\lambda _j^2+2\delta _{i,k}\delta _{i,m}\delta _{j,k}\delta _{j,m}\lambda _i^2\lambda _j^2+3\delta _{i,k}\delta _{i,m}\delta _{j,k}\delta _{j,m}\lambda _i^2\lambda _j^2\\ \hfill \delta _{i,k}\delta _{i,m}\delta _{k,m}\lambda _i^2\lambda _k^22\delta _{i,j}\delta _{i,k}\delta _{i,m}\delta _{k,m}\lambda _i^2\lambda _k^2+\delta _{i,m}\lambda _i^42\delta _{i,j}\delta _{i,m}\lambda _i^4\\ \hfill 4\delta _{i,k}\lambda _i^4+\delta _{i,k}\lambda _j^4+\delta _{i,k}\lambda _m^42\delta _{i,j}\delta _{i,k}\lambda _i^4\delta _{i,k}\delta _{i,m}\delta _{m,k}\lambda _i^2\lambda _m^22\delta _{i,j}\delta _{i,k}\delta _{i,m}\delta _{m,k}\lambda _i^2\lambda _m^2\\ \hfill \delta _{j,k}\delta _{j,m}\delta _{k,m}\lambda _j^2\lambda _k^2+\delta _{j,m}\lambda _i^4+\delta _{j,m}\lambda _k^44\delta _{j,m}\lambda _j^4+\delta _{j,k}\lambda _j^4\delta _{j,k}\delta _{j,m}\delta _{m,k}\lambda _j^2\lambda _m^2)=\\ \hfill \frac{1}{4}(6\delta _{i,k}\delta _{i,m}\delta _{j,k}\delta _{j,m}\lambda _i^4+2\delta _{i,k}\delta _{j,m}\lambda _i^4+2\delta _{i,k}\delta _{j,m}\lambda _j^42\delta _{i,k}\delta _{i,m}\delta _{k,m}\lambda _i^4\\ \hfill +\delta _{i,m}\lambda _i^42\delta _{i,j}\delta _{i,m}\lambda _i^44\delta _{i,k}\lambda _i^4+\delta _{i,k}\lambda _j^4+\delta _{i,k}\lambda _m^42\delta _{i,j}\delta _{i,k}\lambda _i^42\delta _{j,k}\delta _{j,m}\delta _{m,k}\lambda _j^4\\ \hfill +\delta _{j,m}\lambda _i^4+\delta _{j,m}\lambda _k^44\delta _{j,m}\lambda _j^4+\delta _{j,k}\lambda _j^4).\end{array}$$
### 6.2. Orthogonal Hilbert-Schmidt algebra
Note that $`b_{ij}=(e_{ij}e_{ji})/\sqrt{2},i<j`$ is an orthonormal basis for the space of skew-symmetric Hilbert-Schmidt operators $`𝔰𝔬_{HS}`$. Let $`\{\lambda _i\}_{i=1}^{\mathrm{}}`$ be a bounded sequence of strictly positive numbers. In this section we consider the infinite dimensional Lie algebra $`𝔤_{\mathrm{}}`$ generated by the orthonormal basis $`\xi _{ij}=\lambda _i\lambda _jb_{ij}`$. Recall that if the sequence $`\{\lambda _i\}_{i=1}^{\mathrm{}}`$ is bounded, then by Theorem 4.4, Continuity Assumption 3.1 is satisfied for the corresponding norm on $`𝔤_{\mathrm{}}`$.
In what follows the convention is that $`\xi _{ij}=0`$, if $`ij`$. The Lie bracket for these basis elements can be written as
$$[\xi _{ij},\xi _{km}]=\frac{1}{\sqrt{2}}(\delta _{j,k}\lambda _j^2\xi _{im}+\delta _{j,m}\lambda _j^2(\xi _{ki}\xi _{ik})+\delta _{i,k}\lambda _i^2(\xi _{mj}\xi _{jm})\delta _{i,m}\lambda _i^2\xi _{kj})$$
We begin with several computational lemmas.
###### Lemma 6.5.
$$\begin{array}{c}_{ab}\xi _{cd}=\hfill \\ \hfill \frac{1}{2\sqrt{2}}(\lambda _b^2\delta _{b,c}\xi _{a,d}+\lambda _b^2\delta _{b,d}(\xi _{c,a}\xi _{a,c})+\lambda _a^2\delta _{a,c}(\xi _{d,b}\xi _{b,d})\lambda _a^2\delta _{a,d}\xi _{c,b}\\ \hfill \lambda _d^2\delta _{a,c}\xi _{d,b}\lambda _d^2(\delta _{b,c}\xi _{a,d}\delta _{a,c}\xi _{b,d})\lambda _c^2(\delta _{b,d}\xi _{c,a}\delta _{a,d}\xi _{c,b})+\lambda _c^2\delta _{b,d}\xi _{a,c}\\ \hfill +\lambda _a^2\delta _{b,d}\xi _{c,a}+\lambda _b^2(\delta _{a,c}\xi _{d,b}\delta _{a,d}\xi _{c,b})+\lambda _a^2(\delta _{b,c}\xi _{a,d}\delta _{b,d}\xi _{a,c})\lambda _b^2\delta _{a,c}\xi _{b,d}).\end{array}$$
###### Proof.
$$\begin{array}{c}(_{ab}\xi _{cd},\xi _{ef})=\hfill \\ \hfill \frac{1}{2\sqrt{2}}(((\delta _{b,c}\lambda _b^2\xi _{ad}+\delta _{b,d}\lambda _b^2(\xi _{ca}\xi _{ac})+\delta _{a,c}\lambda _a^2(\xi _{db}\xi _{bd})\delta _{a,d}\lambda _a^2\xi _{cb}),\xi _{ef})\\ \hfill (\delta _{d,e}\lambda _d^2\xi _{cf}+\delta _{d,f}\lambda _d^2(\xi _{ec}\xi _{ce})+\delta _{c,e}\lambda _c^2(\xi _{fd}\xi _{df})\delta _{c,f}\lambda _c^2\xi _{ed},\xi _{ab})\\ \hfill +(\delta _{f,a}\lambda _f^2\xi _{eb}+\delta _{f,b}\lambda _f^2(\xi _{ae}\xi _{ea})+\delta _{e,a}\lambda _e^2(\xi _{bf}\xi _{fb})\delta _{e,b}\lambda _e^2\xi _{af},\xi _{cd}))=\\ \hfill \frac{1}{2\sqrt{2}}(\lambda _b^2\delta _{a,e}\delta _{b,c}\delta _{d,f}+\lambda _b^2\delta _{b,d}(\delta _{c,e}\delta _{a,f}\delta _{a,e}\delta _{c,f})+\lambda _a^2\delta _{a,c}(\delta _{d,e}\delta _{b,f}\delta _{d,f}\delta _{b,e})\\ \hfill \lambda _a^2\delta _{a,d}\delta _{b,f}\delta _{c,e}\lambda _d^2\delta _{a,c}\delta _{b,f}\delta _{d,e}\lambda _d^2\delta _{d,f}(\delta _{a,e}\delta _{b,c}\delta _{a,c}\delta _{b,e})\lambda _c^2\delta _{c,e}(\delta _{a,f}\delta _{b,d}\delta _{a,d}\delta _{b,f})\\ \hfill +\lambda _c^2\delta _{a,e}\delta _{b,d}\delta _{c,f}+\lambda _a^2\delta _{a,f}\delta _{b,d}\delta _{c,e}+\lambda _b^2\delta _{b,f}(\delta _{a,c}\delta _{e,d}\delta _{e,c}\delta _{a,d})+\lambda _a^2\delta _{a,e}(\delta _{b,c}\delta _{d,f}\delta _{c,f}\delta _{b,d})\\ \hfill \lambda _b^2\delta _{a,c}\delta _{b,e}\delta _{d,f})\end{array}$$
###### Lemma 6.6.
The Riemannian curvature tensor applied to $`\xi _{ij}`$ is
$$\begin{array}{c}R_{ij,km}\xi _{ij}=\frac{1}{8}\delta _{j,m}\lambda _j^2(\lambda _i^23\lambda _j^2+3\lambda _k^2)\xi _{k,j}+\frac{1}{8}\delta _{i,k}\lambda _i^2(3\lambda _i^2+\lambda _j^2+3\lambda _m^2)\xi _{i,m}\hfill \\ \hfill \frac{1}{4}\delta _{i,m}\lambda _i^2(\lambda _j^2+\lambda _i^2\lambda _k^2)\xi _{k,i}+\frac{1}{8}\delta _{i,m}(\lambda _i^2+\lambda _j^2\lambda _k^2)(\lambda _i^2+\lambda _j^2\lambda _k^2)\xi _{k,i}\\ \hfill +\frac{1}{4}\delta _{j,k}\lambda _j^2(\lambda _i^2\lambda _j^2+\lambda _m^2)\xi _{j,m}+\frac{1}{8}\delta _{j,k}(\lambda _i^2+\lambda _j^2+\lambda _m^2)(\lambda _i^2\lambda _j^2+\lambda _m^2)\xi _{j,m}\end{array}$$
###### Proof.
$$\begin{array}{c}R_{\xi _{ij}\xi _{km}}\xi _{ij}=R_{ij,km}\xi _{ij}=\frac{1}{\sqrt{2}}\delta _{j,k}\lambda _j^2_{im}\xi _{ij}+\frac{1}{\sqrt{2}}\delta _{j,m}\lambda _j^2_{(\xi _{ki}\xi _{ik})}\xi _{ij}\hfill \\ \hfill +\frac{1}{\sqrt{2}}\delta _{i,k}\lambda _i^2_{(\xi _{mj}\xi _{jm})}\xi _{ij}\frac{1}{\sqrt{2}}\delta _{i,m}\lambda _i^2_{kj}\xi _{ij}_{ij}_{km}\xi _{ij}=\end{array}$$
$$\begin{array}{c}\frac{1}{4}\delta _{j,k}\lambda _i^2\lambda _j^2\xi _{j,m}\frac{1}{4}\delta _{j,k}\lambda _j^4\xi _{j,m}+\frac{1}{4}\delta _{j,k}\lambda _j^2\lambda _m^2\xi _{j,m}\hfill \\ \hfill +\frac{1}{\sqrt{2}}\delta _{j,m}\lambda _j^2_{(\xi _{ki}\xi _{ik})}\xi _{ij}+\frac{1}{\sqrt{2}}\delta _{i,k}\lambda _i^2_{(\xi _{mj}\xi _{jm})}\xi _{ij}\frac{1}{\sqrt{2}}\delta _{i,m}\lambda _i^2_{kj}\xi _{ij}\\ \hfill \frac{1}{2\sqrt{2}}(+\lambda _m^2\delta _{i,m}_{ij}\xi _{k,j}+\lambda _m^2\delta _{m,j}_{ij}(\xi _{i,k}\xi _{i,k})+\lambda _k^2\delta _{i,k}_{ij}(\xi _{j,m}\xi _{m,j})\\ \hfill \lambda _k^2\delta _{k,j}_{ij}\xi _{i,m}\lambda _j^2\delta _{i,k}_{ij}\xi _{j,m}\\ \hfill \lambda _j^2(\delta _{i,m}_{ij}\xi _{k,j}\delta _{i,k}_{ij}\xi _{m,j})\lambda _i^2(\delta _{m,j}_{ij}\xi _{i,k}\delta _{k,j}_{ij}\xi _{i,m})+\lambda _i^2\delta _{m,j}_{ij}\xi _{i,k}\\ \hfill +\lambda _k^2\delta _{m,j}_{ij}\xi _{i,k}+\lambda _m^2(\delta _{i,k}_{ij}\xi _{j,m}\delta _{k,j}_{ij}\xi _{i,m})\\ \hfill +\lambda _k^2(\delta _{i,m}_{ij}\xi _{k,j}\delta _{m,j}_{ij}\xi _{i,k})\lambda _m^2\delta _{i,k}_{ij}\xi _{m,j}).\end{array}$$
###### Lemma 6.7.
The sectional curvature is
$$\begin{array}{c}K(\xi _{ij},\xi _{km})=(R_{ij,km}\xi _{ij},\xi _{km})=\hfill \\ \hfill \frac{1}{8}\delta _{j,m}\lambda _j^2(\lambda _i^23\lambda _j^2+3\lambda _k^2)+\frac{1}{8}\delta _{i,k}\lambda _i^2(3\lambda _i^2+\lambda _j^2+3\lambda _m^2)\\ \hfill +\frac{1}{4}\delta _{i,m}\lambda _i^2(\lambda _j^2\lambda _i^2+\lambda _k^2)+\frac{1}{8}\delta _{i,m}(\lambda _i^2+\lambda _j^2\lambda _k^2)(\lambda _i^2+\lambda _j^2\lambda _k^2)\\ \hfill +\frac{1}{4}\delta _{j,k}\lambda _j^2(\lambda _i^2\lambda _j^2+\lambda _m^2)+\frac{1}{8}\delta _{j,k}(\lambda _i^2+\lambda _j^2+\lambda _m^2)(\lambda _i^2\lambda _j^2+\lambda _m^2)\end{array}$$
The third part of the following theorem says that the principal Ricci curvatures for $`𝔰𝔬_{HS}`$ can tend to either $`\mathrm{}`$ or $`\mathrm{}`$ as the dimension $`N\mathrm{}`$ depending on the choice of the scaling $`\{\lambda _i\}_{i=1}^{\mathrm{}}`$.
###### Theorem 6.8.
Let $`N>\mathrm{max}i,j`$. Then
1. the truncated Ricci curvature is
$$R_{ij}^N=\frac{1}{8}(4\lambda _i^2\lambda _j^25(\lambda _i^2\lambda _j^2)^2)N+\frac{3}{8}(\lambda _i^2+\lambda _j^2)\underset{m=1}{\overset{m=N}{}}\lambda _m^2+\frac{1}{4}\underset{m=1}{\overset{m=N}{}}\lambda _m^4.$$
2. For the Hilbert-Schmidt inner product the truncated Ricci curvature is
$$R_{ij}^N=\frac{N}{2}.$$
3. The truncated self-adjoint Ricci curvature is diagonal in the basis $`\{\xi _{km}\}`$. Let $`\{a_{km}\}`$ be its principal Ricci curvatures, that is, $`\widehat{R}^N(\xi _{km})=a_{km}\xi _{km}`$. Then if $`\{\lambda _i\}_{i=1}^{\mathrm{}}\mathrm{}^2`$, the principal Ricci curvatures have the following asymptotics as $`N\mathrm{}`$
$$a_{km}=A_{km}^N+\frac{N}{8}(2\lambda _k^23\lambda _m^2)(\lambda _m^2\lambda _k^2),$$
where $`A_{km}^NA<\mathrm{}`$ as $`N\mathrm{}`$.
###### Proof.
(1)
$$\begin{array}{c}R_{ij}^N=\underset{k,m=1}{\overset{k,m=N}{}}(R_{ij,km}\xi _{ij},\xi _{km})=\hfill \\ \hfill \underset{k=1}{\overset{k=N}{}}\frac{1}{8}\lambda _j^2(\lambda _i^23\lambda _j^2+3\lambda _k^2)+\frac{1}{8}\delta _{i,k}\lambda _i^2(3\lambda _i^2N+\lambda _j^2N+3\underset{m=1}{\overset{m=N}{}}\lambda _m^2)\\ \hfill +\frac{1}{4}\lambda _i^2(\lambda _j^2\lambda _i^2+\lambda _k^2)+\frac{1}{8}(\lambda _i^2+\lambda _j^2\lambda _k^2)(\lambda _i^2+\lambda _j^2\lambda _k^2)\\ \hfill +\frac{1}{4}\delta _{j,k}\lambda _j^2(\lambda _i^2N\lambda _j^2N+\underset{m=1}{\overset{m=N}{}}\lambda _m^2)\\ \hfill +\frac{1}{8}\delta _{j,k}((\lambda _i^2\lambda _j^2)(\lambda _i^2+\lambda _j^2)N2\lambda _i^2\underset{m=1}{\overset{m=N}{}}\lambda _m^2+\underset{m=1}{\overset{m=N}{}}\lambda _m^4)=\\ \hfill +\frac{1}{8}(6\lambda _i^2\lambda _j^25\lambda _i^45\lambda _j^4)N+\frac{3}{8}(\lambda _i^2+\lambda _j^2)\underset{m=1}{\overset{m=N}{}}\lambda _m^2+\frac{1}{4}\underset{m=1}{\overset{m=N}{}}\lambda _m^4.\end{array}$$
(2) follows from (1).
(3) Let $`k<m<N`$, then the truncated self-adjoint Ricci curvature is
$$\begin{array}{c}\widehat{R}^N(\xi _{km})=\underset{i<jN}{}R_{ij,km}\xi _{ij}=\underset{j=1}{\overset{N}{}}\underset{i=1}{\overset{j1}{}}R_{ij,km}\xi _{ij}=\hfill \\ \hfill \underset{j=1}{\overset{k}{}}\underset{i=1}{\overset{j1}{}}R_{ij,km}\xi _{ij}+\underset{j=k+1}{\overset{m}{}}\underset{i=1}{\overset{j1}{}}R_{ij,km}\xi _{ij}+\underset{j=m+1}{\overset{N}{}}\underset{i=1}{\overset{j1}{}}R_{ij,km}\xi _{ij}.\end{array}$$
Thus
$$\begin{array}{c}\widehat{R}^N(\xi _{km})=\frac{1}{8}\underset{j=1}{\overset{k}{}}\underset{i=1}{\overset{j1}{}}\delta _{j,m}\lambda _j^2(\lambda _i^23\lambda _j^2+3\lambda _k^2)\xi _{k,j}\hfill \\ \hfill +2\delta _{j,k}\lambda _j^2(\lambda _i^2\lambda _j^2+\lambda _m^2)\xi _{j,m}+\delta _{j,k}(\lambda _i^2+\lambda _j^2+\lambda _m^2)(\lambda _i^2\lambda _j^2+\lambda _m^2)\xi _{j,m}\\ \hfill +\underset{j=k+1}{\overset{m}{}}\underset{i=1}{\overset{j1}{}}R_{ij,km}\xi _{ij}+\underset{j=m+1}{\overset{N}{}}\underset{i=1}{\overset{j1}{}}R_{ij,km}\xi _{ij}=\end{array}$$
$$\begin{array}{c}\frac{1}{8}[\underset{i=1}{\overset{k1}{}}2\lambda _k^2(\lambda _i^2\lambda _k^2+\lambda _m^2)+(\lambda _i^2+\lambda _k^2+\lambda _m^2)(\lambda _i^2\lambda _k^2+\lambda _m^2)\hfill \\ \hfill +\lambda _m^2(\lambda _k^23\lambda _m^2+3\lambda _k^2)+\underset{j=k+1}{\overset{m}{}}\lambda _k^2(3\lambda _k^2+\lambda _j^2+3\lambda _m^2)\\ \hfill +\underset{i=1,ik}{\overset{m1}{}}\lambda _m^2(\lambda _i^23\lambda _m^2+3\lambda _k^2)+\underset{j=m+1}{\overset{N}{}}\lambda _k^2(3\lambda _k^2+\lambda _j^2+3\lambda _m^2)\\ \hfill +\underset{j=m+1}{\overset{N}{}}2\lambda _m^2(\lambda _j^2+\lambda _m^2\lambda _k^2)+(\lambda _m^2+\lambda _j^2\lambda _k^2)(\lambda _m^2+\lambda _j^2\lambda _k^2)]\xi _{k,m}=\end{array}$$
$$\begin{array}{c}\frac{1}{8}[\underset{l=1}{\overset{k1}{}}2\lambda _k^2(\lambda _l^2\lambda _k^2+\lambda _m^2)+(\lambda _l^2\lambda _k^2\lambda _m^2)(\lambda _l^2+\lambda _k^2\lambda _m^2)\hfill \\ \hfill +\underset{l=k+1}{\overset{N}{}}\lambda _k^2(\lambda _l^23\lambda _k^2+3\lambda _m^2)+\underset{l=1}{\overset{m1}{}}\lambda _m^2(\lambda _l^2+3\lambda _k^23\lambda _m^2)\\ \hfill +\underset{l=m+1}{\overset{N}{}}2\lambda _m^2(\lambda _l^2\lambda _k^2+\lambda _m^2)+(\lambda _l^2\lambda _k^2\lambda _m^2)(\lambda _l^2\lambda _k^2+\lambda _m^2)]\xi _{k,m}=\end{array}$$
$$\begin{array}{c}\frac{1}{8}[(k1)(3\lambda _k^2+\lambda _m^2)3(m1)\lambda _m^2(\lambda _k^2+\lambda _m^2)\hfill \\ \hfill +(Nk)(\lambda _k^2(3\lambda _k^2+3\lambda _m^2))+(Nm)(\lambda _k^23\lambda _m^2)(\lambda _k^2+\lambda _m^2)\\ \hfill +\underset{l=1}{\overset{k1}{}}(2\lambda _k^2\lambda _l^2+\lambda _l^42\lambda _m^2\lambda _l^2)+\underset{l=k+1}{\overset{N}{}}\lambda _k^2\lambda _l^2\\ \hfill +\underset{l=1}{\overset{m1}{}}\lambda _m^2\lambda _l^2+\underset{l=m+1}{\overset{N}{}}\lambda _l^2(2\lambda _m^2+\lambda _l^22\lambda _k^2)]\xi _{k,m}.\end{array}$$
### 6.3. Upper triangular Hilbert-Schmidt algebra
As before let $`\{\lambda _i\}_{i=1}^{\mathrm{}}`$ be a bounded sequence of strictly positive numbers. In this section we consider the infinite dimensional Lie algebra $`𝔤_{\mathrm{}}`$ generated by the orthonormal basis $`\xi _{ij}=\lambda _i\lambda _je_{ij}`$, $`i<j`$. In this case $`𝔤_{\mathrm{}}`$ is a Lie subalgebra of $`𝔥_{HS}`$. Recall that if the sequence $`\{\lambda _i\}_{i=1}^{\mathrm{}}`$ is bounded, then by Theorem 4.4, Continuity Assumption 3.1 is satisfied for the corresponding norm on $`𝔤_{\mathrm{}}`$.
The Lie bracket is
$$[\xi _{ij},\xi _{km}]=\delta _{j,k}\lambda _j^2\xi _{im}\delta _{i,m}\lambda _i^2\xi _{kj},i<j,,k<m.$$
Denote $`_{ab}=_{\xi _{ab}}`$, then as for the general algebra with the convention that $`\xi _{ij}=0`$ if $`ij`$
$$_{ab}\xi _{cd}=\frac{1}{2}\left(\delta _{b,c}\lambda _b^2\xi _{ad}\delta _{a,d}\lambda _a^2\xi _{cb}\delta _{a,c}\lambda _d^2\xi _{db}+\delta _{b,d}\lambda _c^2\xi _{ac}+\delta _{b,d}\lambda _a^2\xi _{ca}\delta _{a,c}\lambda _b^2\xi _{bd}\right).$$
###### Lemma 6.9.
The Riemannian curvature tensor applied to $`\xi _{ij}`$ is
$$\begin{array}{c}R_{\xi _{ij}\xi _{km}}\xi _{ij}=\frac{1}{4}(2\delta _{i,k}\delta _{j,m}\lambda _i^4\xi _{ij}+2\delta _{i,k}\delta _{j,m}\lambda _j^4\xi _{ij}3\delta _{i,m}\lambda _i^4\xi _{ki}3\delta _{j,k}\lambda _j^4\xi _{jm}\hfill \\ \hfill +\delta _{i,k}\lambda _j^4\xi _{im}+\delta _{i,k}\lambda _m^4\xi _{im}+\delta _{j,m}\lambda _i^4\xi _{kj}+\delta _{j,m}\lambda _k^4\xi _{kj})\end{array}$$
The sectional curvature is
$$\begin{array}{c}(R_{ij,km}\xi _{ij},\xi _{km})=\hfill \\ \hfill \frac{1}{4}(2\delta _{i,k}\delta _{j,m}\lambda _i^4+2\delta _{i,k}\delta _{j,m}\lambda _j^43\delta _{i,m}\lambda _i^43\delta _{j,k}\lambda _j^4+\delta _{i,k}\lambda _j^4+\delta _{i,k}\lambda _m^4+\delta _{j,m}\lambda _i^4+\delta _{j,m}\lambda _k^4).\end{array}$$
###### Proof.
The Riemannian curvature tensor applied to $`\xi _{ij}`$ is
$$\begin{array}{c}R_{\xi _{ij}\xi _{km}}\xi _{ij}=R_{ij,km}\xi _{ij}=\hfill \\ \hfill \delta _{j,k}\lambda _j^2_{im}\xi _{ij}\delta _{i,m}\lambda _i^2_{kj}\xi _{ij}_{ij}_{km}\xi _{ij}+_{km}_{ij}\xi _{ij}=\\ \hfill \frac{1}{2}(\delta _{j,k}\lambda _j^4\xi _{jm}\delta _{i,m}\lambda _i^4\xi _{ki})+\frac{1}{4}(+2\delta _{i,k}\delta _{j,m}\lambda _i^4\xi _{ij}+2\delta _{i,k}\delta _{j,m}\lambda _j^4\xi _{ij}\\ \hfill \delta _{i,m}\lambda _i^4\xi _{ki}\delta _{j,k}\lambda _j^4\xi _{jm}+\delta _{i,k}\lambda _j^4\xi _{im}+\delta _{i,k}\lambda _m^4\xi _{im}+\delta _{j,m}\lambda _i^4\xi _{kj}+\delta _{j,m}\lambda _k^4\xi _{kj})=\\ \hfill \frac{1}{4}(2\delta _{i,k}\delta _{j,m}\lambda _i^4\xi _{ij}+2\delta _{i,k}\delta _{j,m}\lambda _j^4\xi _{ij}3\delta _{i,m}\lambda _i^4\xi _{ki}3\delta _{j,k}\lambda _j^4\xi _{jm}\\ \hfill +\delta _{i,k}\lambda _j^4\xi _{im}+\delta _{i,k}\lambda _m^4\xi _{im}+\delta _{j,m}\lambda _i^4\xi _{kj}+\delta _{j,m}\lambda _k^4\xi _{kj}).\end{array}$$
The third part of the following theorem says that the principal Ricci curvatures for $`𝔥_{HS}`$ tend to $`\mathrm{}`$ as the dimension $`N\mathrm{}`$. This can interpreted as the Ricci curvature being $`\mathrm{}`$. Note that the condition on the scaling $`\{\lambda _i\}_{i=1}^{\mathrm{}}\mathrm{}^2`$ corresponds to the condition we assumed in , , , . We needed this condition to construct a heat kernel measure living in $`HS+I`$. This means that such a measure exists even though the Ricci curvature is $`\mathrm{}`$.
###### Theorem 6.10.
Let $`N>\mathrm{max}i,j`$ then
1. the truncated Ricci curvature is
$$\begin{array}{c}R_{ij}^N=\underset{k,m=1}{\overset{k,m=N}{}}(R_{ij,km}\xi _{ij},\xi _{km})=\hfill \\ \hfill \frac{1}{4}((43i+j)\lambda _i^4+(2+3ji2N)\lambda _j^4+\underset{l=i+1}{\overset{N}{}}\lambda _l^4+\underset{l=1}{\overset{j1}{}}\lambda _l^4).\end{array}$$
For the Hilbert-Schmidt inner product the truncated Ricci curvature
$$R_{ij}^N=\frac{1}{4}(55i+5jN).$$
2. The truncated self-adjoint Ricci curvature is diagonal in the basis $`\{\xi _{km}\}`$. Let $`\{a_{km}\}`$ be its principal Ricci curvatures, that is, $`\widehat{R}^N(\xi _{km})=a_{km}\xi _{km}`$. Then if $`\{\lambda _i\}_{i=1}^{\mathrm{}}\mathrm{}^2`$, the principal Ricci curvatures have the following asymptotics as $`N\mathrm{}`$
$$a_{km}=\left(B_{km}^N\frac{N}{2}\lambda _m^4\right)\xi _{k,m},$$
where $`B_{km}^NB_{km}<\mathrm{}`$ as $`N\mathrm{}`$.
###### Proof.
(1) follows directly from the previous lemmas. (2) Let $`k<m<N`$, then the truncated adjoint Ricci curvature is
$$\begin{array}{c}\widehat{R}^N(\xi _{km})=\underset{i<jN}{}R_{ij,km}\xi _{ij}=\underset{j=1}{\overset{N}{}}\underset{i=1}{\overset{j1}{}}R_{ij,km}\xi _{ij}=\hfill \\ \hfill \underset{j=1}{\overset{k}{}}\underset{i=1}{\overset{j1}{}}R_{ij,km}\xi _{ij}+\underset{j=k+1}{\overset{m}{}}\underset{i=1}{\overset{j1}{}}R_{ij,km}\xi _{ij}+\underset{j=m+1}{\overset{N}{}}\underset{i=1}{\overset{j1}{}}R_{ij,km}\xi _{ij}.\end{array}$$
$$\begin{array}{c}\widehat{R}^N(\xi _{km})=\underset{i<jN}{}R_{ij,km}\xi _{ij}=\underset{j=1}{\overset{N}{}}\underset{i=1}{\overset{j1}{}}R_{ij,km}\xi _{ij}=\hfill \\ \hfill \underset{j=1}{\overset{k}{}}\underset{i=1}{\overset{j1}{}}R_{ij,km}\xi _{ij}+\underset{j=k+1}{\overset{m}{}}\underset{i=1}{\overset{j1}{}}R_{ij,km}\xi _{ij}+\underset{j=m+1}{\overset{N}{}}\underset{i=1}{\overset{j1}{}}R_{ij,km}\xi _{ij}=\\ \hfill \frac{(3k+m+4)}{4}\lambda _k^4\xi _{km}+\frac{(2k2N+3m)}{4}\lambda _m^4\xi _{km}+\frac{1}{4}\underset{l=k+1}{\overset{N}{}}\lambda _l^4\xi _{km}+\frac{1}{4}\underset{l=1}{\overset{m1}{}}\lambda _l^4\xi _{km}=\\ \hfill B_{km}^N\xi _{km}\frac{N}{2}\lambda _m^4\xi _{km}.\end{array}$$
###### Corollary 6.11.
Suppose $`\{\lambda _i\}_{i=1}^{\mathrm{}}\mathrm{}^2`$. Then for large $`N`$
$$\widehat{R}^N(\xi _{km})=b_{km}^N\xi _{km},$$
where $`b_{km}^N\mathrm{}`$ as $`N\mathrm{}`$. This can be described as the Ricci curvature being negative infinity for any $`x𝔤`$. |
warning/0506/hep-ph0506051.html | ar5iv | text | # Charm meson resonances and 𝐷→𝑉 semileptonic form factors
## I Introduction
Presently, one of the most important issues in hadronic physics is the extraction of the CKM parameters from exclusive decays. An essential ingredient in this approach is the knowledge of the form factors’ shapes in heavy to light weak transitions. Usually, the attention has been devoted to $`B`$ decays and the determination of the phase of the $`V_{ub}`$ CKM matrix element. At the same time in the charm sector, the most accurate determination of the size of $`V_{cs}`$ and $`V_{cd}`$ matrix elements is not from a direct measurement, mainly due to theoretical uncertainties in the calculations of the relevant form factors’ shapes.
Recently CLEO and FOCUS have published interesting results on $`D^0\pi ^{}\mathrm{}^+\nu _{\mathrm{}}`$ and $`D^0K^{}\mathrm{}^+\nu _{\mathrm{}}`$ decays Link et al. (2005a); Huang et al. (2005a). Their studies indicate that the single pole parametrization of the relevant form factor cannot explain their data very well, leading to unphysical pole masses. Both experimental groups have also attempted a modified pole fit, which was first put forward for $`B`$ decays in Becirevic and Kaidalov (2000), and their results suggest the existence of contributions beyond lowest lying charm meson resonances.
On the other hand we have recently Fajfer and Kamenik (2005) reconsidered $`DP\mathrm{}\nu _{\mathrm{}}`$ decay form factors within a framework which combines heavy meson and chiral symmetries (HM$`\chi `$T) and includes in the interacting Lagrangian contributions coming from excited charm meson states. We have found that a two-poles shape of the relevant form factor can be successfully accommodated within HM$`\chi `$T when excited meson states are included into the model. In our approach the first pole is described by the lowest lying vector resonance, as in the original idea Becirevic and Kaidalov (2000), while for the second one we assume complete saturation by the next vector state which we include in our Lagrangian. In doing this we anticipate the discrepancies from the general two-poles procedure, in which the second effective pole should account for all other excitations that might be exchanged in the $`t`$-channel, to be small and encoded in the parameters of the model. The unknown parameters have been obtained by fitting the experimental results for the branching ratio. The assumed pole behavior agrees well with experimental results confirming our anticipation of small saturation error.
In addition to studies of heavy to light pseudoscalar meson weak transitions ($`HP`$), transitions of heavy pseudoscalar mesons to light vector mesons ($`HV`$) such as $`D_s\varphi \mathrm{}\nu _{\mathrm{}}`$ an $`D_sK^{}\mathrm{}\nu _{\mathrm{}}`$ offer an opportunity to extract the size of the relevant CKM matrix elements. We continue with our study and re-investigate vector and axial-vector form factors in $`DV\mathrm{}\nu _{\mathrm{}}`$ decays within a similar framework as in the case of $`DP\mathrm{}\nu _{\mathrm{}}`$. The $`HV`$ transitions were already carefully investigated within many different frameworks such as perturbative QCD Kurimoto et al. (2001); Mahajan (2004), QCD sum rules Ball (1993); Ball and Zwicky (2005); Bakulev et al. (2000); Wang and Wu (2001); Du et al. (2004); Aliev et al. (2004), lattice QCD Flynn et al. (1996); Del Debbio et al. (1998); Demchuk et al. (1997); Abada et al. (2003), a few attempts to use combined heavy meson and chiral Lagrangians (HM$`\chi `$T) Bajc et al. (1996); Casalbuoni et al. (1997), quark models Wirbel et al. (1985); Scora and Isgur (1995); Faustov et al. (1996); Melikhov and Stech (2000), large energy effective theory (LEET) Charles et al. (1999) and soft collinear effective theory (SCET) Beneke and Feldmann (2001); Bauer et al. (2001); Burdman and Hiller (2001); Ebert et al. (2001); Hill et al. (2004); Hill (2004). Each of these approaches has only a limited range of validity. For example, the QCD sum rules, LEET and SCET are suitable only for the low $`q^2`$ region while lattice QCD and HM$`\chi `$T are successful for maximal $`q^2`$. It is important to note, that currently lattice QCD and QCD sum rules are the only approaches that enable the computation of form factors solely from first principles. On the other hand quark models usually involve parameters which have little physical correspondence to the underlying theory of QCD.
The experimental situation in $`DV\mathrm{}\nu _{\mathrm{}}`$ has not changed a lot in the last few years, but recently it has been gaining pace Link et al. (2002, 2004a, 2004b, 2005b); Coan et al. (2005); Huang et al. (2005b), and hopefully more results on the $`q^2`$ shape of the form factors will be available soon. Unlike in the case of $`HP`$ weak transitions, no general parametrization of the form factors, relevant to $`HV`$ weak decays has yet been proposed. Usually a simple pole behavior of all the form factors is assumed when extracting values of the form factors at $`q^2=0`$ from experiment or extrapolating results of different theoretical approaches.
Recently, the spectroscopy of charm mesons has been enriched by discoveries of many new charm meson resonances. BaBar Aubert et al. (2003) collaboration has announced a new, narrow meson $`D_{sJ}(2317)^+`$. This was confirmed by Focus Vaandering (2004) and CLEO Besson et al. (2004) which also noticed another narrow state, $`D_{sJ}(2463)^+`$. Both states were confirmed by Belle Krokovny et al. (2003). Finally, Selex Evdokimov et al. (2004) has announced a new, surprisingly narrow state $`D_{sJ}^+(2632)`$. The states $`D_{sJ}(2317)^+`$ and $`D_{sJ}(2463)^+`$ are already being identified to belong to the $`(0^+,1^+)`$ spin-parity doublet of the $`D_s`$ mesons while the $`D_{sJ}^+(2632)`$ state has been proposed as the first radial excitation of the $`D_s^{}(2112)`$ with the spin parity assignment $`1^{}`$ Barnes et al. (2004); van Beveren and Rupp (2004); Dai et al. (2004).
The purpose of this study is to (1) devise a general parametrization of all the form factors relevant to $`HV`$ weak transitions which would take into account known experimental results on heavy meson resonances as well as known theoretical limits of heavy quark effective theory (HQET) and LEET relevant to $`HV`$ weak transitions; to (2) investigate contributions of the newly discovered charm mesons to $`DV`$ semileptonic decays within an effective model based on HM$`\chi `$T by incorporating the newly discovered heavy meson fields into the HM$`\chi `$T Lagrangian and utilizing the general form factor parametrization. We restrain our discussion to the leading chiral and $`1/m_H`$ terms in the expansion, but we hope to capture the main physical features about the impact of the nearest poles in the $`t`$-channel to the $`q^2`$-dependence of the form factors.
In Sec. II we revise the common weak current matrix element decomposition relevant to transitions between pseudoscalar and vector mesons and introduce a form factor decomposition, which is independent of the mass of the pseudoscalar meson and thus convenient for studying $`HV`$ weak transitions. In Sec. III we derive a general $`HV`$ form factor parametrization drawing from both known experimental properties of heavy mesons as well as from known theoretical scaling laws and form factor relations in the limit of the infinite heavy meson mass. Sec. IV describes the framework we use in our HM$`\chi `$T calculations: we write down the HM$`\chi `$T Lagrangian for for heavy and light mesons and extend it to incorporate new heavy meson fields. In Sec. V we calculate the values of the $`DV`$ semileptonic form factors near zero recoil within HM$`\chi `$T and extrapolate our results to larger recoils using the general parametrization of Sec. III and by saturating the effective poles with physical masses of experimentally known or theoretically predicted charmed resonances. Finally, a short summary of the results and comparison with other approaches as well as with existing experimental data is given in Sec. VI.
## II Parametrization of $`HV`$ current matrix element
A frequently encountered decomposition of the current matrix elements relevant to semileptonic decays between a heavy pseudoscalar meson state $`|H(p_H)`$ with momentum $`p_H^\nu `$ and a light vector meson state $`|V(p_V,ϵ_V)`$ with momentum $`p_V^\nu `$ and polarization vector $`ϵ_V^\nu `$ is
$`V(ϵ_V,p_V)|\overline{q}\gamma ^\mu Q|H(p_H)`$ $`=`$ $`{\displaystyle \frac{2V(q^2)}{m_H+m_V}}ϵ^{\mu \nu \alpha \beta }ϵ_{V\nu }^{}p_{H\alpha }p_{V\beta },`$
$`V(ϵ_V,p_V)|\overline{q}\gamma ^\mu \gamma ^5Q|H(p_H)`$ $`=`$ $`iϵ_V^{}q{\displaystyle \frac{2m_V}{q^2}}q^\mu A_0(q^2)i(m_H+m_V)\left[ϵ_V^\mu {\displaystyle \frac{ϵ_V^{}q}{q^2}}q^\mu \right]A_1(q^2)`$
$`+i{\displaystyle \frac{ϵ_V^{}q}{(m_H+m_V)}}\left[(p_H+p_V)^\mu {\displaystyle \frac{m_H^2m_V^2}{q^2}}q^\mu \right]A_2(q^2),`$
where $`q=u,d`$ or $`s`$ are the light quark fields, $`Q=b`$ or $`c`$ denote the heavy quark fields, $`q^\nu =(p_Hp_V)^\nu `$ is the exchanged momentum and $`q^2`$ is the exchanged momentum squared. Here $`V`$ denotes the vector form factor and is expected to be dominated by vector meson resonance exchange, the axial $`A_1`$ and $`A_2`$ form factors are expected to be dominated by axial resonances, while $`A_0`$ denotes the pseudoscalar form factor and is expected to be dominated by pseudoscalar meson resonance exchange Wirbel et al. (1985). In order that these matrix elements are finite at $`q^2=0`$, the form factors must also satisfy the well known relation
$$A_0(0)\frac{m_H+m_V}{2m_V}A_1(0)+\frac{m_Hm_V}{2m_V}A_2(0)=0.$$
(2)
We will work in the static limit of HQET where the eigenstates of QCD and HQET Lagrangians are related as
$$\underset{m_H\mathrm{}}{lim}\frac{1}{\sqrt{m}_H}|H(p_H)_{QCD}=|H(v)_{HQET}.$$
(3)
In this limit it is more convenient to use definitions in which the form factors are independent of the heavy meson mass, namely we propose
$`V(ϵ_V,p_V)|\overline{q}\gamma ^\mu Q_v|H(v)`$ $`=`$ $`f_vϵ^{\mu \nu \alpha \beta }ϵ_{V\nu }^{}v_\alpha p_{V\beta },`$
$`V(ϵ_V,p_V)|\overline{q}\gamma ^\mu \gamma ^5Q_v|H(v)`$ $`=`$ $`ia_2(ϵ_V^{}v)\left[p_V^\mu (vp_V)v^\mu \right]`$ (4)
$`ia_1\left[ϵ_V^\mu (vϵ_V^{})v^\mu \right]`$
$`ia_0(vϵ_V^{})v^\mu ,`$
where the HQET heavy quark field $`Q_v`$ is independent of the heavy quark mass. The form factors $`f_v`$, $`a_1`$, $`a_2`$ and $`a_0`$ are functions of the variable
$$vp_V=\frac{m_H^2+m_V^2q^2}{2m_H},$$
(5)
which in the heavy meson rest frame is the energy of the light meson $`E_V`$. In such decomposition, all the form factors ($`f_v`$, $`a_1`$, $`a_2`$ and $`a_0`$) scale as constants with the heavy meson mass. The relation between the two form factor decompositions is obtained by correctly matching QCD and HQET at the scale $`\mu m_Q`$ Eichten and Hill (1990); Broadhurst and Grozin (1995):
$`{\displaystyle \frac{C_{\gamma _1}(m_Q)}{\sqrt{m}_H}}\left[f_V(vp_V)+𝒪(1/m_H)\right]`$ $`=`$ $`{\displaystyle \frac{2V(q^2)}{m_H+m_V}}|_{q^2q_{\mathrm{max}}^2},`$
$`{\displaystyle \frac{C_{\gamma _0\gamma _5}(m_Q)}{\sqrt{m_H}}}\left[a_0(vp_V)+𝒪(1/m_H)\right]`$ $`=`$ $`\{{\displaystyle \frac{(m_HE_V)}{q^2}}[2m_VA_0(q^2)+(m_H+m_V)A_1(q^2)(m_Hm_V)A_2(q^2)]`$
$`+{\displaystyle \frac{(m_H+E_V)}{m_H+m_V}}A_2(q^2){\displaystyle \frac{(m_H+m_V)}{m_H}}A_1(q^2)\}|_{q^2q_{\mathrm{max}}^2},`$
$`C_{\gamma _1\gamma _5}(m_Q)\sqrt{m_H}\left[a_1(vp_V)+𝒪(1/m_H)\right]`$ $`=`$ $`(m_H+m_V)A_1(q^2)|_{q^2q_{\mathrm{max}}^2},`$
$`{\displaystyle \frac{C_{\gamma _1\gamma _5}(m_Q)}{\sqrt{m}_H}}\left[a_2(vp_V)+𝒪(1/m_H)\right]`$ $`=`$ $`\{{\displaystyle \frac{m_H+m_V}{q^2}}[A_1(q^2)+A_0(q^2)]{\displaystyle \frac{m_Hm_V}{q^2}}[A_2(q^2)+A_0(q^2)]`$ (6)
$`{\displaystyle \frac{A_2(q^2)}{m_H+m_V}}\}|_{q^2q_{\mathrm{max}}^2}.`$
In the following we set the matching constants $`C_\mathrm{\Gamma }`$ to their tree level values $`(C_\mathrm{\Gamma }=1)`$. At leading order in $`1/m_Q`$ we thus get
$`V(q^2)|_{q^2q_{\mathrm{max}}^2}`$ $`=`$ $`{\displaystyle \frac{\sqrt{m}_H}{2}}f_v(vp_V),`$
$`A_1(q^2)|_{q^2q_{\mathrm{max}}^2}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{m}_H}}a_1(vp_V),`$
$`A_2(q^2)|_{q^2q_{\mathrm{max}}^2}`$ $`=`$ $`{\displaystyle \frac{\sqrt{m}_H}{2}}a_2(vp_V),`$
$`A_0(q^2)|_{q^2q_{\mathrm{max}}^2}`$ $`=`$ $`{\displaystyle \frac{\sqrt{m}_H}{2m_V}}a_0(vp_V),`$ (7)
which exhibit the usual heavy meson mass scaling laws for the semileptonic form factors Isgur and Wise (1990). This parametrization is especially useful when calculating the form factors within HM$`\chi `$T. The individual contributions of different terms in the HM$`\chi `$T Lagrangian to various form factors can be easily projected out.
## III Parametrization of the form factors
Next we propose a general parametrization of the heavy to light vector form factors, which takes into account all the known scaling and resonance properties of the form factors. As already evident from Eq. (7), there exist the well known HQET scaling laws in the limit of zero recoil Isgur and Wise (1990). On the other hand in the large energy limit $`q^20`$, one obtains the following expressions for the form factors Charles et al. (1999)
$`V(q^2)|_{q^20}`$ $`=`$ $`{\displaystyle \frac{m_H+m_V}{m_H}}\xi _{}(E_V),`$
$`A_1(q^2)|_{q^20}`$ $`=`$ $`{\displaystyle \frac{2E_V}{m_H+m_V}}\xi _{}(E_V),`$
$`A_2(q^2)|_{q^20}`$ $`=`$ $`{\displaystyle \frac{m_H+m_V}{m_H}}\left[\xi _{}(E_V){\displaystyle \frac{m}{E}}\xi _{}(E_V)\right],`$
$`A_0(q^2)|_{q^20}`$ $`=`$ $`\left(1{\displaystyle \frac{m_V^2}{2E_Vm_H}}\right)\xi _{}(E_V)\xi _{}(E_V),`$ (8)
where both universal LEET functions $`\xi _{}`$ and $`\xi _{}`$ scale with the heavy meson mass as $`m_H^{3/2}`$. These scaling laws were subsequently confirmed by means of SCET Beneke and Feldmann (2001); Ebert et al. (2001). This is important since the LEET description breaks down beyond the tree level due to missing soft gluonic degrees of freedom which are however systematically taken into account within SCET.
The starting point is the vector form factor $`V`$, which is dominated by the pole at $`t=m_H^{}^2`$ when considering the part of the phase space that is close to the zero recoil. It is very easy to see, that the residuum at that pole scales as $`m_H^{3/2}`$ with the heavy meson mass Becirevic and Kaidalov (2000). For the heavy to light transitions this situation is expected to be realized near the zero recoil where also the HQET scaling (7) applies. However, since the kinematically accessible region $`q^2(0,q_{\mathrm{max}}^2]`$ is large, the pole dominance can be used only on a small fraction of the phase space, $`i.e.`$ for $`|\stackrel{}{q}|0`$. Even in this region the situation for $`HV`$ form factors is more complex than in the case of $`HP`$ transitions, where $`q_{\mathrm{max}}^2`$ is indeed very close to the vector pole due to low mass of the light pseudoscalar mesons. Here, due to larger masses of the light vector mesons, $`q_{\mathrm{max}}^2`$ is pushed away from the resonance pole and the form factor may not be completely saturated by it. For the sake of clarity and conciseness we, however, in our present study neglect such possible discrepancies and assume complete saturation of the vector form factor in this region by the first physical resonance. On the other hand, in the region of large recoils, LEET dictates the scaling (8). In the full analogy with the discussion made in Refs. Becirevic and Kaidalov (2000); Hill (2005), the vector form factor consequently receives contributions from two poles and can be written as
$$V(q^2)=c_H^{}\frac{1a}{(1x)(1ax)},$$
(9)
where $`x=q^2/m_H^{}^2`$ ensures, that the form factor is dominated by the physical $`H^{}`$ pole, while $`a`$ measures the contribution of higher states which are parametrized by another effective pole at $`m_{\mathrm{eff}}^2=m_H^{}^2/a`$. The parameters $`c_H^{}`$ and $`a`$ scale with the heavy meson mass as $`c_H^{}m_H^{1/2}`$ and $`a1a_0/m_H`$ to ensure the correct form factor scaling in both small and large recoil regions.
An interesting and useful feature one gets from the large energy limit is the relation between $`V`$ and $`A_1`$ Charles et al. (1999)
$$\left[V(q^2)/A_1(q^2)\right]|_{q^20}=\frac{(m_H+m_V)^2}{2E_Vm_H},$$
(10)
which is valid up to terms $`1/m_H^2`$ Ebert et al. (2001). This relation remains valid even when the leading order corrections due to soft gluon exchange are taken into account Burdman and Hiller (2001); Hill (2004). When combined with our result (9), it imposes a single pole structure on $`A_1`$. We can thus continue in the same line of argument and write
$$A_1(q^2)=c_H^{}\xi \frac{1a}{1b^{}x}.$$
(11)
Here $`\xi =m_H^2/(m_H+m_V)^2`$ is the proportionality factor between $`A_1`$ and $`V`$ from (10), while $`b^{}`$ measures the contribution of higher states with spin-parity assignment $`1^+`$ which are parametrized by the effective pole at $`m_{H_{\mathrm{eff}}^{}}^2=m_H^{}^2/b^{}`$. It can be readily checked that also $`A_1`$, when parametrized in this way, satisfies all the scaling constraints.
Next we parametrize the $`A_0`$ form factor, which is completely independent of all the others so far as it is dominated by the pseudoscalar pole and is proportional to a different universal function in LEET. To satisfy both HQET and LEET scaling laws we parametrize it as
$$A_0(q^2)=c_H^{\prime \prime }\frac{1a^{}}{(1y)(1a^{}y)},$$
(12)
where $`y=q^2/m_H^2`$ ensures the physical $`0^{}`$ pole dominance at small recoils. Imposing $`c_H^{\prime \prime }m_H^{1/2}`$ and $`a^{}1a_0^{}/m_H`$ preserves all scaling laws, while $`a^{}`$ again parametrizes the contribution of higher pseudoscalar states by an effective pole at $`m_{H_{\mathrm{eff}}^{}}^2=m_H^2/a^{}`$. The resemblance to $`V`$ is obvious and due to the same kind of analysis Becirevic and Kaidalov (2000) although the parameters appearing in the two form factors are completely unrelated.
Finally for the $`A_2`$ form factor, due to the pole behavior of the $`A_1`$ form factor on one hand and different HQET scaling at $`q_{\mathrm{max}}^2`$ (7) on the other hand, we have to go beyond a simple pole formulation. Thus we impose
$$A_2(q^2)=\frac{(m_H+m_V)\xi c_H^{}(1a)+2m_Vc_H^{\prime \prime }(1a^{})}{(m_Hm_V)(1b^{}x)(1b^{\prime \prime }x)},$$
(13)
which again satisfies all constraints. Due to the relation (2) we only gain one new parameter in this formulation, $`b^{\prime \prime }`$. This however causes the contribution of the $`1^+`$ resonances to be shared between the two effective poles in this form factor.
At the end we have parametrized the four $`HV`$ vector form factors in terms of the six parameters $`c_H^{}`$, $`a`$, $`b^{}`$, $`a^{}`$, $`c_H^{\prime \prime }`$ and $`b^{\prime \prime }`$.
It is convenient to introduce helicity amplitudes for the decays $`HV\mathrm{}\nu `$ as in for example Ball et al. (1991):
$`H_\pm (y)`$ $`=`$ $`+(m_H+m_V)A_1(m_H^2y){\displaystyle \frac{2m_H|\stackrel{}{p}_V(y)|}{m_H+m_V}}V(m_H^2y)`$
$`H_0(y)`$ $`=`$ $`+{\displaystyle \frac{m_H+m_V}{2m_Hm_V\sqrt{y}}}[m_H^2(1y)m_V^2]A_1(m_H^2y)`$ (14)
$`{\displaystyle \frac{2m_H|\stackrel{}{p}_V(y)|}{m_V(m_H+m_V)\sqrt{y}}}A_2(m_H^2y)`$
where $`y=q^2/m_H^2`$ and the three-momentum of the light vector meson is given by:
$$|\stackrel{}{p}_V(y)|^2=\frac{[m_H^2(1y)+m_V^2]^2}{4m_H^2}m_V^2.$$
(15)
As shown in Ref. Ebert et al. (2001) these helicity amplitudes can be related to individual form factors near $`q^2=0`$. Using relations (8), valid in the large energy limit, one can write
$`H_{}(y)|_{y0}`$ $``$ $`2(m_H+m_V)A_1(m_H^2y),`$
$`H_+(y)|_{y0}`$ $``$ $`0,`$ (16)
but also, by using relation (2)
$$H_0(y)|_{y0}\frac{2|\stackrel{}{p}_V(y)|}{\sqrt{y}}A_0(m_H^2y).$$
(17)
Thus in this region we can probe directly for the parameters $`c_H^{}(1a)`$ and $`c_H^{\prime \prime }(1a^{})`$.
On the other hand in the region of small recoil ($`|\stackrel{}{p}_V|0`$ or $`yy_{\mathrm{max}}`$) the helicity amplitudes are saturated by the $`A_1`$ form factor
$`H_\pm (y)|_{yy_{\mathrm{max}}}`$ $``$ $`(m_H+m_V)A_1(m_H^2y),`$
$`H_0(y)|_{yy_{\mathrm{max}}}`$ $``$ $`2(m_H+m_V){\displaystyle \frac{m_V}{m_H}}A_1(m_H^2y).`$ (18)
Consequently we can also directly probe for the value of the $`b^{}`$ parameter determining the position of the first effective axial resonance pole by taking a ratio of $`H_{}`$ helicity amplitude values at small and large recoils
$$\frac{H_{}(y)|_{y0}}{H_{}(y)|_{yy_{\mathrm{max}}}}2\left[1b^{}(m_Hm_V)^2/m_H^2\right].$$
(19)
## IV The Model
### IV.1 Strong interactions
At leading order in chiral and $`1/m_H`$ expansion, strong interactions between lowest lying pseudoscalar and vector heavy meson fields, and light vector meson fields are described by the interaction Lagrangian Bajc et al. (1996); Casalbuoni et al. (1997)
$$_{\mathrm{int}}=i\beta H_bv_\mu \widehat{\rho }_{ba}^\mu \overline{H}_a+i\lambda H_b\sigma ^{\mu \nu }F_{\mu \nu }(\widehat{\rho })_{ba}\overline{H}_a,$$
(20)
where the first term is even under parity transformation, and the second term is parity odd. $`H=1/2(1+/v)[P_\mu ^{}\gamma ^\mu P\gamma _5]`$ is the matrix representation of the heavy meson fields, where $`P_\mu ^{}`$ and $`P`$ are creation operators for heavy vector and pseudoscalar mesons respectively. Light vector meson fields are described by $`\widehat{\rho }_\mu =i\frac{g_V}{\sqrt{2}}\rho _\mu `$, where $`\rho _\mu `$ is the light vector meson field matrix
$$\rho _\mu =\left(\begin{array}{ccc}\frac{1}{\sqrt{2}}(\omega _\mu +\rho _\mu ^0)& \rho _\mu ^+& K_\mu ^+\\ \rho _\mu ^{}& \frac{1}{\sqrt{2}}(\omega _\mu \rho _\mu ^0)& K_\mu ^0\\ K_\mu ^{}& \overline{K}_\mu ^0& \varphi _\mu \end{array}\right).$$
(21)
The gauge field tensor is defined as $`F_{\mu \nu }(\widehat{\rho })=_\mu \widehat{\rho }_\nu _\nu \widehat{\rho }_\mu +[\widehat{\rho }_\mu ,\widehat{\rho }_\nu ]`$. Furthermore, $`\mathrm{}`$ indicate a trace over spinor matrices and summation over light quark flavor indexes.
In order to incorporate positive parity heavy meson states into the model, we introduce the scalar-axial field multiplet $`G=1/2(1+/v)[S_\mu ^{}\gamma ^\mu \gamma _5S]`$ representing axial ($`S_\mu ^{}`$) and scalar ($`S`$) mesons and incorporate it into the interaction Lagrangian by adding additional leading order interaction terms between heavy even and odd parity fields and light vector fields:
$`_{\mathrm{int}}^{}`$ $`=`$ $`i\zeta H_bv_\mu \widehat{\rho }_{ba}^\mu \overline{G}_a+\mathrm{h}.\mathrm{c}.`$ (22)
$`+i\mu H_b\sigma ^{\mu \nu }F_{\mu \nu }(\widehat{\rho })_{ba}\overline{G}_a+\mathrm{h}.\mathrm{c}..`$
There exists another field multiplet in HM$`\chi `$T containing positive parity heavy meson states, $`T^\mu =1/2(1+/v)[T_1^{\mu \nu }\gamma _\nu \sqrt{3/2}T_{2\nu }\gamma _5(g^{\mu \nu }1/3\gamma ^\nu (\gamma ^\mu v^\mu ))]`$, where $`T_1^{\mu \nu }`$ is the tensor field with spin-parity assignment $`2^+`$, while $`T_2^\mu `$ is another $`1^+`$ axial vector meson field. However, as pointed out in Ref. Casalbuoni et al. (1997), the matrix element of the HM$`\chi `$T bosonized currents containing these fields between a single heavy meson state $`|H(p_H)`$ and the vacuum vanishes at leading order in $`1/m_H`$ due to heavy quark spin symmetry. Consequently, such fields do not contribute at leading order to $`HV`$ semileptonic decays.
Finally we also want to include the radially excited states into our discussion and therefore introduce another odd parity heavy meson multiplet field $`H^{}=1/2(1+/v)[P_\mu ^{{}_{}{}^{}}\gamma ^\mu P^{}\gamma _5]`$ containing the radial excitations of ground state pseudoscalar and vector mesons. Such excited states were predicted in Di Pierro and Eichten (2001). The strong interactions between these fields, ground state heavy meson fields $`H`$ and light vector fields can again be described by the lowest order interaction Lagrangian analogous to (22)
$`\stackrel{~}{}_{\mathrm{int}}`$ $`=`$ $`i\stackrel{~}{\zeta }\stackrel{~}{H}_bv_\mu \widehat{\rho }_{ba}^\mu \overline{H}_a`$ (23)
$`+i\stackrel{~}{\mu }\stackrel{~}{H}_b\sigma ^{\mu \nu }F_{\mu \nu }(\widehat{\rho })_{ba}\overline{H}_a+\mathrm{h}.\mathrm{c}..`$
### IV.2 Weak interactions
For the semileptonic decays the weak Lagrangian can be given by the effective current-current Fermi interaction
$$_{\mathrm{eff}}=\frac{G_F}{\sqrt{2}}\left[\overline{\mathrm{}}\gamma ^\mu (1\gamma ^5)\nu _{\mathrm{}}𝒥_\mu \right],$$
(24)
where $`G_F`$ is the Fermi constant and $`𝒥`$ is the effective hadronic current. In heavy to light meson decays it can be written as $`𝒥=K^aJ_a`$, where constants $`K^a`$ parametrize the $`SU(3)`$ flavor mixing, while the leading order weak current $`J_a`$ in chiral and $`1/m_H`$ expansion between heavy ground state pseudoscalar and vector mesons and light vector mesons can be written as Bajc et al. (1996); Casalbuoni et al. (1997)
$`J_a^\mu `$ $`=`$ $`{\displaystyle \frac{1}{2}}i\alpha \gamma ^\mu (1\gamma ^5)H_a`$ (25)
$`+\alpha _1\gamma ^5H_b\widehat{\rho }_{ba}^\mu +\alpha _2\gamma ^\mu \gamma ^5H_bv_\alpha \widehat{\rho }_{ba}^\alpha .`$
For our leading order calculation, we will also need the weak current operator for the scalar and axial heavy mesons’ transition to the vacuum
$$J_a^{}_{}{}^{}\mu =\frac{1}{2}i\alpha ^{}\gamma ^\mu (1\gamma ^5)G_a,$$
(26)
and the same for radially excited pseudoscalar and vector fields
$$\stackrel{~}{J}_a^\mu =\frac{1}{2}i\stackrel{~}{\alpha }\gamma ^\mu (1\gamma ^5)H_a^{}.$$
(27)
## V Form Factor Calculation
### V.1 HM$`\chi `$T calculation at zero recoil
In HM$`\chi `$PT the derived Feynman rules are valid near zero recoil ($`|\stackrel{}{p}_V|0`$). For the heavy meson propagators we use $`i\delta _{ab}/2(vk\mathrm{\Delta })`$ and $`i\delta _{ab}(g_{\mu \nu }v_\mu v_\nu )/2(vk\mathrm{\Delta })`$ for the pseudoscalar(scalar) and vector (axial)mesons respectively, where $`k^\mu =q^\mu m_Hv^\mu `$. $`\mathrm{\Delta }=\mathrm{\Delta }_R`$ is the mass splitting between the heavy resonance meson $`R`$ and the ground state strangeless heavy pseudoscalar meson. It comes from leading order $`1/m_H`$ (spin symmetry breaking), chiral and $`SU(3)`$ breaking corrections, when all heavy meson fields are normalized to physical masses of ground state strangeless heavy pseudoscalar mesons. For the hadronic current matrix element we thus get
$`V(p_V)|J^\mu |H(v)=i\sqrt{2}g_V\left(\alpha _1ϵ_V^\mu \alpha _2vϵ_Vv^\mu \right)`$
$`\sqrt{2}g_V\alpha {\displaystyle \frac{\lambda ϵ^{\mu \nu \alpha \beta }v_\nu p_{V\alpha }ϵ_{V\beta }}{vp_V+\mathrm{\Delta }_H^{}}}\sqrt{2}g_V\stackrel{~}{\alpha }{\displaystyle \frac{\stackrel{~}{\mu }ϵ^{\mu \nu \alpha \beta }v_\nu p_{V\alpha }ϵ_{V\beta }}{vp_V+\mathrm{\Delta }_H^{}}}`$
$`i{\displaystyle \frac{g_V}{\sqrt{2}}}\alpha {\displaystyle \frac{\beta vϵ_Vv^\mu }{vp_V+\mathrm{\Delta }_{H_P}}}i{\displaystyle \frac{g_V}{\sqrt{2}}}\stackrel{~}{\alpha }{\displaystyle \frac{\stackrel{~}{\zeta }vϵ_Vv^\mu }{vp_V+\mathrm{\Delta }_{H_P^{}}}}`$
$`i{\displaystyle \frac{g_V}{\sqrt{2}}}\alpha ^{}{\displaystyle \frac{ϵ_V^\mu \left(\zeta 2\mu vp_V\right)+\left(2\mu p_V^\mu \zeta v^\mu \right)vϵ_V}{vp_V+\mathrm{\Delta }_{H_A}}}.`$ (28)
From this we extract the form factors $`V(q^2)`$, $`A_1(q^2)`$, $`A_2(q^2)`$ and $`A_0(q^2)`$:
$`V(q^2)|_{q^2q_{\mathrm{max}}^2}`$ $`=`$ $`{\displaystyle \frac{g_V}{\sqrt{2}}}\alpha m_H\sqrt{m}_H{\displaystyle \frac{\lambda }{vp_V+\mathrm{\Delta }_H^{}}}`$
$`{\displaystyle \frac{g_V}{\sqrt{2}}}\stackrel{~}{\alpha }m_H\sqrt{m}_H{\displaystyle \frac{\stackrel{~}{\mu }}{vp_V+\mathrm{\Delta }_H^{}}}`$
$`A_1(q^2)|_{q^2q_{\mathrm{max}}^2}`$ $`=`$ $`{\displaystyle \frac{g_V}{\sqrt{2}}}\alpha ^{}{\displaystyle \frac{\sqrt{m}_H}{m_H+m_V}}{\displaystyle \frac{\zeta 2\mu vp_V}{vp_V+\mathrm{\Delta }_{H_A}}}`$
$`\sqrt{2}g_V\alpha _1{\displaystyle \frac{\sqrt{m}_H}{m_H+m_V}}`$
$`A_2(q^2)|_{q^2q_{\mathrm{max}}^2}`$ $`=`$ $`{\displaystyle \frac{g_V}{\sqrt{2}}}\alpha ^{}{\displaystyle \frac{m_H+m_V}{\sqrt{m}_H}}{\displaystyle \frac{\mu }{vp_V+\mathrm{\Delta }_{H_A}}}`$
$`A_0(q^2)|_{q^2q_{\mathrm{max}}^2}`$ $`=`$ $`{\displaystyle \frac{g_V}{2\sqrt{2}}}{\displaystyle \frac{\sqrt{m}_H}{m_V}}(2\alpha _12\alpha _2`$
$`+\alpha {\displaystyle \frac{\beta }{vp_V+\mathrm{\Delta }_{H_P}}}+\stackrel{~}{\alpha }{\displaystyle \frac{\stackrel{~}{\zeta }}{vp_V+\mathrm{\Delta }_{H_P^{}}}})`$
### V.2 Extrapolation to higher recoils
In order to extrapolate our HM$`\chi `$T calculation results at $`q^2q_{\mathrm{max}}^2`$ to higher recoils we employ the general analysis from Sec. III. We model the form factors’ $`q^2`$ behavior using the formulas (9), (11), (12) and (13) with model matching conditions at $`q_{\mathrm{max}}^2`$. In order to reduce the number of free parameters in this extrapolation we employ the same strategy as in our previous work Fajfer and Kamenik (2005). We use the information on the contributions of different resonances to the form factors as suggested by our model. For the vector form factor $`V`$ we thus propose $`a=m_H^{}^2/m_H^{}^2`$ which saturates the effective second pole by the first vector radial excitation $`H^{}`$. Similarly we set $`b^{}=m_H^{}^2/m_{H_A}^2`$ and $`a^{}=m_H^2/m_H^{}^2`$ saturating the poles of the $`A_0`$ and $`A_1`$ form factors and the first pole of the $`A_2`$ form factor with the $`H^{}`$ pseudoscalar radial excitation and the $`H_A`$ orbital axial excitation respectively. Since our model does not contain a second resonance contribution to the $`A_2`$ form factor, we impose $`b^{\prime \prime }=0`$, effectively sending the second pole mass of this form factor to infinity.
At the end we have fixed all the pole parameters appearing in the general form factor parametrization formulas of Sec. III using physical information and model predictions on the resonances contributing to the various form factors. The remaining parameters ($`c_H`$ and $`c_H^{}`$) are on the other hand related to the parameters of HM$`\chi `$T via the model matching conditions at zero recoil.
## VI Numerical results
We now leave the general discussion of $`HV`$ transitions and restrict our present study to $`D`$ decays, although our calculations can readily be applied to semileptonic decays of $`B`$ mesons once more experimental information becomes available on excited $`B`$ meson resonances. In our numerical analysis we use available experimental information and theoretical predictions on charm meson resonances. Particularly for the $`D_s`$ axial resonance we use the mass $`m_{D_{sJ}(2460)}=2.459\mathrm{GeV}`$, while for the first orbital excitation of the $`D^{}`$ meson, we use the mass of $`m_{D_1(2420)}=2.422\mathrm{GeV}`$ Eidelman et al. (2004). For the radially excited vector resonance we then have the Selex $`D_{sJ}^+(2632)`$ state with mass $`m_{D_{sJ}^+(2632)}=2.632\mathrm{GeV}`$ Evdokimov et al. (2004). It is important to note, however, that so-far the Selex discovery has not been confirmed by any other searches Aubert et al. (2004). In $`D`$ decays, the situation is similarly ambiguous. Although the vector $`D^{{}_{}{}^{}}`$ resonance was discovered by Delphi Abreu et al. (1998) with a mass of $`m_D^{{}_{}{}^{}}=2.637\mathrm{GeV}`$ and spin-parity $`1^{}`$, its existence was not confirmed by other searches Rodriguez (1999); Abbiendi et al. (2001). On the other hand recent theoretical studies Di Pierro and Eichten (2001); Vijande et al. (2003) indicate that both radially excited vector states of $`D`$ as well as $`D_s`$ should have slightly larger masses of $`m_D^{{}_{}{}^{}}2.7\mathrm{GeV}`$ and $`m_{D_s^{{}_{}{}^{}}}2.8\mathrm{GeV}`$ Di Pierro and Eichten (2001). We use these theoretically predicted values in our analysis as well as for the radially excited pseudoscalar states, for which currently no experimental indications exist.
In our calculations we use for the heavy meson weak current coupling $`\alpha =f_H\sqrt{m_H}`$ Wise (1992); Becirevic et al. (2003), which we derive from the lattice QCD value of $`f_D=0.225\mathrm{GeV}`$ Wingate (2005) and experimental $`D`$ meson mass $`m_D=1.87\mathrm{GeV}`$ Eidelman et al. (2004) yielding $`\alpha =0.31\mathrm{GeV}^{3/2}`$. The $`\lambda `$ coupling was usually Casalbuoni et al. (1997); Cheng et al. (2005) determined from the value of $`V(0)`$. However, this derivation employed a single pole ansatz for the shape of $`V(q^2)`$. One can instead use data on $`D^{}D\gamma `$ radiative decays. Following discussion in Refs. Prelovsek (2000); Fajfer and Singer (1997), using the most recent data on $`D^{}`$ radiative and strong decays Eidelman et al. (2004), and accounting for the $`SU(3)`$ flavor symmetry breaking effects, we calculate $`\lambda =0.526`$ Gev<sup>-1</sup>. The coupling $`\beta 0.9`$ has been estimated in Ref. Isola et al. (2003) relying on the assumption that the electromagnetic interactions of the light quark within heavy meson are dominated by the exchange of $`\rho ^0`$, $`\omega `$, $`\varphi `$ vector mesons.
We fix the other free parameters ($`\alpha _1,\alpha _2,\alpha ^{},\stackrel{~}{\alpha },\zeta ,\mu ,\stackrel{~}{\zeta },\stackrel{~}{\mu }`$) appearing in the HM$`\chi `$T Lagrangian and weak currents by comparing our model predictions to known experimental values of branching ratios $`(D^0K^{}\mathrm{}^+\nu )`$, $`(D_s\varphi \mathrm{}^+\nu )`$, $`(D^+\rho ^0\mathrm{}^+\nu )`$, $`(D^+K^0\mathrm{}^+\nu )`$, as well as partial decay width ratios $`\mathrm{\Gamma }_L/\mathrm{\Gamma }_T(D^+K^0\mathrm{}^+\nu )`$ and $`\mathrm{\Gamma }_+/\mathrm{\Gamma }_{}(D^+K^0\mathrm{}^+\nu )`$ Eidelman et al. (2004). In order to compare the results of our approach with experimental values, we calculate the decay rates for polarized final light vector mesons. Using helicity amplitudes $`H_{+,,0}`$ defined in Sec. III and by neglecting the lepton masses we get Bajc et al. (1996):
$$\mathrm{\Gamma }_a=\frac{G_F^2m_H^2|K_{HV}|^2}{96\pi ^3}_0^{y_m^V}ydy|H_a(y)|^2|\stackrel{}{p}_V(y)|,$$
(30)
where $`a=+,,0`$ and
$$y_m^V=\left(1\frac{m_V}{m_H}\right)^2.$$
(31)
The constants $`K_{HV}`$ parametrize the flavor mixing relevant to a particular transition, and are given in Table 1 together with the pole mesons.
The transverse, longitudinal and total decay rates are then given trivially by
$`\mathrm{\Gamma }_T`$ $`=`$ $`\mathrm{\Gamma }_++\mathrm{\Gamma }_{},`$
$`\mathrm{\Gamma }_L`$ $`=`$ $`\mathrm{\Gamma }_0,`$
$`\mathrm{\Gamma }_{}`$ $`=`$ $`\mathrm{\Gamma }_T+\mathrm{\Gamma }_L.`$ (32)
Consequently, the $`A_0`$ form factor does not contribute to any decay rate in this approximation and we can not fix the parameters $`\alpha _2`$ and $`\stackrel{~}{\zeta }`$ solely from comparison with experiment. Although $`A_0`$ actually does contribute indirectly through the relation (2) at $`q^2=0`$ as manifested by Eq. (17), this constraint is not automatically satisfied by our model. On the other hand, we can still enforce it ”by hand” after the extrapolation to $`q^2=0`$ to obtain some information on these parameters. Due to the specific combinations in which the parameters appear in Eqs. (LABEL:eq\_ff\_HMcT) we are further restrained to determining only the products $`\stackrel{~}{\alpha }\stackrel{~}{\mu }`$, $`\alpha ^{}\zeta `$ and $`\alpha ^{}\mu `$ using this kind of analysis. Lastly, since the only relevant contribution of $`\alpha _1`$ is to the $`A_1`$ form factor, we cannot disentangle it from the influence of $`\alpha ^{}\zeta `$. Yet again we can impose the large energy limit relation (10) to extract both values independently.
We calculate the result for $`\stackrel{~}{\alpha }\stackrel{~}{\mu }`$, $`\alpha ^{}\zeta `$, $`\alpha ^{}\mu `$ and $`\alpha _1`$ by a weighted average of values obtained from all the measured decay rates and their ratios taking into account for the experimental uncertainties. Furthermore, the values of $`\alpha _1`$ and $`\alpha ^{}\zeta `$ are extracted separately by minimizing the fit function $`(V(0)\xi A_1(0))^2/(V(0)\xi +A_1(0))^2`$. Both minimizations are performed in parallel and the global minimum is sought on the hypercube of dimensions $`[1,1]^4`$ in the hyperspace of the fitted parameters. At the end we obtain the following values of parameters:
$`\stackrel{~}{\alpha }\stackrel{~}{\mu }`$ $`=`$ $`0.090\mathrm{GeV}^{1/2}`$
$`\alpha ^{}\zeta `$ $`=`$ $`0.038\mathrm{GeV}^{3/2}`$
$`\alpha ^{}\mu `$ $`=`$ $`0.066\mathrm{GeV}^{1/2}`$
$`\alpha _1`$ $`=`$ $`0.128\mathrm{GeV}^{1/2}`$ (33)
These values qualitatively agree with the analysis done in Ref. Casalbuoni et al. (1997) using a combination of quark model predictions and single pole experimental fits for all the form factors.
We next use these values in relation (2) to extract information on the the parameters $`\alpha _2`$ and $`\stackrel{~}{\zeta }`$. From Eqs. (LABEL:eq\_ff\_HMcT) it is easy to see that the solutions lie on a straight line in the $`\alpha _2\times \stackrel{~}{\alpha }\stackrel{~}{\zeta }`$ plane. We draw these for the various decay channels used in our analysis in Fig. 1.
We can see that all the decay channels considered fit approximately the same solution in the plane. Consequently, we can use any point on the approximate solution line to obtain the same prediction for the $`q^2`$ dependence of the $`A_0`$ form factor.
It is important to note at this point that due to a high degree of interplay of the various HM$`\chi `$T parameters in the model predictions used in the fit, the values of the new model parameters obtained in such a way are very volatile to changes in the other inputs to the fit. Furthermore these are tree level parameter values and may in addition be very sensitive to chiral and $`1/m_H`$ corrections. Therefore their stated values should be taken cum grano salis. However more importantly, the form factor, branching ratio and polarization width ratio predictions based on this approach are more robust since they are insensitive to particular combinations of parameter values used, as long as they fit the experimental data. We estimate that chiral and heavy quark symmetry breaking corrections could still modify these predictions by as much as $`30\%`$.
We are now ready to draw the $`q^2`$ dependence of all the form factors for the $`D^0K^{}`$, $`D^0\rho ^{}`$ and $`D_s\varphi `$ transitions. The results are depicted in Figs. 23, and 4.
Our predictions for the shapes of the various form factors can also be summarized using the general formulas
$`V(q^2)`$ $`=`$ $`{\displaystyle \frac{V(0)}{(1x)(1ax)}},`$
$`A_0(q^2)`$ $`=`$ $`{\displaystyle \frac{A_0(0)}{(1y)(1a^{}y)}},`$
$`A_1(q^2)`$ $`=`$ $`{\displaystyle \frac{A_1(0)}{1b^{}x}},`$
$`A_2(q^2)`$ $`=`$ $`{\displaystyle \frac{A_2(0)}{(1b^{}x)(1b^{\prime \prime }x)}},`$ (34)
where as before $`x=q^2/m_H^{}^2`$ and $`y=q^2/m_{H_P}^2`$. These expressions are actually simplifications of the form factor parametrizations (9), (12), (11) and (13) respectively. The parameters $`V(0)`$, $`A_0(0)`$, $`A_1(0)`$, $`A_2(0)`$, $`a`$, $`a^{}`$, $`b^{}`$ and $`b^{\prime \prime }`$, which we fix by nearest resonance saturation approximation and HM$`\chi `$T calculation at $`q_{\mathrm{max}}^2`$, are listed in Table 2 for the various decay channels considered.
Finally, we calculate branching ratios and partial decay width ratios for all relevant $`DV\mathrm{}\nu _{\mathrm{}}`$ decays. They are listed in Table 3 together with known experimentally measured values where we have marked those used in the fit of our model parameters.
## VII Conclusion
We have investigated the $`q^2`$ dependence of the heavy to light form factors present in the $`DV\mathrm{}\nu _{\mathrm{}}`$ decays. First we have devised a general parametrization of the $`HV`$ form factors and then used it to compute the $`DV`$ form factors $`q^2`$ dependence in the framework of HM$`\chi `$T. Although we restrict our discussion to $`D`$ decays, the implications of the general form factor parametrization for the semileptonic decays of $`B`$ mesons are obvious, while also our HM$`\chi `$T calculations can easily be applied to these decays once more experimental data becomes available on excited $`B`$ meson resonances. Furthermore, the extension of our parametrization to tensor current form factors is trivial using HQET form factor relations Isgur and Wise (1990); Burdman and Donoghue (1991), which were found to hold even beyond the leading order in heavy quark mass expansion Grinstein and Pirjol (2002, 2004).
In our calculations we do not include chiral corrections and $`1/m_H`$ corrections which might be important in charm decays as we already discussed in Fajfer and Kamenik (2005). As shown in Ref. Boyd and Grinstein (1995) some of these corrections can in fact be absorbed into the leading order parameters while a sizable number of additional parameters remains undetermined. Since our approach is based on a global fit to existing experimental data on $`DV`$ semileptonic decays it does not seem possible at present to disentangle the effects of these new parameters and fix their values. Eventually, once more experimental results become available, it will be possible to learn more about some combinations of these parameters. We note however that while the values of HM$`\chi `$T parameters which are used and obtained in the fit of our model predictions may indeed be affected by chiral and $`1/m_H`$ corrections, the predicted shapes of the form factors are very robust since they are fixed by the underlying general form factor parametrization proposed in the text and by the masses of the involved resonances. Only the overall size of the individual form factors is determined by HM$`\chi `$T calculation at $`q_{\mathrm{max}}^2`$ and could thus be sensitive to chiral and $`1/m_H`$ corrections as well to the input parameters of the experimental fit.
The presence of charm meson resonances in our Lagrangian affects the values of the form factors at $`q_{\mathrm{max}}^2`$ and induces saturation of the second poles in the parametrizations of the $`V(q^2)`$ and $`A_0(q^2)`$ form factors by the next radial excitations of $`D_{(s)}^{}`$ and $`D_{(s)}`$ mesons respectively. The single pole $`q^2`$ behavior of the $`A_1(q^2)`$ form factor is explained by the presence of a single $`1^+`$ state relevant to each decay, while in $`A_2(q^2)`$ in addition to these states one might also account for their next radial excitations. However, due to the lack of data on their presence we assume their masses being much higher than the first $`1^+`$ states and we neglect their effects.
We point out that the single pole parametrization of all the form factors used in previous experimental studies cannot correctly satisfy known HQET and large energy limit constraints. In LEET Charles et al. (1999) and SCET Hill (2004) studies the ratio of $`V(0)`$ and $`A_1(0)`$ obtained for $`B\rho \mathrm{}\nu _{\mathrm{}}`$ within these approaches was compared with experimental results for $`DK^{}\mathrm{}\nu _{\mathrm{}}`$ which were obtained assuming such single pole dependence.
In the current literature there exist a number of studies of the $`q^2`$ shape of the $`DV`$ form factors. In the lattice simulation of Ref. Abada et al. (2003) the scaling behavior of the form factors has been properly included, but the pole/dipole fits of the form factors they used provide no physical information on the $`q^2`$ dependence of the form factors beyond the leading poles. However, our results for the $`DK^{}\mathrm{}\nu _{\mathrm{}}`$ form factors are in good agreement with theirs. Similarly, the fits done in the quark model calculation of Melikhov and Stech (2000) to their particular parametrization do not differ a lot from our $`q^2`$ behavior of the form factors. On the other hand, the authors of Ball et al. (1991) have used QCD sum rules to derive the $`q^2`$ dependence of the form factors. As presented on Fig. 3 and in Table 2 our results are in good agreement with theirs on values of $`V(0)`$ and $`A_1(0)`$ while our results for the $`A_2(0)`$ are somewhat lower. In Ref. Ball (1993) the values of the form factors appearing in $`D\rho \mathrm{}\nu _{\mathrm{}}`$ decay have been investigated in the same approach and they agree well with results of our model calculation and extrapolation.
We hope that the ongoing experimental studies will help to shed more light on the shapes of the $`DV`$ form factors.
###### Acknowledgements.
We are greatly indebted to Damir Bećirević who originally initiated this work and provided useful insight and advice throughout the process of this study. S. F. thanks Alexander von Humboldt foundation for financial support and A. J. Buras for his warm hospitality during her stay at the Physik Department, TU München, where part of this work has been done. This work is supported in part by the Ministry of Higher Education, Science and Technology of the Republic of Slovenia. |
warning/0506/cond-mat0506114.html | ar5iv | text | # Enhancement of superconducting transition temperature by the additional second neighbor hopping 𝑡' in the 𝑡-𝐽 model
\[
## Abstract
Within the kinetic energy driven superconducting mechanism, the effect of the additional second neighbor hopping $`t^{}`$ on the superconducting state of the $`t`$-$`J`$ model is discussed. It is shown that $`t^{}`$ plays an important role in enhancing the superconducting transition temperature of the $`t`$-$`J`$ model. It is also shown that the superconducting-state of cuprate superconductors is the conventional Bardeen-Cooper-Schrieffer like, so that the basic Bardeen-Cooper-Schrieffer formalism is still valid in quantitatively reproducing the doping dependence of the superconducting gap parameter and superconducting transition temperature, and electron spectral function at ($`\pi `$,0) point, although the pairing mechanism is driven by the kinetic energy by exchanging dressed spin excitations.
\]
After intensive investigations over more than a decade, it has now become clear that although the physical properties of cuprate superconductors in the normal-state are fundamentally different from these of the conventional metals , the superconducting (SC)-state of cuprate superconductors is still associated with the formation of the electron Cooper pairs as in the conventional superconductors. In the conventional metals, superconductivity results when electrons pair up into Cooper pairs, which is mediated by the interaction of electrons with phonons . As a result, the pairing in the conventional superconductors is always related with an increase in kinetic energy which is overcompensated by the lowering of potential energy . However, it has been argued that the form of the electron Cooper pairs is determined by the need to reduce the frustrated kinetic energy in doped cuprates , i.e., the strong frustration of the kinetic energy in the normal-state is partially relieved upon entering the SC-state. By virtue of systematic studies using the nuclear magnetic resonance, and muon spin rotation techniques, particularly the inelastic neutron scattering, it has been well established that the antiferromagnetic (AF) short-range correlation (AFSRC) coexists with the SC-state in the whole SC regime , which provide a clear link between the SC pairing mechanism and magnetic excitations. Moreover, it has been shown that although the SC pairing mechanism of cuprate superconductors is beyond the conventional electron-phonon mechanism, the SC-state is the conventional Bardeen-Cooper-Schrieffer (BCS) like , so that the basic BCS formalism is still valid in discussions of the electron spectral properties .
Very soon after the discovery of superconductivity in doped cuprates, Anderson suggested that the essential physics of doped cuprates is contained in the $`t`$-$`J`$ model on a square lattice. This followed from the experiments that cuprate superconductors are doped antiferromagnets, where the common features are the presence of the square lattice CuO<sub>2</sub> planes and a similar phase diagram as a function of the doping concentration . Since then much effort has concentrated on the unusual normal-state and SC mechanism within the $`t`$-$`J`$ model . Based on the charge-spin separation (CSS) fermion-spin theory , we have developed a kinetic energy driven SC mechanism within the $`t`$-$`J`$ model. It is shown that the dressed holons interact occurring directly through the kinetic energy by exchanging the spin excitations, leading to a net attractive force between the dressed holons, then the electron Cooper pairs originating from the dressed holon pairing state are due to the charge-spin recombination, and their condensation reveals the SC ground-state. This SC-state is controlled by both SC gap function and quasiparticle coherence, and the maximal SC transition temperature occurs around the optimal doping, then decreases in both underdoped and overdoped regimes . However, the simple $`t`$-$`J`$ model can not be regarded as a comprehensive model for the quantitative comparison with cuprate superconductors. It has been shown from the angle resolved photoemission spectroscopy (ARPES) experiments that although the highest energy filled electron band is well described by the $`t`$-$`J`$ model in the direction between the $`[0,0]`$ point and the $`[\pi ,\pi ]`$ point in the momentum space, but both experimental data near $`[\pi ,0]`$ point and overall dispersion may be properly accounted by generalizing the $`t`$-$`J`$ model to include the second- and third-nearest neighbors hopping terms $`t^{}`$ and $`t^{\prime \prime }`$. Moreover, the experimental analysis shows that the SC transition temperature for different families of cuprate superconductors is strongly correlated with $`t^{}`$. In this Letter, we discuss the effect of the additional second neighbor hopping $`t^{}`$ on the SC-state of the $`t`$-$`J`$ model within the framework of the kinetic energy driven SC mechanism . Our result shows that the SC-state of cuprate superconductors is the conventional BCS like , so that the basic BCS formalism is still valid in quantitatively reproducing the doping dependence of the effective SC gap parameter and SC transition temperature, and electron spectral function at $`[\pi ,0]`$ point, although the pairing mechanism is driven by the kinetic energy by exchanging dressed spin excitations, and other exotic magnetic properties are beyond the BCS theory. Our result also shows that the additional second neighbor hopping $`t^{}`$ plays an important role in enhancing the SC transition temperature of the $`t`$-$`J`$ model and in determining the correct position of the SC quasiparticle peak of the electron spectral function at $`[\pi ,0]`$ point.
We start from the $`t`$-$`t^{}`$-$`J`$ model on a square lattice ,
$`H`$ $`=`$ $`t{\displaystyle \underset{i\widehat{\eta }\sigma }{}}C_{i\sigma }^{}C_{i+\widehat{\eta }\sigma }+t^{}{\displaystyle \underset{i\widehat{\tau }\sigma }{}}C_{i\sigma }^{}C_{i+\widehat{\tau }\sigma }+\mu {\displaystyle \underset{i\sigma }{}}C_{i\sigma }^{}C_{i\sigma }`$ (1)
$`+`$ $`J{\displaystyle \underset{i\widehat{\eta }}{}}𝐒_i𝐒_{i+\widehat{\eta }},`$ (2)
supplemented by the local constraint $`_\sigma C_{i\sigma }^{}C_{i\sigma }1`$ to avoid the double occupancy, where $`\widehat{\eta }=\pm \widehat{x},\pm \widehat{y}`$, $`\widehat{\tau }=\pm \widehat{x}\pm \widehat{y}`$, $`C_{i\sigma }^{}`$ ($`C_{i\sigma }`$) is the electron creation (annihilation) operator, $`𝐒_i=C_i^{}\stackrel{}{\sigma }C_i/2`$ is spin operator with $`\stackrel{}{\sigma }=(\sigma _x,\sigma _y,\sigma _z)`$ as Pauli matrices, and $`\mu `$ is the chemical potential. The strong electron correlation in the $`t`$-$`t^{}`$-$`J`$ model manifests itself by the electron single occupancy local constraint , which can be treated properly in analytical calculations within the CSS fermion-spin theory , where the constrained electron operators are decoupled as $`C_i=h_i^{}S_i^{}`$ and $`C_i=h_i^{}S_i^+`$, with the spinful fermion operator $`h_{i\sigma }=e^{i\mathrm{\Phi }_{i\sigma }}h_i`$ describes the charge degree of freedom together with some effects of the spin configuration rearrangements due to the presence of the doped hole itself (dressed holon), while the spin operator $`S_i`$ describes the spin degree of freedom (dressed spin), then the electron local constraint for the single occupancy, $`_\sigma C_{i\sigma }^{}C_{i\sigma }=S_i^+h_ih_i^{}S_i^{}+S_i^{}h_ih_i^{}S_i^+=h_ih_i^{}(S_i^+S_i^{}+S_i^{}S_i^+)=1h_i^{}h_i1`$, is satisfied in analytical calculations. It has been shown that these dressed holon and spin are gauge invariant , and in this sense, they are real and can be interpreted as the physical excitations . Although in common sense $`h_{i\sigma }`$ is not a real spinful fermion, it behaves like a spinful fermion. In this CSS fermion-spin representation, the low-energy behavior of the $`t`$-$`t^{}`$-$`J`$ model (1) can be expressed as,
$`H`$ $`=`$ $`t{\displaystyle \underset{i\widehat{\eta }}{}}(h_iS_i^+h_{i+\widehat{\eta }}^{}S_{i+\widehat{\eta }}^{}+h_iS_i^{}h_{i+\widehat{\eta }}^{}S_{i+\widehat{\eta }}^+)`$ (3)
$`+`$ $`t^{}{\displaystyle \underset{i\widehat{\tau }}{}}(h_iS_i^+h_{i+\widehat{\tau }}^{}S_{i+\widehat{\tau }}^{}+h_iS_i^{}h_{i+\widehat{\tau }}^{}S_{i+\widehat{\tau }}^+)`$ (4)
$``$ $`\mu {\displaystyle \underset{i\sigma }{}}h_{i\sigma }^{}h_{i\sigma }+J_{\mathrm{eff}}{\displaystyle \underset{i\widehat{\eta }}{}}𝐒_i𝐒_{i+\widehat{\eta }},`$ (5)
with $`J_{\mathrm{eff}}=(1x)^2J`$, and $`x=h_{i\sigma }^{}h_{i\sigma }=h_i^{}h_i`$ is the hole doping concentration. As a consequence, the kinetic energy terms in the $`t`$-$`t^{}`$-$`J`$ model have been expressed as the dressed holon-spin interactions, which reflects that even the kinetic energy terms in the $`t`$-$`t^{}`$-$`J`$ Hamiltonian have strong Coulombic contributions due to the restriction of no doubly occupancy of a given site, and therefore dominate the essential physics of doped cuprates.
ARPES measurements show that in the real space the gap function and pairing force have a range of one lattice spacing, which indicates that the order parameter for the electron Cooper pair can be expressed as,
$`\mathrm{\Delta }`$ $`=`$ $`C_i^{}C_{i+\widehat{\eta }}^{}C_i^{}C_{i+\widehat{\eta }}^{}`$ (6)
$`=`$ $`h_ih_{i+\widehat{\eta }}S_i^+S_{i+\widehat{\eta }}^{}h_ih_{i+\widehat{\eta }}S_i^{}S_{i+\widehat{\eta }}^+.`$ (7)
In the doped regime without the AF long-range order (AFLRO), the dressed spins form a disordered spin liquid state, where the dressed spin correlation function $`S_i^+S_{i+\widehat{\eta }}^{}=S_i^{}S_{i+\widehat{\eta }}^+`$, then the order parameter for the electron Cooper pair in Eq. (3) can be written as $`\mathrm{\Delta }=S_i^+S_{i+\widehat{\eta }}^{}\mathrm{\Delta }_h`$, with the dressed holon pairing order parameter $`\mathrm{\Delta }_h=h_{i+\widehat{\eta }}h_ih_{i+\widehat{\eta }}h_i`$, which shows that the SC order parameter of the electron Cooper pair is related to the dressed holon pairing amplitude, and is proportional to the number of doped holes, and not to the number of electrons. However, in the extreme low doped regime with AFLRO, where the dressed spin correlation function $`S_i^+S_{i+\widehat{\eta }}^{}S_i^{}S_{i+\widehat{\eta }}^+`$, then the conduct is disrupted by AFLRO, and therefore there is no mixing of superconductivity and AFLRO . In the case without AFLRO, we have shown within the Eliashberg’s strong coupling theory that the dressed holon-spin interaction can induce the dressed holon pairing state (then the electron Cooper pairing state) by exchanging dressed spin excitations in the higher power of the hole doping concentration. Following our previous discussions based on the $`t`$-$`J`$ model , the self-consistent equations that satisfied by the full dressed holon diagonal and off-diagonal Green’s functions in the present $`t`$-$`t^{}`$-$`J`$ model are obtained as,
$`g(k)`$ $`=`$ $`g^{(0)}(k)`$ (9)
$`+`$ $`g^{(0)}(k)[\mathrm{\Sigma }_1^{(h)}(k)g(k)\mathrm{\Sigma }_2^{(h)}(k)\mathrm{}^{}(k)],`$ (10)
$`\mathrm{}^{}(k)`$ $`=`$ $`g^{(0)}(k)[\mathrm{\Sigma }_1^{(h)}(k)\mathrm{}^{}(k)+\mathrm{\Sigma }_2^{(h)}(k)g(k)],`$ (11)
respectively, where the four-vector notation $`k=(𝐤,i\omega _n)`$, and the mean-field (MF) dressed holon diagonal Green’s function $`g^{(0)1}(k)=i\omega _n\xi _𝐤`$, with the MF dressed holon excitation spectrum $`\xi _𝐤=Zt\chi _1\gamma _𝐤Zt^{}\chi _2\gamma _𝐤^{}\mu `$, where $`\gamma _𝐤=(1/Z)_{\widehat{\eta }}e^{i𝐤\widehat{\eta }}`$, $`\gamma _𝐤^{}=(1/Z)_{\widehat{\tau }}e^{i𝐤\widehat{\tau }}`$, $`Z`$ is the number of the nearest neighbor or second-nearest neighbor sites, the spin correlation functions $`\chi _1=S_i^+S_{i+\widehat{\eta }}^{}`$ and $`\chi _2=S_i^+S_{i+\widehat{\tau }}^{}`$, while the dressed holon self-energies are obtained from the dressed spin bubble as,
$`\mathrm{\Sigma }_1^{(h)}(k)`$ $`=`$ $`{\displaystyle \frac{1}{N^2}}{\displaystyle \underset{𝐩,𝐩^{}}{}}(Zt\gamma _{𝐩+𝐩^{}+𝐤}Zt^{}\gamma _{𝐩+𝐩^{}+𝐤}^{})^2`$ (13)
$`\times `$ $`{\displaystyle \frac{1}{\beta }}{\displaystyle \underset{ip_m}{}}g(p+k){\displaystyle \frac{1}{\beta }}{\displaystyle \underset{ip_m^{}}{}}D^{(0)}(p^{})D^{(0)}(p^{}+p),`$ (14)
$`\mathrm{\Sigma }_2^{(h)}(k)`$ $`=`$ $`{\displaystyle \frac{1}{N^2}}{\displaystyle \underset{𝐩,𝐩^{}}{}}(Zt\gamma _{𝐩+𝐩^{}+𝐤}Zt^{}\gamma _{𝐩+𝐩^{}+𝐤}^{})^2`$ (15)
$`\times `$ $`{\displaystyle \frac{1}{\beta }}{\displaystyle \underset{ip_m}{}}\mathrm{}(pk){\displaystyle \frac{1}{\beta }}{\displaystyle \underset{ip_m^{}}{}}D^{(0)}(p^{})D^{(0)}(p^{}+p),`$ (16)
where $`p=(𝐩,ip_m)`$, $`p^{}=(𝐩^{},ip_m^{})`$, $`N`$ is the number of sites, and the MF dressed spin Green’s function , $`D^{(0)1}(p)=[(ip_m)^2\omega _𝐩^2]/B_𝐩`$, with $`B_𝐩=2\lambda _1(A_1\gamma _𝐩A_2)\lambda _2(2\chi _2^z\gamma _𝐩^{}\chi _2)`$, $`\lambda _1=2ZJ_{eff}`$, $`\lambda _2=4Z\varphi _2t^{}`$, $`A_1=ϵ\chi _1^z+\chi _1/2`$, $`A_2=\chi _1^z+ϵ\chi _1/2`$, $`ϵ=1+2t\varphi _1/J_{\mathrm{eff}}`$, the dressed holon’s particle-hole parameters $`\varphi _1=h_{i\sigma }^{}h_{i+\widehat{\eta }\sigma }`$ and $`\varphi _2=h_{i\sigma }^{}h_{i+\widehat{\tau }\sigma }`$, the spin correlation functions $`\chi _1^z=S_i^zS_{i+\widehat{\eta }}^z`$ and $`\chi _2^z=S_i^zS_{i+\widehat{\tau }}^z`$, and the MF dressed spin excitation spectrum,
$`\omega _𝐩^2`$ $`=`$ $`\lambda _1^2[(A_4\alpha ϵ\chi _1^z\gamma _𝐩{\displaystyle \frac{1}{2Z}}\alpha ϵ\chi _1)(1ϵ\gamma _𝐩)`$ (17)
$`+`$ $`{\displaystyle \frac{1}{2}}ϵ(A_3{\displaystyle \frac{1}{2}}\alpha \chi _1^z\alpha \chi _1\gamma _𝐩)(ϵ\gamma _𝐩)]`$ (18)
$`+`$ $`\lambda _2^2[\alpha (\chi _2^z\gamma _𝐩^{}{\displaystyle \frac{3}{2Z}}\chi _2)\gamma _𝐩^{}+{\displaystyle \frac{1}{2}}(A_5{\displaystyle \frac{1}{2}}\alpha \chi _2^z)]`$ (19)
$`+`$ $`\lambda _1\lambda _2[\alpha \chi _1^z(1ϵ\gamma _𝐩)\gamma _𝐩^{}+{\displaystyle \frac{1}{2}}\alpha (\chi _1\gamma _𝐩^{}C_3)(ϵ\gamma _𝐩)`$ (20)
$`+`$ $`\alpha \gamma _𝐩^{}(C_3^zϵ\chi _2^z\gamma _𝐩){\displaystyle \frac{1}{2}}\alpha ϵ(C_3\chi _2\gamma _𝐩)],`$ (21)
with $`A_3=\alpha C_1+(1\alpha )/(2Z)`$, $`A_4=\alpha C_1^z+(1\alpha )/(4Z)`$, $`A_5=\alpha C_2+(1\alpha )/(2Z)`$, and the spin correlation functions $`C_1=(1/Z^2)_{\widehat{\eta },\widehat{\eta ^{}}}S_{i+\widehat{\eta }}^+S_{i+\widehat{\eta ^{}}}^{}`$, $`C_1^z=(1/Z^2)_{\widehat{\eta },\widehat{\eta ^{}}}S_{i+\widehat{\eta }}^zS_{i+\widehat{\eta ^{}}}^z`$, $`C_2=(1/Z^2)_{\widehat{\tau },\widehat{\tau ^{}}}S_{i+\widehat{\tau }}^+S_{i+\widehat{\tau ^{}}}^{}`$, $`C_3=(1/Z)_{\widehat{\tau }}S_{i+\widehat{\eta }}^+S_{i+\widehat{\tau }}^{}`$, and $`C_3^z=(1/Z)_{\widehat{\tau }}S_{i+\widehat{\eta }}^zS_{i+\widehat{\tau }}^z`$. In order to satisfy the sum rule of the correlation function $`S_i^+S_i^{}=1/2`$ in the case without AFLRO, the important decoupling parameter $`\alpha `$ has been introduced in the MF calculation , which can be regarded as the vertex correction.
The self-energy function $`\mathrm{\Sigma }_2^{(h)}(k)`$ describes the effective dressed holon gap function, since both doping and temperature dependence of the pairing force and dressed holon gap function have been incorporated into $`\mathrm{\Sigma }_2^{(h)}(k)`$, while the self-energy function $`\mathrm{\Sigma }_1^{(h)}(k)`$ renormalizes the MF dressed holon spectrum, and therefore it describes the quasiparticle coherence. Moreover, $`\mathrm{\Sigma }_2^{(h)}(k)`$ is an even function of $`i\omega _n`$, while $`\mathrm{\Sigma }_1^{(h)}(k)`$ is not. For the convenience, $`\mathrm{\Sigma }_1^{(h)}(k)`$ can be broken up into its symmetric and antisymmetric parts as, $`\mathrm{\Sigma }_1^{(h)}(k)=\mathrm{\Sigma }_{1e}^{(h)}(k)+i\omega _n\mathrm{\Sigma }_{1o}^{(h)}(k)`$, then both $`\mathrm{\Sigma }_{1e}^{(h)}(k)`$ and $`\mathrm{\Sigma }_{1o}^{(h)}(k)`$ are even functions of $`i\omega _n`$. In this case, the quasiparticle coherent weight can be defined as $`Z_F^1(k)=1\mathrm{\Sigma }_{1o}^{(h)}(k)`$. As in the conventional superconductor , the retarded function $`\mathrm{Re}\mathrm{\Sigma }_{1e}^{(h)}(k)`$ is a constant, independent of ($`𝐤,\omega `$), and it just renormalizes the chemical potential, therefore it can be dropped. Furthermore, we only study the static limit of the effective dressed holon gap function and quasiparticle coherent weight, i.e., $`\mathrm{\Sigma }_2^{(h)}(k)=\overline{\mathrm{\Delta }}_h(𝐤)`$, and $`Z_F^1(𝐤)=1\mathrm{\Sigma }_{1o}^{(h)}(𝐤)`$. Although $`Z_F(𝐤)`$ still is a function of $`𝐤`$, the wave vector dependence is unimportant, since everything happens at the electron Fermi surface. As in the previous discussions within the $`t`$-$`J`$ model , the special wave vector can be estimated qualitatively from the electron momentum distribution as $`𝐤_0=𝐤_𝐀𝐤_𝐅`$ with $`𝐤_𝐀=[\pi ,\pi ]`$ and $`𝐤_𝐅[(1x)\pi /2,(1x)\pi /2]`$, which guarantees $`Z_F=Z_F(𝐤_0)`$ near the electron Fermi surface. In this case, the dressed holon diagonal and off-diagonal Green’s functions in Eqs. (4a) and (4b) can be expressed explicitly as,
$`g(k)`$ $`=`$ $`Z_F{\displaystyle \frac{U_{h𝐤}^2}{i\omega _nE_{h𝐤}}}+Z_F{\displaystyle \frac{V_{h𝐤}^2}{i\omega _n+E_{h𝐤}}},`$ (23)
$`\mathrm{}^{}(k)`$ $`=`$ $`Z_F{\displaystyle \frac{\overline{\mathrm{\Delta }}_{hZ}(𝐤)}{2E_{h𝐤}}}\left({\displaystyle \frac{1}{i\omega _nE_{h𝐤}}}{\displaystyle \frac{1}{i\omega _n+E_{h𝐤}}}\right),`$ (24)
with the dressed holon quasiparticle coherence factors $`U_{h𝐤}^2=(1+\overline{\xi _𝐤}/E_{h𝐤})/2`$ and $`V_{h𝐤}^2=(1\overline{\xi _𝐤}/E_{h𝐤})/2`$, $`\overline{\xi _𝐤}=Z_F\xi _𝐤`$, $`\overline{\mathrm{\Delta }}_{hZ}(𝐤)=Z_F\overline{\mathrm{\Delta }}_h(𝐤)`$, and the dressed holon quasiparticle spectrum $`E_{h𝐤}=\sqrt{\overline{\xi _𝐤^2}+\overline{\mathrm{\Delta }}_{hZ}(𝐤)^2}`$.
Experimentally, some results seem consistent with an s-wave pairing , while other measurements gave the evidence in favor of the d-wave pairing . These experiments reflect a fact that the d-wave gap function $`\mathrm{k}_x^2\mathrm{k}_y^2`$ belongs to the same representation $`\mathrm{\Gamma }_1`$ of the orthorhombic crystal group as does s-wave gap function $`\mathrm{k}_x^2+\mathrm{k}_y^2`$. Within the $`t`$-$`J`$ model, we have shown that the electron Cooper pairs have a dominated d-wave symmetry over a wide range of the doping concentration, around the optimal doping. To make the discussion simpler, we only consider the d-wave case, i.e., $`\overline{\mathrm{\Delta }}_{hZ}(𝐤)=\overline{\mathrm{\Delta }}_{hZ}\gamma _𝐤^{(d)}`$, with $`\gamma _𝐤^{(d)}=(\mathrm{cos}k_x\mathrm{cos}k_y)/2`$. In this case, the dressed holon effective gap parameter and quasiparticle coherent weight in Eqs. (5a) and (5b) satisfy following two equations,
$`1`$ $`={\displaystyle \frac{1}{N^3}}{\displaystyle \underset{𝐤,𝐪,𝐩}{}}(Zt\gamma _{𝐤+𝐪}Zt^{}\gamma _{𝐤+𝐪}^{})^2\gamma _{𝐤𝐩+𝐪}^{(d)}\gamma _𝐤^{(d)}{\displaystyle \frac{Z_F^2}{E_{h𝐤}}}{\displaystyle \frac{B_𝐪B_𝐩}{\omega _𝐪\omega _𝐩}}`$ (26)
$`\times `$ $`\left({\displaystyle \frac{F_1^{(1)}(𝐤,𝐪,𝐩)}{(\omega _𝐩\omega _𝐪)^2E_{h𝐤}^2}}{\displaystyle \frac{F_1^{(2)}(𝐤,𝐪,𝐩)}{(\omega _𝐩+\omega _𝐪)^2E_{h𝐤}^2}}\right),`$ (27)
$`Z_F^1`$ $`=1+{\displaystyle \frac{1}{N^2}}{\displaystyle \underset{𝐪,𝐩}{}}(Zt\gamma _{𝐩+𝐤_\mathrm{𝟎}}Zt^{}\gamma _{𝐩+𝐤_\mathrm{𝟎}}^{})^2Z_F{\displaystyle \frac{B_𝐪B_𝐩}{4\omega _𝐪\omega _𝐩}}`$ (28)
$`\times `$ $`({\displaystyle \frac{F_2^{(1)}(𝐪,𝐩)}{(\omega _𝐩\omega _𝐪E_{h𝐩𝐪+𝐤_\mathrm{𝟎}})^2}}+{\displaystyle \frac{F_2^{(2)}(𝐪,𝐩)}{(\omega _𝐩\omega _𝐪+E_{h𝐩𝐪+𝐤_\mathrm{𝟎}})^2}}`$ (29)
$`+`$ $`{\displaystyle \frac{F_2^{(3)}(𝐪,𝐩)}{(\omega _𝐩+\omega _𝐪E_{h𝐩𝐪+𝐤_\mathrm{𝟎}})^2}}+{\displaystyle \frac{F_2^{(4)}(𝐪,𝐩)}{(\omega _𝐩+\omega _𝐪+E_{h𝐩𝐪+𝐤_\mathrm{𝟎}})^2}}),`$ (30)
respectively, where $`F_1^{(1)}(𝐤,𝐪,𝐩)=(\omega _𝐩\omega _𝐪)[n_B(\omega _𝐪)n_B(\omega _𝐩)][12n_F(E_{h𝐤})]+E_{h𝐤}[n_B(\omega _𝐩)n_B(\omega _𝐪)+n_B(\omega _𝐪)n_B(\omega _𝐩)]`$, $`F_1^{(2)}(𝐤,𝐪,𝐩)=(\omega _𝐩+\omega _𝐪)[n_B(\omega _𝐩)n_B(\omega _𝐪)][12n_F(E_{h𝐤})]+E_{h𝐤}[n_B(\omega _𝐩)n_B(\omega _𝐪)+n_B(\omega _𝐩)n_B(\omega _𝐪)]`$, $`F_2^{(1)}(𝐪,𝐩)=n_F(E_{h𝐩𝐪+𝐤_\mathrm{𝟎}})[n_B(\omega _𝐪)n_B(\omega _𝐩)]n_B(\omega _𝐩)n_B(\omega _𝐪)`$, $`F_2^{(2)}(𝐪,𝐩)=n_F(E_{h𝐩𝐪+𝐤_\mathrm{𝟎}})[n_B(\omega _𝐩)n_B(\omega _𝐪)]n_B(\omega _𝐪)n_B(\omega _𝐩)`$, $`F_2^{(3)}(𝐪,𝐩)=n_F(E_{h𝐩𝐪+𝐤_\mathrm{𝟎}})[n_B(\omega _𝐪)n_B(\omega _𝐩)]+n_B(\omega _𝐩)n_B(\omega _𝐪)`$, and $`F_2^{(4)}(𝐪,𝐩)=n_F(E_{h𝐩𝐪+𝐤_\mathrm{𝟎}})[n_B(\omega _𝐪)n_B(\omega _𝐩)]+n_B(\omega _𝐩)n_B(\omega _𝐪)`$. These two equations must be solved simultaneously with other self-consistent equations , then all order parameters, decoupling parameter $`\alpha `$, and chemical potential $`\mu `$ are determined by the self-consistent calculation. With the above discussions, we now can obtain the dressed holon pair gap function in terms of the off-diagonal Green’s function (7b) as $`\mathrm{\Delta }_h(𝐤)=(1/\beta )_{i\omega _n}\mathrm{}^{}(𝐤,i\omega _n)`$, then the dressed holon pair order parameter can be evaluated as,
$`\mathrm{\Delta }_h={\displaystyle \frac{2}{N}}{\displaystyle \underset{𝐤}{}}[\gamma _𝐤^{(d)}]^2{\displaystyle \frac{Z_F\overline{\mathrm{\Delta }}_{hZ}}{E_{h𝐤}}}\mathrm{tanh}[{\displaystyle \frac{1}{2}}\beta E_{h𝐤}].`$ (31)
This dressed holon pairing state originating from the kinetic energy terms by exchanging dressed spin excitations also leads to form the electron Cooper pairing state , and the SC gap function is obtained from the electron off-diagonal Green’s function $`\mathrm{\Gamma }^{}(ij,tt^{})=C_i^{}(t);C_j^{}(t^{})`$, which is a convolution of the dressed spin Green’s function and dressed holon off-diagonal Green’s function and reflects the charge-spin recombination . In the present case, this electron off-diagonal Green’s function can be evaluated in terms of the MF dressed spin Green’s function and dressed holon off-diagonal Green’s function (7b) as,
$`\mathrm{\Gamma }^{}(k)`$ $`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{𝐩}{}}Z_F{\displaystyle \frac{\overline{\mathrm{\Delta }}_{hZ}^{(a)}(𝐩+𝐤)}{2E_{h𝐩+𝐤}}}{\displaystyle \frac{B_𝐩}{2\omega _𝐩}}\{F_3^{(1)}(𝐤,𝐩)`$ (32)
$`\times `$ $`\left({\displaystyle \frac{1}{i\omega _nE_{h𝐩+𝐤}\omega _𝐩}}{\displaystyle \frac{1}{i\omega _n+E_{h𝐩+𝐤}+\omega _𝐩}}\right)`$ (33)
$``$ $`F_3^{(2)}(𝐤,𝐩)({\displaystyle \frac{1}{i\omega _n+E_{h𝐩+𝐤}\omega _𝐩}}`$ (34)
$``$ $`{\displaystyle \frac{1}{i\omega _nE_{h𝐩+𝐤}+\omega _𝐩}})\},`$ (35)
with $`F_3^{(1)}(𝐤,𝐩)=1n_F(E_{h𝐤+𝐩})+n_B(\omega _𝐩)`$ and $`F_3^{(2)}(𝐤,𝐩)=n_F(E_{h𝐤+𝐩})+n_B(\omega _𝐩)`$, then the SC gap function is obtained from the above electron off-diagonal Green’s function as,
$`\mathrm{\Delta }(𝐤)`$ $`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{𝐩}{}}{\displaystyle \frac{Z_F\overline{\mathrm{\Delta }}_{Zh}(𝐩𝐤)}{2E_{h𝐩𝐤}}}\mathrm{tanh}[{\displaystyle \frac{1}{2}}\beta E_{h𝐩𝐤}]`$ (36)
$`\times `$ $`{\displaystyle \frac{B_𝐩}{2\omega _𝐩}}\mathrm{coth}[{\displaystyle \frac{1}{2}}\beta \omega _𝐩].`$ (37)
From this SC gap function, the SC gap parameter in Eq. (3) is obtained as $`\mathrm{\Delta }=\chi _1\mathrm{\Delta }_h`$. Since both dressed holon (then electron) pairing gap parameter and pairing interaction in cuprate superconductors are doping dependent, therefore the experimental observed SC gap parameter should be an effective SC gap parameter $`\overline{\mathrm{\Delta }}\chi _1\overline{\mathrm{\Delta }}_h`$, which measures the strength of the binding of electrons into electron Cooper pairs. In Fig. 1, we plot the effective dressed holon pairing (a) and effective SC (b) gap parameters in the d-wave symmetry as a function of the hole doping concentration at $`T=0.002J`$ for $`t/J=2.5`$ and $`t^{}/J=0.3`$ (solid line) and $`t/J=2.5`$ and $`t^{}=0`$ (dashed line). For comparison, the experimental result of the upper critical field as a function of the hole doping concentration is also shown in Fig. 1(b). In a given doping concentration, the upper critical field is defined as the critical field that destroys the SC-state at the zero temperature, therefore the upper critical field also measures the strength of the binding of electrons into Cooper pairs like the effective SC gap parameter . In other words, both effective SC gap parameter and upper critical field have a similar doping dependence . In this sense, our result is in good agreement with the experimental data . Our result also shows that the effect of $`t^{}`$ on the SC-state of the $`t`$-$`J`$ model is to enhance the amplitude of the effective dressed holon (then electron) pairing gap parameter, and shift the maximal value of $`\overline{\mathrm{\Delta }}_h`$ (then $`\overline{\mathrm{\Delta }}`$) towards to the low doping regime. In particular, the value of $`\overline{\mathrm{\Delta }}`$ in the $`t`$-$`t^{}`$-$`J`$ model increases with increasing doping in the underdoped regime, and reaches a maximum in the optimal doping $`x_{\mathrm{opt}}0.15`$, then decreases in the overdoped regime. Since the effective dressed holon pairing gap parameter measures the strength of the binding of dressed holons into dressed holon pairs, then our results also show that although the superconductivity is driven by the kinetic energy by exchanging dressed spin excitations, the strength of the binding of electrons into electron Cooper pairs is still suppressed by AFSRC. Based on the numerical simulations, it has been shown that the SC correlation of the $`t`$-$`J`$ model is enhanced by introducing $`t^{}`$, where the particular correlation between the SC gap and electron occupation at $`[\pi ,0]`$ point is the main reason for enhancement of pairs, which is consistent with our present result. However, their result also shows that the SC correlation becomes strongest shifts to the overdoped regime by introducing $`t^{}`$, and therefore the SC correlation is greatly enhanced in the overdoped regime, which is inconsistent with our present result. The reason for this inconsistency is not clear, and the related issue is under investigation now.
Now we turn to discuss the effect of $`t^{}`$ on the SC transition temperature. As in the case of the $`t`$-$`J`$ model , the SC transition temperature $`T_c`$ occurring in the case of the SC gap parameter $`\mathrm{\Delta }=0`$ in Eq. (11) is identical to the dressed holon pair transition temperature occurring in the case of the effective holon pairing gap parameter $`\overline{\mathrm{\Delta }}_{hZ}=0`$. In this case, we have performed a calculation for the doping dependence of the SC transition temperature, and the result of $`T_c`$ as a function of the hole doping concentration in the d-wave symmetry for $`t/J=2.5`$ and $`t^{}/J=0.3`$ (solid line) and $`t/J=2.5`$ and $`t^{}=0`$ (dashed line) is plotted in Fig. 2 in comparison with the experimental result (inset). Our result shows that the maximal SC transition temperature T<sub>c</sub> of the $`t`$-$`t^{}`$-$`J`$ model occurs around the optimal doping $`x_{\mathrm{opt}}0.15`$, and then decreases in both underdoped and overdoped regimes. Furthermore, T<sub>c</sub> in the underdoped regime is proportional to the hole doping concentration $`x`$, and therefore T<sub>c</sub> in the underdoped regime is set by the hole doping concentration . This reflects that the density of the dressed holons directly determines the superfluid density in the underdoped regime. Using an reasonably estimative value of $`J800`$K to 1200K in doped cuprates, the SC transition temperature in the optimal doping is T$`{}_{c}{}^{}0.22J176\mathrm{K}264\mathrm{K}`$, in qualitative agreement with the experimental data . In comparison with the result of the $`t`$-$`J`$ model , our present result also shows that $`t^{}`$ plays an important role in enhancing the SC transition temperature of the $`t`$-$`J`$ model and in shifting the maximal value of $`T_c`$ towards to the low doping regime.
For cuprate superconductors, ARPES experiments have produced some interesting data that introduce important constraints on the SC theory . Since cuprates superconductors are highly anisotropic materials, therefore the electron spectral function $`A(𝐤,\omega )`$ is dependent on the in-plane momentum . Although the electron spectral function in doped cuprates obtained from ARPES is very broad in the normal-state, indicating that there are no quasiparticles . However, in the SC-state, the full energy dispersion of quasiparticles has been observed . According to a comparison of the density of states as measured by scanning tunnelling microscopy and ARPES spectral function at $`[\pi ,0]`$ point on identical samples, it has been shown that the most contributions of the electron spectral function come from $`[\pi ,0]`$ point. In addition, the d-wave gap, and therefore the electron pairing energy scale, is maximized at $`[\pi ,0]`$ point. Although the sharp SC quasiparticle peak at $`[\pi ,0]`$ point in cuprate superconductors has been widely studied, the orgin and its implications are still under debate . As a test of the kinetic energy driven superconductivity in doped cuprates , we now study this issue. For discussions of the electron spectral function, we need to calculate the electron diagonal Green’s function $`G(ij,tt^{})=C_{i\sigma }(t);C_{j\sigma }^{}(t^{})`$, which is a convolution of the dressed spin Green’s function and dressed holon diagonal Green’s function, and can be evaluated in terms of the MF dressed spin Green’s function $`D^{(0)}(p)`$ and dressed holon diagonal Green’s function $`g(k)`$ in Eq. (7a) as,
$`G(𝐤,\omega )`$ $`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{𝐩}{}}Z_F{\displaystyle \frac{B_𝐩}{2\omega _𝐩}}\{({\displaystyle \frac{U_{h𝐩+𝐤}^2}{\omega +E_{h𝐩+𝐤}\omega _𝐩}}`$ (38)
$`+`$ $`{\displaystyle \frac{V_{h𝐩+𝐤}^2}{\omega E_{h𝐩+𝐤}+\omega _𝐩}})[n_F(E_{h𝐩+𝐤})+n_B(\omega _𝐩)]`$ (39)
$`+`$ $`[1n_F(E_{h𝐩+𝐤})+n_B(\omega _𝐩)]`$ (40)
$`\times `$ $`({\displaystyle \frac{U_{h𝐩+𝐤}^2}{\omega +E_{h𝐩+𝐤}+\omega _𝐩}}+{\displaystyle \frac{V_{h𝐩+𝐤}^2}{\omega E_{h𝐩+𝐤}\omega _𝐩}})\},`$ (41)
then from this electron diagonal Green’s function, the electron spectral function $`A(𝐤,\omega )=2\mathrm{I}\mathrm{m}G(𝐤,\omega )`$ is obtained as,
$`A(𝐤,\omega )`$ $`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{𝐩}{}}Z_F{\displaystyle \frac{B_𝐩}{2\omega _𝐩}}\{[n_F(E_{h𝐩+𝐤})+n_B(\omega _𝐩)]`$ (42)
$`\times `$ $`[U_{h𝐩+𝐤}^2\delta (\omega +E_{h𝐩+𝐤}\omega _𝐩)`$ (43)
$`+`$ $`V_{h𝐩+𝐤}^2\delta (\omega E_{h𝐩+𝐤}+\omega _𝐩)]`$ (44)
$`+`$ $`[1n_F(E_{h𝐩+𝐤})+n_B(\omega _𝐩)]`$ (45)
$`\times `$ $`[U_{h𝐩+𝐤}^2\delta (\omega +E_{h𝐩+𝐤}+\omega _𝐩)`$ (46)
$`+`$ $`V_{h𝐩+𝐤}^2\delta (\omega E_{h𝐩+𝐤}\omega _𝐩)]\}.`$ (47)
We have performed the calculation for this electron spectral function, and the result of $`A(𝐤,\omega )`$ at $`[\pi ,0]`$ point in the optimal doping $`x_{\mathrm{opt}}=0.15`$ with $`T=0.002J`$ for $`t/J=2.5`$ and $`t^{}/J=0.3`$ is plotted in Fig. 3 in comparison with the experimental result (inset). Our result shows that there is a sharp SC quasiparticle peak near the electron Fermi surface at $`[\pi ,0]`$ point, and the position of this SC quasiparticle peak is located at $`\omega _{\mathrm{peak}}0.4J0.028`$eV$`0.04`$eV, which is quantitatively consistent with the $`\omega _{\mathrm{peak}}0.03`$eV observed in the cuprate superconductor Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8+x</sub>. Our result also shows that the dressed holon pairs condense with the d-wave symmetry in a wide range of the doping concentration, then the electron Cooper pairs originating from the dressed holon pairing state are due to the charge-spin recombination, and their condensation automatically gives the electron quasiparticle character. Furthermore, we have discussed the temperature dependence of the electron spectral function and overall quasiparticle dispersion, and these and related theoretical results will be presented elsewhere.
Our present result also indicates that the SC-state of cuprate superconductors is the conventional BCS like , this can be understood from the electron diagonal and off-diagonal Green’s functions in Eqs. (12) and (10), which can be rewritten as,
$`G(𝐤,\omega )`$ $`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{𝐩}{}}Z_F{\displaystyle \frac{B_𝐩}{4\omega _𝐩}}\{({\displaystyle \frac{U_{h𝐩+𝐤}^2}{\omega +E_{h𝐩+𝐤}\omega _𝐩}}`$ (49)
$`+`$ $`{\displaystyle \frac{U_{h𝐩+𝐤}^2}{\omega +E_{h𝐩+𝐤}+\omega _𝐩}}+{\displaystyle \frac{V_{h𝐩+𝐤}^2}{\omega E_{h𝐩+𝐤}+\omega _𝐩}}`$ (50)
$`+`$ $`{\displaystyle \frac{V_{h𝐩+𝐤}^2}{\omega E_{h𝐩+𝐤}\omega _𝐩}})\mathrm{coth}[{\displaystyle \frac{1}{2}}\beta \omega _𝐩]`$ (51)
$`+`$ $`\mathrm{tanh}[{\displaystyle \frac{1}{2}}\beta E_{h𝐩+𝐤}]({\displaystyle \frac{U_{h𝐩+𝐤}^2}{\omega +E_{h𝐩+𝐤}+\omega _𝐩}}`$ (52)
$``$ $`{\displaystyle \frac{U_{h𝐩+𝐤}^2}{\omega +E_{h𝐩+𝐤}\omega _𝐩}}+{\displaystyle \frac{V_{h𝐩+𝐤}^2}{\omega E_{h𝐩+𝐤}\omega _𝐩}}`$ (53)
$``$ $`{\displaystyle \frac{V_{h𝐩+𝐤}^2}{\omega E_{h𝐩+𝐤}+\omega _𝐩}})\},`$ (54)
$`\mathrm{\Gamma }^{}(𝐤,\omega )`$ $`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{𝐩}{}}Z_F{\displaystyle \frac{\overline{\mathrm{\Delta }}_{hZ}(𝐩+𝐤)}{2E_{h𝐩+𝐤}}}{\displaystyle \frac{B_𝐩}{4\omega _𝐩}}\{\mathrm{coth}[{\displaystyle \frac{1}{2}}\beta \omega _𝐩]`$ (55)
$`\times `$ $`({\displaystyle \frac{1}{\omega E_{h𝐩+𝐤}\omega _𝐩}}+{\displaystyle \frac{1}{\omega E_{h𝐩+𝐤}+\omega _𝐩}}`$ (56)
$``$ $`{\displaystyle \frac{1}{\omega +E_{h𝐩+𝐤}+\omega _𝐩}}{\displaystyle \frac{1}{\omega +E_{h𝐩+𝐤}\omega _𝐩}})`$ (57)
$`+`$ $`\mathrm{tanh}[{\displaystyle \frac{1}{2}}\beta E_{h𝐩+𝐤}]({\displaystyle \frac{1}{\omega E_{h𝐩+𝐤}\omega _𝐩}}`$ (58)
$``$ $`{\displaystyle \frac{1}{\omega E_{h𝐩+𝐤}+\omega _𝐩}}{\displaystyle \frac{1}{\omega +E_{h𝐩+𝐤}+\omega _𝐩}}`$ (59)
$`+`$ $`{\displaystyle \frac{1}{\omega +E_{h𝐩+𝐤}\omega _𝐩}})\},`$ (60)
respectively. Since the dressed spins center around $`[\pm \pi ,\pm \pi ]`$ in the Brillouin zone in the mean-field level , therefore the above electron diagonal and off-diagonal Green’s functions can be approximately reduced in terms of $`\omega _{𝐩=\pm \pi ,\pm \pi }0`$ and the equation $`1/2=S_i^zS_i^z=1/N_𝐩B_𝐩\mathrm{coth}(\beta \omega _𝐩/2)/(2\omega _𝐩)`$ as,
$`g(𝐤,\omega )`$ $``$ $`Z_F{\displaystyle \frac{V_𝐤^2}{\omega E_𝐤}}+Z_F{\displaystyle \frac{U_𝐤^2}{\omega +E_𝐤}},`$ (62)
$`\mathrm{\Gamma }^{}(𝐤,\omega )`$ $``$ $`Z_F{\displaystyle \frac{\overline{\mathrm{\Delta }}_{hZ}(𝐤)}{2E_𝐤}}\left({\displaystyle \frac{1}{\omega E_𝐤}}+{\displaystyle \frac{1}{\omega +E_𝐤}}\right),`$ (63)
with the electron quasiparticle coherence factors $`U_𝐤^2V_{𝐤+𝐤_𝐀}^2`$ and $`V_𝐤^2U_{𝐤+𝐤_𝐀}^2`$, and electron quasiparticle spectrum $`E_𝐤E_{h𝐤+𝐤_𝐀}`$, with $`𝐤_𝐀=[\pi ,\pi ]`$, i.e., the hole-like dressed holon quasiparticle coherence factors $`U_{h𝐤}`$ and $`V_{h𝐤}`$ have been transferred into the electron quasiparticle coherence factors $`U_𝐤`$ and $`V_𝐤`$ by the convolution of the dressed spin Green’s function and dressed holon diagonal Green’s function due to the charge-spin recombination, this is why the basic BCS formalism is still valid in discussions of the doping dependence of the effective SC gap parameter and SC transition temperature, and electron spectral function , although the pairing mechanism is driven by the kinetic energy by exchanging dressed spin excitations, and other exotic properties are beyond the BCS theory.
The essential physics of superconductivity in the present $`t`$-$`t^{}`$-$`J`$ model is the same as that in the $`t`$-$`J`$ model . The antisymmetric part of the self-energy function $`\mathrm{\Sigma }_{1o}^{(h)}(𝐤)`$ (then $`Z_F`$) describes the dressed holon (then electron) quasiparticle coherence, and therefore $`Z_F`$ is closely related to the SC quasiparticle density, while the self-energy function $`\mathrm{\Sigma }_2^{(h)}(𝐤)`$ describes the effective dressed holon (then electron) pairing gap function. In particular, both $`Z_F`$ and $`\mathrm{\Sigma }_2^{(h)}(𝐤)`$ are doping and temperature dependent. Since the SC-order is established through an emerging quasiparticle , therefore the SC-order is controlled by both gap function and quasiparticle coherence, and is reflected explicitly in the self-consistent equations (8a) and (8b). The dressed holons (then electrons) interact by exchanging the dressed spins and that this interaction is attractive. This attractive interaction leads to form the dressed holon pairs (then electron Cooper pairs). The perovskite parent compound of doped cuprate superconductors is a Mott insulator, when holes are doped into this insulator, there is a gain in the kinetic energy per hole proportional to $`t`$ due to hopping, but at the same time, the spin correlation is destroyed, costing an energy of approximately $`J`$ per site, therefore the doped holes into the Mott insulator can be considered as a competition between the kinetic energy ($`xt`$) and magnetic energy ($`J`$), and the magnetic energy decreases with increasing doping. In the underdoped and optimally doped regimes, the magnetic energy is rather too large, and the dressed holon (then electron) attractive interaction by exchanging the dressed spin is also rather strong to form the dressed holon pairs (then electron Cooper pairs) for the most dressed holons (then electrons), therefore the number of the dressed holon pairs (then electron Cooper pairs), SC transition temperature , and quasiparticle coherent weight are proportional to the hole doping concentration. However, in the overdoped regime, the magnetic energy is relatively small, and the dressed holon (then electron) attractive interaction by exchanging the dressed spin is also relatively weak, in this case, not all dressed holons (then electrons) can be bounden as dressed holon pairs (then electron Cooper pairs) by the weak attractive interaction, and therefore the number of the dressed holon pairs (then electron Cooper pairs), SC transition temperature , and quasiparticle coherent weight decrease with increasing doping. To show this point clearly, we plot the quasiparticle coherent weight $`Z_F(T_c)`$ as a function of the hole doping concentration for $`t/J=2.5`$ and $`t^{}/t=0.3`$ (solid line) and $`t/J=2.5`$ and $`t^{}=0`$ (dashed line) in Fig. 4. As seen from Fig. 4, the doping dependent behavior of the quasiparticle coherent weight resembles that of the superfluid density in cuprate superconductors, i.e., $`Z_F(T_c)`$ grows linearly with the hole doping concentration in the underdoped and optimally doped regimes, and then decreases with increasing doping in the overdoped regime, which leads to that the SC transition temperature reaches a maximum in the optimal doping, and then decreases in both underdoped and overdoped regimes. The behavior of the doping dependence of $`Z_F`$ in Fig. 4 is consistent with the experimental result , where the quasiparticle coherent weight increases monotonically with increasing doping in the underdoped and optimally doped regimes , and then decreases with increasing doping in the overdoped regime . On the other hand, the electronic structure becomes asymmetric and hole doping shifts the Fermi surface to the van Hove singularity when the additional second neighbor hopping $`t^{}`$ is introduced in the $`t`$-$`J`$ model , which leads to increase the density of states at the Fermi energy, then the SC correlation is enhanced. Furthermore, the additional second neighbor hopping $`t^{}`$ in the $`t`$-$`J`$ model is equivalent to increase the kinetic energy. These are also why $`t^{}`$ plays an important role in enhancing the SC transition temperature of the $`t`$-$`J`$ model under the kinetic energy driven SC mechanism.
In summary, we have discussed the effect of the additional second neighbor hopping $`t^{}`$ on the SC-state of the $`t`$-$`J`$ model based on the kinetic energy driven SC mechanism. Our result shows that $`t^{}`$ plays an important role in enhancing the SC transition temperature of the $`t`$-$`J`$ model. Within the $`t`$-$`t^{}`$-$`J`$ model, we show that the SC-state of cuprate superconductors is the conventional BCS like, so that the basic BCS formalism is still valid in quantitatively reproducing the doping dependence of the effective SC gap parameter and SC transition temperature, and electron spectral function, although the pairing mechanism is driven by the kinetic energy by exchanging dressed spin excitations, and other exotic magnetic properties are beyond the BCS theory.
Superconductivity in cuprates emerges when charge carriers, holes or electrons, are doped into Mott insulators . Both hole-doped and electron-doped cuprate superconductors have the layered structure of the square lattice of the CuO<sub>2</sub> plane separated by insulating layers . In particular, the symmetry of the SC order parameter is common in both case , manifesting that two systems have similar underlying SC mechanism. On the other hand, the strong electron correlation is common for both hole-doped and electron-doped cuprates, then it is possible that superconductivity in electron-doped cuprates is also driven by the kinetic energy as in hole-doped case. Within the $`t`$-$`t^{}`$-$`J`$ model, we have discussed this issue, and found that in analogy to the phase diagram of the hole-doped case, superconductivity appears over a narrow range of the electron doping concentration in the electron-doped side, and the maximum achievable SC transition temperature in the optimal doping in the electron-doped case is much lower than that of the hole-doped side due to the electron-hole asymmetry.
###### Acknowledgements.
The author would like to thank Dr. Huaiming Guo, Professor Y.J. Wang, and Professor H.H. Wen for the helpful discussions. This work was supported by the National Natural Science Foundation of China under Grant Nos. 10125415 and 90403005. |
warning/0506/quant-ph0506070.html | ar5iv | text | # Distributed measurement-based quantum computation
## 1 Introduction
Measurement-based models provide an intriguing new framework for thinking about quantum computation. While quantum circuits are still widely considered as a convenient formalism for describing algorithms, using measurements to steer quantum computation is considered a serious alternative. Due to their inherently probabilistic nature, measurements were long thought to be a disturbance to quantum computations – unavoidable though they are when wanting to read out the final output of a computation. That they can be an active component of a computation has been known for quite some time through the teleportation protocol. Only much later was it realized that in fault-tolerant constructions, measurements can be quite useful. Soon thereafter, with the advent of models such as the one-way quantum computer \[RBB03\] and the teleportation model \[Nie03, Leu04\], it was established that measurements could not only be a recurring component of a computation, but the actual driving force behind it. Moreover, the measurement paradigm throws a whole new light on the strategies for actual physical implementations of a quantum computer \[WRR<sup>+</sup>05\].
However, measurements are not the only crucial ingredient of these models: they are also inherently *distributed*. Indeed, it is the realization that a variation on the teleportation protocol not only transports but also *transforms* quantum information, which is the basis of the teleportation model. Likewise, the one-way quantum computer is all about transformation via measurement and transportation, this time by way of a generic entangled state, the *graph state*. One-qubit measurements on this state transform the logical qubits, i.e. the quantum inputs, while transporting them via a path of graph state qubits. Again, non-local correlations provided by particular entanglement properties, together with measurements, steer the computation. Of course quantum measurements remain intrinsically probabilistic, but this can be solved by applying corrections dependent on previous measurement outcomes, rendering computations effectively deterministic. Note that measurement outcome dependencies are crucial in order to obtain universality of the model. As a result, the typically distributed notion of classical communication is also naturally present in these models. All of this is very nicely captured by the *measurement calculus* \[DKP04\], a formal framework for one-way computations. Measurement patterns are defined essentially by sequences of commands allowed in the one-way model. From this one can define operational and denotational semantics and show their equivalence, and prove that notions of composition are well-defined. More importantly, there is an associated *rewrite system* which allows one to put any pattern into a standard form. The measurement calculus, which can be seen as an *assembly language*, proves to be a valuable tool for formal investigations into all measurement-based models; for example, one can easily show how the teleportation model reduces to the one-way model, via a conversion between the associated calculi \[DKP05\].
Because of the inherently distributed aspect, measurement-based models for quantum computation are well-suited as a starting point for a formal model for distributed quantum computations. By this we mean *macroscopically* distributed, i.e. we are talking about coordinated actions between different parties. Of these there are many examples in quantum computation \[NC00\]: teleportation, of course, but also entanglement swapping, logic gate teleportation, cryptographic protocols, and also quantum versions of classical distributed applications such as leader election \[DP04\]. However, a formal language for distributed quantum computation is lacking. There have recently been interesting developments based on classical process calculi \[LJ04, GN04\], which have focused mostly on the concurrency aspects. While the distributed nature of computations was introduced via types in Ref. \[GN04\], the aim there is to develop formal verification techniques. In this work, we take an assembly-language point of view, assuming that computations are well-defined. This results in a compact model, with which we can explore properties of distributed protocols, such as the coordination between agents.
We define an assembly language for distributed applications, directly built on the most basic distributed model of all: the one-way quantum computer. We adopt a *local* view and describe the system as a set of *agents* communicating synchronously and operating on a globally entangled quantum state; this is explained in Sec. 2. In Sec. 2.1 we develop a formal semantics for systems of agents as probabilistic transition systems. Operational equivalence is defined in a way such that it corresponds to the notion of bisimilarity. We then prove that quantum teleportation is bisimilar to a direct quantum channel, and this also within the context of other, possibly entangled, agents, in Sec. 3. While the correctness of teleportation has been proved within other formal frameworks before \[GN04, AC04\], the bisimilarity approach, and specifically, taking into account larger contexts as well, is new. We conclude in Sec. 4.
Some familiarity with the measurement calculus model is assumed; for an in-depth exposition we refer to Ref. \[DKP04\].
## 2 Networks of agents
The main concept in our model for distributed measurement-based quantum computations is that of an *agent*. Agents are localized processes which, executing concurrently, make up a distributed system. Formally, we define agents in the following way.
###### Definition 2.1
An *agent* $`𝐀(𝐢,𝐨):Q.`$, with classical input $`𝐢`$ and output $`𝐨`$, and sort given by a set of qubit references $`Q`$, is defined by a finite event sequence $``$ composed of
1. patterns command sequences $`𝒜`$, with input qubits in $`Q`$;
2. classical message reception $`𝚌\mathrm{?}x`$ and sending $`𝚌!y`$, where $`𝚌`$ is a classical channel, and $`x`$ and $`y`$ are names;
3. qubit reception $`\mathrm{𝚚𝚌}\mathrm{?}x`$ and sending $`\mathrm{𝚚𝚌}!q`$, where $`\mathrm{𝚚𝚌}`$ is a quantum channel and $`q`$ a qubit reference.
An agent’s *state* is given by a classical environment $`\mathrm{\Gamma }`$, which is a partial mapping from names, i.e. classical variables and qubit references, to values.
Notice that any pattern $`𝒫(V,I,O,𝒜)`$ trivially corresponds to agent $`𝐀:I.𝒜`$. In general the sort $`Q`$ equals $`II^s`$, where $`I`$ is the local quantum input and $`I^s`$ are qubits of a shared entangled state supplied by the network – see Def. 2.3 below. The classical input $`𝐢`$ and output $`𝐨`$ allows us to model protocols such as superdense coding, in which an agent wishes to send classical values to another agent with the help of a shared Bell state. The local state is used to store measurement outcomes resulting from local pattern executions, input from other agents over classical channels $`𝚌`$, and classical input bindings from the local environment. Further bindings are added to the state as required; that is, for example, whenever a qubit is measured, only then the signal name $`s_i`$ corresponding to the qubit reference $`q_i`$ is added to the domain of $`\mathrm{\Gamma }`$ and bound to the classical measurement outcome $`v`$, denoted $`\mathrm{\Gamma }[s_iv]`$. The classical output set $`𝐨`$ determines which bindings in $`\mathrm{\Gamma }`$ have to be preserved for the final output of a computation. We denote the local state restricted to the classical output by $`\mathrm{\Gamma }_𝐨`$.
Our interpretation of agents is different than the usual process approach, in that agents in our setting always correspond to actual parties in a distributed network, denoted by the label $`𝐀`$. Therefore, an expression of the form $`𝐀(𝐢,𝐨):Q.`$ should be read as: the agent with *name* $`𝐀`$ (Alice) runs the *program* $``$ with qubits and in- and outputs as specified. An agents thus is a piece of code running on a particular processor. In this context, it only makes sense to compose agents if the agent names are the same. We refer to this as agent composition, which is formally defined below. Note that, since we allow subsequent agent programs to supply extra classical and quantum inputs as well as further process outputs from a previous program, some care is required in determining the classical type and quantum sort of the resulting agent. Other than that, agent composition is just straightforward concatenation of event sequences.
###### Definition 2.2
The composition of agents $`𝐀(𝐢_1,𝐨_1):Q_1._1`$ and $`𝐀(𝐢_2,𝐨_2):Q_2._2`$, which is only defined for agents with identical agent names, is denoted
$$𝐀[(𝐢_2,𝐨_2):Q_2._2][(𝐢_1,𝐨_1):Q_1._1]$$
and given by
$$𝐀(𝐢,𝐨):Q._2_1with\{\begin{array}{cc}𝐢=𝐢_1(𝐢_2\backslash 𝐨_1)\hfill & \\ 𝐨=𝐨_1𝐨_2\hfill & \\ Q=Q_1(Q_2\backslash Q_1^{})\hfill & \end{array}\text{ ,}$$
(1)
where $`Q_1^{}`$ is the output sort of the first agent.
Note that output sorts can be determined by inspecting an agent’s event sequence. The general idea is that $`𝐨_1𝐢_2`$ and $`Q_1^{}Q_2`$, so that subsequent programs are defined at least on the outputs of previous ones. This can always be arranged by assuming identity transformations for those names in $`𝐨_1`$ and $`Q_1^{}`$ not appearing in $`_2`$. In particular this means that $`Q^{}=Q_2^{}`$. Extra inputs in subsequent programs are encountered for example when several qubits are teleported one after each other. In this case, Alice needs to execute her side of the protocol several times in a row, supplying a new local quantum input every time she initiates the protocol. The same argument holds for classical input when sequencing several dense coding protocols. Of course this is best understood within the context of *network* rather than agent composition, which is defined below.
A network of agents consist of several agents executing their event sequence concurrently, together with a global shared entangled state. As the network quantum state is inherently nonlocal, there is no other option than to regard it as some kind of global memory – even though we wish to adhere to a local view. This leads to the following definition.
###### Definition 2.3
A *network of agents* $`𝒩`$ is defined by a set of concurrently acting agents together with a shared quantum state, that is
$$\begin{array}{cc}\hfill 𝒩& =𝐀_1(𝐢_1,𝐨_1):Q_1._1|\mathrm{}|𝐀_m(𝐢_m,𝐨_m):Q_m._m\sigma \hfill \\ & =|_i𝐀_i(𝐢_i,𝐨_i):Q_i._i\sigma \text{ ,}\hfill \end{array}$$
(2)
where $`\sigma 𝒟(_{_iI_i^s})`$, with $`Q_i=I_iI_i^s`$ for all $`i`$.
The network state $`\sigma `$ in the definition is the initial entanglement resource which is distributed among agents. Local quantum inputs specified in $`I_i`$ are added to the network state $`\sigma `$ during initialization. In this way we can keep initial shared entanglement as a first-class primitive in our model. In this paper, we do not describe the actual procedure for producing $`\sigma `$. Note that agents in a network need to have different names, since they correspond to different parties that make up the distributed system. In other words, concurrency comes *only* from distribution; we do not consider parallel composition of processes in the context of one party. Finally, individual agents $`𝐀(𝐢,𝐨):Q.`$ trivially correspond to a network $`𝐀(𝐢,𝐨):Q.\mathbf{\hspace{0.17em}0}`$. Therefore statements about networks affect individual agents as well.
We define two different ways of composing networks of agents, namely sequential and parallel composition. Because of our interpretation of agents as distributed processes, there are some constraints on these operations. Sequential composition is only defined for networks containing the same agents; the idea is that agents carry out event sequences of both networks one after the other. Furthermore, agent composition must be defined as per Def. 2.2, that is, inputs and initial sorts of the second network must contain outputs and final sorts of the first. Formally, as follows.
###### Definition 2.4
The sequential composition of networks $`𝒩_1=|_{i=1}^m𝐀_i(𝐢_{1,i},𝐨_{1,i}):Q_{1,i}._{1,i}\sigma _1`$ and $`𝒩_2=|_{i=1}^m𝐀_i(𝐢_{2,i},𝐨_{2,i}):Q_{2,i}._{2,i}\sigma _2`$ is defined as
$$𝒩_2𝒩_1=|_{i=1}^m𝐀_i[(𝐢_{2,i},𝐨_{2,i}):Q_{2,i}._{2,i}][(𝐢_{1,i},𝐨_{1,i}):Q_{1,i}._{1,i}]\sigma _1\sigma _2\text{ .}$$
(3)
Note that we have overloaded the notation $``$ to denote both agent and network composition. As long as both networks have the same number of agents, one can always arrange for them two be sequentially composable by renaming agents. In fact one can even compose networks with a different number of agents by adding null agents, i.e. just agent names, to the network with fewer agents. The preparations of both networks are assumed to be defined on disjoint Hilbert spaces; the composed network’s preparation is given by the tensor product of these. They are to be interpreted as different entanglement resources that are used by both networks. Parallel composition, which expresses that two networks are operating in parallel and independent of each other, is only defined for networks containing different agents. Again, one can always rename agents so that this is well-defined.
###### Definition 2.5
The composition of networks $`𝒩_1=|_{i=1}^m𝐀_i(𝐢_{1,i},𝐨_{1,i}):Q_{1,i}._{1,i}`$ $`\sigma _1`$ and $`𝒩_2=|_{i=1}^n𝐁_i(𝐢_{2,i},𝐨_{2,i}):Q_{2,i}._{2,i}\sigma _2`$ is defined as
$$\begin{array}{cc}\hfill 𝒩_1𝒩_2=& |_{i=1}^m𝐀_i(𝐢_{1,i},𝐨_{1,i}):Q_{1,i}._{1,i}\hfill \\ & |_{i=1}^n𝐁_i(𝐢_{2,i},𝐨_{2,i}):Q_{2,i}._{2,i}\hfill \\ & \sigma _1\sigma _2\text{ .}\hfill \end{array}$$
(4)
By combining sequential and parallel composition, one can express a broad range of network combinations. In Sec. 2.3 we show that the semantics of networks is preserved under these operations.
The definitions given in the above can be rendered concisely under the form of an abstract grammar. We use $`[]`$ instead of $`|`$ to separate choices for expressions, as the latter is already in use to denote parallel composition.
$$\begin{array}{cc}\hfill 𝒜& ::=\mathrm{nil}[]E[]M[]C[]𝒜𝒜[]𝒜.𝒜\hfill \\ \hfill & ::=𝚌\mathrm{?}x[]𝚌!x[]\mathrm{𝚚𝚌}\mathrm{?}q[]\mathrm{𝚚𝚌}!q[]𝒜[].\hfill \\ \hfill 𝐚& ::=𝐀(𝐢,𝐨):Q.[]𝐀[(𝐢_2,𝐨_2):Q_2._2][(𝐢_1,𝐨_1):Q_1._1]\hfill \\ \hfill 𝒩& ::=|_i𝐚_i\sigma []𝒩𝒩[]𝒩𝒩\hfill \end{array}$$
(5)
In the above $`E`$, $`M`$ and $`C`$ stand for any entanglement, measurement or correction commands, and nil is the null command. We refrain from giving a grammar for names; rather, we use conventions for these described below. Notice that both pattern command sequences and networks can be viewed as processes in the traditional process algebra sense, with a composition operation transforming any two patterns, respectively networks, into a new pattern or network. The definition of event sequences, which can be composed sequentially, is also common. Agents, however, have a less clear status in the process algebra framework. This is exactly because they are constructs that formalize *distributed* notions, and therefore there are constraints on how they may be composed. This is not a flaw of our model but rather a requirement if one wants to make distribution explicit.
There are several notational conventions we adhere to throughout this paper. When the context is clear, as is often the case, we do not explicitly mention inputs and outputs, but instead just write $`𝐀:Q`$. Moreover, we do not write inputs, outputs, sorts, or preparations if there are none, sometimes writing $``$ for empty input or output sets. For example, $`𝐀(,\{x\}).𝚌\mathrm{?}x`$ is an agent with no input and no qubits, while $`𝐀(,\{x\}).𝚌\mathrm{?}x|𝐁(\{y\},).𝚌!y`$ is a network with no preparation in which a classical value is exchanged between two agents. We write $`𝚌\mathrm{?}xy`$ for $`𝚌\mathrm{?}x𝚌\mathrm{?}y`$ and similarly for other communicating channels. In Def. 2.1 we have used pattern command sequences rather than full patterns for brevity, with the convention that the input of a pattern is always given by those qubits in its computation space that belong to the agent’s sort at the moment when the pattern is executed. The output of the pattern is simply given by those qubits that are not measured. For example, in the single-agent network $`𝐀:\{1,2\}.X_4^{s_1}M_1^0E_{14}`$ qubit 1 is an input to the pattern and 4 is an output; we also write this as $`𝐀:\{1,2\}.(1,4)`$. Pattern qubits not in an agent’s sort are assumed to be computation qubits initialized to $`|+`$ as before, and need not be mentioned explicitly in an agent’s sort. On the other hand, agent qubits not mentioned in a pattern event, as qubit 2 above, are always assumed to be left alone. That is, we do not explicitly write the identity pattern $``$ applied to $`𝐀`$’s remaining qubit $`2`$. We consider pattern command *sequences* rather than singular commands so that we can use the big-step semantics of patterns. As can be seen from these examples, we often refer to qubits via numbers; specifically we mostly refer to qubit $`q_i`$ as qubit $`i`$, and $`s_i`$ instead of $`s_{q_i}`$ for the corresponding signal variable. This has the advantage that patterns look just as they do in Ref. \[DKP04\]. In what follows below, we use specific letters for specific names, in particular we use $`q_i`$ or just $`i`$ for qubit references, $`x_i`$ and $`y_i`$ for ordinary classical variables, $`v_i`$ for classical values, and $`s_i`$ for classical signal variables.
We subject networks of agents to definiteness conditions, which ensure that the computation is well-defined.
an agent’s communication events operate only on qubits in its sort $`Q`$. Pattern events with computation space $`V`$ have inputs in $`QV`$;
an agent’s events depend only on values in its state $`\mathrm{\Gamma }`$;
each quantum and classical message reception event has a corresponding quantum, respectively classical message sending event;
all names, i.e. classical variables and quantum references, are unique.
### 2.1 Operational semantics
Before we provide concrete evaluation rules for distributed computations, we give some clarifying explanations and examples as to how execution proceeds in the local view. Throughout network evolution, each agent has access to the network state $`\sigma `$ via the qubits it owns, transforming $`\sigma `$ whenever a pattern event is executed. While these patterns are local, $`\sigma `$ is *not*, and to preserve all information on the correlations we need to keep $`\sigma `$ unreduced at all times. As an example, suppose that an agent $`𝐀`$ owns the first two qubits of the system’s state $`\sigma (1,2,3,4,5)`$. Its next event is to execute the Hadamard pattern $`(1,6)`$ on its first qubit. As stated above, we do not explicitly write the identity pattern $``$ applied to $`𝐀`$’s remaining qubit $`2`$. More importantly, neither do we write explicit identity patterns for the qubits not belonging to $`𝐀`$, in this case $`3`$, $`4`$ and $`5`$. In full, this means that we have an evaluation step as follows
$$\sigma ,𝐀:\{1,2\}.(1,6)\sigma ^{},𝐀:\{6,2\}\text{ ,}$$
(6)
where $`\sigma ^{}=(1,6)^4(2,3,4,5)\sigma `$, that is, $`\sigma ^{}=(HI^4)\sigma (HI^4)`$. We denote the transition relation by $``$. Execution of $``$ occurs in a single transition step by relying on its big-step semantics. In this way we avoid getting into the actual details of pattern execution, which is not what this paper is about. Note especially that the sort of $`𝐀`$ has changed. This is because measurements are destructive, so that the input qubit $`1`$ has disappeared, and the output $`6`$ has taken its place. Because of this, we sometimes explicitly write $`\sigma ^{}(6,2,3,4,5)`$ in the right-hand side of the above evaluation rule. The only remnant of $`1`$ is its corresponding measurement outcome $`s_1`$, which is recorded in the state $`\mathrm{\Gamma }`$, via the added binding $`\mathrm{\Gamma }[s_1v]`$, where $`v`$ is the measurement outcome. $`𝐀`$ might subsequently send $`s_1`$ to other agents by an event $`𝚌!s_1`$, so in general we cannot delete this entry from $`\mathrm{\Gamma }`$. Essentially this means that when new qubit references are generated by the execution of subsequent local patterns, or by quantum communications from other agents, these names need to be *unique*. Because we do not want to get into the actual details of this naming procedure, as we are considering our model to be at the level of an assembly language, i.e. at a stage where naming conflicts have been resolved, we have imposed this as a definiteness condition in the above.
A final point concerns an agent’s classical input and output, received from, respectively sent to, its own local system. In the local view, pattern events can depend also on classical inputs, rather than only on measurement outcomes. Because of the uniform structure of the local state $`\mathrm{\Gamma }`$ we can piggyback input dependencies onto the signal dependency structure of MC. This is, as mentioned before, one of the reasons why MC is such a good basis on which to built a framework for distributed quantum computations. With these concrete examples in mind, we are now ready to develop the operational semantics of the local view. As usual, we first define rules in terms of small-step transitions, after which we switch over to the big-step framework.
The small-step transitions for distributed computations essentially describe how agents, and the network with them, evolve over different time steps. We adopt a shorthand notation for agents, leaving out classical inputs and output, which do not change with small-step reductions.
$$\begin{array}{cc}\hfill 𝐚_i& =𝐀_i:Q_i._i\hfill \\ \hfill 𝐚_i.E& =𝐀_i:Q_i.[_i.E]\hfill \\ \hfill 𝐚^q& =𝐀:Q\backslash \{q\}.\hfill \\ \hfill 𝐚^{+q}& =𝐀:Q\{q\}.[q/x]\text{ ,}\hfill \end{array}$$
(7)
where $`E`$ is some event, and $`_i`$ and $`_i^{}`$ are event sequences. A *configuration* is given by the system state $`\sigma `$ together with a set of agent programs, and their states, specifically
$$\sigma ,|_i\mathrm{\Gamma }_i,𝐚_i=\sigma ,\mathrm{\Gamma }_1,𝐚_1|\mathrm{\Gamma }_2,𝐚_2|\mathrm{}|\mathrm{\Gamma }_m,𝐚_m\text{ .}$$
(8)
The small-step rules for configuration transitions, denoted $``$, are specified below; we give some explanations afterwards. When the system state is not changed in an evaluation step, we stress this by preceding a rule by $`\sigma `$.
$$\frac{\sigma ,𝒫(V,I,O,𝒜)_\lambda \sigma ^{},\mathrm{\Gamma }^{}}{\sigma ,\mathrm{\Gamma },𝐀:IR.[.𝒫]_\lambda \sigma ^{},\mathrm{\Gamma }\mathrm{\Gamma }^{},𝐀:OR.}$$
(9)
$$\frac{\mathrm{\Gamma }_2(y)=v}{\sigma (\mathrm{\Gamma }_1,𝐚_1.𝚌\mathrm{?}x|\mathrm{\Gamma }_2,𝐚_2.𝚌!y\mathrm{\Gamma }_1[xv],𝐚_1|\mathrm{\Gamma }_2,𝐚_2)}$$
(10)
$$\frac{}{\sigma (\mathrm{\Gamma }_1,𝐚_1.\mathrm{𝚚𝚌}\mathrm{?}x|\mathrm{\Gamma }_2,𝐚_2.\mathrm{𝚚𝚌}!q\mathrm{\Gamma }_1,𝐚_1^{+q}|\mathrm{\Gamma }_2,𝐚_2^q)}$$
(11)
$$\frac{L_\lambda R}{L|L^{}_\lambda R|L^{}}$$
(12)
Here, $``$ denotes the union of outcome maps. Implicit in these rules is a sequential composition rule, which ensures that all events in an agent’s event sequence are executed one after the other. The first rule is for local operations; we have written the full pattern instead of only its command sequence here to make pattern input and output explicit. Because a pattern’s big-step semantics is given by a probabilistic transition system described by $``$, we pick up a probability $`\lambda `$ here. Furthermore, an agent changes its sort depending on pattern’s output $`O`$, as explained in the examples above. The next rule is for classical rendez-vous and is straightforward. For quantum rendez-vous, we need to substitute $`q`$ for $`x`$ in the event sequence of the receiving agent, and furthermore adapt qubit sorts. The last rule is a metarule, which is required to express that any of the other rules may fire in the context of a larger system. $`L`$ and $`R`$ stand for any of the possible left-, respectively right-hand sides of any of the previous rules, while $`L^{}`$ is an arbitrary configuration. Note that we might need to rearrange terms in the parallel composition of agents in order to be able to apply the context rule. This can always be done since the order of agents in a configuration is arbitrary. In derivations of network execution, we often do not explicitly write reductions as specified by (12), but rather specify which in which order the other rules fire in the context of the network at hand. It is precisely in this last rule that introduces nondeterminism at the network level, that is, several agent transitions may be possible within the context of a network at the same time.
Starting from the small-step rules above we can now define the big-step semantics of a system of agents. We first define computation paths, which run from initial to final configurations via small-step transitions. In the *initial configuration*, all local states are given by the map containing the classical input bindings $`\mathrm{\Gamma }_{𝐢_i}`$, while the network state is determined by the entanglement resource together with local quantum inputs. A *final configuration* is one in which all agents have an empty command sequence, and in which all local states have been restricted to classical output bindings, $`\mathrm{\Gamma }_{i,𝐨_i}`$. Because of the definiteness conditions we imposed each computation always ends with a final configuration. Supposing the initial quantum state of agent $`𝐀_i`$ is given by $`\rho _iI_i`$, paths are defined as follows.
###### Definition 2.6
Given a network of agents $`𝒩=|_i𝐀_i(𝐢_i,𝐨_i):Q_i._i\sigma `$ and quantum inputs $`\rho _i`$, a *path* $`\gamma `$ is a *maximal* sequence of configurations $`\{C_j=\sigma _j,|_i\mathrm{\Gamma }_i^j,𝐚_i^j,j=1,\mathrm{},k1\}`$, i.e.
$$\begin{array}{cc}\hfill C_1& =\sigma _{i=1}^m\rho _i,|_i\mathrm{\Gamma }_{𝐢_i},𝐀_i:Q_i._i\hfill \\ \hfill C_j& _{\lambda _j}C_{j+1}\hfill \\ \hfill C_k& =\sigma _k,|_i\mathrm{\Gamma }_{i,𝐨_i}^k,𝐀_i:Q_i^k\hfill \end{array}$$
(13)
We write $`C_1\underset{\lambda _\gamma }{\overset{\gamma }{}}C_k`$ where $`\lambda _\gamma =_{j=1}^k\lambda _j`$, and call $`C_k`$ a final configuration of $`𝒩`$.
Notice also that paths always terminate since event sequences are finite.
One could straightforwardly define the operational semantics of a system $`𝒩`$ to be the probabilistic transition system (PTS) defined by of all its paths. However, we choose to identify those paths leading to the same observable behavior of an agent network. Concretely, for particular inputs $`𝐢`$ and $`\rho I`$, we identify final configurations for which only internal bindings in the local state of agents are different, that is, bindings for names not in the classical output set $`𝐨`$. These bindings correspond either to outcomes of measurements appearing in pattern events, or result from classical rendez-vous events. As long as these are not part of the classical output, their actual values are unimportant. We cannot trace out these measurement outcomes after each pattern event, since a subsequent event may depend on these values. Notice that the final sorts of agents *do* need to be identical, as well as, obviously, the final network quantum state. Identifying such paths, in the style of the measurement calculus, then gives us the semantics of a network of agents. However, because of the nondeterminism in the order in which concurrent agents execute their event sequence, we have to define the operational semantics of a network with respect to a particular *schedule*. A schedule is precisely a particular order in which agents execute events. For example, in the network $`𝐀:\{1\}.(1,2)|𝐁:\{3\}.(3,4)`$, possible schedules are $`\mathrm{𝐀𝐁}`$ and $`\mathrm{𝐁𝐀}`$. If we do not take schedules into account, we would add probabilities of all paths and all schedules resulting in identical final configurations, which for the example above leads to a probability of 2 for computing $`HH`$ for arbitrary inputs, which is clearly not what we intend to say. We refrain from giving a formal definition for schedules, as we will find in Sec. 2.2 that the semantics of a network is independent of the schedule. Putting all this together, we obtain the following definition.
###### Definition 2.7
The *operational semantics* of a network $`𝒩=|_i𝐀_i(𝐢_i,𝐨_i):Q_i._i\sigma `$, with respect to a particular schedule, is a probabilistic transition system relating initial with final configurations,
$$𝒩_{op}:_iQ_i_iQ_i^{}._i\rho _i,|_i\mathrm{\Gamma }_{𝐢_i}_\lambda \sigma ^{},|_i\mathrm{\Gamma }_{𝐨_i}$$
(14)
with $`\lambda =_\gamma \lambda _\gamma `$ and the sum runs over all paths $`\gamma `$ such that
$$\sigma _i\rho _i,|_i(\mathrm{\Gamma }_{𝐢_i},𝐀_i(𝐢_i,𝐨_i):Q_i._i)\underset{\lambda _\gamma }{\overset{\gamma }{}}\sigma ^{},|_i(\mathrm{\Gamma }_{𝐨_i},𝐀_i(𝐢_i,𝐨_i):Q_i^{})\text{ .}$$
(15)
We call $`_iQ_i_iQ_i^{}`$ the *type* of the network.
From now on, we denote $`I=_iI_i`$ for the set of input qubits and $`O=_iQ_i^{}`$ for the set of output qubits, and call $`𝒟(_I)`$ and $`𝒟(_O)`$ the quantum input and output space respectively. The semantics of a network with respect to a schedule is thus that it relates quantum states in $`𝒟(_I)`$ plus classical input to quantum states in $`𝒟(_O)`$ and classical output with particular probabilities. Note that the type of the transition system is a mapping from initial to final sorts; this component is identical in the denotational semantics we develop in Sec. 2.2.
We say that two networks $`𝒩_1`$ and $`𝒩_2`$ are *operationally equivalent* if their operational semantics, given by a PTS, is identical, and write $`𝒩_1_{op}𝒩_2_{op}`$. In fact, we identify operational equivalence with the notion of *bisimilarity*. This is indeed sensible since, by identifying computation paths as in Def. 2.7, we actually impose a bisimilarity relation on final configurations. As we shall see in Sec. 3, it is by doing exactly this that we can show, among others, that teleportation is bisimilar to a direct quantum channel.
### 2.2 Denotational semantics
Because of its more abstract character, in many situations it is more adequate to work with the denotational semantics. This is why we develop this notion for networks of agents. As before, it is closely related to the operational semantics, as well as to the semantics of ordinary patterns. Indeed, upon inspection of Def. 2.7, we see that for any schedule, the PTS associated with a network of agents decomposes in several parts. First, there is a map from initial to final sorts, which determines how qubit ownership evolves for each of the agents from initial to final configurations. This is formalized as a type signature, exactly as we did for the operational semantics. The sort mapping is independent of the classical input: indeed, classical inputs appear only in classical communications, Pauli corrections and measurement angles, none of which affect qubit sorts. Furthermore, they can be read of statically from the network definition and are also schedule independent. The denotational semantics is then a mapping from classical inputs to classical outputs and a quantum operation, which in turn determines how quantum states evolve in the network. However, these two components are not independent of each other, since classical outputs can be measurement outcomes, which occur with probabilities that depend on the quantum operation applied. For simplicity, let us first consider the case where there are no classical outputs. In this case, for each classical input $`𝐢=_i𝐢_i`$ we have map $``$ which describes how the network quantum state evolves. This map is a multilocal quantum operation, because if we throw away all distributed information, that is, sorts and communication events, we just have an ordinary pattern, i.e. a quantum operation. There is one caveat: since computation occurs asynchronously, there is usually some choice in the order in which different agents execute events in their program, i.e. there are different possible schedules. However, since at each instance of the computation local events operate on disjoint sets of qubits, it does not actually matter in which order these operations are applied, or, in fact, whether they are executed at the same time. This statement is proved formally in Sec. 2.4; we postpone the full proof until then since it has bearing on other situations that are covered below. Therefore, any schedule of the computation leads to the same quantum operation. So to determine the operation elements of $``$, we choose a particular schedule, and then compose patterns in the order in which they are executed, tensoring with identity patterns where necessary and ignoring communication commands. Each operation element $`L_j`$ then corresponds to a sequential composition of actualizations for each of these patterns.
Suppose now that the network contains classical outputs $`𝐨=_i𝐨_i`$. We need to make a distinction between *signal* outputs, which are measurement outcomes, and *external* outputs, which are values that were originally input by some agent and sent around the network. By definition, the external outputs $`𝐨_e=_i𝐨_{i,e}`$ depend only on the classical input $`𝐢`$; these constant values are sent around the network via classical channels. It is precisely the signal outputs $`𝐨_s=_i𝐨_{i,s}`$ that depend on the quantum operation and vice versa. Indeed, when there are signal outputs particular measurement outcomes are preserved, therefore excluding actualizations of $``$ that do not correspond to that outcome. This essentially means that a different quantum operation is applied for each possible signal output. For example, suppose one of the signal outputs, corresponding to a measurement on qubit 3 is equal to $`1`$. Then only those actualizations containing the operator $`_\alpha |_3`$ are compatible with this output. We denote actualizations compatible with output $`𝐨_s`$ by $`L_i^{𝐨_s}`$, and the quantum operation with these operation elements by $`^{𝐨_s}`$, and call these *restricted*. Note, however, that this is a *trace-decreasing* operation, and that $`\mathrm{Tr}()^{𝐨_s(\rho )}`$ is precisely the probability with which the output $`𝐨_s`$ occurs. So, whereas the operational semantics gives explicit probabilities for each path, in the denotational semantics these are contained within the quantum operations. It is this abstraction, together with schedule-independence which makes the denotational framework advantageous. Indeed, classical inputs are the same for all schedules, and classical outputs, depending on classical inputs and measurement values, thus occur with the same probabilities for all schedules, since $``$ is schedule-independent. Putting all of this together, we arrive at the following definition.
###### Definition 2.8
The *denotational semantics* of a network of agents $`𝒩=|_i𝐀_i(𝐢_i,𝐨_i):Q_i._i\sigma `$ is given by
$$𝒩_{de}:_iQ_i_iQ_i^{}.𝐢\{(𝐨,^{𝐨_s}),𝐨_s\}$$
(16)
with
$$:𝒟(_I)𝒟(_O):_i\rho _i\underset{j}{}L_j(\sigma _i\rho _i)L_j^{}\text{ ,}$$
(17)
where $`𝐨=𝐨_e𝐨_s`$, $`\rho _i`$ is the quantum input, $`Q_i^{}`$ the final sort of agent $`𝐀_i`$, and $`I`$ and $`O`$ are quantum input and output spaces respectively. In case there are no outputs, we have $`𝒩_{de}:_iQ_i_iQ_i^{}.𝐢`$, or just $`.`$ if there are no inputs either.
Note that the $`𝐨_e`$ part of $`𝐨`$ in each of the above tuples is identical. As an example, consider the pattern $`X_1^{s_2}M_2^\alpha `$, which implements the bit-flip channel \[NC00\]. It can be interpreted as a one-agent network $`𝐀(,):\{1\}.X_1^{s_2}M_2^\alpha `$, with has as denotational semantics the quantum operation $`(\rho )=p\rho +(1p)X\rho X`$. However, the one-agent network $`𝐀(,\{s_2\}):\{1\}.X_1^{s_2}M_2^\alpha `$ has a different semantics, namely
$$\{(0,p\rho ),(1,(1p)X\rho X)\}\text{ ,}$$
(18)
for all $`\rho `$, where $`p`$ is a function of $`\alpha `$. While this example may seem contrived, it is actually crucial that the semantics of these kind of networks are different, as they describe different states of *knowledge* of the output $`s_2`$, and hence, of what actual computation path was taken.
In case the underlying quantum operation $``$ is deterministic, all actualizations lead to the same quantum output. However, even in that case we require the above formulation with different trace-decreasing quantum operations $`^{𝐨_s}`$, since these determine the probabilities of outputting $`𝐨_s`$. A network is deterministic only if for any input $`𝐢`$, the quantum operation $``$ is deterministic, i.e. it implements a unitary, *and* the classical output $`𝐨`$ is identical in all actualizations.
As before, two networks are called denotationally equivalent if they have the same denotational semantics. Note that equivalent networks may have different preparation states, and also that agent names may be different, though the number of agents must be identical in both networks. With the above definitions in place, we can now prove the following result.
###### Proposition 2.9
There is a *precise correspondence* between the operational and the denotational semantics of networks of agents, that is to say
$$𝒩_1,𝒩_2:𝒩_1_{op}𝒩_2𝒩_1_{de}𝒩_2$$
(19)
Proof. Suppose $`𝒩_1_{de}=𝒩_2_{de}`$. This means that for all classical and quantum inputs, both the the type and classical external outputs are identical. Since for all signal outputs $`_1^{𝐨_s}=_2^{𝐨_s}`$, we have that $`_1=_2`$ and therefore computation paths are also the same for both networks \[D’H05\]. Hence, both networks are operationally equivalent.
Since both semantics’ are equivalent, we can choose the operational or denotational framework at out convenience. We usually derive the denotational semantics via operational computation paths, relying also on the fact that the semantics is schedule independent.
### 2.3 Compositionality
An important goal is to prove that the semantics of networks is conserved with respect to network composition as in Defs. 3 and 4. Knowing this, we can compose any two networks and be ensured that the resulting network carries out the intended computation. In other words, we need to prove Prop. 2.10 below – notice that a similar result exists for patterns \[DKP04\], though the operations are of course defined differently. However, in order to do this we require a proper definition for composing mathematical objects of the form of Def. 2.8. The reasonable way to this is to gather types, inputs and outputs, eliminating those that are fed from one network into the other in case of sequential composition, exactly like we did for agents in Def. 2.2. Next, we need to combine quantum operations by tensoring with the identity map were necessary. As we shall see below, in this way we already recover most of Defs. 3 and 4 . In fact, we only need to be careful in checking whether the quantum operations combine in the correct way, but, for this we can rely on the abovementioned Prop. 2.10.
###### Proposition 2.10
The semantics of networks is *compositional*, i.e.
$`𝒩_2.𝒩_1`$ $`=𝒩_2.𝒩_1`$ (20)
$`𝒩_1𝒩_2`$ $`=𝒩_2𝒩_1`$ (21)
Proof. Suppose we have two networks $`𝒩_1=|_i𝐀_i(𝐢_{1,i},𝐨_{1,i}):Q_{1,i}._{1,i}\sigma _1`$ and $`𝒩_2=|_i𝐀_i(𝐢_{2,i},𝐨_{2,i}):Q_{2,i}._{2,i}\sigma _2`$ with semantics given by
$$\begin{array}{cc}\hfill 𝒩_1& :_iQ_{1,i}_iQ_{1,i}^{}.𝐢_1\{(𝐨_1,_1^{𝐨_s})\}\hfill \\ \hfill 𝒩_2& :_iQ_{2,i}_iQ_{2,i}^{}.𝐢_2\{(𝐨_2,_2^{𝐨_s})\}\text{ ,}\hfill \end{array}$$
(22)
with
$$\begin{array}{cc}\hfill _1& :𝒟(_{I_1})𝒟(_{O_1})_i\rho _{1,i}\underset{j}{}L_{1,j}(\sigma _1_i\rho _{1,i})L_{1,j}^{}\hfill \\ \hfill _2& :𝒟(_{I_2})𝒟(_{O_2})_i\rho _{2,i}\underset{j}{}L_{2,j}(\sigma _2_i\rho _{2,i})L_{2,j}^{}\text{ .}\hfill \end{array}$$
(23)
We then find that
$$𝒩_2.𝒩_1:_iQ_i_iQ_i^{}.𝐢\{(𝐨,_2^{𝐨_{2,s}}.(_1^{𝐨_{1,s}})),𝐨_{1,s},𝐨_{2,s}\}$$
(24)
where type, classical input and output are found by pointwise application of the rules in Eq.(1), and $``$ is the identity operation on $`I_2\backslash O_2`$. Quantum operations in the above map states of the form $`_i\rho _{1,i}_i\rho _{2,i}`$ in $`I=I_1(I_2\backslash O_2)`$ to states in $`O=O_2`$, after tensoring them with $`\sigma _1\sigma _2`$; On the other hand, the semantics of $`𝒩_2𝒩_1`$ is given by
$$𝒩_2𝒩_1:_iQ_i_iQ_i^{}.𝐢\{(𝐨,[_2.(_1)]^{𝐨_s}),𝐨_s\}\text{ ,}$$
(25)
where $`𝐨_s=𝐨_{1,s}𝐨_{2,s}`$, and the quantum operations operate on the same $`I`$ and $`O`$ as above. So we only need to check that first composing the quantum operations and then restricting them is the same as first restricting and then composing. However, this follows from the analogous result for ordinary patterns \[DKP04\].
For parallel composition, consider instead $`𝒩_2=|_{i=1}^n𝐁_i(𝐢_{2,i},𝐨_{2,i}):Q_{2,i}._{2,i}\sigma _2`$, with semantics as above. We then find that
$$𝒩_1𝒩_2:_iQ_i_iQ_i^{}.𝐢\{(𝐨,_1^{𝐨_{1,s}}_2^{𝐨_{2,s}}),𝐨_{1,s},𝐨_{2,s}\}\text{ ,}$$
(26)
where types, inputs and outputs are found by taking the (disjoint) union of those of the composing networks. Again, via the analogous result for ordinary patterns \[DKP04\], we find that $`(_1_2)^{𝐨_s}=_1^{𝐨_{1,s}}_2^{𝐨_{2,s}}`$, and therefore it follows that the above expression equals $`𝒩_1𝒩_2`$.
### 2.4 Entanglement contexts
In the previous section, the networks $`𝒩_1`$ and $`𝒩_2`$ individually operate on tensor product states of the form $`_i\rho _i`$, in accordance with the fact that these are local inputs provided by each of the agents, as is made explicit in Def. 2.7. While this is sensible when considering a network operating in isolation – in which case entangled input states are specified as preparations – it is less so when a network is only one factor in a complex compositional structure. Indeed, already for sequential composition the input state is in general no longer disentangled over agents, since it is fed in partly as output of a previous network computation, and there is no guarantee whatsoever that this output is a product state. Another subtle difference lies in the fact that input spaces of individual agents are combined when composing networks. Concretely, while each agent supplies local inputs $`\rho _1`$ when $`𝒩_1`$ is run separately, and $`\rho _2`$ when $`𝒩_2`$ is run separately, inputs to the composed network $`𝒩_2𝒩_1`$ are generally *not* of the form $`\rho _1\rho _2`$. Yet another situation is that where agents have mixed state inputs because their inputs are entangled into a state on a larger system than that which the network operates on – in fact one view is that this is the *only* way in which mixed states arise. These are of course all typical manifestations of entangled states, but we want to stress the different situations in the context of distributed networks in which these arise. The point is, is entanglement preserved when applying operations to only part of the entangled state? We show below that it does, by proving the a quantum operation applied to the $`A`$-system of a state $`\rho `$ existing on system $`AC`$ does not touch the part of this state in $`C`$. Each of the aforementioned situations can be cast in this form, since $`A`$ may be a system of several agents entangled with a group of other agents described by system $`C`$, as well as an input system of one agent only entangled with another input system $`C`$. The quantum operation itself is defined to map states on $`A`$ to states on $`B`$, because we want to consider situations where a network’s output spaces is different than its input space.
###### Proposition 2.11
Suppose $``$ is a quantum operation on system $`A`$ to system $`B`$ such that
$$:𝒟(_A)𝒟(_B):\rho _A\rho _B=\underset{k}{}L_k\rho _AL_k^{}\text{ .}$$
(27)
Then for all quantum states $`\rho _{AC}`$ living on a system $`AC`$, applying $``$ to one half of $`\rho _{AC}`$ results in
$$\rho _{BC}=\underset{k}{}(L_kI_C)\rho _{AC}(L_k^{}I_C)$$
(28)
Proof. Note first that any complex matrix can be written as a linear combination of Hermitian matrices, which in turn can be written as a sum of density matrices. Then by the spectral decomposition and linearity, it follows that for all complex matrices $`Z`$ we have $`(Z)=_kL_kZL_k^{}`$. Writing $`\rho _{AC}=_{ijkl}\alpha _{ijkl}|i_A|j_Ck|_Al|_C`$, we find that
$$\begin{array}{cc}\hfill (𝒩I_C)\rho _{AC}& =\underset{ijkl}{}\alpha _{ijkl}(𝒩I_C)(|ik|^A|jl|^C)\hfill \\ & =\underset{ijkl}{}\alpha _{ijkl}(\underset{k}{}(L_k|ik|^AL_k^{})|jl|^C)\hfill \\ & =\underset{k}{}(L_kI_C)\rho _{AC}(L_k^{}I_C)\text{ ,}\hfill \end{array}$$
(29)
which proves the theorem.
The proof, while easy, is not trivial, and has several important consequences. First of all, it shows that our statement in Sec. 2.2 on the fact that networks are schedule-independent is true. Indeed, agents transform parts of a shared entangled state via local operations, and the above results shows that any order results in the same multilocal quantum operation. Next, while we have proved compositionality in the previous section considering only product state, we are now ensured that it holds also for arbitrary input states. We rely on Prop. 28 also in the next section, when we discuss sequencing several teleportation networks to transfer an entangled state from one agent to another.
## 3 Teleportation vs. quantum channels
In this section we prove that the teleportation protocol, described within our framework of networks of agents, is bisimilar to a direct quantum communication of the qubit to be teleported. We first give the network specification for direct quantum communication. Next, we give the network for the teleportation protocol, rederiving the correctness of the protocol by developing its semantics in the local view elaborated in Sec. 2, which gives us a handle on qubit locations. A second goal of this section is to prove that it is bisimilar to a direct quantum channel. That is, by comparing the semantics of a direct quantum channel and that of teleportation – evolved in the correctness proof – we conclude that they do indeed define identical PTS’s. By Prop. 28, this holds in arbitrary entanglement contexts. This is a nontrivial result because these agents may be entangled with the qubit to be teleported, and it is not a priori clear whether this entanglement is preserved throughout TP. This result has several consequences that we mention below. While the correctness of teleportation has been proved within other formal frameworks for distributed systems before \[GN04, AC04\], the bisimilarity approach, and specifically, taking into account arbitrary contexts, is new.
Consider a network of two agents named $`𝐀`$ and $`𝐁`$. A direct quantum communication of a qubit is implemented simply by the network
$$𝒩=𝐀:\{1\}.(\mathrm{𝚚𝚌}!1)|𝐁.(\mathrm{𝚚𝚌}\mathrm{?}1)\mathbf{\hspace{0.17em}0}\text{ ,}$$
(30)
where $`\mathrm{𝟎}`$ is the null state. Note that $`𝐁`$ has an empty sort. The small-step semantics, given input $`|\psi `$, is derived in one step, and leads to the operational semantics
$$𝒩:(\{1\},)(,\{1\}).|\psi |\psi \text{ ,}$$
(31)
which immediately is the denotational semantics as well.
Consider the following network definition, which, as we derive explicitly below, implements the teleportation protocol.
$$\begin{array}{cc}\hfill 𝐚& =𝐀:\{1,2\}.[(𝚌!s_2s_1).M_{12}^{0,0}]\hfill \\ \hfill 𝐛& =𝐁:\{3\}.[X_3^{x_2}Z_3^{x_1}.(𝚌\mathrm{?}x_2x_1)]\hfill \\ \hfill 𝒩_{TP}& =𝐚|𝐛E_{23}\hfill \end{array}$$
(32)
To derive the semantics of the above network, note that for the first step there is only one possibility, namely that agent $`𝐀`$ executes the local Bell measurement $`M_{12}^{0,0}`$. The latter requires a local quantum input from $`𝐀`$, namely the qubit $`|\psi `$ that needs to be teleported. This is clearly the case since the pattern applies to qubit $`1`$, which is not part of the system state $`E_{23}`$. So by rule (9) we need to apply the pattern to the first two qubits of $`|\psi E_{23}`$. Using first rule (9)and writing $`\mathrm{\Gamma }_𝐚=\mathrm{}[s_2s_1j_2j_1]`$, we derive
$$\frac{|\psi _1E_{23},M_{12}^{0,0}_{1/4}X^{j_2}Z^{j_1}|\psi _3,\mathrm{\Gamma }_𝐚}{E_{23},\mathrm{},𝐚_{1/4}X^{j_2}Z^{j_1}|\psi _3,\mathrm{\Gamma }_𝐚,𝐀.(𝚌!s_2s_1)}\text{ ,}$$
(33)
where for each of the values of $`j_2j_1`$ the transition occurs with the same probability of $`1/4`$. This reduction fires within the context of (12), which we do not write out explicitly. The next step is a classical rendez-vous between both agents (i.e. Alice calls Bob), as per rule (10). Defining $`\mathrm{\Gamma }_𝐛=\mathrm{}[x_2x_1j_2j_1]`$, we get
$$\begin{array}{cc}\hfill X^{j_2}Z^{j_1}|\psi _3(\mathrm{\Gamma }_𝐚,𝐀& .(𝚌!s_2s_1)|\mathrm{},𝐁:\{3\}.[X_3^{x_2}Z_3^{x_1}.(𝚌\mathrm{?}x_2x_1)]\hfill \\ & \mathrm{\Gamma }_𝐚,𝐀|\mathrm{\Gamma }_𝐛,𝐁:\{3\}.X_3^{x_2}Z_3^{x_1})\text{ .}\hfill \end{array}$$
(34)
The last step of the computation is the execution of a local pattern by agent $`𝐁`$, as follows.
$$X^{j_2}Z^{j_1}|\psi _3,\mathrm{\Gamma }_𝐚,𝐀|\mathrm{\Gamma }_𝐛,𝐁:\{3\}.X_3^{x_2}Z_3^{x_1}|\psi _3,\mathrm{\Gamma }_𝐚,𝐀|\mathrm{\Gamma }_𝐛,𝐁:\{3\}$$
(35)
The only probabilistic transition is the first one, due to the Bell measurement. However, we see that the four branches lead to identical final system state and agents specifications. Furthermore, since there is no classical input or output, we can trace out the different local states. Thus, adding the probabilities as specified in Def. 2.7, we find that for any input $`|\psi `$, we have
$$𝒩_{TP}:(\{1,2\},\{3\})(,\{3\}).|\psi |\psi $$
(36)
In other words, we find that $`𝒩_{TP}𝒩`$.
Note that by linearity, the above derivation also works for mixed states.
We have shown that a direct quantum channel is operationally equivalent to the teleportation protocol. Suppose however, that two agents wish to exchange a qubit whilst they are contained in a larger network of agents. Can we say anything about the equivalence of both procedures in this context? By Prop. 28, we know that the entanglement with the larger system is conserved. Indeed, since teleportation just implements an identity channel, applying teleportation to one half of the mixed state $`\rho _{AC}`$ results in the state $`\rho _{BC}`$. Furthermore, suppose agent $`𝐀`$ wants to send an $`n`$-qubit entangled state to $`𝐁`$. Then by the same result, $`𝐀`$ can just apply the teleportation protocol $`n`$ times, and, since entanglement is conserved, the state is transferred unchanged. It remains to be seen whether the conservation of correlations between agents can also be employed in higher-level applications. Considering the fact that shared entanglement provides much of the extra power in distributed quantum systems, this behavior seems a promising primitive.
## 4 Conclusion
In this paper, we develop a formal model for distributed measurement-based quantum computations. We adopt an agent-based view, such that computations are described locally where possible. Because the network quantum state is in general entangled, we need to model it as a global structure, reminiscent of global memory in classical agent systems. Local quantum computations are described as measurement patterns. Since measurement-based quantum computation is inherently distributed, this allows us to extend naturally several concepts of the measurement calculus (MC) \[DKP04\], a formal model for such computations. Just as in MC, we aim at defining an *assembly language*, i.e. we assume that computations are well-defined and do not concern ourselves with verification techniques. The operational semantics for systems of agents is given by a probabilistic transition system, and we define operational equivalence in a way that it corresponds to the notion of *bisimilarity*. The denotational semantics is given by a set of quantum operations, together with type information which determines the localization of resources. Both forms of semantics are proved to be equivalent, and we define a notion of network composition such that the semantics is preserved with respect to this operation. Moreover, we show that within the larger entanglement contexts, the semantics is also conserved. With this in place, we prove that teleportation is bisimilar to a direct quantum channel, and, by the abovementioned result, this also holds within the context of a larger network. Though the proof is quite simple, it is important within the context of agent systems. Indeed, the possibility of inheriting agent correlations via teleportation means that, for example, collaborating agents are not cut off from each other when part of the shared data is transferred. Rather, the correlations are preserved in an oblivious manner. That this is the case if qubits are transported physically is clear, but that this remains so even if a general protocol is employed is maybe more surprising. It remains to be investigated how this preservation of the entanglement context can be exploited in more general situations. Also, there are other interesting situations to be investigated, such as the so-called *channel inequalities*, and more elaborate communication protocols. This is the subject of current investigations. |
warning/0506/hep-th0506029.html | ar5iv | text | # General Non-Extremal Rotating Black Holes in Minimal Five-Dimensional Gauged Supergravity
(June 2, 2005)
## Abstract
We construct the general solution for non-extremal charged rotating black holes in five-dimensional minimal gauged supergravity. They are characterised by four non-trivial parameters, namely the mass, the charge, and the two independent rotation parameters. The metrics in general describe regular rotating black holes, providing the parameters lie in appropriate ranges so that naked singularities and closed timelike curves (CTC’s) are avoided. We calculate the conserved energy, angular momenta and charge for the solutions, and show how supersymmetric solutions arise in a BPS limit. These have naked CTC’s in general, but for special choices of the parameters we obtain new regular supersymmetric black holes or smooth topological solitons.
preprint: MIFP-05-13 UPR-1125-T hep-th/0506029
The discovery of the remarkable AdS/CFT correspondence showed that bulk properties of solutions in the five-dimensional gauged supergravities that result from compactification of the type IIB string are related to properties of strongly-coupled conformal field theories on the four-dimensional boundary of five-dimensional anti-de Sitter spacetime mald ; guklpo ; wit . It therefore becomes of great importance to study the solutions of the five-dimensional gauged supergravity theories. One of the most important classes of such solutions are those that describe black holes in five dimensions. In particular, it has been argued that the boundary conformal field theory dual to rotating five-dimensional black holes should describe a system in a four-dimensional rotating Einstein universe hawhuntay .
The rotating five-dimensional solutions found in hawhuntay were neutral Kerr-(anti)-de Sitter black holes. In order to be able to make contact with supersymmetric BPS configurations, for which the AdS/CFT correspondence is more solidly founded, it is of considerable interest to generalise the neutral solutions to include electric charge too. In the analogous problem in ungauged supergravity, it is straightforward to generate charged solutions from neutral ones, by using the global symmetries of the ungauged supergravities as solution-generating transformations. By this means, the general charged rotating black holes of five-dimensional ungauged supergravity were obtained in cvetyoum , starting from the neutral rotating Ricci-flat black holes found in myerperr . For the solutions in gauged supergravity there are no surviving global symmetries that can be used to provide solution-generating transformations, and one has little option but to resort to brute-force calculations, starting from an appropriate ansatz, to construct the charged rotating solutions. One way to simplify the problem is to specialise to the case where the two independent rotation parameters of the generic five-dimensional rotating black hole are set equal, since this reduces the problem from cohomogeneity-2, with partial differential equations, to cohomogeneity-1, with ordinary differential equations. Supersymmetric rotating black holes with two equal angular momenta were obtained in gutreal , and it was shown in gutreal that the rotation is necessary for the solution to be free of naked singularities and CTC’s. The non-extremal charged rotating solutions of gauged five-dimensional supergravity with equal rotation parameters were constructed in d5gauge1 ; d5gauge2 . Recently, some special cases involving unequal rotation parameters were also constructed, in d5gauge3 . However, these latter arose as solutions of $`𝒩=2`$ gauged supergravity coupled to two vector multiplets, with a specific relation between the three electric charges, and did not, in general, admit a specialisation to solutions of pure minimal $`𝒩=2`$ gauged supergravity. The purpose of this letter is to present the general solution for charged rotating non-extremal black holes in minimal five-dimensional gauged supergravity, with independent rotation parameters in the two orthogonal 2-planes.
We have found the general solution for charged rotating black holes in five-dimensional minimal gauged supergravity, with unequal angular momenta, by a process involving a considerable amount of trial and error, followed by an explicit verification that the equations of motion are satisfied. In doing this, we have been guided by the previously-obtained special case found in d5gauge1 , where the two angular momenta were set equal, and the general charged rotating solutions in ungauged minimal supergravity, which are contained within the results in cvetyoum . In this letter we begin by presenting our new solutions, and then we calculate the conserved angular momenta and electric charge. By integrating the first law of thermodynamics, we also obtain the conserved mass, or energy, of the solutions. By considering the conditions under which the anticommutator of supercharges in the AdS superalgebra has zero eigenvalues, we then show how a BPS limit of our general non-extremal solutions gives rise to new supersymmetric configurations. These include new supersymmetric rotating black holes, with two independently-specifiable angular momenta, and new topological solitons that are non-singular on complete manifolds.
In terms of Boyer-Lindquist type coordinates $`x^\mu =(t,r,\theta ,\varphi ,\psi )`$ that are asymptotically static (i.e. the coordinate frame is non-rotating at infinity), we find that the metric and gauge potential for our new rotating solutions can be expressed as
$`ds^2`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }_\theta [(1+g^2r^2)\rho ^2dt+2q\nu ]dt}{\mathrm{\Xi }_a\mathrm{\Xi }_b\rho ^2}}+{\displaystyle \frac{2q\nu \omega }{\rho ^2}}`$ (1)
$`+{\displaystyle \frac{f}{\rho ^4}}\left({\displaystyle \frac{\mathrm{\Delta }_\theta dt}{\mathrm{\Xi }_a\mathrm{\Xi }_b}}\omega \right)^2+{\displaystyle \frac{\rho ^2dr^2}{\mathrm{\Delta }_r}}+{\displaystyle \frac{\rho ^2d\theta ^2}{\mathrm{\Delta }_\theta }}`$
$`+{\displaystyle \frac{r^2+a^2}{\mathrm{\Xi }_a}}\mathrm{sin}^2\theta d\varphi ^2+{\displaystyle \frac{r^2+b^2}{\mathrm{\Xi }_b}}\mathrm{cos}^2\theta d\psi ^2,`$
$`A`$ $`=`$ $`{\displaystyle \frac{\sqrt{3}q}{\rho ^2}}\left({\displaystyle \frac{\mathrm{\Delta }_\theta dt}{\mathrm{\Xi }_a\mathrm{\Xi }_b}}\omega \right),`$ (2)
where
$`\nu `$ $`=`$ $`b\mathrm{sin}^2\theta d\varphi +a\mathrm{cos}^2\theta d\psi ,`$
$`\omega `$ $`=`$ $`a\mathrm{sin}^2\theta {\displaystyle \frac{d\varphi }{\mathrm{\Xi }_a}}+b\mathrm{cos}^2\theta {\displaystyle \frac{d\psi }{\mathrm{\Xi }_b}},`$
$`\mathrm{\Delta }_\theta `$ $`=`$ $`1a^2g^2\mathrm{cos}^2\theta b^2g^2\mathrm{sin}^2\theta ,`$
$`\mathrm{\Delta }_r`$ $`=`$ $`{\displaystyle \frac{(r^2+a^2)(r^2+b^2)(1+g^2r^2)+q^2+2abq}{r^2}}2m,`$
$`\rho ^2`$ $`=`$ $`r^2+a^2\mathrm{cos}^2\theta +b^2\mathrm{sin}^2\theta ,`$
$`\mathrm{\Xi }_a`$ $`=`$ $`1a^2g^2,\mathrm{\Xi }_b=1b^2g^2,`$
$`f`$ $`=`$ $`2m\rho ^2q^2+2abqg^2\rho ^2.`$ (3)
A straightforward calculation shows that these configurations solve the equations of motion of minimal gauged five-dimensional supergravity, which follow from the Lagrangian
$$=(R+12g^2)\text{1}\mathrm{l}\frac{1}{2}FF+\frac{1}{3\sqrt{3}}FFA,$$
(4)
where $`F=dA`$, and $`g`$ is assumed to be positive, without loss of generality.
For some purposes, it is useful to note that the non-vanishing metric components are given by
$`g_{00}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }_\theta (1+g^2r^2)}{\mathrm{\Xi }_a\mathrm{\Xi }_b}}+{\displaystyle \frac{\mathrm{\Delta }_\theta ^2(2m\rho ^2q^2+2abqg^2\rho ^2)}{\rho ^4\mathrm{\Xi }_a^2\mathrm{\Xi }_b^2}},`$
$`g_{03}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }_\theta [a(2m\rho ^2q^2)+bq\rho ^2(1+a^2g^2)]\mathrm{sin}^2\theta }{\rho ^4\mathrm{\Xi }_a^2\mathrm{\Xi }_b}},`$
$`g_{04}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }_\theta [b(2m\rho ^2q^2)+aq\rho ^2(1+b^2g^2)]\mathrm{cos}^2\theta }{\rho ^4\mathrm{\Xi }_b^2\mathrm{\Xi }_a}},`$
$`g_{33}`$ $`=`$ $`{\displaystyle \frac{(r^2+a^2)\mathrm{sin}^2\theta }{\mathrm{\Xi }_a}}+{\displaystyle \frac{a[a(2m\rho ^2q^2)+2bq\rho ^2]\mathrm{sin}^4\theta }{\rho ^4\mathrm{\Xi }_a^2}},`$
$`g_{44}`$ $`=`$ $`{\displaystyle \frac{(r^2+b^2)\mathrm{cos}^2\theta }{\mathrm{\Xi }_b}}+{\displaystyle \frac{b[b(2m\rho ^2q^2)+2aq\rho ^2]\mathrm{cos}^4\theta }{\rho ^4\mathrm{\Xi }_b^2}},`$
$`g_{34}`$ $`=`$ $`{\displaystyle \frac{[ab(2m\rho ^2q^2)+(a^2+b^2)q\rho ^2]\mathrm{sin}^2\theta \mathrm{cos}^2\theta }{\rho ^4\mathrm{\Xi }_a\mathrm{\Xi }_b}},`$
$`g_{11}`$ $`=`$ $`{\displaystyle \frac{\rho ^2}{\mathrm{\Delta }_r}},g_{22}={\displaystyle \frac{\rho ^2}{\mathrm{\Delta }_\theta }}.`$ (5)
The Killing vector
$$\mathrm{}=\frac{}{t}+\mathrm{\Omega }_a\frac{}{\varphi }+\mathrm{\Omega }_b\frac{}{\psi }$$
(6)
becomes null on the outer Killing horizon at $`r=r_+`$, the largest positive root of $`\mathrm{\Delta }_r=0`$, where the angular velocities on the horizon are given by
$`\mathrm{\Omega }_a`$ $`=`$ $`{\displaystyle \frac{a(r_+^2+b^2)(1+g^2r_+^2)+bq}{(r_+^2+a^2)(r_+^2+b^2)+abq}},`$
$`\mathrm{\Omega }_b`$ $`=`$ $`{\displaystyle \frac{b(r_+^2+a^2)(1+g^2r_+^2)+aq}{(r_+^2+a^2)(r_+^2+b^2)+abq}}.`$ (7)
One can then easily evaluate the surface gravity
$$\kappa =\frac{r_+^4[(1+g^2(2r_+^2+a^2+b^2)](ab+q)^2}{r_+[(r_+^2+a^2)(r_+^2+b^2)+abq]},$$
(8)
and hence the Hawking temperature $`T=\kappa /(2\pi )`$. The entropy is given by
$$S=\frac{\pi ^2[(r_+^2+a^2)(r_+^2+b^2)+abq]}{2\mathrm{\Xi }_a\mathrm{\Xi }_br_+}.$$
(9)
The angular momenta can be evaluated from the Komar integrals $`J=1/(16\pi )_{S^3}dK`$, where $`K=/\varphi `$ or $`K=/\psi `$, yielding
$`J_a`$ $`=`$ $`{\displaystyle \frac{\pi [2am+qb(1+a^2g^2)]}{4\mathrm{\Xi }_a^2\mathrm{\Xi }_b}},`$
$`J_b`$ $`=`$ $`{\displaystyle \frac{\pi [2bm+qa(1+b^2g^2)]}{4\mathrm{\Xi }_b^2\mathrm{\Xi }_a}}.`$ (10)
The electric charge follows from the Gaussian integral $`Q=1/(16\pi )_{S^3}(FFA/\sqrt{3})`$, yielding
$$Q=\frac{\sqrt{3}\pi q}{4\mathrm{\Xi }_a\mathrm{\Xi }_b}.$$
(11)
Using the technique introduced in gibperpop , the easiest way to calculate the conserved mass, or energy, is to integrate the first law of thermodynamics $`dE=TdS+\mathrm{\Omega }_adJ_a+\mathrm{\Omega }_bdJ_b+\mathrm{\Phi }dQ`$, where $`\mathrm{\Phi }=\mathrm{}^\mu A_\mu `$ is the electrostatic potential on the horizon. Doing this, we find
$$E=\frac{m\pi (2\mathrm{\Xi }_a+2\mathrm{\Xi }_b\mathrm{\Xi }_a\mathrm{\Xi }_b)+2\pi qabg^2(\mathrm{\Xi }_a+\mathrm{\Xi }_b)}{4\mathrm{\Xi }_a^2\mathrm{\Xi }_b^2}.$$
(12)
The BPS limit can be found by looking at the eigenvalues of the Bogomol’nyi matrix coming from the anticommutators of the supercharges, as discussed in cvgilupo . Thus we have BPS solutions if
$$EgJ_agJ_b\sqrt{3}Q=0.$$
(13)
From the expressions derived above for $`(E,J_a,J_b,Q)`$, we find that the BPS limit is achieved if
$$q=\frac{m}{1+(a+b)g}.$$
(14)
The supersymmetry of the solutions in this limit can be confirmed by calculating the norm of the Killing vector
$$K_+\frac{}{t}+g\frac{}{\varphi }+g\frac{}{\psi },$$
(15)
which, as discussed in cvgilupo , arises as the square of the Killing spinor $`\eta `$, in the sense that $`K_+^\mu =\overline{\eta }\gamma ^\mu \eta `$. We find that its norm is given by
$$K_+^2=\frac{[hm(1+ag\mathrm{cos}^2\theta +bg\mathrm{sin}^2\theta )]^2}{h^2},$$
(16)
where
$$h=(1+ag)(1+bg)[1+(a+b)g]\rho ^2.$$
(17)
Thus indeed the norm of $`K_+`$ is, as it should be since it has a spinorial square root, manifestly negative definite. The fraction of supersymmetry preserved is in general $`\frac{1}{4}`$, except when $`a=b`$, in which case, the preserved supersymmetry is doubled to become $`\frac{1}{2}`$. The latter solution was previously obtained in ks .
We now discuss the global structure of the rotating AdS<sub>5</sub> black hole. To do this, we first note that the metric can be expressed as
$`ds^2`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }_r\mathrm{\Delta }_\theta r^2\mathrm{sin}^22\theta }{4\mathrm{\Xi }_a^2\mathrm{\Xi }_b^2B_\varphi B_\psi }}dt^2+\rho ^2({\displaystyle \frac{dr^2}{\mathrm{\Delta }_r}}+{\displaystyle \frac{d\theta ^2}{\mathrm{\Delta }_\theta }})`$
$`+B_\psi (d\psi +v_1d\varphi +v_2dt)^2+B_\varphi (d\varphi +v_3dt)^2,`$
where the functions $`B_\varphi `$, $`B_\psi `$, $`v_1`$, $`v_2`$ and $`v_3`$ can be straightforwardly found by comparing (General Non-Extremal Rotating Black Holes in Minimal Five-Dimensional Gauged Supergravity) with the metric in (1). The absence of naked closed timelike curves (CTC’s) requires that $`B_\varphi `$ and $`B_\psi `$ be non-negative outside the horizon. We shall focus on the discussion of supersymmetric solutions, satisfying the condition (14). It can be seen from (16) that the identity
$`{\displaystyle \frac{\mathrm{\Delta }_r\mathrm{\Delta }_\theta r^2\mathrm{sin}^22\theta }{4\mathrm{\Xi }_a^2\mathrm{\Xi }_b^2B_\varphi B_\psi }}+B_\psi (v_2+g+gv_1)^2+B_\varphi (v_3+g)^2`$
$`={\displaystyle \frac{[hm(1+ag\mathrm{cos}^2\theta +bg\mathrm{sin}^2\theta )]^2}{h^2}}`$ (19)
holds. It follows that in general, at the Killing horizon where $`\mathrm{\Delta }_r=0`$, we have $`B_\varphi B_\psi <0`$, implying the existence of naked CTC’s. There are two special cases where naked CTC’s can be avoided, leading to either supersymmetric black holes or topological solitons:
Supersymmetric black holes: The first way to avoid naked CTC’s is if the right-hand side of (19) vanishes on the Killing horizon. This occurs when the parameters in the supersymmetric solutions satisfy the further restriction
$$gm=(a+b)(1+ag)(1+bg)(1+ag+bg).$$
(20)
Remarkably, when this extra condition is satisfied, the function $`\mathrm{\Delta }_r`$ has a double root; $`\mathrm{\Delta }_r`$ is now given by
$$\mathrm{\Delta }_r=r^2(r^2r_0^2)^2[g^2r^2+(1+ag+bg)^2],$$
(21)
where $`r_0^2=g^1(a+b+abg)`$. At the Killing horizon $`r=r_0`$, we find that the determinant of the metric in the $`(\theta ,\varphi ,\psi )`$ directions is given by
$$detg_{(\theta ,\varphi ,\psi )}=\frac{(a+b)^2(a+b+abg)\mathrm{sin}^22\theta }{4g^3(1ag)^2(1bg)^2}$$
(22)
This implies that naked CTC’s are avoided if the remaining free parameters $`a`$ and $`b`$ satisfy the inequality
$$a+b+abg>0.$$
(23)
The Killing horizon $`r=r_0`$ is then the event horizon of a well-defined supersymmetric black hole that is regular on and outside the event horizon. The occurrence of the double-root of $`\mathrm{\Delta }_r`$ at $`r=r_0`$ implies that the black hole has zero temperature. The various conserved and thermodynamic quantities for these new supersymmetric black holes are given by
$`E`$ $`=`$ $`{\displaystyle \frac{\pi (a+b)}{4g(1ag)^2(1bg)^2}}((1ag)(1bg)`$
$`+(1+ag)(1+bg)(2agbg)),`$
$`S`$ $`=`$ $`{\displaystyle \frac{\pi ^2(a+b)\sqrt{a+b+abg}}{2g^{3/2}(1ag)(1bg)}},`$
$`J_a`$ $`=`$ $`{\displaystyle \frac{\pi (a+b)(2a+b+abg)}{4g(1ag)^2(1bg)}},`$
$`J_b`$ $`=`$ $`{\displaystyle \frac{\pi (a+b)(a+2b+abg)}{4g(1ag)(1bg)^2}},`$
$`Q`$ $`=`$ $`{\displaystyle \frac{\pi \sqrt{3}(a+b)}{4g(1ag)(1bg)}}.`$ (24)
Note that supersymmetric black holes cannot arise when $`a=b`$. For $`a=b`$, our new supersymmetric black hole solutions, which for general $`a`$ and $`b`$ have cohomogeneity 2, become cohomogeneity 1; these special cases were previously obtained in gutreal .
Topological solitons: The second way to avoid naked CTC’s is if $`B_\varphi =0`$ at $`r=r_0`$. This can happen when the free parameters in the general supersymmetric solutions obey the further restriction
$`m`$ $`=`$ $`(1+ag)(1+bg)(1+ag+bg)`$ (25)
$`\times (2a+b+abg)(a+2b+abg).`$
Now $`r_0`$, the outer root of $`\mathrm{\Delta }_r`$, is given by
$$r_0^2=(a+b+abg)^2.$$
(26)
Defining a new radial coordinate $`R=r^2r_0^2`$, we find that the metric describes a smooth topological soliton, with $`R`$ running from 0 to $`\mathrm{}`$. The requirement of the absence of a conical singularity when $`B_\varphi `$ vanishes at $`R=0`$ implies the quantisation condition
$$\frac{(a+b+abg)(3+5ag+5bg+3abg^2)}{(1ag)(a+2b+abg)}=1.$$
(27)
In the cohomogeneity-1 special cases $`a=b`$ or $`a=b`$, these toplogical solitons are encompassed within the soliton solutions obtained in cvgilupo .
Aside from the above two possibilities, the supersymmetric solutions in general have naked CTC’s. As in the examples discussed in cvgilupo ; d5gauge3 , a conical singularity at the Killing horizon can be avoided by periodically identifying the asymptotic time coordinate $`t`$ with an appropriate period. However, if the Killing horizon is associated with a double root of $`\mathrm{\Delta }_r`$, then such an identification is unnecessary, analogous to the ungauged rotating solution obtained in bmpv . The geodesic analysis of analogous time machines can be found in gibher ; cks .
In the general case where the charged rotating metrics that we have found are non-extremal, they describe regular black holes provided the parameters lie in appropriate ranges that are easily determinable using the same techniques we have used above for analysing the BPS limits.
As discussed in hawhuntay , rotating black hole solutions in five-dimensional gauged supergravity provide backgrounds whose AdS/CFT duals describe four-dimensional field theories in the rotating Einstein universe on the boundary of anti-de Sitter spacetime. With the general solutions in minimal gauged supergravity that we have now found, this aspect of the AdS/CFT correspondence can be studied in a framework that also allows one to take a BPS or near-BPS limit, where the mapping from the bulk to the boundary is better controlled. In particular, it is of great interest to provide the microscopic interpretation from the boundary CFT for the entropy (24) of the supersymmetric black holes with two general rotations. We plan to report further on these considerations in forthcoming work.
Acknowledgements:
We thank Gary Gibbons for useful discussions. C.N.P. thanks the Relativity and Cosmology group in DAMTP, Cambridge, for hospitality during the course of this work. Research supported in part by DOE grants DE-FG02-95ER40893 and DE-FG03-95ER40917, NSF grant INTO3-24081, and (M.C.) the University of Pennsylvania Research Foundation Award and the Fay R. and Eugene L.Langberg Chair. |
warning/0506/gr-qc0506120.html | ar5iv | text | # Distortion of Schwarzschild-anti-de Sitter black holes to black strings
## I INTRODUCTION
Stationary black holes in spacetimes with a negative cosmological constant $`\mathrm{\Lambda }`$ (which we call hereafter anti-de Sitter black holes for abbreviation) have been a subject of current interest, in particular, motivated by the conjecture of the anti-de Sitter/Conformal field theory-correspondence Mal98 , according to which the emergence of black hole thermodynamics is interpreted in terms of thermal states of the dual conformal field theory Witt1 ; Witt2 . It is also interesting that the vacuum Einstein equations with the cosmological term admit stationary exact solutions which have the event horizon with various topologies even in four dimensions Lemos1 ; Lemos2 ; Huang ; Lemos3 ; Cai ; Ami ; Brill ; Vanzo97 ; Klemn98 . Namely, the spatial section of the event horizon can be Einstein’s manifold with positive, zero or negative scalar curvature corresponding to $`k=+1,0,1`$. The thermodynamic properties of such anti-de Sitter (AdS) black holes have been investigated in HP for $`k=+1`$ and in Brill ; Vanzo97 for $`k=0,1`$. Though the higher-dimensional generalization has been extensively studied (see, for example, Mann ; Bir ), in this paper we focus our attention on the four-dimensional black holes.
For $`k=1`$ the event horizon has a spherical topology ($`S^2`$), and the Schwarzschild-anti-de Sitter (SAdS) solution is well-known as a typical example of spherically symmetric anti-de Sitter black holes. On the other hand, the horizon topology corresponding to zero scalar curvature (i.e., $`k=0`$) may be planar ($`R^2`$), cylindrical ($`R^1\times S^1`$) or toruslike ($`S^1\times S^1`$) according to the compactification scheme for the two-dimensional spatial section. Hyperbolic (sometimes called topological) black holes represented by the $`k=1`$ solutions may have a negative mass, and the (in)stability becomes a subtle problem Gib ; Neu .
It is remarkable that conformal (compactified) spatial infinity of a stationary AdS black hole spacetime has the same topology as the event horizon. If the horizon topology is fixed, a black hole solution satisfying the stationary vacuum Einstein equations with $`\mathrm{\Lambda }<0`$ may be unique And ; Kod . However, this uniqueness theorem holds only if the boundary metric at conformal spatial infinity is required to be the Einstein metric with a constant scalar curvature. In fact, axisymmetric static perturbations of the SAdS solution have been explicitly presented in Yos . The key result of the perturbative analysis is that a static small distortion from spherical symmetry (which is regular at the event horizon) does not vanish even at spatial infinity. Namely, we obtain a non-uniform spherical 2-surface at spatial infinity, where an arbitrary function $`f(\theta )`$ dependent on the zenithal angle $`\theta `$ appears. This should be compared to vacuum black hole (namely, Schwarzschild and Kerr) solutions with $`\mathrm{\Lambda }=0`$, for which any static distortion regular at the horizon must diverge at spatial infinity (see Ger ; Tom for the exact solutions of distorted black holes). The role of a negative cosmological constant is to give a finite distortion to the boundary metric on $`S^2`$ at spatial infinity. Interestingly, we note the one-to-one correspondence between the boundary data represented by the function $`f(\theta )`$ and the bulk spacetime geometry from the horizon to spatial infinity.
In this paper we consider a large distortion of a SAdS black hole as an extension of the analysis in Yos to a non-perturbative case. Our purpose is to construct a family of distorted anti-de Sitter black hole solutions connecting the $`k=1`$ (SAdS) black hole solution to the $`k=0`$ black string solution with a cylindrical horizon topology. This will be useful for studying the important area of AdS black hole physics concerning the quasi-static or thermodynamic evolution which may be accompanied with a large distortion from spherical symmetry.
For mathematical convenience our investigation is limited to the large mass domain such that $`mr_A`$, where $`m`$ is the SAdS mass parameter and $`r_A\sqrt{3/\mathrm{\Lambda }}`$ is the anti-de Sitter radius. (Hereafter we use units such that $`c=G=\mathrm{}=k_B=1`$.) Because in the large mass domain the inequality $`r_Ar_Hm`$ holds for the SAdS horizon radius $`r_H`$, the effect of $`\mathrm{\Lambda }`$ becomes significant even near the horizon, and the horizon 2-surface is distorted in the nearly same manner as the 2-surface at spatial infinity, as was shown in Yos . This allows us to obtain easily a non-perturbative distortion of the metric in Sec. II. Then, in Sec. III, we estimate the mass of distorted AdS black holes, using the definition proposed by Ashtekar and Magnon AM applicable to spacetimes with a distorted boundary metric at spatial infinity. We find the condition for distortion keeping the Astekar-Magnon mass equal to the SAdS mass parameter $`m`$. It is also possible to obtain the Euclidean gravitational action, using the boundary counterterm technique Emp ; Mann2 . The leading-order calculation in the large mass domain clearly shows that the thermodynamic mass energy is equal to $`m`$. The entropy and the temperature are given as functions of $`m`$ in the same way as the SAdS case. We can conclude that any effect due to the horizon distortion becomes thermally insignificant in the large mass domain. In Sec. IV, we present an explicit example representing a black hole distortion connecting to a uniform black string. Such a distortion is shown to satisfy the so-called Penrose inequality for the mass and the horizon area, of which the validity has been discussed in asymptotically flat spacetimes (see, for example, Bray ). We also mention, in terms of the hoop conjecture proposed by Thorne Thorne , the possible existence of an extremal state of distorted (non-rotating) AdS black holes.
## II STATIC AXISYMMETRIC DISTORTION
Without loss of generality the static axisymmetric metric for describing a black hole with a spherical horizon topology is given by
$$ds^2=e^{2\nu }dt^2+e^{2\mu }dr^2+r^2e^{2\psi }(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2),$$
(1)
where the functions $`\nu `$, $`\mu `$ and $`\psi `$ depend only on $`r`$ and $`\theta `$. For the SAdS metric we obtain
$$e^{2\nu _0}=e^{2\mu _0}=1\frac{2m}{r}+(\frac{r}{r_A})^2,e^{2\psi _0}=1.$$
(2)
If we consider a static axisymmetric distortion of the SAdS metric according to a usual perturbative scheme, we have up to the first order
$$\nu =\nu _0+ϵ\nu _1,\mu =\mu _0+ϵ\mu _1,\psi =ϵ\psi _1,$$
(3)
with a small parameter $`ϵ`$. Using the vacuum Einstein equations $`R_{ab}=\mathrm{\Lambda }g_{ab}`$ with a negative cosmological constant $`\mathrm{\Lambda }`$, the first-order perturbations $`\nu _1(r,\theta )`$, $`\mu _1(r,\theta )`$ and $`\psi _1(r,\theta )`$ have been studied in Yos . For each multiple component given by Legendre’s polynomial $`P_l(\mathrm{cos}\theta )`$ with $`ł2`$ we obtain
$$\nu _1=\mu _1=H^{(1)}(r)P_l(\mathrm{cos}\theta ),\psi _1=K^{(1)}(r)P_l(\mathrm{cos}\theta ).$$
(4)
In this paper we focus our interest on the large mass domian such that $`mr_A`$, in which the radial functions are given by
$$H^{(1)}=\frac{3x}{x^2+x+1}\frac{1}{x}+O(\delta ),K^{(1)}=\frac{6}{\delta (l^2+l2)}+O(1),$$
(5)
where $`\delta r_H/2m1`$, $`xr/r_H`$, and the horizon radius $`r_H`$ is approximately written by $`r_H(2mr_A^2)^{1/3}`$ (see Yos for the derivation). The key point of the result (5) is that the functional form of $`H^{(1)}`$ does not depend on $`l`$ in the leading-order calculation with respect to the new small parameter $`\delta `$. This allows us to write the first-order perturbation $`\nu _1`$ as follows,
$$\nu _1=\mu _1=(\frac{3x}{x^2+x+1}\frac{1}{x})h(\theta ),$$
(6)
where $`h(\theta )`$ is an arbitrary function of $`\theta `$. Note also that the radial function $`K^{(1)}`$ does not depend on $`r`$, though the constant value depends on $`l`$. The first-order perturbation $`\psi _1`$ can be written by
$$\psi _1=\frac{1}{\delta }f(\theta ),$$
(7)
where the function $`f(\theta )`$ must satisfy the relation
$$6h=\frac{d^2f}{d\theta ^2}+\mathrm{cot}\theta \frac{df}{d\theta }+2f.$$
(8)
This linear analysis clearly shows that one arbitrary function (i.e., $`f`$ or $`h`$) of $`\theta `$ appears as a hair of distorted AdS black holes. In particular, for SAdS black holes with large mass $`m`$, any axisymmetric distortion of a spherical 2-surface ($`S^2`$) is represented by the function $`f(\theta )`$ independent of the radial coordinate $`r`$. Namely, we obtain the same distortion of $`S^2`$ in the whole range from the horizon $`r=r_H`$ to spatial infinity $`r\mathrm{}`$.
Now we consider a non-perturbative extension of the results obtained by the linear analysis. It is clear from Eqs. (6) and (7) that if the distortion parameter $`ϵ`$ is chosen to be $`ϵ=\delta 1`$, the linear approximation for $`ϵ\psi _1`$ breaks down, while the perturbation $`ϵ\nu _1=ϵ\mu _1`$ remains small. Hence, the expansion of the metric (1) with respect to the small parameter $`\delta `$ will allow us to give a non-perturbative distortion only to the metric function $`\psi `$ defined on $`S^2`$.
Note that the SAdS metric (2) can be written by
$$e^{2\nu _0}=e^{2\mu _0}=\frac{1}{\delta }(x^2\frac{1}{x})+1$$
(9)
in the range $`x1`$. Then the expansion of the metric function $`\nu `$ with respect to $`\delta `$ (i.e., $`\nu =\nu _0+\delta \nu _1`$) should be rewritten into the form
$$e^{2\nu }e^{2\mu }\frac{1}{\delta }(x^2\frac{1}{x})\{1+2\delta \nu _1+\delta (x^2\frac{1}{x})^1\}.$$
(10)
The approximation such that $`e^{2\delta \nu _1}1+2\delta \nu _1`$ is still possible in the same way as the linear analysis. However, the leading term for the metric function $`\psi `$ should be
$$e^{2\psi }e^{2f(\theta )},$$
(11)
which does not allow the approximation such that $`e^{2\delta \psi _1}1+2\delta \psi _1`$. Nevertheless it is straightforward to see that the vacuum Einstein equations with a negative cosmological constant are satisfied up to the next-to leading-order calculation in the small $`\delta `$ domain, only if the nonlinear relation
$$6h1=e^{2f}(\frac{d^2f}{d\theta ^2}+\mathrm{cot}\theta \frac{df}{d\theta }1),$$
(12)
is required instead of Eq. (8). Of course, for a small distortion corresponding to the approximation $`|f|1`$ Eq. (12) reduces to Eq. (8).
We can treat any non-perturbative distortion of $`S^2`$ by the use of the arbitrary metric function $`f(\theta )`$, which also determines the small correction $`\delta h(\theta )`$ to the metric function $`\nu =\mu `$ through Eq. (12). If such a distortion is required to preserve the area of $`S^2`$, the function $`f`$ should satisfy the additional condition
$$\frac{1}{2}_0^\pi e^{2f(\theta )}\mathrm{sin}\theta d\theta =1,$$
(13)
from which Eq. (12) leads to the result that the mean value $`\overline{h}`$ vanishes if the distortion perturbation $`h(\theta )`$ is averaged over $`S^2`$, namely,
$$\overline{h}\frac{1}{2}_0^\pi e^{2f(\theta )}h(\theta )\mathrm{sin}\theta d\theta =0.$$
(14)
Let us remark that the area-preserving condition (13) is allowed without loss of generality. For some choice of $`f(\theta )`$ the integral in Eq. (13) may become equal to $`e^{2\overline{f}}`$ with a nonzero value $`\overline{f}`$. Then, we can consider the new distortion functions $`f^{}(\theta )`$ and $`h^{}(\theta )`$, using the transformations $`ff^{}=f\overline{f}`$ and $`hh^{}=e^{2\overline{f}}h+(1e^{2\overline{f}})/6`$. It is easy to check that these functions $`f^{}`$ and $`h^{}`$ can satisfy Eqs. (13) and (14). Further, as was previously mentioned, for a small distortion of $`S^2`$ (i.e., for $`|f|1`$) the distortion functions $`f(\theta )`$ and $`h(\theta )`$ may be given by Legendre’s polynomial $`P_l(\mathrm{cos}\theta )`$ with $`l2`$. We note that if the distortion is treated as a linear perturbation, Eqs. (13) and (14) should be regarded as the necessary conditions for $`f`$ and $`h`$. Hence, in the following section we will discuss mass and thermodynamic properties of distorted AdS black holes under the condition (13).
## III THE ASHTEKAR-MAGNON MASS AND THE ENTROPY
The distortion given by $`f(\theta )`$ may induce a black hole mass different from the SAdS mass parameter $`m`$. Though the mass is a basic quantity characterizing black hole states, the definition in spacetimes with a negative cosmological constant remains ambiguous as a problem to be investigated from various viewpoints. If the distortion from sherical symmetry vanishes at spatial infinity, one may use the so-called Abbott-Desser mass $`M_{AD}`$ AD . Unfortunately, as was shown in Yos , the calculation based on the Abbott-Deser method which depends on the choice of the background metric is not applicable to distorted AdS black holes. Hence, we adopt here the evaluation method proposed by Ashtekar and Magnon AM , for which any background subtraction is unnecessary. The Ashtelar-Magnon mass denoted by $`M_{AM}`$ is a conserved quantity defined in the conformally transformed spacetime with the metric $`\overline{g}_{ab}=\mathrm{\Omega }^2g_{ab}`$, where $`g_{ab}`$ is the physical metric (1) representing a AdS black hole. If the conformal factor $`\mathrm{\Omega }`$ is chosen to be $`\mathrm{\Omega }=1/r`$, it is given by
$$M_{AM}=\frac{r_A}{4}\underset{r\mathrm{}}{lim}_0^\pi \xi ^tr^2e^\nu [e^{2\mu }\{\nu ,_{rr}+\nu ,_r(\nu ,_r\mu ,_r)\}+\frac{e^{2\psi }}{r^2}\nu ,_\theta \mu ,_\theta \frac{1}{r_A^2}]e^{2\psi }\mathrm{sin}\theta d\theta ,$$
(15)
where $`\xi ^t`$ is a time component of the timelike Killing vector $`\xi ^a`$ on the conformal boundary at $`r\mathrm{}`$.
For distorted black holes obtained under the approximation $`\delta 1`$ in the previous section, we find that $`\nu _1=\mu _10`$ in the limit $`r\mathrm{}`$, and the conformal boundary metric is given by
$$d\overline{s}^2=\frac{1}{r_A^2}dt^2+e^{2f(\theta )}(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2).$$
(16)
Though the normalization of the Killing vector $`\xi ^a`$ on the boundary metric distorted by $`f`$ has been discussed in Yos , here we choose $`\xi ^t=1`$ in the same way as the SAdS spacetime, considering the condition (13) for $`f`$. To calculate Eq. (15) to the leading order in the small $`\delta `$ domain, it is sufficient to use the approximated SAdS form
$$e^{2\nu }e^{2\mu }\frac{1}{\delta }(x^2\frac{1}{x})$$
(17)
for $`\nu `$ and $`\mu `$ in the integrand, while we have $`\psi =f0`$ for the metric function on $`S^2`$. Then, the final result is simply
$$M_{AM}=m.$$
(18)
In spite of a large distortion of black hole geometry represented by $`f`$ a change of the black hole mass can remain very small, namely, $`(M_{AM}m)/m=O(\delta )`$.
Let us also check the effect of black hole distortion at the thermodynamic level. For this purpose we calculate the Euclidean gravitational action $`I`$, using the counterterm prescription applicable to spacetimes with a negative cosmological constant Emp ; Mann2 . For the leading-order calculation of thermodynamic quantities in the small $`\delta `$ domain, the gravitational action $`I`$ may be evaluated by giving the Euclidean version of the approximate metric (17) such as
$$ds^2=(\frac{r^2}{r_A^2}\frac{2m}{r})d\tau ^2+(\frac{r^2}{r_A^2}\frac{2m}{r})^1dr^2+r^2e^{2f(\theta )}(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2),$$
(19)
where $`\tau =it`$ is the Euclidean time with the period $`\beta `$ given by
$$\beta =\frac{2\pi }{3}(\frac{4r_A^4}{m})^{1/3}.$$
(20)
The Euclidean gravitational action $`I`$ contains three contributions denoted by
$$I=I_{bulk}+I_{surf}+I_{ct}.$$
(21)
The first two terms $`I_{bulk}`$ and $`I_{surf}`$ in Eq. (21) are the familiar classical action corresponding to the volume integral in the range $`r_Hrr_0`$ and the boundary integral at $`r=r_0`$, respectively, and we have
$$I_{bulk}=\frac{3}{8\pi r_A^2}d^4x\sqrt{g}=\frac{\beta }{2r_A^2}(r_0^3r_H^3),$$
(22)
and
$$I_{surf}=\frac{1}{8\pi }d^3x\sqrt{h}K=\frac{3\beta }{2r_A^2}(r_0^3+\frac{r_H^3}{2}),$$
(23)
where $`K`$ is the trace of the extrinsic curvature of the boundary with the metric $`h_{ab}`$, giving the determinant
$$\sqrt{h}=(\frac{r_0^2}{r_A^2}\frac{2m}{r_0})^{1/2}r_0^2e^{2f}\mathrm{sin}\theta .$$
(24)
Here, the counterterm $`I_{ct}`$ necessary for canceling the divergence of $`I`$ is given by the boundary integral
$$I_{ct}=\frac{1}{4\pi r_A}d^3x\sqrt{h}=\frac{\beta r_0^2}{r_A}(\frac{r_0^2}{r_A^2}\frac{2m}{r_0})^{1/2}.$$
(25)
Then, in the limit $`r_0\mathrm{}`$, we obtain
$$I=\frac{\beta }{2}m=\frac{\pi }{3}(2r_A^2m)^{2/3}.$$
(26)
By virtue of the condition (13) it is clear that no correction due to the distortion function $`f`$ appears in this formula for $`I`$. Hence, for the thermodynamic mass energy $`E`$ and the entropy $`S`$ defined by
$$E=\frac{I}{\beta },S=\beta MI,$$
(27)
we arrive at the well-known results
$$E=m,S=\pi r_H^2=\pi (2r_A^2m)^{2/3}.$$
(28)
Thermodynamic properties of distorted AdS black holes become almost the same as SADS black holes in the small $`\delta `$ domain.
## IV CONNECTION TO BLACK STRINGS
In the previous section we have seen that the thermodynamic evolution of black holes to a distorted configuration can occur without changing the mass energy and the entropy. Here we present an explicit example of a family of solutions describing a distortion from a SAdS black hole to a uniform black string which has the cylindrical metric Lemos1
$$ds^2=(\frac{r^2}{r_A^2}\frac{2m}{r})dt^2+(\frac{r^2}{r_A^2}\frac{2m}{r})^1dr^2+r^2(\frac{dz^2}{z_0^2}+d\varphi ^2),$$
(29)
where $`r`$ is interpreted as a cylindrical radius, and $`z_0`$ is an arbitrary parameter with the dimension of length. The parameter $`m`$ differs from the black hole mass which will become divergent for this cylindrical system with the infinite horizon area. Hence, to obtain such a black string solution, we must consider the increase of the SAdS mass in addition to the cylindrical distortion.
Note that the SAdS metric rewritten into the form
$$ds^2=(a^2\frac{2m}{r}+\frac{r^2}{r_A^2})dt^2+(a^2\frac{2m}{r}+\frac{r^2}{r_A^2})^1dr^2+\frac{r^2}{a^2}(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2),$$
(30)
represents a spherical black hole with mass $`m/a^3`$. If the distortion procedure explained in Sec. II is applied to this spherical metric (30), the distorted metric is given by
$$e^{2\nu }e^{2\mu }\frac{1}{\delta }(x^2\frac{1}{x})\{1+2\delta \nu _1+a^2\delta (x^2\frac{1}{x})^1\},$$
(31)
and
$$e^{2\psi }\frac{e^{2f(\theta )}}{a^2},$$
(32)
instead of Eqs. (10) and (11). Then it is easily found that
$$\frac{6h}{a^2}1=e^{2f}(\frac{d^2f}{d\theta ^2}+\mathrm{cot}\theta \frac{df}{d\theta }1)$$
(33)
as a modified version of Eq. (12). Though the condition (13) is assumed to be preserved, the area of $`S^2`$ (and the mass $`m/a^3`$) can increase as the positive parameter $`a`$ decreases.
The interesting example of the distortion function $`f`$ is
$$e^{2f}=\frac{b}{\mathrm{sin}^2\theta +\gamma },$$
(34)
where $`\gamma `$ is an arbitrary positive parameter, and the condition (13) leads to
$$b=2\sqrt{1+\gamma }\times (\mathrm{ln}\frac{\sqrt{1+\gamma }+1}{\sqrt{1+\gamma }1})^1.$$
(35)
Further, from Eq. (33) we obtain
$$\frac{6h}{a^2}=1+\frac{\gamma }{b}[1\frac{2(1+\gamma )}{\mathrm{sin}^2\theta +\gamma }].$$
(36)
We find that in the large $`\gamma `$ domain (where $`a`$ is assumed to be $`a1`$) the function $`f`$ represents a small quadrupole distortion given by
$$f\frac{1}{3\gamma }P_2(\mathrm{cos}\theta ),$$
(37)
using the approximated relation $`b\gamma +(2/3)`$. On the other hand, in the small $`\gamma `$ domain, we obtain
$$b(\mathrm{ln}\frac{2}{\sqrt{\gamma }})^11.$$
(38)
We must keep the metric function $`e^{2\psi }`$ finite even in the limit $`\gamma 0`$. Hence, from Eq. (32) the parameter $`a`$ is chosen to be $`a^2b0`$ in the small $`\gamma `$ limit, for which it is easy to see that the metric functions given by Eqs. (31) and (32) reduce to the cylindrical metric given by Eq. (29), using the coordinate transformation
$$\frac{z}{z_0}=_\theta ^{\pi /2}\frac{d\theta }{\mathrm{sin}\theta }.$$
(39)
In fact, the limit $`a^20`$ for the metric (30) means a transition to the horizon with zero scalar curvature $`k=0`$. Further, according to Eq. (36) the function $`h`$ giving the perturbed metric $`\nu _1`$ also vanishes in the same limit. By virtue of the disappearance of $`h`$ we obtain a non-perturbed uniform black string solution. Thus, we can claim that the distorted black hole metric written by Eq. (34) is a one-parameter family of solutions giving a SAdS black hole in the limit $`\gamma \mathrm{}`$ (i.e., $`a1`$) and a uniform black string in the limit $`\gamma 0`$ (i.e., $`a0`$).
Finally, let us discuss the validity of the Penrose inequality Bray and the hoop conjecture Thorne for black hole geometry in the distortion process to the black string. The Penrose inequality states that the black hole mass given by $`M=m/a^3`$ is not less than $`\sqrt{A/16\pi }`$, where $`A`$ is the horizon area given by $`A=4\pi (2mr_A^2/a^3)^{2/3}`$. The parameter $`a`$ is required to decrease from unity to zero as $`\gamma `$ decreases to zero. Then the validity of the Penrose inequality is apparent, because the ratio given by
$$\frac{M}{\sqrt{A/16\pi }}=\frac{1}{\delta a^2}$$
(40)
cannot be smaller than unity in the range $`a1`$.
On the other hand the hoop conjecture states that any circumference $`C`$ on the horizon is bounded by $`C4\pi M`$. Here, the circumference $`C`$ should be given as the length of a closed loop at a constant $`\varphi `$, namely, we obtain
$$C=\frac{(2mr_A^2)^{1/3}}{a}_0^\pi e^f𝑑\theta .$$
(41)
Then, in the limit $`\gamma 0`$, we can roughly evaluate the ratio $`C/4\pi M`$ as
$$\frac{C}{4\pi M}\delta (\mathrm{ln}\frac{1}{\gamma })^{1/2}1,$$
(42)
using $`a^2b1/\mathrm{ln}(1/\gamma )`$. We find that the hoop conjecture remains valid even for the extremely prolonged distortion $`\gamma 0`$.
It should be noted that the hoop conjecture becomes consistent with the existence of such a prolonged horizon owing to the rapid increase of mass $`M=m/a^3`$ as $`a^2`$ decreases in proportion to $`b`$. Hence, one may expect the breakdown of the hoop conjecture for a prolonged distortion (i.e., $`\gamma 1`$) keeping $`M=m`$ (i.e., $`a=1`$), which gives the ratio
$$\frac{C}{4\pi M}\delta (\mathrm{ln}\frac{1}{\gamma })^{1/2}.$$
(43)
However, the circumference $`C`$ can become larger than $`4\pi M`$ only for $`\mathrm{ln}(1/\gamma )`$ larger than $`1/\delta ^2`$. This is the case such that the perturbation $`\delta \nu _1\delta h`$ which is assumed to be very small in our analysis has the amplitude of the order of $`\delta b1/\delta 1`$. The non-perturbative change of the metric functions $`\nu `$ and $`\mu `$ may make the horizon structure disappear, consistently with the hoop conjecture. Unfortunately, it is impossible to find solutions describing the disappearance of the event horizon within the framework of our analysis, in which the distortion of the metric functions $`\nu `$ and $`\mu `$ is treated perturbatively, though the spherical surface $`S^2`$ is non-perturbatively distorted by the function $`\psi =f(\theta )`$. To construct an extremal state of distorted ”static” black holes is an interesting problem to be studied in future works.
In summary, we have succeeded in providing a simple non-perturbative scheme to obtain strongly distorted AdS black holes and showing explicitly the existence of a family of solutions connecting spherically symmetric black holes to uniform black strings. From the leading-order calculations in the large mass domain (i.e., $`mr_A`$), we have found that the black hole mass is identical with the SAdS mass parameter $`m`$, and the usual thermodynamic relation between the mass and the entropy holds, irrespective of the large distortion of the metric on $`S^2`$.
We would like to emphasize that the above-mentioned results have been obtained under the large mass approximation $`\delta r_H/2m1`$. If the parameter $`\delta `$ is not so small, the amplitude of the first-order metric perturbation $`\psi _1`$ distorting geometry on $`S^2`$ has the same order as $`\nu _1`$ and $`\mu _1`$, as was discussed in Yos . Namely, no large distortion of $`S^2`$ is allowed, unless the metric functions $`\nu _0`$ and $`\mu _0`$ of SAdS black holes are also strongly disturbed. In particular, in the small mass domain corresponding to $`mr_A`$ (i.e., $`\delta 1`$), the second-order perturbations $`\nu _2`$ and $`\mu _2`$ can become significantly large even if the first-order perturbations $`\nu _1`$ and $`\mu _1`$ has the small amplitude of the order of $`(m/r_A)^{1/2}`$. In relation to the uniqueness theorem (as well as the hoop conjecture) for black hole solutions in four-dimensional $`\mathrm{\Lambda }=0`$ spacetimes, it is a remaining important task to check non-perturbatively whether the horizon distorted from spherical symmetry can exist for the small mass domain.
###### Acknowledgements.
The author would like to thank Dr. H. Yoshino for helpful discussions. |
warning/0506/cond-mat0506297.html | ar5iv | text | # Fluctuations in the coarsening dynamics of the O(N) model: are they similar to those in glassy systems?
## 1 Introduction
Many extended systems which consist in interacting microscopic degrees of freedom exhibit non-trivial slow dynamics at low temperatures. Macroscopic observables such as density-density or other relevant correlations have extremely slow relaxations. Magnetic, dielectric or other susceptibilities slowly evolve in time. A large amount of experimental and numerical data allow for a qualitative, and sometimes also quantitative, description of these macroscopic observables in a number of well-studied materials. A satisfying understanding of the mechanism leading to such dramatic slowing down is, however, still lacking. In order to get a better insight on the relaxation of glassy systems it is important to investigate the dynamics at length/times scales that range from the microscopic to the macroscopic, through proper experimental -, numerical -, and theoretical tools -.
Systems with a clear mechanism for slow relaxations are the ones that evolve through coarsening of domains. They may thus provide a useful guideline to understand the dynamics of, in principle, more complicated systems. After a transient, systems undergoing phase-ordering kinetics enter a scaling regime in which the order-parameter morphology and its correlation functions depend on time only through a time-dependent length $`L(t)`$, that characterises the mean size of the domains . Interestingly enough, all microscopic details are absorbed in $`L(t)`$. It is tempting to speculate that such space-time scaling also exists asymptotically in glassy systems. This is the starting point, for example, in the dynamic droplet theory of spin-glasses (see for a detailed numerical examination).
Independently, analytical studies of dynamical mean-field theories of glassy systems demonstrated that the relaxation of global two-time correlation functions follows a self-similar structure, with a long-times scaling given by a ratio between a function of time evaluated at the two times involved, $`C(t,t^{})f_C[h(t^{})/h(t)]`$ . In these models, there is no interpretation of the function $`h(t)`$ as a length-scale. Even more generally, one can argue that any monotonic two-time correlation, independently of the origin of the slow dynamics, should depend on times only through a ratio $`h(t^{})/h(t)`$ within a given correlation scale .
It has been noticed by several authors - that the dynamic equations for the slow decay of the global correlations and responses of mean-field disordered models with glassy features acquire time-reparametrisation invariance once the time-derivatives (and other irrelevant terms) are dropped in the long times limit, in which the scaling in $`h(t^{})/h(t)`$ actually holds. This symmetry is not exactly realised since one function $`h(t)`$ is selected by the dynamic evolution; in other words, the time-derivative and other irrelevant terms, act as (asymptotically vanishing) pinning fields that select the time-scaling $`h(t)`$. The development, at long times, of an approximate invariance under generic reparameterisation of time has hindered the complete solution of the dynamic problem, for fixing the choice of reparametrisation involves a proper matching of the short-time and long-time dynamics that should be done by taking into account the effect of the time-derivative – and other terms.
More recently it has been suggested that a global time-reparametrisation invariance may also exist in finite dimensional glassy systems and that it may be responsible for the main spatio-temporal fluctuations -. In this way, the inconvenience generated by the time-reparametrisation invariance was transformed into a tool with predictive power. Some consequences of this proposal were listed in these articles together with their numerical checks in finite dimensional spin-glasses and kinetically facilitated models . Interestingly enough, a kind of ‘universality’ emerged in the sense that the time evolution and form of the distributions of local correlations and responses followed a similar pattern for these rather different systems.
The global time-reparametrisation $`th(t)`$ we are referring to acts on all spatial positions in identical way and it does not involve transforming space simultaneously. It is then simpler in form than the usual space-time rescaling that holds in coarsening systems at long times and large scales. In several stages of this paper we compare the time reparametrization invariance to the usual space-time rescaling. We also stress that time-reparametrisation invariance is a different transformation from Henkel’s local scale invariance hypothesis (see for a discussion on the validity of the latter).
The aim of this paper is to investigate the similarities and differences between fluctuations in simple coarsening and glassy systems. Specifically, we study analytically the coarsening dynamics of the $`d`$ dimensional $`O(N)`$ model in the large $`N`$ limit. This model has been studied in a large number of papers, see e.g. - and references therein. In Sect. 2 we review its static and dynamic behaviour. We explain in special detail the separation of the field in two components, as presented by Corberi, Lippiello and Zannetti , and how this helps understanding the condensation phenomenon and thermal fluctuations. Next, we analyse the fluctuating dynamics. In Sect. 3 we derive the dynamic generating functional and write it in terms of the slow and fast fields introduced in Sect. 2. We also derive closed dynamic equations for the global correlation and linear response of any $`O(N)`$ model in the large $`N`$ limit or spherical model. Then, in Sect. 4 we examine the symmetries of the dynamic equations for the global correlation and response, and the dynamical generating functional, under global transformations of time. We compare with the time-reparametrisation invariance suggested for glassy systems. Section 5 is devoted to the study of the probability distributions of the fluctuations at various mesoscopic length/time scales through several dynamical observables. We confront the latter to the results obtained for disordered spin - and kinetically constrained models and with the usual space-time scaling invariance of pure ferromagnetic coarsening. In Sect. 6 we compute a four-point correlation function similar to the one that is usually used in the context of super-cooled liquids -, -, to extract a dynamic growing length. We study its behaviour as a function of the two times involved and discuss its relation to a response. Finally, in Sect. 7 we present our conclusions together with some speculations.
## 2 The $`O(N)`$ model
The $`d`$-dimensional $`O(N)`$ non-linear sigma model is a coarse-grained approximation to a lattice spin model with nearest-neighbour ferromagnetic interactions. Its Hamiltonian reads
$$H=_Vd^dx\left[\frac{1}{2}(\stackrel{}{\varphi }(\stackrel{}{x}))^2+\frac{g}{4N}(\varphi ^2(\stackrel{}{x}))^2+\frac{r}{2}\varphi ^2(\stackrel{}{x})\stackrel{}{h}(\stackrel{}{x},t)\stackrel{}{\varphi }(\stackrel{}{x})\right].$$
The spatial dependence is given by the continuous $`d`$-dimensional vector $`\stackrel{}{x}=(x_1,\mathrm{},x_d)`$ and $`V`$ is the volume of the system. The field $`\stackrel{}{\varphi }`$ is an $`N`$-dimensional vector, $`\stackrel{}{\varphi }=(\varphi _1,\mathrm{},\varphi _N)`$ with $`\mathrm{}<\varphi _\alpha <\mathrm{}`$. A subindex $`\alpha `$ labels its $`N`$ components, $`\alpha =1,\mathrm{},N`$. The interplay between the quadratic and quartic terms (with couplings $`r`$ and $`g>0`$, respectively) favours the $`\varphi ^2(\stackrel{}{x},t)_{\alpha =1}^N\varphi _\alpha ^2(\stackrel{}{x},t)=Nr/g`$ configurations for $`r<0`$. $`h_\alpha `$ is a magnetic field coupled linearly to the field. In the infinitesimal limit $`\stackrel{}{h}`$ serves to compute the linear response, see eq. (20). A soft Ising, XY or Heisenberg model correspond to $`N=1,2`$ and $`3`$, respectively. In principle, the large $`N`$ limit is the starting point for a systematic $`1/N`$ expansion, although this may be difficult to control .
In the absence of the magnetic field $`\stackrel{}{h}`$, the Hamiltonian $`H`$ is invariant under uniform rotations of $`\stackrel{}{\varphi }`$:
$$\varphi _\alpha (\stackrel{}{x})\stackrel{~}{\varphi }_\alpha (\stackrel{}{x})=_{\alpha \beta }\varphi _\beta (\stackrel{}{x}),\stackrel{}{x},$$
$`O(N)`$. The summation convention over repeated indeces is used here and in what follows.
Dynamics is attributed to the field via the Langevin equations of motion:
$`\gamma \dot{\varphi }_\alpha (\stackrel{}{x},t)`$ $`=`$ $`^2\varphi _\alpha (\stackrel{}{x},t)\left({\displaystyle \frac{g}{N}}\varphi ^2(\stackrel{}{x},t)+r\right)\varphi _\alpha (\stackrel{}{x},t)+h_\alpha (\stackrel{}{x},t)+\eta _\alpha (\stackrel{}{x},t).`$
Henceforth we measure time in units of the inverse of the friction coefficient $`\gamma `$. $`\eta _\alpha (\stackrel{}{x},t)`$ is a spatially uncorrelated Gaussian white noise with zero mean, $`\eta _\alpha (\stackrel{}{x},t)=0`$ for all $`\stackrel{}{x}`$ and $`t`$, and variance
$`\eta _\alpha (\stackrel{}{x},t)\eta _\beta (\stackrel{}{x}^{},t^{})`$ $`=`$ $`2k_BT\delta _{\alpha \beta }\delta ^d(\stackrel{}{x}\stackrel{}{x}^{})\delta (tt^{}),`$
where $`T`$ is the temperature of the bath and $`k_B`$ is the Boltzmann constant. It is convenient to regularise the spatial correlations of the noise including a finite short-distance cut-off
$`\eta _\alpha (\stackrel{}{x},t)\eta _\beta (\stackrel{}{x}^{},t^{})`$ $`=`$ $`2k_BT\delta _{\alpha \beta }{\displaystyle \frac{e^{\frac{1}{4}(\stackrel{}{x}\stackrel{}{x}^{})^2\mathrm{\Lambda }^2}}{(4\pi \mathrm{\Lambda }^2)^{d/2}}}\delta (tt^{}),`$
that introduces correlations over a typical length $`1/\mathrm{\Lambda }`$ simulating the lattice spacing and cures some short distance divergences. $`1/(2\mathrm{\Lambda }^2)`$ will define a microscopic time scale $`t_0`$ that regularises divergent equal-time correlations. Hereafter the angular brackets indicate an average over the thermal noise and we set $`k_B=1`$.
The stochastic evolution has to be supplemented with the initial condition $`\stackrel{}{\varphi }(\stackrel{}{x},0)`$. Since we are interested in phase-ordering dynamics, we typically choose initial conditions that are uncorrelated in the $`N`$ dimensional space, $`[\varphi _\alpha (\stackrel{}{x},0)\varphi _\beta (\stackrel{}{x},0)]_{ic}\delta _{\alpha \beta }`$, and in real space, and have a Gaussian distribution
$$P[\stackrel{}{\varphi }(\stackrel{}{x},0)]=(2\pi \mathrm{\Delta }^2)^{NV/2}e^{\frac{1}{2\mathrm{\Delta }^2}_\alpha {\scriptscriptstyle d^dx\varphi _\alpha ^2(\stackrel{}{x},0)}}.$$
(1)
Hereafter we use square brackets, $`[\mathrm{}]_{ic}`$, to represent an average over initial conditions.
In the large $`N`$ limit one expects that the sum over components in $`\varphi ^2(\stackrel{}{x},t)`$ averages away the $`\stackrel{}{x}`$ dependence. One then looks for a solution such that
$$z(\stackrel{}{x},t)\frac{g}{N}\varphi ^2(\stackrel{}{x},t)+rz(t)\frac{g}{N}[\varphi ^2(\stackrel{}{x},t)]_{ic}+r,$$
(2)
where the average in the last term is taken over thermal histories and initial conditions. The functional form of $`z(t)`$ has to be determined self-consistently. As we shall see below the time-dependence of $`z(t)`$ determines the scaling in time of most of the interesting dynamic quantities. Note that we are implicitly assuming that $`N\mathrm{}`$ in that we are not letting $`z`$ fluctuate. All results in this paper have been derived in this limit. As discussed by Newman and Bray , fluctuations of $`z(t)`$ appear at order $`1/N`$.
Under the assumption (2), that has to be verified a posteriori, one can Fourier transform the Langevin equation and the noise-noise correlation. We use the following conventions:
$`f(\stackrel{}{k})={\displaystyle d^dxe^{i\stackrel{}{k}\stackrel{}{x}}f(\stackrel{}{x})},f(\stackrel{}{x})={\displaystyle \frac{d^dk}{(2\pi )^d}e^{i\stackrel{}{k}\stackrel{}{x}}f(\stackrel{}{k})},`$
and we obtain
$`\dot{\varphi }_\alpha (\stackrel{}{k},t)=k^2\varphi _\alpha (\stackrel{}{k},t)z(t)\varphi _\alpha (\stackrel{}{k},t)+\eta _\alpha (\stackrel{}{k},t),`$ (3)
$`\eta _\alpha (\stackrel{}{k},t)\eta _\beta (\stackrel{}{k}^{},t^{})=2T\delta _{\alpha \beta }e^{\frac{k^2}{\mathrm{\Lambda }^2}}(2\pi )^d\delta ^d(\stackrel{}{k}+\stackrel{}{k}^{})\delta (tt^{}).`$
In terms of the Fourier components $`\stackrel{}{\varphi }(\stackrel{}{k},0)`$, the initial conditions are distributed according to
$$P[\stackrel{}{\varphi }(\stackrel{}{k},0)]=(2\pi \mathrm{\Delta }^2)^{NV/2}e^{\frac{1}{2\mathrm{\Delta }^2}{\scriptscriptstyle {\scriptscriptstyle \frac{d^dk}{(2\pi )^d}}\stackrel{}{\varphi }(\stackrel{}{k},0)\stackrel{}{\varphi }(\stackrel{}{k},0)}}.$$
(4)
Thus, the coupled dynamics in $`x`$ space transforms into a set of $`N`$ independent first-order differential equations for the $`k`$-components of the field. The label $`\alpha `$ is now superfluous and we omit it unless otherwise stated.
The $`O(N)`$ model is intimately related to the spherical ferromagnet on a lattice and the fully-connected spherical spin-glass with two-body interactions. The main difference between these models is the form of the density of states of the quadratic interaction matrix and how it decays to zero at its edge. In the case of the $`O(N)`$ model the density of states is given by
$$g(ϵ)ϵ^\nu \nu =d/21.$$
(5)
at low energies $`ϵ`$. Many papers have been devoted to the study of the relaxation dynamics and global properties of the $`O(N)`$ model -, the spherical ferromagnet -, and the fully-connected spin-glass with two-body interactions , -. In the rest of this section we recall the main features of the statics and dynamics of the $`O(N)`$ model while in the rest of the paper we focus on the study of fluctuations and of symmetries under time transformations.
### 2.1 Statics
Let us briefly review the static behaviour of the $`O(N)`$ model (see Refs. for more details). If the volume $`V`$ is finite, the system equilibrates in finite time and the probability distribution function (pdf) of the order parameter approaches the Gibbs-Boltzmann form
$`P_{eq}(\stackrel{}{\varphi })=Z^1e^{\frac{\beta }{2V}_\stackrel{}{k}(k^2+\xi ^2)\stackrel{}{\varphi }(\stackrel{}{k})\stackrel{}{\varphi }(\stackrel{}{k})},`$ (6)
$`Z{\displaystyle 𝒟\stackrel{}{\varphi }e^{\frac{\beta }{2V}_\stackrel{}{k}(k^2+\xi ^2)\stackrel{}{\varphi }(\stackrel{}{k})\stackrel{}{\varphi }(\stackrel{}{k})}},`$ (7)
meaning that the Fourier components are independent Gaussian random variables. The path-integral measure is $`𝒟\stackrel{}{\varphi }_\alpha _kd\varphi _\alpha (\stackrel{}{k})`$. $`\xi `$ is the static correlation length
$$\xi ^2=\frac{g}{N}\varphi ^2(\stackrel{}{x})_{eq}+r$$
(8)
where the subindex ‘eq’ indicates that the average has to be computed using the measure (6)-(7). (We shall see below that $`z(t)`$ plays a similar role to $`\xi ^2`$.) All modes have vanishing thermal average, $`\stackrel{}{\varphi }(\stackrel{}{k})_{eq}=0`$ for all $`\stackrel{}{k}`$. The static structure factor
$`C_{\mathrm{eq}}(\stackrel{}{k}){\displaystyle \frac{1}{N}}\stackrel{}{\varphi }(\stackrel{}{k})\stackrel{}{\varphi }(\stackrel{}{k})_{eq}={\displaystyle \frac{TV}{k^2+\xi ^2}}`$ (9)
shows the ordering in the low temperature phase. The correlation length, $`\xi `$, is determined by eq. (8) with $`\varphi ^2(\stackrel{}{x})_{eq}=\varphi ^2(\stackrel{}{0})_{eq}`$ replaced by $`V^1`$ times the sum over $`\stackrel{}{k}`$ of (9). The detailed analysis of this equation has been presented elsewhere (see e.g. ). One finds that in $`2<d`$ there is a finite critical temperature $`T_c`$ defined by
$$r+gT_c\frac{d^dk}{(2\pi )^d}\frac{e^{\frac{k^2}{\mathrm{\Lambda }^2}}}{k^2}=0,$$
(10)
where the correlation length changes from a volume independent value at $`T>T_c`$ to a volume-dependent one at $`TT_c`$. In $`d=2`$ the integral over $`k`$ has a logarithmic divergence and the critical temperature is pushed down to zero. Above but near criticality $`\xi `$ behaves as
$`\xi `$ $`\left({\displaystyle \frac{TT_c}{T_c}}\right)^\nu ,\nu =\{\begin{array}{cc}1/2\hfill & d>4,\hfill \\ (d2)^1\hfill & d<4,\hfill \end{array}`$
with logarithmic corrections in $`d=4`$. At $`T_c`$, $`\xi V^\zeta `$ with $`\zeta =1/4`$ for $`d>4`$ and $`\zeta =d^1`$ for $`d<4`$, again with logarithmic corrections in $`d=4`$. Below $`T_c`$, the order parameter $`m_{eq}`$,
$$V^2m_{eq}^2N^1\varphi ^2(\stackrel{}{k}=\stackrel{}{0})_{eq}.$$
(12)
becomes non-zero and one finds,
$`m_{eq}^2`$ $`=`$ $`{\displaystyle \frac{r}{g}}{\displaystyle \frac{T_cT}{T_c}}\text{and}\xi ^2m_{eq}^2{\displaystyle \frac{V}{T}}.`$ (13)
The temperature and volume dependence of $`\xi `$ dictates that of the structure factor. When $`T>T_c`$ the variance of all modes grows linearly with the volume. Instead, when $`TT_c`$, $`\xi ^2`$ is negligible with respect to $`k^2`$ except at $`\stackrel{}{k}=0`$ yielding
$`C_{eq}(\stackrel{}{k})=\{\begin{array}{cc}VT_ck^2(1\delta _{\stackrel{}{k},\stackrel{}{0}})+\overline{c}T_cV^{2\zeta +1}\delta _{\stackrel{}{k},\stackrel{}{0}}\hfill & T=T_c\hfill \\ VTk^2(1\delta _{\stackrel{}{k},\stackrel{}{0}})+m_{\mathrm{eq}}^2V^2\delta _{\stackrel{}{k},\stackrel{}{0}}\hfill & T<T_c\hfill \end{array}`$
where $`\overline{c}`$ is a constant and $`m_{eq}`$ is given in eq. (13). The transition is characterised by a zero wave-vector mode that condenses and has a variance, $`\varphi ^2(\stackrel{}{k}=\stackrel{}{0})_{eq}`$, that grows as $`V^2`$. Below $`T_c`$ the equilibrium susceptibility $`\chi _{eq}`$ per unit volume is
$$\chi _{eq}=\frac{m_0^2m_{eq}^2}{T}=\frac{r}{g}T_c^1=(4\pi )^{d/2}\frac{2}{d2}\mathrm{\Lambda }^{d2}$$
(15)
where
$$m_0^2=r/g.$$
(16)
In conclusion, the $`O(N)`$ model has a phase transition from a paramagnetic to a ferromagnetic phase. The low temperature phase is characterised by a condensation phenomenon that signals ordering. The upper critical dimension is $`d=4`$ and the lower critical dimension is $`d=2`$.
#### 2.1.1 Separation of the field
The nature of the phase transition and low-temperature phase can be well understood by splitting the real-space order parameter in a constant contribution and a space-varying one :
$`\stackrel{}{\varphi }(\stackrel{}{x})=\stackrel{}{\sigma }+\stackrel{}{\psi }(\stackrel{}{x})\text{with}`$
$`\stackrel{}{\sigma }V^1\stackrel{}{\varphi }(\stackrel{}{k}=0)\text{and}\stackrel{}{\psi }(\stackrel{}{x})V^1{\displaystyle \underset{\stackrel{}{k}0}{}}\stackrel{}{\varphi }(\stackrel{}{k})e^{i\stackrel{}{k}\stackrel{}{x}}.`$
It is clear that the fields $`\stackrel{}{\sigma }`$ and $`\stackrel{}{\psi }`$ are independent. The Gibbs-Boltzmann measure factorises:
$`P[\stackrel{}{\varphi }(\stackrel{}{x})]=P(\stackrel{}{\sigma })P[\stackrel{}{\psi }(\stackrel{}{x})]`$
$`P(\stackrel{}{\sigma })=(2\pi m_{\mathrm{eq}}^2)^{N/2}e^{\sigma ^2/(2m_{\mathrm{eq}}^2)},P[\stackrel{}{\psi }(\stackrel{}{x})]=Z_\psi ^1e^{\beta /2_Vd^dx[\stackrel{}{}\stackrel{}{\psi }(\stackrel{}{x})]^2},`$
with $`Z_\psi =𝒟\stackrel{}{\psi }e^{\beta /2_Vd^dx[\stackrel{}{}\stackrel{}{\psi }(\stackrel{}{x})]^2}`$. The first factor describes the condensate with macroscopic variance $`\sigma _\alpha ^2_{\mathrm{eq}}=m_{\mathrm{eq}}^2`$. The second factor describes thermal fluctuations about the condensate. Consequently, the static correlation function separates in two terms:
$$C_{\mathrm{eq}}(\stackrel{}{r})N^1\stackrel{}{\varphi }(\stackrel{}{x})\stackrel{}{\varphi }(\stackrel{}{x}+\stackrel{}{r})_{eq}=m_{\mathrm{eq}}^2+N^1\stackrel{}{\psi }(\stackrel{}{x})\stackrel{}{\psi }(\stackrel{}{x}+\stackrel{}{r})_{eq}$$
(17)
where the first term represents the macroscopic variance of the condensate and the second one is the correlation of thermal fluctuations.
### 2.2 Dynamics
The set of linear differential equation (3) can be easily solved:
$`\varphi (\stackrel{}{k},t)`$ $`=`$ $`e^{k^2t_0^t𝑑t^{}z(t^{})}\varphi (\stackrel{}{k},0)`$ (18)
$`+{\displaystyle _0^t}𝑑t^{}e^{k^2(tt^{})_t^{}^t𝑑t^{}z(t^{})}[\eta (\stackrel{}{k},t^{})+h(\stackrel{}{k},t^{})],`$
where we dropped the component index $`\alpha `$, since all components satisfy the same equations due to rotational symmetry.
The function $`z(t)`$ is self-consistently determined. Indeed,
$$Y^2(t)e^{\mathrm{\hspace{0.33em}2}_0^t𝑑t^{}z(t^{})}$$
satisfies the differential equation
$$\frac{dY^2(t)}{dt}=2\left(\frac{g}{N}[\varphi ^2(\stackrel{}{x},t)]_{ic}+r\right)Y^2(t)$$
(19)
that, using the equation of motion to represent $`[\varphi ^2(\stackrel{}{x},t)]_{ic}`$, transforms into a closed first-order differential equation for $`Y^2(t)`$ complemented by the initial condition $`Y(0)=1`$. One finds that $`Y^2`$ grows exponentially at high temperatures, as a power law at criticality, and it decays as a power law, $`Y^2(t)t^{d/2}`$, below the critical temperature (see for a careful study of the preasymptotic behaviour of $`Y(t)`$ at low temperature).
The solution (18) takes a specially appealing form when written in terms of
$`{\displaystyle \frac{\delta \varphi _\alpha (\stackrel{}{k},t)}{\delta h_\beta (\stackrel{}{k}^{},t^{})}}|_{h=0}=\delta _{\alpha \beta }\delta ^d(\stackrel{}{k}+\stackrel{}{k}^{}){\displaystyle \frac{Y(t^{})}{Y(t)}}e^{k^2(tt^{})}\theta (tt^{}),`$ (20)
which can be Fourier transformed to give
$`{\displaystyle \frac{\delta \varphi _\alpha (\stackrel{}{x},t)}{\delta h_\beta (\stackrel{}{x}^{},t^{})}}|_{h=0}={\displaystyle \frac{d^dk}{(2\pi )^d}\frac{d^dk^{}}{(2\pi )^d}e^{i\stackrel{}{k}\stackrel{}{x}}e^{i\stackrel{}{k}^{}\stackrel{}{x}^{}}(2\pi )^d\frac{\delta \varphi _\alpha (\stackrel{}{k},t)}{\delta h_\beta (\stackrel{}{k}^{},t^{})}}|_{h=0}`$
Note that these quantities depend on the noise realisation and the initial condition only through the value of $`Y(t)`$. They are also identical to the linear response function,
$`R_{\alpha \beta }(\stackrel{}{x},\stackrel{}{x}^{};t,t^{}){\displaystyle \frac{\delta \varphi _\alpha (\stackrel{}{x},t)}{\delta h_\beta (\stackrel{}{x}^{},t^{})}}|_{h=0}.`$ (21)
This property is special of (quasi) quadratic models. Calling now
$$r(k;t,t^{})\frac{Y(t^{})}{Y(t)}e^{k^2(tt^{})}$$
the solution (18) can be rewritten as
$$\varphi (\stackrel{}{k},t)=r(k;t,0)\varphi (\stackrel{}{k},0)+_0^t𝑑t^{}r(k;t,t^{})\left[\eta (\stackrel{}{k},t^{})+h(\stackrel{}{k},t^{})\right].$$
(22)
### 2.3 Evolution of the distribution of Fourier components
Let us consider the evolution of initial configurations distributed according to the Gaussian law (1) \[and (4)\] in the absence of the perturbing field $`\stackrel{}{h}`$. Expression (22) indicates that the field configuration at time $`t`$ is in linear relation with the initial condition and the thermal noise. Since these fields are independent and Gaussian distributed, $`\stackrel{}{\varphi }(\stackrel{}{k},t)`$ is also Gaussian distributed with zero mean and time-dependent variance
$`[\varphi (\stackrel{}{k},t)\varphi (\stackrel{}{k}^{},t)]_{ic}=r(k,t,0)r(k^{},t,0)[\varphi (\stackrel{}{k},0)\varphi (\stackrel{}{k}^{},0)]_{ic}`$
$`+{\displaystyle _0^t}𝑑t^{}{\displaystyle _0^t}𝑑t^{\prime \prime }r(k,t,t^{})r(k^{},t,t^{\prime \prime })\eta (\stackrel{}{k},t^{})\eta (\stackrel{}{k}^{},t^{\prime \prime })`$
$`=(2\pi )^d\delta ^d(\stackrel{}{k}+\stackrel{}{k}^{})\left[{\displaystyle \frac{\mathrm{\Delta }^2}{Y^2(t)}}e^{2k^2t}+2T{\displaystyle _0^t}𝑑t^{}{\displaystyle \frac{Y^2(t^{})}{Y^2(t)}}e^{2k^2(tt^{}+t_0)}\right],`$
where we set $`t_0(2\mathrm{\Lambda }^2)^1`$. Note that this result holds for each component in the $`N`$-dimensional space while $`[\varphi _\alpha (\stackrel{}{k},t)\varphi _\beta (\stackrel{}{k}^{},t^{})]_{ic}=0`$ for all $`t`$ and $`t^{}`$ if $`\alpha \beta `$.
### 2.4 Correlations and responses
Let us discuss the relaxation of the correlations
$$C_{\alpha \beta }(\stackrel{}{x},\stackrel{}{x^{}};t,t^{})[\varphi _\alpha (\stackrel{}{x},t)\varphi _\beta (\stackrel{}{x^{}},t^{})]_{ic},$$
and the linear response defined in eq. (21). Due to the decorrelation of the initial conditions and noise in the $`N`$ dimensional space, these quantities are proportional to $`\delta _{\alpha \beta }`$ and we henceforth omit the internal indeces assuming that we take $`\alpha =\beta `$ in all quantities studied.
#### 2.4.1 Asymptotic behaviour.
In $`2<d`$ one finds a dynamic phase transition at the static critical temperature, $`T_c`$ given in eq. (10), where the asymptotic behaviour of $`Y(t)`$ changes. At high temperature each mode and, hence, the global correlation decay exponentially
$`C(t,t^{})V^1{\displaystyle d^dxC(x,x;t,t^{})}C_{\mathrm{eq}}e^{(tt^{})/t_{\mathrm{eq}}}\text{with}`$
$`C_{\mathrm{eq}}\varphi ^2(\stackrel{}{x})_{eq}\text{and}t_{\mathrm{eq}}=2\xi ^2,`$
where $`\xi `$ is the static correlation length given in eq. (8). The linear response is related to the correlation by the fluctuation dissipation theorem, $`R(t,t^{})=T^1_t^{}C(t,t^{})\theta (tt^{})`$. At the transition one finds interrupted ageing. Since we shall not discuss the critical dynamics, we do not give a detailed description of the scaling laws at criticality. Below the transition, after a transient (i.e. $`t^{}\tau _t1`$),
$`R(\stackrel{}{k},\stackrel{}{k}^{};t,t^{})`$ $``$ $`(2\pi )^d\delta ^d(\stackrel{}{k}+\stackrel{}{k}^{})\left({\displaystyle \frac{t}{t^{}}}\right)^{d/4}e^{k^2(tt^{})}\theta (tt^{}),`$
$`C(\stackrel{}{k},\stackrel{}{k}^{};t,t^{})`$ $``$ $`(2\pi )^d\delta ^d(\stackrel{}{k}+\stackrel{}{k}^{})(tt^{})^{d/4}e^{k^2tk_{}^{}{}_{}{}^{2}t^{}}`$
$`\times \left[\mathrm{\Delta }^2+2T{\displaystyle _0^{\mathrm{min}(t,t^{})}}𝑑t^{\prime \prime }Y^2(t^{\prime \prime })e^{(k^2+k_{}^{}{}_{}{}^{2})(t^{\prime \prime }2t_0)}\right].`$
Note that the asymptotic linear response does not depend on temperature. The first term in the correlation represents the decay of the initial conditions while the second one has its origin in the thermal noise. Each Fourier component with $`k>0`$ decays exponentially in time (with power law corrections). The $`k=0`$ component behaves differently since the exponential factor disappears. The slow decay of the low wave-vector components generates the non-trivial dynamics of the global correlation and response.
From the expressions above one easily recovers the real-space behaviour of the response and correlation. Using the large wave-vector cut-off $`\mathrm{\Lambda }`$, \[see eq. (2)\] – that will be important in $`d4`$,
$`R(\stackrel{}{x},\stackrel{}{x^{}},t,t^{}){\displaystyle \frac{d^dk}{(2\pi )^d}\frac{d^dk^{}}{(2\pi )^d}R(\stackrel{}{k},\stackrel{}{k}^{};t,t^{})e^{i\stackrel{}{k}\stackrel{}{x}i\stackrel{}{k}^{}\stackrel{}{x}^{}}e^{(k^2+k_{}^{}{}_{}{}^{2})/\mathrm{\Lambda }^2}},`$ (23)
and similarly for the correlation, one finds
$`R(\stackrel{}{x},\stackrel{}{x}^{};t,t^{})\left({\displaystyle \frac{t}{t^{}}}\right)^{d/4}(tt^{}+t_0)^{d/2}e^{(\stackrel{}{x}\stackrel{}{x}^{})^2/[4(tt^{}+t_0)]},`$ (24)
$`C(\stackrel{}{x},\stackrel{}{x}^{};t,t^{})(tt^{})^{d/4}[{\displaystyle \frac{\mathrm{\Delta }^2}{(t+t^{})^{d/2}}}e^{(xx^{})^2/[4(t+t^{})]}`$
$`+2T{\displaystyle _0^{\mathrm{min}(t,t^{})}}dt^{\prime \prime }{\displaystyle \frac{Y^2(t^{\prime \prime })}{[t+t^{}+2(t_0t^{\prime \prime })]^{d/2}}}e^{(xx^{})^2/\{4[t+t^{}+2(t_0t^{\prime \prime })]\}}],`$ (25)
The local response and correlation on the same spatial point, $`\stackrel{}{x}=\stackrel{}{x}^{}`$, are independent of $`\stackrel{}{x}`$; thus they are also equal to the global values, $`R(\stackrel{}{x},\stackrel{}{x};t,t^{})=R(t,t^{})`$ and $`C(\stackrel{}{x},\stackrel{}{x};t,t^{})=C(t,t^{})`$ with
$`R(t,t^{})\left({\displaystyle \frac{t}{t^{}}}\right)^{d/4}(tt^{}+t_0)^{d/2},`$ (26)
$`C(t,t^{})\left({\displaystyle \frac{tt^{}}{(t+t^{})^2}}\right)^{d/4}\left[\mathrm{\Delta }^2+2T{\displaystyle _0^{\mathrm{min}(t,t^{})}}𝑑t^{\prime \prime }{\displaystyle \frac{Y^2(t^{\prime \prime })}{\left[1\frac{2t^{\prime \prime }}{(t+t^{})}\right]^{d/2}}}\right],`$ (27)
where we assumed that $`t+t^{}`$ is larger than $`t_0=1/(2\mathrm{\Lambda }^2)`$ and we neglected the dependence of the correlation on this time-scale. The contribution from the small wave vectors lead to a non-trivial dynamics of the global correlation and response, with no exponential decay, and a separation of time scales shown below.
The case $`d=2`$, the lower critical dimension, may seem slightly different . Interesting dynamics occurs only at zero temperature, i.e. at the critical point. However, the dynamics is not typically critical but it corresponds to the zero temperature limit of the coarsening phenomena observed in higher dimensions. More precisely, there is still an additive separation of time-scales, as opposed to what occurs in critical relaxations where the separation is multiplicative and the ageing contribution to the correlation progressively disappears as time elapses.
One can check that these results are valid for all initial conditions with short-range correlations. Initial configurations with long-range correlations (as for an ordered configuration) lead to different scaling forms .
#### 2.4.2 Separation of time-scales.
At low temperature, $`T<T_c`$, and for very long waiting-time, $`t^{}\tau _t`$, the global linear response and correlation have two distinct two-time regimes depending on the relation between the times $`t`$ and $`t^{}`$. These are defined by
$`tt^{}t^{}`$ $`\text{stationary regime},`$
$`\lambda {\displaystyle \frac{t^{}}{t}}[0,1)`$ $`\text{ageing regime},`$
In the limit $`t^{}\tau _t`$ the global (and local) correlation and linear response are well described by an additive separation
$`C(t,t^{})=C_{\mathrm{st}}(tt^{})+C_{\mathrm{ag}}(t,t^{}),`$ (28)
$`C_{\mathrm{st}}(tt^{})(m_0^2m_{eq}^2)(\mathrm{\Lambda }^2(tt^{})+1)^{1d/2},`$ (29)
$`C_{\mathrm{ag}}(t,t^{})m_{eq}^2\left[{\displaystyle \frac{4\lambda }{(1+\lambda )^2}}\right]^{d/4},`$ (30)
and
$`R(t,t^{})=R_{\mathrm{st}}(tt^{})+R_{\mathrm{ag}}(t,t^{}),`$ (31)
$`R_{\mathrm{st}}(tt^{})=(4\pi )^{d/2}(tt^{}+t_0)^{d/2}=T^1_t^{}C_{\mathrm{st}}(tt^{})`$ (32)
$`R_{ag}(t,t^{})(4\pi )^{d/2}t^{d/2}\left[\left({\displaystyle \frac{t^{}}{t}}\right)^{d/4}1\right]\left(1{\displaystyle \frac{t^{}}{t}}+{\displaystyle \frac{t_0}{t}}\right)^{d/2}.`$ (33)
$`m_{\mathrm{eq}}`$ is the equilibrium magnetisation given in (13) and $`m_0`$ is its value at $`T=0`$ given in (16). The relation (15) between $`m_{eq}^2m_0^2`$ and $`t_0`$ ensures the validity of the fluctuation-dissipation theorem in the stationary regime. For long-time differences, such that the ratio between the two times $`t`$ and $`t^{}`$ is held fixed, $`t^{}/t=\lambda [0,1)`$ the correlation and response “age”, i.e. they depend on the waiting-time $`t^{}`$.
The detailed scaling of the correlation of the field evaluated at different times and spatial points for the (simpler) Gaussian scalar model was presented in . Here we just recall that eqs. (24)-(25) can be put in a scaling form :
$`C(\stackrel{}{x},\stackrel{}{x}^{}=\stackrel{}{x}+\stackrel{}{r};t,t^{})f_{C_r}({\displaystyle \frac{r}{L(t^{})}},{\displaystyle \frac{L(t^{})}{L(t)}}),`$ (34)
$`R(\stackrel{}{x},\stackrel{}{x}^{}=\stackrel{}{x}+\stackrel{}{r};t,t^{})t^{d/2}f_{R_r}({\displaystyle \frac{r}{L(t^{})}},{\displaystyle \frac{L(t^{})}{L(t)}},{\displaystyle \frac{t_0}{t}}),`$
with $`r=|\stackrel{}{x}\stackrel{}{x}^{}|`$. $`L(t)`$ is the ‘domain length’ at time $`t`$, which in the relaxation $`O(N)`$ model with non-conserved order parameter is given by $`L(t)\sqrt{t}`$. In particular, the ageing contribution to the global correlation (30) scales as
$$C_{\mathrm{ag}}(t,t^{})=f_C\left(\frac{L(t^{})}{L(t)}\right)\text{with}f_C(x)=m_{eq}^2\left(\frac{2x}{1+x^2}\right)^{d/2},$$
(35)
$`f_C(x)x^{\frac{d}{2}}`$ when $`x0`$ and $`f_C(x)m_{eq}^2(1dϵ^2/4)`$ when $`x1ϵ`$ and $`ϵ1`$. Note that the regime of very separated times, $`tt^{}`$ or $`x0`$, is characterised by a power law decay with an exponent $`\overline{\lambda }=d/2`$ . The global linear response (33) also takes a scaling form :
$`R_{ag}(t,t^{})t^{d/2}f_R(x,y),\text{with}`$
$`f_R(x,y)=(4\pi )^{d/2}(1x)^{d/4}(1x+y)^{d/2},\text{and}`$
$`x=t^{}/t=L^2(t^{})/L^2(t),y=t_0/t.`$ (36)
In the coarsening regime the global correlation and response are not related by the fluctuation dissipation theorem. One defines the ratio
$$X(t,t^{})\frac{TR_{ag}(t,t^{})}{_t^{}C_{ag}(t,t^{})}t^{1d/2}f_X(\lambda ).$$
(37)
Note that this is a decreasing function of time that tends to zero in all $`d>2`$ and to a function of the times ratio, $`f_X(\lambda )`$, taking finite values when $`d2^+`$.
The dc susceptibility or zero field cooled magnetisation, defined as the integral of the linear response over a time period:
$$\chi (t,t^{})=_t^{}^t𝑑t^{\prime \prime }R(t,t^{\prime \prime })$$
(38)
can be expressed as a sum of two terms, a stationary and an ageing contribution,
$$\chi (t,t^{})\chi _{st}(t,t^{})+\chi _{ag}(t,t^{})=_t^{}^t𝑑t^{\prime \prime }R_{st}(tt^{\prime \prime })+_t^{}^t𝑑t^{\prime \prime }R_{ag}(t,t^{\prime \prime })$$
(39)
given by
$`\chi _{st}(tt^{})`$ $`=`$ $`\chi _{\mathrm{eq}}\{1[(tt^{})/t_0+1]^{1d/2}\},`$ (40)
$`\chi _{ag}(t,t^{})`$ $``$ $`\{\begin{array}{cc}t^{1d/2}\hfill & d<4,\hfill \\ t^{1d/2}\mathrm{ln}(t/t_0)\hfill & d=4,\hfill \\ t^1t_0^{1d/2}\hfill & d>4.\hfill \end{array}`$ (44)
where $`\chi _{\mathrm{eq}}=(4\pi )^{d/2}t_0^{1d/2}/(d/21)`$ is the equilbrium susceptibility given in eq. (15). There are several features to be noticed in these expressions. The first one is that the stationary integrated response approaches a value proportional to $`t_0^{1d/2}`$ in the long $`tt^{}`$ limit. If one takes the cut-off $`\mathrm{\Lambda }`$ to infinity this value diverges as a power law, $`t_0^{1d/2}`$ in all $`d>2`$. Instead, in $`d=2`$ $`\chi _{st}`$ diverges as a logarithm of the time difference, $`\chi _{st}(t,t^{})\mathrm{ln}[(tt^{})/t_0]`$, for $`tt^{}t_0`$. The approach to this asymptotic value is given by a power law, $`[(tt^{})/t_0]^{1d/2}/(1d/2)`$, that will play an important role in the analysis of the invariances of the slow dynamics. Above $`d=4`$ the decay of the ageing part of the integrated linear response does not depend on dimensionality any longer. As discussed by Corberi et al a similar upper dimension $`d_\chi `$ is expected to exist in other coarsening systems . In all $`d>2`$ the ageing contribution to the total susceptibility vanishes at long times, i.e. when $`t\mathrm{}`$.
$`d=2`$ is the lower critical dimension. But the dynamic behaviour at zero temperature can be reached as the zero temperature limit of the finite temperature coarsening dynamics just described . The additive separation of the correlation and response also holds in this case. In particular, the Edwards-Anderson parameter, $`q_{ea}=m_{eq}^2`$, that separates the stationary from the ageing regime in the correlation, remains finite (and equal to $`m_0^2`$ at zero temperature) in the limit of long waiting-time, $`t^{}\mathrm{}`$. Note that the stationary response in (32) does not depend on temperature and thus these soft ‘spins’ respond even at zero temperature. In particular, one has
$$R_{st}(tt^{})=\underset{T0}{lim}T^1\frac{d}{dt}C_{st}(tt^{}),\text{for}tt^{}>0.$$
(45)
#### 2.4.3 Separation of the field
Interestingly enough, Corberi, Lippiello and Zannetti showed that in $`d>2`$ the above results can also be found by using a splitting of the space and time dependent field $`\stackrel{}{\varphi }(\stackrel{}{x},t)`$ in two components :
$`\stackrel{}{\varphi }(\stackrel{}{x},t)`$ $`=`$ $`\stackrel{}{\sigma }(\stackrel{}{x},t)+\stackrel{}{\psi }(\stackrel{}{x},t).`$
Indeed, the solution (22) can be rewritten as
$`\stackrel{}{\varphi }(\stackrel{}{k},t)`$ $`=`$ $`\stackrel{}{\sigma }_h(\stackrel{}{k},t)+\stackrel{}{\psi }(\stackrel{}{k},t),`$
$`\stackrel{}{\sigma }_h(\stackrel{}{k},t)`$ $``$ $`r(k;t,t_1)\stackrel{}{\varphi }(\stackrel{}{k},t_1)+{\displaystyle _{t_w}^t}𝑑t^{}r(k;t,t^{})\stackrel{}{h}(\stackrel{}{k},t^{}),`$
$`\stackrel{}{\psi }(\stackrel{}{k},t)`$ $``$ $`{\displaystyle _{t_1}^t}𝑑t^{}r(k;t,t^{})\eta (\stackrel{}{k},t^{}).`$
$`t_1`$ is an arbitrary time satisfying $`t_1t^{}t`$ and sufficiently long so that the scaling limit has been established between the initial quench and $`t_1`$, i.e. $`t_1\tau _t`$. The second term in $`\stackrel{}{\sigma }`$ represents the effect of an external field applied from $`t_w`$ (another long time $`t_wt_1`$) on. Note that, in the absence of the field $`\stackrel{}{h}`$, the “slow” component $`\stackrel{}{\sigma }`$ can also be written as the evolution of the initial condition since
$`\stackrel{}{\sigma }(\stackrel{}{k},t)`$ $`r(k;t,t_1)\stackrel{}{\varphi }(\stackrel{}{k},t_1)`$
$`=r(k;t,t_1)r(k;t_1,0)\stackrel{}{\varphi }(\stackrel{}{k},0)+{\displaystyle _0^{t_1}}𝑑t^{}r(k;t,t_1)r(k;t_1,t^{})\stackrel{}{\eta }(\stackrel{}{k},t^{})`$
$`=r(k;t,0)\stackrel{}{\varphi }(\stackrel{}{k},0)+{\displaystyle _0^{t_1}}𝑑t^{}r(k;t,t^{})\stackrel{}{\eta }(\stackrel{}{k},t^{}).`$
Any typical initial condition that is the result of a quench from high temperatures can be thought of as a random noise of Gaussian type. Thus, the first term is statistically ‘identical’ to the contribution of the lower limit of the integral in the second term.
The field $`\stackrel{}{\sigma }`$ is associated to local condensation of the order parameter while $`\stackrel{}{\psi }`$ describes thermal fluctuations within the domains. These fields are statistically independent ($`\stackrel{}{\sigma }(\stackrel{}{x},t)\stackrel{}{\psi }(\stackrel{}{x}^{},t)=0`$) and have zero average. The explicit calculations in demonstrate that, in the long times limit, $`tt^{}t_1`$ with $`t_1`$ itself diverging, the global correlation of $`\stackrel{}{\sigma }`$, $`NC_\sigma (t,t^{},t_1)=\stackrel{}{\sigma }(t)\stackrel{}{\sigma }(t^{})`$, yields the ageing component of the global correlation of the field $`\stackrel{}{\varphi }`$, while the global correlation of $`\stackrel{}{\psi }`$, $`NC_\psi (t,t^{})=\stackrel{}{\psi }(t)\stackrel{}{\psi }(t^{})`$, yields the stationary components of the global correlation of the field $`\stackrel{}{\varphi }`$. More precisely, for $`t,t^{}t_1`$ one finds
$`C_\psi (t,t^{};t_1)`$ $`=`$ $`(m_0^2m_{eq}^2)[(4t_1/t_0)^{1d/2}\left({\displaystyle \frac{4tt^{}}{(t+t^{})^2}}\right)^{d/4}`$ (46)
$`+(2(tt^{})/t_0+1)^{1d/2}],`$
$`C_\sigma (t,t^{};t_1)`$ $`=`$ $`\left[m_{eq}^2(m_0^2m_{eq}^2)(4t_1/t_0)^{1d/2}\right]\left({\displaystyle \frac{4tt^{}}{(t+t^{})^2}}\right)^{d/4}.`$ (47)
In the limit $`t_1/t_01`$ and for $`d>2`$ the first term in (46) vanishes and $`C_\psi `$ describes the time-difference variation of the correlation in the stationary approach to the plateau at $`m_0^2m_{eq}^2`$. Similarly, $`C_\sigma (t,t^{},t_1)`$ becomes $`C_{\mathrm{ag}}(t,t^{})`$. Indeed, in the stationary regime, $`tt^{}t^{}`$, $`C_\psi `$ varies from $`m_0^2m_{eq}^2`$ to zero while $`C_\sigma `$ takes the constant value $`m_{eq}^2`$. In the ageing regime $`t^{}/t=\lambda `$, $`C_\psi `$ has already decayed to zero while $`C_\sigma `$ varies from $`m_{eq}^2`$ to zero.
The linear response is simply obtained as
$$\frac{\delta \varphi _\alpha (\stackrel{}{k},t)}{\delta h_\beta (\stackrel{}{k}^{},t^{})}|_{h=0}=\frac{\delta \sigma _{h}^{}{}_{\alpha }{}^{}(\stackrel{}{k},t)}{\delta h_\beta (\stackrel{}{k}^{},t^{})}|_{h=0}.$$
A similar separation has been used by Franz and Virasoro in a more general context .
## 3 The action
Let us now write the dynamic generating functional in terms of a path integral over the fields $`\stackrel{}{\sigma }`$ and $`\stackrel{}{\psi }`$. This will be useful to identify the symmetries of the “slow” action under transformations of time.
The dynamic generating functional is
$`Z={\displaystyle 𝒟\stackrel{}{\varphi }(\stackrel{}{k},t)𝒟i\stackrel{}{\widehat{\varphi }}(\stackrel{}{k},t)𝒟\stackrel{}{\eta }(\stackrel{}{k},t)e^{{\scriptscriptstyle {\scriptscriptstyle \frac{d^dk}{(2\pi )^d}}_0^{\mathrm{}}𝑑tS_{kt}^{\left(1\right)}}}},`$
$`S_{kt}^{(1)}i\stackrel{}{\widehat{\varphi }}(\stackrel{}{k},t)[(_t+k^2+z(t))\stackrel{}{\varphi }(\stackrel{}{k},t)\stackrel{}{\eta }(\stackrel{}{k},t)`$
$`\stackrel{}{h}(\stackrel{}{k},t)\theta (tt_w)](4T)^1\stackrel{}{\eta }(\stackrel{}{k},t)\stackrel{}{\eta }(\stackrel{}{k},t)`$ (48)
where, for simplicity, we took the cut-off $`\mathrm{\Lambda }`$ in the noise-noise correlation to infinity. The external field $`\stackrel{}{h}`$ is applied from $`t_w`$ onwards. We call $`R^1(k;t,t^{})`$ the differential operator $`\delta (tt^{})\left[_t^{}+k^2+z(t^{})\right]`$ whose inverse is the retarded linear response function, see eq. (21),
$$𝑑t^{}\delta (tt^{})\left(_t^{}+k^2+z(t^{})\right)R(k;t^{},t^{\prime \prime })=\delta (tt^{\prime \prime }),$$
$`R(k;t,t^{})=r(k;t,t^{})\theta (tt^{})`$, for each $`k`$. The action in eq. (48) can be rewritten as
$`{\displaystyle \frac{d^dk}{(2\pi )^d}_0^{\mathrm{}}𝑑tS_{kt}^{(1)}}={\displaystyle \frac{d^dk}{(2\pi )^d}_0^{\mathrm{}}𝑑t_0^{\mathrm{}}𝑑t^{}S_{ktt^{}}^{(2)}}`$ (49)
$`S_{ktt^{}}^{(2)}=i\stackrel{}{\widehat{\varphi }}(\stackrel{}{k},t)R^1(k;t,t^{})[\stackrel{}{\varphi }(\stackrel{}{k},t^{}){\displaystyle _0^{t_1}}dt^{\prime \prime }R(k;t^{},t^{\prime \prime })\stackrel{}{\eta }(\stackrel{}{k},t^{\prime \prime })`$
$`{\displaystyle _{t_1}^{\mathrm{}}}dt^{\prime \prime }R(k;t^{},t^{\prime \prime })\stackrel{}{\eta }(\stackrel{}{k},t^{\prime \prime }){\displaystyle _0^{\mathrm{}}}dt^{\prime \prime }R(k;t^{},t^{\prime \prime })\stackrel{}{h}(\stackrel{}{k},t^{\prime \prime })\theta (t^{\prime \prime }t_w)]`$
$`(4T)^1\stackrel{}{\eta }(\stackrel{}{k},t)\delta (tt^{})\stackrel{}{\eta }(\stackrel{}{k},t^{}).`$
Defining
$`\stackrel{}{\sigma }(\stackrel{}{k},t)`$ $``$ $`{\displaystyle _0^{t_1}}𝑑t^{\prime \prime }R(k;t,t^{\prime \prime })\stackrel{}{\eta }(\stackrel{}{k},t^{\prime \prime })+{\displaystyle _{t_w}^{\mathrm{}}}𝑑t^{\prime \prime }R(k;t,t^{\prime \prime })\stackrel{}{h}(\stackrel{}{k},t^{\prime \prime })`$
$`=`$ $`{\displaystyle _0^{t_1}}𝑑t^{\prime \prime }r(k;t,t^{\prime \prime })\stackrel{}{\eta }(\stackrel{}{k},t^{\prime \prime })+{\displaystyle _{t_w}^t}𝑑t^{\prime \prime }r(k;t,t^{\prime \prime })\stackrel{}{h}(\stackrel{}{k},t^{\prime \prime })\theta (tt_w),`$
$`\stackrel{}{\psi }(\stackrel{}{k},t)`$ $``$ $`{\displaystyle _{t_1}^{\mathrm{}}}𝑑t^{\prime \prime }R(k;t,t^{\prime \prime })\stackrel{}{\eta }(\stackrel{}{k},t^{\prime \prime })={\displaystyle _{t_1}^t}𝑑t^{\prime \prime }r(k;t,t^{\prime \prime })\stackrel{}{\eta }(\stackrel{}{k},t^{\prime \prime }),`$
along the lines of what has been reviewed in Sect. 2.4.3, and introducing these definitions with delta functions in the generating functional one has
$`Z={\displaystyle 𝒟\stackrel{}{\varphi }𝒟i\stackrel{}{\widehat{\varphi }}𝒟\stackrel{}{\eta }𝒟\stackrel{}{\sigma }𝒟\stackrel{}{\widehat{\sigma }}𝒟\stackrel{}{\psi }𝒟\stackrel{}{\widehat{\psi }}e^{{\scriptscriptstyle {\scriptscriptstyle \frac{d^dk}{(2\pi )^d}}_0^{\mathrm{}}𝑑t_0^{\mathrm{}}𝑑t^{}S_{ktt^{}}^{\left(3\right)}}}}`$
$`S_{ktt^{}}^{(3)}=i\stackrel{}{\widehat{\varphi }}(\stackrel{}{k},t)R^1(k;t,t^{})\left[\stackrel{}{\varphi }(\stackrel{}{k},t^{})\stackrel{}{\sigma }(\stackrel{}{k},t^{})\stackrel{}{\psi }(\stackrel{}{k},t^{})\right]`$
$`+i\stackrel{}{\widehat{\sigma }}(\stackrel{}{k},t)\delta (tt^{})[\stackrel{}{\sigma }(\stackrel{}{k},t^{}){\displaystyle _0^{t_1}}dt^{\prime \prime }r(k;t^{},t^{\prime \prime })\stackrel{}{\eta }(\stackrel{}{k},t^{\prime \prime })+`$
$`{\displaystyle _{t_w}^{\mathrm{}}}dt^{\prime \prime }R(k;t^{},t^{\prime \prime })\stackrel{}{h}(\stackrel{}{k},t^{\prime \prime })]`$
$`+i\stackrel{}{\widehat{\psi }}(\stackrel{}{k},t)\delta (tt^{})\left[\stackrel{}{\psi }(\stackrel{}{k},t^{}){\displaystyle _{t_1}^{\mathrm{}}}𝑑t^{\prime \prime }R(k;t^{},t^{\prime \prime })\stackrel{}{\eta }(\stackrel{}{k},t^{\prime \prime })\right]`$
$`(4T)^1\stackrel{}{\eta }(\stackrel{}{k},t)\delta (tt^{})\stackrel{}{\eta }(\stackrel{}{k},t^{}).`$
In the end we shall focus on the behaviour of the two-time action for long times; more explicitly we shall take $`t`$ and $`t^{}`$ to be longer than a long but otherwise arbitrary time $`t_1`$ but, for the moment, $`t_1`$ is just an arbitrary time scale.
The integration over the $`\stackrel{}{\widehat{\varphi }}`$ field yields a functional delta-function. Integrating next over the thermal noise (note that the noises multiplying $`\stackrel{}{\sigma }`$ and $`\stackrel{}{\psi }`$ are independent since they are evaluated at different times), one finds
$`Z={\displaystyle 𝒟\stackrel{}{\varphi }𝒟\stackrel{}{\sigma }𝒟\stackrel{}{\widehat{\sigma }}𝒟\stackrel{}{\psi }𝒟\stackrel{}{\widehat{\psi }}e^{{\scriptscriptstyle {\scriptscriptstyle \frac{d^dk}{(2\pi )^d}}_0^{\mathrm{}}𝑑t_0^{\mathrm{}}𝑑t^{}S_{ktt^{}}^{\left(4\right)}}}}`$
$`\times \delta \left[{\displaystyle _0^t}𝑑t^{}R^1(k;t,t^{})\left(\stackrel{}{\varphi }(\stackrel{}{k},t^{})\stackrel{}{\sigma }(\stackrel{}{k},t^{})\stackrel{}{\psi }(\stackrel{}{k},t^{})\right)\right]`$
$`S_{ktt^{}}^{(4)}=\delta (tt^{})[i\stackrel{}{\widehat{\sigma }}(\stackrel{}{k},t)\stackrel{}{\sigma }(\stackrel{}{k},t^{})+i\stackrel{}{\widehat{\psi }}(\stackrel{}{k},t)\stackrel{}{\psi }(\stackrel{}{k},t^{})`$
$`i\stackrel{}{\widehat{\sigma }}(\stackrel{}{k},t){\displaystyle _{t_w}^{\mathrm{}}}dt^{\prime \prime }R(k;t^{},t^{\prime \prime })\stackrel{}{h}(\stackrel{}{k},t^{\prime \prime })]`$
$`+T{\displaystyle _0^{t_1}}𝑑t^{\prime \prime }i\stackrel{}{\widehat{\sigma }}(\stackrel{}{k},t)r(k;t,t^{\prime \prime })i\stackrel{}{\widehat{\sigma }}(\stackrel{}{k},t^{})r(k;t^{},t^{\prime \prime })`$
$`+T{\displaystyle _{t_1}^{\mathrm{}}}𝑑t^{\prime \prime }i\stackrel{}{\widehat{\psi }}(\stackrel{}{k},t)R(k;t,t^{\prime \prime })i\stackrel{}{\widehat{\psi }}(\stackrel{}{k},t^{})R(k;t^{},t^{\prime \prime }).`$
From this expression one can relate the linear response of the field $`\stackrel{}{\sigma }`$ to the correlation between the fields $`\stackrel{}{\widehat{\sigma }}`$ and $`\stackrel{}{\sigma }`$:
$`{\displaystyle \frac{\delta \sigma _\alpha (\stackrel{}{k},t)}{\delta h_\beta (\stackrel{}{k}^{},t^{})}}|_{h=0}={\displaystyle _{t_w}^{\mathrm{}}}𝑑t^{\prime \prime }R(k^{},t^{\prime \prime },t^{})\sigma _\alpha (\stackrel{}{k},t)i\widehat{\sigma }_\beta (\stackrel{}{k}^{},t^{\prime \prime })\theta (t^{}t_w).`$
The average has to be computed with the action $`S_{ktt^{}}^{(4)}`$ evaluated at zero field ($`\stackrel{}{h}=\stackrel{}{0}`$). Since this action is diagonal in $`\stackrel{}{k}`$ and time the cross-correlation between the fields is just $`i\widehat{\sigma }_\beta (\stackrel{}{k}^{},t^{\prime \prime })\sigma _\alpha (\stackrel{}{k},t)=\delta ^d(\stackrel{}{k}+\stackrel{}{k}^{})\delta _{\alpha \beta }\delta (tt^{\prime \prime })`$. Thus,
$`{\displaystyle \frac{\delta \sigma _\alpha (\stackrel{}{k},t)}{\delta h_\beta (\stackrel{}{k}^{},t^{})}}|_{h=0}=r(k,t,t^{})\theta (tt^{})\delta ^d(\stackrel{}{k}+\stackrel{}{k}^{})\theta (t^{}t_w)`$ (50)
as it should.
The action $`S_{ktt^{}}^{(4)}`$ is quadratic in the fields $`\stackrel{}{\widehat{\sigma }}`$ and $`\stackrel{}{\widehat{\psi }}`$. Integrating them out one finds
$`Z{\displaystyle 𝒟\stackrel{}{\varphi }𝒟\stackrel{}{\sigma }𝒟\stackrel{}{\psi }e^{{\scriptscriptstyle {\scriptscriptstyle \frac{d^dk}{(2\pi )^d}}_0^{\mathrm{}}𝑑t_0^{\mathrm{}}𝑑t^{}S_{ktt^{}}^{\left(5\right)}}}}`$
$`\times \delta \left[{\displaystyle _0^t}𝑑t^{}R^1(k;t,t^{})\left(\stackrel{}{\varphi }(\stackrel{}{k},t^{})\stackrel{}{\sigma }(\stackrel{}{k},t^{})\stackrel{}{\psi }(\stackrel{}{k},t^{})\right)\right],`$
$`S_{ktt^{}}^{(5)}=(4T)^1\stackrel{}{\sigma }(\stackrel{}{k},t)\left[{\displaystyle _0^{t_1}}𝑑t^{\prime \prime }r(k;t,t^{\prime \prime })r(k;t^{},t^{\prime \prime })\right]^1\stackrel{}{\sigma }(\stackrel{}{k},t^{})`$
$`(4T)^1\stackrel{}{\psi }(\stackrel{}{k},t)\left[{\displaystyle _{t_1}^{\mathrm{}}}𝑑t^{\prime \prime }R(k;t,t^{\prime \prime })R(k;t^{},t^{\prime \prime })\right]^1\stackrel{}{\psi }(\stackrel{}{k},t^{}).`$ (51)
(The proportionality sign is due to the fact that the integration over the fields $`\stackrel{}{\widehat{\sigma }}`$ and $`\stackrel{}{\widehat{\psi }}`$ also yields a determinant that is just a ‘numerical constant’ that we ignore.) From this action one easily derives the correlations of the Fourier components of the $`\stackrel{}{\sigma }`$ and $`\stackrel{}{\psi }`$ fields:
$`N^1\stackrel{}{\sigma }(\stackrel{}{k},t)\stackrel{}{\sigma }(\stackrel{}{k},t^{})`$ $`=`$ $`2T{\displaystyle _0^{t_1}}𝑑t^{\prime \prime }r(k;t,t^{\prime \prime })r(k;t^{},t^{\prime \prime })`$ (52)
$`=`$ $`r(k;t,t_1)r(k;t^{},t_1)\stackrel{}{\sigma }(\stackrel{}{k},t_1)\stackrel{}{\sigma }(\stackrel{}{k},t_1),`$
$`N^1\stackrel{}{\psi }(\stackrel{}{k},t)\stackrel{}{\psi }(\stackrel{}{k},t^{})`$ $`=`$ $`2T{\displaystyle _{t_1}^{\mathrm{min}(t,t^{})}}𝑑t^{\prime \prime }r(k;t,t^{\prime \prime })r(k;t^{},t^{\prime \prime }).`$ (53)
In Sect. 2.4 we showed that the separation in fast and slow time-scales is clear in the spatial domain. Going back to real space one has
$`S_{xytt^{}}^{(5)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\stackrel{}{\sigma }(\stackrel{}{x},t)K_\sigma (\stackrel{}{x},\stackrel{}{y};t,t^{})\stackrel{}{\sigma }(\stackrel{}{y},t^{})`$ (54)
$`{\displaystyle \frac{1}{2}}\stackrel{}{\psi }(\stackrel{}{x},t)K_\psi (\stackrel{}{x},\stackrel{}{y};t,t^{})\stackrel{}{\psi }(\stackrel{}{y},t^{}),`$
with $`K_\sigma ^1(\stackrel{}{x},\stackrel{}{y};t,t^{})=N^1\stackrel{}{\sigma }(\stackrel{}{x},t)\stackrel{}{\sigma }(\stackrel{}{y},t^{})`$ and $`K_\psi ^1(\stackrel{}{x},\stackrel{}{y};t,t^{})=N^1\stackrel{}{\psi }(\stackrel{}{x},t)\stackrel{}{\psi }(\stackrel{}{y},t^{})`$ the Fourier transforms of (52) and (53), respectively.
Focusing on equal space points, $`\stackrel{}{x}=\stackrel{}{y}`$, and taking the limit $`tt^{}t_1\tau _t`$ as in , eq. (54) is the action in the generating functional for the slow and fast components of the global correlation given in eqs. (29) and (30), respectively.
## 4 Time transformations
Long ago it was realised that the dynamic equations of motion of mean-field disordered models acquire, in the long waiting time limit and for large separations of times, an invariance under generic reparameterisation of time -. This symmetry initially appeared as a nuisance since it was related to the impossibility of determining the equivalent of the scaling function $`L(t)`$ analytically. More recently, we tried to use this symmetry as a guideline to predict the main fluctuations in finite dimensional systems undergoing glassy dynamics -. With this aim we first analysed the symmetry properties of the action of the $`d`$-dimensional Edwards-Anderson spin-glass . Let us here recall the definition of the time-reparametrisation, how it acts on the fields, and check whether this invariance exists in the $`O(N)`$ model.
### 4.1 Global time-reparametrisation
Global monotonic time-reparametrisation is defined as
$$t\stackrel{~}{t}h(t)$$
(55)
with $`h(t)`$ any monotonic function of time. A particular subset of transformations are re-scalings of time
$$t\zeta t\text{that correspond to}h(t)=\zeta t.$$
(56)
The transformation (55) acts on the fields $`\stackrel{}{\varphi }(\stackrel{}{x},t)`$ and $`\stackrel{}{\widehat{\varphi }}(\stackrel{}{x},t)`$ as
$`\stackrel{}{\varphi }(\stackrel{}{x},t)`$ $``$ $`\stackrel{~}{\stackrel{}{\varphi }}(\stackrel{}{x},t)\stackrel{}{\varphi }(\stackrel{}{x},h(t)),`$ (57)
$`\stackrel{}{\widehat{\varphi }}(\stackrel{}{x},t)`$ $``$ $`\stackrel{~}{\stackrel{}{\widehat{\varphi }}}(\stackrel{}{x},t){\displaystyle \frac{dh(t^{})}{dt^{}}}\stackrel{}{\widehat{\varphi }}(\stackrel{}{x},h(t)).`$ (58)
Consequently, the space and time dependent two-point functions transform as
$`C(\stackrel{}{x},\stackrel{}{x}^{};t,t^{})`$ $``$ $`\stackrel{~}{C}(\stackrel{}{x},\stackrel{}{x}^{};t,t^{})C(\stackrel{}{x},\stackrel{}{x}^{};h(t),h(t^{})),`$ (59)
$`R(\stackrel{}{x},\stackrel{}{x}^{};t,t^{})`$ $``$ $`\stackrel{~}{R}(\stackrel{}{x},\stackrel{}{x}^{};t,t^{}){\displaystyle \frac{dh(t^{})}{dt^{}}}R(\stackrel{}{x},\stackrel{}{x}^{};h(t),h(t^{})).`$ (60)
All spatial positions transform in the same way under the simultaneous transformation of the two-times. The Fourier components transform in an identical way.
The choice of the transformation of the fields is such that the integrated linear response transforms as the correlation under these reparametrisations of time:
$`\chi (\stackrel{}{x},\stackrel{}{x}^{};t,t^{})\stackrel{~}{\chi }(\stackrel{}{x},\stackrel{}{x}^{};t,t^{})={\displaystyle _t^{}^t}𝑑t^{\prime \prime }\stackrel{~}{R}(\stackrel{}{x},\stackrel{}{x}^{};t,t^{})`$
$`={\displaystyle _t^{}^t}𝑑t^{\prime \prime }\left({\displaystyle \frac{dh(t^{\prime \prime })}{dt^{\prime \prime }}}\right)R(\stackrel{}{x},\stackrel{}{x}^{};h(t),h(t^{\prime \prime }))={\displaystyle _h^{}^h}𝑑h^{\prime \prime }R(\stackrel{}{x},\stackrel{}{x}^{};h,h^{\prime \prime })`$
$`=\chi (\stackrel{}{x},\stackrel{}{x}^{};h,h^{}).`$
It is interesting to notice that the transformation in (57) and (58) does not leave all terms in the Martin-Siggia-Rose action invariant. If we write this action in its most general form
$$Sd^dx𝑑t\left[T(i\widehat{\varphi }(\stackrel{}{x},t))^2+i\widehat{\varphi }(\stackrel{}{x},t)\frac{\varphi (\stackrel{}{x},t)}{t}+i\widehat{\varphi }(\stackrel{}{x},t)\frac{\delta V[\varphi (\stackrel{}{x},t)]}{\delta \varphi (\stackrel{}{x},t)}\right]$$
we note that the first and second terms are not invariant while the last one is. This is not surprising since a particular evolution, i.e. a particular $`h(t)`$, has to be chosen by the dynamic action. It is only the slow dynamics, which is generated in some models, that may acquire full time-reparametrisation (or a reduced) invariance. We shall come back to this important point below.
### 4.2 Symmetries in the dynamic equations
Following the same route as in the study of the dynamics of disordered spin models, let us first examine whether the dynamic equations for the global correlation and response of the $`O(N)`$ model become invariant under generic reparametrisation of time in the scaling regime of long waiting-time ($`t^{}\tau _t`$) and for very separated times ($`tt^{}t^{}`$).
With this aim, we first derive closed-form dynamic equations for the global correlation and response of the $`O(N)`$ model. We then show that these are not invariant under the most generic time-reparametrisation defined in eqs. (59) and (60) but only under the subgroup of time-re-scalings given in eq. (56).
#### 4.2.1 Dynamic equations for the global correlation and response
In A we show that the dynamic equation for the global linear response takes the form
$`{\displaystyle \frac{R(t,t^{})}{t}}=z(t)R(t,t^{})`$
$`+{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}A_n{\displaystyle 𝑑t_n𝑑t_{n1}\mathrm{}𝑑t_1R(t,t_1)R(t_1,t_2)\mathrm{}R(t_n,t^{})}`$ (61)
for all spherical models with arbitrary two-body interactions and all $`O(N)`$ models in the limit $`N\mathrm{}`$ with arbitrary two-body elastic energy. The only requirement for this result to hold is that the energy band must have a finite edge. The coefficients $`A_n`$ are determined by the density of states of the interaction matrix or elastic ‘coefficients’ and thus depend on dimensionality. In particular, for the $`p=2`$ spherical spin-glass with interactions chosen from a Gaussian distribution with zero mean and variance of order $`1/N`$, $`A_0=0`$, $`A_10`$, $`A_{n2}=0`$ and the series truncates at $`n_{max}=1`$. For a general density of states $`n_{max}\mathrm{}`$. Dilute models are excluded from this family since their densities of states have long-tails . It is interesting to note that the dynamics of the response is decoupled from that of the correlation for all these models.
Putting this equation in the Schwinger-Dyson form
$`{\displaystyle \frac{R(t,t^{})}{t}}=z(t)R(t,t^{})+{\displaystyle 𝑑t_n\mathrm{\Sigma }(t,t_n)R(t_n,t^{})}`$ (62)
allows us to identify the self-energy:
$$\mathrm{\Sigma }(t,t^{})=\underset{n=0}{\overset{\mathrm{}}{}}A_n𝑑t_{n1}𝑑t_{n2}\mathrm{}𝑑t_1R(t,t_1)\mathrm{}R(t_{n1},t^{}).$$
(63)
Since we are not considering the possibility of applying non-potential forces, the Schwinger-Dyson equation for the global correlation should read
$`{\displaystyle \frac{C(t,t^{})}{t}}=z(t)C(t,t^{})+{\displaystyle 𝑑t_n[\mathrm{\Sigma }(t,t_n)C(t_n,t^{})+D(t,t_n)R(t^{},t_n)]}`$ (64)
with $`D(t,t^{})`$ the vertex kernel. If the model has an equilibrium high temperature phase the vertex should be related to the self-energy in such a way that the solution verifies the fluctuation-dissipation theorem. This is achieved by
$`\mathrm{\Sigma }(tt^{})={\displaystyle \frac{1}{T}}{\displaystyle \frac{D(tt^{})}{t^{}}}\theta (tt^{})`$
in the high $`T`$ phase. One can then guess that
$$D(t,t_n)=\underset{n=0}{\overset{\mathrm{}}{}}A_n𝑑t_1\mathrm{}𝑑t_{n1}R(t,t_1)\mathrm{}R(t_{n2},t_{n1})C(t_{n1},t_n).$$
(65)
Note that
$$\mathrm{\Sigma }(t,t_n)=𝑑t_a𝑑t_b\frac{\delta D(t,t_n)}{\delta C(t_a,t_b)}R(t_a,t_b).$$
The equal-time global correlation $`C(t,t)d^dx[\varphi ^2(\stackrel{}{x},t)]_{ic}`$ may not, in general, be fixed to a constant value. The dynamic equation determining its time-evolution is obtained by writing $`d_tC(t,t)lim_{t^{}t}[_tC(t,t^{})+_t^{}C(t,t^{})]`$ using eq. (64). The Lagrange multiplier $`z(t)`$ is in general determined by eq. (19) while one should use the equation for $`C(t,t)`$ to compute the average $`[\varphi ^2(\stackrel{}{x},t)]_{ic}`$.
#### 4.2.2 Solution in the ordered phase
Let us assume that the global correlation and response on the one side, and the self-energy and the vertex on the other, separate in a fast and a slow component as in eqs. (28)-(33). We then introduce this Ansatz in eqs. (62) and (64) to derive dynamic equations for the fast and slow parts . The equations of motion for the slow parts have the form
$`{\displaystyle \frac{C_{ag}(t,t^{})}{t}}=z(t)C_{ag}(t,t^{})+int_C`$
$`{\displaystyle \frac{R_{ag}(t,t^{})}{t}}=z(t)R_{ag}(t,t^{})+int_R`$
with $`int_C`$ and $`int_R`$ being two series of rather complicated terms involving $`n`$-order convolutions of the response and the correlation over the times. Clearly, the time-derivatives on the left-hand-side are not invariant under a generic reparametrisation of time. A necessary step in trying to prove time-reparametrisation invariance is to assume that asymptotically they are much smaller than each term on the right-hand-side, drop them, and check the invariance of the remaining terms. This is an assumption that should be checked a posteriori once the solution for $`C_{ag}`$ and $`R_{ag}`$ is derived from the remaining equations. In the case of the $`O(N)`$ model we already know the exact solution for all times, from which we can derive the approximate form that holds in the scaling limit of very long times and separations among them, and check whether this form allows for the time-reparametrisation invariance of the equations.
Let us first focus on the equation for the global response which is easier to deal with. The slow ageing part of the linear response behaves asymptotically as $`t^{d/2}f_R(\lambda )`$, see eq. (36). Its time-derivative is
$$\frac{R_{ag}(t,t^{})}{t}t^{d/21}\left[\frac{d}{2}f_R(\lambda )+\lambda f_R^{^{}}(\lambda )\right].$$
The ‘mass’ $`z(t)`$ decays as
$$z(t)=\frac{d}{dt}\left(\frac{1}{2}\mathrm{ln}Y^2(t)\right)\frac{1}{2}\frac{d}{dt}\mathrm{ln}t^{d/2}=\frac{d}{4}t^1,$$
(66)
consequently, the first term in the right-hand-side goes as
$`z(t)R_{ag}(t,t^{}){\displaystyle \frac{d}{4}}t^{d/21}f_R(\lambda )`$
and it is of the same order as the time-derivative in the left-hand-side.
The terms in the series can also be analysed by separating the stationary and ageing contributions to the integrals in eq. (61) (see for a detailed explanation). Such separation, as carried out in B, leads to
$`{\displaystyle \frac{R_{ag}(t,t^{})}{t}}=z(t)R_{ag}(t,t^{})`$ (67)
$`+{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\stackrel{~}{A}_n{\displaystyle 𝑑t_n𝑑t_{n1}\mathrm{}𝑑t_1R_{ag}(t,t_1)R_{ag}(t_1,t_2)\mathrm{}R_{ag}(t_n,t^{})},`$
where the coefficients $`\stackrel{~}{A}_n`$ are given by \[see eq. (104)\]
$$\stackrel{~}{A}_n\frac{1}{(n+1)!}\left(\frac{d}{d\chi _{st}}\right)^nϵ(\chi _{st}),$$
(68)
with $`\chi _{st}`$ the integrated stationary response evaluated at time scales of order $`\mathrm{\Delta }t=tt^{}`$. The function $`ϵ(h)`$ is the inverse of the function
$$h(ϵ)=_0^{\mathrm{}}𝑑ϵ^{}\frac{g(ϵ^{})}{ϵ^{}ϵ}$$
which is obtained from the density of states $`g(ϵ)`$ of the model (see A and B).
In the $`O(N)`$ model the coefficients $`\stackrel{~}{A}_n`$ scale as a power of $`\mathrm{\Delta }t`$, and the precise power is controlled by the form of the density of states at low energies, $`g(ϵ)ϵ^\nu `$, with $`\nu =d/21`$ \[see eq. (5)\]. The scaling of $`\stackrel{~}{A}_n`$ at long time differences can be obtained as follows. First, notice that $`h(0)h(ϵ)ϵ^\nu `$ by taking into account the power law dependence of the density of states at low energies. Then, the inverse function $`ϵ(h)`$ is $`ϵ=h^1(h(ϵ))(h(0)h(ϵ))^{1/\nu }`$. On the other hand, the stationary susceptibility $`\chi _{st}(\mathrm{\Delta }t)`$ is intimately related to the density of states. Using eq. (A) one finds,
$$\chi _{st}(\mathrm{\Delta }t)=_{\mathrm{}}^{\mathrm{\Delta }t}𝑑\mathrm{\Delta }t^{}G(\mathrm{\Delta }t^{})=_0^{\mathrm{}}𝑑ϵg(ϵ)\frac{1e^{ϵ\mathrm{\Delta }t}}{ϵ}=\chi _{st}(\mathrm{})_0^{\mathrm{}}𝑑ϵ\frac{g(ϵ)}{ϵ}e^{ϵ\mathrm{\Delta }t}.$$
Using now the power law decay of the density of states at low energies one finds
$$\chi _{st}(\mathrm{})\chi _{st}\mathrm{\Delta }t^\nu \mathrm{or}1/\mathrm{\Delta }t[\chi _{st}(\mathrm{})\chi _{st}]^{1/\nu },$$
(69)
with $`\chi _{st}\chi _{st}(\mathrm{\Delta }t)`$. From eq. (98) one notices that $`\chi _{st}(\mathrm{})=h(0)`$. Thus,
$$ϵ(\chi _{st})[\chi _{st}(\mathrm{})\chi _{st}]^{1/\nu },$$
(70)
so that $`d^nϵ/d\chi _{st}^n[\chi _{st}(\mathrm{})\chi _{st}]^{1/\nu n}\mathrm{\Delta }t^{1+n\nu }`$ which yields
$$\stackrel{~}{A}_n\mathrm{\Delta }t^{1+n\nu }.$$
(71)
Notice that the coefficients $`\stackrel{~}{A}_n`$ keep a power law dependence on time differences. This anomalous dependence is in remarkable contrast to glassy systems, such as the $`p`$-spin disordered models with $`p3`$, in which the coefficients $`\stackrel{~}{A}_n`$ reach finite constants as $`\mathrm{\Delta }t\mathrm{}`$ and $`\chi _{st}\chi _{st}(\mathrm{})`$. We discuss the consequences of this difference below, when we consider the scaling dimensions of the global correlation and response under time re-scalings.
#### 4.2.3 Scaling dimensions
Consider the situation in which under, say, a scale transformation, the global correlation and response in the ageing regime transform according to
$`t\zeta t`$ (72)
$`C_{ag}(t,t^{})\stackrel{~}{C}_{ag}(t,t^{})=C_{ag}(\zeta t,\zeta t^{})`$ (73)
$`R_{ag}(t,t^{})\stackrel{~}{R}_{ag}(t,t^{})=\zeta ^{\mathrm{\Delta }_R}R_{ag}(\zeta t,\zeta t^{}),`$ (74)
where $`\mathrm{\Delta }_R`$ is the retarded dimension of the response (the advanced dimensions for both response and correlation, as well as the retarded dimension for the correlation, are zero in the above). In systems in which correlation and response are related by an off-equilibrium fluctuation-dissipation relation with a finite effective temperature , the retarded dimension takes the value $`\mathrm{\Delta }_R=1`$. In this case, one can show that if the equations of motion are invariant under these scale transformations, they are also invariant under time-reparametrisation $`th(t)`$. In other words, for the special case $`\mathrm{\Delta }_R=1`$, scale invariance implies reparametrisation invariance . The situation is similar in character to what happens with scale invariant field theories, where in the special case of two dimensional systems local scale invariance implies conformal invariance, a much larger symmetry.
Let us discuss concisely why the symmetry is larger when $`\mathrm{\Delta }_R=1`$. Consider the ageing contribution to a generic term $`I_n`$ with $`n`$ integrals and $`n+1`$ responses as an example:
$`\stackrel{~}{I}_n(t,t^{})`$ $``$ $`{\displaystyle 𝑑t_n𝑑t_{n1}\mathrm{}𝑑t_1\stackrel{~}{R}_{ag}(t,t_1)\stackrel{~}{R}_{ag}(t_1,t_2)\mathrm{}\stackrel{~}{R}_{ag}(t_n,t^{})}`$ (75)
$`=`$ $`{\displaystyle 𝑑t_n\mathrm{}𝑑t_1\zeta ^{(n+1)\mathrm{\Delta }_R}R_{ag}(\zeta t,\zeta t_1)\mathrm{}R_{ag}(\zeta t_n,\zeta t^{})}`$
$`=`$ $`\zeta ^{(n+1)\mathrm{\Delta }_Rn}{\displaystyle d(\zeta t_n)\mathrm{}d(\zeta t_1)R_{ag}(\zeta t,\zeta t_1)\mathrm{}R_{ag}(\zeta t_n,\zeta t^{})}`$
$`=`$ $`\zeta ^{n(\mathrm{\Delta }_R1)}\zeta ^{\mathrm{\Delta }_R}I_n(\zeta t,\zeta t^{}).`$
Thus, under the rescaling transformation, $`I_n(t,t^{})`$ has dimension
$$\mathrm{\Delta }_{I_n}=n(\mathrm{\Delta }_R1)+\mathrm{\Delta }_R.$$
There are two special features that arise when $`\mathrm{\Delta }_R=1`$. The first is that all the $`I_n(t,t^{})`$ have the same dimension, $`\mathrm{\Delta }_{I_n}=1`$ for all $`n`$. The second is that the change of variables inside the integrals can be carried out for a more general change of variables $`th(t)`$ with an arbitrary monotonic function $`h(t)`$, because
$$𝑑t_i\left(\frac{dh(t_i)}{dt_i}\right)^{\mathrm{\Delta }_R=1}\mathrm{}=𝑑t_i\frac{dh(t_i)}{dt_i}\mathrm{}=𝑑h_i\mathrm{}$$
holds for each of the times that are being integrated over. Therefore, for $`\mathrm{\Delta }_R=1`$ the rescaling of the correlations and responses can be absorbed in the Jacobian for the changes of integration variables.
Hence for $`\mathrm{\Delta }_R=1`$ scale invariance implies reparametrisation invariance.
#### 4.2.4 Fixing the retarded dimension $`\mathrm{\Delta }_R`$
Let us discuss now how the asymptotic behaviour of $`z(t)`$ fixes the retarded dimension $`\mathrm{\Delta }_R`$.
The case of glassy dynamics
In mean-field spherical models displaying glassy dynamics, such as for example the spherical $`p`$-spin model for $`p3`$, the right-hand-side of the dynamic equation for the response is much simpler that eq. (61) in that the series actually has only one term of the form $`I_1`$. A supplementary difficulty arises from the fact that $`C`$ enters the integral but this is not very difficult to deal with if we assume that the scaling dimensions of $`C`$ vanish.
The function $`z(t)z_{\mathrm{}}0`$ as $`t\mathrm{}`$. In this case, the term
$$z(t)R_{ag}(t,t^{})z_{\mathrm{}}R_{ag}(t,t^{})$$
has the same scaling dimension as the response itself, i.e., $`\mathrm{\Delta }_R`$. In these glassy systems, the coefficients $`\stackrel{~}{A}_n`$ in front of the integral terms are finite constants in the limit of $`\mathrm{\Delta }t=tt^{}\mathrm{}`$. (In the specific case of the $`p`$-spin models, only $`\stackrel{~}{A}_0`$ and $`\stackrel{~}{A}_1`$ are non-vanishing.)
The time-derivative term
$$\frac{}{t}R_{ag}(t,t^{})$$
has scaling dimension $`\mathrm{\Delta }_{\mathrm{deriv}.}=\mathrm{\Delta }_R+1`$, because the $`\frac{}{t}`$ has dimension 1. So in this case the time-derivative term is irrelevant, and it can be dropped in the long-time limit, as long as one finds a non-trivial solution to the remaining equations. Indeed, there can be a non-trivial solution of the long-time dynamical equations if at least one of the integral terms, $`I_n(t,t^{})`$ for some $`n0`$, can balance the $`z_{\mathrm{}}R_{ag}(t,t^{})`$ term. As we discussed previously, the contribution to a generic integral $`I_n`$ with $`n0`$ has scaling dimension $`\mathrm{\Delta }_{I_n}=n(\mathrm{\Delta }_R1)+\mathrm{\Delta }_R`$. The cancellation can be achieved if and only if
$$\mathrm{\Delta }_{I_n}=n(\mathrm{\Delta }_R1)+\mathrm{\Delta }_R=\mathrm{\Delta }_R\text{for some}n0.$$
There are two ways of achieving this scope. The contribution with $`n=0`$ trivially satisfies this identity for any value of $`\mathrm{\Delta }_R`$. Besides, terms of the same order arise from $`n1`$ only if $`\mathrm{\Delta }_R=1`$, in which case the $`n`$ dependence disappears and the condition is actually satisfied for all $`n1`$. The second possibility is realized by the $`p>2`$ spherical Gaussian spin-glass, a model with
$$\mathrm{\Delta }_R=1$$
and for which reparametrisation invariance develops. The $`p=2`$ spherical spin-glass with Gaussian interactions is discussed in detail in C.
One can argue that the scaling dimensions zero for the correlation (both retarded and advanced dimensions) and retarded dimension $`\mathrm{\Delta }_R=1`$ for the response are consistent with a factor $`X(t,t^{})=T/T_{\mathrm{eff}}`$ that remains finite for fixed $`C_{ag}`$ in the long time limit. Consider the out-of-equilibrium fluctuation-dissipation relation:
$$R_{ag}(t,t^{})=\frac{X(t,t^{})}{T}\frac{}{t^{}}C_{ag}(t,t^{})\theta (tt^{}).$$
(76)
If the factor $`X(t,t^{})X(C_{ag})`$ for fixed $`C_{ag}`$ in the large $`t,t^{}`$ limit, without vanishing with some anomalous extra powers of $`t`$, then it follows that the retarded dimension of the response is one more than that of the correlation, because of the $`/t^{}`$. So if the correlation has retarded dimension zero, the response will have retarded dimension $`\mathrm{\Delta }_R=1`$ as long as $`X`$ remains finite and has no anomalous power law dependence on $`t`$ for fixed $`C_{ag}`$. This is the situation in glassy systems, where finite factors $`X=T/T_{\mathrm{eff}}`$ have been observed in experiments and simulations (see for a review).
The case of the $`O(N)`$ model
Consider a case in which the function $`z(t)t^{\mathrm{\Delta }_z}0`$ as $`t\mathrm{}`$. Particularly, $`\mathrm{\Delta }_z=1`$ for the $`O(N)`$ model \[see eq. (66)\]. Now, the term
$$z(t)R_{ag}(t,t^{})$$
has scaling dimension $`\mathrm{\Delta }_R+\mathrm{\Delta }_z`$. The time-derivative term, just as in the case of glassy systems above, has scaling dimension $`\mathrm{\Delta }_{\mathrm{deriv}}=\mathrm{\Delta }_R+1`$. Therefore, as opposed to the cases discussed above, one cannot naively neglect the time-derivative term, because it has the same scaling dimension as the $`z(t)R_{ag}(t,t^{})`$ term.
In the $`O(N)`$ model the series in the r.h.s. of eq. (61) does not truncate. The prefactor of the integral term $`I_n(t,t^{})`$ depends on $`\mathrm{\Delta }t=tt^{}`$, $`\stackrel{~}{A}_n\mathrm{\Delta }t^{1+n\nu }`$ with $`\nu =1d/2`$ \[see eq. (71)\], and there is an additional scaling dimension arising from the anomalous scaling of the prefactors $`\stackrel{~}{A}_n`$:
$$\mathrm{\Delta }_{\stackrel{~}{A}_n}=1n\nu .$$
In order to determine the dimension $`\mathrm{\Delta }_R`$, one must balance at least one of the integral terms, $`\stackrel{~}{A}_n(\mathrm{\Delta }t)I_n(t,t^{})`$ for some $`n0`$, against the $`z(t)R_{ag}(t,t^{})`$ term and the time-derivative $`R_{ag}(t,t^{})/t`$. This can be achieved if and only if
$$\mathrm{\Delta }_{\stackrel{~}{A}_n}+\mathrm{\Delta }_{I_n}=n(\mathrm{\Delta }_R1\nu )+\mathrm{\Delta }_R+1=\mathrm{\Delta }_R+1\text{for some}n0.$$
(77)
Notice that this condition is satisfied in particular by $`n=0`$, but it can also be satisfied for any $`n`$ if
$$\mathrm{\Delta }_R=\nu +1=d/2,$$
which is indeed consistent with the exact result given in eq. (33).
Notice that all terms in the equation of motion of $`R_{\mathrm{ag}}(t,t^{})`$ have the same scaling dimension $`\mathrm{\Delta }_R+1`$ as the time derivative term, which thus cannot be dropped in any dimension $`d`$, in contrast to the case in glassy systems. Notice also that $`\mathrm{\Delta }_R1`$ for $`d>2`$, so reparametrisation invariance does not develop; only scale invariance is a symmetry of the long-time dynamical equations of motion. A retarded scaling dimension $`\mathrm{\Delta }_R>1`$ implies, using eq. (76), that the factor $`X(t,t^{})0`$ for long times and fixed $`C_{ag}`$, if the correlation has retarded and advanced dimensions zero. This result is in agreement with the direct calculation of the factor $`X`$ in eq. (37).
In $`d=2`$, one obtains that $`\mathrm{\Delta }_R=1`$, but in contrast to the case of glassy dynamics where the prefactors $`\stackrel{~}{A}_n`$ were constant, $`\mathrm{\Delta }_{\stackrel{~}{A}_n}=1`$. For reparametrisation invariance to develop, it is necessary that $`\mathrm{\Delta }_R=1`$ and that $`\mathrm{\Delta }_{\stackrel{~}{A}_n}=0`$. Hence, reparametrisation invariance does not develop even in the $`d=2`$ case. Notice, however, that $`\mathrm{\Delta }_R=1`$ in $`d=2`$ implies a non-trivial $`X(t,t^{})`$ which is actually found in the exact solution ; there is still an additive separation of correlation and linear response in a stationary and an ageing part at $`T=0`$ (as opposed to the multiplicative scaling found in critical relaxations ). Nevertheless, it is important to remark that this $`X(t,t^{})`$ depends continuously on the ratio $`t/t^{}`$ \[see eq. (37)\], which implies a $`C_{ag}`$ dependent effective temperature instead of a constant effective temperature; the latter is expected for a problem with a single correlation scale.
### 4.3 Conjecture
We argued that $`\mathrm{\Delta }_R=1`$ is a necessary condition for having an asymptotic time-reparametrization invariance - though this condition is not sufficient, as shown by the $`d=2`$ $`O(N)`$ case. In addition, $`\mathrm{\Delta }_R=1`$ implies a finite integrated linear response and a finite effective temperature, as can be derived from eq. (76).
The $`O(N)`$ model has a weaker response than that of glassy models, as for example the $`p`$-spin spherical disordered system. Indeed, in the $`O(N)`$ model the ageing contribution to the integrated response vanishes asymptotically in all $`d>d_L=2`$ – and this can be related to the development of an infinite effective temperature at long times; while it approaches a finite $`C`$-dependent value in $`d=d_L=2`$ – and this cannot be interpreted in terms of an effective temperature since one would have a $`C`$ dependent value within a single correlation scale . Other solvable coarsening problems have a similar integrated response (see e.g. ).
On the basis of the discussion above, we conjecture that models with a finite and well-defined effective temperature, such as the $`p`$ spin spherical disordered system with $`p3`$ or the more complex Sherrington-Kirkpatrick spin-glass, develop time-reparametrisation invariance asymptotically, while this does not occur in systems with a diverging or ill-defined effective temperature, such as the $`O(N)`$ model.
#### 4.3.1 Space-time rescaling
For the sake of comparison, in the following we consider the standard dynamical scaling which consists of simultaneous rescaling of time and space.
So far we discussed time rescaling and time-reparametrisation invariance properties in the real space representation. This is because we have been interested in making contact with glassy systems for which composite fields, that are related to the two-time correlation and response functions at equal space points, might be the natural order parameters . In the case of the $`O(N)`$ model, however, one knows that the original field $`\varphi _\alpha (\stackrel{}{x},t)`$ is already the natural order parameter. Especially, one expects its Fourier space representation $`\varphi _\alpha (\stackrel{}{k},t)`$ to be easier to handle.
Let us take the response function $`\delta \varphi _\alpha (\stackrel{}{k},t)/\delta h_\beta (\stackrel{}{k}^{},t^{})|_{h=0}=r(k;t,t^{})\delta _{\alpha \beta }\delta (\stackrel{}{k}+\stackrel{}{k}^{})`$ and the two-time composite field $`\varphi _\alpha (\stackrel{}{k},t)\varphi _\beta (\stackrel{}{k}^{},t^{})=c(k;t,t^{})\delta _{\alpha \beta }\delta (\stackrel{}{k}+\stackrel{}{k}^{})`$ where $`k=|\stackrel{}{k}|`$. In the $`T0`$ limit <sup>2</sup><sup>2</sup>2This is general since the domain growth scaling is controlled by a “zero temperature fixed point” . these quantities satisfy the same evolution equation,
$`{\displaystyle \frac{}{t}}r(k;t,t^{})=[k^2+z(t)]r(k;t,t^{}),`$ (78)
$`{\displaystyle \frac{}{t}}c(k;t,t^{})=[k^2+z(t)]c(k;t,t^{}).`$ (79)
Now let us suppose that these equations admit a set of asymptotic solutions with the following scaling forms<sup>3</sup><sup>3</sup>3The simple analysis that we present here yields the correct behaviour for non-conserved but it is does not in the case of dynamics with conserved order-parameter.
$`r(k;t,t^{})=\zeta ^{\mathrm{\Delta }_r^R+\mathrm{\Delta }_r^Ad\mathrm{\Delta }_s}r(k\zeta ^{\mathrm{\Delta }_s},\zeta t,\zeta t^{}),`$ (80)
$`c(k;t,t^{})=\zeta ^{\mathrm{\Delta }_c^R+\mathrm{\Delta }_c^Ad\mathrm{\Delta }_s}c(k\zeta ^{\mathrm{\Delta }_s},\zeta t,\zeta t^{}),`$ (81)
and
$$z(t)=\zeta ^{\mathrm{\Delta }_z}z(\zeta t).$$
$`\mathrm{\Delta }_r^R`$ and $`\mathrm{\Delta }_r^A`$ are the retarded and advanced dimensions of the response $`r`$; similarly, $`\mathrm{\Delta }_c^R`$ and $`\mathrm{\Delta }_c^A`$ are the retarded and advanced dimensions of the correlation $`c`$. The scaling dimensions are then fixed by inserting this Ansatz in eqs. (78) and (79). First, focusing on the $`k=0`$ component, one finds $`\mathrm{\Delta }_z=1`$. Then considering the $`k>0`$ components, one finds two possibilities. The first one is that the $`k^2`$ term has the same scaling dimension as the other two terms. In this case one finds the exponent for spatial scaling $`\mathrm{\Delta }_s=1/2`$ <sup>4</sup><sup>4</sup>4Its inverse corresponds to the dynamical exponent $`z`$ in critical dynamics.. The other possibility is $`\mathrm{\Delta }_s>1/2`$ which means that the $`k^2`$ term becomes irrelevant. Which of the two cases appear in the large times regime depends on the initial conditions. For usual random initial condition of the form given in eq. (4) with the same statistical weight on all $`k`$ components, the exact solution summarised in Sect. 2.4 tells us that $`\mathrm{\Delta }_s=1/2`$ is actually selected. In the following we only consider this case.
In the asymptotic regime all the terms in eqs. (78) and (79) have the same scaling dimensions. Thus none of them can be dropped irrespective of the scaling dimensions of the response and correlation functions which will be determined below.
We still need to determine the retarded and advanced scaling dimensions of the response and correlation functions. Since the solution of the Langevin equation at $`T=0`$ can be written as \[see eq. (22)\]
$$\varphi (\stackrel{}{k},t)=r(k;t,t^{})\varphi (\stackrel{}{k},t^{}),$$
the overall scaling factor in $`r`$ must be identical to one, and one has that retarded and advanced scaling dimensions of the response functions must satisfy $`\mathrm{\Delta }_r^R+\mathrm{\Delta }_r^A=d\mathrm{\Delta }_s=d/2`$. This is achieved by $`\mathrm{\Delta }_r^R=d\mathrm{\Delta }_s=d/2`$ and $`\mathrm{\Delta }_r^A=0`$. The analysis of the self-consistent equation for $`z(t)`$ fixes the scaling dimensions $`\mathrm{\Delta }_c^R`$ and $`\mathrm{\Delta }_c^A`$. Indeed, eq. (2) reads
$$z(t)=g\frac{d^dk}{(2\pi )^d}c(k,t,t)+r.$$
In the large time limit $`z(t)0`$ and the integral converges to $`r/g`$. Therefore $`\mathrm{\Delta }_c^R+\mathrm{\Delta }_c^A=0`$ that implies the natural choice $`\mathrm{\Delta }_c^R=\mathrm{\Delta }_c^A=0`$.
It is instructive to consider the inverse Fourier transform of the scaling Ansatz in eqs. (80) and (81) which reads
$`R(|\stackrel{}{x}\stackrel{}{y}|;t,t^{})`$ $`=`$ $`\zeta ^{\mathrm{\Delta }_r^A+\mathrm{\Delta }_r^R}R(|\stackrel{}{x}\stackrel{}{y}|\zeta ^{\mathrm{\Delta }s},\zeta t,\zeta t^{}),`$
$`C(|\stackrel{}{x}\stackrel{}{y}|;t,t^{})`$ $`=`$ $`\zeta ^{\mathrm{\Delta }_c^A+\mathrm{\Delta }_c^R}C(|\stackrel{}{x}\stackrel{}{y}|\zeta ^{\mathrm{\Delta }s},\zeta t,\zeta t^{}).`$
Thus the solution can be written in the scaling form \[see (34)\],
$`R(|\stackrel{}{x}\stackrel{}{y}|;t,t^{})`$ $``$ $`{\displaystyle \frac{1}{L^d(t)}}f_{R_r}({\displaystyle \frac{|\stackrel{}{x}\stackrel{}{y}|}{L(t)}},{\displaystyle \frac{L(t)}{L(t^{})}}),`$
$`C(|\stackrel{}{x}\stackrel{}{y}|;t,t^{})`$ $``$ $`f_{C_r}({\displaystyle \frac{|\stackrel{}{x}\stackrel{}{y}|}{L(t)}},{\displaystyle \frac{L(t)}{L(t^{})}}),`$
with the domain growth law
$$L(t)t^{\mathrm{\Delta }_s}=\sqrt{t}.$$
Thus, the analysis of the invariance of the dynamic equations for the global $`C`$ and $`R`$ under rescaling of time presented in the previous sections, that serves to fix the scaling dimension of the global response $`\mathrm{\Delta }_R`$, can be extended to study the scaling dimensions of the space-dependent correlation and response $`C(r;t,t^{})`$ and $`R(r;t,t^{})`$ under simultaneously rescaling of space and time. We find $`\mathrm{\Delta }_{R_r}^R=d/2`$ and $`\mathrm{\Delta }_{R_r}^A=0`$, and $`\mathrm{\Delta }_{C_r}^A=0`$ and $`\mathrm{\Delta }_{C_r}^R=0`$, for any $`d`$.
### 4.4 Symmetries of the (long-times) action
We have already mentioned that a generic Martin-Siggia-Rose action is not invariant under reparametrisations of the times and fields defined in eqs. (57)-(58). Let us now analyse the symmetries of the action in the long times limit in which there is a separation of time-scales in the global correlation and response.
Let us first focus on the case $`d>2`$. Using the simple transformations described in Sect. 3 the dynamic generating function of the $`O(N)`$ model can be expressed as a path integral over the fields $`\stackrel{}{\sigma }(\stackrel{}{x},t)`$ and $`\stackrel{}{\psi }(\stackrel{}{x},t)`$ only. On the one hand one can argue that the ‘fast’ $`\stackrel{}{\psi }`$-part ‘renormalises’ to zero under generic time-reparametrisation and becomes asymptotically irrelevant. On the other hand, while the $`\stackrel{}{\sigma }`$ field transforms in such a way that it ensures the correct transformation of the integration measure, the kernel in its action is just the inverse of the global correlation itself that is time dependent and transforms in a non-trivial manner under generic reparametrisations of time. The quadratic action for $`\stackrel{}{\sigma }(\stackrel{}{x},t)`$ at equal space points is not invariant under generic reparametrisations of time.
The local action of the $`\stackrel{}{\sigma }`$ field is, however, invariant under a reduced subset of transformations, namely time re-scalings. Since the kernel is a function of $`t^{}/t`$ one finds that the slow action written in terms of the fields only, see eq. (51), is invariant under
$`t\zeta t,`$
$`\stackrel{}{\sigma }(\stackrel{}{x},t)\stackrel{}{\sigma }(\stackrel{}{x},\zeta t).`$
But one can also go one step back and check whether the action for the field, $`\stackrel{}{\sigma }(\stackrel{}{x},t)`$, and response field, $`i\stackrel{}{\widehat{\sigma }}(\stackrel{}{x},t)`$, is invariant under time re-scalings that change the response field as
$$i\stackrel{}{\widehat{\sigma }}(\stackrel{}{x},t)\zeta i\stackrel{}{\widehat{\sigma }}(\stackrel{}{x},\zeta t),$$
(82)
i.e. the reduction of (58) to time re-scalings. Transforming the ‘slow’ part of $`S_{ktt^{}}^{(4)}`$ into spatial coordinates and writing explicitly all integrals one has
$`S^{(4)}`$ $`=`$ $`{\displaystyle }d^dx{\displaystyle }d^dy{\displaystyle }dt{\displaystyle }dt^{}[i\stackrel{}{\widehat{\sigma }}(\stackrel{}{x},t)K_\sigma (\stackrel{}{x}\stackrel{}{y};t,t^{})i\stackrel{}{\widehat{\sigma }}(\stackrel{}{y},t^{})`$ (83)
$`+i\stackrel{}{\widehat{\sigma }}(\stackrel{}{x},t)\delta (tt^{})\delta ^d(\stackrel{}{x}\stackrel{}{y})\stackrel{}{\sigma }(\stackrel{}{y},t^{})]`$
with
$$K_\sigma (\stackrel{}{x}\stackrel{}{y};t,t^{})T\frac{d^dk}{(2\pi )^d}e^{i\stackrel{}{k}(\stackrel{}{x}\stackrel{}{y})}e^{k^2/\mathrm{\Lambda }^2}_0^{t_1}𝑑t^{\prime \prime }r(k,t,t^{\prime \prime })r(k,t^{},t^{\prime \prime }).$$
(84)
One can then easily check that the action at equal space points, $`\stackrel{}{x}=\stackrel{}{y}`$, remains invariant under the time re-scalings proposed above. Indeed, in terms of the transformed fields the local action reads
$`\stackrel{~}{S}_x^{(4)}`$ $`=`$ $`{\displaystyle 𝑑t𝑑t^{}\left[i\stackrel{~}{\stackrel{}{\widehat{\sigma }}}(\stackrel{}{x},t)K_\sigma (\stackrel{}{0};t,t^{})i\stackrel{~}{\stackrel{}{\widehat{\sigma }}}(\stackrel{}{x},t^{})+i\stackrel{~}{\stackrel{}{\widehat{\sigma }}}(\stackrel{}{x},t)\delta (tt^{})\stackrel{~}{\stackrel{}{\sigma }}(\stackrel{}{x},t^{})\right]}`$ (85)
$`=`$ $`{\displaystyle }dt{\displaystyle }dt^{}[\zeta i\stackrel{}{\widehat{\sigma }}(\stackrel{}{x},\zeta t)K_\sigma (\stackrel{}{0};\zeta t,\zeta t^{})\zeta i\stackrel{}{\widehat{\sigma }}(\stackrel{}{x},\zeta t^{})`$
$`+\zeta i\stackrel{}{\widehat{\sigma }}(\stackrel{}{x},\zeta t)\zeta \delta (\zeta t\zeta t^{})\stackrel{}{\sigma }(\stackrel{}{x},\zeta t^{})]`$
$`=`$ $`S_x^{(4)}`$
where we used the fact that $`K_\sigma (\stackrel{}{0};t,t^{})`$ is just identical to the global correlation function, $`C(t,t^{})`$. In the limit $`t,t^{}t_0`$ this is a function of $`t^{}/t`$ and thus invariant under time-rescaling. The last identity follows simply from changing the integration variables form $`t`$ to $`\zeta t`$.
A similar treatment in $`d=2`$ is much more delicate. The explicit calculations in show that the separation of the field is achieved by taking advantage of the fact that a factor proportional to $`(\mathrm{\Lambda }^2t_1)^{1d/2}`$ vanishes \[see eqs. (46) and (47)\]. This, however, is no longer true in $`d=2`$. Besides, the interesting dynamics in this case arises only at $`T=0`$, another non-trivial limit to be taken in the asymptotic expressions. For these reasons, we cannot simply carry through the arguments above to $`d=2`$. Another way to attack the same problem would be to write an action in terms of $`R(\stackrel{}{x},\stackrel{}{x};t,t^{})`$ and $`C(\stackrel{}{x},\stackrel{}{x};t,t^{})`$ and use a similar reasoning to the one we used for the analysis of the equations of motion for $`R(t,t^{})`$ and $`C(t,t^{})`$. We shall not pursue this study here.
Let us note that the invariance under time rescaling discussed above can also be understood as a part of the usual space-time scaling invariance discussed in Sect. 4.3.1. To this end, we perform a renormalisation group (RG) analysis on the Fourier space representation of the slow part of the action,
$`S_{k<\mathrm{\Lambda }}^{(4)}[i\widehat{\sigma },\sigma ]`$ $`=`$ $`{\displaystyle _0^\mathrm{\Lambda }}d^dk{\displaystyle _0^{\mathrm{}}}dt{\displaystyle _0^{\mathrm{}}}dt^{}[i\widehat{\sigma }(k,t)K_\sigma (k,t,t^{})i\widehat{\sigma }(k,t^{})`$ (86)
$`+\delta (tt^{})i\widehat{\sigma }(k,t)\sigma (k,t^{})]`$
$`K_\sigma `$ the Fourier transform of (84). First, by integrating out the “fast modes” in $`\mathrm{\Lambda }/b<k<\mathrm{\Lambda }`$ we obtain $`S_{k<\mathrm{\Lambda }/b}^{(4)}`$. Next, we choose a set of rescaled variables
$`\begin{array}{cc}\stackrel{~}{k}=kb,\hfill & \stackrel{~}{t}=t/b^2,\hfill \\ i\stackrel{~}{\widehat{\sigma }}(\stackrel{~}{k},\stackrel{~}{t})=b^{d/2+2}i\widehat{\sigma }(k,t),\hfill & \stackrel{~}{\sigma }(\stackrel{~}{k},\stackrel{~}{t})=b^{d/2}\sigma (k,t),\hfill \\ \stackrel{~}{z}(\stackrel{~}{t})=b^2z(t).\hfill & \end{array}`$ (88)
In terms of the new variables the cut-off is put back to $`\mathrm{\Lambda }`$ and the action of the original form is recovered. Converting the above results to the real space representation and equating the scaling parameter of space $`b`$ and time $`\zeta `$ as $`b^2=\zeta `$. The space dependent slow action is invariant under simultaneous rescaling of space and time.
### 4.5 How the $`O(N)`$ escapes reparametrisation invariance but displays scale invariance
It was shown in that under rather mild assumptions (namely, causality, a separation of time-scales as the one discussed in Sect. 2, the fact that the remaining free field action does not lead itself to slow dynamics and the use of the naive scaling dimensions of the fields) the slow part of the action of the $`3d`$ Edwards-Anderson spin-glass, when written in terms of the two-time dependent dynamic order parameters, remains invariant under global time-reparametrisation. We have shown above that the action for the slow evolution of the $`O(N)`$ model is not invariant under these transformations. One would like to identify which of the assumptions used in is (are) violated in the $`O(N)`$ case. Indeed, in we used the naive dimensions for $`Q_R`$, that is to say $`\mathrm{\Delta }_{Q_R}^A=0`$ and $`\mathrm{\Delta }_{Q_R}^R=1`$. This assumption should be correct for systems that develop a finite and well-defined effective temperature in the ageing regime. The O(N) falls out of this class and this assumption does not apply to it.
## 5 The distribution of local two-time observables
During the ageing relaxation of glassy systems one expects important temporal and spatial fluctuations. The distribution of local coarse-grained correlations and linear responses in spin-glasses and kinetically facilitated models were computed numerically. The comparison of these probability distribution functions (pdfs) with the theoretical framework developed in - was also discussed. In short, the main features of these distributions are:
i. The pdf of coarse-grained local two-time correlations is a function that depends on the two times and, when these are chosen to lie in the ageing regime, the pdf scales in time just as the global correlation itself.
ii. The functional form of the pdf of coarse-grained local two-time correlations changes with the two times. It can be approximately described with a Gumbel-like function with a two-time dependent parameter, which in the ageing regime is simply a function of the global correlation . The parameter $`a`$ characterising the Gumbel-like form is positive for values of $`C`$ that are relatively large and close to the maximum given by $`q_{ea}`$: the distribution is negatively skewed. The parameter $`a`$ increases when decreasing $`C`$ and diverges at some value of $`C`$ signalling and approximately symmetric and Gaussian-like distribution. For still lower values of $`C`$ the pdf becomes positively skewed and this can be described with a negative value of the parameter $`a`$.
iii. The joint pdf of local two-time correlations and linear responses follows the global $`\chi (C)`$ curve in the ageing regime. This means that the longitudinal fluctuations that take the points out of this “master” curve become rare when the coarse-graining size increases while the transverse fluctuations along the master curve become more and more important when the waiting-time increases.
In this Section we compute these and similar distributions and we check whether the same features are observed in the $`O(N)`$ model. For the sake of simplicity we work at $`T=0`$ and we analyse the fluctuations induced by a Gaussian distribution of initial conditions keeping in mind that all calculations can be generalised to the finite temperature case. We take the $`N\mathrm{}`$ limit strictly and we do not let the “constraint” $`N^1_\alpha \varphi _\alpha ^2(\stackrel{}{x},t)`$ fluctuate (see for more details).
### 5.1 Coarse-graining the field
The $`O(N)`$ model yields a mean-field description of ferromagnetic ordering. It is then worth starting by studying the distribution of the local magnetisations coarse-grained within a region of volume $`V_{x_0}\mathrm{}^d`$ around $`\stackrel{}{x}_0`$. This quantity is defined as,
$$\stackrel{}{m}_{\mathrm{}}(\stackrel{}{x}_0,t)\frac{d^dx}{(2\pi \mathrm{}^2)^{d/2}}e^{\frac{|\stackrel{}{x}\stackrel{}{x}_0|^2}{2\mathrm{}^2}}\stackrel{}{\varphi }(\stackrel{}{x},t).$$
In terms of the Fourier transform $`\stackrel{}{\varphi }(\stackrel{}{k},t)`$ of the original field we find
$$\stackrel{}{m}_{\mathrm{}}(\stackrel{}{x}_0,t)=\frac{d^dk}{2\pi ^{d/2}}e^{i\stackrel{}{k}\stackrel{}{x_0}}\stackrel{}{\varphi }(\stackrel{}{k},t)e^{k^2\mathrm{}^2/2}.$$
Using the solution of the equation of motion at $`T=0`$ we find
$`\stackrel{}{m}_{\mathrm{}}(\stackrel{}{x}_0,t)=m_{\mathrm{}}(t)\stackrel{}{\varphi }(\stackrel{}{x}_0,t+\mathrm{}^2/2)\text{with}m_{\mathrm{}}(t)={\displaystyle \frac{Y(t+\mathrm{}^2/2)}{Y(t)}}.`$
The coarse-grained local magnetisation has the same statistical properties as the original field $`\stackrel{}{\varphi }(\stackrel{}{x},t)`$ but with time increased from $`t`$ to $`t+\mathrm{}^2/2`$ and amplitude reduced from $`1`$ to $`m_{\mathrm{}}(t)`$. Namely, the fluctuations of each of its components obey a Gaussian distribution but the amplitude of the vector in $`N`$ space does not fluctuate at all due to the limit $`N\mathrm{}`$. <sup>5</sup><sup>5</sup>5We do not find a non-trivial Gumbel-like distribution $`P(m)`$ of the amplitude of the magnetisation as found for the finite-volume two dimensional XY model in equilibrium in the KT phase . This is again due to the $`N\mathrm{}`$ limit.
The dependence of the amplitude $`m_{\mathrm{}}(t)`$ on time $`t`$ and coarse-graining size $`\mathrm{}`$ is consistent with what one expects for a domain growth system. Firstly, if one fixes the coarse-graining size $`\mathrm{}`$, the amplitude approaches $`1`$ as time $`t`$ is increased. So the system looks “more ordered” at longer time scales. On the other hand, if the time $`t`$ is held fixed while the coarse-graining size $`\mathrm{}`$ is increased, the amplitude of the magnetisation decreases meaning that the system looks “more disordered” at larger length scales.
### 5.2 The distribution of coarse-grained two-time correlations
In the $`O(N)`$ model we can define local coarse-grained “correlations” in the following way:
$$q_{V𝒩}(t,t^{})\frac{1}{𝒩}\underset{\alpha =1}{\overset{𝒩}{}}\frac{1}{V_x}\underset{\stackrel{}{y}V_x}{}\varphi _\alpha (\stackrel{}{y},t)\varphi _\alpha (\stackrel{}{y},t^{}).$$
Strictly speaking this is not a correlation function but rather a composite field. We shall use, however, both names in the following. The first sum is an average over components of the vector $`\stackrel{}{\varphi }`$ in its internal space. Clearly, when $`𝒩=1`$ we test the single component correlation while when $`𝒩=N`$ we sum over all the components of the $`\stackrel{}{\varphi }`$ vector. In the following we shall discuss these two limiting cases and the intermediate cases of finite $`𝒩`$. The second sum is a coarse-graining in real space and it runs over a neighbouring region of the point $`\stackrel{}{x}`$. If $`V_x=1`$ we have a strictly local quantity while for $`V_x=V`$ we recover the global correlation.
### 5.3 Local composite field
Let us start by studying the strictly local composite field
$$q_𝒩q_{V_x=1,𝒩}(t,t^{})=\frac{1}{𝒩}\underset{\alpha =1}{\overset{𝒩}{}}\varphi _\alpha (\stackrel{}{x},t)\varphi _\alpha (\stackrel{}{x},t^{}).$$
The pdf of $`q_𝒩`$ is given by
$`p(q_𝒩)={\displaystyle \frac{1}{Z_0}}{\displaystyle 𝒟\varphi \delta \left(q_𝒩\frac{1}{𝒩}\underset{\alpha =1}{\overset{𝒩}{}}\varphi _\alpha (\stackrel{}{x},t)\varphi _\alpha (\stackrel{}{x},t^{})\right)}`$
$`\times e^{\frac{1}{2\mathrm{\Delta }^2}_\alpha _k\varphi _\alpha ^{}(\stackrel{}{k},0)\varphi _\alpha (\stackrel{}{k},0)}`$
$`={\displaystyle \frac{d\eta }{2\pi }e^{i\eta q_𝒩}\frac{1}{Z_0}𝒟\varphi e^{\frac{i\eta }{𝒩}_\alpha \varphi _\alpha (x,t)\varphi _\alpha (x,t^{})}e^{\frac{1}{2\mathrm{\Delta }^2}_\alpha _k\varphi _\alpha ^{}(\stackrel{}{k},0)\varphi _\alpha (\stackrel{}{k},0)}}`$
$`={\displaystyle \frac{d\eta }{2\pi }e^{i\eta q_𝒩}\frac{Z_\eta }{Z_0}},`$
where
$$Z_\eta D\varphi e^{\frac{1}{2\mathrm{\Delta }^2}_\alpha _{\stackrel{}{k}_1\stackrel{}{k}_2}\varphi _\alpha ^{}(k_1,0)e^{i\stackrel{}{k}_1\stackrel{}{x}}_\eta (\stackrel{}{k}_1,\stackrel{}{k}_2)e^{i\stackrel{}{k}_2\stackrel{}{x}}\varphi _\alpha (\stackrel{}{k}_2,0)},$$
with the symmetric matrix
$$_\eta (\stackrel{}{k}_1,\stackrel{}{k}_2)=\delta _{\stackrel{}{k}_1\stackrel{}{k}_2}+\frac{i\eta \mathrm{\Delta }^2}{𝒩}\left[r(k_1,t,0)r(k_2,t^{},0)+r(k_1,t^{},0)r(k_2,t,0)\right].$$
(89)
Notice that the second term in $`_\eta `$ contains the $`r(k,t_1,t_2)`$ terms for the time evolution of the $`k`$-component. Also notice that $`Z_0=Z_{\eta =0}`$.
The calculation of the function $`Z`$ can be done as follows. It is convenient to rescale the field $`\stackrel{}{\varphi }`$,
$$\stackrel{~}{\stackrel{}{\varphi }}(k)=\mathrm{\Delta }^1e^{i\stackrel{}{k}\stackrel{}{x}}\stackrel{}{\varphi }(\stackrel{}{k},0),$$
in such a way that
$$Z_\eta =𝒟\stackrel{~}{\stackrel{}{\varphi }}e^{\frac{1}{2}_\alpha _{k_1k_2}\stackrel{~}{\stackrel{}{\varphi }}_\alpha ^{}(\stackrel{}{k}_1,0)_\eta (\stackrel{}{k}_1,\stackrel{}{k}_2)\stackrel{~}{\stackrel{}{\varphi }}_\alpha (\stackrel{}{k}_2,0)},$$
up to a trivial (independent of $`\eta `$) multiplicative constant coming from the change of measure. It follows that
$$\frac{Z_\eta }{Z_0}=\left(\frac{det_\eta }{det_0}\right)^{𝒩/2},$$
where we used that the $`\alpha =1,\mathrm{},𝒩`$ components are independent.
In appendix D.1 the eigenmodes of the matrix $`_\eta (\stackrel{}{k}_1,\stackrel{}{k}_2)`$ defined in eq. (89) are obtained; one finds two non-trivial eigenvalues $`\lambda _\pm =1+i\eta [C(t,t^{})\pm 1]`$ and $`2L^d2`$ trivial eigenvalues $`\lambda =1`$. Using these results we have
$$\frac{Z_\eta }{Z_0}=\left\{\left[1+\frac{i\eta }{𝒩}(C(t,t^{})+1)\right]\left[1+\frac{i\eta }{𝒩}(C(t,t^{})1)\right]\right\}^{\frac{𝒩}{2}}$$
(90)
where $`C(t,t^{})`$ is the global correlation function. Thus, the pdf $`p(q_𝒩)`$ is solely parametrised by the value of the the global correlation function.
Lastly, let us note that it is straightforward to generalise the above result to the case of a composite field associated with two different points in space, $`q_𝒩(\stackrel{}{x},\stackrel{}{y})=𝒩^1_{\alpha =1}^𝒩\varphi _\alpha (\stackrel{}{x},t)\varphi _\alpha (\stackrel{}{y},t^{})`$. One finds the same result as the one in eq. (90) but with the global correlation function being replaced by the global two-point function $`C(\stackrel{}{x},\stackrel{}{y};t,t^{})`$.
#### 5.3.1 $`N`$-component averaged local composite field.
When the average over all components of the $`\stackrel{}{\varphi }`$ is considered, i.e. when $`𝒩=N`$, and the $`N\mathrm{}`$ limit is also taken, eq. (90) implies
$$\frac{Z_\eta }{Z_0}e^{i\eta C(t,t^{})}$$
and
$$p(q_N)=\frac{d\eta }{2\pi }e^{i\eta q_N}\frac{Z_\eta }{Z_0}=\frac{d\eta }{2\pi }e^{i\eta q_N}e^{i\eta C(t,t^{})}=\delta [q_NC(t,t^{})].$$
As expected, the average over all internal components of the field erases all fluctuations and the local composite field is forced to take the value of the global correlation function on each site. Note that this is a special feature of the $`N\mathrm{}`$ limit. It is clear that a further coarse-graining on real space will have no effect on the form of the distribution. Thus, the scaling in time is trivially dictated by the global correlation function in this case.
We emphasise that $`p(q_N)`$ computed above is valid in $`N\mathrm{}`$ limit. For some purposes $`O(1/N)`$ corrections of $`p(q_N)`$ can be important. For example one may consider a spin-glass susceptibility-like quantity $`\chi _{SG}=N(q_N^2q_N^2)`$ in an analogous way to the equilibrium one. In equilibrium, the spin-glass susceptibility may be defined as follows. Consider two replicas, say A and B, coupled by an interaction term $`Nϵq`$ in the Hamiltonian where $`q=(1/N)_\alpha \varphi _\alpha ^A\varphi _\alpha ^B`$ is the overlap between the two replicas. The equilibrium spin-glass susceptibility is then defined as $`\chi _{SG}^{eq}=q_{eq}/ϵ=N(q^2_{eq}q_{eq}^2)`$. Note that if $`q^2_{eq}q_{eq}^2=O(1/N)`$, $`\chi _{SG}^{eq}`$ does not vanish in the $`N\mathrm{}`$ limit. We may expect a similar non-trivial result out-of-equilibrium.
#### 5.3.2 One component local composite field
Let us now consider the opposite limit in which we take $`𝒩=1`$ and look at the distribution, $`p(q_\alpha )`$, of the $`x`$-dependent composite field assembled from a single component of $`\stackrel{}{\varphi }`$:
$$q_\alpha \varphi _\alpha (\stackrel{}{x},t)\varphi _\alpha (\stackrel{}{x},t^{}).$$
By setting $`𝒩=1`$ in eq. (90)
$`p(q_\alpha )`$ $`={\displaystyle \frac{d\eta }{2\pi }e^{i\eta q_\alpha }\left[1+i\eta (C+1)\right]^{1/2}\left[1+i\eta (C1)\right]^{1/2}}.`$
The distribution is non-Gaussian and it is a function of times only through the value of the global correlation function $`C=C(t,t^{})`$. In the above integral, the integrand has two branch points: one at $`\eta =i/(1+C)`$ and the other at $`\eta =i/(1C)`$. Performing the integral we obtain
$`p(q_\alpha )`$ $`=`$ $`e^{\frac{|q_\alpha |}{1+\mathrm{sgn}(q_\alpha )C}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dr}{\pi }}{\displaystyle \frac{e^{|q_\alpha |r}}{\sqrt{2r+(1C^2)r^2}}}={\displaystyle \frac{e^{\frac{Cq_\alpha }{1C^2}}}{\pi \sqrt{1C^2}}}K_0\left({\displaystyle \frac{|q_\alpha |}{1C^2}}\right),`$
with $`K_0(x)`$ the modified Bessel function which can be expressed as $`K_0(z)=_1^{\mathrm{}}𝑑xe^{zx}(x^21)^{1/2}`$. This function does not depend on the dimension of space explicitly, it only does through the form of $`C`$. It is sketched in Fig. 1 for four values of the global correlation that are given in the key.
In the special limit $`C1`$ we find $`p(q_\alpha )=e^{q_\alpha /2}/\sqrt{2\pi q_\alpha }`$ for $`q_\alpha >0`$ and $`p(q_\alpha )=0`$ for $`q_\alpha <0`$. In the extreme limit of very separated times in which $`C0`$, $`p(q_\alpha )`$ becomes a symmetric function with respect to $`q_\alpha =0`$ which is not, however, a delta function.
Note that this form is very similar to the result found by Fusco and Zannetti for the equilibrium overlap distribution, $`P(q)`$, of the mean-spherical model at zero temperature .
#### 5.3.3 Finite M-component averaged local composite field
The distribution of finite $`M`$ component averaged local field, $`q_MM^1_\alpha \varphi _\alpha (x,t)\varphi _\alpha (x,t^{})`$, can be studied similarly by setting $`𝒩=M`$ in eq. (90). The integral can be transformed into multiple convolutions of the result for $`M=𝒩=1`$. With the purpose of presenting the result graphically we prefer to perform the integral explicitly. For simplicity we consider only even $`M`$, i.e. $`M=2n`$ with integer $`n=1,2,3,\mathrm{}`$. The integrand has two simple poles at $`\eta =iM/(1+C)`$ and $`\eta =iM/(1C)`$. We then obtain:
$`p(q_M)`$ $`=`$ $`{\displaystyle \frac{n}{4^{n1}}}(1C^2)^{n1}e^{\frac{2n|q_M|}{1+\mathrm{sgn}(q_M)C}}`$
$`\times {\displaystyle \underset{l=0}{\overset{n1}{}}}\left({\displaystyle \frac{4n|q_M|}{1C^2}}\right)^{n1l}{\displaystyle \frac{(n1+l)!}{(n1l)!l!(n1)!}},`$
whose mean is $`C`$ and the variance $`\sigma `$ is given by
$`\sigma ^2`$ $`=`$ $`q_+^2+q_{}^2C^2`$
$`q_\pm ^2`$ $`=`$ $`{\displaystyle \frac{n}{4^{n1}}}\left({\displaystyle \frac{1\pm C}{2n}}\right)^3{\displaystyle \underset{l=0}{\overset{n1}{}}}{\displaystyle \frac{(n1+l)!}{(n1l)!l!(n1)!}}\left({\displaystyle \frac{2}{1C}}\right)^{n1l}\mathrm{\Gamma }(nl+2)`$
where $`\mathrm{\Gamma }(x)`$ is the gamma function. Again, we see that the distribution function is parametrised solely by the global correlation function $`C=C(t,t^{})`$.
Although the mean value of $`p(q_M)`$ is independent of $`M`$ and identical to the global correlation $`C`$, the functional form of this pdf depends strongly on $`M`$. In Fig. 2 we show the functional form for six values of the number of components $`M`$ given in the key and fixed global correlation $`C`$. It can be noticed that for relatively small $`M`$, the position of the peak is different from $`C`$. As $`M`$ increases the position of the peak approaches $`C`$ and the width of the peak shrinks in such a way that the pdf becomes the delta function $`\delta (qC)`$ obtained in Sect. 5.3.1 in the $`M=N\mathrm{}`$ limit.
It is interesting to study the form of these pdfs in more detail. Figure 3 shows $`p(q_M)`$, with $`q_M=M^1_{\alpha =1}^M\varphi _\alpha (\stackrel{}{x},t)\varphi _\alpha (\stackrel{}{x},t^{})`$, against $`x=(q_Mq_M)/\sigma _{q_M}`$. The number of components is $`M=4`$ and different curves correspond to several values of the global correlation given in the key. The plot is in double logarithmic scale. One sees that the curves are positively skewed with the right tail of the distribution being approximately independent of the value of $`C`$ while the left one is not. In Fig. 4 we compare the form of $`p(q_M)`$ to a Gaussian $`e^{x^2}/\sqrt{2\pi }`$ and a Gumbel curve with positive parameter $`a`$ in such a way to make it positively skewed. The normalised Gumbel distribution with mean zero and variance $`1`$ is given by
$`\mathrm{\Phi }_a(x)={\displaystyle \frac{|\alpha |}{\mathrm{\Gamma }(a)}}e^{a\mathrm{log}a}e^{a(\alpha (xx_0))e^{\alpha (xx_0)}},\text{with}`$
$`\alpha =\sqrt{\mathrm{\Psi }^{}(a)}\text{and}\alpha x_0=\mathrm{log}a\mathrm{\Psi }(a),`$ (91)
where $`\mathrm{\Gamma }(x)`$ is the gamma function and $`\mathrm{\Psi }(x)=\mathrm{\Gamma }^{}(x)/\mathrm{\Gamma }(x)`$ is the digamma function. For this intermediate value of $`M`$ the pdf is clearly not Gaussian. The right tail is well fitted with the Gumbel form while the left tail is not.
#### 5.3.4 Effect of coarse-graining the “correlation”
With a similar analysis one shows that coarse-graining has no effect if $`V_x[L(t^{})^d,L(t)^d]`$ while for $`V_x[L(t^{})^d,L(t)^d]`$ the probability distribution function becomes a Gaussian as in the case in which we averaged over all components of the field. Consider again the case $`𝒩=1`$, but some general $`V_x=\mathrm{}^d`$:
$$q_{V_x}q_{V_x,𝒩=1}(\stackrel{}{x};t,t^{})=\frac{1}{V_x}\underset{\stackrel{}{y}V_x}{}\varphi _\alpha (\stackrel{}{y},t)\varphi _\alpha (\stackrel{}{y},t^{}).$$
The pdf can be computed similarly as before but now with
$`_\eta (\stackrel{}{k}_1,\stackrel{}{k}_2)=\delta _{\stackrel{}{k}_1\stackrel{}{k}_2}+{\displaystyle \frac{i\eta \mathrm{\Delta }^2}{V_x}}{\displaystyle \underset{\stackrel{}{y}V_x}{}}[`$ $`r(k_1,t,0)e^{i\stackrel{}{k}_1\stackrel{}{y}}r(k_2,t^{},0)e^{i\stackrel{}{k}_2\stackrel{}{y}}`$
$`+r(k_1,t^{},0)e^{i\stackrel{}{k}_1\stackrel{}{y}}r(k_2,t,0)e^{i\stackrel{}{k}_2\stackrel{}{y}}].`$
The eigenmodes of this matrix are studied in D.2. Diagonalising is not easy for the general case but the following two limiting cases can be considered.
* $`\mathrm{}L(t),L(t^{})`$
The eigenvalues are the same as in the $`V_x=1`$ case: one finds two non-trivial eigenvalues $`\lambda _\pm =1+i\eta [C(t,t^{})\pm 1]`$ and $`2V_x2`$ trivial eigenvalues $`\lambda =1`$. Thus, the pdf is the same as for $`V_x=1`$ (see Fig. 1).
* $`\mathrm{}L(t)L(t^{})`$
For simplicity we consider $`L=L(t)L(t^{})`$. Two non-trivial eigenvalues $`\lambda _\pm 1+i\eta (L/l)^d(C(t,t^{})\pm 1)`$ and $`2V_x2`$ trivial eigenvalues $`\lambda =1`$ are obtained. Note that this is equivalent to the case $`V_x=1`$ if one substitutes
$$𝒩(\mathrm{}/L)^d.$$
Thus, the pdf is the same as those corresponding to composite fields averaged over this number of field components (see Fig. 2-4).
### 5.4 The distribution of coarse-grained linear responses
The distribution of local linear responses is surprisingly trivial in quasi-quadratic systems such as the $`O(N)`$ model. Indeed, the linear response of each thermal run in Fourier space is given by
$$\frac{\delta \varphi _\alpha (\stackrel{}{k},t)}{\delta h_\beta (\stackrel{}{k}^{},t^{})}|_{h=0}=\delta _{\alpha \beta }\delta ^d(\stackrel{}{k}+\stackrel{}{k}^{})r(k;t,t^{})\theta (tt^{}).$$
This implies
$`{\displaystyle \frac{\delta \varphi _\alpha (\stackrel{}{x},t)}{\delta h_\beta (\stackrel{}{x},t^{})}}`$ $`=`$ $`\delta _{\alpha \beta }{\displaystyle \frac{d^dk}{(2\pi )^d}e^{i\stackrel{}{k}\stackrel{}{x}}\frac{d^dk^{}}{(2\pi )^d}e^{i\stackrel{}{k}^{}\stackrel{}{x}}\delta ^d(\stackrel{}{k}+\stackrel{}{k}^{})r(k;t,t^{})\theta (tt^{})}`$
$`=`$ $`\delta _{\alpha \beta }{\displaystyle \frac{d^dk}{(2\pi )^d}r(k;t,t^{})}=R_{\alpha \beta }(t,t^{}),`$
i.e. a uniform result in space that is just equal to the global value. Again, this is independent of the spatial dimension $`d`$.
### 5.5 The joint distribution of local correlations and responses
Using the results above one concludes that the projection of the joint pdf on the $`(C_x,\chi _x)`$ plane at fixed pair of times $`(t,t^{})`$ such that they fall in the ageing regime (this is the extended fdt plot studied in ) is such that there are no fluctuations in the vertical direction while there are in the horizontal one.
### 5.6 Summary
In the strict $`N\mathrm{}`$ limit in which we have not taken into account fluctuations of the constraint $`N^1_\alpha \varphi _\alpha ^2(\stackrel{}{x},t)`$, we found similarities and differences with the distributions of coarse-grained local correlations and responses found in glassy systems. Let us discuss the points enumerated in the introduction to this Section in detail.
i. In all cases the pdfs of local composite fields in the ageing regime depend on times only through the global correlation. This property appears to be independent of there being time-reparametrisation invariance.
ii. The form of the pdfs of one-component composite fields is definitely non-Gaussian but different from the one observed in the $`3d`$ ea model and kinetically constrained particle systems on the lattice . Before coarse-graining in real or internal space the pdf has a maximum at $`q=0`$ for all values of $`C`$ (see Fig. 1). This is simply due to the fact that in the $`O(N)`$ model with $`N\mathrm{}`$ the configurations with many vanishing components are very favourable. <sup>6</sup><sup>6</sup>6The configuration at each point in real space is a vector in $`N`$ dimensions with fixed length. For example, any such vector chosen from a flat distribution typically has a few large components and many \[$`O(N)`$\] components with vanishing value. This can be easily worked out when $`N=2`$, i.e. for the XY model, for which the pdf of the $`x`$ and $`y`$ components are Gaussians centred at zero. To understand the role of the large $`N`$ limit one should compare the above results to, for example, the same pdfs in the XY problem.
Averageing over components or over real space washes out the weight on negative values ($`q<0`$) just as found in the spin models. For finite value of $`𝒩`$ or for coarse-graining boxes that do not go beyond the domain length, the pdfs remain, though, positively skewed for all values of $`C`$ even those corresponding to times that are close to each other (see Fig. 3).
We have also checked whether the distributions of $`q_M`$ can be approximated by a Gumbel-like form with negative parameter. We find that while the tail on the right is quite well described with this functional form, the tail on the left is not (see Fig. 4).
In the large coarse-graining volume, $`\mathrm{}[L(t),L(t^{})]`$, or averaging over a diverging number of components, $`𝒩=N\mathrm{}`$, the pdf becomes a delta function, $`\delta (qC)`$.
iii. There are no fluctuations of the linear responses. This is intimately related to the quasi quadratic nature of the model in the limit $`N\mathrm{}`$. This result is clearly different from what found in glassy models, in which the local response functions do fluctuate form site to site though constrained to follow the global $`\chi (C)`$ curve. In the $`O(N)`$ model the projection of the joint pdf of local correlations and responses also follows the global $`\chi (C)`$ curve but in a trivial manner, since the local responses take a single value.
It would be interesting to study whether these results are modified by $`1/N`$ corrections when the constraint is allowed to fluctuate and yields an additional contribution to the linear response .
## 6 Four point correlation function
A coarsening system is one in which the growing length is easily identified as the typical domain length. A scaling theory then predicts that all correlations should depend on distance and on times only through the value of the typical domain length. This is explicitly realised by the $`O(N)`$ model and an example of such scaling law is given in eq. (34).
In spin-glasses and structural glasses the observation of such a growing length has been elusive. A growing correlation length in the super-cooled liquid has been extracted from the analysis of the connected correlation of fluctuating local composite operators in a number of model systems . The analysis of numerical simulations of several models as well as some experiments indicate that this length takes very small values, of the order of a few nanometres in the super-cooled liquid. A summary of these results appeared recently in .
In an out of equilibrium system, such as the problem at hand, this “four-point” correlation function is naturally defined as
$`C_4(\stackrel{}{x},\stackrel{}{x}^{};t,t^{})`$ $``$ $`[\varphi _\alpha (\stackrel{}{x},t)\varphi _\alpha (\stackrel{}{x},t^{})\varphi _\alpha (\stackrel{}{x}^{},t)\varphi _\alpha (\stackrel{}{x}^{},t^{})]_{ic}`$ (92)
$`[\varphi _\alpha (\stackrel{}{x},t)\varphi _\alpha (\stackrel{}{x},t^{})]_{ic}[\varphi _\alpha (\stackrel{}{x}^{},t)\varphi _\alpha (\stackrel{}{x}^{},t^{})]_{ic}.`$
Note that this quantity is nothing but the connected spatial correlation function of the composite field $`q_\alpha \varphi _\alpha (\stackrel{}{x},t)\varphi _\alpha (\stackrel{}{x},t^{})`$ (see Sect. 5). Since noise and initial condition averaged quantities are expected to be invariant under translations of the space coordinates, this quantity should be equal to
$`C_4(\stackrel{}{r};t,t^{})`$ $``$ $`{\displaystyle \frac{1}{V}}{\displaystyle d^dx[\varphi _\alpha (\stackrel{}{x},t)\varphi _\alpha (\stackrel{}{x},t^{})\varphi _\alpha (\stackrel{}{x}^{},t)\varphi _\alpha (\stackrel{}{x}^{},t^{})]_{ic}}`$
$`{\displaystyle \frac{1}{V^2}}{\displaystyle d^dx[\varphi _\alpha (\stackrel{}{x},t)\varphi _\alpha (\stackrel{}{x},t^{})]_{ic}d^dx^{}[\varphi _\alpha (\stackrel{}{x}^{},t)\varphi _\alpha (\stackrel{}{x}^{},t^{})]_{ic}}`$
$`=`$ $`{\displaystyle \frac{1}{V}}{\displaystyle d^dx[\varphi _\alpha (\stackrel{}{x},t)\varphi _\alpha (\stackrel{}{x},t^{})\varphi _\alpha (\stackrel{}{x}^{},t)\varphi _\alpha (\stackrel{}{x}^{},t^{})]_{ic}}C_{\alpha \alpha }^2(t,t^{})`$
with $`\stackrel{}{r}\stackrel{}{x}\stackrel{}{x}^{}`$. $`C_4(\stackrel{}{x},\stackrel{}{x}^{};t,t^{})`$ measures the probability that similar decorrelations taking place between $`t^{}`$ and $`t`$ occur at a spatial distance $`\stackrel{}{r}`$ in the sample.
The volume integral of $`A`$ defines the quantity
$$\chi _4(t,t^{})d^drC_4(\stackrel{}{r};t,t^{}).$$
that is loosely called a “susceptibility” advocating the use of a fluctuation-dissipation theorem to relate the correlation of composite operators to a linear response. When the operators and the perturbations are composite ones depending on two (or more) times this is however highly non-trivial. Some examples have been exhibited in . In particular, $`C_4`$ is not equal to the response of the composite observable $`[\varphi _\alpha (\stackrel{}{x},t)\varphi _\alpha (\stackrel{}{x},t^{})]_{ic}`$ to an infinitesimal field that couples linearly to $`\varphi _\alpha (\stackrel{}{x}^{},t)\varphi _\alpha (\stackrel{}{x}^{},t^{})`$ (see E) as one would naively propose. In the low-temperature phase, where equilibrium dynamics is lost, the relation between spontaneous and induced fluctuations is still more complicated due to the fact that these are not determined by the equilibrium measure.
With the aim of comparing to the results found in super-cooled liquids we study the behaviour of $`\chi _4`$ during coarsening. The four point correlation function eq. (92) is easily obtained using the solution to the equation of motion, eq. (18). Again, for simplicity we work at $`T=0`$ and we find
$`C_4(\stackrel{}{r};t,t^{})`$ $`=`$ $`[\varphi _\alpha (\stackrel{}{x},t)\varphi _\alpha (\stackrel{}{x}^{},t)]_{ic}[\varphi _\alpha (\stackrel{}{x},t^{})\varphi _\alpha (\stackrel{}{x}^{},t^{})]_{ic}`$ (93)
$`+[\varphi _\alpha (\stackrel{}{x},t)\varphi _\alpha (\stackrel{}{x}^{},t^{})]_{ic}[\varphi _\alpha (\stackrel{}{x}^{},t)\varphi _\alpha (\stackrel{}{x},t^{})]_{ic}`$
$`=`$ $`e^{(r/L(t))^2}e^{(r/L(t^{}))^2}+C^2(t,t^{})e^{2(r/L(t+t^{}))^2}`$
where $`r=|\stackrel{}{x}\stackrel{}{x}^{}|`$, $`C(t,t^{})`$ is the global correlation function and $`L(t)\sqrt{t}`$ is the usual domain size. The first term is a rather trivial contribution since it is just the product of the (average) equal-time spatial correlation functions at $`t`$ and $`t^{}`$. The last term depends on the domain length evaluated at the sum of the two times involved, $`L(t+t^{})`$. Note that if “reciprocity” holds the last term becomes $`[\varphi _\alpha (\stackrel{}{x},t)\varphi _\alpha (\stackrel{}{x}^{},t^{})]_{ic}^2`$. In the ageing regime the length scales $`L(t)`$ and $`L(t^{})`$ are of the same order. Moreover, since $`L(t)t^{1/2}`$, $`L(t+t^{})`$ is also of the same order. Using $`t^{}=\lambda t`$ with $`\lambda [0,1]`$, $`L(t^{})=\lambda ^{1/2}L(t)`$ and $`L(t+t^{})(1+\lambda )^{1/2}L(t)`$. Thus, for distances $`r`$ of the order of $`L(t)`$ all terms contribute. Note that $`C_4(\stackrel{}{r};t,t)`$ does not vanish.
Using eq. (30) for the global correlation in the ageing regime we note that $`C_4(\stackrel{}{r};t,t^{})`$ can be put into the scaling form
$$C_4(\stackrel{}{r};t,t^{})=f_{C_4}(\frac{L(t)}{L(t^{})},\frac{r}{L(t^{})})=\stackrel{~}{f}_{C_4}(\frac{t}{t^{}},\frac{r}{L(t^{})})$$
as expected from simple scaling arguments and found for the one-dimensional Ising chain .
From expression (93) we easily compute $`\chi _4(t,t^{})`$:
$`\chi _4(t,t^{})L^d(t^{})f_{\chi _4}\left({\displaystyle \frac{\tau }{t^{}}}\right)\text{with}f_{\chi _4}(x)=2^{\frac{d}{2}+2^{\frac{d}{2}}}\left({\displaystyle \frac{1+x}{1+x/2}}\right)^{\frac{d}{2}}.`$
This function has the form shown in Fig. 5. It does not have a maximum as a function of $`\tau tt^{}`$ but it monotonically increases towards a finite $`t^{}`$-dependent asymptote. In this respect the behaviour is rather different from what has been found in the supercooled-liquid phase of a number of glassy systems and in the coarsening foam studied in .
It is interesting to analyse the behaviour of the second term too. If one assumes, based on scaling arguments, that
$$C_{ag}(t,t^{})\left(\frac{L(t^{})}{L(t)}\right)^{\overline{\lambda }},$$
i.e. that the very last decay is characterised by the $`\overline{\lambda }`$ exponent , the behaviour of this term at very long time-differences depends on whether $`\overline{\lambda }`$ is larger or equal than $`d/2`$, the lower bound conjectured by Fisher and Huse . When $`\overline{\lambda }=d/2`$, as in the $`O(N)`$ model, this term is also finite and contributes to the asymptotic value of $`\chi _4`$. For other systems in which $`\overline{\lambda }`$ is larger than $`d/2`$ this terms vanishes asymptotically (as implicitly assumed in ).
Let us mention that the alternative definition of $`C_4(t,t^{})`$ proposed in also has a finite asymptotic ($`t^{}\mathrm{}`$) value in the $`O(N)`$ model. This is due to the fact that the last added term is equal to the second term discussed in the previous paragraph and does not vanish.
One could also define a connected spatio-temporal correlation function of the original field $`\varphi _\alpha (\stackrel{}{x},t)`$. Due to the factorisation rules for $`N\mathrm{}`$, these quantities vanish. However, there are $`O(1/N)`$ corrections which yield essential contributions to some related integral susceptibilities .
We can now compare to what has been observed in numerical simulations of the $`3d`$ Edwards-Anderson model using a slightly different expression for the four-point correlation that differs from (92) just in a normalisation. In $`C_4`$ was normalised to be one at $`r=0`$ for all times. For the $`O(N)`$ model this normalisation factor is $`1+C^2(t,t^{})`$. Thus, the space integral of the $`O(N)`$ normalised four-point correlation also approaches a finite limit when $`\tau \mathrm{}`$ and $`t^{}`$ is held fixed.
The normalised four point correlation in the $`3d`$ Edwards-Anderson model was rather well described with the form $`e^{r/\xi (t,t^{})}`$. Even if the space and time dependence in the $`O(N)`$ model is more complicated than a simple exponential decay, the qualitative behaviour of $`\xi (t,t^{})`$ in is similar to that of $`\chi _4(t,t^{})`$ for the $`O(N)`$ model in that $`\xi (t,t^{})`$ increases with both $`t^{}`$ and $`tt^{}`$.
Finally, let us compare the four-point correlation function $`C_4`$ to the integrated response of the composite field, $`\varphi _\alpha (\stackrel{}{x},t)\varphi _\alpha (\stackrel{}{x},t^{})`$ to a composite perturbation $`h_\alpha (\stackrel{}{x}^{},t^{\prime \prime },t^{\prime \prime \prime })`$ . In E we show
$`\left[{\displaystyle \frac{\delta \varphi (\stackrel{}{x},t)\varphi (\stackrel{}{x},t^{})}{\delta h(\stackrel{}{x}^{},t^{\prime \prime },t^{\prime \prime \prime })}}\right]_{ic}=`$
$`R(\stackrel{}{r};t,t^{\prime \prime })C(\stackrel{}{r};t^{},t^{\prime \prime \prime })\theta (tr^{\prime \prime \prime })+R(\stackrel{}{r};t,t^{\prime \prime \prime })C(\stackrel{}{r};t^{},t^{\prime \prime })\theta (tt^{\prime \prime })`$
$`+R(\stackrel{}{r};t^{},t^{\prime \prime })C(\stackrel{}{r};t,t^{\prime \prime \prime })\theta (t^{}t^{\prime \prime \prime })+R(\stackrel{}{r};t^{},t^{\prime \prime \prime })C(\stackrel{}{r};t,t^{\prime \prime })\theta (t^{}t^{\prime \prime })`$
It is clear that this is not simply related to $`C_4`$.
In summary, we found that a model undergoing coarsening has a $`\chi _4(t,t^{})`$ that, as expected, depends on times only through $`L(t)`$ and $`L(t^{})`$, but does not decay to zero at long-time differences. We conclude that the existence of a maximum in $`\chi _4(t,t^{})`$ cannot be taken as evidence for a growing correlation length but other features of this quantity have to be analysed.
## 7 Conclusions
Neither the dynamic equations of the slow (coarsening) contributions to the global correlation and response nor their effective action in the large $`N`$ $`O(N)`$ model are invariant under generic time reparametrisations. This symmetry is reduced to uniform time re-scalings $`t\zeta t`$, with the advanced and retarded scaling dimensions of the global correlation and response, $`\mathrm{\Delta }_C^A=\mathrm{\Delta }_C^R=0`$ and $`\mathrm{\Delta }_R^A=0`$, $`\mathrm{\Delta }_R^R=d/2`$, respectively, and similarly for the corresponding fluctuating fields.
The breakdown of time-reparametrisation invariance seems to be intimately related to the absence of a finite or well-defined effective temperature in $`d>2`$ and $`d=2`$, respectively. Indeed, the retarded scaling dimension $`\mathrm{\Delta }_R^R=d/2`$ in $`d>2`$ implies that the fluctuation-dissipation ratio vanishes asymptotically in the low-temperature phase. Instead, $`\mathrm{\Delta }_R^R=1`$ in $`d=2`$ implies that the fluctuation-dissipation ratio takes a non-trivial $`L(t^{})/L(t)`$ dependent form but the evolution occurs in a single ageing scale in which the correlation itself varies as a function of this ratio. This is inconsistent with the natural requirement of having a single value of the effective temperature per correlation scale.
If we were to use the remaining time-rescaling symmetry to characterise the fluctuations of local correlations of the $`O(N)`$ model when $`N\mathrm{}`$, as we did when we used time-reparametrisation invariance as a guideline to characterise fluctuations in glassy systems , we should introduce the spatial dependence by rescaling time with a space-dependent parameter: $`t\eta _xt`$. However, a simple multiplicative rescaling of time disappears from the correlations:
$$C_x(t,t^{})f_C\left(\frac{\sqrt{\eta _xt^{}}}{\sqrt{\eta _xt}}\right)=f_C\left(\frac{\sqrt{t^{}}}{\sqrt{t}}\right)=f_C\left(\frac{L(t^{})}{L(t)}\right).$$
This means that no such fluctuations are generated. Therefore spatio-temporal fluctuations in the $`O(N)`$ model have a different origin.
We analysed several distribution functions with the aim of identifying similarities and differences with the ones generated by time-reparametrisation invariance. For simplicity we focused on the zero temperature dynamics and we analysed the fluctuations induced by random initial conditions. We concentrated in times such that the dynamics is in the coarsening – ageing – regime. Let us now summarise and discuss our findings.
Each $`\varphi _\alpha (\stackrel{}{x},t)`$, with $`\alpha `$ any component of the $`N`$-dimensional vector $`\stackrel{}{\varphi }`$, obeys a Gaussian pdf. The local one-component composite field, $`\varphi _\alpha (\stackrel{}{x},t)\varphi _\alpha (\stackrel{}{x},t^{})`$, has a non-Gaussian distribution. We derived the functional form of this pdf and we showed how it crosses over to a delta function under coarsening over a sufficient large volume, of linear size larger than the typical domain lengths, $`\mathrm{}L(t^{}),L(t)`$. We also found that the two-time observable made of a sum over a number, $`𝒩`$, of components of the composite field has a similar behaviour to the distribution of the one-component quantity coarse-grained over a volume of linear size $`\mathrm{}^dL^d(t^{})𝒩`$ when $`L(t)`$ and $`L(t^{})`$ are of the same order.
In all these cases, the pdf of local composite fields scales in time just as the global correlation itself; that is to say, it is a function of the ratio between the two characteristic scales $`L(t^{})`$ and $`L(t)`$:
$$p(q_{V_x𝒩};t,t^{})=p(q_{V_x𝒩};C(t,t^{}))=p(q_{V_x𝒩};f_C\left(\frac{L(t^{})}{L(t)}\right)).$$
In - we argued that uniform time-reparametrisation invariance and the simplest choices of effective action for the local reparametrisations, $`h_x(t)`$, imply this kind of scaling and, using numerical simulations, we found it in the $`3d`$ Edwards-Anderson model and a kinetically constrained lattice gas . The solution of the $`O(N)`$ model when $`N\mathrm{}`$ shows that this property is not unique to models with time-reparametrisation invariance.
The form of the pdf of these local two-time quantities is not the Gumbel-like form that we argued should describe the fluctuations of local correlations of spin-like variables that are associated to global time-reparametrisation. In particular, before any coarse-graining – or even after averaging over a small number of components or a small coarse-graining box – the pdf has a peak at very small values of the argument, $`q_{V_x𝒩}0`$. Under further coarse-graining the peak moves towards positive values of $`q_{V_x𝒩}`$ until reaching a Gaussian form centred on the average – global – value $`C`$, that eventually becomes a delta function. Note that the form of the pdfs does not depend on the dimension of space explicitly – it does only through $`C`$.
We may then conjecture that the reason for finding a strong weight at small values of $`q_{V_x𝒩}`$ is the continuous character of the order parameter and its large dimensionality. Indeed, a peak at small values of the two-time composite field should be present in the pdf of $`q_{V_x𝒩}`$ with $`V_x\mathrm{}`$ and $`𝒩<N`$ for all models with a continuous order parameter and a spherical constraint. But this peak should not be necessarily unique. Indeed, preliminary numerical simulations of the dynamics of the $`2d`$ XY model starting from a random initial condition and in the low temperature phase show that the pdf of, say, the horizontal component of the local composite field has a second peak at one, when the two times are not very far away and the global correlation takes a large value. The fate of the two peaks, and thus of the full pdf, under coarse-graining needs to be analysed in more detail but it is not excluded that it may then take a form and evolution similar to the one observed in the $`3d`$ Edwards-Anderson and kinetically constrained lattice gas. Note, however, that the $`2d`$ XY model is critical in the full low temperature phase; its non-equilibrium dynamics is then typical of a critical point with a multiplicative separation of time-scales and an ageing regime that eventually disappears in the long waiting-time limit . While in the ageing regime, this model has the very appealing feature of having a finite integrated response. It belongs to yet another class of models and it is then a very interesting case to study in the context of our discussion.
The pdf of local linear responses is deceptively trivial in the quasi-quadratic large $`N`$ $`O(N)`$ model: these quantities do not fluctuate at all and are just identical to the global value. We do not expect this result to survive in such a trivial manner when including $`1/N`$ corrections or for other coarsening problems that are not (almost) quadratic. In particular, if full time-reparametrisation invariance is broken there is no obvious reason why the joint probability distribution of local responses and correlations should follow the global $`\chi (C)`$ parametric curve between integrated response and correlation. This is a problem that deserves to be addressed analytically and/or numerically in other coarsening models.
We computed the four-point correlation, $`C_4`$, that is usually used to identify a growing correlation length in super-cooled liquids, now during coarsening. Not surprisingly we found that it satisfies a scaling relation in which times enter only through the typical domain length, $`L`$. Contrary to what found in super-cooled liquids and glasses, the integral over space of $`C_4`$ does not vanish at very long time-differences, $`tt^{}\mathrm{}`$ for any fixed $`t^{}`$. The reason for this is the fact that in coarsening systems the spatial correlation, $`C(r,t)`$, does not vanish. The same feature was signalled in in the context of the ferromagnetic Ising chain. We also stressed the fact that this quantity is not trivially related to a susceptibility (see for a similar discussion).
It is interesting to compare the pure time transformations studied in this paper to the common space-time invariance of domain growth . The exact solution of the $`O(N)`$ model is invariant, in the long times and large scales limit, under simultaneously rescaling of time and space \[see eqs. (34)\], and the slow part of the dynamic action is invariant under the related renormalisation group transformation \[see Sect. 4.4\]. However, it is not this space-time invariance that is the relevant symmetry if one is interested in fluctuations within a given domain. One should consider separations $`r`$ that are held fixed while the long-time limit is taken. More precisely, one should consider fixed ratios $`t/t^{}`$ while $`L(t)\mathrm{}`$, and thus $`r/L(t)0`$. It is in this limit that reparametrisation invariance should be investigated, and there are a number of issues that one must consider specifically in the case of the $`O(N)`$ model. First, reparametrisation invariance cannot be an exact symmetry of the solution to the $`O(N)`$ – or any other similar dynamic problem – since a particular function $`h(t)`$ is bound to be chosen by the evolution. It may only arise as an approximate invariance in the asymptotic limit in which the non-invariant terms – that act as “pinning fields” and fix the time scaling $`h(t)`$ – become less and less important. This is indeed what happens in mean-field disordered models of the $`p`$ spin type and, we argued , in the $`3d`$ Edwards-Anderson spin-glass. Second, we showed in this paper that reparametrisation invariance does not develop in the $`O(N)`$ model, only a smaller symmetry, simple time scale invariance, does. We arrived at these results by studying the equations of motion for the global $`C(t,t^{})`$ and $`R(t,t^{})`$ and the action for the slow flucting fields. Similar results would be obtained for the two space-point correlation $`C(r;t,t^{})`$ and response $`R(r;t,t^{})`$, when $`r`$ is held fixed while the long-time limit is taken.
Let us finally stress the main issue arising from this study, i.e. the conjecture that an extreme violation of the fluctuation-dissipation theorem is intimately related to the breakdown of time-reparametrisation invariance at long times in general. If this is correct, systems with a finite or an asymptotically infinite ’effective temperature’ belong to different ‘universality’ classes, as non-equilibrium fluctuations are concerned . It would be interesting to put this conjecture to the test in other solvable models. In particular, by comparing to similar fluctuations in the XY model one should be able to identify the peculiar features due to the $`N\mathrm{}`$ limit. The special $`d=2`$ case should be particularly interesting. Another route is to analyse other coarsening systems with a discrete order parameter: one then should be able to disentangle the features that are due to $`X0`$ from those that are due to the continuous character of the field.
Acknowledgements
We thank G. Biroli and M. Picco for very useful discussions. L.F.C. is a member of the Institut Universitaire de France. This research was supported in part by NSF grants DMR-0305482, DMR-0403997, and INT-0128922 (C.C.), an NSF-CNRS collaboration, the ACI-France “Algorithmes d’optimisation et systèmes desordonnés quantiques”, the STIPCO European Community Network, and NSF Grant No. PHY99-07949 (L.F.C.). H. Yoshino acknowledges financial support from the Japanese Society of Promotion of Science and CNRS.
## Appendix A The equation of motion for the global linear response
In this appendix we show how to obtain the equations of motion for the correlation and response, expressed in terms of $`R(t,t^{})`$ and $`C(t,t^{})`$ themselves, starting from the exact expressions for $`R(t,t^{})`$ and $`C(t,t^{})`$ obtained from the equations of motion for the field $`\stackrel{}{\varphi }(t)`$.
The exact solution for the correlation and response follows from the self-consistent solution of eqs. (18) and (19):
$`C(t,t^{})`$ $`=`$ $`Y^1(t)Y^1(t^{})[\mathrm{\Delta }^2e^{ϵ_k(t+t^{})}`$
$`+2T{\displaystyle _0^{\mathrm{min}(t,t^{})}}dt^{\prime \prime }e^{ϵ_k(t+t^{}2t^{\prime \prime })}Y^2(t^{\prime \prime })]`$
$`R(t,t^{})`$ $`=`$ $`Y^1(t)Y(t^{})e^{ϵ_k(tt^{})}\theta (tt^{}),`$ (95)
where $`ϵ_k`$ is the dispersion and $`f(k)=\frac{d^dk}{(2\pi )^d}f(k)`$.
For $`t>t^{}`$ we can write
$`{\displaystyle \frac{C(t,t^{})}{t}}`$ $`=`$ $`z(t)C(t,t^{})+Y^1(t)Y^1(t^{})[\mathrm{\Delta }^2ϵ_ke^{ϵ_k(t+t^{})}`$
$`+2T{\displaystyle _0^{\mathrm{min}(t,t^{})}}dt^{\prime \prime }ϵ_ke^{ϵ_k(t+t^{}2t^{\prime \prime })}Y^2(t^{\prime \prime })]`$
$`{\displaystyle \frac{R(t,t^{})}{t}}`$ $`=`$ $`z(t)R(t,t^{})+Y^1(t)Y(t^{})ϵ_ke^{ϵ_k(tt^{})}\theta (tt^{}),`$ (97)
where $`z(t)=\frac{}{t}\mathrm{ln}Y(t)`$.
In order to express eqs. (A) and (97) in terms of $`R`$’s and $`C`$’s, all one needs to do is express the right hand side of these equations in terms of convolutions of $`R`$’s and $`C`$’s. The first step to do so is to express the function
$$D(t)=ϵ_ke^{ϵ_kt}\theta (t)=_0^{\mathrm{}}𝑑ϵg(ϵ)(ϵ)e^{ϵt}\theta (t)$$
in terms of convolutions of the function
$`G(t)=e^{ϵ_kt}\theta (t)`$ $`={\displaystyle _0^{\mathrm{}}}𝑑ϵg(ϵ)e^{ϵt}\theta (t)={\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{d\omega }{2\pi }}e^{i\omega t}\stackrel{~}{G}(\omega ),`$
where
$$\stackrel{~}{G}(\omega )_0^{\mathrm{}}𝑑ϵ\frac{g(ϵ)}{ϵi\omega }$$
(98)
and $`g(ϵ)`$ is the density of states with $`ϵ=ϵ_k`$. In other words, we basically need to cast
$$D(t)=\underset{n=1}{\overset{\mathrm{}}{}}A_{n1}\underset{n}{\underset{}{GG\mathrm{}G}}(t).$$
(99)
We start by writing
$`\underset{n}{\underset{}{GG\mathrm{}G}}(t)`$ $`={\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\tau _1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\tau _2\mathrm{}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\tau _{n1}G(t\tau _1)G(\tau _1\tau _2)`$
$`sG(\tau _{n2}\tau _{n1})G(\tau _{n1})`$
$`={\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{d\omega }{2\pi }}e^{i\omega t}\left[\stackrel{~}{G}(\omega )\right]^n`$
$`={\displaystyle _0^{\mathrm{}}}𝑑ϵ_1s{\displaystyle _0^{\mathrm{}}}𝑑ϵ_n{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{d\omega }{2\pi }}e^{i\omega t}{\displaystyle \underset{a=1}{\overset{n}{}}}{\displaystyle \frac{g(ϵ_a)}{ϵ_ai\omega }}`$
$`=n{\displaystyle _0^{\mathrm{}}}𝑑ϵg(ϵ)e^{ϵt}\left[{\displaystyle _0^{\mathrm{}}}𝑑ϵ^{}{\displaystyle \frac{g(ϵ^{})}{ϵ^{}ϵ}}\right]^{n1}\theta (t).`$
Thus, in short, we have
$$\underset{n}{\underset{}{GG\mathrm{}G}}(t)=n_0^{\mathrm{}}𝑑ϵg(ϵ)e^{ϵt}\left[h(ϵ)\right]^{n1}\theta (t),$$
where the function $`h(ϵ)`$ is defined as
$$h(ϵ)=_0^{\mathrm{}}𝑑ϵ^{}\frac{g(ϵ^{})}{ϵ^{}ϵ}.$$
(100)
Next, let us expand $`ϵ`$ as a function of $`h(ϵ)`$:
$$ϵ=h^1h(ϵ)=\underset{n=0}{\overset{\mathrm{}}{}}a_n[h(ϵ)]^n,\mathrm{with}a_n=\frac{1}{n!}\frac{d^nh^1}{dz^n}|_{z=0}.$$
Therefore, we can write
$`D(t)`$ $`=\theta (t){\displaystyle _0^{\mathrm{}}}𝑑ϵg(ϵ)e^{ϵt}(ϵ)`$
$`=\theta (t){\displaystyle _0^{\mathrm{}}}𝑑ϵg(ϵ)e^{ϵt}\left({\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}a_{n1}[h(ϵ)]^{n1}\right)`$
$`={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{a_{n1}}{n}}\underset{n}{\underset{}{GG\mathrm{}G}}(t),`$
which is exactly eq. (99), with $`A_n=a_n/(n+1)`$.
Now that we have the expression for $`D(t)`$, let us show how one can write, for example, an integral-differential equation for $`R(t,t^{})`$ \[eq. (97)\]. First, notice that from eq. (95)
$$G(tt^{})=\frac{Y(t)}{Y(t^{})}R(t,t^{}).$$
Hence,
$`\underset{n}{\underset{}{GG\mathrm{}G}}(tt^{})`$ $`={\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\tau _1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\tau _2\mathrm{}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\tau _{n1}{\displaystyle \frac{Y(t)}{Y(\tau _1)}}R(t,\tau _1)`$
$`s{\displaystyle \frac{Y(\tau _{n1})}{Y(t^{})}}R(\tau _{n1},t^{})`$
$`={\displaystyle \frac{Y(t)}{Y(t^{})}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\tau _1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\tau _2\mathrm{}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\tau _{n1}R(t,\tau _1)`$
$`sR(\tau _{n1},t^{})`$
$`={\displaystyle \frac{Y(t)}{Y(t^{})}}\underset{n}{\underset{}{RR\mathrm{}R}}(t,t^{}),`$
which allows us to write the last term in eq. (97) as
$`{\displaystyle \frac{Y(t^{})}{Y(t)}}D(tt^{})`$ $`={\displaystyle \frac{Y(t^{})}{Y(t)}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}A_{n1}\underset{n}{\underset{}{GG\mathrm{}G}}(tt^{})`$
$`={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}A_{n1}\underset{n}{\underset{}{RR\mathrm{}R}}(t,t^{}).`$
Thus finally we have
$$\frac{R(t,t^{})}{t}=z(t)R(t,t^{})+\underset{n=0}{\overset{\mathrm{}}{}}A_n\underset{n+1}{\underset{}{RR\mathrm{}R}}(t,t^{}).$$
(101)
Lastly, let us note that the above equations can be easily extended to the describe the evolution of the two-time two-point correlation function $`C(r;t,t^{})C(\stackrel{}{x},\stackrel{}{y};t,t^{})`$ and response function $`R(r;t,t^{})R(\stackrel{}{x},\stackrel{}{y};t,t^{})`$, with $`r=|\stackrel{}{x}\stackrel{}{y}|>0`$. One can easily verify that the generalisation can be done by formally replacing the density of states $`g(ϵ)`$ by
$$g(ϵ;r)cg(ϵ)_1^1𝑑y(1y^2)^{(d3)/2}\mathrm{cos}(\sqrt{ϵ}ry)$$
(102)
in $`d3`$. Here $`c^1=_1^1𝑑y(1y^2)^{(d3)/2}=2^{(d2)}B((d1)/2,(d1)/2)`$ is the normalization constant. In $`d=1`$ and $`2`$, one simply has to use $`g(ϵ)\mathrm{cos}(\sqrt{ϵ}r)`$ and $`g(ϵ)_0^\pi 𝑑\theta \mathrm{cos}(\sqrt{ϵ}r\mathrm{cos}(\theta ))/\pi `$, respectively. The closed set of equations of motion for $`C(r;t,t^{})`$ and $`R(r;t,t^{})`$ are series expansions with coefficients $`A_n(r)`$ which now depend on the distance $`r`$ explicitly.
## Appendix B The ageing limit of the equations of motion
To obtain the equations of motion for the response in the ageing limit, one substitutes in eq. (61) \[or eq. (101)\]
$$R(t,t^{})=R_{st}(tt^{})+R_{ag}(t,t^{})$$
and use that the stationary response decays to zero in time scales in which the ageing component remains roughly constant. (See for a detailed explanation of this separation). For example, in the term
$`A_1{\displaystyle 𝑑t^{\prime \prime }R(t,t^{\prime \prime })R(t^{\prime \prime },t^{})}`$
$`A_1[{\displaystyle _t^{}^t}dt^{\prime \prime }R_{ag}(t,t^{\prime \prime })R_{ag}(t^{\prime \prime },t^{})+{\displaystyle _t^{}^{t^+}}dt^{\prime \prime }R_{ag}(t,t^{\prime \prime })R_{st}(t^{\prime \prime }t^{})`$
$`+{\displaystyle _t^{}^t}dt^{\prime \prime }R_{st}(tt^{\prime \prime })R_{ag}(t^{\prime \prime },t^{})]`$
$`A_1\left[{\displaystyle _t^{}^t}𝑑t^{\prime \prime }R_{ag}(t,t^{\prime \prime })R_{ag}(t^{\prime \prime },t^{})+2\chi _{st}R_{ag}(t,t^{})\right]`$
with
$$\chi _{st}=_t^{}^{t^+}𝑑t^{\prime \prime }R_{st}(t^{\prime \prime }t^{})=_t^{}^t𝑑t^{\prime \prime }R_{st}(tt^{\prime \prime }).$$
Notice that if we start from a term with one time integral (the term with coefficient $`A_1`$), then we collect in the ageing regime, in addition to the term with one time integral, a term with no time integrals. Similarly, starting from a term with $`n`$ integrals (the term with coefficient $`A_n`$), we would generate in the ageing limit terms with $`n,n1,\mathrm{},0`$ integrals. We can collect all these terms into a new series
$$\underset{n=0}{\overset{\mathrm{}}{}}\stackrel{~}{A}_n𝑑t_n𝑑t_{n1}\mathrm{}𝑑t_1R_{ag}(t,t_1)R_{ag}(t_1,t_2)\mathrm{}R_{ag}(t_n,t^{}),$$
where the coefficients $`\stackrel{~}{A}_n`$ are related to the original $`A_n`$ by a simple combinatorial argument, that goes as follows. Terms with $`n`$ time integrals and $`n+1`$ $`R_{ag}`$’s are obtained starting with terms with $`pn`$ integrals and $`p+1`$ $`R`$’s, where $`pn`$ of the $`R`$’s are replaced by $`R_{st}`$ and the remaining $`n+1`$ $`R`$’s are replaced by $`R_{ag}`$’s. This allows us to write
$`\stackrel{~}{A}_n`$ $`={\displaystyle \underset{p=n}{\overset{\mathrm{}}{}}}A_p\left(\begin{array}{cc}p+1\hfill & \\ pn\hfill & \end{array}\right)\chi _{st}^{pn}`$
$`={\displaystyle \frac{1}{(n+1)!}}{\displaystyle \underset{p=n}{\overset{\mathrm{}}{}}}A_p(p+1)p(p1)\mathrm{}(pn+1)\chi _{st}^{pn}`$
$`={\displaystyle \frac{1}{(n+1)!}}\left({\displaystyle \frac{d}{d\chi _{st}}}\right)^n{\displaystyle \underset{p=n}{\overset{\mathrm{}}{}}}A_p(p+1)\chi _{st}^p`$
$`={\displaystyle \frac{1}{(n+1)!}}\left({\displaystyle \frac{d}{d\chi _{st}}}\right)^n{\displaystyle \underset{p=0}{\overset{\mathrm{}}{}}}A_p(p+1)\chi _{st}^p`$
Now, from A, $`a_p=A_p(p+1)`$ are the coefficients of the series expansion of the function $`ϵ(h)`$. Therefore we can simply write
$$\stackrel{~}{A}_n=\frac{1}{(n+1)!}\left(\frac{d}{d\chi _{st}}\right)^nϵ(\chi _{st}).$$
(104)
## Appendix C The spherical spin-glass with Gaussian interactions
The spherical spin-glass model with Gaussian distributed two-body interactions has been studied in a series of papers , -. It was there shown that the asymptotic solution in the ageing regime scales as
$$R_{ag}(t,t^{})t^{3/2}f_R\left(\lambda \right),C_{ag}(t,t^{})f_C\left(\lambda \right),$$
(105)
and $`0\lambda t^{}/t1`$. Here, we look at this problem from a different angle, motivated by the generic discussion presented in Sect. 4.2.2. Let us analyse each term in the equations for the global response and correlation by evaluating them in the ageing regime using the scaling forms in (105). The equation for the global response reads
$`{\displaystyle \frac{R(t,t^{})}{t}}`$ $`=`$ $`z(t)R(t,t^{})+{\displaystyle _t^{}^t}𝑑t^{\prime \prime }R(t,t^{\prime \prime })R(t^{\prime \prime },t^{}),`$
with the Lagrange multiplier $`z(t)`$ being fixed by the condition $`C(t,t)=1`$ that yields:
$$z(t)=T+2_0^t𝑑t^{}C(t,t^{})R(t,t^{}).$$
(106)
In the aging regime, the left-hand-side scales as
$$\left[\frac{3}{2}f_R(\lambda )+f_R^{}(\lambda )\right]t^{5/2}.$$
(107)
The asymptotic scaling of the Lagrange multiplier is know from the exact solution to be
$$z(t)2+ct^1$$
(108)
with $`c`$ a numerical coefficient. Let us derive this result from eq. (106) using the forms in (105). If, proceeding as usual, we separate the integral in (106) into a stationary and an aging part and we keep the leading contributions to each of these, we find
$$\underset{t\mathrm{}}{lim}z(t)=z_{\mathrm{}}+\alpha t^{1/2}T+\frac{1}{T}(1q_{ea}^2)+t^{1/2}_0^1𝑑\lambda ^{}f_R(\lambda ^{})f_C(\lambda ^{}).$$
If one uses the relation between $`q_{ea}`$ and $`T`$, the time independent term is consistent with $`z_{\mathrm{}}=2`$. However, the approach to the asymptotic value is incorrect. The mistake we have done is that we neglected the correction to the constant value of the stationary contribution that cancels the leading aging one, and we neglected the correction to the leading aging contribution that yields the correct $`t^1`$ decay.
The easiest and most general way of deriving the result above is to go back to the general representation of the solution for $`C`$ and $`R`$ and plug these into the integral term in (106). After some algebra, and working at $`T=0`$ for simplicity, one finds
$$_0^t𝑑t^{}C(t,t^{})R(t,t^{})=Y^2(t)𝑑ϵe^{2ϵt}g(ϵ)h(ϵ)$$
(109)
where $`g(ϵ)`$ is a generic density of states and $`h(ϵ)`$ is the function defined in eq. (100). Now, a density of states with a finite support in $`[0,1]`$ and power law decays on the two ends can be mimicked by the form
$$g(ϵ)ϵ^\nu (1ϵ)^{1\nu }$$
(110)
that allows us to do the calculations explicitly. In particular, the semicircle case is mimicked by $`\nu =1/2`$. Close to $`ϵ0`$ the function $`h(ϵ)`$ then reads
$$h(ϵ)\frac{\pi }{\mathrm{sin}\pi \nu }\left[(1\nu )ϵ\mathrm{cos}(\pi \nu )ϵ^\nu +\mathrm{}\right].$$
(111)
Replacing in (109) and using the asymptotic form of $`Y(t)`$ one has
$$z(t)a(1\nu )\alpha \mathrm{cos}\pi \nu t^\nu ct^1+\mathrm{}$$
(112)
Thus, for the special case $`\nu =1/2`$ the prefactor of the $`t^\nu `$ term vanishes and one recovers the correct behaviour in $`t^1`$. A similar phenomenon occurs in the integral over the two responses. The stationary contributions yields a term that is $`O(t^{3/2})`$ and its cancellation with the constant asymptotic value of $`z_{\mathrm{}}`$ fixes the Edwards-Anderson order parameter as a function of temperature:
$$T+\frac{1}{T}(1q_{ea}^2)=\frac{2(1q_{ea})}{T}$$
(113)
that is equivalent to $`T^2=(1q_{ea})^2q_{ea}=1T`$ ($`TT_c=1`$). The next-to-leading order terms are $`O(t^2)`$ but their prefactor vanishes. Finally, one is left with a term that is $`O(t^{5/2})`$, just as the time-derivative and the another term left from $`z(t)R_{ag}(t,t^{})`$. This non-trivial equation fixes the functions $`f_C`$ and $`f_R`$.
The analysis of the equation for $`C`$ is similar. The leading terms are $`O(1)`$; their cancellation leads to an equation identical to (113). The next-to-leading order terms are $`O(t^{1/2})`$ but their overall prefactor vanishes. The time-derivative term is $`O(t^1)`$ and it combines with the remaining terms to yield a non-trivial equation.
Note that in the analysis above we used the correct asymptotic behaviour of $`R`$ and $`C`$ in the ageing regime, that we know from the direct solution to the (linear set of) Langevin equations. If $`p3`$ one cannot solve the dynamics exactly and one is forced to do an asymptotic analysis of the equations for $`R`$ and $`C`$ assuming a decay of the linear response and searching for a consistent solution. When $`p3`$ one proposes $`R_{ag}(t,t^{})t^1f_R(\lambda )`$ and $`C_{ag}(t,t^{})f_C(\lambda )`$. In this case, the stationary and ageing contributions to the Lagrange multiplier are both finite. Moreover, all terms in the right-hand-side of the equations for $`R`$ and $`C`$ are of the same order, $`O(t^1)`$ and $`O(1)`$, respectively, while the time derivatives are much smaller, $`O(t^2)`$ and $`O(t^1)`$, respectively. Dropping the time-derivatives one finds a solution that is consistent this the scaling assumption. In the $`p=2`$ one could have proposed a similar (wrong) scaling and look for its consequences. It is interesting to notice that if one naively pursues this calculation one finds $`X=0`$ as the unique possible asymptotic solution \[see eq. (37) for the definition of $`X`$\] which is consistent with the exact result, $`Xt^{1/2}`$, in the $`t\mathrm{}`$ limit.
## Appendix D Diagonalising the matrix $`_\eta (\stackrel{}{k}_1,\stackrel{}{k}_2)`$
In this appendix we study the eigenmodes the matrix $`_\eta (\stackrel{}{k}_1,\stackrel{}{k}_2)`$ defined in eq. (89) for the case of $`V_x=1`$ (without coarse-graining) and eq. (5.3.4) with finite coarse-graining volume $`V_x`$.
### D.1 Case $`V_x=1`$ (without coarse-graining)
First we study the the matrix $`_\eta (\stackrel{}{k}_1,\stackrel{}{k}_2)`$ defined in eq. (89) for the case of $`V_x=1`$. We show that two and only two eigenvalues of $`_\eta `$ depend on $`\eta `$ and all the others are fixed to one. For convenience, let $`_k(t)r(k,t,0)`$ and $`_k(t^{})=r(k,t^{},0)`$ label the $`k`$-indexed row of column vectors $`𝐑(t)`$ and $`𝐑(t^{})`$ (these vectors live in $`𝒟=L^d`$ dimensions). Note that the length of this vector is constant in time $`|𝐑(t)|^2=R^2`$ (which ensures conservation of the length of the $`N`$-component vector field in the $`N\mathrm{}`$ limit). Let $`𝐯^\lambda `$ be, within the same notation, an eigenvector of $`_\eta `$ with eigenvalue $`\lambda `$. Then
$`\lambda v_{\stackrel{}{k}_1}^\lambda ={\displaystyle \underset{\stackrel{}{k}_2}{}}_\eta (\stackrel{}{k}_1,\stackrel{}{k}_2)v_{\stackrel{}{k}_2}^\lambda `$
$`={\displaystyle \underset{\stackrel{}{k}_2}{}}\delta _{\stackrel{}{k}_1\stackrel{}{k}_2}v_{\stackrel{}{k}_2}^\lambda +i\eta {\displaystyle \frac{\mathrm{\Delta }^2}{𝒩}}\left[_{\stackrel{}{k}_1}(t){\displaystyle \underset{\stackrel{}{k}_2}{}}_{\stackrel{}{k}_2}(t^{})v_{\stackrel{}{k}_2}^\lambda +_{\stackrel{}{k}_1}(t^{}){\displaystyle \underset{\stackrel{}{k}_2}{}}_{\stackrel{}{k}_2}(t)v_{\stackrel{}{k}_2}^\lambda \right]`$
$`=v_{\stackrel{}{k}_1}^\lambda +i\eta {\displaystyle \frac{\mathrm{\Delta }^2}{𝒩}}\left[(𝐑(t^{})𝐯^\lambda )_{\stackrel{}{k}_1}(t)+(𝐑(t)𝐯^\lambda )_{\stackrel{}{k}_1}(t^{})\right].`$
This equation is equivalent to
$$(\lambda 1)𝐯^\lambda =i\eta \frac{\mathrm{\Delta }^2}{𝒩}\left[(𝐑(t^{})𝐯^\lambda )𝐑(t)+(𝐑(t)𝐯^\lambda )𝐑(t^{})\right]$$
(114)
and has $`𝒟`$ solutions.
Only two eigenvalues of $`_\eta `$ are changed by the presence of the $`\eta `$ term. One of them is
$$𝐯^\lambda 𝐑(t)+𝐑(t^{})\mathrm{with}\lambda _+=1+\frac{i\eta }{𝒩}\left(C(t,t^{})+1\right).$$
the other is
$$𝐯^\lambda 𝐑(t)𝐑(t^{})\mathrm{with}\lambda _{}=1+\frac{i\eta }{𝒩}\left(C(t,t^{})1\right).$$
Here we used that expression for the global correlation at $`T=0`$,
$$C(t,t^{})=\mathrm{\Delta }^2𝐑(t^{})𝐑(t)=\mathrm{\Delta }^2\underset{k}{}r(k,t^{},0)r(k,t,0),$$
and
$$\mathrm{\Delta }^2R^2=1.$$
The other $`𝒟2`$ solutions are such that
$$𝐯^\lambda 2\mathrm{d}\text{plane spanned by the above two eigenvectors with}\lambda =1.$$
### D.2 Case of finite $`V_x`$
Next we study the the matrix $`_\eta (\stackrel{}{k}_1,\stackrel{}{k}_2)`$ defined in eq. (5.3.4) for the case of finite coarse-graining volume $`V_x`$. Let $`_k(t,\stackrel{}{y})r(k,t,0)e^{i\stackrel{}{k}\stackrel{}{y}}`$ label the $`k`$-indexed row of the column vector $`𝐑(t,\stackrel{}{y})`$. This allows us to write an eigenvalue equation for $`_\eta `$, similarly to what we have done above for the case $`l=1`$,
$`(\lambda 1)𝐯^\lambda =i\eta {\displaystyle \frac{\mathrm{\Delta }^2}{V_x}}{\displaystyle \underset{\stackrel{}{y}V_x}{}}[`$ $`(𝐑(t^{},\stackrel{}{y})𝐯^\lambda )𝐑(t,\stackrel{}{y})`$ (115)
$`+(𝐑(t,\stackrel{}{y})𝐯^\lambda )𝐑(t^{},\stackrel{}{y})],`$
where the inner (dot) product is here defined as $`\mathrm{𝐚𝐛}=_ka_k^{}b_k`$.
This eigenvalue equation has $`𝒟2V_x`$ trivial solutions with $`\lambda =1`$. The eigenvectors for such solutions satisfy $`𝐑(t^{},\stackrel{}{y})𝐯^\lambda =0`$ and $`𝐑(t,\stackrel{}{y})𝐯^\lambda =0`$, for $`\stackrel{}{y}V_x`$, and hence span the orthogonal subspace to that spanned by the $`2V_x`$ vectors $`𝐑(t,\stackrel{}{y})`$ and $`𝐑(t^{},\stackrel{}{y})`$ ($`\stackrel{}{y}V_x`$).
The remaining (non-trivial) eigenvectors can be written as
$$𝐯^\lambda =\underset{\stackrel{}{y}V_x}{}\alpha ^\lambda (\stackrel{}{y})𝐑(t,\stackrel{}{y})+\beta ^\lambda (\stackrel{}{y})𝐑(t^{},\stackrel{}{y})$$
for some $`2V_x`$ expansion coefficients $`\alpha ^\lambda (\stackrel{}{y})`$ and $`\beta ^\lambda (\stackrel{}{y})`$ for $`\stackrel{}{y}V_x`$. Plugging this into eq. (115) leads to
$`(\lambda 1)\alpha ^\lambda (\stackrel{}{y})=i\eta {\displaystyle \frac{\mathrm{\Delta }^2}{V_x}}{\displaystyle \underset{\stackrel{}{y^{}}V_x}{}}[`$ $`\alpha ^\lambda (\stackrel{}{y^{}})(𝐑(t^{},\stackrel{}{y})𝐑(t,\stackrel{}{y^{}}))`$ (116)
$`+\beta ^\lambda (\stackrel{}{y^{}})(𝐑(t^{},\stackrel{}{y})𝐑(t^{},\stackrel{}{y^{}}))],`$
$`(\lambda 1)\beta ^\lambda (\stackrel{}{y})=i\eta {\displaystyle \frac{\mathrm{\Delta }^2}{V_x}}{\displaystyle \underset{\stackrel{}{y^{}}V_x}{}}[`$ $`\alpha ^\lambda (\stackrel{}{y^{}})(𝐑(t,\stackrel{}{y})𝐑(t,\stackrel{}{y^{}}))`$ (117)
$`+\beta ^\lambda (\stackrel{}{y^{}})(𝐑(t,\stackrel{}{y})𝐑(t^{},\stackrel{}{y^{}}))].`$
Using that
$$\mathrm{\Delta }^2𝐑(t,\stackrel{}{y})𝐑(t^{},\stackrel{}{y^{}})=C(\stackrel{}{y},\stackrel{}{y^{}};t,t^{})=C(t,t^{})\mathrm{exp}\left[\frac{|\stackrel{}{y}\stackrel{}{y^{}}|^2}{L^2(t)+L^2(t^{})}\right],$$
with the length scales $`L(t)=2\sqrt{t}`$ and $`L(t^{})=2\sqrt{t^{}}`$, and substituting in eqs. (116) and (117), one obtains
$`(\lambda 1)\alpha ^\lambda (\stackrel{}{y})=i\eta {\displaystyle \frac{1}{V_x}}{\displaystyle \underset{\stackrel{}{y^{}}V_x}{}}`$ $`\{\alpha ^\lambda (\stackrel{}{y^{}})C(t,t^{})\mathrm{exp}\left[{\displaystyle \frac{|\stackrel{}{y}\stackrel{}{y^{}}|^2}{L^2(t)+L^2(t^{})}}\right]`$ (118)
$`+\beta ^\lambda (\stackrel{}{y^{}})\mathrm{exp}\left[{\displaystyle \frac{|\stackrel{}{y}\stackrel{}{y^{}}|^2}{2L^2(t^{})}}\right]\},`$
$`(\lambda 1)\beta ^\lambda (\stackrel{}{y})=i\eta {\displaystyle \frac{1}{V_x}}{\displaystyle \underset{\stackrel{}{y^{}}V_x}{}}`$ $`\{\alpha ^\lambda (\stackrel{}{y^{}})\mathrm{exp}\left[{\displaystyle \frac{|\stackrel{}{y}\stackrel{}{y^{}}|^2}{2L^2(t)}}\right]`$ (119)
$`+\beta ^\lambda (\stackrel{}{y^{}})C(t,t^{})\mathrm{exp}\left[{\displaystyle \frac{|\stackrel{}{y}\stackrel{}{y^{}}|^2}{L^2(t)+L^2(t^{})}}\right]\}.`$
These equations are difficult to solve for generic ratios of the length scales $`L(t)`$ and $`L(t^{})`$ to the coarse-graining box size $`\mathrm{}`$. In the following we consider some limiting cases.
#### D.2.1 Case $`\mathrm{}L(t),L(t^{})`$
One simple situation is given by $`\mathrm{}L(t),L(t^{})`$, in which case $`|\stackrel{}{y}\stackrel{}{y^{}}|L(t),L(t^{})`$ and eqs. (118) and (119) simplify to
$$(\lambda 1)\alpha ^\lambda (\stackrel{}{y})=i\eta \frac{1}{V_x}\underset{\stackrel{}{y^{}}V_x}{}\left[C(t,t^{})\alpha ^\lambda (\stackrel{}{y^{}})+\beta ^\lambda (\stackrel{}{y^{}})\right],$$
(120)
$$(\lambda 1)\beta ^\lambda (\stackrel{}{y})=i\eta \frac{1}{V_x}\underset{\stackrel{}{y^{}}V_x}{}\left[\alpha ^\lambda (\stackrel{}{y^{}})+C(t,t^{})\beta ^\lambda (\stackrel{}{y^{}})\right].$$
(121)
The eigenvalues can now be found if one adds and subtracts eqs. (120) and (121) and then sums both sides over $`\stackrel{}{y}`$, obtaining
$`(\lambda 1)\left[{\displaystyle \underset{\stackrel{}{y}V_x}{}}\alpha ^\lambda (\stackrel{}{y})\pm {\displaystyle \underset{\stackrel{}{y}V_x}{}}\beta ^\lambda (\stackrel{}{y})\right]`$
$`=i\eta [C(t,t^{})\pm 1]\left[{\displaystyle \underset{\stackrel{}{y}V_x}{}}\alpha ^\lambda (\stackrel{}{y})\pm {\displaystyle \underset{\stackrel{}{y}V_x}{}}\beta ^\lambda (\stackrel{}{y})\right],`$ (122)
which has two non-trivial solutions
$$\lambda _\pm =1+i\eta [C(t,t^{})\pm 1],$$
and $`2V_x2`$ trivial solutions such that $`\lambda =1`$ and $`_{\stackrel{}{y}V_x}\alpha ^\lambda (\stackrel{}{y})=_{\stackrel{}{y}V_x}\beta ^\lambda (\stackrel{}{y})=0`$. Thus, in the case $`\mathrm{}L(t),L(t^{})`$ we recover the same eigenvalues, and hence the same distribution as in the case $`V_x=1`$. This result was to be expected since coarse-graining of completely correlated regions should not affect the distribution obtained for a single site.
#### D.2.2 Case $`\mathrm{}L(t),L(t^{})`$
One can seek approximate solutions of eqs. (118) and (119) in this limit if one assumes that the $`\alpha ^\lambda (\stackrel{}{y})`$ and $`\beta ^\lambda (\stackrel{}{y})`$ are slowly varying functions of $`\stackrel{}{y}`$, in which case one must solve the approximate equations
$$(\lambda 1)\alpha ^\lambda (\stackrel{}{y})i\eta \frac{L^d}{V_x}\left[C(t,t^{})\alpha ^\lambda (\stackrel{}{y})+\beta ^\lambda (\stackrel{}{y})\right],$$
(123)
$$(\lambda 1)\beta ^\lambda (\stackrel{}{y})i\eta \frac{L^d}{V_x}\left[\alpha ^\lambda (\stackrel{}{y})+C(t,t^{})\beta ^\lambda (\stackrel{}{y})\right],$$
(124)
where for simplicity we considered $`L=L(t)L(t^{})`$. These equations admit non-trivial solutions
$$\lambda _\pm 1+i\eta \left(\frac{L}{\mathrm{}}\right)^d(C(t,t^{})\pm 1).$$
Naively, there are as many of these solutions as the number of $`\stackrel{}{y}`$ points in $`V_x`$, for each of $`\lambda _\pm `$. However, the assumption that the $`\alpha ^\lambda (\stackrel{}{y})`$ and $`\beta ^\lambda (\stackrel{}{y})`$ are slowly varying correlates them, and thus one cannot expect that the non-trivial solutions span the whole of the $`2V_x`$ dimensional space. The number of independent non-trivial solutions should be only order $`V_x/L^d=(\mathrm{}/L)^d`$ for each of $`\lambda _\pm `$.
## Appendix E The response of composite operators
In this Appendix we compute the response of the composite operator $`\varphi (\stackrel{}{x},t)\varphi (\stackrel{}{x},t^{})`$ to a perturbation that couples to the same composite operator evaluated at a different spatial point and the same times . In the Langevin equation such a perturbation is represented by an additional deterministic time-dependent force:
$$F(\stackrel{}{x},t)=_0^t𝑑t^{\prime \prime }[h(\stackrel{}{x};t^{\prime \prime },t)+h(\stackrel{}{x};t,t^{\prime \prime })]\varphi (\stackrel{}{x},t^{\prime \prime }).$$
In the following we work at zero temperature. The perturbed field $`\varphi _h`$ is
$$\varphi _h(\stackrel{}{k},t)=r(k;t,0)\varphi (\stackrel{}{k},0)+_0^t𝑑t^{\prime \prime }r(k;t,t^{\prime \prime })F(\stackrel{}{k},t^{\prime \prime }).$$
The variation of the force $`F`$ with respect to the perturbation $`h`$ is
$$\frac{\delta F(\stackrel{}{k},t)}{\delta h(\stackrel{}{k}^{},t_1,t_2)}=\varphi (\stackrel{}{k}\stackrel{}{k}^{},t_2)\delta (tt_1)\theta (tt_2)+\varphi (\stackrel{}{k}\stackrel{}{k}^{},t_1)\delta (tt_2)\theta (tt_1)$$
The response we are interested in is given by
$`\left[{\displaystyle \frac{\delta \varphi (\stackrel{}{x},t)\varphi (\stackrel{}{x},t^{})}{\delta h(\stackrel{}{x}^{},t^{\prime \prime },t^{\prime \prime \prime })}}\right]_{ic}=\left[\varphi (\stackrel{}{x},t){\displaystyle \frac{\delta \varphi (\stackrel{}{x},t^{})}{\delta h(\stackrel{}{x}^{},t^{\prime \prime },t^{\prime \prime \prime })}}\right]_{ic}+\left[{\displaystyle \frac{\delta \varphi (\stackrel{}{x},t)}{\delta h(\stackrel{}{x}^{},t^{\prime \prime },t^{\prime \prime \prime })}}\varphi (\stackrel{}{x},t^{})\right]_{ic}.`$
After some rather straightforward calculations one finds
$`\left[{\displaystyle \frac{\delta \varphi (\stackrel{}{x},t)\varphi (\stackrel{}{x},t^{})}{\delta h(\stackrel{}{x}^{},t^{\prime \prime },t^{\prime \prime \prime })}}\right]_{ic}`$ $`=`$ $`R(\stackrel{}{x}\stackrel{}{x}^{};t,t^{\prime \prime })C(\stackrel{}{x}\stackrel{}{x}^{};t^{},t^{\prime \prime \prime })\theta (tr^{\prime \prime \prime })`$
$`+R(\stackrel{}{x}\stackrel{}{x}^{};t,t^{\prime \prime \prime })C(\stackrel{}{x}\stackrel{}{x}^{};t^{},t^{\prime \prime })\theta (tt^{\prime \prime })`$
$`+R(\stackrel{}{x}\stackrel{}{x}^{};t^{},t^{\prime \prime })C(\stackrel{}{x}\stackrel{}{x}^{};t,t^{\prime \prime \prime })\theta (t^{}t^{\prime \prime \prime })`$
$`+R(\stackrel{}{x}\stackrel{}{x}^{};t^{},t^{\prime \prime \prime })C(\stackrel{}{x}\stackrel{}{x}^{};t,t^{\prime \prime })\theta (t^{}t^{\prime \prime })`$
where $`R`$ and $`C`$ are the usual two-point, two-time linear response and correlation. One can readily verify that this expression is not simple related to time-variations of the four-point correlation $`C_4`$ contrary to what one might have naively expected. Note that this expression has the expected $`t=t^{}`$ and $`t^{\prime \prime }=t^{\prime \prime \prime }`$ limit. |
warning/0506/astro-ph0506399.html | ar5iv | text | # Hot Neutron and Quark Star Evolution Lectures delivered at the Helmholtz International Summer School and Workshop on Hot points in Astrophysics and Cosmology, JINR, Dubna, Russia, August 2 - 13, 2004.
## 1 Introduction
The interiors of compact stars are considered as systems where high-density phases of strongly interacting matter do occur in nature, see Shapiro and Teukolsky , Glendenning and Weber for textbooks. The consequences of different phase transition scenarios for the cooling behaviour of compact stars have been reviewed recently in comparison with existing X-ray data .
The Einstein Observatory was the first that started the experimental study of surface temperatures of isolated neutron stars (NS). Upper limits for some sources have been found. Then ROSAT offered first detections of surface temperatures. Next $`X`$-ray data came from Chandra and XMM/Newton. Appropriate references to the modern data can be found in recent works by , devoted to the analysis of the new data. More upper limits and detections are expected from satellites planned to be sent in the nearest future. In general, the data can be separated in three groups. Some data show very “slow cooling” of objects, other demonstrate an “intermediate cooling” and some show very “rapid cooling”. Now we are at the position to carefully compare the data with existing cooling calculations.
The “standard” scenario of neutron star cooling is based on the main process responsible for the cooling, which is the modified Urca process (MU) $`nnnpe\overline{\nu }`$ calculated using the free one pion exchange between nucleons, see . However, this scenario explains only the group of slow cooling data. To explain a group of rapid cooling data “standard” scenario was supplemented by one of the so called “exotic” processes either with pion condensate, or with kaon condensate, or with hyperons, or involving the direct Urca (DU) reactions, see and refs therein. All these processes may occur only for the density higher than a critical density, $`(2÷6)n_0`$, depending on the model, where $`n_0`$ is the nuclear saturation density. An other alternative to ”exotic” processes is the DU process on quarks related to the phase transition to quark matter.
Particularly the studies of cooling evolution of compact objects can give an opportunity for understanding of properties of cold quark gluon plasma. In dense quark matter at temperatures below $`50`$ MeV, due to attractive interaction channels, the Cooper pairing instability is expected to occur which should lead to a variety of possible quark pair condensates corresponding to color superconductivity (CSC) phases, see for a review.
Since it is difficult to provide low enough temperatures for CSC phases in heavy-ion collisions, only precursor phenomena are expected under these conditions.
CSC phases may occur in neutron star interiors and could manifest themselves, e.g., in the cooling behavior .
However, the domain of the QCD phase diagram where neutron star conditions are met is not yet accessible to Lattice QCD studies and theoretical approaches have to rely on non-perturbative QCD modeling. The class of models closest to QCD are Dyson-Schwinger equation (DSE) approaches which have been extended recently to finite temperatures and densities . Within simple, infrared-dominant DSE models early studies of quark stars and diquark condensation have been performed.
Estimates of the cooling evolution have been performed for a self-bound isothermal quark core neutron star (QCNS) which has a crust but no hadron shell, and for a quark star (QS) which has neither crust nor hadron shell. It has been shown there in the case of the 2SC (3SC) phase of QCNS that the consequences of the occurrence of gaps for the cooling curves are similar to the case of usual hadronic neutron stars (enhanced cooling). However, for the CFL case it has been shown that the cooling is extremely fast since the drop in the specific heat of superconducting quark matter dominates over the reduction of the neutrino emissivity. As has been pointed out there, the abnormal rate of the temperature drop is the consequence of the approximation of homogeneous temperature profiles the applicability of which should be limited by the heat transport effects. Page et al. (2000) estimated the cooling of hybrid neutron stars (HNS) where heat transport effects within the superconducting quark core have been disregarded. Neutrino mean free path in color superconducting quark matter have been discussed in where a short period of cooling delay at the onset of color superconductivity for a QS has been conjectured in accordance with the estimates of in the CFL case for small gaps.
A completely new situation might arise if the scenarios suggested for (color) superconductivity besides of bigger pairing gaps ($`\mathrm{\Delta }_q50÷100`$ MeV) will allow also small diquark pairing gaps ($`\mathrm{\Delta }_q<1`$ MeV) in quark matter.
The questions which should be considered within these models are the following: (i) Is strange quark matter relevant for structure and evolution of compact stars? (ii) Are stable hybrid stars with quark matter interior possible? (iii) What can we learn about possible CSC phases from neutron star cooling? Further on in this lectures we discuss the scheme and the results of realization of the these points in relation with the cooling evolution of compact objects.
In the consideration of the scenario for the thermal evolution of NS and HNS we include the heat transport in both the quark and the hadronic matter. We will demonstrate the influence of the diquark pairing gaps and the hadronic gaps on the evolution of the surface temperature.
The main strategy of the simulation of the cooling evolution of compact objects is presented in Fig 1. On the top of scheme we have the general theoretical background of QCD models as it has been discussed in the introduction. On the second level of the scheme we separate two branches one for the structure of the compact objects and the other for the thermal properties of the stellar matter. On the bottom of the scheme those two branches are combined in the code of the cooling simulations project and entail the comparison of the theoretical and observational results.
## 2 EoS of stellar matter and Compact Star configurations
### 2.1 Hadronic matter EoS and DU threshold
For the modeling of the dense hadronic matter different approaches, like relativistic and non-relativistic Brückner-Hartree-Fock approach including the three body correlations and different parameterizations of relativistic mean field models, are discussed successfully in the literature .
For the appropriate cooling simulations we exploit the EoS of (specifically the Argonne $`V18+\delta v+UIX^{}`$ model), which is based on the most recent models for the nucleon-nucleon interaction including a parameterized three-body force and relativistic boost corrections. Actually we adopt a simple analytic parameterization of this model by Heiselberg and Hjorth-Jensen , hereafter HHJ.
The latter uses the compressional part with the compressibility $`K240`$ MeV, a symmetry energy fitted to the data around nuclear saturation density and smoothly incorporates causality at high densities. The density dependence of the symmetry energy is very important since it determines the value of the threshold density for the DU process. The HHJ EoS fits the symmetry energy to the original Argonne $`V18+\delta v+UIX^{}`$ model yielding $`n_c^{\mathrm{DU}}5.0n_0`$ ($`M_c^{\mathrm{DU}}1.839M_{}`$).
Fig. 2 demonstrates the partial densities of $`p`$, $`e`$ and $`\mu ^{}`$ in HHJ and relativistic non-linear Walecka (NLW) model in the parameterization of , adjusted to the following bulk parameters of the nuclear matter at saturation: $`n_0=0.16`$ fm<sup>-3</sup>. Anyhow, in the given NLW model the threshold density for the DU process is $`n_c^{\mathrm{DU}}2.7n_0`$.
### 2.2 Quark matter EoS with 2SC superconductivity
The quark matter models predict, that the diquark pairing condensate is possible for temperatures and densities relevant for compact objects. The order of magnitude of pairing gaps is 100 MeV and a remarkably rich phase structure of the matter has been identified .
This leads to expectations of observable effects of color superconducting phases, particularly in the compact star cooling behaviour . Generally color superconductivity is involved in all aspects of neutron star studies, such as magnetic field evolution or burst-type phenomena .
The applications of the nonlocal chiral quark model developed in for the case of neutron star constraints shows that the relevant CSC phase is a 2SC phase while the omission of the strange quark flavor is justified by the fact that chemical potentials in central parts of the stars do barely reach the threshold value at which the mass gap for strange quarks breaks down and they appear in the system .
It has been shown in that work that the Gaussian formfactor ansatz of quark interaction (hereafter we call it SM model) leads to an early onset of the deconfinement transition and such a model is therefore suitable to discuss hybrid stars with large quark matter cores .
The resulting quark matter EoS within this nonlocal chiral model can be represented in a form reminiscent of a bag model
$$P^{(s)}=P_{id}(\mu _B)B^{(s)}(\mu _B),$$
(1)
where $`P_{id}(\mu _B)`$ is the ideal gas pressure of quarks and $`B^{(s)}(\mu _B)`$ a density dependent bag pressure, see Fig. 3. The occurrence of diquark condensation depends on the value of the ratio $`\eta =G_2/G_1`$ of coupling constants and the superscript $`s\{S,N\}`$ indicating whether we consider the matter in the superconducting mixed phase ($`\eta =1`$) or in the normal phase ($`\eta =0`$), respectively.
### 2.3 EoS of hybrid star matter
In some density interval above the onset of the first order phase transition, there may appear a mixed phase region between the hadronic (confined quark phase) and deconfined quark phases, see . Ref. disregarded finite size effects, such as surface tension and charge screening. On the example of the hadron-quark mixed phase Refs demonstrated that finite size effects might play a crucial role substantially narrowing the region of the mixed phase or even forbidding its appearance in a hybrid star configuration. Therefore we omit the possibility of the hadron-quark mixed phase in our model assuming that the quark phase arises by the Maxwell construction.
To demonstrate the EoS with quark-hadron phase transition applicable for hybrid star configurations in the next Section, we show in Fig. 4 results using the relativistic mean field (RMF) model of asymmetric nuclear matter including a non-linear scalar field potential and the $`\rho `$ meson (nonlinear Walecka model) ( see ) for the cases $`\eta =1`$ (left panel) when the quark matter phase is superconducting and for $`\eta =0`$ (right panel) when it is normal.
In case the hadonic EoS is chosen to be HHJ model and the quark one is SM model we found a tiny density jump at the phase-boundary from $`n_c^{\mathrm{hadr}}0.44\mathrm{fm}^3`$ to $`n_c^{\mathrm{quark}}0.46\mathrm{fm}^3`$. The critical mass, when the branch of stable hybrid stars starts is $`M_c^{\mathrm{quark}}=1.214M_{}`$.
As shown in Refs. quark matter itself can exist in a mixed phase of 2SC and normal states. The presence or absence of the 2SC - normal quark mixed phase instead of only one of those phases is not so important for the hybrid star cooling problem since the latter is governed by processes involving either normal excitations or excitations with the smallest gap.
### 2.4 Stability of Hybrid star configurations
Not all compositions of hadronic and quark matter EoS’s give stable configu-rations for compact objects, which could be demonstrated using the HHJ model EoS for hadronic and the two flavor nonlocal chiral model EoS for quark matter.
The mechanical structure of the spherically symmetric, static gravitational self-bound configurations of dense matter can be calculated with the well-known Tolman-Oppenheimer-Volkoff equations. These are the conditions of hydrostatic equilibrium of self- gravitating matter, see also ,
$$\frac{dP(r)}{dr}=\frac{[\epsilon (r)+P(r)][m(r)+4\pi r^3P(r)]}{r[r2m(r)]}.$$
(2)
Here $`\epsilon (r)`$ is the energy density and $`P(r)`$ the pressure at distance $`r`$ from the center of the star. The mass enclosed in a sphere with radius $`r`$ is defined by
$$m(r)=4\pi _0^r\epsilon (r^{})r^2𝑑r^{}.$$
(3)
These equations are solved for given central baryon number densities, $`n_B(r=0)`$, thereby defining a sequence of star configurations.
The configurations are stable if they are on the rising branch of the mass - central density relation.
The relation between pressure and energy density is given by the choice of the corresponding EoS model, which generally can be a function of temperature. The temperature is either constant for isothermal configurations, see , or has some profile, when the cooling evolution is assumed to be a hydrodynamically quasi-stationary process.
Hot quark stars have been discussed, e.g., in .
For late cooling, when the temperatures are below $`T<1\mathrm{MeV}`$ the effects of temperature on the distribution of matter are negligible and the calculations for the star structure can be performed for the $`T=0`$ case.
Dots in Fig. 2 indicate threshold densities for the DU process. The possibility of charged pion condensation is suppressed. Otherwise for $`n>n_c^{\mathrm{DU}}`$ the isotopic composition is changed in favor of increasing proton fraction and a smaller critical density for the DU reaction. Deviations in the $`M(n)`$ relation for HHJ and NLW EoS are minor, whereas the DU thresholds are quite distinct.
In Fig. 5 we show the mass-radius relation for hybrid stars with HHJ EoS vs. Gaussian nonlocal chiral quark separable model (SM) EoS. Configurations with possible 2SC phase given by the solid line, are stable, whereas the hybrid star configurations without 2SC large gap color superconductivity (dash-dotted lines) are unstable. In case of “HHJ-SM with 2SC” the maximum neutron star mass proves to be $`1.793M_{}`$. For an illustration of the constraints on the mass-radius relation which can be derived from compact star observations we show in Fig. 5 the compactness limits from the thermal emission of the isolated neutron star RX J1856.5-3754 as given in and from the redshifted absorption lines in the X-ray burst spectra of EXO 0748-676 given in . These are, however, rather weak constraints.
## 3 Heat transport, Neutrino production and diffusion
### 3.1 Boltzmann equation in curved space - time
For study of cooling evolution, we will use the diffusion approximation for neutrino transport, because it allows us easily to assess how global characteristics of the neutrino emission such as average energies, integrated fluxes, and time scales change in response to various input parameters. This approach is sufficient to establish connections between the microphysical ingredients and the duration of deleptonization and cooling time scales, which are necessary to estimate possible effects on neutrino signals. However, the detailed results of simulations for early cooling will not be presented in this lecture. The reader is referred to the works and citations within.
We will assume that the neutron star has an interior where the matter is spherically symmetric distributed. The evolution of the star will be assumed to be quasi stationary, which means that for each time step the velocity of the matter is zero.
The Boltzmann equation (BE) for massless particles is
$$p^\beta \left(\frac{f}{x^\beta }\mathrm{\Gamma }_{\beta \gamma }^\alpha p^\gamma \frac{f}{p^\alpha }\right)=\left(\frac{df}{d\tau }\right)_{coll}$$
(4)
where $`f`$ is the invariant neutrino distribution function, $`p^\alpha `$ is the neutrino 4-momentum and $`\mathrm{\Gamma }_{\beta \gamma }^\alpha `$ are the Christoffel symbols for the metric
$$d\tau ^2=e^{2\varphi }dt^2+e^{2\mathrm{\Lambda }}dr^2+r^2d\theta ^2+r^2\mathrm{sin}^2\theta d\mathrm{\Phi }^2$$
(5)
of static spherically symmetric space-time manifold.
For simplicity one can use the comoving basis and rewrite the BE
$$p^b\left(e_b^\beta \frac{f}{x^\beta }\mathrm{\Gamma }_{bc}^ap^c\frac{f}{p^a}\right)=\left(\frac{df}{d\tau }\right)_{coll},$$
(6)
where $`e_b^\beta `$ are basis vectors in a comoving frame to an observer and in the static case they are diagonal. Indices $`a,b,c`$ are running from 0 to 3. The $`\mathrm{\Gamma }_{bc}^a`$ are Ricci rotational coefficients with the non-zero components
$`\mathrm{\Gamma }_{00}^1`$ $`=`$ $`\mathrm{\Gamma }_{01}^0=e^\mathrm{\Lambda }{\displaystyle \frac{\varphi }{r}},`$ (7)
$`\mathrm{\Gamma }_{21}^2`$ $`=`$ $`\mathrm{\Gamma }_{22}^1=\mathrm{\Gamma }_{33}^1=\mathrm{\Gamma }_{31}^3={\displaystyle \frac{e^\mathrm{\Lambda }}{r}},`$
$`\mathrm{\Gamma }_{33}^2`$ $`=`$ $`\mathrm{\Gamma }_{32}^3={\displaystyle \frac{\mathrm{cot}\theta }{r}}.`$
The neutrino 4-momentum is
$$p^a=(\omega ,\omega \mu ,\omega (1\mu ^2)^{1/2}\mathrm{cos}\mathrm{\Phi },\omega (1\mu ^2)^{1/2}\mathrm{sin}\mathrm{\Phi }),$$
where $`\mu `$ is the cosine of the angle between the neutrino momentum and the radial direction, $`\omega `$ is the neutrino energy in a comoving frame. Using the previous definitions we will have
$`\omega e_0^t{\displaystyle \frac{f}{t}}`$ $`+`$ $`\omega \mu e_1^r{\displaystyle \frac{f}{r}}\omega ^2\mu \mathrm{\Gamma }_{00}^1{\displaystyle \frac{f}{\omega }}\omega (1\mu ^2)\left(\mathrm{\Gamma }_{00}^1+\mathrm{\Gamma }_{22}^1\right){\displaystyle \frac{f}{\mu }}=\left({\displaystyle \frac{df}{d\tau }}\right)_{coll}.`$ (8)
#### 3.1.1 Equations in Linear response Approximation
After applying the operator
$$\frac{1}{2}_1^{+1}𝑑\mu \mu ^i,i=0,1,2,\mathrm{},$$
to equation (8) and defining the $`i^{\mathrm{th}}`$ moments
$$M_i=\frac{1}{2}_1^{+1}𝑑\mu \mu ^if,Q_i=\frac{1}{2}_1^{+1}𝑑\mu \mu ^i\left(\frac{df}{d\tau }\right)_{coll}.$$
we get for $`i=0`$
$`\omega \left(e_0^t{\displaystyle \frac{M_0}{t}}+e_1^r{\displaystyle \frac{M_1}{r}}\right)\omega ^2\left(\mathrm{\Gamma }_{00}^1{\displaystyle \frac{M_1}{\omega }}\right)2\omega (\mathrm{\Gamma }_{00}^1+\mathrm{\Gamma }_{22}^1)M_1=Q_0,`$ (9)
and for $`i=1`$
$`\omega \left(e_0^t{\displaystyle \frac{M_1}{t}}+e_1^r{\displaystyle \frac{M_2}{r}}\right)\omega ^2\left(\mathrm{\Gamma }_{00}^1{\displaystyle \frac{M_2}{\omega }}\right)+\omega (\mathrm{\Gamma }_{00}^1+\mathrm{\Gamma }_{22}^1)(M_03M_2)=Q_1.`$ (10)
Let us introduce $`N_\nu `$, $`F_\nu `$, and $`S_N`$ as the number density, number flux and the number source term, respectively,
$$N_\nu =_0^{\mathrm{}}\frac{d\omega }{2\pi ^2}M_0\omega ^2,F_\nu =_0^{\mathrm{}}\frac{d\omega }{2\pi ^2}M_1\omega ^2,S_N=_0^{\mathrm{}}\frac{d\omega }{2\pi ^2}Q_0\omega ,$$
while $`J_\nu `$, $`H_\nu `$, $`P_\nu `$, and $`S_E`$ are the neutrino energy density, energy flux, pressure, and the energy source term
$`J_\nu `$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d\omega }{2\pi ^2}}M_0\omega ^3,H_\nu ={\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d\omega }{2\pi ^2}}M_1\omega ^3,`$
$`P_\nu `$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d\omega }{2\pi ^2}}M_2\omega ^3,S_E={\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d\omega }{2\pi ^2}}Q_0\omega ^2.`$ (11)
After integration over the neutrino energy and utilizing the continuity equation under assumption of a quasi-static evolution, one can recover the well known neutrino transport equation
$$\frac{(N_\nu /n_B)}{t}+\frac{(e^\varphi 4\pi r^2F_\nu )}{a}=e^\varphi \frac{S_N}{n_B},$$
(12)
$$\frac{(J_\nu /n_B)}{t}+P_\nu \frac{(1/n_B)}{t}+e^\varphi \frac{(e^{2\varphi }4\pi r^2H_\nu )}{a}=e^\varphi \frac{S_E}{n_B}.$$
(13)
The distribution function in the diffusion approximation could be represented in the following way
$$f(\omega ,\mu )=f_0(\omega )+\mu f_1(\omega ),f_0=[1+e^{\left(\frac{\omega \mu _\nu }{kT}\right)}]^1,$$
(14)
where $`f_0(\omega )`$ is the distribution function in equilibrium ($`T=T_{mat}`$, $`\mu _\nu =\mu _\nu ^{eq}`$) with the neutrino energy $`\omega `$ and the chemical potential $`\mu _\nu `$ respectively. Hereafter the dependence of $`f_0`$ and $`f_1`$ on $`\omega `$ and non explicit dependence on space-time coordinates will be assumed without listing in arguments.
Thus the moments $`M_i`$ of the distribution function $`f`$ are
$$M_0=f_0,M_1=\frac{1}{3}f_1,M_2=\frac{1}{3}f_0,\mathrm{and}M_3=\frac{1}{5}f_1.$$
(15)
Therefore the equation (10) now reads
$$e^\mathrm{\Lambda }\left(\frac{f_0}{r}\omega \frac{\varphi }{r}\frac{f_0}{\omega }\right)=3\frac{Q_1}{\omega }.$$
(16)
The collision term $`Q_1`$ can be represented as
$$\left(\frac{df}{d\tau }\right)_{coll}=\omega \left(j_a(1f)\frac{f}{\lambda _a}+j_s(1f)\frac{f}{\lambda _s}\right).$$
(17)
Here $`j_a`$ is the emissivity, $`\lambda _a`$ is the absorptivity, $`j_s`$ and $`\lambda _s`$ are the scattering contributions.
Namely
$$j_s=\frac{1}{(2\pi )^3}_0^{\mathrm{}}𝑑\omega ^{}\omega ^2_1^1𝑑\mu ^{}_0^{2\pi }𝑑\mathrm{\Phi }f(\omega ^{},\mu ^{})R_s^{in}(\omega ,\omega ^{},\mathrm{cos}\theta ),$$
(18)
$$\frac{1}{\lambda _s}=\frac{1}{(2\pi )^3}_0^{\mathrm{}}𝑑\omega ^{}\omega ^2_1^1𝑑\mu ^{}_0^{2\pi }𝑑\mathrm{\Phi }[1f(\omega ^{},\mu ^{})]R_s^{out}(\omega ,\omega ^{},\mathrm{cos}\theta ),$$
(19)
where $`\theta `$ is the scattering angle. The relation between emissivities and absorptivities ($`R_s^{in}`$ and $`R_s^{out}`$ are scattering kernels) is given by
$$\frac{1}{\lambda _a(\omega )}=e^{\beta (\omega \mu _\nu ^{eq})}j_a(\omega )\mathrm{and}R_s^{in}=e^{\beta (\omega ^{}\omega )}R_s^{out}.$$
(20)
Using the Legendre expansion for the moments one has
$$R_l^{out}=_1^1d\mathrm{cos}\theta P_l(\mathrm{cos}\theta )R_s^{out}(\omega ,\omega ^{},\mathrm{cos}\theta ).$$
(21)
Performing the angular integrations of Eq. (17) one can define the relation between $`Q_0`$, $`Q_1`$ and $`R_0^{out}`$, $`R_1^{out}`$.
After substitution of the expression for $`Q_1`$ into Eq.(16) the relation between $`f_0`$ and $`f_1`$ reads
$`f_1=D(\omega )\left[{\displaystyle \frac{f_0}{r}}\omega {\displaystyle \frac{\varphi }{r}}{\displaystyle \frac{f_0}{\omega }}\right]e^\mathrm{\Lambda },D(\omega )=\left(j_a+{\displaystyle \frac{1}{\lambda _a}}+\kappa _1^s\right)^1,`$ (22)
where $`D(\omega )`$ is the diffusion coefficient. Using the relation between $`\frac{f_0}{r}`$ and $`\frac{f_0}{\omega }`$ and the notation $`\eta =\mu _\nu /T`$ to be the neutrino degeneracy parameter, we obtain
$$f_1=D(\omega )e^\mathrm{\Lambda }\left[T\frac{\eta }{r}+\frac{\omega }{Te^\varphi }\frac{(Te^\varphi )}{r}\right]\left(\frac{f_0}{\omega }\right).$$
(23)
Now the energy-integrated lepton and energy fluxes are
$`F_\nu `$ $`=`$ $`{\displaystyle \frac{e^\mathrm{\Lambda }e^\varphi T^2}{6\pi ^2}}\left[D_3{\displaystyle \frac{(Te^\varphi )}{r}}+(Te^\varphi )D_2{\displaystyle \frac{\eta }{r}}\right]`$
$`H_\nu `$ $`=`$ $`{\displaystyle \frac{e^\mathrm{\Lambda }e^\varphi T^3}{6\pi ^2}}\left[D_4{\displaystyle \frac{(Te^\varphi )}{r}}+(Te^\varphi )D_3{\displaystyle \frac{\eta }{r}}\right].`$ (24)
Here the coefficients $`D_2`$, $`D_3`$, and $`D_4`$ are defined through the integrals over $`\omega `$ as
$$D_n=T^{(n+1)}_0^{\mathrm{}}𝑑\omega \omega ^nD(\omega )f_0(\omega )(1f_0(\omega )).$$
The explicit expression for the absorption mean free path is
$`{\displaystyle \frac{1}{\lambda _a}}={\displaystyle \frac{G_F^2}{\pi ^2}}{\displaystyle _0^{\mathrm{}}}𝑑E_eE_e^2\left[1f_{eq}(E_e)\right]{\displaystyle _1^{+1}}d\mathrm{cos}\theta {\displaystyle \frac{(\mathrm{cos}\theta 1)}{1e^z}}\left[AR_1+R_2+BR_3\right]`$ (25)
$$R_s^{out}=4G_F^2\frac{(\mathrm{cos}\theta 1)}{1e^z}\left[AR_1+R_2+BR_3\right]$$
(26)
The response functions written in terms of the polarization functions are
$`R_1`$ $`=`$ $`(𝒱^2+𝒜^2)[\mathrm{Im}\mathrm{\Pi }_L^R(q_0,q)+\mathrm{Im}\mathrm{\Pi }_T^R(q_0,q)]`$ (27)
$`R_2`$ $`=`$ $`(𝒱^2+𝒜^2)\mathrm{Im}\mathrm{\Pi }_T^R(q_0,q)𝒜^2\mathrm{Im}\mathrm{\Pi }_A^R(q_0,q)`$ (28)
$`R_3`$ $`=`$ $`2𝒱𝒜\mathrm{Im}\mathrm{\Pi }_{VA}^R(q_0,q),`$ (29)
where the coupling constants for the absorption are $`𝒱=g_V\mathrm{cos}(\theta _c)`$ and $`𝒜=g_A\mathrm{cos}(\theta _c)`$ and for the scattering $`𝒱=g_V/2`$ and $`𝒜=g_A/2`$. $`\theta _c`$ is the Cabibbo angle, $`g_V`$ and $`g_A`$ are the vector and axial-vector couplings.
Introducing the $`s`$ entropy per baryon and $`Y_L=(N_\nu +N_e+N_\mu )/n_B`$ lepton number fraction and with help of conservation laws for energy and baryon number combined with Eqs. (12) and (13) we finally will obtain the energy transport equation
$$Te^\varphi \frac{s}{t}+\mu _\nu e^\varphi \frac{Y_L}{t}+\frac{(e^{2\varphi }4\pi r^2H_\nu )}{a}=0$$
(30)
and the equation for lepton diffusion
$`{\displaystyle \frac{Y_L}{t}}+{\displaystyle \frac{(e^\varphi 4\pi r^2F_\nu )}{a}}=0.`$ (31)
The energy $`H_\nu `$ and lepton $`F_\nu `$ fluxes are defined in Eq.(24).
These coupled equations for the unknowns $`T`$ and $`\mu _\nu `$ with the Eq.(2) for the star structure completely describe the early cooling evolution of the compact objects in quasi-stationary regime. The solutions and more detailed information one can find in . In the next section we will use the energy transport equation for consideration of late cooling evolution when the neutrinos are already untrapped, i.e. the mean free path $`\lambda _a`$ is larger than the star radius $`R`$ and $`\mu _\nu =0`$.
### 3.2 Late cooling evolution
The cooling scenario is described with a set of cooling regulators. In the regime without neutrino diffusion in the righthand side of Eq. (13) and therefore also in Eq. (30) for energy transport we introduce the energy loss term $`ϵ_\nu `$ due to neutrino emission additional to $`S_E`$ and instead of the $`H_\nu `$ we use the total energy flux function $`l(r,t)`$. The temperature profile $`T(r,t)`$ is the one time dependent unknown dynamical quantity and during the cooling process due to an inhomogeneous distribution of the matter inside the star and the finite heat conductivity it can differ from the isothermal one. The heat conductivity $`\kappa `$ (the corresponding term to $`D_4`$ in Eq. (30)) determines the relation between $`l(r,t)`$ and gradient of $`T(r,t)`$. In isothermal age the temperature on the inner crust boundary ($`T_m`$) and the central temperature ($`T_c`$) are connected by the relation $`T_m=T_c\mathrm{exp}[\varphi (0)\varphi (R)]`$. Here $`\varphi (0)\varphi (R)`$ is the difference of the gravitational potentials in the center and at the surface of the star, respectively.
The flux of energy $`l(r)`$ per unit time through a spherical slice at the distance $`r`$ from the center, is proportional to the gradient of the temperature on both sides of this slice,
$$l(r)=4\pi r^2\kappa (r)\frac{(T\mathrm{e}^\varphi )}{r}\mathrm{e}^\varphi \sqrt{1\frac{2M}{r}},$$
(32)
where the factor $`\mathrm{e}^\varphi \sqrt{1\frac{2M}{r}}`$ corresponds to the relativistic correction of the time scale and the unit of thickness. The equations for energy balance and thermal energy transport are
$`{\displaystyle \frac{}{A}}\left(l\mathrm{e}^{2\varphi }\right)`$ $`=`$ $`{\displaystyle \frac{1}{n}}\left(ϵ_\nu \mathrm{e}^{2\varphi }+c_V{\displaystyle \frac{}{t}}(T\mathrm{e}^\varphi )\right),`$ (33)
$`{\displaystyle \frac{}{A}}\left(T\mathrm{e}^\varphi \right)`$ $`=`$ $`{\displaystyle \frac{1}{\kappa }}{\displaystyle \frac{l\mathrm{e}^\varphi }{16\pi ^2r^4n}},`$ (34)
where $`n=n(r)`$ is the baryon number density, $`A=A(r)`$ is the total baryon number within a sphere of radius $`r`$. One has
$$\frac{r}{A}=\frac{1}{4\pi r^2n}\sqrt{1\frac{2M}{r}}.$$
(35)
The total neutrino emissivity $`ϵ_\nu `$ and the total specific heat $`c_V`$ are given as the sum of the corresponding partial contributions defined in the next subsections. The density profiles $`n_i(r)`$ of the constituents $`i`$ of the matter are under the conditions of the actual temperature profile $`T(r,t)`$. The accumulated mass $`M=M(r)`$ and the gravitational potential $`\varphi =\varphi (r)`$ can be determined by
$`{\displaystyle \frac{M}{A}}`$ $`=`$ $`{\displaystyle \frac{\epsilon }{n}}\sqrt{1{\displaystyle \frac{2M}{r}}},`$ (36)
$`{\displaystyle \frac{\varphi }{A}}`$ $`=`$ $`{\displaystyle \frac{4\pi r^3p+M}{4\pi r^2n}}{\displaystyle \frac{1}{\sqrt{1\frac{2M}{r}}}},`$ (37)
where the energy density profile $`\epsilon =\epsilon (r)`$ and the pressure profile $`p=p(r)`$ is defined by the condition of hydrodynamical equilibrium (see Eq. 2)
$$\frac{p}{A}=(p+\epsilon )\frac{\varphi }{A}.$$
(38)
The boundary conditions for the solution of (33) and (34) read $`l(r=0)=l(A=0)=0`$ and $`T(A(r_m)=A,t)=T(r_m=R,t)=T_m(t)`$, respectively.
In our examples we choose the initial temperature to be 1 MeV. This is a typical value for the temperature $`T_{\mathrm{opacity}}`$ at which the star becomes transparent for neutrinos. Simplifying we disregard the neutrino influence on transport. These effects dominate for $`t<1÷100`$ min, when the star cools down to $`TT_{\mathrm{opacity}}`$ and become unimportant for later times.
#### 3.2.1 Neutrino Processes in dense matter
We compute the NS thermal evolution adopting our fully general relativistic evolutionary code. This code was originally constructed for the description of hybrid stars by . The main cooling regulators are the thermal conductivity, the heat capacity and the emissivity.
#### Emissivity
The luminosities are calculated using the corresponding matrix element of the neutrino production process
$$L_\nu =(2\pi )^4\frac{d^3p_n}{(2\pi )^32E_n}\mathrm{}\frac{d^3p_\nu }{(2\pi )^32E_\nu }\delta ^3(\stackrel{}{p}_i)\delta (E_i)\left|M_{fi}\right|^2f_n(1f_p)(1f_e).$$
(39)
The most effective process of neutrino production is the direct Urca (DU) one, $`np+e^{}+\overline{\nu }_e`$, which is allowed when the momentum conservation holds $`\stackrel{}{p}_{F,n}=\stackrel{}{p}_{F,p}+\stackrel{}{p}_{F,e}`$ $``$ $`|\stackrel{}{p}_{F,n}||\stackrel{}{p}_{F,p}|+|\stackrel{}{p}_{F,e}|`$. Since DU process is so intensive that as soon as it takes place the star cools too fast for better correspondence between the cooling simulations and the observational data, it is likely to assumed DU to be forbidden in the interior of the star.
The other processes, i.e. the modified Urca (MU), pair breaking and formation (PBF) and bremsstrahlung are the main ones governing the cooling in hadronic matter. The modified Urca (MU) $`nnnpe\overline{\nu }`$, $`npppe\overline{\nu }`$, has to be medium modified Urca (MMU) process due to the softening of in-medium pion propagator. The final emissivity is given by
$`ϵ_\nu ^{\mathrm{nMU}}`$ $`=`$ $`8.6\times 10^{21}m_{\mathrm{nMU}}^4(Y_eu)^{1/3}\zeta _{\mathrm{nMU}}T_9^8\mathrm{ergcm}^3\mathrm{s}^1,`$ (40)
$`ϵ_\nu ^{\mathrm{pMU}}`$ $`=`$ $`8.5\times 10^{21}m_{\mathrm{pMU}}^4(Y_eu)^{1/3}\zeta _{\mathrm{pMU}}T_9^8\mathrm{ergcm}^3\mathrm{s}^1.`$ (41)
Here $`m_i^{}=\sqrt{m_{\mathrm{rel},i}^2+p_{\mathrm{F},i}^2}`$ is the non-relativistic quasiparticle effective mass related to the in-medium one-particle energies from a given relativistic mean field model for $`i=n,p`$. We have introduced the abbreviations $`m_{\mathrm{nMU}}^4=(m_n^{}/m_n)^3(m_p^{}/m_p)`$ and $`m_{pMU}^4=(m_p^{}/m_p)^3(m_n^{}/m_n)`$. The suppression factors are $`\zeta _{\mathrm{nMU}}=\zeta _n\zeta _p\mathrm{exp}\{[\mathrm{\Delta }_n(T)+\mathrm{\Delta }_p(T)]/T\}`$, $`\zeta _{\mathrm{pMU}}\zeta _p^2`$, and should be replaced by unity for $`T>T_{\mathrm{crit},i}`$, when for given species $`i`$ the corresponding gap vanishes. For neutron and proton $`S`$-wave pairing is $`\mathrm{\Delta }_i(0)=1.76T_{\mathrm{crit},i}`$ and for the $`P`$-wave pairing of neutrons $`\mathrm{\Delta }_n(0)=1.19T_{\mathrm{crit},n}`$ (see Fig. 6). The gap as a function of temperature is given by the interpolation formula $`\mathrm{\Delta }_N(T)=\mathrm{\Delta }(0)\sqrt{1T/T_{\mathrm{crit},N}}`$.
To be conservative we have used in (40) the free one-pion exchange estimate of the $`NN`$ interaction amplitude. Restricting ourselves to a qualitative analysis we use here simplified exponential suppression factors $`\zeta _i`$. In a more detailed analysis these $`\zeta _i`$-factors have prefactors with rather strong temperature dependences . At temperatures $`TT_c`$ their inclusion only slightly affects the resulting cooling curves. For $`TT_c`$ the MU process gives in any case a negligible contribution to the total emissivity and thereby corresponding modifications can again be omitted. Also for the sake of simplicity the general possibility of a <sup>3</sup>P$`{}_{2}{}^{}(|m_J|=2)`$ pairing which may result in a power-law behaviour of the specific heat and the emissivity of the MU process is disregarded since mechanisms of this type of pairing are up to now not elaborated. Even more essential modifications of the MU rates may presumably come from in-medium effects which could result in extra prefactors of $`10^2÷10^3`$ already at $`TT_c`$.
In order to estimate the role of the in-medium effects in the $`NN`$ interaction for the HNS cooling we have also performed calculations for the so-called medium modified Urca (MMU) process by multiplying the rates (40) by the appropriate prefactor
$$ϵ_\nu ^{\mathrm{MMU}}/ϵ_\nu ^{\mathrm{MU}}10^3\left[\mathrm{\Gamma }^6(g^{})/\stackrel{~}{\omega }^8(kp_F)\right]u^{10/3},$$
(42)
where the value $`\mathrm{\Gamma }(g^{})1/[1+1.4u^{1/3}]`$ is due to the dressing of $`\pi NN`$ vertices and $`\stackrel{~}{\omega }m_\pi `$ is the effective pion gap which we took as function of density from Fig. 2 of .
For $`T<T_{\mathrm{crit}}`$ the most important contribution comes from the neutron and the proton pair breaking and formation processes. We take their emissivities from Ref. , which is applicable for both cases, $`S`$\- and $`P`$-wave nucleon pairing
$`ϵ_\nu ^{\mathrm{nPBF}}`$ $`=`$ $`6.6\times 10^{28}(m_n^{}/m_n)(\mathrm{\Delta }_n(T)/\mathrm{MeV})^7u^{1/3}`$ (43)
$`\times \xi I(\mathrm{\Delta }_n(T)/T)\mathrm{erg}\mathrm{cm}^3\mathrm{s}^1,`$
$`ϵ_\nu ^{\mathrm{pPBF}}`$ $`=`$ $`0.8\times 10^{28}(m_p^{}/m_p)(\mathrm{\Delta }_p(T)/\mathrm{MeV})^7u^{2/3}`$ (44)
$`\times I(\mathrm{\Delta }_p(T)/T)\mathrm{erg}\mathrm{cm}^3\mathrm{s}^1,`$
where
$$I(\mathrm{\Delta }_i(T)/T)0.89\sqrt{T/\mathrm{\Delta }_i(T)}\mathrm{exp}[2\mathrm{\Delta }_i(T)/T],$$
(45)
$`\xi 0.5`$ for $`{}_{}{}^{1}S_{0}^{}`$ pairing and $`\xi 1`$ for $`{}_{}{}^{3}P_{2}^{}`$ pairing.
A significant contribution of the proton channel is due to the $`NN`$ correlation effects, which are taken into account in .
The phonon contribution to the emissivity of the $`3P_2`$ superfluid phase is negligible. The main emissivity regulators are the MMU Eq. (42) and the neutron (nPBF) and proton (pPBF) pair breaking and formation processes Eq. (43), see above for a rough estimation. Finally, we include the effect of a pion condensate (PU process), where the PU emissivity is about 1-2 orders of magnitude smaller than the DU one.
All emissivities are corrected by correlation effects. We adopt the same set of partial emissivities as in the work of .
In quark matter the process $`du+e^{}+\overline{\nu }_e`$ is always allowed unless it is suppressed by huge energy gaps due to quark pairing.
$`ϵ_\nu ^{\mathrm{QDU}}9.4\times 10^{26}\alpha _suY_e^{1/3}\zeta _{\mathrm{QDU}}T_9^6\mathrm{erg}\mathrm{cm}^3\mathrm{s}^1,`$ (46)
where at a compression $`u=n/n_02`$ the strong coupling constant is $`\alpha _s1`$ and decreases logarithmically at still higher densities. The nuclear saturation density is $`n_0=0.17\mathrm{fm}^3`$, $`Y_e=n_e/n`$ is the electron fraction, and $`T_9`$ is the temperature in units of $`10^9`$ K. If, for a somewhat higher density, the electron fraction was too small ($`Y_e<Y_{ec}10^8`$), then all the QDU processes would be completely switched off and the neutrino emission would be governed by two-quark reactions like the quark modified Urca (QMU) and the quark bremsstrahlung (QB) processes $`dquqe\overline{\nu }`$ and $`q_1q_2q_1q_2\nu \overline{\nu }`$, respectively. The emissivities of the QMU and QB processes have been estimated as
$`ϵ_\nu ^{\mathrm{QMU}}ϵ_\nu ^{\mathrm{QB}}9.0\times 10^{19}\zeta _{\mathrm{QMU}}T_9^8\mathrm{erg}\mathrm{cm}^3\mathrm{s}^1.`$ (47)
Due to the pairing, the emissivities of QDU processes are suppressed by a factor $`\zeta _{\mathrm{QDU}}\text{exp}(\mathrm{\Delta }_q/T)`$ and the emissivities of QMU and QB processes are suppressed by a factor $`\zeta _{\mathrm{QMU}}\text{exp}(2\mathrm{\Delta }_q/T)`$ for $`T<T_{\mathrm{crit},q}0.4\mathrm{\Delta }_q`$ whereas for $`T>T_{\mathrm{crit},q}`$ these factors are equal to unity. The modification of $`T_{\mathrm{crit},q}(\mathrm{\Delta }_q)`$ relative to the standard BCS formula is due to the formation of correlations as, e.g., instanton- anti-instanton molecules . For the temperature dependence of the gap below $`T_{\mathrm{crit},q}`$ we use the interpolation formula $`\mathrm{\Delta }(T)=\mathrm{\Delta }(0)\sqrt{1T/T_{\mathrm{crit},q}}`$, with $`\mathrm{\Delta }(0)`$ being the gap at zero temperature.
The contribution of the reaction $`eeee\nu \overline{\nu }`$ is very small
$$ϵ_\nu ^{ee}=2.8\times 10^{12}Y_e^{1/3}u^{1/3}T_9^8\mathrm{erg}\mathrm{cm}^3\mathrm{s}^1,$$
(48)
but can become important, when quark processes are blocked out for large values of $`\mathrm{\Delta }_q/T`$ in superconducting quark matter.
#### Thermal conductivity
The heat conductivity of the matter is the sum of the partial contributions
$$\kappa =\underset{i}{}\kappa _i,\frac{1}{\kappa _i}=\underset{j}{}\frac{1}{\kappa _{ij}},$$
(49)
where $`i,j`$ denote the components (particle species).
The contribution of neutrons and protons is
$`\kappa _{nn}`$ $`=`$ $`8.3\times 10^{22}\left({\displaystyle \frac{m_n}{m_n^{}}}\right)^4{\displaystyle \frac{z_n^3\zeta _n}{S_{kn}T_9}}\mathrm{erg}\mathrm{s}^1\mathrm{cm}^1\mathrm{K}^1,`$ (50)
$`S_{kn}`$ $`=`$ $`0.38z_n^{7/2}+3.7z_n^{2/5},`$ (51)
$`\kappa _{np}`$ $`=`$ $`8.9\times 10^{16}\left({\displaystyle \frac{m_n}{m_n^{}}}\right)^2{\displaystyle \frac{z_n^2\zeta _pT_9}{z_p^3S_{kp}}}\mathrm{erg}\mathrm{s}^1\mathrm{cm}^1\mathrm{K}^1,`$ (52)
$`S_{kp}`$ $`=`$ $`1.83z_n^2+1.43z_n^2(0.4+z_n^8)^1.`$ (53)
Here we have introduced the appropriate suppression factors $`\zeta _i`$ which act in the presence of gaps for superfluid hadronic matter (see Fig. 6) and we have used the abbreviation $`z_i=(n_i/(4n_0))^{1/3}`$. The heat conductivity of electrons is given by Eq. (56). The total contribution related to electrons is then
$`1/\kappa _e=1/\kappa _{ee}+1/\kappa _{ep}.`$ (54)
Similar expressions we have for neutrons and protons using Eqs. (50)-(53).
The total thermal conductivity is the straight sum of the partial contributions $`\kappa _{tot}=\kappa _e+\kappa _n+\mathrm{}`$ . Other contributions to this sum are smaller than those presented explicitly ($`\kappa _e`$ and $`\kappa _n`$).
For quark matter $`\kappa `$ is the sum of the partial conductivities of the electron, quark and gluon components
$$\kappa =\kappa _e+\kappa _q+\kappa _g,$$
(55)
where $`\kappa _e\kappa _{ee}`$ is determined by electron-electron scattering processes since in superconducting quark matter the partial contribution $`1/\kappa _{eq}`$ (as well as $`1/\kappa _{gq}`$ ) is additionally suppressed by a $`\zeta _{\mathrm{QDU}}`$ factor, as for the scattering on impurities in metallic superconductors. For $`\kappa _{ee}`$ we have
$`\kappa _{ee}`$ $`=`$ $`5.5\times 10^{23}uY_eT_{9}^{}{}_{}{}^{1}\mathrm{erg}\mathrm{s}^1\mathrm{cm}^1\mathrm{K}^1,`$ (56)
and
$`\kappa _q\kappa _{qq}`$ $`1.1\times 10^{23}\sqrt{{\displaystyle \frac{4\pi }{\alpha _s}}}u\zeta _{\mathrm{QDU}}T_{9}^{}{}_{}{}^{1}\mathrm{erg}\mathrm{s}^1\mathrm{cm}^1\mathrm{K}^1,`$ (57)
where we take into account the suppression factor. We estimate the contribution of massless gluons as
$`\kappa _g\kappa _{gg}6.0\times 10^{17}T_9^2\mathrm{erg}\mathrm{s}^1\mathrm{cm}^1\mathrm{K}^1.`$ (58)
#### Heat capacity
The heat capacity contains nucleon, electron, photon, phonon, and other contributions. The main in-medium modification of the nucleon heat capacity is due to the density dependence of the effective nucleon mass. We use the same expressions as . The main regulators are the nucleon and the electron contributions. For the nucleons ($`i=n,p`$), the specific heat is
$$c_i1.6\times 10^{20}(m_i^{}/m_i)(n_i/n_0)^{1/3}\zeta _{ii}T_9\mathrm{erg}\mathrm{cm}^3\mathrm{K}^1,$$
(59)
and for the electrons it is
$$c_e6\times 10^{19}(n_e/n_0)^{2/3}T_9\mathrm{erg}\mathrm{cm}^3\mathrm{K}^1.$$
(60)
Near the phase transition point the heat capacity acquires a fluctuation contribution. For the first order pion condensation phase transition this additional contribution contains no singularity, unlike the second order phase transition, see . Finally, the nucleon contribution to the heat capacity may increase up to several times in the vicinity of the pion condensation point. The effect of this correction on global cooling properties is rather unimportant.
The symmetry of the $`3P_2`$ superfluid phase allows a Goldstone boson (phonon) contribution of
$`c_G610^{14}T_9^3\mathrm{erg}\mathrm{cm}^3\mathrm{K}^1,`$ (61)
for $`T<T_{cn}(3P_2)`$, $`n>n_{cn}(3P_2)`$. We include this contribution in our study too, although its effect on the cooling is rather minor.
For the quark specific heat we use the expression
$`c_q10^{21}u^{2/3}\zeta _\mathrm{S}T_9\mathrm{erg}\mathrm{cm}^3\mathrm{K}^1,`$ (62)
where $`\zeta _\mathrm{S}3.1(T_{\mathrm{crit},q}/T)^{5/2}\text{exp}(\mathrm{\Delta }_q/T)`$. Besides, one should add the gluon-photon contribution
$`c_{g\gamma }=3.0\times 10^{13}N_{g\gamma }T_9^3\mathrm{erg}\mathrm{cm}^3\mathrm{K}^1,`$ (63)
where $`N_{g\gamma }`$ is the number of available massless gluon-photon states (which are present even in the color superconducting phase), as well as the electron contribution,
$$c_e=5.7\times 10^{19}Y_e^{2/3}u^{2/3}T_9\mathrm{erg}\mathrm{cm}^3\mathrm{K}^1.$$
(64)
## 4 Results of simulations for Cooling evolution
We compute the neutron star thermal evolution adopting our fully general relativistic evolutionary code. The code originally constructed for the description of hybrid stars by has been developed and updated according to the modern knowledge of inputs in .
The density $`n0.5÷0.7n_0`$ is the boundary of the neutron star interior and the inner crust. The latter is constructed of a pasta phase discussed by , see also recent works of .
Further on we need a relation between the crust and the surface temperature for neutron star. A sharp change of the temperature occurs in the envelope. This $`T_\mathrm{s}T_{\mathrm{in}}`$ relation has been calculated in several works, see , depending on the assumed value of the magnetic field at the surface and some uncertainties in our knowledge of the structure of the envelope. For applications we use three different approximations: the simplified “Tsuruta law” $`T_\mathrm{s}^{\mathrm{Tsur}}=(10T_{\mathrm{in}})^{2/3}`$ used in many old cooling calculations, models used by Ref. , and our fit formula interpolating two extreme curves describing the borders of region of available $`T_\mathrm{s}T_{\mathrm{in}}`$ relation, taken from (see Fig. 4 of Ref. ).
### 4.1 Cooling Evolution of Hadronic Stars
Here we will shortly summarize the results on hadronic cooling.
In framework of a ”minimal cooling” scenario, the pair breaking and formation (PBF) processes may allow to cover an ”intermediate cooling” group of data (even if one artificially suppressed medium effects). These processes are very efficient for large pairing gaps and temperatures being not much less than the value of the gap.
Gaps, which we have adopted in the framework of the ”nuclear medium cooling” scenario, see , are presented in Fig. 6. Thick dashed lines show proton gaps which were used in the work of performed in the framework of the “standard plus exotics” scenario. We will call the choice of the “3nt” model from the model I. Thin lines show $`1S_0`$ proton and $`3P_2`$ neutron gaps from , for the model AV18 by (we call it the model II). Recently has argued for a strong suppression of the $`3P_2`$ neutron gaps, down to values $`10`$ keV, as the consequence of the medium-induced spin-orbit interaction.
These findings motivated to suppress values of $`3P_2`$ gaps shown in Fig. 6 by an extra factor $`f(3P_2,n)=0.1`$. Further possible suppression of the $`3P_2`$ gap is almost not reflected on the behavior of the cooling curves.
Contrary to expectations of a more recent work of argued that the $`3P_2`$ neutron pairing gap should be dramatically enhanced, as the consequence of the strong softening of the pion propagator. According to their estimate, the $`3P_2`$ neutron pairing gap is as large as $`1÷10`$ MeV in a broad region of densities, see Fig. 1 of their work. Thus results of calculations of and , which both had the same aim to include medium effects in the evaluation of the $`3P_2`$ neutron gaps, are in a deep discrepancy with each other.
* Including superfluid gaps we see, in agreement with recent microscopic findings of , that the $`3P_2`$ neutron gap should be as small as $`10`$ keV or less. So the “nuclear medium cooling” scenario of supports results of and fails to appropriately fit the neutron star cooling data when a strong enhancement of the $`3P_2`$ neutron gaps is assumed as shown by .
* Medium effects associated with the pion softening are called for by the data. As the result of the pion softening the pion condensation may occur for $`nn_c^{\mathrm{PU}}`$ ($`n3n_0`$ in our model). Its appearance at such rather high densities does not contradict to the cooling data (see Fig. 7), but also the data are well described using the pion softening but without assumption on the pion condensation. The similar argumentation holds also for DU threshold density. That puts restrictions on the density dependence of the symmetry energy. Both statements might be important in the discussion of the heavy ion collision experiments.
* We demonstrated a regular mass dependence for neutron stars with masses $`M>1M_{}`$, where less massive neutron stars cool down slower and more massive neutron stars cool faster. This feature is more general and concerning also to hybrid stars .
### 4.2 Cooling Evolution of Hybrid Stars with 2SC Quark Matter Core
For the calculation of the cooling of the quark core in the hybrid star we use the model . We incorporate the most efficient processes: the quark direct Urca (QDU) processes on unpaired quarks, the quark modified Urca (QMU), the quark bremsstrahlung (QB), the electron bremsstrahlung (EB), and the massive gluon-photon decay (see ). Following we include the emissivity of the quark pair formation and breaking (QPFB) processes, too. The specific heat incorporates the quark contribution, the electron contribution and the massless and massive gluon-photon contributions. The heat conductivity contains quark, electron and gluon terms.
The calculations are based on the hadronic cooling scenario presented in Fig. 7 and we add the contribution of the quark core. For the Gaussian form-factor the quark core occurs already for $`M>1.214M_{}`$ according to the model , see Fig. 5. Most of the relevant neutron star configurations (see Fig. 7) are then affected by the presence of the quark core.
First we check the possibility of the 2SC+ normal quark phases, see Fig. 8.
The variation of the gaps for the strong pairing of quarks within the 2SC phase and the gluon-photon mass in the interval $`m_{g\gamma }20÷200`$MeV only slightly affects the results. The main cooling process is the QDU process on normal quarks. We see that the presence of normal quarks entails too fast cooling. The data could be explained only, if all the masses lie in a very narrow interval ($`1.21<M/M_{}<1.22`$ in our case). In case of the other two crust models the resulting picture is similar.
The existence of only a very narrow mass interval, in which the data can be fitted seems to us unrealistic by itself. Moreover the observations show existence of neutron stars in binary systems with very different masses, e.g., $`M_{\mathrm{B1913}+16}1.4408\pm 0.0003M_{}`$ and $`M_{\mathrm{J0737}3039\mathrm{B}}1.250\pm 0.005M_{}`$, cf. . Thus the data can’t be satisfactorily explained.
Then we assume a possibility of an yet unknown X-pairing channel with a small $`\mathrm{\Delta }_X`$ gap and first check the case $`\mathrm{\Delta }_X`$ to be constant. For $`\mathrm{\Delta }_X1`$MeV the cooling is too slow . This is true for all three crust models. Thus the gaps for formerly unpaired quarks should be still smaller in order to obtain a satisfactory description of the cooling data.
For the $`\mathrm{\Delta }_X=30`$ keV the cooling data can be fitted but have a very fragile dependence on the gravitational mass of the configuration. Namely, we see that all data points, except the Vela, CTA 1 and Geminga, correspond to hybrid stars with masses in the narrow interval $`M=1.21÷1.22M_{}`$
Therefore we would like to explore whether a density-dependent X-gap could allow a description of the cooling data within a larger interval of compact star masses.
We employ an ansatz for the X-gap as a decreasing function of the chemical potential
$$\mathrm{\Delta }_X(\mu )=\mathrm{\Delta }_c\mathrm{exp}[\alpha (\mu \mu _c)/\mu _c],$$
(65)
where the parameters are chosen such that at the critical quark chemical potential $`\mu _c=330`$ MeV the onset of the deconfinement phase transition for the X-gap has its maximal value of $`\mathrm{\Delta }_c=1.0`$ MeV and at the highest attainable chemical potential $`\mu _{\mathrm{max}}=507`$ MeV, i.e. in the center of the maximum mass hybrid star configuration, it falls to a value of the order of $`10`$ keV. We choose the value $`\alpha =10`$ for which $`\mathrm{\Delta }_X(\mu _{\mathrm{max}})=4.6`$ keV. In Fig. 9 we show the resulting cooling curves for the gap model II with gap ansatz eq. (65), which we consider as the most realistic one.
We observe that the mass interval for compact stars which obeys the cooling data is constrained between $`M=1.32M_{}`$ for slow coolers and $`M=1.75M_{}`$ for fast coolers such as Vela. This results we obtain are based on a purely hadronic model with different choices of the parameters . Note that according to a recently suggested independent test of cooling models by comparing results of a corresponding population synthesis model with the Log N - Log S distribution of nearby isolated X-ray sources the cooling model I did not pass the test. Thereby it would be interesting to see, whether our quark model within the gap ansatz II could pass the Log N - Log S test.
## 5 Conclusions
In this lecture devoted to the modern problems of the cooling evolution of neutron stars we have discussed different successful scenarios aiming to explain the known temperature - age data from observations of compact objects. Using up today known theoretical and experimental constraints on the structure and cooling regulators of compact star we end up with alternative explanations. The neutron star can be either pure hadronic or hybrid one with a quark core in superconducting state of matter.
Particularly, for the hybrid stars, which are more intrigued alternative of superdense compact objects we conclude that:
* Within a nonlocal, chiral quark model the critical densities for a phase transition to color superconducting quark matter can be low enough for these phases to occur in compact star configurations with masses below $`1.3M_{}`$.
* For the choice of the Gaussian form-factor the 2SC quark matter phase arises at $`M1.21M_{}`$.
* Without a residual pairing the 2SC quark matter phase could describe the cooling data only if compact stars had masses in a very narrow band around the critical mass for which the quark core can occur.
* Under assumption that formally unpaired quarks can be paired with small gaps $`\mathrm{\Delta }_X<1`$MeV (2SC+X pairing), which values we varied in wide limits, only for density dependent gaps the cooling data can be appropriately fitted.
So the present day cooling data could be still explained by hybrid stars, assuming a complex pairing pattern, where quarks are partly strongly paired within the 2SC channel, and partly weakly paired with gaps $`\mathrm{\Delta }_X<1`$MeV, which are rapidly decrease with an increase of the density.
It remains to be investigated which microscopic pairing pattern could fulfill the constraints obtained in this work. Another indirect check of the model could be the Log N - Log S test.
### Acknowledgments
The research has been supported by the Virtual Institute of the Helmholtz Association under grant No. VH-VI-041 and by the DAAD partnership programm between the Universities of Yerevan and Rostock. In particular I acknowledge D. Blaschke for his active collaboration and support. I thank my colleagues D. Blaschke, D.N. Voskresensky, D.N. Aguilera, J. Berdermann, and A. Reichel for collaboration and discussions. I also thank the organizers of the organizers of the Helmholtz International Summer School and Workshop on Hot points in Astrophysics and Cosmology for their invitation to present these lectures.
———— |
warning/0506/math0506120.html | ar5iv | text | # Moduli Stacks of Polarized K3 Surfaces in Mixed Characteristic
## Introduction
In this note we will consider moduli spaces of K3 surfaces with a polarization. For a natural number $`d`$ and an algebraically closed field $`k`$, a K3 surface with a polarization of degree $`2d`$ over $`k`$ is a pair $`(X,)`$ consisting of a K3 surface $`X`$ over $`k`$ and an ample line bundle $``$ on $`X`$ with self intersection number $`(,)=2d`$. The moduli space of polarized K3 surfaces with certain level structure over $``$ is constructed as an open subspace of the Shimura variety associated with $`\mathrm{SO}(2,19)`$. Over $``$ we use techniques developed by Artin to show the existence of such spaces.
In various places in the literature one finds detailed accounts on coarse moduli schemes of primitively polarized complex K3 surfaces. We outline in Section 4.3 two approaches to the theory, one via geometric invariant theory (\[Vie95\]) and another via periods of complex K3 surfaces (\[BBD85, Exposé XIII\] and \[Fri84, §1\]). Here we take up a different point of view and work with moduli stacks rather than with coarse moduli schemes. In this way, our exposition is closer to \[Ols04\] where moduli stacks of primitively polarized K3 surfaces and their compactifictions over $``$ are constructed. We define the categories $`_{2d}`$ and $`_{2d}`$ of primitively polarized (respectively polarized) K3 surfaces of degree $`2d`$ over $``$ and show that they are Deligne-Mumford stacks over $``$.
For various technical reasons we will need to work with algebraic spaces rather than with Deligne-Mumford stacks. In the case of abelian varieties one introduces level $`n`$-structures using Tate modules and considers moduli functors of polarized abelian varieties with level $`n`$-structure for $`n,n3`$. These functors are representable by schemes. We adopt a similar strategy in order to define moduli functors which are representable by algebraic spaces. For a certain class of compact open subgroups $`𝕂`$ of $`\mathrm{SO}(2,19)(𝔸_f)`$ we introduce the notion of a level $`𝕂`$-structure on K3 surfaces using their second étale cohomology groups. Further, we introduce moduli spaces $`_{2d,𝕂}`$ of primitively polarized K3 surfaces with level $`𝕂`$-structure and show that these are smooth algebraic spaces over $`\mathrm{Spec}([1/N_𝕂])`$ where $`N_𝕂`$ depends on $`𝕂`$. These moduli spaces are finite unramified covers of $`_{2d}`$. Important examples of level structures are spin level $`n`$-structures. These are level structures defined by the images of some principal level $`n`$-subgroups of $`\mathrm{CSpin}(2,19)(𝔸_f)`$ under the adjoint representation homomorphism $`\mathrm{CSpin}(2,19)\mathrm{SO}(2,19)`$. We denote the corresponding moduli space by $`_{2d,n^{\mathrm{sp}}}`$.
Let us outline briefly the contents of this note. In the first few sections we review some basic properties of K3 surfaces. Then we continue with the study of the representability of Picard and automorphism functors arising from K3 surfaces. The core of the problems discussed here is Section 4.3 in which we define various moduli functors of polarized K3 surfaces and prove that those define Deligne-Mumford stacks. In Section 5.1 we define level structures on K3 surfaces associated to compact open subgroups of $`\mathrm{SO}(2,19)(𝔸_f)`$. In the last section we show that the moduli functors of primitively polarized K3 surfaces with level structure are representable by algebraic spaces.
Notations
We write $`\widehat{}`$ for the profinite completion of $``$. We denote by $`𝔸`$ the ring of adèles of $``$ and by $`𝔸_f=\widehat{}`$ the ring of finite adèles of $``$. Similarly, for a number field $`E`$ we denote by $`𝔸_E`$ and $`𝔸_{E,f}`$ the ring of adèles and the ring of finite adèles of $`E`$.
If $`A`$ is a ring, $`AB`$ a ring homomorphism then for any $`A`$-module ($`A`$-algebra etc.) $`V`$ we will denote by $`V_B`$ the $`B`$-module ($`B`$-algebra etc.) $`V_AB`$.
For a ring $`A`$ we denote by $`(\mathrm{Sch}/A)`$ the category of schemes over $`A`$. We will write $`\mathrm{Sch}`$ for the category of schemes over $``$.
By a variety over a field $`k`$ we will mean a separated, geometrically integral scheme of finite type over $`k`$. For a variety $`X`$ over $``$ we will denote by $`X^{\mathrm{an}}`$ the associated analytic variety. For an algebraic stack $``$ over a scheme $`S`$ and a morphism of schemes $`S^{}S`$ we will denote by $`_S^{}`$ the product $`\times _SS^{}`$ and consider it as an algebraic stack over $`S^{}`$.
A superscript <sup>0</sup> indicates a connected component for the Zariski topology. For an algebraic group $`G`$ will denote by $`G^0`$ the connected component of the identity. We will use the superscript <sup>+</sup> to denote connected components for other topologies.
Let $`V`$ be a vector space over $``$ and let $`G\mathrm{GL}(V)`$ be an algebraic group over $``$. Suppose given a full lattice $`L`$ in $`V`$ (i.e., $`L=V`$). Then $`G()`$ and $`G(\widehat{})`$ will denote the abstract groups consisting of the elements in $`G()`$ and $`G(𝔸_f)`$ preserving the lattices $`L`$ and $`L_\widehat{}`$ respectively.
Acknowledgments
This note contains the results of Chapter 1 of my Ph.D. thesis \[Riz05\]. I thank my advisors, Ben Moonen and Frans Oort for their help, their support and for everything I have learned from them. I would like to thank Bas Edixhoven and Gerard van der Geer for pointing out some mistakes and for their valuable suggestions. I thank the Dutch Organization for Research N.W.O. for the financial support with which my thesis was done.
## 1. Basic Results
### 1.1. Definitions and Examples
We will briefly recall some basic notions concerning families of K3 surfaces.
###### Definition 1.1.1.
Let $`k`$ be a field. A non-singular, proper surface $`X`$ over $`k`$ is called a *K3 surface* if $`\mathrm{\Omega }_{X/k}^2𝒪_X`$ and $`H^1(X,𝒪_X)=0`$.
Note that a K3 surface is automatically projective. Let us give some basic examples one can keep in mind:
###### Example 1.1.2.
Let $`S`$ be a non-singular sextic curve in $`_k^2`$ where $`k`$ is a field and consider a double cover i.e., a finite generically étale morphism, $`\pi :X_k^2`$ which is ramified along $`S`$. Then $`X`$ is a K3 surface.
###### Example 1.1.3.
Complete intersections: Let $`X`$ be a smooth surface which is a complete intersection of $`n`$ hypersurfaces of degree $`d_1,\mathrm{},d_n`$ in $`^{n+2}`$ over a field $`k`$. The adjunction formula shows that $`\mathrm{\Omega }_{X/k}^2𝒪_X(d_1+\mathrm{}+d_nn3)`$. So a necessary condition for $`X`$ to be a K3 surface is $`d_1+\mathrm{}+d_n=n+3`$. The first three possibilities are:
$$\begin{array}{cc}n=1& d_1=4\\ n=2& d_1=2,d_2=3\\ n=3& d_1=d_2=d_3=2.\end{array}$$
For a complete intersection $`M`$ of dimension $`n`$ one has that $`H^i(M,𝒪_M(m))=0`$ for all $`m`$ and $`1in1`$. Hence in those three cases we have $`H^1(X,𝒪_X)=0`$ and therefore $`X`$ is a K3 surface.
###### Example 1.1.4.
Let $`A`$ be an abelian surface over a field $`k`$ of characteristic different from 2. Let $`A[2]`$ be the kernel of the multiplication by-2-map, let $`\pi :\stackrel{~}{A}A`$ be the blow-up of $`A[2]`$ and let $`\stackrel{~}{E}`$ be the exceptional divisor. The automorphism $`[1]_A`$ lifts to an involution $`[1]_{\stackrel{~}{A}}`$ on $`\stackrel{~}{A}`$. Let $`X`$ be the quotient variety of $`\stackrel{~}{A}`$ by the group of automorphisms $`\{\mathrm{id}_{\stackrel{~}{A}},[1]_{\stackrel{~}{A}}\}`$ and denote by $`\iota :\stackrel{~}{A}X`$ the quotient morphism. It is a finite map of degree 2. We have the following diagram
of morphisms over $`k`$. The variety $`X`$ is a K3 surface and it is called the *Kummer surface* associated to $`A`$.
###### Definition 1.1.5.
By a *K3 scheme* over a base scheme $`S`$ we will mean a scheme $`X`$ and a proper and smooth morphism $`\pi :XS`$ whose geometric fibers are K3 surfaces. A *K3 space over a scheme $`S`$* is an algebraic space $`X`$ together with a proper and smooth morphism $`\pi :XS`$ such that there is an étale cover $`S^{}S`$ of $`S`$ for which $`\pi ^{}:X^{}=X\times _SS^{}S^{}`$ is a K3 scheme.
If $`\pi :XS`$ is a K3 space, then $`\pi _{}𝒪_X=𝒪_S`$. Indeed, this is true since $`\pi `$ is proper and its geometric fibers are reduced and connected.
###### Remark 1.1.6.
A K3 space $`X`$ over $`S`$ is usually defined as an algebraic space $`X`$ together with a proper and smooth morphism $`\pi :XS`$ such that for every geometric point $`sS`$ the fiber $`X_s`$ is a K3 surface. In this note we will restrict ourselves to Definition 1.1.5 above. The reason is that for this class of K3 spaces one can easily see that certain automorphism functors of K3 spaces are representable by schemes (cf. Theorem 3.3.1). We do not know if this holds in general.
### 1.2. Ample Line Bundles on K3 Surfaces
In order to construct the moduli stacks of polarized K3 spaces one needs a number of results on ample line bundles. We give them below.
###### Definition 1.2.1.
Let $`X`$ be a K3 surface over a field $`k`$. The self-intersection index $`(,)_X`$ of a line bundle $``$ on $`X`$ will be called its *degree*. A line bundle $``$ on $`X`$ is called *primitive* if $`\overline{k}`$ is is not a positive power of a line bundle on $`X_{\overline{k}}`$.
###### Theorem 1.2.2.
Let $`X`$ be a K3 surface over a field $`k`$.
1. If $``$ is a line bundle on $`X`$, then $`(,)`$ is even. If $``$ is ample and $`d:=(,)/2`$, then the Hilbert polynomial of $``$ is given by $`h_{}(t)=dt^2+2`$.
2. Suppose $``$ is an ample bundle. Then $``$ is effective and $`H^i(X,)=0`$ for $`i>0`$. Further, $`^n`$ is generated by global sections if $`n2`$ and is very ample if $`n3`$.
###### Proof.
(a) First note that, by Serre duality, $`h^2(𝒪_X)=h^0(\mathrm{\Omega }_{X/k}^2)=h^0(𝒪_X)=1`$. Since $`h^0(𝒪_X)=1`$ we find that $`\chi (𝒪_X)=2`$. Hirzebruch-Riemann-Roch gives
$$\begin{array}{cc}\hfill \chi ()& =\chi (𝒪_X)+\frac{1}{2}\left((,)(,\mathrm{\Omega }_{X/k}^2)\right)\hfill \\ & =2+\frac{1}{2}(,)\hfill \end{array}$$
as $`\mathrm{\Omega }_{X/k}^2`$ is trivial. Hence $`(,)=2d`$ is even. If $``$ is ample then its Hilbert polynomial is $`h_{}(t)=dt^2+2`$
(b) By Serre duality and the fact that $`\mathrm{\Omega }_{X/k}^2𝒪_X`$ we have $`h^i()=h^{2i}(^1)`$. In particular $`h^2()=h^0(^1)=0`$ as an anti-ample bundle is not effective. Since $`d:=(,)/2>0`$ it follows that $`h^0()=d+2+h^1()>0`$, so $``$ is effective. For the remaining assertions we refer to \[SD74\], Section 8. ∎
###### Example 1.2.3.
Let $`\pi :X^2`$ be a double cover of $`^2`$ as in Example 1.1.2. The line bundle $`=\pi ^{}𝒪_^2(1)`$ is ample and one has that $`(,)_X=2(𝒪_^2(1),𝒪_^2(1))_^2=2`$. Hence any K3 surface $`X`$ which is a double cover of $`^2`$ ramified along a non-singular sextic curve has an ample line bundle $``$ of degree 2.
###### Example 1.2.4.
Let $`X^{n+2}`$ be a K3 surface which is obtained as a complete intersection of multiple degree $`(d_1,d_2,\mathrm{},d_n)`$; see Example 1.1.3. Then $`𝒪_X(1)`$ degree $`d_1d_2\mathrm{}d_n`$. Note that the equality $`d_1+d_2+\mathrm{}+d_n=n+3`$ implies that at least one of the $`d_i`$ is even.
Note that if $`\pi :XS`$ is a K3 scheme over a connected base $`S`$ then for a line bundle $``$ on $`X`$ the intersection index $`(_{\overline{s}},_{\overline{s}})_{X_{\overline{s}}}`$ is constant for any $`\overline{s}`$. This follows from the fact that $`\pi `$ is flat and the relation $`(_{\overline{s}},_{\overline{s}})_{X_{\overline{s}}}=2\chi (_{\overline{s}})4`$.
###### Lemma 1.2.5.
Let $`\pi :XS`$ be a K3 scheme and let $``$ be a line bundle on $`X`$ which is fiberwise ample on $`X`$ i.e., $`_{\overline{s}}`$ is ample on $`X_{\overline{s}}`$ for every geometric point $`\overline{s}S`$. Let $`2d=(_{\overline{s}},_{\overline{s}})_{X_{\overline{s}}}`$ for any point $`\overline{s}S`$. Then $`\pi _{}^n`$ is a locally free sheaf of of rank $`dn^2+2`$ and $`^n`$ is relatively very ample over $`S`$ if $`n3`$.
###### Proof.
By Theorem 1.2.2 (b) we have that for all $`\overline{s}S`$ the group $`H^1(X_{\overline{s}},_{\overline{s}}^n)`$ is trivial. It follows from \[GD67, Ch. III, §7\], that $`\pi _{}^n`$ is a locally free sheaf and that $`\pi _{}_{\overline{s}}H^0(𝒳_{\overline{s}},_{\overline{s}}^n)`$. The rank statement follows from Theorem 1.2.2 (a). By part (a) of Theorem 1.2.2 one sees that for every geometric point $`\overline{s}S`$ and any $`n3`$ the line bundle $`_{\overline{s}}^n`$ gives a closed immersion $`X_{\overline{s}}(\pi _{}_{\overline{s}}^n)`$ over $`\kappa (\overline{s})`$. Hence the morphism $`X(\pi _{}^n)`$ induced by $`^n`$ is a closed immersion. This finishes the proof. ∎
## 2. Cohomology Groups of K3 Surfaces
### 2.1. Quadratic Lattices Related to Cohomology Groups of K3 Surfaces
In this section we introduce some notations which will be used in the sequel. Let $`U`$ be the hyperbolic plane and denote by $`E_8`$ the *positive* quadratic lattice associated to the Dynkin diagram of type $`E_8`$ (cf. \[Ser73, Ch. V, 1.4 Examples\]).
###### Notation 2.1.1.
Denote by $`(L_0,\psi )`$ the quadratic lattice $`U^3E_8^2`$. Further, let $`(V_0,\psi _0)`$ be the quadratic space $`(L_0,\psi )_{}`$.
We have that $`L_0`$ is a free $``$-module of rank 22. The form $`\psi _{}`$ has signature $`(19+,3)`$ on $`L_0`$.
Let $`\{e_1,f_1\}`$ be a basis of the first copy of $`U`$ in $`L_0`$ such that
$$\psi (e_1,e_1)=\psi (f_1,f_1)=0\text{and}\psi (e_1,f_1)=1.$$
For a positive integer $`d`$ we consider the vector $`e_1df_1`$ of $`L_0`$. It is a primitive vector i.e., the module $`L_0/e_1df_1`$ is free and we have that $`\psi (e_1df_1,e_1df_1)=2d`$. The orthogonal complement of $`e_1df_1`$ in $`L_0`$ with respect to $`\psi `$ is $`e_1+df_1U^2E_8^2`$.
###### Notation 2.1.2.
Denote the quadratic sublattice $`e_1+df_1U^2E_8^2`$ of $`L_0`$ by $`(L_{2d},\psi _{2d})`$. Further, we denote by $`(V_{2d},\psi _{2d})`$ the quadratic space $`(L_{2d},\psi _{2d})_{}`$.
The signature of the form $`\psi _{2d,}`$ is $`(19+,2)`$. We have that $`e_1df_1L_{2d}`$ is a sublattice of $`L_0`$ of index $`2d`$. The inclusion of lattices $`i:L_{2d}L_0`$ defines injective homomorphisms of groups
(1)
$$i^{\mathrm{ad}}:\{g\mathrm{O}(V_0)()|g(e_1df_1)=e_1df_1\}\mathrm{O}(V_{2d})()$$
and
(2)
$$i^{\mathrm{ad}}:\{g\mathrm{SO}(V_0)()|g(e_1df_1)=e_1df_1\}\mathrm{SO}(V_{2d})().$$
Let $`L_{2d}^{}`$ denote the dual lattice $`\mathrm{Hom}(L_{2d},)`$. Then the bilinear form $`\psi _{2d}`$ defines an embedding $`L_{2d}L_{2d}^{}`$ and we denote by $`A_{2d}`$ the factor group $`L_{2d}^{}/L_{2d}`$. It is an abelian group of order $`2d`$ (\[LP81, §2, Lemma\]). One can extend the bilinear form $`\psi _{2d}`$ on $`L_{2d}`$ to a $``$-valued form on $`L_{2d}^{}`$ and define
$$q_{2d}:A_{2d}/2$$
defined by
$$q_{2d}(x+L_{2d})=\psi _{2d}(x,x)+2$$
for any $`xL_{2d}^{}`$. Let $`\mathrm{O}(q_{2d})`$ denote the group of isomorphisms of $`A_{2d}`$ preserving the form $`q_{2d}`$. Then one has a natural homomorphism $`\tau :\mathrm{O}(V_{2d})()\mathrm{O}(q_{2d})`$. It is shown in \[Nik80\] that
$$i^{\mathrm{ad}}\left(\{g\mathrm{O}(V_0)()|g(e_1df_1)=e_1df_1\}\right)=\mathrm{ker}(\tau ).$$
### 2.2. De Rham Cohomology
Let $`X`$ be a K3 surface over a field $`k`$. The following proposition will play an essential role when studying deformations of K3 surfaces (Section 4.1). We will use it also to show that the automorphism group $`\mathrm{Aut}(X)`$ of a K3 surface is reduced (see Theorem 3.3.1 below).
###### Proposition 2.2.1.
If $`X`$ is a K3 surface over a field $`k`$, then
1. The Hodge-de Rham spectral sequence
$$E_1^{i,j}=H^j(X,\mathrm{\Omega }_{X/k}^i)H_{DR}^{i+j}(X,k)$$
degenerates at $`E_1`$. For the Hodge numbers $`h^{i,j}=dim_kH^j(X,\mathrm{\Omega }_{X/k}^i)`$ of $`X`$ we have
$$h^{1,0}=h^{0,1}=h^{2,1}=h^{1,2}=0$$
$$h^{0,0}=h^{2,0}=h^{0,2}=h^{2,2}=1$$
$$h^{1,1}=20.$$
2. Let $`\mathrm{\Theta }_{X/k}=\mathrm{\Omega }_{X/k}^1`$ be the tangent bundle of $`X`$. Then $`H^i(X,\mathrm{\Theta }_{X/k})=0`$ for $`i=0`$ and $`2`$ and $`dim_kH^1(X,\mathrm{\Theta }_{X/k})=20`$.
###### Proof.
If $`k`$ has characteristic zero, then one may assume that $`k=`$ and the proposition follows from \[LP81, §1, Prop. 1.2\]. The case $`\mathrm{char}(k)=p>0`$ is treated in \[Del81b, Prop. 1.1\]. ∎
###### Remark 2.2.2.
Part (b) of the proposition is classical in the case $`k=`$. The proof in the general case is due to Rudakov and Shafarevich. It can be reformulated in following way: There exist no non-trivial regular vector fields on a K3 surface (cf. \[RS76, §6, Thm. 7\]).
### 2.3. Betti Cohomology
Let $`X`$ be a complex K3 surface. Then the Betti cohomology groups $`H_B^i(X,)`$ are free $``$-modules of rank $`1,0,22,0,1`$ for $`i=0,1,2,3,4`$ respectively. One has a non-degenerate bilinear form (given by the Poincaré duality pairing):
$$\psi :H_B^2(X,)(1)\times H_B^2(X,)(1)$$
given by
$$\psi (x,y)=\mathrm{tr}(xy)$$
where $`xy`$ is the cup product of $`x`$ and $`y`$ and $`\mathrm{tr}:H_B^4(X,(2))`$ is the trace map. It has signature $`(19+,3)`$ over $``$. The quadratic lattice $`(H_B^2(X,)(1),\psi )`$ is isometric to $`(L_0,\psi )`$ (cf. Section 2.1). For proofs of those results we refer to \[LP81, §1, Prop. 1.2\].
The group $`H_B^2(X,)`$ carries a natural $``$-Hodge structure (which we will abbreviate as $``$-HS) of type $`\{(2,0),(1,1),(0,2)\}`$ with $`h^{2,0}=h^{0,2}=1`$ and $`h^{1,1}=20`$ as we see from Proposition 2.2.1.
For a complex K3 surface $`H^1(X,𝒪_X)`$ is trivial so the first Chern class map
$$c_1:\mathrm{Pic}(X)H_B^2(X,)(1)$$
is injective. Exactly in the same way we see that for a K3 space $`\pi :XS`$, where $`S`$ is a scheme over $``$, one has a short exact sequence of sheaves
$$0R^1\pi _{}^{\mathrm{an}}𝒪_X^{}R^2\pi _{}^{\mathrm{an}}(1)$$
as $`R^1\pi _{}^{\mathrm{an}}𝒪_X`$ is trivial.
###### Notation 2.3.1.
Let $``$ be an ample line bundle on $`X`$. We denote by $`P_B^2(X,)(1)`$ the orthogonal complement of $`c_1()`$ with respect to $`\psi `$. It is a free $``$-module of rank 21 called the *primitive part* (or the *primitive cohomology group*) of $`H_B^2(X,)(1)`$ with respect to $`c_1()`$. The restriction of $`\psi `$ defines a non-degenerate bilinear form:
$$\psi _{}:P_B^2(X,)(1)\times P_B^2(X.)(1).$$
The group $`P_B^2(X,(1))`$ carries a natural $``$-HS induced by the one on $`H_B^2(X,(1))`$ of type $`\{(1,1),(0,0),(1,1)\}`$ with $`h^{1,1}=h^{1,1}=1`$ for which $`\psi _{}`$ is a polarization.
###### Remark 2.3.2.
Let $``$ be an ample line bundle for which $`(,)_X=2d`$ and assume that it is primitive. Let $`\{e_1,f_1\}`$ be a basis of the first copy of $`U`$ in $`L_0`$ as in Section 2.1. By \[BBD85, Exp. IX, §1, Prop. 1\] one can find an isometry
$$a:(H_B^2(X,(1)),\psi )L_0$$
such that $`a(c_1())=e_1df_1`$. Therefore $`a`$ induces an isometry
$$a:(P_B^2(X,(1)),\psi _{})(L_{2d},\psi _{2d}).$$
### 2.4. Étale Cohomology
Let $`k`$ be a field of characteristic $`p0`$ and fix a prime $`l`$ which is different from $`p`$. Suppose given a K3 surface $`X`$ over $`k`$. Then the étale cohomology group $`H_{\mathrm{et}}^i(X_{\overline{k}},_l)`$ is a free $`_l`$-module of rank $`1,0,22,0,1`$ for $`i=0,1,2,3,4`$. One sees this in the following way: If $`k`$ has characteristic zero, then the claim follows from the corresponding result for Betti cohomology and the comparison theorem between Betti and étale cohomology (\[Mil80, Ch. III, §3, Thm. 3.12\]). Assume that $`p>0`$. By \[Del81b, §1, Cor. 1.8\] there exists a discrete valuation ring $`R`$ with residue field $`\overline{k}`$ and a smooth lift $`𝒳`$ over $`R`$ of $`X`$. If $`\eta `$ is the generic point of $`\mathrm{Spec}(R)`$, then by the smooth base change theorem for étale cohomology (\[Mil80, Ch. VI, §4, Cor. 4.2\]) one has that
(3)
$$H_{\mathrm{et}}^i(X_{\overline{k}},/l^n)H_{\mathrm{et}}^i(𝒳_{\overline{\eta }},/l^n)$$
for every $`i=0,\mathrm{},4`$ and every $`n`$. Hence $`H_{\mathrm{et}}^i(X_{\overline{k}},_l)H_{\mathrm{et}}^i(𝒳_{\overline{\eta }},_l)`$ and we deduce the claim from the characteristic zero result.
Further, one has a non-degenerate bilinear form
$$\psi __l:H_{\mathrm{et}}^2(X_{\overline{k}},_l)(1)\times H_{\mathrm{et}}^2(X_{\overline{k}},_l)(1)_l$$
given by
$$\psi __l(x,y)=\text{tr}__l(xy)$$
where $`\text{tr}__l:H_{\mathrm{et}}^4(X_{\overline{k}},_l)(2)_l`$ is the trace isomorphism. This is simply Poincaré duality for étale cohomology (\[Mil80, Ch. VI, §11, Cor. 11.2\]).
The Kummer short exact sequence of étale sheaves on $`X`$
$$1𝝁_{l^n}𝔾_m𝔾_m1$$
gives an exact sequence of cohomology groups
$$H_{\mathrm{et}}^1(X_{\overline{k}},𝝁_{l^n})H_{\mathrm{et}}^1(X_{\overline{k}},𝔾_m)H_{\mathrm{et}}^1(X_{\overline{k}},𝔾_m)H_{\mathrm{et}}^2(X_{\overline{k}},𝝁_{l^n}).$$
By (3) the group $`H_{\mathrm{et}}^1(X_{\overline{k}},𝝁_{l^n})`$ is trivial we have an injection
$$0\mathrm{Pic}(X)/l^n\mathrm{Pic}(X)H_{\mathrm{et}}^2(X_{\overline{k}},𝝁_{l^n}).$$
Taking the projective limit over $`n`$ one sees that the first Chern class map
$$c_1:\mathrm{Pic}(X)_{}_lH_{\mathrm{et}}^2(X_{\overline{k}},_l)(1)$$
is injective. In particular, since $`H_{\mathrm{et}}^2(X_{\overline{k}},_l(1))`$ is free, $`\mathrm{Pic}(X)`$ has no $`l`$-torsion for any $`l`$ different from $`p`$.
Similarly, if $`\pi :XS`$ is a K3 space then one can consider the long exact sequence of higher direct images, coming from the Kummer sequence
$$R_{\mathrm{et}}^1\pi _{}𝝁_{l^n}R_{\mathrm{et}}^1\pi _{}𝔾_mR_{\mathrm{et}}^1\pi _{}𝔾_mR_{\mathrm{et}}^2\pi _{}𝝁_{l^n}.$$
Further, since the stalk of $`R_{\mathrm{et}}^1\pi _{}𝝁_{l^n}`$ at any geometric point of $`S`$ is zero (one uses here the proper base change theorem), the sheaf itself is zero (\[Mil80, Ch. II, §2, Prop. 2.10\]). Hence passing again to the projective limit over $`n`$ we obtain the exact sequence of $`_l`$-sheaves
$$0R_{\mathrm{et}}^1\pi _{}𝔾_m_lR_{\mathrm{et}}^2\pi _{}_l(1).$$
###### Notation 2.4.1.
Let $``$ be a primitive ample line bundle on $`X`$ with $`(,)_X=2d`$. Denote by $`P_{\mathrm{et}}^2(X_{\overline{k}},_l(1))`$ the *primitive part* of $`H_{\mathrm{et}}^2(X_{\overline{k}},_l)(1)`$ with respect to $`c_1()`$ i.e., the orthogonal complement of $`c_1()`$ in $`H_{\mathrm{et}}^2(X_{\overline{k}},_l)(1)`$ with respect to $`\psi __l`$. Denote the restriction of $`\psi __l`$ to $`P_{\mathrm{et}}^2(X_{\overline{k}},_l(1))`$ by $`\psi _{,_l}`$.
If $`k`$ has characteristic 0, then by the comparison theorem between Betti and étale cohomology one has that $`(H_{\mathrm{et}}^2(X_{\overline{k}},_l(1)),\psi __l)`$ is isometric to $`(H_B^2(X_{},(1)),\psi )_{}_l`$ which is isometric to $`(L_0,\psi )_{}_l`$. Moreover since the comparison isomorphism respects algebraic cycles, the same holds for the primitive parts with respect to $``$ i.e., we have that $`(P_{\mathrm{et}}^2(X_{\overline{k}},_l(1)),\psi _{,_l})(L_{2d},\psi _{2d})_{}_l`$.
Assume that $`\mathrm{char}(k)=p>0`$. Then the pair $`(X,)\overline{k}`$ has a lift $`(𝒳,)`$ over a discrete valuation ring $`R`$ with $`\mathrm{char}(R)=0`$ and with residue field $`\overline{k}`$ (see \[Del81b, §1, Cor. 1.8\]). Using the same argument as above one concludes that
$$H_{\mathrm{et}}^i(X_{\overline{k}},_l)(m)H_{\mathrm{et}}^i(𝒳_{\overline{\eta }},_l)(m)$$
and that $`(H_{\mathrm{et}}^2(X_{\overline{k}},_l)(m),\psi __l)`$ is isometric to $`(H_{\mathrm{et}}^2(𝒳_{\overline{\eta }},_l)(m),\psi __l)`$, where $`\eta `$ is the generic point of $`\mathrm{Spec}(R)`$. Consequently the two quadratic lattices $`(P_{\mathrm{et}}^2(X_{\overline{k}},_l(1)),\psi _{,_l})`$ and $`(P_{\mathrm{et}}^2(𝒳_{\overline{\eta }},_l(1)),\psi _{_{\overline{\eta }}})_{}_l`$ are also isometric. Thus, if $``$ is primitive, then there is an isometry
$$a:(H_{\mathrm{et}}^2(X_{\overline{k}},_l(1)),\psi _l)L_0_l$$
such that $`a(c_1())=e_1df_1`$. It induces an isometry
$$a:(P_{\mathrm{et}}^2(X_{\overline{k}},_l(1)),\psi _{,_l})(L_{2d},\psi _{2d})_l.$$
###### Remark 2.4.2.
Let $`k`$ be a field of characteristic $`p`$. We make the following notations
$$\widehat{}^{(p)}:=\underset{lp}{}_l\text{and}𝔸_f^{(p)}=\widehat{}^{(p)}.$$
In the sequel we will be considering étale cohomology with $`\widehat{}^{(p)}`$ or $`𝔸_f^{(p)}`$ coefficients. Then we have that for a K3 surface over a field $`k`$ one has isometries
$$(H_{\mathrm{et}}^2(X_{\overline{k}},\widehat{}^{(p)}(1)),\psi _f)(L_0,\psi )_{}\widehat{}^{(p)}$$
and for a primitive ample line bundle $``$ of degree $`2d`$ on $`X`$ one has
$$(P_{\mathrm{et}}^2(X_{\overline{k}},\widehat{}^{(p)}(1)),\psi _{,f})(L_{2d},\psi _{2d})_{}\widehat{}^{(p)}.$$
Here $`\psi _f`$ and $`\psi _{,f}`$ are the corresponding bilinear forms coming from the Poincaré duality on $`H_{\mathrm{et}}^2(X_{\overline{k}},\widehat{}^{(p)}(1))`$.
### 2.5. Crystalline Cohomology
Let $`k`$ be a perfect field of characteristic $`p>0`$ and let $`W=W(k)`$ be the ring of Witt vectors with coefficients in $`k`$. Consider a K3 surface $`X`$ over $`k`$. Then by \[Del81b, Prop. 1.1\] the crystalline cohomology group $`H_{\mathrm{cris}}^i(X/W)`$ is a free $`W`$-module of rank 1, 0, 22, 0, 1 for $`i=`$ 0, 1, 2, 3, 4 respectively. We consider next the crystalline Chern class map
$$c_1:\mathrm{Pic}(X)H_{\mathrm{cris}}^2(X/W).$$
As pointed out in \[Del81b, Appendice, Rem. 3.5\] the Chern class map defines an injection
$$c_1:\mathrm{NS}(X_{\overline{k}})_{}_pH_{\mathrm{cris}}^2(X/W(\overline{k}))$$
where $`\mathrm{NS}(X_{\overline{k}})=\mathrm{Pic}(X_{\overline{k}})/\mathrm{Pic}^0(X_{\overline{k}})`$ is the Néron-Severi group of $`X_{\overline{k}}`$. In particular this means that the Néron-Severi group of $`X_{\overline{k}}`$ has no $`p`$-torsion.
If $`K`$ is the fraction field of $`W`$ then we shall denote by $`H_{\mathrm{cris}}^i(X/K)`$ the $`K`$-vector space $`H_{\mathrm{cris}}^i(X/W)_WK`$.
## 3. Picard Schemes and Automorphisms of K3 Surfaces
### 3.1. Picard and Néron-Severi Groups of K3 Surfaces.
In this section we will study Picard functors of K3 spaces. Those functors will play an important role in two aspects in the construction of moduli spaces of (primitively) polarized K3 surfaces. First, we will define (quasi-) polarizations on K3 surfaces using Picard spaces (cf. Definition 3.2.2 below). Later, in Section 4.2, we will use Picard spaces in the construction of the Hilbert scheme parameterizing K3 subschemes of $`^N`$.
For a separated algebraic space $`X`$ over a scheme $`S`$ we denote by $`\mathrm{Pic}(X)`$ the group of isomorphism classes of invertible sheaves on $`X`$. Let $`\pi :XS`$ be a K3 space and consider the *relative Picard functor*
$$\mathrm{Pic}_{X/S}:(\mathrm{Sch}/S)^0\mathrm{Groups}.$$
By definition it is the fppf-sheafification of the functor
$$P_{X/S}:(\mathrm{Sch}/S)^0\mathrm{Groups}\text{given by}T\mathrm{Pic}(X\times _ST).$$
For every $`g:TS`$ we have that $`\mathrm{Pic}_{X/S}(T)=H^0(T,R^1\pi _{}^{}𝔾_m)`$ where $`\pi ^{}:X\times _STT`$ is the product morphism and all derived functors are taken with respect to the fppf-topology.
###### Theorem 3.1.1.
For a K3 space $`\pi :XS`$ the relative Picard functor $`\mathrm{Pic}_{X/S}`$ is represented by a separated algebraic space locally of finite presentation over $`S`$.
###### Proof.
The representability follows form \[Art69, §7, Thm. 7.3\]. The proof of the separatedness property goes exactly in the same way as the proof of Theorem 3 in \[BLR90, Ch. 8, §8.4\]. ∎
Let $`S=\mathrm{Spec}(k)`$ be a spectrum of a field. Then $`\mathrm{Pic}_{X/k}`$ is represented by a group scheme (cf. \[Oor62\] or Lemma 3.1.2 below) and shall denote by $`\mathrm{Pic}_{X/k}^0`$ its identity component. We set further
$$\mathrm{Pic}_{X/k}^\tau =\underset{n>0}{}n^1\left(\mathrm{Pic}_{X/k}^0\right)$$
where $`n:\mathrm{Pic}_{X/k}\mathrm{Pic}_{X/k}`$ is the multiplication by $`n`$.
###### Lemma 3.1.2.
Let $`X`$ be a K3 surface over a field $`k`$. Then $`\mathrm{Pic}_{X/k}`$ is represented by a separated, smooth, zero dimensional scheme over $`k`$. In particular $`\mathrm{Pic}_{X/k}^0`$ is trivial. Further, we have also that $`\mathrm{Pic}_{X/k}^\tau `$ is trivial.
###### Proof.
Combining Theorem 3 and Theorem 1, with $`S=\mathrm{Spec}(k)`$, of \[BLR90, Ch. 8, §8.2\] one concludes that $`\mathrm{Pic}_{X/k}`$ is representable by a separated scheme, locally of finite type over $`k`$.
By Theorem 1 of \[BLR90, Ch. 8, §8.4\] one has that
$$dim_k\mathrm{Pic}_{X/k}dim_kH^1(X,𝒪_X)=0$$
and as the equality holds in this case, $`\mathrm{Pic}_{X/k}`$ is smooth over $`k`$. This shows the validity of all assertions except for the claim about $`\mathrm{Pic}_{X/k}^\tau `$.
The scheme $`\mathrm{Pic}_{X/k}^\tau `$ is proper and of finite type over $`k`$ (cf. \[BLR90, Ch. 8, Thm. 4\]). Since its dimension is zero it is a finite commutative group scheme over $`k`$. The injectivity of the étale Chern class map shows that $`\mathrm{Pic}(X)`$ has no $`l`$-torsion for $`lp`$. By the first part of the lemma we have that $`\mathrm{NS}(X)=\mathrm{Pic}(X)`$. Then the injectivity of the crystalline Chern class map shows that $`\mathrm{Pic}(X)`$ has no $`p`$-torsion either. Thus $`\mathrm{Pic}(X)`$ is torsion free and therefore $`\mathrm{Pic}_{X/k}^\tau (\overline{k})`$ is trivial. Since in this case $`\mathrm{Pic}_{X/k}^\tau `$ is reduced we conclude it is trivial. ∎
If $`X`$ is a K3 surface over a field $`k`$, then $`\mathrm{NS}(X)=\mathrm{Pic}(X)`$, which follows from the fact that in this case $`\mathrm{Pic}^0(X)`$ is trivial. Hence $`\mathrm{Pic}(X)`$ is a free abelian group of rank at most $`22`$ (use \[Mil80, Ch. V, §3, Cor 3.28\]). If the characteristic of the ground field is zero, then $`\mathrm{rk}_{}\mathrm{Pic}(X)20`$.
Let $`\pi :XS`$ be a K3 scheme. Define $`\mathrm{Pic}_{X/S}^0`$ and $`\mathrm{Pic}_{X/S}^\tau `$ as the subfunctors of $`\mathrm{Pic}_{X/S}`$ consisting of all elements whose restrictions to all fibers $`X_s`$ belong to $`\mathrm{Pic}_{X_s/\kappa (s)}^0`$ and $`\mathrm{Pic}_{X_s/\kappa (s)}^\tau `$ respectively.
###### Proposition 3.1.3.
For a K3 scheme $`\pi :XS`$ over a quasi-compact base $`S`$ one has that $`\mathrm{Pic}_{X/S}`$ is an algebraic space which is unramified over $`S`$. Further, we have that $`\mathrm{Pic}_{X/S}^0`$ and $`\mathrm{Pic}_{X/S}^\tau `$ are trivial.
###### Proof.
The first part of the proposition follows from the preceding lemma as it is enough to check the $`\mathrm{Pic}_{X/S}`$ is unramified in the case $`S`$ is a spectrum of a field. To prove the second part we notice that according to \[BLR90, Ch. 8, §8.3, Thm. 4\] we have open immersions $`\mathrm{Pic}_{X/S}^0\mathrm{Pic}_{X/S}`$ and $`\mathrm{Pic}_{X/S}^\tau \mathrm{Pic}_{X/S}`$. By Lemma 3.1.2 above for every geometric point $`\overline{s}S`$ the subspaces $`\mathrm{Pic}_{X_{\overline{s}}/\kappa (\overline{s})}^0`$ and $`\mathrm{Pic}_{X_{\overline{s}}/\kappa (\overline{s})}^\tau `$ are trivial hence $`\mathrm{Pic}_{X/S}^0`$ and $`\mathrm{Pic}_{X/S}^\tau `$ are trivial. ∎
###### Remark 3.1.4.
Let $`\pi :XS`$ be a K3 scheme and let $``$ and $``$ be two line bundles on $`X`$. If $`^n=^n`$ for some $`n`$, then $``$ is isomorphic to $`\pi ^{}𝒩`$ where $`𝒩`$ is a line bundle on $`S`$. Indeed, we have that $`cl()^n=cl()^n`$ in $`\mathrm{Pic}_{X/S}`$. Since $`\mathrm{Pic}_{X/S}^\tau `$ is trivial we have that the multiplication by $`n`$-morphism $`[n]:\mathrm{Pic}_{X/S}\mathrm{Pic}_{X/S}`$ is an injective homomorphism of group schemes. Since $`cl()^n=cl()^n`$ we conclude that $`cl(^1)`$ is trivial, so $``$ and $``$ differ by an invertible sheaf coming from the base $`S`$ (\[BLR90, Ch. 8, §8.1, Prop. 4\]).
###### Remark 3.1.5.
It is easy to see that the statement of Proposition 3.1.3 remains true for K3 spaces.
A morphism of schemes $`\pi :XS`$ is called *strongly projective* (respectively *strongly quasi-projective*) if there exists a locally free sheaf $``$ on $`S`$ of constant finite rank such that $`X`$ is $`S`$-isomorphic to a closed subscheme (respectively a subscheme) of $`()`$.
###### Lemma 3.1.6.
Let $`S`$ be a noetherian scheme and suppose given a K3 scheme $`\pi :XS`$. If $`\pi `$ is a strongly projective morphism, then we have that
1. for any $`n`$ the multiplication by $`n`$-morphism
$$[n]:\mathrm{Pic}_{X/S}\mathrm{Pic}_{X/S}$$
is a closed immersion of group schemes over $`S`$.
2. for any $`\lambda \mathrm{Pic}_{X/S}(S)`$ the set of points
$$S^o=\{sS|\lambda _s\mathrm{is}\mathrm{primitive}\mathrm{on}X_s\}$$
is open in $`S`$.
###### Proof.
(i): By definition we have a closed immersion $`X()`$ for some locally free sheaf $``$ on $`S`$. Let $`𝒪_X(1)`$ denote the pull-back of the canonical bundle $`𝒪(1)`$ on $`()`$ via this inclusion. For a polynomial $`\mathrm{\Phi }[t]`$ let $`\mathrm{Pic}_{X/S}^\mathrm{\Phi }`$ be the subfunctor of $`\mathrm{Pic}_{X/S}`$ which is induced by the line bundles $``$ on $`X`$ with a given Hilbert polynomial $`\mathrm{\Phi }`$ (with respect to $`𝒪_X(1)`$) on the fibers of $`X`$ over $`S`$. Then $`\mathrm{Pic}_{X/S}^\mathrm{\Phi }`$ is representable by a strongly quasi-projective scheme over $`S`$ and $`\mathrm{Pic}_{X/S}`$ is the disjoint union of the open and closed subschemes $`\mathrm{Pic}_{X/S}^\mathrm{\Phi }`$ for all $`\mathrm{\Phi }[t]`$. For a proof of this result we refer to \[BLR90, Ch. 8, §8.2, Thm. 5\].
Since all schemes $`\mathrm{Pic}_{X/S}^\mathrm{\Phi }`$ are quasi-compact we have that for a given $`\mathrm{\Phi }`$ the image $`[n](\mathrm{Pic}_{X/S}^\mathrm{\Phi })`$ is contained in a finite union $`_{i𝒞_\mathrm{\Phi }^n}\mathrm{Pic}_{X/S}^{\mathrm{\Phi }_i}`$. We will show first that for a given $`\mathrm{\Phi }[t]`$ the morphism
$$[n]:\mathrm{Pic}_{X/S}^\mathrm{\Phi }\underset{i𝒞_\mathrm{\Phi }^n}{}\mathrm{Pic}_{X/S}^{\mathrm{\Phi }_i}$$
is proper. As all schemes involved are noetherian we can apply the valuative criterion for properness. We may assume that $`S`$ is a spectrum of a discrete valuation ring $`R`$ and that $`X`$ admits a section over $`S`$ and let $`\eta `$ and $`s`$ be the generic and the special point of $`S`$. Under those assumptions any element of $`\mathrm{Pic}_{X/S}`$ comes from a class of a line bundle (\[BLR90, Ch. 8, §8.1, Prop. 4\]). To show that the restriction of $`[n]`$ to $`\mathrm{Pic}_{X/S}^\mathrm{\Phi }`$ is proper we have to show that if $``$ is a line bundle over the generic fiber $`X_\eta `$ of $`X`$, then $`^n`$ extends uniquely to a line bundle on $`X`$ which is a $`n`$-th power of a line bundle. This follows from \[BLR90, Ch. 8, §8.4, Thm. 3\] as both $``$ and $`^n`$ extend uniquely over $`X`$.
Further, the morphism $`[n]:\mathrm{Pic}_{X/S}\mathrm{Pic}_{X/S}`$ is an immersion of the corresponding topological spaces and as it is proper on every open and closed $`\mathrm{Pic}_{X/S}^\mathrm{\Phi }`$, the image $`[n](\mathrm{Pic}_{X/S})`$ is closed in $`\mathrm{Pic}_{X/S}`$. We are left to show that the natural homomorphism of sheaves $`𝒪_{\mathrm{Pic}_{X/S}}[n]_{}𝒪_{\mathrm{Pic}_{X/S}}`$ is surjective. As this can be checked on stalks we see further that it is enough to show the surjectivity assuming that $`S`$ is a spectrum of a field. But under this condition the claim follows from Lemma 3.1.2. Indeed, $`\mathrm{Pic}_{X/k}`$ is a reduced, zero dimensional scheme. Hence all subschemes $`\mathrm{Pic}_{X/k}^\mathrm{\Phi }`$ being reduced, quasi-projective and zero dimensional, are finite unions of points. Then the restrictions $`[n]:\mathrm{Pic}_{X/k}^\mathrm{\Phi }_{i𝒞_\mathrm{\Phi }^n}\mathrm{Pic}_{X/k}^{\mathrm{\Phi }_i}`$ are closed immersions and hence $`[n]:\mathrm{Pic}_{X/k}\mathrm{Pic}_{X/k}`$ is also a closed immersion. Therefore $`𝒪_{\mathrm{Pic}_{X/k}}[n]_{}𝒪_{\mathrm{Pic}_{X/k}}`$ is surjective.
(ii): We may assume that $`S`$ is connected. Then the intersection index $`(\lambda _{\overline{s}},\lambda _{\overline{s}})`$ is constant on $`S`$, say $`(\lambda _{\overline{s}},\lambda _{\overline{s}})=2d`$. For any natural number $`n`$ consider the closed subscheme $`S_n`$ of $`S`$ defined by the following Cartesian diagram
Then the subset $`S^o`$ of $`S`$ can be identified with $`S_nS_n`$ where the union is taken over all $`n`$ such that $`n^2`$ divides $`d`$. So it has a structure of an open subscheme of $`S`$. ∎
###### Remark 3.1.7.
Note that if $`\pi :XS`$ is a K3 scheme, then the Picard functor $`\mathrm{Pic}_{X/S}`$ can be constructed using the étale topology on $`S`$ instead of the fppf-topology. In other words $`\mathrm{Pic}_{X/S}`$ is also the étale sheafification of $`P_{X/S}`$. This follows from the fact that $`\pi `$ is a proper morphism, using the Leray spectral sequence for $`\pi `$ and the sheaf $`𝔾_m`$. For a proof we refer to the comments on p. 203 in \[BLR90, Ch. 8, §8.1\].
###### Example 3.1.8.
Let $`A`$ be an abelian surface over an algebraically closed field $`k`$ of characteristic different from 2 and let $`X`$ be the associated Kummer surface. Then one has that
$$\mathrm{Pic}(X)_{}=\mathrm{NS}(X)_{}\mathrm{NS}(A)_{}^{[1]_A}^{16}$$
where $`\mathrm{NS}(A)^{[1]_A}`$ denotes the elements of $`\mathrm{NS}(A)`$ invariant under the action of $`[1]_A`$. We refer to \[Shi79, §3, Prop. 3.1\] for a proof.
### 3.2. Polarizations of K3 Surfaces
Here we will define the notion of a polarization on a K3 space.
###### Definition 3.2.1.
Let $`k`$ be a field. A *polarization* on a K3 surface $`X/k`$ is a global section $`\lambda \mathrm{Pic}_{X/k}(k)`$ which over $`\overline{k}`$ is the class of an ample line bundle $`_{\overline{k}}`$. The degree of $`_{\overline{k}}`$ is called the *polarization degree* of $`\lambda `$. A *quasi-polarization* on $`X`$ is a global section $`\lambda \mathrm{Pic}_{X/k}(k)`$ which over $`\overline{k}`$ comes from a line bundle $`_{\overline{k}}`$ with the following property:
1. $`_{\overline{k}}`$ is nef i.e., $`(_{\overline{k}},𝒪_{X_{\overline{k}}}(C))0`$ for all irreducible curves in $`X_{\overline{k}}`$,
2. if $`(_{\overline{k}},𝒪_{X_{\overline{k}}}(C))=0`$ for a curve $`C`$ in $`X_{\overline{k}}`$ then $`(C,C)_{X_{\overline{k}}}=(2)`$.
If $`(X,\lambda )`$ is a polarized K3 surface over $`k`$, then one can find a finite separable extension $`k^{}`$ of $`k`$ such that $`\lambda `$ comes from a line bundle $`_k^{}`$ over $`k^{}`$. Indeed, this follows either from Remark 3.1.7 or from Proposition 4 in \[BLR90, Ch. 8, §8.1\] taking $`T=\mathrm{Spec}(k^{\mathrm{sp}})`$ and the fact that $`\text{Br}(k^{\mathrm{sp}})`$ is trivial.
###### Definition 3.2.2.
Let $`S`$ be scheme. A *polarization* on a K3 space $`\pi :XS`$ is a global section $`\lambda \mathrm{Pic}_{X/S}(S)`$ such that for every geometric point $`\overline{s}`$ of $`S`$ the section $`\lambda _{\overline{s}}\mathrm{Pic}_{X_{\overline{s}}/\kappa (\overline{s})}(\kappa (\overline{s}))`$ is a polarization of $`X_{\overline{s}}`$. A *quasi-polarization* on $`X/S`$ is a global section $`\lambda \mathrm{Pic}_{X/S}(S)`$ such that for every geometric point $`\overline{s}`$ of $`S`$ the section $`\lambda _{\overline{s}}\mathrm{Pic}_{X_{\overline{s}}/\kappa (\overline{s})}(\kappa (\overline{s}))`$ is a quasi-polarization of $`X_{\overline{s}}`$.
###### Definition 3.2.3.
A polarization (respectively quasi-polarization) $`\lambda `$ on a K3 space $`\pi :XS`$ is called *primitive* if for every geometric point $`\overline{s}`$ of $`S`$ the polarization (respectively the quasi-polarization) $`\lambda _{\overline{s}}\mathrm{Pic}_{X_{\overline{s}}/\kappa (\overline{s})}(\kappa (\overline{s}))`$ is primitive i.e., it is not a positive power of any element in $`\mathrm{Pic}_{X_{\overline{s}}/\kappa (\overline{s})}(\kappa (\overline{s}))`$.
###### Lemma 3.2.4.
Let $`(\pi :XS,\lambda )`$ be a K3 space over $`S`$ with a polarization $`\lambda `$. Then one can find an étale covering $`S^{}S`$ such that $`\pi _S^{}:X_S^{}S^{}`$ is a K3 scheme and $`\lambda _S^{}`$ is the class of a relatively ample line bundle $`_S^{}`$ on $`X_S^{}`$.
###### Proof.
By definition one can find an étale covering $`S_1S`$ such that $`\pi _1:X_{S_1}S_1`$ is a K3 scheme. The pull-back $`\lambda _{S_1}`$ of $`\lambda `$ is a polarization on $`X_{S_1}`$. By Remark 3.1.7 the Picard functor $`\mathrm{Pic}_{X_{S_1}/S_1}`$ can be computed using the étale topology on $`S_1`$. Hence one can find an étale covering $`S^{}S`$ such that $`\lambda _S^{}`$ is equal to the class of a line bundle $`_S^{}`$ on $`X_S^{}`$. By definition $`_S^{}`$ is pointwise ample hence using Lemma 1.2.5 we conclude that it is relatively ample. This finishes the proof. ∎
The self-intersection $`(_{\overline{s}^{}},_{\overline{s}^{}})`$ for a geometric point $`\overline{s}^{}`$ on $`S^{}`$ is constant on every connected component of $`S^{}`$. We say that $`\lambda `$ is a *polarization of degree* $`2d`$ if $`(_{\overline{s}^{}},_{\overline{s}^{}})=2d`$ for every geometric point $`\overline{s}^{}`$ of $`S^{}`$.
### 3.3. Automorphism Groups
Let $`S`$ be a scheme and $`\pi :XS`$ be an algebraic space over $`S`$. Define the automorphism functor in the following way:
$$\mathrm{Aut}_S(X):(\mathrm{Sch}/S)^0\mathrm{Groups}$$
$$\mathrm{Aut}_S(X)(T)=\mathrm{Aut}_T(X_T)$$
for every $`S`$-scheme $`T`$.
###### Theorem 3.3.1.
If $`\pi :XS`$ is a polarized K3 space over $`S`$, then $`\mathrm{Aut}_S(X)`$ is representable by a separated group scheme which is unramified and locally of finite type over $`S`$.
###### Proof.
Let $`S^{}S`$ be an étale cover such that $`\pi ^{}:X^{}=X\times _SS^{}S^{}`$ is a projective K3 scheme over $`S^{}`$. The existence of such an étale covering $`S^{}`$ follows from Lemmas 1.2.5 and 3.2.4. Let $`S^{\prime \prime }`$ be the product $`S^{}\times _SS^{}`$. Denote by $`\pi _i`$ the projection morphisms $`\pi _i:X^{}\times _XX^{}X^{}XS`$ for $`i=1,2`$. By definition $`X^{}\times _XX^{}`$ is representable by a quasi-compact subscheme of $`X^{}\times _SX^{}`$.
Using Proposition 1.4 in \[Knu71, Ch. II\] we can see that we have an exact sequence of groups
(4)
It follows from \[Gro62, Exp. 221, §4.c\] that the functors $`\mathrm{Aut}_S^{}(X^{})`$ and $`\mathrm{Aut}_{S^{\prime \prime }}(X^{}\times _XX^{})`$ are representable by group schemes locally of finite type over $`S`$. For simplicity we denote them by $`𝒴`$ and $`𝒲`$ respectively. Then from the exact sequence (4) we see that $`\mathrm{Aut}_S(X)`$ is representable by the fiber product
where $`\mathrm{\Delta }:𝒲𝒲\times _S𝒲`$ is the diagonal morphism.
The fact that the $`\mathrm{Aut}_S(X)`$ is separated follows directly from the valuative criterion for separatedness.
To check that $`\mathrm{Aut}_S(X)`$ is unramified we may take $`S`$ to be the spectrum of an algebraically closed field $`k`$. A point in $`\mathrm{Aut}_k(X)(k[ϵ]/(ϵ^2))`$, which under the natural homomorphism maps to the the identity in $`\mathrm{Aut}_k(X)(k)`$, may be identified with a vector field on $`X`$. By Proposition 2.2.1 (1) a K3 surface has no non-trivial vector fields hence we conclude that $`\mathrm{Aut}_k(X)`$ is reduced. ∎
###### Remark 3.3.2.
The proof of the theorem shows that $`\mathrm{Aut}_S(X)`$ is $`0`$-dimensional over $`S`$. Its fibers are constant group schemes.
Let $`\pi :XS`$ be a K3 space and let $`\lambda `$ be a polarization of $`X`$. Define the subfunctor $`\mathrm{Aut}_S(X,\lambda )`$ of $`\mathrm{Aut}_S(X)`$ in the following way
$$\mathrm{Aut}_S(X,\lambda ):(\mathrm{Sch}/S)^0\mathrm{Groups}$$
$$\mathrm{Aut}_S(X,\lambda )(T)=\{\alpha \mathrm{Aut}_S(X)(T)|\alpha ^{}\lambda =\lambda \mathrm{Pic}_{X/S}(T)\}$$
for every $`S`$-scheme $`T`$.
###### Proposition 3.3.3.
The functor $`\mathrm{Aut}_S(X,\lambda )`$ is a closed subfunctor of $`\mathrm{Aut}_S(X)`$. It is represented by a separated group scheme which is unramified and of finite type over $`S`$. Its relative dimension over $`S`$ is zero.
###### Proof.
The functor $`\mathrm{Aut}_S(X,\lambda )`$ is a closed subfunctor of $`\mathrm{Aut}_S(X)`$. It is representable by the subgroup scheme of $`G=\mathrm{Aut}_S(X)`$ (locally of finite type over $`S`$) given by the following (Cartesian) diagram:
Here we have that $`\lambda :S\mathrm{Pic}_{X/S}`$ is the section given by $`\lambda `$ and $`\psi `$ is the composition $`\sigma (\mathrm{id},\lambda )`$ where
$$\sigma :G\times \mathrm{Pic}_{X/S}\mathrm{Pic}_{X/S}$$
is the action of $`G`$ on $`\mathrm{Pic}_{X/S}`$.
Just as in the proof of the preceding theorem we may take $`S`$ to be the spectrum of an algebraically closed field $`k`$ in order to check that $`\mathrm{Aut}_S(X,\lambda )`$ is unramified. If $`\alpha \mathrm{Aut}_k(X,\lambda )(k[ϵ]/ϵ^2)`$ which is the identity in $`\mathrm{Aut}_k(X,\lambda )(k)`$, then by Theorem 3.3.1 above we see that $`\alpha `$ is the identity element of the group $`\mathrm{Aut}_k(X)(k[ϵ]/ϵ^2)`$. Since by definition we have an inclusion
$$\mathrm{Aut}_k(X,\lambda )(k[ϵ]/ϵ^2)\mathrm{Aut}_k(X)(k[ϵ]/ϵ^2)$$
we conclude that $`\mathrm{Aut}_S(X,\lambda )`$ is unramified over $`S`$.
Let $`\overline{s}:\mathrm{Spec}(\mathrm{\Omega })S`$ be a geometric point. Then by \[Mat58\] (see also Corollary 2 in \[MM64\]) the set $`\mathrm{Aut}_S(X,\lambda )(\mathrm{\Omega })`$ is finite. Hence $`\mathrm{Aut}_S(X,\lambda )`$ is of finite type over $`S`$. ∎
Note that in general, for a K3 surface $`X`$ over a field $`k`$, the group $`\mathrm{Aut}_k(X)(k)`$ might be infinite.
###### Example 3.3.4.
For any complex K3 surface $`X`$ with $`\mathrm{rk}_{}\mathrm{Pic}(X)=20`$ one has that $`\mathrm{Aut}_{}(X)()`$ is infinite. For a proof see \[SI77, §5, Thm. 5\].
There are also examples of K3 surfaces $`X`$ having a finite group of automorphisms. An example of a complex K3 surface with $`\mathrm{rk}_{}\mathrm{Pic}(X)=18`$ and finite automorphism group is given in the remark on page 132 in \[SI77\].
### 3.4. Automorphisms of Finite Order
In this section $`k`$ will be an algebraically closed field. If it is a field of characteristic $`p`$, then we will denote by $`W`$ the ring of Witt vectors with coefficients in $`k`$ and $`K`$ will be the field of fractions of $`W`$.
Let $`X`$ be a K3 surface over $`k`$. If $`k=`$, then it is a well-known theorem that $`\mathrm{Aut}_{}(X)()`$ acts faithfully on $`H_B^2(X,)`$. Here we prove a similar result for the automorphisms of finite order of $`X`$ acting trivially on $`H_{\mathrm{et}}^2(X,_l)`$ where $`l`$ is a prime number different from $`\mathrm{char}(k)`$. The only restriction we impose is that $`\mathrm{char}(k)2`$. Later on in Section 5.1 we will introduce level structures on K3 surfaces and we will use this result to show that the corresponding moduli stacks are algebraic spaces.
###### Lemma 3.4.1.
Let $`X`$ be a K3 surface over $`k`$ and assume that char$`(k)=0`$. Then $`\mathrm{Aut}_k(X)(k)`$ acts faithfully on $`H_{\mathrm{et}}^2(X,_l)`$ for every prime $`l`$.
###### Proof.
Without loss of generality we may assume that the field $`k`$ can be embedded into $``$. Fix an embedding $`\sigma :k`$. By the comparison theorem between Betti and étale cohomology we have an isomorphism $`H_{\mathrm{et}}^2(X,_l)H_B^2(X_\sigma ,)_{}_l`$. Let $`\alpha \mathrm{Aut}_k(X)(k)`$ be an automorphism acting trivially on $`H_{\mathrm{et}}^2(X,_l)`$. Then $`\alpha _{}`$ acts trivially on $`H_B^2(X_\sigma ,)_l`$. Since $`H_B^2(X_\sigma ,)`$ is a free $``$-module we conclude from \[LP81, Prop. 7.5\] that $`\alpha =\mathrm{id}_X`$. ∎
###### Proposition 3.4.2.
Let $`(X,\lambda )`$ be a polarized K3 surface over $`k`$ and assume that $`\mathrm{char}(k)=p`$ is different from $`2`$. Then the finite group $`\mathrm{Aut}_k(X,\lambda )(k)`$ acts faithfully on $`H_{\mathrm{et}}^2(X,_l)`$ for any $`lp`$.
###### Remark 3.4.3.
This result can be viewed as an analogue of Theorem 3 in \[Mum74, Ch. IV\] for (polarized) K3 surfaces.
We will reduce the proof of Proposition 3.4.2 to the preceding lemma. To do so we will use crystalline cohomology and compare the action of an element in $`\mathrm{Aut}_k(X,\lambda )(k)`$ on $`H_{\mathrm{et}}^2(X,_l)`$ and $`H_{\mathrm{cris}}^2(X/K)`$.
Let $`X`$ be a K3 surface over a field $`k`$. We denote by $`H^n(X)`$ and $`H^n(X\times X)`$ either $`H_{\mathrm{et}}^n(X,_l)`$ and $`H_{\mathrm{et}}^n(X\times X,_l)`$ for any $`l`$ prime to char$`(k)`$ or $`H_{\mathrm{cris}}^n(X/K)`$ and $`H_{\mathrm{cris}}^n(X\times X/K)`$. Note that we will be working with classes of certain algebraic cycles on $`X`$ and $`X\times X`$ so we should consider some Tate twists of these cohomology groups. But since $`k`$ is algebraically closed and the Galois action does not play any role in our consideration (we shall only consider some characteristic polynomials of automorphisms of $`X`$) we will omit these twists.
For an isomorphism $`\alpha :XY`$ we will denote by $`\alpha _l^{}`$ and $`\alpha _{\mathrm{cris}}^{}`$ the isomorphisms induced on $`H_{\mathrm{et}}^2(X,_l)`$ and $`H_{\mathrm{cris}}^2(X/K)`$ respectively.
###### Lemma 3.4.4.
The Künneth components of the class $`cl(u)H^4(X\times X)`$ of any algebraic cycle on $`X\times X`$ are algebraic.
###### Proof.
We have that $`H_{\mathrm{et}}^1(X,_l)=H_{\mathrm{et}}^3(X,_l)=0`$ and $`H_{\mathrm{cris}}^1(X/W)=H_{\mathrm{cris}}^3(X/W)=0`$. Then the Künneth isomorphism reads
$$H^4(X\times X)=\left(H^4(X)H^0(X)\right)\left(H^2(X)H^2(X)\right)\left(H^0(X)H^4(X)\right).$$
Using this decomposition we write
$$cl(u)=u_0u_2u_4.$$
Every element of the one dimensional spaces $`H^4(X)H^0(X)`$ and $`H^0(X)H^4(X)`$ is algebraic. These are rational multiple of the classes of $`\{pt\}\times X`$ and $`X\times \{pt\}`$. Hence $`u_0`$ and $`u_4`$ are algebraic. It follows that $`u_2`$ is expressed as a linear combination of algebraic classes, hence it is algebraic. ∎
In particular, if $`\mathrm{\Delta }=\delta (X)X\times X`$ is the diagonal, then its Künneth components $`cl(\mathrm{\Delta })=\pi _0\pi _2\pi _4H^4(X\times X)`$ are algebraic. Denote by $`,`$ the intersection pairing on $`\mathrm{CH}^2(X\times X)_{}`$.
###### Corollary 3.4.5.
Let $`u\mathrm{CH}^2(X\times X)_{}`$ be a rational cycle and let $`cl(u)H^4(X\times X)`$ be its algebraic class. Then its characteristic polynomial $`det\left(1tcl(u)|H^2(X)\right)`$ has rational coefficients which are independent of $`l`$ and $`p`$ (i.e., of $`H_{\mathrm{et}}^2(X,_l)`$ and $`H_{\mathrm{cris}}^2(X/K)`$). The coefficient in front of $`t^i`$ is given by
$$s_i=u_i,\pi _2$$
for $`i=1,\mathrm{},22`$.
###### Proof.
The proof follows from the preceding lemma and by Theorem 3.1 in \[Tat95\]. ∎
###### Theorem 3.4.6 (Ogus).
If $`p>2`$ then the natural morphism of groups
$$\mathrm{Aut}_k(X)(k)\mathrm{Aut}\left(H_{\mathrm{cris}}^2(X/W)\right)$$
is injective.
###### Proof.
This is a result of A. Ogus and can be found in his paper on Supersingular K3 crystals \[Ogu79, §2, Cor. 2.5\]. ∎
Proof of Proposition 3.4.2. Take an element $`\alpha \mathrm{Aut}_k(X,\lambda )(k)`$. According to Proposition 3.3.3 it has finite order. Denote by $`u=\mathrm{\Gamma }_\sigma X\times X`$ the graph of $`\alpha `$. Then the automorphism of $`H_{\mathrm{et}}^2(X,_l)`$ induced by $`cl(u)H_{\mathrm{et}}^4(X\times X,_l)`$ is the one induced by $`\alpha `$. By assumption it is the identity hence its characteristic polynomial is $`(t1)^{22}`$. By Corollary 3.4.5 it is exactly the characteristic polynomial of the automorphism $`\alpha _{\mathrm{cris}}^{}`$ of $`H_{\mathrm{cris}}^2(X/K)`$ induced by $`\alpha `$. Since $`\alpha `$ is an automorphism of finite order the induced map $`\alpha _{\mathrm{cirs}}^{}`$ on the crystalline cohomology is semi-simple ($`K`$ has characteristic zero). Hence $`\alpha _{\mathrm{cris}}^{}`$ acts trivially on $`H_{\mathrm{cris}}^2(X/K)`$ and by Theorem 3.4.6 it is the identity automorphism as $`H_{\mathrm{cris}}^2(X/W)`$ is torsion free. ∎
###### Remark 3.4.7.
Note that the only property of $`\alpha `$ which we used in the proof of Proposition 3.4.2 is that it has finite order. This is really essential as in general the characteristic polynomial of $`\alpha _l^{}`$ will not give enough information to conclude that the action of $`\alpha _{\mathrm{cris}}`$ on $`H_{\mathrm{cris}}^2(X/W)`$ is trivial. The proof given above shows actually that any automorphism of finite order $`\alpha `$ of $`X`$ acting trivially on $`H_{\mathrm{et}}^2(X,_l)`$ for some $`lp`$ is the identity automorphism $`\mathrm{id}_X`$.
## 4. The Moduli Stack of Polarized K3 Surfaces
We are ready to define moduli functors of (primitively) polarized K3 surfaces over $`\mathrm{Spec}()`$. We will follow the line of thoughts in \[DM68\] in order to prove that these functors define Deligne-Mumford stacks. Shortly, this can be given in three steps.
1. Describe the deformations of primitively polarized K3 surfaces.
2. Construct a Hilbert scheme parameterizing K3 surfaces embedded in $`^N`$ for some appropriate $`N`$.
3. Construct a “Hilbert morphism” $`\pi _{\mathrm{Hib}}`$ from the Hilbert scheme to the moduli stack which is surjective and smooth. Use this morphism to conclude that the moduli stack is a Deligne-Mumford stack.
These steps are spelled out in detail in Sections 4.1-4.3.
### 4.1. Deformations of K3 Surfaces
Let $`k`$ be an algebraically closed field. Denote by $`W`$ the ring of Witt vectors $`W(k)`$ in case $`\mathrm{char}(k)=p>0`$ and $`W=k`$ otherwise. Let $`\underset{¯}{A}`$ be the category of local artinian $`W`$-algebras $`(A,𝔪_A)`$ together with an isomorphism $`A/𝔪_Ak`$ compatible with the isomorphism $`W/pWk`$.
Let $`X_0`$ be a K3 surface over $`k`$. Consider the covariant functor
$$\mathrm{Def}_{\mathrm{Sch}}(X_0):\underset{¯}{A}\mathrm{Sets}$$
given by
$$\begin{array}{cc}\hfill \mathrm{Def}_{\mathrm{Sch}}(X_0)(A)=\{\mathrm{isom}.\mathrm{classes}\mathrm{of}\mathrm{pairs}(X,\varphi _0)|& \mathrm{where}X\mathrm{Spec}(A)\hfill \\ & \mathrm{is}\mathrm{a}\mathrm{K3}\mathrm{scheme}\mathrm{and}\varphi _0\mathrm{is}\hfill \\ & \mathrm{an}\mathrm{isom}.\varphi _0:X_AkX_0\}.\hfill \end{array}$$
###### Proposition 4.1.1.
The functor $`\mathrm{Def}_{\mathrm{Sch}}(X_0)`$ is pro-representable by a formal scheme $`S`$ over $`\mathrm{Spf}(W)`$ which is formally smooth of relative dimension $`20`$ i.e., it is (non-canonically) isomorphic to $`\mathrm{Spf}(W[[t_1,\mathrm{},t_{20}]])`$.
###### Proof.
This is Corollary 1.2 in \[Del81b\] in case $`\mathrm{char}(k)=p>0`$ and \[LP81, Cor. 5.7\] in case $`\mathrm{char}(k)=0`$. ∎
Let $`_0`$ be a line bundle on $`X_0`$. For moduli problems one should study the deformations of the pair $`(X_0,_0)`$. Define
$$\mathrm{Def}_{\mathrm{Sch}}(X_0,_0):\underset{¯}{A}\mathrm{Sets}$$
to be the functor sending an object $`A`$ of $`\underset{¯}{A}`$ to the isomorphism classes of triples $`(X,,\varphi _0)`$ of flat deformations $`X`$ of $`X_0`$ over $`A`$, an invertible sheaf $``$ on $`X`$ and an isomorphism $`\varphi _0:(X,)_Ak(X_0,_0)`$. We have a morphism
(5)
$$\mathrm{Def}_{\mathrm{Sch}}(X_0,_0)\mathrm{Def}_{\mathrm{Sch}}(X_0).$$
###### Theorem 4.1.2.
If the line bundle $`_0`$ is non-trivial, then the functor $`\mathrm{Def}_{\mathrm{Sch}}(X_0,_0)`$ is pro-representable by a formally flat scheme of relative dimension 19 over $`W`$ and the morphism (5) is a closed immersion, defined by a single equation.
###### Proof.
See \[Del81b, Prop. 1.5 and Thm. 1.6\]. ∎
Deligne proves that if $`_0`$ is an ample line bundle over $`X_0`$ then one can find a discrete valuation ring $`R`$ which is a finite $`W`$ module and a lift $`(X\mathrm{Spec}(R),)`$ of $`(X_0,_0)`$ over $`R`$. In general one needs ramified extensions of $`W`$ in order to find a lift of $`(X_0,_0)`$. The next lemma shows that one can find a lift over $`W`$ if the self-intersection of $`_0`$ is prime to the characteristic of $`k`$. More precisely one has:
###### Lemma 4.1.3.
Let $`_0`$ be an ample line bundle over $`X_0`$. If the polarization degree $`(_0,_0)_{X_0}=2d`$ is prime to the characteristic of $`k`$, then $`\mathrm{Def}_{\mathrm{Sch}}(X_0,_0)`$ is formally smooth.
###### Proof.
According to \[Ogu79, §2, Prop. 2.2 \] (see also Lemma 2.2.6 in \[Del81a\]) it is enough to see that $`c_1(_0)F^2H_{DR}^2(X_0/k)`$. Since we have that $`(c_1(_0),c_1(_0))=2d0`$ in $`k`$ it follows that $`c_1(_0)F^2H_{DR}^2(X_0/k)`$. For the proof in the case $`k`$ has characteristic zero we refer to \[PSS72, §2, Thm. 1\]. ∎
### 4.2. The Hilbert Scheme
Recall that if $`X`$ is a K3 surface over a field $`k`$ with an ample line bundle $``$, then the Hilbert polynomial of $``$ is $`h_{}(x)=dx^2+2`$, where $`(,)=2d`$.
We fix two natural numbers $`n`$ and $`d`$ assuming that $`n3`$. Let $`P_{d,n}(x)`$ be the polynomial $`n^2dx^2+2`$ and let $`N=P_{d,n}(1)1`$. Denote by $`\mathrm{𝐇𝐢𝐥𝐛}_N^{P_{d,n}}`$ the *Hilbert scheme* over $``$ representing the subvarieties of $`^N`$ with Hilbert polynomial $`P_{d,n}(x)`$. Let
$$\pi :𝒵\mathrm{𝐇𝐢𝐥𝐛}_N^{P_{d,n}}$$
be the universal family over the Hilbert scheme. For any morphism of schemes $`f:S\mathrm{𝐇𝐢𝐥𝐛}_N^{P_{d,n}}`$ we consider the following (Cartesian) diagram:
(6)
###### Proposition 4.2.1.
There is a unique subscheme $`H_{d,n}`$ of $`\mathrm{𝐇𝐢𝐥𝐛}_N^{P_{d,n}}`$ with the property:
A morphism of schemes $`f:S\mathrm{𝐇𝐢𝐥𝐛}_N^{P_{d,n}}`$ factors through $`H_{d,n}`$ if and only if the following conditions are satisfied.
1. The pull-back $`𝒳`$ of the universal family over $`\mathrm{𝐇𝐢𝐥𝐛}_N^{P_{d,n}}`$ is a K3 scheme over $`S`$ (see Diagram (6) above),
2. the line bundle $`f_{}^{}{}_{}{}^{}𝒪_^N(1)`$ is isomorphic to $`^n\pi _{}^{}{}_{}{}^{}`$ for some ample line bundle $``$ on $`𝒳`$ and some line bundle $``$ on $`S`$,
3. for every geometric point $`\overline{s}:\mathrm{Spec}(\mathrm{\Omega })S`$ the natural homomorphism
$$H^0(^N,𝒪_^N(1))\mathrm{\Omega }H^0(𝒳_{\overline{s}},_{\overline{s}}^n)$$
is an isomorphism.
There exists an open subscheme $`H_{d,n}^{pr}`$ of $`H_{d,n}`$ such that: A morphism of schemes $`f:S\mathrm{𝐇𝐢𝐥𝐛}_N^{P_{d,n}}`$ factors through $`H_{d,n}^{pr}`$ if and only if conditions (i), (ii) and (iii) are satisfied and in addition for every geometric point $`\overline{s}`$ of $`S`$ the line bundle $`_{\overline{s}}`$ from (ii) is primitive.
###### Proof.
The proof of the proposition is standard and can be found in the case of curves in Mumford’s book \[Mum65, Ch. 5, §2, Prop. 5.1\]. We shall sketch only the additional arguments needed in our situation.
There is a maximal open subscheme $`U_1`$ of $`\mathrm{𝐇𝐢𝐥𝐛}_N^{P_{d,n}}`$ such that every fiber of the pull-back $`𝒳_1`$ of the universal family $`𝒵`$ over $`U_1`$ is a non-singular variety. Let $`U_2`$ be the open subscheme of $`U_1`$ consisting of the points $`s`$ for which $`H^1(𝒳_{1,s},𝒪_{𝒳_{1,s}})=0`$ (see \[Har77, Ch. III, §12, Thm. 12.8\]). Denote by $`𝒳_2`$ the pull-back of the universal family over $`U_2`$.
Let $`\mathrm{Pic}_{𝒳_2/U_2}`$ be the relative Picard scheme of $`𝒳_2`$ over $`U_2`$. The two line bundles $`\mathrm{\Omega }_{𝒳_2/U_2}^2`$ and $`𝒪_{𝒳_2}`$ define two morphisms: $`\omega ,\lambda :U_2\mathrm{Pic}_{𝒳_2/U_2}`$. Define $`U_2`$ to be the fiber product:
where $`\mathrm{\Delta }`$ is the diagonal morphism. Since $`\mathrm{Pic}_{𝒳_2/U_2}`$ is separated $`U_3`$ is a closed subscheme of $`U_2`$. The pull-back $`𝒳_3U_3`$ of the universal family over $`\mathrm{𝐇𝐢𝐥𝐛}_N^{P_{d,n}}`$ is a K3 scheme.
Let $`[n]:\mathrm{Pic}_{𝒳_3/U_3}\mathrm{Pic}_{𝒳_3/U_3}`$ be the multiplication by-$`n`$-morphism. The pull-back of $`𝒪_^N(1)`$ over $`U_3`$ defines a morphism $`\lambda :U_3\mathrm{Pic}_{𝒳_3/U_3}`$. Define $`U_4`$ to be the fiber product
By Lemma 3.1.6 the morphism $`[n]`$ is a closed immersion hence $`U_4`$ is a closed subscheme of $`U_3`$. Clearly, $`U_4`$ is the subscheme of $`\mathrm{𝐇𝐢𝐥𝐛}_N^{P_{d,n}}`$ for which properties (i) and (ii) hold. One takes $`H_{d,n}`$ to be the (closed) subscheme of $`U_4`$ obtained as in the end of the proof of Proposition 5.1 in \[Mum65, Ch. 5, §2\] (where instead of $`\mathrm{\Omega }_{\mathrm{\Gamma }/U_2}^1`$ one works with the pull-back $`^{}`$ of the bundle $`𝒪_^N(1)`$). It satisfies all conditions of the proposition.
To show the existence of $`H_{n,d}^{pr}`$ one has to take the open subscheme $`U_4^0`$ of $`U_4`$ above corresponding to the points in $`U_4`$ over which the class of the pull-back of $`𝒪_^N(1)`$ in $`\mathrm{Pic}_{𝒳_4/U_4}`$ is only divisible by $`n`$. The existence of such a subscheme can be seen, in a way similar to the proof of Lemma 3.1.6 (ii), using the fact that the homomorphisms $`[n]:\mathrm{Pic}_{𝒳_4/U_4}\mathrm{Pic}_{𝒳_4/U_4}`$ are closed immersions. ∎
We will use the schemes $`H_{d,n}`$ and $`H_{d,n}^{pr}`$ to construct moduli stacks of polarized K3 surfaces over $``$.
### 4.3. The Moduli Stack
One way to construct the coarse moduli space of complex K3 surfaces with a primitive polarization of degree $`2d`$ is to use period maps. This approach is taken up in \[BBD85, Exposé XIII, §3\]. Here we will use rather different techniques to deal with this problem in positive and more generally in mixed characteristic.
###### Definition 4.3.1.
Let $`d`$ be a natural number. Consider the category $`_{2d}`$ defined in the following way:
1. The objects of $`_{2d}`$ are pairs $`(\pi :XS,\lambda )`$ consisting of a K3 space $`\pi :XS`$ with a primitive polarization $`\lambda `$ of degree $`2d`$ over $`S\mathrm{Sch}`$.
2. For two objects $`𝒳_1=(\pi _1:X_1S_1,\lambda _1)`$ and $`𝒳_2=(\pi _2:X_2S_2,\lambda _2)`$ we define the morphisms to be
$$\begin{array}{cc}\hfill \mathrm{Hom}(𝒳_1,𝒳_2)=\{\mathrm{pairs}(f_S,f)|& f_S:S_1S_2\mathrm{is}\mathrm{a}\mathrm{morph}.\mathrm{of}\hfill \\ & \mathrm{schemes}\mathrm{and}f:X_1X_2\times _{S_2,f_S}S_1\hfill \\ & \mathrm{is}\mathrm{an}\mathrm{isom}.\mathrm{over}S_1\mathrm{with}f^{}\lambda _2=\lambda _1\}.\hfill \end{array}$$
The functor $`p_{_{2d}}:_{2d}\mathrm{Sch}`$ sending a pair $`(\pi :XS,\lambda )`$ to $`S`$ makes $`_{2d}`$ into a category over $`\mathrm{Sch}`$. We will denote by $`_{2d,S}`$ the full subcategory of $`_{2d}`$ consisting of the objects over $`S`$.
###### Definition 4.3.2.
For a natural number $`d`$ we define the category $`_{2d}`$ of K3 spaces with a polarization of degree $`2d`$ in the same way as in Definition 4.3.1 but taking as objects pairs of polarized K3 spaces $`(\pi :XS,\lambda )`$ over a scheme $`S`$.
We have that $`_{2d}`$ is a full subcategory of $`_{2d}`$. Those two categories are the same if and only if $`d`$ is square-free.
###### Theorem 4.3.3.
The categories $`_{2d}`$ and $`_{2d}`$ are separated Deligne-Mumford stacks of finite type over $``$. The inclusion $`_{2d}_{2d}`$ is an open immersion.
###### Definition 4.3.4.
We will call $`_{2d}`$ *the moduli stack of primitively polarized K3 surfaces of degree* $`2d`$ and $`_{2d}`$ *the moduli stack of polarized K3 surfaces of degree* $`2d`$.
###### Remark 4.3.5.
Let us explain first why we want to consider moduli of *primitively* polarized K3 surfaces. For various reasons we will have to work with algebraic spaces rather than with algebraic stacks. Just like in the case of abelian varieties one can introduce level structures on K3 surfaces and hope that the corresponding moduli problems are representable by algebraic spaces. We will define level structures on a polarized K3 surface $`(X,\lambda )`$ using its primitive cohomology groups $`P_{\mathrm{et}}^2(X_{\overline{k}},_l(1))`$ for certain primes $`l`$ (see Section 5.1). To be able to do that we will need that $`P_{\mathrm{et}}^2(X_{\overline{k}},_l(1))`$ belongs to a single isometry class of quadratic lattices, which is the case, if $`\lambda `$ is primitive.
We will prove the theorem in a sequence of steps.
###### Lemma 4.3.6.
The categories $`_{2d}`$ and $`_{2d}`$ are groupoids.
###### Proof.
We have to check two axioms. See for instance \[LMB00, Ch. 2, Def. 2.1\] or p. 96 of \[DM68\]. One sees immediately that the usual notions of pull-backs satisfy these two axioms. ∎
###### Lemma 4.3.7.
The groupoids $`_{2d}`$ and $`_{2d}`$ are stacks for the étale topology.
###### Proof.
The proofs for $`_{2d}`$ and $`_{2d}`$ are exactly the same so we will prove the lemma for $`_{2d}`$. We have to check two properties. Namely, first we will show that for any scheme $`S\mathrm{Sch}`$ and any two objects $`𝒳`$ and $`𝒴`$ over $`S`$ the functor
$$\mathrm{Isom}_S(𝒳,𝒴):(\mathrm{Sch}/S)\mathrm{Sets}$$
defined by
$$(\pi :S^{}S)\mathrm{Hom}(\pi ^{}𝒳,\pi ^{}𝒴)$$
is a sheaf for the étale topology on $`S`$. Then we prove that descent data are effective (cf. \[LMB00, Ch. 2, Def. 3.1\] or Definition 4.1 in \[DM68\]).
The functor $`\mathrm{Isom}_S(𝒳,𝒴)`$ is an étale sheaf: Take two objects $`𝒳=(XS,\lambda _X)`$ and $`𝒴=(YS,\lambda _Y)`$ over $`S`$. Let $`S^{}`$ be an $`S`$-scheme.
Let $`\{S_i^{}\}_{iI}`$ be an étale covering of $`S^{}`$ and $`f_j\mathrm{Isom}_S(𝒳,𝒴)(S^{})`$ for $`j=1,2`$ are two elements such that $`f_1|_{S_i^{}}=f_2|_{S_i^{}}`$. Then clearly $`f_1=f_2`$ as isomorphisms of the pair $`(𝒳_S^{},𝒴_S^{})`$.
Let $`\{S_i^{}\}_{iI}`$ be an étale covering of $`S^{}`$. Suppose given elements $`f_i\mathrm{Isom}_S(𝒳,𝒴)(S_i^{})`$ such that $`f_i|_{S_{ij}^{}}=f_j|_{S_{ij}^{}}`$ where $`S_{ij}^{}=S_i^{}\times _S^{}S_j^{}`$. We have to show that those come from a global “isomorphism”. Note that without loss of generality we may assume that $`X_iS_i^{}`$ are K3 schemes. Combining \[Knu71, Ch. II, Prop. 1.4\] and effectiveness of descent for morphisms of schemes (see \[BLR90, Ch. 6, §1, Thm. 6(a)\]) we conclude that $`f^{}`$ descends to a morphism $`f:X_S^{}Y_S^{}`$ such that $`f_{S_i^{}}=f_i`$. Since $`\mathrm{Pic}_{X/S}`$ and $`\mathrm{Pic}_{Y/S}`$ are algebraic spaces (in particular sheaves for the étale topology on $`S`$) and $`f^{}\lambda _{Y_S^{}}|_{S_i^{}}=\lambda _{X_S^{}}|_{S_i^{}}`$ we see that $`f^{}\lambda _{Y_S^{}}=\lambda _{X_S^{}}`$. Hence we have that $`f\mathrm{Isom}_S(𝒳,𝒴)(S^{})`$ and $`f|_{S_i^{}}=f_i`$. This shows that $`\mathrm{Isom}_S(𝒳,𝒴)`$ is an étale sheaf.
Effectiveness of descent: Suppose given an étale cover $`S^{}`$ of $`S`$ and an object $`𝒳^{}=(\pi ^{}:X^{}S^{},\lambda ^{})`$ with descent datum over $`S`$. Without loss of generality we may assume that the algebraic space $`X^{}`$ is actually a scheme (by refining the étale covering $`S^{}`$ if needed). We have to show that $`(\pi ^{}:X^{}S^{},\lambda ^{})`$ descends to a polarized K3 space $`(\pi :XS,\lambda )`$ over $`S`$.
Denote by $`S^{\prime \prime }`$ the product $`S^{}\times _SS^{}`$ and let $`pr_i`$ for $`i=1,2`$ be the two projection maps. The descent datum on $`X^{}S^{}`$ over $`S`$ identifies the two schemes $`pr_1^{}X^{}`$ and $`pr_2^{}X^{}`$. Denote this scheme by $`R`$. Then we have two étale morphisms
which make $`RX^{}\times _SX^{}`$ into an étale equivalence relation. Following the constructions of \[Knu71, Ch. I, §5, 5.4\] we obtain an algebraic space $`X`$ over $`S`$ such that $`X\times _SS^{}`$ is isomorphic to $`X^{}`$. Hence $`\pi :XS`$ is a K3 space.
Since $`\mathrm{Pic}_{X/S}`$ is an étale sheaf the local section $`\lambda ^{}`$ over $`S^{}`$ together with descent datum over $`S`$ give rise to a global section $`\lambda \mathrm{Pic}_{X/S}(S)`$ such that $`\lambda _S^{}=\lambda ^{}`$. Clearly, $`\lambda `$ is a polarization of $`XS`$. ∎
Next we deal with the representability of the isomorphism functors of polarized K3 surfaces. For two algebraic spaces $`X`$ and $`Y`$ over a base scheme $`S`$ define the contravariant isomorphism functor
$$\mathrm{Isom}_S(X,Y):(\mathrm{Sch}/S)\mathrm{Sets}$$
by
$$\mathrm{Isom}_S(X,Y)(T)=\{f:X_TY_T|f\mathrm{is}\mathrm{an}\mathrm{isomorph}.\mathrm{of}\mathrm{alg}.\mathrm{spaces}\mathrm{over}T\}$$
for any $`S`$-scheme $`T`$.
###### Lemma 4.3.8.
For any $`S\mathrm{Sch}`$ and two objects $`𝒳`$ and $`𝒴`$ of $`_{2d}`$ (respectively $`_{2d}`$) over $`S`$, the functor $`\mathrm{Isom}_S(𝒳,𝒴)`$ is representable by a separated scheme which is unramified and of finite type over $`S`$.
###### Proof.
Let $`𝒳`$ and $`𝒴`$ be the objects $`(XS,\lambda _X)`$ and $`(YS,\lambda _Y)`$ respectively.
Step 1: We can find an étale cover $`S^{}`$ of $`S`$ such that $`X^{}=X\times _SS^{}`$ and $`Y^{}=Y\times _SS^{}`$ are projective K3 schemes over $`S^{}`$. Denote by $`S^{\prime \prime }`$ the product $`S^{}\times _SS^{}`$. By \[Gro62, Exp. 221, §4.c\]) the functors $`\mathrm{Isom}_S^{}(X^{},Y^{})`$ and $`\mathrm{Isom}_{S^{\prime \prime }}(X^{}\times _XX^{},Y^{}\times _YY^{})`$ are representable by schemes $`𝒰`$ and $`𝒱`$, locally of finite type over $`S`$. By Proposition 1.4 in \[Knu71, Ch. II\] one has an exact sequence of sets
Then we see that $`\mathrm{Isom}_S(X,Y)`$ is representable by the scheme defined by the following Cartesian diagram
where $`\mathrm{\Delta }:𝒱𝒱\times _S𝒱`$ is the diagonal morphism.
Step 2: By Step 1 the functor $`\mathrm{Isom}_S(X,Y)`$ is represented by a scheme locally of finite type over $`S`$. Then the functor $`\mathrm{Isom}_S(𝒳,𝒴)`$ is represented by the scheme defined by the following Cartesian diagram:
where the bottom-right arrow is just the pull back morphism.
Step 3: We are left to show that $`\mathrm{Isom}_S(𝒳,𝒴)`$ is unramified over $`S`$. As in the proof of Theorem 3.3.1 it is enough to check the properties of $`\mathrm{Isom}_S(𝒳,𝒴)`$ when $`S`$ is a spectrum of an algebraically closed field. In this case $`\mathrm{Isom}_S(𝒳,𝒴)`$ is either empty or it is isomorphic to $`\mathrm{Aut}_k(X,\lambda )`$. As the latter is separated, reduced and of finite type over $`k`$ we conclude that the same holds for $`\mathrm{Isom}_S(𝒳,𝒴)`$. ∎
Proof of Theorem 4.3.3: We will give the proof for $`_{2d}`$ in several steps. For the proof that $`_{2d}`$ is a Deligne-Mumford stack one should only replace $`H_{d,3}^{pr}`$ by $`H_{d,3}`$ below.
Step 1: We saw in Proposition 4.2.1 that there exists a Hilbert scheme $`H_{3,d}^{pr}`$, of finite type over $``$, classifying K3 surfaces with a polarization of degree $`2d`$ which are embedded in a projective space via the third power of the polarization. One has then the universal family $`f:𝒳H_{3,d}^{pr}`$ and we know that $`𝒪_𝒳(1)^3f^{}`$ for some ample line bundle $``$ on $`𝒳`$ of degree $`2d`$ and an invertible sheaf $``$ on $`H_{3,d}^{pr}`$. Although the line bundle $``$ with this property is not unique, its class $`\lambda _𝒳=cl()\mathrm{Pic}_{𝒳/H_{3,d}^{pr}}`$ is uniquely determined as $`\lambda _𝒳^3=cl(𝒪_𝒳(1))`$. Define the morphism of stacks
$$\pi _{\mathrm{Hilb}}:H_{3,d}^{pr}_{2d}.$$
sending $`H_{3,d}^{pr}`$ to the pair $`(f:𝒳H_{3,d}^{pr},\lambda _𝒳)`$. By construction the self-intersection index $`(\lambda _{𝒳,h},\lambda _{𝒳,h})`$ is $`2d`$ for any $`hH_{3,d}^{pr}`$ and $`\lambda _𝒳`$ is primitive so this morphism is correctly defined.
Step 2: The morphism $`\pi _{\mathrm{Hilb}}`$ is surjective. This follows form the definition (cf. \[LMB00, Def. 3.6\]) and Lemma 3.2.4. Indeed, for any $`(\pi :XS,\lambda )_{2d}(S)`$ one can find an étale cover $`S^{}S`$ such that $`\pi _S^{}:X_S^{}S^{}`$ is a K3 scheme and $`\lambda _S^{}`$ is equal to the class of a relatively ample line bundle $`^{}`$ on $`X_S^{}`$. By Lemma 1.2.5 the line bundle $`_{}^{}{}_{}{}^{3}`$ defines a closed immersion $`X_S^{}(\pi _{S^{}}^{}{}_{}{}^{}_{}^{}{}_{}{}^{3})`$. Refining $`S^{}`$ further if needed we may assume that $`(\pi _{S^{}}^{}{}_{}{}^{}_{}^{}{}_{}{}^{3})`$ is isomorphic with $`_S^{}^{9d+1}`$. Then the inclusion $`X_S^{}(\pi _{S^{}}^{}{}_{}{}^{}_{}^{}{}_{}{}^{3})`$ satisfies the conditions of Proposition 4.2.1 by construction. Hence it corresponds to a morphism $`f_X:S^{}H_{3,d}^{pr}`$ and we have that
$$\pi _{\mathrm{Hilb}}(f_X:S^{}H_{3,d}^{pr})=(\pi _S^{}:X_S^{}S^{},\lambda _S^{}).$$
Step 3: The morphism $`\pi _{\mathrm{Hilb}}`$ is representable and smooth. Let $`S`$ be a scheme and suppose given a morphism $`S_{2d}`$ corresponding to a primitively polarized K3 space $`(\pi :XS,\lambda )`$. We have to show that the product $`S\times _{_{2d}}H_{3d}^{pr}`$ is representable by an algebraic space which is smooth over $`S`$ (via $`pr_1`$). By the surjectivity of $`\pi _{\mathrm{Hilb}}`$ one can find an étale cover $`S^{}`$ of $`S`$ and a projective embedding $`X_S^{}_S^{}^{9d+1}`$, defined by a very ample line bundle $`^3`$. It gives rise to a morphism $`S^{}H_{3,d}^{pr}`$ with
$$\pi _{\mathrm{Hilb}}(S^{}H_{3,d}^{pr})=(X_S^{}S^{},\lambda _S^{})_{2d}(S^{}).$$
We claim that the product $`S^{}\times _{_{2d}}H_{3,d}^{pr}`$ is representable by a scheme isomorphic to $`\mathrm{PGL}(9d+2)_S^{}`$. For any $`S^{}`$-scheme $`U`$ we have that
$$\begin{array}{cc}\hfill S^{}\times _{_{2d}}H_{3,d}^{pr}(U)=& \left\{\right((US^{}),(UH_{3,d}^{pr}),g\left)\right|\hfill \\ & g\mathrm{Hom}((X_UU,\lambda _U),\pi _{\mathrm{Hilb}}(UH_{3,d}^{pr}))\mathrm{in}_{2d}\}\hfill \end{array}$$
where $`\pi _{\mathrm{Hilb}}(UH_{3,d}^{pr})=(𝒳_UU,\lambda _𝒳|_U)`$. Any such morphism $`g`$ gives rise to an isomorphism $`^3𝒪_{𝒳_U}(1)f_U^{}`$ for some invertible sheaf $``$ on $`U`$ and hence an isomorphism
$$(\pi _U^3)(f_U𝒪_{𝒳_U}(1)).$$
But by condition $`(iii)`$ of Proposition 4.2.1 we have an isomorphism
$$(f_U𝒪_{𝒳_U}(1))(pr_U𝒪_{_U^{9d+1}}(1))=^{9d+1}_U$$
and hence we obtain an isomorphism $`(\pi _U^3)_U^{9d+1}`$. This correspondence gives a bijection
$$S^{}\times _{_{2d}}H_{3,d}^{pr}(U)\left\{\mathrm{isomorphisms}(\pi _U^3)_U^{9d+1}\right\}$$
and the right hand set can be identified with $`\mathrm{PGL}(9d+1)_S^{}(U)`$. For this we refer to the arguments given on pp. 101-103 in \[Mum65\]. Hence $`S^{}\times _{_{2d}}H_{3,d}^{pr}`$ is representable by the scheme $`\mathrm{PGL}(9d+1)_S^{}`$ which is smooth over $`S^{}`$.
We will show next that $`S\times _{_{2d}}H_{3,d}^{pr}`$ is a smooth algebraic space over $`S`$. We have a surjective map of étale sheaves
$$S^{}\times _{_{2d}}H_{3,d}^{pr}S\times _{_{2d}}H_{3,d}^{pr}.$$
The product
$$R:=\left(S^{}\times _{_{2d}}H_{3,d}^{pr}\right)\times _{\left(S\times _{_{2d}}H_{3,d}^{pr}\right)}\left(S^{}\times _{_{2d}}H_{3,d}^{pr}\right)$$
can be identified with the smooth $`S`$-scheme $`(S^{}\times _SS^{})\times _{_{2d}}H_{3,d}^{pr}`$. The natural morphism
$$R(S^{}\times _{_{2d}}H_{3,d}^{pr})\times (S^{}\times _{_{2d}}H_{3,d}^{pr})$$
is quasi-compact and the two projection maps
are étale as they correspond to the two étale projection morphisms . Hence $`S\times _{_{2d}}H_{3,d}^{pr}`$ is an algebraic space, which is moreover smooth over $`S`$ as it possesses a smooth atlas $`S^{}\times _{_{2d}}H_{3,d}^{pr}`$ (over $`S`$).
Step 4: Using Remark 4.1.2 (i) in \[LMB00, Ch. 4\] (or Prop. 4.4 in \[DM68\]) and Lemma 4.3.8 we see that the diagonal morphism $`\mathrm{\Delta }:_{2d}_{2d}\times _{}_{2d}`$ is representable, separated and quasi-compact. Then we can apply Theorem 4.21 of \[DM68\] to the morphism $`\pi _{\mathrm{Hilb}}:H_{3,d}^{pr}_{2d}`$ and conclude that $`_{2d}`$ is a Deligne-Mumford stack of finite type over $``$.
Step 5: We will show that the algebraic stack $`_{2d}`$ is separated. As $`_{2d}`$ is of finite type over $``$ one can use the valuative criterion for separateness from \[DM68, Thm. 4.18\] (cf. \[LMB00, Prop. 7.8 and Thm. 7.10\]). It reduces to showing that if $`(\pi _i:X_iS,\lambda _i)`$, for $`i=1,2`$, are two primitively polarized K3 spaces over the spectrum $`S`$ of a discrete valuation ring $`R`$ with field of fractions $`K`$, then every isomorphism $`f:(X_1K,\lambda _1K)(X_2K,\lambda _2K)`$ extends to a $`S`$-isomorphism between $`(X_1,\lambda _1)`$ and $`(X_2,\lambda _2)`$. Note that after taking a finite étale covering of $`S`$ we may assume that:
1. $`X_i`$ are schemes,
2. $`\lambda _i=c_1(_i)`$ for some ample line bundle $`_i`$,
3. $`f`$ gives an isomorphism of pairs $`f:(X_1K,_1K)(X_2K,_2,K)`$.
Then using \[MM64, Thm. 2\] (as a K3 surface is non-ruled) we see that $`f`$ extends uniquely to an isomorphism between $`(X_1,_1)`$ and $`(X_2,_2)`$.
Step 6: We are left to show that the natural inclusion $`_{2d}_{2d}`$ is an open immersion. Take a noetherian scheme $`S`$ and suppose given a morphism $`S_{2d}`$ corresponding to a polarized K3 space $`(\pi :XS,\lambda )`$. Let $`f:S^{}S`$ be an étale covering such that $`\pi _S^{}:X_S^{}S^{}`$ is strongly projective (cf. Step 2 in the proof of Theorem 4.3.3). According to Lemma 3.1.6 the set of points
$$S_{}^{}{}_{}{}^{o}=\{sS^{}|\mathrm{such}\mathrm{that}\lambda _{S^{},s}\mathrm{is}\mathrm{primitive}\}$$
is an open subscheme of $`S^{}`$. The morphism $`f`$ is étale and hence $`f(S_{}^{}{}_{}{}^{o})S`$ is also an open subscheme which represents $`S\times _{_{2d}}_{2d}`$. ∎
###### Remark 4.3.9.
Another possible proof of Theorem 4.3.3 is to use Artin’s criterion (\[LMB00, Cor. 10.11\]). This approach is taken up in \[Ols04, Thm. 6.2\] where M. Olsson constructs a compact stack of “polarized log K3 spaces” over $``$.
An immediate consequence of Theorem 4.3.3 is the existence of a coarse moduli space of polarized K3 surfaces. More precisely Corollary 1.3 in \[KM97\] says
###### Corollary 4.3.10.
The moduli stacks $`_{2d}`$ and $`_{2d}`$ have coarse moduli spaces which are separated algebraic spaces.
Note that this argumet shows that $`_{2d}`$ and $`_{2d}`$ are global quotient stacks.
Before going on we will shortly outline how one can obtain stronger results on coarse moduli schemes of polarized K3 surfaces in characteristic zero.
Approach via periods of K3 surfaces. As we mentioned in the beginning of this section one can use analytic methods to construct a coarse moduli scheme of primitively polarized K3 surfaces. Consider the complex space
$$\mathrm{\Omega }^\pm =\{\omega (L_{2d})|\psi _{2d}(\omega ,\omega )=0\text{and}\psi _{2d}(\omega ,\overline{\omega })>0\}$$
which consists of two connected components. It can be identified with the space
$$\mathrm{SO}(2,19)()/\left(\mathrm{SO}(2)()\times \mathrm{SO}(19)()\right).$$
Let $`\mathrm{\Omega }^+`$ denote one of its connected components, say corresponding to
$$\mathrm{SO}(2,19)()^+/\left(\mathrm{SO}(2)()\times \mathrm{SO}(19)()\right),$$
where $`\mathrm{SO}(2,19)()^+`$ is the connected component of $`\mathrm{SO}(2,19)()`$ containing the identity. It is a bounded symmetric domain of type IV and of dimension 19. Let $`\mathrm{\Gamma }`$ be the group $`\{g\mathrm{O}(V_0)()|g(e_1df_1)=e_1df_1\}`$ and denote by $`\mathrm{\Gamma }^+`$ the subgroup of $`\mathrm{\Gamma }`$ of index 2 which consists of isometries preserving the connected components of $`\mathrm{\Omega }^\pm `$. Then $`\mathrm{\Gamma }^+`$ acts on $`\mathrm{\Omega }^+`$ properly discontinuously and the space $`\mathrm{\Omega }^+/\mathrm{\Gamma }^+`$ is a coarse moduli scheme for primitively quasi-polarized complex K3 surfaces of degree $`2d`$. There is an open part $`\mathrm{\Omega }^0`$ of $`\mathrm{\Omega }^+`$ such that $`\mathrm{\Omega }^0/\mathrm{\Gamma }^+`$ is a coarse moduli scheme for primitively polarized complex K3 surfaces of degree $`2d`$. For details and proofs we refer to \[BBD85, Exp. XIII\]. The existence of a coarse moduli scheme is Proposition 8 in loc. cit..
Approach via geometric invariant theory. Let $`k`$ be an algebraically closed field of characteristic zero. Then using the techniques of \[Vie95, Ch. 8\], and more precisely §8.2 (see Theorem 8.23), one can prove that the moduli functor $`_{2d}k`$ (respectively $`_{2d}k`$) has a quasi-projective coarse moduli scheme over $`k`$. Indeed, one has that Assumptions 8.22 in \[Vie95, §8.2\] are satisfied:
1. The functor is locally closed. This follows from the proof of Proposition 4.2.1.
2. The separateness property is shown in Step 2 of the proof of Theorem 4.3.3.
3. The functor is bounded by Theorem 1.2.2. See also Remark 8.24 in loc. cit. and note that the condition ‘$`\omega ^2`$ is trivial’ is a locally closed condition.
One actually shows that the scheme in question is $`H_{3,d}^{pr}k/\mathrm{PGL}(N)_k`$ (respectively $`H_{3,d}k/\mathrm{PGL}(N)_k`$) for a suitable $`N`$.
Combining the approach to coarse moduli schemes via geometric invariant theory and Corollary 4.3.10 we conclude that $`_{2d,}`$ (respectively $`_{2d,}`$) has a quasi-projective coarse moduli scheme.
###### Proposition 4.3.11.
The moduli stacks $`_{2d}`$ and $`_{2d}`$ are smooth of relative dimension 19 over $`[\frac{1}{2d}]`$.
###### Proof.
According to \[LMB00, Prop. 4.15\] we have to show that for any strictly henselian local ring $`R`$ and surjection $`\mathrm{Spec}(R)\mathrm{Spec}(R^{})`$ defined by a nilpotent sheaf of ideals one has that the natural map
$$\mathrm{Hom}(\mathrm{Spec}(R^{}),_{2d,[1/2d]})\mathrm{Hom}(\mathrm{Spec}(R),_{2d,[1/2d]})$$
is surjective. Since $`R`$ is strictly henselian every K3 space over $`\mathrm{Spec}(R)`$ is a K3 scheme and the same holds for spaces over $`\mathrm{Spec}(R^{})`$ (see \[GD67, EGA IV, 18.1.2\]). Hence by Lemma 4.1.3 we conclude that $`_{2d,[1/2d]}`$ is smooth over $`[1/2d]`$.
The same argument applies also to the dimension claim. Since every K3 space over $`\mathrm{Spec}(k[ϵ]/ϵ^2)`$ is a K3 scheme we conclude from Theorem 4.1.2 that the dimension of $`_{2d,[1/2d]}`$ at every point is 19.
This proof also shows that $`_{2d}`$ is smooth of relative dimension 19. ∎
###### Remark 4.3.12.
Since smoothness will be essential for all our further considerations, unless explicitly stated, by $`_{2d}`$ (respectively $`_{2d}`$) we will mean the smooth stack $`_{2d}_{}[\frac{1}{2d}]`$ (respectively $`_{2d}_{}[\frac{1}{2d}]`$) over $`\mathrm{Spec}([\frac{1}{2d}])`$.
We will end this section speculating about other possible moduli spaces and functors of polarized K3 surfaces. Note first that one could have started with a moduli functor $`_{2d}^{}`$ of (primitively) polarized K3 schemes of degree $`2d`$. The problem we came up with restricting only to schemes was proving effectiveness of descent for K3 schemes. For this reason one takes the “étale sheafification” of $`_{2d}^{}`$ considering (primitively) polarized K3 spaces. This makes the descent obstruction essentially trivial.
Next, one can consider deformations of polarized K3 surfaces as in Section 4.1 by algebraic spaces and not only schemes. For a polarized K3 surface $`(X_0,\lambda _0)`$ over an algebraically closed field $`k`$ define
$$\mathrm{Def}_{\mathrm{AlgSp}}(X_0,\lambda _0):\underset{¯}{A}\mathrm{Sets}$$
to be the functor sending an object $`A`$ of $`\underset{¯}{A}`$ to the isomorphism classes of triples $`(𝒳,\lambda ,\varphi _0)`$ where $`(𝒳\mathrm{Spec}(A),\lambda )`$ is a polarized K3 space and $`\varphi _0`$ is an isomorphism $`\varphi _0:(𝒳,)_Ak(X_0,_0)`$. Combining Theorem 4.3.3, Lemma 4.3.11 and \[LMB00, Cor. 10.11\] we conclude that $`\mathrm{Def}_{\mathrm{AlgSp}}`$ is pro-representable, formally smooth and of dimension $`19`$.
## 5. Level Structures of Polarized K3 Surfaces
Recall that for an abelian scheme $`(A,\lambda )`$ over a base scheme $`S`$ and a natural number $`n`$ which is invertible in $`S`$ one defines a (Jacobi) level $`n`$-structure on $`A`$ to be an isomorphism $`\theta :A[n](/n)_S`$ of étale sheaves on $`S`$ satisfying some further properties. In other words, one uses the Tate module of an abelian variety in order to define level structures. For a K3 surface $`X`$ we will use the same idea applied to $`H_{\mathrm{et}}^2(X_{\overline{k}},_l(1))`$. More precisely, we will introduce the notion of level structures on primitively polarized K3 surfaces of degree $`2d`$ corresponding to open compact subgroups of $`\mathrm{SO}(V_{2d},\psi _{2d})(𝔸_f)`$ (see below) and define moduli spaces of primitively polarized K3 surfaces with level structures. We set up some notations first.
* All schemes in this section will be assumed to be locally noetherian.
* For a finite set of primes $`\mathrm{B}=\{\mathrm{p}_1,\mathrm{},\mathrm{p}_\mathrm{r}\}`$ we denote by $`_\mathrm{B}`$ the product $`_{p\mathrm{B}}_p`$ and by $`N_\mathrm{B}`$ the product of the primes in $`\mathrm{B}`$.
* We fix a natural number $`d`$. We shall use the notations $`L_{2d,\mathrm{B}}`$ and $`L_{0,\mathrm{B}}`$ for the quadratic lattices $`L_{2d}_\mathrm{B}`$ and $`L_0_\mathrm{B}`$ (cf. Section 2.1).
* Let $`𝕂\mathrm{SO}(V_{2d})(\widehat{})`$ be a subgroup of finite index and let $`\mathrm{B}=\{\mathrm{p}_1,\mathrm{},\mathrm{p}_\mathrm{r}\}`$ be the set of prime divisors of $`2d`$ and primes $`p`$ for which $`𝕂_p\mathrm{SO}(V_{2d})(_p)`$. We denote by $`𝕂_\mathrm{B}`$ the product $`_{p\mathrm{B}}𝕂_p`$.
### 5.1. Level Structures
Let $`S`$ be a connected scheme over $`[\frac{1}{p_1\mathrm{}p_r}]`$ and suppose given a polarized K3 space $`(\pi :XS,\lambda )`$ of degree $`2d`$. Let $`P_{\mathrm{et}}^2\pi _{}_\mathrm{B}(1)`$ be the sheaf of primitive cohomology i.e., the orthogonal complement of $`c_1(\lambda )`$ in $`R_{\mathrm{et}}^2\pi _{}_\mathrm{B}(1)`$. Take a geometric point $`\overline{b}`$ of $`S`$ and let $`\overline{b}:\mathrm{Spec}(k(\overline{b}))S`$ be the corresponding morphism of schemes. Consider the free $`_\mathrm{B}`$-module of rank 21
$$P^2(\overline{b}):=\overline{b}^{}P_{\mathrm{et}}^2\pi _{}_\mathrm{B}(1)$$
i.e., the fiber of $`P_{\mathrm{et}}^2\pi _{}_\mathrm{B}(1)`$ at $`\overline{b}`$ with its action of $`\pi _1^{\mathrm{alg}}(S,\overline{b})`$ and the bilinear form $`\psi _{\lambda ,_\mathrm{B}}`$.
Suppose given an class $`\alpha _{\overline{b}}`$ in the set
$$\left\{𝕂_\mathrm{B}\backslash \text{Isometry}(L_{2d,_\mathrm{B}},P^2(\overline{b}))\right\}^{\pi _1^{\mathrm{alg}}(S,\overline{b})}$$
where $`𝕂_\mathrm{B}`$ acts on $`\text{Isometry}(L_{2d,_\mathrm{B}},P^2(\overline{b}))`$ on the right via its action on $`L_{2d,_\mathrm{B}}`$ and $`\pi _1^{\mathrm{alg}}(S,\overline{b})`$ acts on the left via its action on $`P^2(\overline{b})`$. Let $`\overline{b}^{}`$ be another geometric point in $`S`$. The $`\alpha _{\overline{b}}`$ determines uniquely a class in
$$\left\{𝕂_\mathrm{B}\backslash \text{Isometry}(L_{2d,_\mathrm{B}},P^2(\overline{b}^{}))\right\}^{\pi _1^{\mathrm{alg}}(S,\overline{b}^{})}$$
in the following way: One can find an isomorphism
(7)
$$\delta _\pi :\pi _1^{\mathrm{alg}}(S,\overline{b})\pi _1^{\mathrm{alg}}(S,\overline{b}^{})$$
and an isometry
$$\delta _{\mathrm{et}}:H_{\mathrm{et}}^2(X_{\overline{b}},_\mathrm{B}(1))H_{\mathrm{et}}^2(X_{\overline{b}^{}},_\mathrm{B}(1))$$
determined uniquely by $`\delta _\pi `$, mapping $`c_1(\lambda _{\overline{b}})`$ to $`c_1(\lambda _{\overline{b}^{}})`$, such that $`\delta _{\mathrm{et}}(\gamma x)=\delta _\pi (\gamma )\delta _{\mathrm{et}}(x)`$ for every $`xH_{\mathrm{et}}^2(X_{\overline{b}},_\mathrm{B}(1))`$ and $`\gamma \pi _1^{\mathrm{alg}}(S,\overline{b})`$. The isometry $`\delta _{\mathrm{et}}`$ defines an isometry between $`P^2(\overline{b})`$ and $`P^2(\overline{b}^{})`$ which we will denote again by $`\delta _{\mathrm{et}}`$. Let $`\stackrel{~}{\alpha }`$ be a representative of the class $`\alpha _{\overline{b}}`$. Then the class $`\alpha _{\overline{b}^{}}`$ of $`\delta _{\mathrm{et}}\stackrel{~}{\alpha }`$ in $`𝕂_\mathrm{B}\backslash \text{Isometry}(L_{2d,\mathrm{B}},P^2(\overline{b}^{}))`$ is $`\pi _1^{\mathrm{alg}}(S,\overline{b}^{})`$-invariant. Any other representative $`\stackrel{~}{\alpha }_1`$ of $`\alpha _{\overline{b}}`$ differs by an element in $`𝕂_\mathrm{B}`$ and hence gives rise to the same class $`\alpha _{\overline{b}^{}}`$ in $`𝕂_\mathrm{B}\backslash \text{Isometry}(L_{2d,\mathrm{B}},P^2(\overline{b}^{}))`$.
Any two isomorphisms (7) differ by an inner automorphism of $`\pi _1^{\mathrm{alg}}(S,\overline{b})`$ and therefore we see that that class of $`\delta _{\mathrm{et}}\stackrel{~}{\alpha }`$ is independent of the choice of an isomorphism (7).
This remark allows us to make the following definition.
###### Definition 5.1.1.
A *level $`𝕂`$-structure* on a primitively polarized K3 space $`(\pi :XS,\lambda )`$ over a connected scheme $`S(\mathrm{Sch}/[1/p_1\mathrm{}p_r])`$ is an element of the set
$$\left\{𝕂_\mathrm{B}\backslash \text{Isometry}(L_{2d,_\mathrm{B}},P^2(\overline{b}))\right\}^{\pi _1^{\mathrm{alg}}(S,\overline{b})}.$$
The group $`𝕂_\mathrm{B}`$ acts on $`\text{Isometry}(L_{2d,_\mathrm{B}},P^2(\overline{b}))`$ on the right via its action on $`L_{2d,_\mathrm{B}}`$ and $`\pi _1^{\mathrm{alg}}(S,\overline{b})`$ acts on the left via its action on $`P^2(\overline{b})`$. In general, a level $`𝕂`$-structure on $`(\pi :XS,\lambda )`$ is a level $`𝕂`$-structure on each connected component of $`S`$.
If $`\stackrel{~}{\alpha }:L_{2d,\mathrm{B}}P_{et}^2(\overline{b})`$ is a representative of the class $`\alpha `$, then via the isomorphism
$$\stackrel{~}{\alpha }^{\mathrm{ad}}:\mathrm{O}(V_{2d})(_\mathrm{B})\mathrm{O}(P^2(\overline{b}))(_\mathrm{B})$$
the monodromy action
$$\rho :\pi _1^{\mathrm{alg}}(S,\overline{b})\mathrm{O}(P^2(\overline{b}))(_\mathrm{B})$$
factorizes through $`\stackrel{~}{\alpha }^{\mathrm{ad}}(𝕂_\mathrm{B})`$.
###### Remark 5.1.2.
If all residue fields of the points in $`S`$ in Definition 5.1.1 are of characteristic zero, then one can define a level $`𝕂`$-structure to be an element of set
$$\left\{𝕂\backslash \text{Isometry}(L_{2d,\widehat{}},P^2(\overline{b}))\right\}^{\pi _1^{\mathrm{alg}}(S,\overline{b})}$$
where $`P^2(\overline{b}):=\overline{b}^{}P_{\mathrm{et}}^2\pi _{}\widehat{}(1)`$.
We will consider two important examples of level structures on primitively polarized K3 spaces.
###### Example 5.1.3.
Fix a natural number $`n`$ and consider the group
$$𝕂_n=\left\{\gamma \mathrm{SO}(V_{2d})(\widehat{})\right|\gamma 1(modn)\}.$$
Then the set $`\mathrm{B}`$ consists of the prime divisors of $`2dn`$. We will give a direct interpretation of level $`𝕂_n`$-structures.
Let $`S`$ be a scheme over $`[1/2dn]`$ and consider a primitively polarized K3 space $`(\pi :XS,\lambda )`$ of degree $`2d`$. As usual we denote by $`P_{\mathrm{et}}^2\pi _{}(/n)(1)`$ the orthogonal complement of $`c_1(\lambda )`$ in $`R_{\mathrm{et}}^2\pi _{}(/n)(1)`$ with respect to the bilinear form $`\psi _n=\psi _{}/n`$. Then a level $`𝕂_n`$-structure amounts to giving an isomorphism
$$\alpha _N:(P_{\mathrm{et}}^2\pi _{}(/n)(1),\psi _{,n})(L_{2d,/n},\psi _{2d,/n})_S$$
of étale sheaves on $`S`$, where $`(L_{2d,/n},\psi _{2d,/n})_S`$ is the constant polarized étale sheaf over $`S`$ with fibers $`(L_{2d},\psi _{2d})/n`$.
We will call level a $`𝕂_n`$-structure on $`X`$ simply a *level $`n`$-structure*.
###### Example 5.1.4.
Let $`G`$ be the algebraic group $`\mathrm{SO}(V_{2d})`$ over $``$. Consider the *even Clifford algebra* $`C^+(V_{2d},\psi _{2d})`$ over $``$ and let $`G_1`$ be the *even Clifford group* over $``$. In other words we set
$$G_1=\mathrm{CSpin}(V_{2d})=\left\{gC^+(V_{2d})^{}\right|gV_{2d}g^1=V_{2d}\}.$$
The natural homomorphism of linear algebraic groups $`G_1G`$ given by $`g(vgvg^1)`$ fits into an exact sequence (see \[Del72, §3.2\])
$$0𝔾_mG_1G0.$$
Set $`G_1()`$ to be $`G_1()C^+(L_{2d})^{}`$. We have an exact sequence (see \[And96, §4.4\])
(8)
$$0/2G_1()G().$$
For a natural number $`n`$ denote
$$\mathrm{\Gamma }_n=\left\{\gamma G()\right|\gamma 1(modn)\}$$
and
$$\mathrm{\Gamma }_n^{\mathrm{sp}}=\{\gamma G_1()|\gamma 1(modn)\mathrm{in}C^+(L_{2d})\}.$$
If $`n>2`$, then $`\mathrm{\Gamma }_n`$ and $`\mathrm{\Gamma }_n^{\mathrm{sp}}`$ are torsion free. Hence one sees from the exact sequence (8) that $`\mathrm{\Gamma }_n^{\mathrm{sp}}`$ is isomorphic with its image $`\mathrm{\Gamma }_n^\mathrm{a}`$ in $`G()`$.
Consider the group
$$𝕂_n^{\mathrm{sp}}=\{\gamma G_1(\widehat{})|\gamma 1(modn)\mathrm{in}C^+(L_{2d,\widehat{}})\}.$$
We have that $`𝕂_n^{\mathrm{sp}}G_1()=\mathrm{\Gamma }_n^{\mathrm{sp}}`$. Moreover the image $`𝕂_n^a`$ of $`𝕂_n^{\mathrm{sp}}`$ in $`G(\widehat{})`$ is of finite index. Indeed, for every $`l`$ not dividing $`2nd`$, the $`l`$-component of $`𝕂_n^\mathrm{a}`$ is $`G(_l)`$ as shown in \[And96, §4.4\]. Hence the set $`\mathrm{B}`$ for $`𝕂_n^\mathrm{a}`$ is the set of prime divisors of $`2dn`$.
We consider polarized K3 surfaces with level $`𝕂_n^\mathrm{a}`$-structure. Note that this level structure is in general finer than level $`𝕂_n`$-structure as $`𝕂_n^\mathrm{a}𝕂_n`$ is of finite index. We will call it *spin level $`n`$-structure*.
### 5.2. Motivation
We will pause here and give a motivation for the rest of the definitions we make in this section. So far we have defined level $`𝕂`$-structures using the primitive second étale cohomology group of a polarized K3 surface. Using these level structures one can define moduli stacks $`_{2d,𝕂}`$ of primitively polarized K3 surfaces of degree $`2d`$ with a level $`𝕂`$-structure and show that they are algebraic spaces (cf. Theorem 6.1.2 below). Over $``$, we can relate these spaces to the orthogonal Shimura variety associated to the group $`\mathrm{SO}(2,19)`$. More precisely in Chapter 3, Section 3.4.2 of \[Riz05\] we define a period morphism
$$j_{d,𝕂,}:_{2d,𝕂,}Sh_𝕂(\mathrm{SO}(2,19),\mathrm{\Omega }^\pm )_{}$$
which is étale. This is similar to the case of moduli of abelian varieties where one can identify $`𝒜_{g,1,n}`$ with $`Sh_{\mathrm{\Lambda }_n}(\mathrm{CSp}_{2g},_g^\pm )_{}`$. In general, due to the fact that the injective homomorphism (2)
$$i^{\mathrm{ad}}:\{g\mathrm{SO}(V_0)()|g(e_1df_1)=e_1df_1\}\mathrm{SO}(V_{2d})()$$
defined in Section 2.1 is not surjective, the period map $`j_{d,𝕂,}`$ need not be injective. In order to construct an injective period morphism we will define level structures using the “full” second étale cohomology group of a K3 surface. We use these full level structures in \[Riz05, Ch. 3\] to show that every complex K3 surface with complex multiplication by a CM-field $`E`$ is defined over an abelian extension of $`E`$.
### 5.3. Full Level Structures
The inclusion of lattices $`i:L_{2d}L_0`$ (see Section 2.1) defines injective homomorphisms of groups
$$i^{\mathrm{ad}}:\{g\mathrm{O}(V_0)(\widehat{})|g(e_1df_1)=e_1df_1\}\mathrm{O}(V_{2d})(\widehat{})$$
and
$$i^{\mathrm{ad}}:\{g\mathrm{SO}(V_0)(\widehat{})|g(e_1df_1)=e_1df_1\}\mathrm{SO}(V_{2d})(\widehat{}).$$
###### Definition 5.3.1.
A subgroup $`𝕂\mathrm{SO}(V_{2d})(\widehat{})`$ of finite index is called *admissible* if it is contained in the image
$$i^{\mathrm{ad}}\left(\{g\mathrm{SO}(V_0)(\widehat{})|g(e_1df_1)=e_1df_1\}\right)\mathrm{SO}(V_{2d})(\widehat{}).$$
If $`𝕂`$ is an admissible subgroup of $`\mathrm{SO}(V_{2d})(\widehat{})`$ then all its subgroups of finite index $`𝕂^{}𝕂`$ are also admissible.
###### Example 5.3.2.
The group $`𝕂_{2d}`$ is admissible. Hence all its subgroups of finite index are admissible, as well.
###### Example 5.3.3.
If $`d=1`$ then $`𝕂_n`$ is admissible for any $`n2`$.
Let $`𝕂`$ be an admissible subgroup of $`\mathrm{SO}(V_{2d})(\widehat{})`$ and let $`\mathrm{B}`$ be the set, consisting of all prime divisors of $`2d`$ and, of the primes $`p`$ for which $`𝕂_p\mathrm{SO}(V_{2d})(_p)`$. Using the notations introduced before Definition 5.1.1 we set
$$H^2(\overline{b}):=\overline{b}^{}R_{\mathrm{et}}^2\pi _{}_\mathrm{B}(1).$$
In order to simplify the notations we will identify a subgroup of $`\{g\mathrm{SO}(V_0)(\widehat{})|g(e_1df_1)=e_1df_1\}`$ with its image in $`\mathrm{SO}(V_{2d})(\widehat{})`$ under the injective homomorphism $`i^{\mathrm{ad}}`$.
###### Definition 5.3.4.
A *full level $`𝕂`$-structure* on a primitively polarized K3 space $`(\pi :XS,\lambda )`$ over a connected scheme $`S(\mathrm{Sch}/[1/p_1\mathrm{}p_r])`$ is an element of the set
$$\left\{𝕂_\mathrm{B}\backslash \left\{g\text{Isometry}(L_{0,_\mathrm{B}},H^2(\overline{b}))\right|g(e_1df_1)=c_1(\lambda _{\overline{b}})\}\right\}^{\pi _1^{\mathrm{alg}}(S,\overline{b})}.$$
The group $`𝕂_\mathrm{B}`$ acts on $`\left\{g\text{Isometry}(L_{0,_\mathrm{B}},H^2(\overline{b}))\right|g(e_1df_1)=c_1(\lambda _{\overline{b}})\}`$ on the right via its action on $`L_{0,_\mathrm{B}}`$ and $`\pi _1^{\mathrm{alg}}(S,\overline{b})`$ acts on the left via its action on $`H^2(\overline{b})`$. A full level $`𝕂`$-structure on $`(\pi :XS,\lambda )`$ over a general base $`S`$ is a full level $`𝕂`$-structure on each connected component of $`S`$.
Again, a class $`\alpha _{\overline{b}}`$ for a geometric point $`\overline{b}`$ as above determines uniquely a class $`\alpha _{\overline{b}^{}}`$ for any other geometric point $`\overline{b}^{}`$. If $`\stackrel{~}{\alpha }:L_{0,\mathrm{B}}H^2(\overline{b})`$ is a representative of the class $`\alpha `$, then via the isomorphism
$$\stackrel{~}{\alpha }^{\mathrm{ad}}:\mathrm{O}(V_0)(_\mathrm{B})\mathrm{O}(H^2(\overline{b}))(_\mathrm{B})$$
the monodromy action $`\rho :\pi _1^{\mathrm{alg}}(S,\overline{b})\mathrm{O}(H^2(\overline{b}))(_\mathrm{B})`$ factorizes through $`\stackrel{~}{\alpha }^{\mathrm{ad}}(𝕂_\mathrm{B})`$.
###### Example 5.3.5.
Let $`n3`$ be an integer. Define the group
$$𝕂_n^{\mathrm{full}}=\left\{g\mathrm{SO}(V_0)(\widehat{})\right|g(e_1df_1)=e_1df_1\mathrm{and}g1(modn)\}.$$
By definition it is an admissible subgroup of $`\mathrm{SO}(V_{2d})(\widehat{})`$. Let $`S`$ be a scheme over $`[1/2dn]`$ and consider a K3 space $`(\pi :XS,\lambda )`$ with a primitive polarization of degree $`2d`$. Then a full level $`𝕂_n^{\mathrm{full}}`$-structure amounts to giving an isomorphism
$$\alpha _N:(R_{\mathrm{et}}^2\pi _{}(/n)(1),\psi )(L_{0,/n},\psi _{0,/n})_S$$
of étale sheaves on $`S`$, where $`(L_{0,/n},\psi _{0,/n})_S`$ is the constant polarized étale sheaf over $`S`$ with fibers $`(L_0,\psi _0)/n`$.
We will call a full level $`𝕂_n^{\mathrm{full}}`$-structure on $`X`$ simply a *full level $`n`$-structure*.
## 6. Moduli Spaces of Polarized K3 Surfaces with a Level Structure
In this section we will use the notion of a (full) level structure level structure to define moduli functors of primitively polarized K3 spaces with a (full) level structure. Using Artin’s criterion and Proposition 3.4.2 we will show that these functors are representable by algebraic spaces over open parts of $`\mathrm{Spec}()`$.
We shall be using the notations established in the beginning of Section 5.1. In particular we fix a natural number $`d`$. To a subgroup $`𝕂`$ of $`\mathrm{SO}(V_{2d})(\widehat{})`$ we associated a finite set of primes $`\mathrm{B}`$ and $`N_\mathrm{B}`$ will denote the product of these primes.
### 6.1. Moduli of K3 Surfaces with Level Structure
Let $`𝕂`$ be a subgroup of $`\mathrm{SO}(\widehat{})`$ of finite index. We will assume further that it is contained in $`𝕂_n`$ for some $`n3`$. Let $`𝒳_1=(\pi _1:X_1S_1,\lambda _1)`$ and $`𝒳_2=(\pi _2:X_2S_2)`$ be two objects of $`_{2d}`$. Suppose that $`S_1`$ and $`S_2`$ are connected and let $`(f,f_S)\mathrm{Hom}(𝒳_1,𝒳_2)`$ (in $`_{2d}`$). Let $`\overline{b}_1`$ and $`\overline{b}_2`$ be two geometric points of $`S_1`$ and $`S_2`$ such that $`f_S(\overline{b}_1)=\overline{b}_2`$. Then the morphism $`f`$ defines a homomorphism $`f_{\mathrm{et}}^{}:P^2(\overline{b}_2)P^2(\overline{b}_1)`$. Hence we obtain a map
$$f^{}:𝕂_\mathrm{B}\backslash \text{Isometry}(L_{2d,_{0,\mathrm{B}}},P^2(\overline{b}_2))𝕂_\mathrm{B}\backslash \text{Isometry}(L_{2d,_\mathrm{B}},P^2(\overline{b}_1))$$
given by $`\alpha f_{\mathrm{et}}^{}\alpha `$ and commuting with the monodromy actions on both sides.
###### Definition 6.1.1.
For $`d`$ and $`𝕂`$ as above consider the category $`_{2d,𝕂}`$ defined in the following way:
1. Triples $`(\pi :XS,\lambda ,\alpha )`$ of a K3 space $`\pi :XS`$ with a primitive polarization $`\lambda `$ of degree $`2d`$ and with a level $`𝕂`$-structure $`\alpha `$ on $`(\pi :XS,\lambda )`$.
2. Suppose given two triples $`𝒳_1=(\pi _1:X_1S_1,\lambda _1,\alpha _1)`$ and $`𝒳_2=(\pi _2:X_2S_2,\lambda _2,\alpha _2)`$. Let $`f_S:S_1S_2`$ be a morphism of schemes. Choose base geometric points $`\overline{b}_1^{}`$ and $`\overline{b}_2^{}`$ on any two connected components $`S_1^{}`$ and $`S_2^{}`$ of $`S_1`$ and $`S_2`$ for which $`f:S_1^{}S_2^{}`$ such that $`f_S(\overline{b}_1^{})=\overline{b}_2^{}`$. Define the morphisms between $`𝒳_1`$ and $`𝒳_2`$ in the following way
$$\begin{array}{cc}\hfill \mathrm{Hom}(𝒳_1,𝒳_2)=\{\mathrm{pairs}(f_S,f)|& f_S:S_1S_2\mathrm{is}\mathrm{a}\mathrm{morph}.\mathrm{of}\mathrm{spaces},\hfill \\ & f:X_1X_2\times _{S_2,f_S}S_1\mathrm{is}\mathrm{an}\mathrm{isom}.\mathrm{of}\hfill \\ & S_1\mathrm{spaces}\mathrm{with}f^{}\lambda _2=\lambda _1\mathrm{and}\hfill \\ & f^{}(\alpha _1)=\alpha _2\mathrm{on}\mathrm{any}\mathrm{conn}.\mathrm{cmpt}.\mathrm{of}S_1\}\hfill \end{array}$$
Next we define three projection functors.
1. Consider the following forgetful functor
$$pr_{_{2d,𝕂}}:_{2d,𝕂}(\mathrm{Sch}/[1/N_\mathrm{B}])$$
sending a triple $`(\pi :XS,\lambda ,\alpha )`$ to $`S`$. It makes $`_{2d,𝕂}`$ into a category over $`(\mathrm{Sch}/[1/N_\mathrm{B}])`$.
2. For any $`𝕂`$, satisfying the assumptions of the beginning of the section, one has a projection functor
(9)
$$pr_𝕂:_{2d,𝕂}_{2d,[1/N_\mathrm{B}]}$$
sending a triple $`(\pi :XS,\lambda ,\alpha )`$ to $`(\pi :XS,\lambda )`$ and an element $`(f,f_S)\mathrm{Hom}(𝒳,𝒴)`$ of $`_{2d,𝕂}`$ to $`(f,f_S)`$.
3. For any two subgroups $`𝕂_1𝕂_2`$ of finite index in $`\mathrm{SO}(V_{2d})(\widehat{})`$ (contained in some $`𝕂_n`$ for $`n3`$) one has a projection functor
(10)
$$pr_{(𝕂_1,𝕂_2)}:_{2d,𝕂_1,[1/N_{\mathrm{B}_1\mathrm{B}_2}]}_{2d,𝕂_2,[1/N_{\mathrm{B}_1\mathrm{B}_2}]}.$$
It sends an object $`(XS,\lambda ,\alpha _{𝕂_1})`$ to $`(XS,\lambda ,\alpha _{𝕂_2})`$ where $`\alpha _{𝕂_2}`$ is the class of $`\alpha _{𝕂_1}`$ in $`𝕂_{2,\mathrm{B}}\backslash \text{Isometry}(L_{2d,_\mathrm{B}},P^2(\overline{b}))`$. Morphism of $`_{2d,𝕂_1,[1/N_{\mathrm{B}_1\mathrm{B}_2}]}`$ are mapped to morphism of $`_{2d,𝕂_2,[1/N_{\mathrm{B}_1\mathrm{B}_2}]}`$ in the obvious way.
From the definitions of the functors we see that $`pr_{𝕂_1}=pr_{(𝕂_1,𝕂_2)}pr_{𝕂_2}`$ over $`[1/N_{\mathrm{B}_1\mathrm{B}_2}]`$.
###### Theorem 6.1.2.
The category $`_{2d,𝕂}`$ is a separated algebraic space over $`[1/N_\mathrm{B}]`$. It is smooth of relative dimension 19 and the forgetful morphism (9)
$$pr_𝕂:_{2d,𝕂}_{2d,[1/N_\mathrm{B}]}$$
is finite and étale.
###### Proof.
We divide the proof into several steps.
Step 1: The category $`_{2d,𝕂}`$ is a stack. The proof goes exactly in the same lines as the one of Lemma 4.3.7. We will use Artin’s criterion (cf. \[LMB00, Cor. 10.11\]) to show that $`_{2d,𝕂}`$ is an algebraic space.
We claim that the diagonal morphism $`\mathrm{\Delta }:_{2d,𝕂}_{2d,𝕂}\times _{[1/N_\mathrm{B}]}_{2d,𝕂}`$ is representable, separated and of finite type. By Remark 4.1.2 in \[LMB00\] it is equivalent to showing that for any two objects $`𝒳=(XS,\lambda _X,\alpha _X)`$ and $`𝒴=(YS,\lambda _Y,\alpha _Y)`$ the functor $`\mathrm{Isom}_S(𝒳,𝒴)`$ has these properties. We will prove first the following result.
###### Lemma 6.1.3.
For any object $`𝒳`$ of $`_{2d,𝕂}`$ we have that $`\mathrm{Aut}_S(𝒳)=\{\mathrm{id}_𝒳\}`$.
###### Proof.
By assumption the group $`𝕂`$ is contained in $`𝕂_n`$ for some $`n3`$. Hence a level $`𝕂`$-structure on a primitively polarized K3 space $`(XS,\lambda )`$ defines in a natural way (using the functor $`pr_{(𝕂,𝕂_n)}`$) a level $`n`$-structure $`\alpha _n`$ on $`X`$. We have that
$$\mathrm{Aut}_S\left((XS,\lambda ,\alpha )\right)(U)\mathrm{Aut}_S\left((XS,\lambda ,\alpha _n)\right)(U)$$
for an $`S`$-scheme $`U`$ hence it is enough to prove the lemma assuming that $`𝕂=𝕂_n`$.
Let $`𝒳=(XS,\lambda ,\alpha )`$ be an object in $`_{2d,𝕂}`$, let $`f\mathrm{Aut}_S(𝒳)(U)`$ and assume that $`U`$ is connected. Take a geometric point $`\overline{b}:\mathrm{Spec}(\mathrm{\Omega })U`$. Then for the finite set $`\mathrm{B}=\{\mathrm{the}\mathrm{prime}\mathrm{divisors}\mathrm{of}`$n$`\}`$ the morphism $`f`$ induces an automorphism
$$f_{\mathrm{et}}^{}:H_{\mathrm{et}}^2(X_{\overline{b}},_\mathrm{B}(1))H_{\mathrm{et}}^2(X_{\overline{b}},_\mathrm{B}(1))$$
fixing $`c_1(\lambda _{\overline{b}})`$ and such that
$$f_{\mathrm{et}}^{}:P_{\mathrm{et}}^2(X_{\overline{b}},/n(1))P_{\mathrm{et}}^2(X_{\overline{b}},/n(1))$$
is the identity (cf. Example 5.1.3). As the automorphism $`f`$ is of finite order we have that $`f_{\mathrm{et}}^{}\mathrm{O}\left(P_{\mathrm{et}}^2(X_{\overline{b}},_\mathrm{B}(1))\right)`$ is semi-simple and its eigenvalues are roots of unity. We have further that $`f_{\mathrm{et}}^{}1(modn)`$ so we conclude by \[Mum74, Ch. IV, Application II, p. 207, Lemma\] that $`f_{\mathrm{et}}^{}`$ is the identity automorphism of $`P_{\mathrm{et}}^2(X_{\overline{b}},_\mathrm{B}(1))`$. As it fixes $`c_1(\lambda _{\overline{b}})`$ we see that it acts as the identity on $`H_{\mathrm{et}}^2(X_{\overline{b}},_\mathrm{B}(1))`$. Therefore by Proposition 3.4.2 we that $`f=\mathrm{id}_{X_{\overline{b}}}`$. As the geometric point $`\overline{b}`$ can be chosen arbitrary we have that $`f=\mathrm{id}_{𝒳_U}`$. ∎
We see from the lemma that for a $`S`$-scheme $`U`$ the set $`\mathrm{Isom}_S(𝒳,𝒴)(U)`$ is either empty or it consists of one element. Indeed, suppose that $`f_i\mathrm{Isom}_S(𝒳,𝒴)(U)`$ for $`i=1,2`$. Then the composition $`f_2^1f_1`$ belongs to $`\mathrm{Aut}_S(𝒳)(U)`$ and hence it is the identity. This shows that $`\mathrm{Isom}_S(𝒳,𝒴)`$ is representable and of finite type. The fact that it is unramified and separated over $`S`$ follows from Lemma 4.3.8 as one has that
$$\mathrm{Isom}_S(𝒳,𝒴)(U)\mathrm{Isom}_S((XS,\lambda _X),(YS,\lambda _Y))(U).$$
Next we claim that the stack $`_{2d,𝕂}`$ is locally of finite presentation. This follows from \[AGV71, Exposé IX, 2.7.4\] and the fact that $`_{2d}`$ is locally of finite presentation. Conditions (iii) and (iv) of \[LMB00, Cor. 10.11\] follow from the corresponding properties of $`_{2d}`$ and the fact that for any small surjection of rings $`RR^{}`$ the category of étale schemes over $`R`$ is equivalent to the category of étale schemes over $`R^{}`$ (\[GD67, EGA IV, 18.1.2\]).
Thus $`_{2d,𝕂}`$ is an algebraic stack. As $`\mathrm{Aut}_S(𝒳)=\{\mathrm{id}_𝒳\}`$ for any object we have that $`_{2d,𝕂}`$ is an algebraic space (\[LMB00, Cor. 8.1.1\]).
Step 2: We will show that the morphism of algebraic stacks $`pr_𝕂:_{2d,𝕂}_{2d,[1/N_\mathrm{B}]}`$ is representable and étale. Indeed, let $`S`$ be a connected scheme and suppose given a morphism $`S_{2d}`$ i.e., a polarized K3 space $`(\pi :XS,\lambda )`$ over $`S`$. Let $`\overline{b}:\mathrm{Spec}(\mathrm{\Omega })S`$ be a geometric point of $`S`$. Let $`\rho :\pi ^{\mathrm{alg}}(S,\overline{b})\mathrm{O}(P^2(\overline{b}))`$ be the monodromy representation and let $`\stackrel{~}{a}:L_{2d,\mathrm{B}}P^2(\overline{b})`$ be an isometry. Then the preimage $`\rho ^1\alpha ^{\mathrm{ad}}(𝕂_\mathrm{B})`$ is an open subgroup of $`\pi _1^{\mathrm{alg}}(S,\overline{b})`$ (of finite index) and hence it defines an étale cover $`S_{\stackrel{~}{\alpha }}`$ of $`S`$. One has that the class $`\alpha `$ of $`\stackrel{~}{\alpha }`$ in $`𝕂_\mathrm{B}\backslash \text{Isometry}(L_{2d,_\mathrm{B}},P^2(\overline{b}))`$ is $`\pi _1^{\mathrm{alg}}(S_{\stackrel{~}{\alpha }},\overline{b})`$-invariant by construction (for a fixed geometric point $`\overline{b}S_{\stackrel{~}{\alpha }}`$ over $`\overline{b}`$). Therefore we obtain a primitively polarized K3 space $`(X_{S_{\stackrel{~}{\alpha }}}S_{\stackrel{~}{\alpha }},\lambda _{S_{\stackrel{~}{\alpha }}},\alpha )`$ with a level $`𝕂`$-structure $`\alpha `$. For two markings $`\stackrel{~}{\alpha }_1`$ and $`\stackrel{~}{\alpha }_2`$ we have that $`\stackrel{~}{\alpha }_1^{\mathrm{ad}}(𝕂_\mathrm{B})=\stackrel{~}{\alpha }_2^{\mathrm{ad}}(𝕂_\mathrm{B})`$ if and only if $`\stackrel{~}{\alpha }_2^1\stackrel{~}{\alpha }_1`$ is an element of the normalizer $`N_{\mathrm{O}(V_{2d})(_\mathrm{B})}(𝕂_\mathrm{B})`$ of $`𝕂_\mathrm{B}`$ in $`\mathrm{O}(V_{2d})(_\mathrm{B})`$.
Denote by $`S^{}`$ the disjoint union of $`S_{\stackrel{~}{\alpha }}`$ where $`\stackrel{~}{\alpha }`$ runs over all (finitely many) classes in $`\mathrm{O}(V_{2d})(_\mathrm{B})/N_{\mathrm{O}(V_{2d})(_\mathrm{B})}(𝕂_\mathrm{B})`$. Let $`(X^{}S^{},\lambda _S^{},\alpha )`$ be the primitively polarized K3 space with a level $`𝕂`$-structure given by the triple $`(X_{S_{\stackrel{~}{\alpha }}}S_{\stackrel{~}{\alpha }},\lambda _{S_{\stackrel{~}{\alpha }}},\alpha )`$ on the $`\stackrel{~}{\alpha }`$-th connected component $`S_{\stackrel{~}{\alpha }}`$ of $`S^{}`$. Then by construction we have a morphism of algebraic spaces
$$\pi :S^{}S\times _{_{2d,[1/N_\mathrm{B}]}}_{2d,𝕂}$$
over $`S`$. This morphism is surjective. Indeed, by \[LMB00, Prop. 5.4\] this condition can be checked on points, in which case it is obvious by construction. The morphism $`S^{}S`$ is étale and therefore we conclude that $`pr_{𝕂,S}:S\times _{_{2d,[1/N_\mathrm{B}]}}_{2d,𝕂}S`$ and $`\pi `$ are also étale. Hence $`pr_𝕂`$ is étale.
Step 3: By Step 2 and Theorem 4.3.3 the algebraic space $`_{2d,𝕂}`$ is smooth and of relative dimension 19 over $`[1/N_\mathrm{B}]`$. ∎
###### Remark 6.1.4.
Let $`𝕂_1𝕂_2𝕂_n`$ be subgroups of finite index in $`\mathrm{SO}(V_{2d})(\widehat{})`$ and suppose that $`n3`$. Then the morphism (10) of algebraic spaces
$$pr_{(𝕂_1,𝕂_2)}:_{2d,𝕂_1,[1/N_{\mathrm{B}_1\mathrm{B}_2}]}_{2d,𝕂_2,[1/N_{\mathrm{B}_1\mathrm{B}_2}]}$$
is finite and étale. This follows from the theorem above and the relation $`pr_{𝕂_1}=pr_{𝕂_2}pr_{(𝕂_1,𝕂_2)}`$.
###### Example 6.1.5.
Let $`n3`$ be a natural number. Consider the group $`𝕂_n`$ defined in Example 5.1.3. We define $`_{2d,n}=_{2d,𝕂_n}`$ to be *the moduli space of primitively polarized K3 surfaces with level $`n`$-structure* over $`[1/2dn]`$.
###### Example 6.1.6.
Fix a natural number $`n3`$ and consider the group $`𝕂_n^\mathrm{a}`$ defined in Example 5.1.4. We define $`_{2d,n^{\mathrm{sp}}}=_{2d,𝕂_n^\mathrm{a}}`$ to be *the moduli space of polarized K3 surfaces with spin level $`n`$-structure* over $`[1/2dn]`$.
### 6.2. Moduli K3 Spaces with Full Level Structures
Suppose that $`𝕂𝕂_n`$ for some $`n3`$ is an admissible subgroup of $`\mathrm{SO}(V_{2d})(\widehat{})`$. Let $`𝒳_1=(\pi _1:X_1S_1,\lambda _1)`$ and $`𝒳_2=(\pi _2:X_2S_2)`$ be two objects of $`_{2d}`$. Suppose that $`S_1`$ and $`S_2`$ are connected and let $`(f,f_S)\mathrm{Hom}(𝒳_1,𝒳_2)`$ (in $`_{2d}`$). Let $`\overline{b}_1`$ and $`\overline{b}_2`$ be two geometric points of $`S_1`$ and $`S_2`$ such that $`f_S(\overline{b}_1)=\overline{b}_2`$. Then the morphism $`f`$ defines a homomorphism $`f_{\mathrm{et}}^{}:H^2(\overline{b}_2)H^2(\overline{b}_1)`$ sending the class of $`\lambda _{\overline{b}_2}`$ to the class of $`\lambda _{\overline{b}_1}`$. Hence we obtain a map
$$f^{}:𝕂_\mathrm{B}\backslash \left\{g\text{Isometry}(L_{0,_\mathrm{B}},H^2(\overline{b}_2))\right|g(e_1df_1)=c_1(\lambda _{2,\overline{b}_2})\}$$
$$𝕂_\mathrm{B}\backslash \left\{g\text{Isometry}(L_{0,_\mathrm{B}},P^2(\overline{b}_1))\right|g(e_1df_1)=c_1(\lambda _{1,\overline{b}_1})\}$$
given by $`\alpha f_{et}^{}\alpha `$ and commuting with the monodromy actions on both sides.
###### Definition 6.2.1.
For a natural number $`d`$ and an admissible subgroup $`𝕂`$ of $`\mathrm{SO}(V_{2d})(\widehat{})`$ as above consider the category $`_{2d,𝕂}^{\mathrm{full}}`$ defined in the following way:
1. Triples $`(\pi :XS,\lambda ,\alpha )`$ of a K3 space $`\pi :XS`$ over $`S`$ with a primitive polarization $`\lambda `$ of degree $`2d`$ and with a full level $`𝕂`$-structure $`\alpha `$ on $`(\pi :XS,\lambda )`$.
2. Suppose given two triples $`𝒳_1=(\pi _1:X_1S_1,\lambda _1,\alpha _1)`$ and $`𝒳_2=(\pi _2:X_2S_2,\lambda _2,\alpha _2)`$. Let $`f_S:S_1S_2`$ be a morphism of schemes. Choose base geometric points $`\overline{b}_1^{}`$ and $`\overline{b}_2^{}`$ on any two connected components $`S_1^{}`$ and $`S_2^{}`$ of $`S_1`$ and $`S_2`$ for which $`f:S_1^{}S_2^{}`$ such that $`f_S(\overline{b}_1^{})=\overline{b}_2^{}`$. Define the morphisms between $`𝒳_1`$ and $`𝒳_2`$ in the following way
$$\begin{array}{cc}\hfill \mathrm{Hom}(𝒳_1,𝒳_2)=\{\mathrm{pairs}(f_S,f)|& f_S:S_1S_2\mathrm{is}\mathrm{a}\mathrm{morph}.\mathrm{of}\mathrm{spaces},\hfill \\ & f:X_1X_2\times _{S_2,f_S}S_1\mathrm{is}\mathrm{an}\mathrm{isom}.\mathrm{of}\hfill \\ & S_1\mathrm{spaces}\mathrm{with}f^{}\lambda _2=\lambda _1\mathrm{and}\hfill \\ & f^{}(\alpha _1)=\alpha _2\mathrm{on}\mathrm{any}\mathrm{conn}.\mathrm{cmpt}.\mathrm{of}S_1\}\hfill \end{array}$$
A full level $`𝕂`$-structure $`\alpha `$ on a primitively polarized K3 space $`(XS,\lambda )`$ defines in a natural way a level $`𝕂`$-structure via the injective morphism
$$i__\mathrm{B}^{}:𝕂_\mathrm{B}\backslash \left\{g\text{Isometry}(L_{0,_\mathrm{B}},H^2(\overline{b}))\right|g(e_1df_1)=c_1(\lambda _{\overline{b}})\}$$
$$𝕂_\mathrm{B}\backslash \text{Isometry}(L_{2d,_\mathrm{B}},P^2(\overline{b}))$$
commuting with the monodromy action. This morphism is defined by the embedding of lattices $`i:L_{2d}L_0`$ (see (1) in Section 2.1). Using this, just like in the case of moduli of primitively polarized K3 surfaces with a level structure, we define natural functors.
1. Define a functor
$$i_𝕂:_{2d,𝕂}^{\mathrm{full}}_{2d,𝕂}$$
sending $`(XS,\lambda ,\alpha )`$ to $`(X,S,\lambda ,i^{}(\alpha ))`$ which makes $`_{2d,𝕂}^{\mathrm{full}}`$ into a full subcategory of $`_{2d,𝕂}`$ over $`(\mathrm{Sch}/[1/N_\mathrm{B}])`$.
2. One has the forgetful functor
(11)
$$pr_𝕂:_{2d,𝕂}^{\mathrm{full}}_{2d,[1/N_\mathrm{B}]}$$
sending a triple $`(\pi :XS,\lambda ,\alpha )`$ to $`(\pi :XS,\lambda )`$ and an element $`(f,f_S)\mathrm{Hom}(𝒳,𝒴)`$ of $`_{2d,𝕂}^{\mathrm{full}}`$ to $`(f,f_S)`$.
3. For any two admissible subgroups $`𝕂_1𝕂_2`$ of $`\mathrm{SO}(V_{2d})(\widehat{})`$, contained in some $`𝕂_n`$ for $`n3`$, one has a projection functor
(12)
$$pr_{(𝕂_1,𝕂_2)}:_{2d,𝕂_1,[1/N_{\mathrm{B}_1\mathrm{B}_2}]}^{\mathrm{full}}_{2d,𝕂_2,[1/N_{\mathrm{B}_1\mathrm{B}_2}]}^{\mathrm{full}}$$
defined in a similar way as the corresponding morphism (10) in 3.
The functors $`pr_𝕂`$ and $`pr_{(𝕂_1,𝕂_2)}`$ defined above are the restrictions of the corresponding functors (9) and (10) to the category of primitively polarized K3 surfaces with full level $`𝕂_j`$-structures via $`i_{𝕂_j}`$ for $`j=1,2`$.
###### Theorem 6.2.2.
Let $`𝕂`$ be an admissible subgroup of $`\mathrm{SO}(V_{2d})(\widehat{})`$ contained in $`𝕂_n`$ for some $`n3`$. The category $`_{2d,𝕂}^{\mathrm{full}}`$ is a separated, smooth algebraic space of relative dimension 19 over $`[1/N_\mathrm{B}]`$. The morphism $`p_{2d,𝕂}:_{2d,𝕂}^{\mathrm{full}}_{2d,[1/N_\mathrm{B}]}`$ is étale and the morphism $`i_𝕂:_{2d,𝕂}^{\mathrm{full}}_{2d,𝕂}`$ is an open immersion.
###### Proof.
To prove that $`_{2d,𝕂}^{\mathrm{full}}`$ is representable by an algebraic space of finite type over $`[1/N_\mathrm{B}]`$ one follows the steps of the proof of Theorem 6.1.2. In this way we also see that the projection morphism $`p_{2d,𝕂}:_{2d,𝕂}^{\mathrm{full}}_{2d,[1/N_\mathrm{B}]}`$ is finite and étale. Therefore we have a commutative diagram
where the two morphisms $`p_{2d,𝕂}`$ are étale and surjective. Hence $`i_𝕂`$ is also étale and therefore it is open. ∎
###### Remark 6.2.3.
Let $`𝕂_1𝕂_2`$ be two admissible subgroups of $`\mathrm{SO}(V_{2d})(\widehat{})`$. Then the morphism of algebraic spaces
$$pr_{(𝕂_1,𝕂_2)}:_{2d,𝕂_1,[1/N_{\mathrm{B}_1\mathrm{B}_2}]}^{\mathrm{full}}_{2d,𝕂_2,[1/N_{\mathrm{B}_1\mathrm{B}_2}]}^{\mathrm{full}}$$
is finite and étale. This follows from the theorem above and the relation $`pr_{𝕂_1}=pr_{𝕂_2}pr_{(𝕂_1,𝕂_2)}`$.
###### Example 6.2.4.
Let $`n3`$ be a natural number. Consider the group $`𝕂_n^{\mathrm{full}}`$ defined in Example 5.3.5. We define $`_{2d,n}^{\mathrm{full}}=_{2d,𝕂_n^{\mathrm{full}}}^{\mathrm{full}}`$ to be *the moduli space of primitively polarized K3 surfaces with full level $`n`$-structure* over $`[1/2dn]`$. |
warning/0506/q-bio0506038.html | ar5iv | text | # Some protein interaction data do not exhibit power law statistics
## 1 Introduction
Experimental data on protein-protein interaction (PPI) networks have been extensively gathered with the aim of acquiring a system-level understanding of biological processes . Various statistical features of complex graphical structures have received attention, including the size of the largest connected component, the node degree distribution, the graph diameter, the characteristic path length, and the clustering coefficient. However, the feature that has attracted the most attention is the distribution of node degree (the number of links from a node) and whether or not the distribution follows a power law (linear plot on log-log scale). The degree distribution of PPI networks was claimed to follow a power law in , and thus PPI networks are considered to be ”scale-free” (SF) , a generic property of network topologies common to various networks in different domains, from social networks and biological systems to the Internet.
Although “scale-free” has not been clearly defined in the existing literature , most treatments assume that a power law node degree distribution is an important, and sometimes defining feature. Other characteristics described in the SF literature include failure tolerance but attack vulnerability at hubs (nodes possessing high degree) and various kinds of self-similarity. A recent attempt at a more theoretically rigorous treatment shows however that no additional features follow from power law node degree sequence alone, and require additional restrictions, such as high likelihood of occurrence by random generation (e.g. by preferential attachment). Other work has highlighted important differences between PPI networks and ”SF networks” constructed by a stochastic growth model. Moreover, one may question the rigor with which the power law node degree distribution, the primary feature of SF networks, has been demonstrated in certain examples.
This letter shows that the node degree sequences of some published PPI networks are better described by an exponential function when properly plotted and analyzed. The problem with previous work is that data were plotted using frequency-degree plots, as is common in papers purporting to discover power laws in complex biological systems, which lead to systematic errors compared with rank-degree plots. We demonstrate here that data plotted on a loglog scale frequency-degree plot may appear to be linear, but when the same data are plotted on a loglog scale rank-degree plot, they are clearly shown not to be power law. Thus, the data for some PPI networks lack even the minimal features of scale-free networks.
## 2 Materials and Methods
Publicly available data for PPI networks represent only an approximation of the real interaction network because of the large number of false positive and false negative interactions. However, because of the assumed self-similarity features of SF networks, it has been claimed that if the real PPI network is SF, then any appropriately sampled subnetwork is also SF . Thus, we might still gain valuable information by examining whether the publicly available PPI network data possess a power law node degree distribution characteristic of SF networks.
A finite sequence of node degrees $`y=(y_1,y_2,\mathrm{},y_n`$) of integers, assumed without loss of generality always to be ordered such that $`y_1y_2\mathrm{}y_n`$, is said to follow a power law if
$$kcy_{k}^{}{}_{}{}^{\alpha },$$
(1)
where $`k`$ is (by definition) the rank of $`y_k`$, $`c>0`$ is a constant, and $`\alpha >0`$ is called the scaling index. Because of the ordering, the rank $`k`$ is the number of nodes with the degree equal or larger than $`y_k`$. Since $`\mathrm{log}k=\mathrm{log}c\alpha \mathrm{log}y_k`$, the rank $`k`$ versus the node degree $`y_k`$ plot on a loglog scale appears as a straight line of slope $`\alpha `$. In contrast, $`y`$ is said to follow an exponential if
$$ka\mathrm{exp}^{by_k},$$
(2)
where $`a>0`$ and $`b>0`$ are constants. The $`k`$ versus $`y_k`$ plot on a semilog scale approximates a straight line of slope of $`b`$ since $`\mathrm{log}k=\mathrm{log}aby_k`$.
Note that the rank-degree relationships (1) and (2) are non-stochastic, in the sense that there need be no assumption of an underlying probability distribution for the sequence $`y`$. Indeed, no coherent justification has been given for why biological networks should be viewed as samples from a random ensemble. On the contrary, what is known of evolution would suggest that it yields extremely nonrandom structure at every level of organization. Nevertheless, random graphs have been remarkably popular models for biological networks, but have led to substantial confusion, particularly with regard to power laws. Suppose a non-negative random variable $`X`$ has cumulative distribution function (CDF) $`F(x)=P[Xx]`$. In this stochastic context, a random variable $`X`$ or its corresponding distribution function $`F`$ is said to follow a power law with index $`\alpha >0`$ if, as $`x\mathrm{}`$,
$`P[X>x]=1F(x)cx^\alpha ,`$ (3)
for some constant $`c>0`$ and a tail index $`\alpha >0`$, where $`f(x)g(x)`$ as $`x\mathrm{}`$ if $`f(x)/g(x)1`$ as $`x\mathrm{}`$. We call (3) the stochastic form of power law rank-degree relationship. The loglog plot of $`P[X>x]`$ versus $`x`$ appears as a straight line of slope $`\alpha `$ for large $`x`$. If the CDF $`F(x)`$ satisfying (3) is differentiable, then its derivative, the probability density function $`f(x)=\frac{d}{dx}F(x)`$, satisfies
$`f(x)c^{}x^{(1+\alpha )}.`$ (4)
The loglog plot of $`f(x)`$ versus $`x`$ also would be a line of slope $`(1+\alpha )`$. In contrast to the rank-degree relationships (1) and (2), the definitions in (3) and (4) are stochastic and require an underlying probability model. As is standard in physics, the SF literature almost exclusively assumes some underlying stochastic models, and power law node degree distributions are typically investigated in terms of the frequency-degree relationship based on the probability density function $`f(x)`$.
In the case of node degree of graphs the data is inherently discrete. Even if the data were sampled from some ensemble, $`F(x)`$ is not differentiable and the frequency-degree plots simply do not make sense and can easily lead to mistakes. Furthermore, differentiation of noisy data, such as PPI data, amplifies errors, making frequency-based data uninformative and ambiguous. A typical approach to overcome these problems is to smooth the data or to group individual data values into a small number of bins, and then plot the relative number of data values in each bin. The problem is that this smoothing or binning process can dramatically change the nature of frequency-based statistics as will be shown below (Figs.1 and 2). This use of ad hoc statistical analysis can lead to concluding incorrectly that a power law relationship is present (or absent). This problem is easily avoided if one were to make rank-degree plots of raw data instead of using frequency-degree plots to check the power law or exponential relationships in (1) and (2).
From among many publicly available studies on PPI networks, we used the filtered yeast interactome (FYI) data set and the predicted human protein-interaction (HPI) map to illustrate these points. Much of the original data suffers from numerous false positives and false negatives, but more recent investigations have sought to refine the data. For example, the FYI data set contains high-confidence interactions for yeast, each observed by at least two different methods, thereby enriching for genuine positives. The HPI map was generated using data from seven experimental and four computationally predicted protein-interaction maps from Saccharomyces cerevisiae , Drosophila melanogaster and Caenorhabditis elegans . The idea is that a human protein interaction can be predicted if orthologs in a model organism show an interaction. Its accuracy has been assessed in . We consider both FYI and HPI to be refined data sets, and investigate whether their node degree sequences follow a power law, a defining feature of scale-free networks, by rank-degree plots.
## 3 Results and Discussion
The rank-degree plots of the HPI and FYI data are shown in (a) loglog scale and (b) semilog scale in Figs. 1 and 2, respectively. The straight lines and the dotted curve in loglog scale (a) show least-squares fitting of data to a power law with the value of its slope and to an exponential, respectively. The same fittings are depicted as the curve and the dotted straight line in semilog scale (b). From these figures, we can clearly conclude that the node degree sequences of HPI and FYI data are much closer to an exponential (2), and are clearly not power laws (1). More sophisticated statistical analysis can be used to confirm these conclusions. In addition, the rank-degree plots show raw data and readers can easily judge at a glance the relative suitability of various models.
However, using frequency-degree plots (c) in Figs.1 and 2 could lead to the erroneous conclusion that the node degree sequence appears to follow a power law, although the correct rank-degree plot clearly shows that this is not the case. Furthermore, even if the PPI data were a power law, the slope for frequency-degree plot $`\beta `$ is simply not related to the slope for the rank-degree plot $`\alpha `$ by $`\beta =\alpha +1`$, as holds for differentiable distributions. These results conclusively demonstrate that these two refined PPI data sets are not power laws, and thus certainly not scale-free, no matter how this is defined.
It is in principle possible that the data studied here is misleading and real PPI networks might have some features attributed to scale-free networks. At this time we only can draw conclusions about (noisy) subgraphs of the true network since the data sets are incomplete and presumably contain errors. However, the fact that these subgraphs exhibit an exponential node degree sequences suggests that the entire network is not SF. Appropriately sampled subraphs of a SF graph should be SF, and hence possess a power law node degree sequence. Furthermore, a SF network possessing significant non-SF subnetworks could not be considered to be self-similar, a typically assumed though as yet unproven feature of scale-free networks. Finally, since essentially all claims that biological networks are scale-free are based on error-prone frequency-degree analysis, this analysis must be completely redone to determine the correct form of the degree sequences.
It has also been shown that the Internet and cell metabolism, the two most prominent examples of SF networks, might have power laws for some degree sequences, but have none of the other features attributed to scale-free networks. One important feature of the Internet and metabolic networks is the complete absence of centrally located high-degree hubs which are responsible for global network connectivity and whose removal would fragment the network, in contrast to what has been claimed in the SF literature. Metabolic networks have also been shown to be scale-rich (SR), but not SF, in the sense that they are far from self-similar despite some power laws in certain node degree sequence. Their power law node degree sequence is a result of the mixture of exponential distributions in each functional module. In principle, PPI networks could have this SR structure as well, and perhaps power laws could emerge at higher levels of organization. This will be revealed only when a more complete network is elucidated. Still, the most important point is not whether the node degree sequence follows a power law, but whether the variability of the node degree sequences is high or low , and the biological protocols that necessitate this high or low variability. These issues will be explored in future publications.
The authors thank Nicolas Bertin and Marc Vidal for providing FYI data. |
warning/0506/cond-mat0506468.html | ar5iv | text | # Universality Classes of Metal-Insulator Transitions in Strongly Correlated Electron Systems and Mechanism of High-Temperature Superconductivity
## I Introduction
Mott transition belongs to one of the metal-insulator transitions ubiquitous in various compounds RMP . Physical properties of Mott transition and its nature are a long-standing subject of research with many controversial issues. The problem has first been postulated in 1937 by Peierls and Mott PeierlsMott . Although the band theory of metals and insulators established soon after the foundation of quantum mechanics is quite successful, it was pointed out by de Boer and Verway deBoer that insulating behaviors of many transition metal compounds as NiO cannot be explained by a simple band picture, because the bands are partially filled. Peierls pointed out a crucial role of electron correlations as the mechanism of the insulating behavior. Mott developed this idea and introduced a concept which we nowadays call the Mott insulator Mott ; Mott2 . Since then, it has been recognized for a long time that the Mott insulator stabilized by the electron-electron Coulomb interaction and metallic states near it provide us with fruitful physics with various novel concepts. This has become very popular after the discovery of the high-$`T_c`$ cuprate superconductors, which was indeed discovered in doped Mott insulators Bednorz . However, relationship of this fruitful outcome to the nature of the Mott transition itself has not been fully clarified.
The Mott transition can be realized basically by two routes: Bandwidth-control and filling-control. In the first route, the bandwidth is controlled relative to the amplitude of the local electron-electron interaction by keeping the electron density fixed at a commensurate value (namely, the electron density per unit cell is kept at an integer value). This route may be experimentally realized by applying pressure or by substituting elements with a different ionic radius and the same valence. In this route, the overlap of the wavefunctions between neighboring electronic atomic orbitals forming conduction electron bands is controlled. In the second route, electron filling is changed from the Mott insulator. This is typically realized by substituting with elements, which have a different valence from the substituted elements in the reservoir structure of the Mott insulator.
Originally, Mott has considered only the first route, the bandwidth-control transition. In this category, there exist several typical Mott transitions including those observed in V<sub>2</sub>O<sub>3</sub>, $`R`$NiO<sub>3</sub> and $`\kappa `$-ET type organic compounds, where $`R`$ represents a rare earth element. Many of them show first-order transitions between the Mott insulator and metals at low temperatures RMP . In many cases, some magnetic and/or orbital order exist at low temperature of the Mott insulator phase. Although sometimes concurrent magnetic transitions occur with the Mott transition at low temperatures, in the typical phase diagram illustrated in Fig. 1, the first-order Mott transitions occur even at high temperatures, where the magnetic and orbital orders are not involved. This typical phase diagram clearly indicates that the Mott transition is inherently independent of the magnetic and orbital-ordering transitions. The statement that the mechanism of the Mott transition is primarily independent of the symmetry breaking of spins or orbitals is also corroborated by the existence of a quantum spin liquid phase in the Mott insulator recently found numerically Kashima1 ; Morita as well as experimentally KanodaSpinliquid ; Fukuyama ; Ishimoto . The Mott transition between a quantum spin-liquid and a metal does not accompany magnetic transitions. In fact, as we see in this paper, the Mott transition is not driven by the spin or orbital degrees of freedom but by the density degrees of freedom.
Mott argued that the long-range part of the Coulomb interaction is necessary to reproduce the first-order transition Mott1storder . The argument by Mott was the following: We can identify the Mott insulating phase as that where the two electrons sitting on the same atomic orbitals (we call it a doublon) and an empty site (we call it a holon) make a bound state. Metals are characterized by the phase where the bound state disappears. In terms of the binding of the doublon and holon, the binding energy is controlled by the screening of the attractive Coulomb interaction between the doublon and holon. Since the screening relies on other doublons and holons moving in between, the screening becomes rapidly poor with decreasing density of free doublons and holons. Then at some threshold concentration of free doublons and holons, the screening becomes too weak to keep a doublon and a holon free and as a consequence of this feedback, they suddenly form a bound state, which leads to a first-order transition to a Mott insulating state. Another argument for the origin of the first-order transition emphasizes couplings to lattice distortion. When the metal is stabilized, the lattice constant diminishes further to gain the kinetic energy of electrons, which in general strongly favors the first-order transition through the electron-lattice coupling.
However, recent detailed numerical studies on the Hubbard model on the square lattice overturned these speculations Kashima1 ; Morita ; WatanabeGPIRG : Although the long-range Coulomb force and the coupling to the lattice may have some roles in the realistic Mott transition, even the Hubbard model with only the local onsite interaction without any coupling to lattice distortions is enough to reproduce the first-order Mott transition. The $`N`$-site Hubbard model is defined as
$``$ $`=`$ $`_t+{\displaystyle \underset{i}{}}H_{Ui}\mu MN`$ (1)
$`_t`$ $`=`$ $`{\displaystyle \underset{ij}{}}t_{ij}(c_{i\sigma }^{}c_{j\sigma }+h.c.)`$ (2)
and
$$_{Ui}=U(n_i\frac{1}{2})(n_i\frac{1}{2}).$$
(3)
Here, $`M_{i\sigma }n_{i\sigma }/N`$ and $`n_{i\sigma }=c_{i\sigma }^{}c_{i\sigma }`$ with the creation (annihilation) operator $`c_{i\sigma }^{}(c_{i\sigma })`$ of an electron at the site $`i`$ with the spin $`\sigma `$. The chemical potential is $`\mu `$ and $`U`$ is the onsite Coulomb repulsion. The phase diagram of the Hubbard model at zero temperature $`T=0`$ on the square lattice with the nearest-neighbor and next-nearest-neighbor transfers $`t`$ and $`t^{}`$, respectively, studied by the path-integral renormalization group method ImadaKashima ; KashimaImadaPIRG is shown in the plane of $`U/t`$ and the chemical potential $`\mu `$ at zero temperature in Fig. 2 WatanabeGPIRG . The filling-control transition occurs across the edge of the boundary in Fig. 2, while the bandwidth-control transition is realized through the corner of the phase boundary at the bottom. The first-order transition through the bandwidth-control route is indicated by the jump of the averaged doublon density $`Dn_in_i=_in_in_i/N`$ as in Fig. 3 for the case of the anisotropic triangular lattice Morita . When the first-order transition takes place at zero temperature, its jump decreases with raising temperatures and closes at the critical endpoint. In fact, the Hubbard model on this anisotropic triangular lattice is a relevant effective model for the $`\kappa `$-ET type organic compound and the first-order Mott transition with the finite-temperature critical end point was observed experimentally Fournier ; Kanoda ; Lefebvre .
Another remarkable feature in Fig. 3 is that the jump in $`D`$ decreases with increasing the next nearest neighbor transfer $`t^{}`$. This means that the so-called frustration effects reduce the first-order jump and drives to more continuous type transitions. In fact, recent experimental studies on pyrochlore compounds Tokura appear to show a continuous Mott transition by the bandwidth control with a signature of the Anderson localization in the vicinity of the boundary. We note that the pyrochlore lattice has the fully frustrated structure, where the antiferromagnetic order is severely suppressed. In addition, recent studies on $`\kappa `$-(ET)<sub>2</sub>Cu<sub>2</sub>(CN)<sub>3</sub> suggest that the critical temperature of the Mott transition becomes lower for the more frustrated structure, namely for compounds with larger effective $`t^{}`$ in the corresponding theoretical model. Kanoda3 All of these consistently suggest that the order of the Mott transition may be systematically controlled from the first order to continuous.
When the filling is controlled, the first-order transition appears as the phase separation. Experimentally, the existence of the phase separation or the electronic inhomogeneity is a controversial issue as we discuss later. From theoretical side, numerical studies on the Hubbard model show the marginal result, where the phase separation does not occur while the charge susceptibility shows a critical divergence at the Mott transition. FurukawaImada0 ; FurukawaImada ; FurukawaImada2 In the terminology in this paper, we use the charge susceptibility and the density susceptibility as the same quantity $`dM/d\mu `$ for the electrons with charge.
Since the theoretical and experimental results suggest the controllability of the order of the Mott transition and its critical temperature, it is desired to understand the whole feature of the Mott transition from the universality classes of the finite-temperature critical point to the zero-temperature critical phenomena on the same grounds. We will show in this paper that three regimes of the Mott transition exist. One is the classical transition at a high temperature, which is described by the Ising universality class of the critical point accompanied by the first-order transition below the critical temperature. The second is the quantum transition, where the transition appears only at zero temperature, and the density susceptibility remains finite for the spatial dimension $`d2`$. The third regime is the marginally quantum one, which emerges at the crossing point of the classical and quantum transitions. The marginally quantum regime is characterized by the diverging density susceptibility at small wavenumber for $`d2`$ at low temperatures. In the second and third regimes, the conventional scheme of the Ginzburg-Landau-Wilson theory does not apply.
The new universality class at the marginally quantum transition has a deep consequence on the induced non-Fermi-liquid behavior in the metallic side, electron differentiation in the momentum space and the sensitivity toward electronic inhomogeneity. One of the most remarkable consequences is the superconductivity emerging from this marginal quantum Mott criticality. We show that the high temperature superconductivity of the $`d_{x^2y^2}`$ symmetry is obtained under the realistic choice of the parameter values for the cuprate superconductors, where the density (charge) fluctuations at small wavenumber play the crucial role for the Cooper pairing. The energy scale of the fluctuation is characterized by the Mott gap, which can be by far larger than the energy scale of magnetic and orbital fluctuations for the filling-control transition. This solves many puzzling experimental results in strongly correlated electron systems particularly in transition metal compounds.
A part of the discussions in this article is already given Imada2004 ; Imada2005 . We summarize the previous results and further extend the discussion on the quantum Mott criticality and its consequences in greater detail. In particular, detailed analyses on the breakdown of the Ginzburg-Landau-Wilson scheme for the quantum Mott transition are presented together with the scaling analysis. The validity of the mean-field exponents and the compatibility with the hyperscaling description are discussed in detail. Two-dimensional systems are especially analyzed and are compared with the experimental results for the organic compounds and the cuprate superconductors. Basic finite temperature effects are also obtained and discussed in connection with the experimental results. An important issue for the filling-control transition is the effects of the long-range Coulomb interaction. We discuss how the present results are modified in the presence of the long-range interaction and also discuss the experimental relevance. The non-Fermi-liquid properties and the anisotropic Cooper pairing originating from the nonperturbative enhancement of the density susceptibility at a small wavenumber is an important subject we study in this article in detail.
In Sec. II, we summarize the conventional Ginzburg-Landau-Wilson scheme of the Mott transition. In Sec. III, the nature of the quantum Mott transition is examined with emphasis on the breakdown of the Ginzburg-Landau-Wilson scheme. Section IV is devoted to the resultant non-Fermi-liquid behavior expected in the metallic side of the critical region of the Mott transition. In Sec. V we discuss the mechanism of high-temperature superconductivity arising from the marginal quantum Mott criticality. Section VI concludes and summarizes the paper.
## II Conventional Ginzburg-Landau-Wilson Scheme
Recently, the critical endpoint of the first-order transition line of the Mott transition in Fig. 1 has been a subject of intensive studies. In case of V<sub>2</sub>O<sub>3</sub>, from the detailed study of the conductance, it has been suggested that the criticality of the transition follows the Ising-type universality class Limelette . In an organic conductor of $`\kappa `$-ET salt, the diverging electronic compressibility at the critical end point has been probed by the ultrasound velocity Fournier .
In prior to these experimental studies, several theoretical studies have focused on the nature of the transition. The Mott transition by itself does not change any symmetry. Therefore, from theoretical point of view, it has an analogy with the text-book gas-liquid transition, which is known to be equivalent to the ferromagnetic transition in the Ising model under magnetic fields. The first-order metal-insulator transition corresponds to the Ising transition between spin-up and down phases taking place with switching the direction of magnetic fields below the critical temperature. In fact, Castellani Castellani has discussed the Ising nature of the first-order Mott transition by extending the Blume-Emery-Griffiths model BEG for the phase separation of <sup>3</sup>He-<sup>4</sup>He mixture. It has further been considered in the dynamical mean-field theory by Kotliar et al. Kotliar ; Kotliar2 , where the low-energy part of the single-particle Green’s function appears to follow Ginzburg-Landau scheme in accordance with the mean-field theory of the Ising model.
In the Ginzburg-Landau-Wilson scheme for the Ising-type transition, the free energy may be expanded by the spatially dependent scalar order parameter $`X(r)`$ integrated over space coordinate $`r`$ as
$$F=𝑑r[\frac{1}{2}a_0(TT_c)X(r)^2+\frac{1}{4}bX(r)^4\mu X(r)]$$
(4)
near the critical temperature $`T_c`$ with $`a_0`$ and $`b`$ being positive constants. In the Ising model, $`X`$ is indeed the order parameter, namely, the magnetization $`m`$. In the mapping to the gas-liquid transition, $`X`$ is interpreted as the density of particles, $`n`$, measured from the critical density. At the critical temperature $`T_c`$ of the gas-liquid transition, the uniform density susceptibility $`\chi _n=[d^2F/dn^2]^1`$ diverges.
When it is further extended to the mapping to the Mott transition, for the filling-control transition, $`X`$ is identified indeed as the electron doping concentration $`X`$ measured from the critical density at the critical point of the Mott transition Imada2004 . This is a natural consequence, because, in the filling-control transition, the control parameter is the chemical potential, which is conjugate to the carrier density.
When the bandwidth is controlled, the control parameter is $`U/t`$ in the Hubbard model and the conjugate quantity to $`U`$ is the doulon density $`D`$. Therefore, the order parameter $`X`$ in this case is the doublon density $`D`$, which indeed jumps at the first-order transition as in Fig. 3. We note that in this case the holon density should be the same as the doublon density, because the holon and doublon densities must be the same to keep the density for the route of the bandwidth-control transition.
To describe the both types of the transitions, we take the natural order parameter as $`\zeta `$, where $`\zeta =X`$ for the filling-control transition and $`\zeta =D`$ for the bandwidth-control transition. Namely, we take
$$F=𝑑r[\frac{1}{2}a_0(TT_c)\zeta (r)^2+\frac{1}{4}b\zeta (r)^4\mu _\zeta \zeta (r)].$$
(5)
Here, $`\mu _\zeta `$ is $`U`$ for the bandwidth-control transition whereas is the chemical potential $`\mu `$ conjugate to the doping concentration $`X`$ for the filling-control transition. The Ising universality is resulted from this Ginzburg-Landau-Wilson functional Goldenfeld . When the Ising universality is correct, we obtain the critical susceptibility exponent defined by $`\chi _c[^2F/\zeta ^2]^1(TT_c)^\gamma `$ with $`\gamma _=7/4`$ for two-dimensional systems, $`d=2`$ and $`\gamma 1.24`$ for $`d=3`$. The order parameter exponent defined below $`T_c`$ as $`\zeta |TT_c|^\beta `$ obtained from $`F/\zeta =0`$ at $`\mu _\zeta =0`$ satisfies $`\beta =1/8`$ and $`\beta 0.325`$ for $`d=2`$ and 3, respectively. The exponent with varying $`\mu _\zeta `$ defined as $`\zeta \mu _\zeta ^{1/\delta }`$ at $`T=T_c`$ is given by $`\delta =15`$ and 4.8 for $`d=2`$ and 3, respectively. These exponents for $`d=3`$ were indeed claimed to be observed at the critical point of the bandwidth-control Mott transition for V<sub>2</sub>O<sub>3</sub> Limelette .
In the Ising-transition picture Goldenfeld , the transition is characterized by these simple exponents with the hyperscaling assumption being satisfied below the upper critical dimension $`d_u=4`$. The system has a single length scale $`\xi `$ which diverges at $`T=T_c`$. In the present context, $`\xi `$ expresses the density correlation length or doublon density correlation length.
## III Quantum Mott Criticality
### III.1 General remark on quantum effect
A nontrivial question arises when the critical temperature of the Mott transition, $`T_c`$ can be lowered. With lowering of $`T_c`$, how do the quantum effects emerge? If $`T_c`$ becomes zero, then one might naively expect that conventional quantum critical phenomena would appear. A naive expectation would be that the transition might be described by the Ising universality class in $`d+1`$ dimensions, where the additional one dimension emerges from the dimension in the imaginary time in the path integral formalism. It turns out later that this is not the case. In any case, the criteria for the existence of the non-negligible quantum effect should be determined from the existence of the Fermi degeneracy of the electrons. When the Fermi degeneracy temperature becomes comparable or higher than $`T_c`$, the Mott transition has to be treated fully quantum mechanically.
One may argue that even when the bare Fermi temperature is high, the effective Fermi temperature would be suppressed near the Mott critical point, because, at the continuous transition point to the insulator, the Fermi degeneracy temperature should be zero. However, the Fermi degeneracy may still coexist with the critical fluctuation in the metallic side near $`T_c`$. In fact, even when $`T_c`$ is zero, the Fermi degeneracy always appears at temperatures sufficiently close to zero, if the parameter infinitesimally deviates from the critical point. Then the quantum effect should become relevant when one approaches the transition point by keeping the temperature sufficiently low.
In the quantum region, we have to consider quantum dynamics. This can be done by considering the path-integral formalism, where the imaginary time direction must be additionally considered in addition to the real spatial dimension. The time scale $`\omega ^1`$ diverges as $`\omega ^1\xi ^z`$ in addition to the divergence of the spatial correlation length $`\xi `$. The quantum dynamics is characterized by the dynamical exponent $`z`$. We note here that the Mott transition can be characterized by two different dynamical exponents in principle. This is because at the Mott transition, the single-particle spectra given by the single-particle Green’s function and the two-particle correlations represented by the density (charge) correlation functions both have singular behaviors with diverging time scale, while these two may in principle follow different scalings. Therefore we can define two dynamical exponents $`z`$ and $`z_t`$ for single-particle and two-particle spectra, respectively. We will show below that these two coincide each other.
Here, we discuss how the quantum effect alters the transition by assuming the region where $`T_c`$ is low or even zero. When $`T_c`$ becomes low, the Ginzburg-Landau expansion tells that the charge fluctuation becomes diverging accompanied by the quantum degeneracy, which becomes beyond the scope of the form (5). One might expect that the quantum region could be described by the Ising universality in $`d+z`$ dimensions. We will show that this does not apply and the quantum effect is more profound.
### III.2 What is different from the conventional quantum critical phenomena?
In contrast to phase transitions with simple spontaneous symmetry breakings, the transitions from metals to the band insulators and the Mott insulators have no spontaneous symmetry breaking by themselves. Therefore as in the gas-liquid transition, the metals and insulators have no clear distinction at nonzero temperatures, if the first-order transition would be absent. However, at zero temperature, insulators are always clearly distinguished from metals by the vanishing conductivity. Among various types of insulators, band insulators and Mott insulators both have clear distinction from metals by vanishing Drude weight and vanishing charge susceptibility (compressibility) at zero temperature RMP ; Imada1995 . The Drude weight is the stiffness to the twist of the phase of the wavefunction in the spatial direction, while the charge susceptibility is the stiffness to the twist in the temporal direction in the path-integral formalism Kohn ; Fisher ; RMP . Both of the two quantities have nonzero values only in metals. Therefore, metals under the perfectly periodic potential of ions are regarded as a state where symmetry of the phase of the spatially-extended electron wavefunction is broken. Then, at zero temperature, insulators cannot be adiabatically continued to a metal. Two phases have to be clearly separated by a phase boundary.
From these facts, we notice that the Mott transition at zero temperature may have a quite different universality class. In fact, if one could lower $`T_c`$ in Eq. (5) by controlling some microscopic parameter, one may also expect that it could pass through zero and even to a negative temperature. This implies that the transition would become quantum critical and then the Mott transition would disappear as in the conventional scenario of the emergence of quantum critical phenomena as schematically illustrated in Fig. 4(a). However, we have seen above that this cannot happen because of the clear distinction between metals and insulators at $`T=0`$. This by itself indicates that the Ginzburg-Landau-Wilson scheme has to break down when $`T_c`$ becomes zero. When $`T_c`$ is lowered to zero, and if one tries to drive the control parameter further to lower $`T_c`$, it in reality keeps $`T_c`$ at zero, namely the quantum transition at $`T=0`$ continues as we see in Fig. 4(b). We call this continuation line, the $`T_c=0`$ boundary. When $`T_c`$ becomes just zero from nonzero values, we call this point the marginal quantum critical point, which is indicated by the solid circle in Fig. 4(b).
We have discussed already that the first-order transitions with the critical end point are indeed found in experiments. We have also discussed and will discuss later that the metal-insulator transition through the $`T_c=0`$ boundary appears to exist. Then the marginal quantum critical region is the crossing point of these two regions. We will clarify that many strongly correlated systems including the high-$`T_c`$ cuprates may be located in this marginal quantum critical region. Numerical results of the filling-control transition in the two-dimensional Hubbard model indeed suggest a continuous transition at zero temperature with the diverging charge susceptibility (density susceptibility) FurukawaImada0 ; FurukawaImada ; FurukawaImada2 , which is consistent with what is expected in the marginal quantum critical region. Therefore, this quantum criticality is certainly a realistic possibility. From experimental point of view, the quantum parameter $`g`$ may be controlled by the lattice structure, particularly, by the geometrical frustration as already discussed in §I. The control of the frustration parameter was actually achieved by the choice of anions in $`\kappa `$-type ET compounds while uniaxial pressure may also be used to control the frustration effects in general.
### III.3 Single-particle quantum dynamics
The metallic phase except for one-dimensional systems has the adiabatic continuity with the Fermi liquid. Therefore, low-energy part of the free energy can be described by the fermionic operators of renormalized single particles, where the higher-order terms are renormalized to the single-particle coefficient. In the insulating side as well, single-particle Green’s function $`G`$ describes the charge dynamics and may be given from a quasiparticle description with a gap $`\mathrm{\Delta }_c(q)`$ as
$$G(q,k,\omega _n)^1=i\omega _n+E(q,k)\mu ,$$
(6)
where the Matsubara frequency is $`\omega _n`$, and $`E(q,k)=\pm \sqrt{\mathrm{\Delta }_c(q)^2+\epsilon (q,k)^2}`$ with the bare dispersion $`\epsilon `$. In this expression, $`k`$ is the momentum coordinate perpendicular to the locus of $`\epsilon (q,k)=0`$ and $`q`$ denotes that parallel to $`\epsilon (q,k)=0`$. Here $`E(q,k)`$ is assumed to satisfy $`E(q,k)0(0)`$ in the electron-doped (hole-doped) region (, namely for the pole of the upper (lower) Hubbard branch ). The imaginary part of the self-energy $`\mathrm{Im}\mathrm{\Sigma }`$ and the renormalization factor $`Z`$ is not considered here, because the singularities of the Mott transition are our main interest in this article while we assume that the singularities are not altered by $`\mathrm{Im}\mathrm{\Sigma }`$ and $`Z`$ AssaadZ .
Now we take the hole picture for the hole doping side so that $`E`$ always takes $`E0`$ both in electron and hole doped regions. Then, aside from the rigorous validity of the details of the above form for $`E`$, around the Mott gap edge, one can assume that the dispersion is expanded in terms of $`k`$ as
$$E(q,k)=a(q)k^2+b(q)k^4+\mathrm{}..$$
(7)
We have shifted the chemical potential to cancel the gap $`\mathrm{\Delta }_c`$. Here, the gap edge is not necessarily isolated points in the momentum space, but may be a line or a surface, which may evolve to the Fermi surface in the Fermi liquid. The $`k`$-linear term does not exist because we have assumed that $`k=0`$ is the gap edge: The $`k`$-linear term violates the requirement $`E0`$ for small negative $`k`$. The cubic term should also be vanishing in the region of our interest because it becomes relevant only when $`a`$ becomes sufficiently small, while then the cubic term also violates the initial assumption of $`E(q,k)0`$ for negative $`k`$. The analyticity of the dispersion at small $`k`$ is our assumption. This is actually plausible when the transfer energy for the distant pair of Wannier orbitals well converges to zero with increasing distance in the Hubbard-type models.
The coefficients $`a`$ and $`b`$ are obtained as renormalized values after eliminating the higher order terms of the quasiparticle operator. In the metallic side, the rigid band picture is not justified. However, it is still legitimate to consider the quasiparticle dispersion around the Fermi level and the coefficients $`a`$ and $`b`$ as effective quasiparticle coefficients obtained in the evolution process of the metallic phase. This means that $`a`$ and $`b`$ may depend on the distance from the Mott transition point, whereas they still behave continuously. It should be noted that the variations of $`a`$ and $`b`$ can again be renormalized, which leads to $`\zeta `$ independent $`a`$ and $`b`$ near the transition point.
If the gap edge is given by isolated points and the coefficients $`a(q)`$ at these gap edges approach nonzero positive constants on the verge of the transition, the dynamical exponent $`z`$ characterized by the single-particle dispersion is given by $`z=2`$ as in the generic transition to the band insulator RMPX with a finite effective mass of quasiparticles. We see below that this does not hold any more at the marginal quantum critical point.
From this quasiparticle description, the transition between metals and Mott insulators are described by the change in the quasiparticle dispersions. One of our central statements in this paper is that the criticality of the transition has one to one correspondence with the singularity of the quasiparticle dispersions. When the coefficient $`a`$ stays positive through the transition at finite number of isolated points of gap edge, the free energy has a similar singularity with the transition between the band insulator and metals, which is the continuous transition at zero temperature along the $`T=0`$ boundary in Fig 4(b). The first-order transition may evolve only when $`a`$ becomes zero. In the following, we clarify how the character of the transition and the single-particle dispersion are related each other.
### III.4 Relation to Free Energy Form
Here we relate the quasiparticle dynamics and the singularity of the free energy at the Mott transition in the quantum region. When we take the path integral formalism with the imaginary time $`\tau `$, the singular part of the free energy density is formally written by using the quasiparticle dispersion $`E`$ as
$`F`$ $`=`$ $`(T/N)\mathrm{ln}Z`$ (8)
$`Z`$ $`=`$ $`{\displaystyle \underset{i}{}𝒟\varphi _i(\tau )𝒟\varphi _i^{}(\tau )\mathrm{e}^{S/\mathrm{}}},`$ (9)
$`S`$ $`=`$ $`{\displaystyle _0^{\mathrm{}/T}}d\tau {\displaystyle \underset{i}{}}\varphi _i^{}\mathrm{}_\tau \varphi _i+{\displaystyle _0^{\mathrm{}/T}}d\tau H(\varphi ^{},\varphi ),`$ (10)
where the effective hamiltonian $`H`$ generating the single particle excitation $`E`$ in Eq. (7) is rewritten by using the Grassmann variables $`\varphi _i`$ and $`\varphi _i^{}`$ at the site $`i`$.
The above quasiparticle form leads to the effective action at the chemical potential $`\mu `$ as
$$S=\underset{i,n}{}\varphi ^{}(q,k,\omega _n)(i\omega _n\mu +E(q,k))\varphi (q,k,\omega _n)$$
(11)
with the Matsubara frequency $`\omega _n`$.
Let us first study the filling-control transition. Although the expansion (5) does not hold in the quantum region, the singular part of the free energy at zero temperature still has an expansion with respect to the doping concentration $`X`$. From the Matsubara-frequency and wavenumber dependent path integral form of the quantum dynamics leading from Eqs.(8) to (11), we obtain for the singular part of the free energy density at the transition,
$$F=X\mu T_0^{\mathrm{}}𝑑ED(E)\mathrm{log}(1+e^{(E\mu )/T})$$
(12)
with $`D(E)`$ being the singular part of the density of states of quasiparticles. The Boltzmann constant is taken to be unity for our temperature scale. The particle density measured from the insulating phase is given by
$$X=_0^{\mathrm{}}𝑑Ef(E)D(E)$$
(13)
with the Fermi distribution function $`f(E)1/(e^{(E\mu )/T}+1)`$.
We first consider zero temperature and the case where the dispersion $`E`$ has minima at finite number of isolated points in the momentum space with the dispersion given by Eq.(7). Then the particle density is given by
$$X=A_dk_F^d,$$
(14)
where $`A_d`$ is a dimensionality-dependent constant and $`k_F`$ is the Fermi wavenumber measured from the gap edge at the dispersion minima. Here we ignored the possible anisotropy of $`k_F`$ because it does not alter the essential part of the results below for the scaling properties.
From the above relations, at $`T=0`$, we have
$$F=X\mu +A_d[\frac{2a}{d+2}(\frac{X}{A_d})^{\frac{2}{d}+1}+\frac{4b}{d+4}(\frac{X}{A_d})^{\frac{4}{d}+1}+\mathrm{}].$$
(15)
We note that we are considering only the singular part of the free energy at the transition.
When the dispersion proportional to $`a`$ around these points $`q_1`$ in Eq.(7) is present, the quasiparticle picture predicts that the total free energy measured from the insulator has the lowest order term at $`T=0`$ as
$$F+X\mu aX^{(d+2)/d}.$$
(16)
From this free energy form, the charge susceptibility shows the scaling $`\chi _c(^2F/X^2)^1X^{12/d}`$, which is the same as the transition to the band insulator. Then the first-order transition does not take place, because the charge susceptibility does not diverge for $`d2`$. The metal insulator transition is well defined only at zero temperature. This means that the transition occurs across the $`T_c=0`$ boundary illustrated in Fig. 4.
By starting from this continuous transition at zero temperature, the first-order transition can evolve in two fashions. One possibility is the case where a large Fermi surface satisfying the Luttinger theorem appears immediately upon doping. In this case, one has to take that Fermi surface as the locus $`E(q,k)\mu =0`$ and one gets
$$F+X\mu X^3$$
(17)
by eliminating $`k_F`$ from
$$X_0^{k_F}dk$$
(18)
and
$$F+X\mu _0^{k_F}k^2dk,$$
(19)
where we perform the integrations in the region around the locus $`E\mu =0`$ with the assumption $`a>0`$ everywhere. In this case, one gets the charge susceptibility $`\chi _c(^2F/X^2)^1X^1`$, which means that the transition occurs at the marginal quantum critical point, where $`T_c`$ is still zero. This is because $`\chi _c`$ diverges only at $`X=0`$ at zero temperature. To realize the phase separation (namely, the first-order transition), we need a further additional degeneracy of the dispersion at the gap edge.
The second possibility is the marginal quantum critical point emerging with the vanishing $`a`$ term at some isolated points $`q_0`$ of the gap edge. In this case, we have the lowest order term
$$F+X\mu bX^{(d+z)/d},$$
(20)
with $`z=4`$, which yields
$$\chi _cX^{1z/d}.$$
(21)
The exponent $`z=4`$ appears because we are left with the quartic term proportional to $`b`$ when the quadratic term proportional to $`a`$ vanishes.
### III.5 Electron differentiation
Now it turns out that the two possible ways of realizing the first-order transitions eventually become merged to a unified picture, because even when we have the locus of $`E\mu =0`$ with Eq.(17) being satisfied, a further flattening of the dispersion at the gap edge with vanishing $`a`$ is required to realize the first-order transition. It is unlikely that such a flattening emerges uniformly on the locus $`E=0`$. Instead, it generically occurs from particular points of the locus because of the initial anisotropy of the Fermi surface in the band structure and the anisotropic correlation effects as well. Namely, the point with the smallest amplitude of $`a`$ becomes zero first as a special point of the $`E\mu =0`$ surface. Then a quartic dispersion appears at this special point of the $`E\mu =0`$ surface when the system becomes marginal, which results in $`z=4`$.
Therefore, the Mott criticality of the marginally quantum critical point is also characterized by an inevitable evolution of the electron differentiation, if the large Fermi surface is involved in the metallic side. The singular differentiation generates a quartic dispersion at particular points of the Fermi surface coexisting with dispersive generic part.
### III.6 Comparisons with numerical and experimental results
This large dynamical exponent $`z=4`$ was suggested in several independent numerical calculations for the filling-control transition of the Hubbard and $`t`$-$`J`$ models in two dimensions at $`T=0`$ RMP ; Imada1995 . These are the exponent estimated from the compressibility in the form (21FurukawaImada0 ; FurukawaImada ; FurukawaImada2 ; Kohno , the Drude weight $`D`$ estimated from the form $`DX^{1+(z2)/d}`$ Tsunetsugu ; Nakano , single-particle dispersion Assaad99 and the localization length $`\xi (\mu \mu _c)^{1/z}`$ in the insulator side estimated from Green’s function Eq.(6), where $`\xi `$ is obtained from the Fourier-transformed spatially-dependent Green’s function $`G(r,\omega =0)\mathrm{exp}(r/\xi )`$ Assaad96 . These imply that the Hubbard and the $`t`$-$`J`$ models are located close to this marginal quantum critical point.
Although the quasiparticle picture does not hold, it has numerically been shown that the dynamical exponent indeed becomes $`z=4`$ at the marginal quantum critical point in the one-dimensional Hubbard model with next-nearest-neighbor transfers 25 . This suggests that the form (20) is universally valid irrespective of the applicability of the quasiparticle picture.
It is insightful to compare experimental results obtained for the high-$`T_c`$ cuprates with the present picture of the electron differentiation. In the high-$`T_c`$ cuprates, flat dispersions are universally observed near $`(\pi ,0)`$ and $`(0,\pi )`$ points in the angle-resolved photoemission experiments 7 ; 8 . This flatness is beyond the conventional expectation obtained from the van Hove singularity, while it is a natural consequence, if the Mott critical point at $`T=T_c`$ is located at low temperatures. At and below $`T_c`$, vanishing quadratic dispersion, given by $`a=0`$ should emerge in a region of the expected Fermi surface.
The distance from the marginal quantum critical point in the phase diagram in Fig. 4 may depend on details of materials and models. In addition, the actual $`q_0`$ positions responsible for the marginal quantum Mott criticality may also depend on materials and models. For example, the change in relative amplitude of $`t^{}`$ to $`t`$ in the Hubbard model (3) may change the location of $`q_0`$. This may even change the nature of the transition from the route across the $`T_c=0`$ boundary to the route through the first-order transition illustrated in Fig. 4. This change may indeed occur from one high-$`T_c`$ compounds to another Andersen ; Imada2004 . Actually the singular points $`q_0`$ may deviate from $`(\pi ,0)`$ and $`(0,\pi )`$ points for larger $`t^{}`$. Width of the critical region may be influenced by the amplitude of $`t^{}`$ as well.
For the moment, definite assignments are not possible, but we infer two possibilities for the cuprates. One is that the system is indeed close to the marginal quantum critical point and the Mott criticality is controlled by the flat dispersion. Even when the Mott criticality is controlled by these flat points, the experimental observation of the flat part of the dispersion at the Fermi level may have some difficulty because the strong damping is inevitably accompanied. The arc structure is observed in the angle-resolved photoemission experiments for the underdoped cuprates, which literally means that the Fermi surface around the flat-dispersion region is missing. This implies that the experimental resolution might not allow the detection of the Fermi surface around the flat part, because of the strong damping, while this flat part may govern the criticality.
The other possibility is that the part of the flat dispersion is slightly away from the Mott gap edge. This is inferred from the fact that the flat-dispersion level in the hole-doped cuprates is slightly lower than the arc part around $`(\pi /2,\pi /2)`$, which is the dispersion minimum (in the hole picture) at the real Mott gap edge indicated by experiments 8 ; Ino and by model calculations as well AssaadZ . In this case, with the lowering doping concentration, the system first shows the marginal quantum critical behavior reflecting the flat dispersion. However, with further approaching the real critical point, it crossovers from the marginal quantum criticality characterized by $`z=4`$ to the ordinary class $`z=2`$, which eventually flows to the criticality for the $`T_c=0`$ boundary. The real high-$`T_c`$ cuprates appears to have a variety between these two possibilities depending on the compounds. We will discuss consequences of the latter case further in Sec. III.11.
Aside from these details and uncertainties, the overall structure of very different evolutions of the Fermi surface depending on the momenta is completely consistent with the picture that the transition metal compounds including the cuprates show electron differentiation arising from the proximity of the marginal quantum Mott criticality. Electron differentiation should become prominent when the system becomes closer to the marginal quantum critical point, while the electrons are more or less uniform along the $`T_c=0`$ boundary.
### III.7 Breakdown of Ginzburg-Landau-Wilson scheme
From the results in Secs. III.3 and III.4, the free energy near the marginal quantum critical point is generally expressed as
$$F=\mu X+aX^{(d+2)/d}+bX^{(d+4)/d},$$
(22)
where $`a`$ and $`b`$ have absorbed numerical constants in Eq.(15) as well as effects of the renormalization factor $`Z`$ being less than unity, and the constants have properly been rescaled.
Even in the case of the bandwidth-control transition, when one can regard the closing of the gap by hole doping around a point $`q_{0h}`$ and simultaneous particle doping around $`q_{0p}`$ with the constraint of keeping the electron density $`n=1`$, the above relation Eq.(22) may be replaced with $`D`$ as
$$F=UD+aD^{(d+2)/d}+bD^{(d+4)/d},$$
(23)
because the “doublon” and “holon” concentrations are nothing but the above self-doping concentration of particles and holes.
Now instead of Eq.(5), the free energy at zero temperature is expanded by $`\zeta `$ and obtained after rescaling of the parameters as
$$F=\mu _\zeta \zeta +a\zeta ^{(d+2)/d}+b\zeta ^{(d+4)/d}+c\zeta ^{(d+6)/d}\mathrm{}$$
(24)
with the constraint $`\zeta 0`$. It should be noted that this expansion of the free energy in terms of $`\zeta `$ is obtained from the path integral form with the spatial as well as imaginary time dependence explicitly taken into account. This form of the free energy clearly violates the Ginzburg-Landau-Wilson scheme because the form of the free energy itself has $`d`$ dependent nonanalytic expansion.
The metal-insulator transition across the $`T_c=0`$ boundary with $`a>0`$ is driven by the $`\mu _\zeta `$ term in Eq.(24). When $`\mu _\zeta `$ is negative, the free-energy minimum exists at $`\zeta =0`$, which corresponds to an insulator. Whereas the metallic phase is represented by the minimum at a nonzero positive $`\zeta `$, which is realized by a positive $`\mu _\zeta `$. The $`T_c=0`$ boundary is determined from $`\mu _\zeta =0`$. The criticality of this $`T_c=0`$ boundary is given for the order parameter as
$$\zeta |\mu _\zeta |^\beta .$$
(25)
In the “mean-field approximation”, $`\beta `$ is obtained from the spatially uniform derivative $`F/\zeta =0`$ as
$$\beta =d/2.$$
(26)
The susceptibility is given by
$$\chi _\zeta =[\frac{d^2F}{d\zeta ^2}]^1[\frac{2(d+2)}{d^2}a\zeta ^{2/d1}]^1,$$
(27)
leading to the susceptibility exponent defined by
$$\chi _\zeta |\mu _\zeta |^\gamma $$
(28)
in the metallic side $`\mu _\zeta <0`$ as
$$\gamma =1d/2.$$
(29)
At $`\mu =0`$, we obtain another exponent from
$$\chi _\zeta a^1\zeta ^{1\delta }$$
(30)
with
$$\delta =2/d$$
(31)
in the metallic side. These are all “mean-field exponents”, although we will show these exponents are indeed correct. In the Ginzburg-Landau mean-field theory, the critical exponents do not depend on the dimensionality. However, this $`d`$-dependent form of the free energy leads to the $`d`$ dependent exponents even in the mean-field treatment.
When $`a`$ becomes zero, the critical point $`\mu _\zeta =0`$ becomes marginal and the first-order transition evolves if $`a`$ becomes negative. This point with $`\mu _\zeta =a=0`$ at $`T=0`$ is nothing but the marginal quantum critical point. Namely, the marginal quantum critical point at $`T=0`$ may be reached at a control parameter $`g=g_c`$, for $`a=a_0(gg_c(T))`$, $`a_0>0`$ and $`b>0`$.
One might argue that the marginal quantum critical point looks similar to the conventional tricritical point Tricritical because the continuous transition converts to the first-order transition at $`T=0`$ at this point. However, it is qualitatively different because the Mott transition contains only the insulator and metal phases and no additional competitions a priori exist.
Similarly to the transition across the $`T_c=0`$ boundary, the critical exponents of the marginal quantum critical point have $`d`$ dependent forms even in the mean-field treatment. The exponent $`\beta `$ defined by the order parameter at $`g<g_c`$ and $`\mu _\zeta =0`$ as
$$\zeta |gg_c|^\beta $$
(32)
is given by
$$\beta =d/2.$$
(33)
Near $`g=g_c`$, $`\chi _\zeta `$ is expressed as
$$\chi _\zeta =[\frac{d^2F}{d\zeta ^2}]^1[\frac{2(d+2)}{d^2}a_0(gg_c)\zeta ^{2/d1}+\frac{4(d+4)}{d^2}b\zeta ^{4/d1}+\frac{6(d+6)}{d^2}c\zeta ^{6/d1}]^1.$$
(34)
Then
$$\chi _\zeta |gg_c|^\gamma $$
(35)
holds in the metallic side $`g<g_c`$, yielding the “mean-field” exponent
$$\gamma =2d/2.$$
(36)
At $`g=g_c`$,
$$\chi _\zeta b^1\zeta ^{1\delta }$$
(37)
with
$$\delta =4/d$$
(38)
is obtained. We note that the divergence of the susceptibility is stronger at the marginal quantum critical point than along the $`T_c=0`$ boundary. It is also stronger for lower spatial dimensions. For example, $`\gamma `$ is 1 and 1/2 in two and three dimensions, respectively, for the marginal quantum criticality. On the contrary,$`\gamma `$ is not positive in two and three dimensions along the $`T_c=0`$ boundary.
Even though the Ginzburg-Landau-Wilson scheme breaks down, the scaling relations $`\beta \delta =\gamma +\beta `$ and $`\alpha +2\beta +\gamma =2`$ are satisfied both for the $`T_c=0`$ boundary and the marginal quantum critical point. Along the $`T_c=0`$ boundary we obtain $`\alpha =1d/2`$ while $`\alpha =d/2`$ for the marginal quantum criticality. They are consistent with the hyperscaling law $`2\alpha =(d+z_t)\nu `$ with $`\nu =1/2`$ and $`z_t=4`$ for the marginal quantum criticality and $`z_t=2`$ along the $`T_c=0`$ boundary. Here, $`z_t`$ should be the dynamical exponent of this Mott transition, namely the dynamical exponent for the density or doublon density correlations, which turns out to coincide with the dynamical exponent for the single particle excitations, $`z`$ determined from the quasiparticle dispersion.
We can also confirm that the hyperscaling relation is satisfied in the following way: The scaling relation and the exponents are derived from the scaling form of the free energy,
$$F(a,\mu _\zeta )=\xi ^{dz_t}f(a\xi ^{y_g},\mu _\zeta \xi ^{y_\mu })$$
(39)
with a scaling function $`f`$ and the correlation length
$$\xi (gg_c)^{1/2}\zeta ^{1/d},$$
(40)
which implies $`y_g=2`$. Here the crossover exponent $`y_\mu =4`$ ($`y_\mu =2`$) is derived from the dynamical exponent of the density fluctuations given by $`z_t=4`$ ($`z_t=2`$) for the marginal quantum criticality ($`T_c=0`$ boundary), respectively. The hyperscaling relation holds because this scaling form (39) is satisfied. In fact this scaling form is derived from the single length scale $`\xi `$ which diverges at the transition point. This correlation length is indeed proportional to the mean carrier distance $`X^{1/d}`$ in the filling-control transition. This is obviously the single length scale which diverges at the transition. Satisfaction of the hyperscaling form (24) also clearly shows that this criticality is the consequence of the spatio-temporal quantum dynamics of two-particle excitations in the path integral form, where $`d`$ and $`z_t`$ represent spatial and imaginary-time fluctuations, respectively.
### III.8 Two-dimensional case
In one and two dimensions, the powers of expansions in Eq.(24) stay at integers (for example, in 1D, $`(d+2)/d`$ reduces to 3 and in 2D, $`(d+2)/d`$ reduces to 2). In two dimensions, the free energy is reduced to
$$F=\mu _\zeta \zeta +a_0(gg_c)\zeta ^2+b\zeta ^3+c\zeta ^4.$$
(41)
This again does not belong to the conventional scheme of the Ginzburg-Landau-Wilson formalism, because the odd order term (the cubic term here) is not allowed in the conventional Landau expansion from the constraint of the symmetry around the critical point. This has some resemblance to the breakdown of the Ginzburg-Landau-Wilson scheme at the Lifshitz point of the structural transition although the physics contained here is quite different. Here the asymmetry is allowed because the part of negative $`\zeta `$ does not exist. This can be easily understood in the analogy to the trivial metal to the band-insulator transition, where the carrier density cannot be negative either. Even in the transition between metals and band insulators in the noninteracting systems, the free energy has a similar form to Eq.(24) and the Ginzburg-Landau-Wilson scheme does not hold.
In the Mott transition in two dimensions, the susceptibility is given by
$$\chi _\zeta =\left(\frac{d^2F}{d\zeta ^2}\right)^1\frac{1}{2a_0(gg_c)+6b\zeta +12c\zeta ^2}.$$
(42)
Then, $`\gamma =1,\beta =1`$ and $`\delta =2`$ hold Imada2004 ; Imada2005 . Remarkably, this agrees with recent experimental results on a $`\kappa `$-ET compound, $`\kappa `$-(ET)<sub>2</sub>Cu\[N(CN)<sub>2</sub>\]Cl by Kagawa, Miyagawa and Kanoda Kanoda ; Kanoda3 . The exponents of the finite-temperature critical point estimated by the conductance are indeed consistent with these values for $`\beta ,\gamma `$ and $`\delta `$ within the experimental accuracy. We note that this compound has a structure of highly two-dimensional anisotropy. For the moment, it is not well clarified how the crossover to the three dimensionality arising from weak three-dimensional coupling should appear experimentally. We will discuss below that these unusual exponents are also obtained practically even at finite temperatures, which is relevant in the realistic experimental condition of the finite-temperature Mott transition.
### III.9 Validity of the mean-field theory
Very close to the marginal quantum critical point, the present mean-field exponent only marginally breaks down, because the Ginzburg criterion Goldenfeld $`d+z_t(2\beta +\gamma )/\nu `$ with $`\nu =1/2`$ being the correlation length exponent indicates that the system is always at the upper critical dimension irrespective of the dimensionality, because the equality $`d+z_t=(2\beta +\gamma )/\nu `$ always holds. Here, $`\nu =1/2`$ is a direct consequence of $`y_g=2`$. The fluctuation beyond the mean-field theory becomes irrelevant above the upper critical dimension. Although logarithmic corrections may exist, the mean-field description is thus basically correct at any dimension in this case. This explains why the “mean-field exponents” are observed in the $`\kappa `$-ET compound. Below we restrict our analysis to the mean-field study because the primary exponents are correct. Although the mean-field exponents are correct, the hyperscaling also holds at any dimension. This peculiar compatibility is explained by the fact that the system is always at the upper critical dimension at any dimension. Detailed analysis of corrections based on the renormalization group study will be reported elsewhere. The same argument for the validity of the mean-field exponents is applied for the $`T_c=0`$ boundary.
Although the diverging density fluctuation is an inevitable consequence of the Mott critical point for the marginal quantum criticality and the classical Ising criticality, it is highly nontrivial effect from the viewpoint of the weak coupling picture. In fact, naive perturbation expansions result in suppressions of the density fluctuations when the Mott transition is approached and the available perturbative treatment fails in reproducing this criticality. The one-loop calculation does not account for the Mott criticality, which in principle has to be explained in the real part of the self-energy to be calculated from the higher order loops, while it is so far an open issue to be derived in the future. In this sense, even the mean-field theory is not a straightforward framework in contrast with most of the mean-field theories as those in magnetic transitions.
### III.10 Finite temperature effect
The parameter $`g`$ also has temperature dependence, because the quasiparticle dispersion in general has a temperature dependence. As we have clarified above, $`g`$ is determined from the quasiparticle dispersions. Therefore $`gg_c`$ can be replaced by $`TT_c`$, where the critical temperature becomes nonzero. This is the dominant finite-temperature effect at low temperatures. However, we have quite independent origin of the finite-temperature effects originating from the entropy term $`T\zeta \mathrm{ln}\zeta `$ in addition to the above temperature dependence in the free energy $`F`$. This generates an essentially singular contribution for $`\mu _\zeta <0`$, because the free energy has the extremum at $`\zeta _0\mathrm{exp}[\mu _\zeta /T]`$. This is an exponentially small contribution at low temperatures for $`\mu _\zeta <0`$. Therefore, it does not contribute to the present scaling behavior near $`T=0`$ in the power of $`T`$. At high temperatures, however, without the Fermi degeneracy, the expansion around the extremum reproduces the regular Ginzburg-Landau form
$$F_T=\frac{A_0(TT_c)\zeta ^2}{2}+\frac{B\zeta ^4}{4}\mu _\zeta \zeta ,$$
when we redefine as $`\zeta \zeta _0\zeta `$.
When $`T`$ is nonzero, strictly speaking, the Ising universality class may appear in the region extremely close to the critical point even at low temperatures. However, this real critical region becomes exponentially narrow with decreasing temperatures and the quantum criticality governs out of this region. This crossover of the criticality is due to the essentially singular contribution of the entropy term.
The results obtained for V<sub>2</sub>O<sub>3</sub> Limelette indicates that $`T_c`$ is above 400 K and is high enough so that the Ising classical universality is well observed. On the other hand, $`T_c`$ for $`\kappa `$-ET<sub>2</sub>Cu\[N(CN)<sub>2</sub>\]Cl is below 40K and low enough Kanoda , so that the quantum region dominates in the experimental results. All of these are consistent with what were observed. More quantitative analyses are left for future studies.
### III.11 Effects of long-range part of Coulomb interaction for filling-control transitions
The role of the long-range part of the Coulomb interaction between electrons is sometimes controversial. In this subsection, we clarify how the long-range interaction modifies the conclusions obtained for the models only with the short-range interaction. Because the phase separation occurs for models with only the short-range force as the first-order transition in the classical region, we restrict ourselves to the case of the Ising classical universality in this subsection, Sec.III.11. Because of the long-range part and the resultant electrostatic condition, the real phase separation into two different electron densities is not allowed in the filling-control Mott transition. Reflecting the electrostatic condition, the diverging charge susceptibility at strict zero wavenumber is eventually suppressed for the filling control transitions. However, it still causes critical fluctuations at nonzero and small wavenumbers.
The Poisson equations for the external test charge $`\rho _{ext}`$ and the induced charge $`\rho _{ind}`$ are given by
$`^2\varphi _{ext}`$ $`=`$ $`{\displaystyle \frac{\rho _{ext}}{ϵ}},`$
$`^2\varphi _{ind}`$ $`=`$ $`{\displaystyle \frac{\rho _{ind}}{ϵ}},`$ (43)
where $`\varphi _{ext}`$ is the electrostatic potential generated by the external test charge, while the induced charge generates additional potential $`\varphi _{ind}`$. The dielectric constant is $`ϵ`$. When the external test charge $`\rho _{ext}`$ is at the origin with the unit charge, in metals, the Thomas-Fermi screening occurs as $`\rho _{ind}=\alpha \varphi _{tot}`$ for weakly $`k`$-dependent term with the definition of the total electrostatic potential $`\varphi _{tot}=\varphi _{ext}+\varphi _{ind}`$. In fact, by the Thomas-Fermi screening, $`\alpha `$ is estimated as $`\alpha =2meq_F/(\pi \mathrm{}^2)`$, where $`m,e`$ and $`q_F`$ are the electron effective mass, charge and the Fermi wavenumber, respectively. Here we have assumed isotropic sphere of the Fermi surface for simplicity. Since the external test charge at the origin leads to a wavenumber independent form $`\rho _{ext}(q)=1`$ in the momentum representation, we obtain
$$\varphi _{tot}=\frac{1}{\alpha +ϵq^2}$$
(44)
and the resultant induced charge as
$$\rho _{ind}=\frac{\alpha }{\alpha +ϵq^2},$$
(45)
which indicates the Yukawa-type normal screening, $`\rho _{ind}(r)=\alpha \mathrm{exp}[2\sqrt{\pi e\alpha }r]/r`$ in the real space. When one considers the susceptibility as the linear response induced by the imposed external charge as the perturbation, the susceptibility $`\chi `$ defined by $`\rho _{ind}=\chi \rho _{ext}`$ is now given by
$$\chi =\frac{\alpha }{\alpha +ϵq^2}.$$
(46)
Whereas, if the Coulomb potential is regarded as the perturbation, the linear response defined by $`\rho _{ind}=\chi _c\varphi _{ext}`$ is given by
$$\chi _c=\frac{ϵ\alpha q^2}{\alpha +ϵq^2},$$
(47)
which shows that the charge susceptibility is suppressed near $`q=0`$ because of the long-range Coulomb force. With this screening, the free energy in the classical region is modified to
$$F(r)=𝑑r^{}\frac{1}{2}\chi _c^1(rr^{})X(r^{})^2+\frac{1}{4}bX(r)^4e\varphi X,$$
(48)
where the Fourier transformed susceptibility is given by
$`\chi _c^1(q)`$ $`=`$ $`\chi _{c\mathrm{H}}^1(q)+{\displaystyle \frac{1}{ϵq^2}},`$ (49)
$`\chi _{c\mathrm{H}}^1(q)`$ $`=`$ $`{\displaystyle \frac{1}{\alpha }}+c(\alpha )q^2.`$ (50)
Here, $`\chi _{c\mathrm{H}}`$ is the charge susceptibility for the Hubbard model without the long-range Coulomb part. The second term in the right hand side of Eq.(49) represents the long-range Coulomb interaction. The scaling of $`\alpha `$ near the critical point is determined from the effective mass $`m`$ and is given as $`\alpha \xi ^{4d}`$ and the function $`c(\alpha )`$ scales as $`c(\alpha )\xi ^{2d}`$. The charge susceptibility is enhanced except for extremely small wavenumber.
In the realistic condition, the minimum of Eq.(48) appears at $`q_{min}1/(ϵc)^{1/4}`$. In the cuprates, this is roughly estimated to be $`q_{min}0.1\pi `$, since $`ϵ`$ is the order of 10-100, and $`c10`$ as we see below. Namely, the charge susceptibility becomes strongly enhanced in the region of small wavenumbers, which are in the order of one tenth of the Brillouin-zone size, $`2\pi `$. Even around the marginal quantum Mott criticality, this induces a “softening” of the charge response with poor screening in the nm length scale, which causes strong dynamical fluctuations of electron density at these small wavenumbers. This provides us with mechanisms of various unusual properties for metals near the Mott insulator. If $`F`$ becomes negative at $`q_{min}`$, an instability toward the charge ordering occurs.
Even when $`T`$ is below $`T_c`$ in the Hubbard-type models, where only the short-range force is considered, the phase separation dynamics toward the $`q=0`$ mode in reality freezes at a stage of a finite $`q`$ constrained from the long-range repulsion for the filling-control transition.
It should be noted that in the bandwidth-control transition, the divergence at $`q=0`$ indeed occurs. This is because the electrostatic condition is not violated for the diverging doublon susceptibility.
We now realize that the larger dielectric constant induces the larger enhancement of the charge susceptibility as obtained by substituting $`q_{min}`$ into Eq. (49). Actually, if good metallic carriers independently exist in addition to the strongly correlated electrons which yields the doped Mott insulator, such good metallic carriers efficiently screen the long-range Coulomb interaction of the correlated electrons. This allows more possibility for the part of the correlated electrons to approach the marginal quantum critical point. In fact, when the good metallic carriers perfectly screen the motion of the correlated electrons, it corresponds to the limit $`ϵ\mathrm{}`$ in the above argument and the diverging density susceptibility for the correlated part of electrons can indeed occur at $`q=0`$, because the density fluctuations are completely compensated by the good-metallic carriers except for the onsite interaction. This may occur in a two-band system, where one band has large bandwidth supplying good metallic carriers, and the electrons on the other band are strongly correlated near the Mott insulating phase. Although it is not a simple two-band system, the situation in the hole-doped cuprate superconductors is in a sense ideal from this viewpoint, because upon the carrier doping, “good metallic carrier” first appears around the $`(\pi /2,\pi /2)`$ region of the Fermi surface, while the carriers near $`(\pi ,0)`$ region with strong correlation effects become doped in the presence of these itinerant carriers. Then the instability for the phase separation or the inhomogeneity actually occurs as the inhomogeneity of the “$`(\pi ,0)`$ carriers” compensated by the ”$`(\pi /2,\pi /2)`$ carriers”. To realize this situation, it appears to be important to recognize that the cuprates are located close to the marginal quantum critical point, but strictly speaking located slightly in the side of the $`T_c=0`$ boundary.
Even when the charge ordering or the phase separation does not occur, the spatial inhomogeneity is easily driven by impurity potential or lattice distortions, because of the underlying enhanced density susceptibility. This may be relevant as the mechanism of the structure observed by scanning tunnel microscopy (STM) in the cuprates and manganites 9 ; 10 ; 11 ; 12 ; 13 . The spatial inhomogeneity in the long length scale (typically at 1 to 10 nm scale) observed experimentally cannot be explained by the naive Thomas-Fermi screening length $`\lambda _{\mathrm{TF}}`$, since nominally $`\lambda _{\mathrm{TF}}`$ is the order of $`1/q_F0.3`$nm. The inhomogeneous structrure may also be enhanced by the experimental condition probed at the surface. The present results obtained from the Mott criticality has a tight connection to the approach from dynamical stripe fluctuations 20 ; 21 , while the importance of the underlying Mott criticality has not been recognized in the literature.
## IV Non-Fermi-Liquid Properties
### IV.1 Mode coupling theory
Now we discuss the consequences of the enhanced density susceptibility near the marginal quantum critical point. Within the mean-field theory, we assume the dominant part of the susceptibility in the vicinity of a small and nonzero momentum $`Q`$ and around the zero frequency as
$$\chi _\zeta (q,\omega )=\frac{\mathrm{\Gamma }^1}{i\omega +D_s(K^2+(qQ)^2+\mathrm{})},$$
(51)
where $`D_s`$ is the diffusion constant of the density fluctuations, $`K`$ is the distance from the marginal quantum critical point, and $`\mathrm{\Gamma }`$ is a constant. Near the critical point, they follow the critical scaling as $`\mathrm{\Gamma }^1\xi ^d,K\xi ^1`$ and $`D_s\xi ^2`$, which reproduces the above scaling Eq.(35) with Eq.(40) at $`\omega =0`$. This again satisfies the dynamical exponent $`z_t=4`$ because $`\omega `$ scales as $`D_sK^2\xi ^4`$, while the scaling of $`K`$ reproduces $`\nu =1/2`$. Since it is at the upper critical dimension, this Ornstein-Zernike-type form is justified. We note that the enhancement is much stronger for $`d=2`$ than $`d=3`$.
It should be noted that the characteristic energy scale of this charge fluctuation is much larger than that of spin and orbital fluctuations, because $`c(\alpha )q^2`$ should typically have the energy scale of the Mott gap at the boundary of the Brillouin zone. This is clear because at the zone boundary, the density fluctuation requires generation of the spatially alternating doublons and holons. Then $`D_s`$ has the energy scale comparable to the Mott gap. For the filling control transition, it may have an energy scale as large as several eV as in the case of the cuprates ($`2`$ eV), while for the bandwidth-control transition, this energy scale is not necessarily large. The large energy scale of these Mott fluctuations explains many puzzling properties of metals near the Mott insulator as we will discuss below.
Since the dominant fluctuation occurs at small but nonzero wavenumbers, the conservation law as in the density conservation at $`Q=0`$ does not exist and the dynamical exponent stays at $`z_t=4`$ for the filling-control transition. For the bandwidth-control transition, the real divergence at $`q=0`$ again does not involve the conservation law because the doublon density does not commute with the Hamiltonian and is not conserved. Therefore, the dynamical exponent again stays at four and the form Eq.(51) is justified. This is in contrast with the ferromagnetic fluctuation at $`Q=0`$ for the spin fluctuations, where $`i\omega `$ in Eq. (51) is replaced with $`i\omega /\mathrm{\Xi }`$ with $`\mathrm{\Xi }q`$ so that the dynamical exponent increases by one 28 .
Now we formulate a mode coupling scheme for the electron-density and doublon-density fluctuations originating from the Mott criticality around small $`Q`$. The inverse of the static susceptibility is obtained from Eq.(42) as
$$\chi _\zeta (Q,0)^1=2a_0(TT_c)+6b\zeta _Q+12c\underset{q}{}\zeta _q^2,$$
(52)
where we introduced the average of the density fluctuation as $`\zeta _q^2`$. This is given from the fluctuation dissipation theorem as
$$\zeta _q^2=\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}𝑑\omega \mathrm{coth}(\omega /2T)\mathrm{Im}\chi _\zeta (q,\omega ).$$
(53)
In the inverse of the bare susceptibility, the mass term is renormalized by the mode-coupling term proportional to $`_q\zeta _q^2`$ in Eq.(52). By renormalizing the zero temperature value, $`T_c`$ is renormalized to
$$T_c^{}=T_c\frac{6c}{a_0}\underset{q}{}\zeta _q^2(T=0),$$
(54)
where
$$\zeta _q^2(T=0)\frac{1}{\pi }_0^{\mathrm{}}𝑑\omega \mathrm{Im}\chi _\zeta (q,\omega )$$
(55)
is the zero-temperature fluctuation. Then the renormalized inverse susceptibility is given by
$$\chi _\zeta (Q,0)^1=2a_0(TT_c^{})+6b\zeta _Q+12c\underset{q}{}\delta \zeta _q^2$$
(56)
with the definition for the finite-temperature correction;
$$\delta \zeta _q^2=\frac{2}{\pi }_0^{\mathrm{}}𝑑\omega \frac{1}{\mathrm{exp}(\omega /T)1}\mathrm{Im}\chi _\zeta (q,\omega ).$$
(57)
We can solve Eqs.(51), (56) and (57) selfconsistently, where we take into account the Gaussian fluctuations of the density through the mode coupling. In the mode-coupling scheme, the critical exponents stay at the mean-field values, which is justified from the above arguments in Sec. III.9. However, effects of the fluctuations are taken into account in a selfconsistent fashion. A similar phenomenological theory has been formulated by Moriya for spin fluctuations 28 .
### IV.2 Perturbative self-energy
In the following, we employ a perturbative treatment to understand the non-Fermi-liquid behavior as well as superconducting instability. The perturbative scheme has a limited applicability for these highly nonperturbative phenomena of the density fluctuations. However, when the Mott criticality is properly taken into account through the phenomenological treatment of the density fluctuations, we expect that the essence can be captured even when we employ a perturbation theory for the other part of calculations.
In the actual calculations for the realistic choice of parameters, a reasonable set of parameters can be derived in the following: We choose the parameter values of a filling-control transition for a two-dimensional system appropriate for the copper oxide superconductors inferred from the frequency dependence of the optical conductivity 29 ; 30 , characteristic size of the observed inhomogeneity 9 ; 10 and the doping dependence of the density susceptibility in numericalFurukawaImada and experimental results, which suggest $`a_00,b0.7,c100,\mathrm{\Gamma }^13X`$ and $`D_s30X`$ by taking the energy unit $`t(0.4`$eV) and the lattice constant as the length unit. The characteristic wavenumber and energy of the density fluctuations are roughly $`\pi /10`$ and 0.5-1eV, respectively, which determines $`b`$ and $`D_s`$, respectively. In fact, the dielectric function $`ϵ(q=0,\omega )`$ obtained from the optical conductivity provides us with $`\mathrm{Im}\chi (q=0,\omega )=\mathrm{Im}[1/ϵ(q=0,\omega )]1`$. The obtained results for $`\mathrm{Im}\chi `$ deduced from the experimental data for the optical conductivity have a prominent peak structure around 0.5-1 eV 30 , which indicates the characteristic energy scale of the density fluctuations. After considering the screening effect by the long-range Coulomb part, this suggests $`D_s`$ has the order of 1eV. The uncertainty of the parameters remains because of the lack of accurate experimental probes to estimate the frequency and wavenumber dependences of dynamical density fluctuations. Basically, all the results presented here do not depend on $`Q`$ within the choice $`0<Q<0.2\pi `$.
As an example, we consider the carrier doping in the Hubbard model with the dispersion of the square lattice $`E(q)=2t(\mathrm{cos}q_x+\mathrm{cos}q_y)`$ with additional input of the density fluctuations given by Eq.(51). Through the mode coupling, the solution of the selfconsistent equations (51),(56) and (57) shows that the Curie-Weiss type behavior $`\chi _\zeta (T+\mathrm{\Theta })^1`$ holds in an extended temperature region with small Weiss temperature $`\mathrm{\Theta }`$ near the marginal quantum critical point. Figure 5 shows calculated results of such Curie-Weiss bahaviors in an extended temperature region. The criticality stays at the mean-field form, while the Curie-Weiss form is retained over very large temperature region with renormalized values of coefficients. Even when $`a_0=0`$ is employed, we obtain the linear temperature dependence of $`\chi _\zeta ^1`$, namely the Curie-Weiss form in a wide temperature region because of the linear temperature dependence of $`\delta \zeta _q^2`$.
Now we calculate the electron self-energy. The electron self-energy in the perturbation expansion up to the second order of the interaction for the filling control is given as
$$\mathrm{\Sigma }(q,\omega _n)=\frac{TU^2}{2N}\underset{k,n}{}G(k,i\omega _n)\chi _X(qk,i(\omega _n\omega _m)).$$
(58)
Here the imaginary part of the self-energy $`\mathrm{Im}\mathrm{\Sigma }`$ is governed by the Curie-Weiss behavior of $`\chi _X`$ through Eq.(58). Because of the linear temperature dependence in $`\chi _\zeta ^1`$ in an extended temperature region, we obtain the linear temperature dependence also for $`\mathrm{Im}\mathrm{\Sigma }`$ in a wide temperature region. In the marginal-critical region with $`\chi _XX^1`$, the standard Fermi-liquid behavior $`\mathrm{Im}\mathrm{\Sigma }T^2`$ is replaced with the non-Fermi-liquid form $`\mathrm{Im}\mathrm{\Sigma }T`$, which may cause various unusual properties. The resistivity in two-dimensional systems becomes nearly proportional to $`T`$ as $`\rho \mathrm{Im}\mathrm{\Sigma }T`$ in contrast to the standard Fermi-liquid scaling $`\rho \mathrm{Im}\mathrm{\Sigma }T^2`$ .
A long-standing puzzle in the doped Mott insulators is widely observed long-tail structures in the optical conductivity extending up to the order of 1 eV in various transition metal oxides and organic conductors RMP ; 29 ; 30 ; Lee . The tail has a structure of power law decay in the optical conductivity as $`\sigma (\omega )\omega ^p`$ with $`p`$ ranging between 0.3 and 1. Origin of such long tail structure has to be attributed to fluctuations in the energy scale of 1 eV and cannot be accounted for by the spin and orbital fluctuations, since they have much lower energy scale typically less than 0.1 eV. We note that the density fluctuation mechanism examined in this article naturally accounts for such fluctuations at large energy scale.
## V Superconductivity emerging from Mott criticality
### V.1 Pairing originating from the marginal quantum Mott criticality
It is widely recognized that the origin of the high-temperature superconductivity in the copper oxides Bednorz has to be explained by considering the strong electron correlation effects, although the mechanism is still puzzling and not definitely figured out. After the discovery, various aspects of magnetic mechanisms were extensively examined. From the weak coupling picture Moriya ; Pines , the spin fluctuation theories were considered, where strong antiferromagnetic fluctuations were assumed to mediate the Cooper pairing. Then the origin of the high-$`T_c`$ superconductivity was assumed to arise from the criticality of the antiferromagnetic quantum critical point. In the strong-coupling expansion represented by the $`t`$-$`J`$ model AndersonBaskaran , it was claimed that the pairing is basically through the singlet formation stabilized by the superexchange term proportional to $`J`$ in the $`t`$-$`J`$ model. In both of the approaches, the mechanism of the Cooper pairing is more or less the same and they are categorized as the magnetic mechanism.
On the other hand, from the initial stage of the studies on the cuprate superconductors, it has been well recognized that the superconductivity occurs in the region of the doped Mott insulator near the Mott transition Anderson . However, since the criticality of the Mott transition was not well identified until recently, the role of the Mott criticality for the mechanism of the superconductivity was not well appreciated. In fact, as we already clarified in the previous sections, the Mott transition itself is a transition driven by the order parameter of the electron (or doublon-holon) density, and has nothing to do with the magnetic degrees of freedom by itself. Although the antiferromagnetic fluctuations occur at $`(\pi ,\pi )`$ in the magnetic Brillouin zone for the square lattice, the Mott criticality occurs independently through quite a different fluctuation, namely through the singular density fluctuations at small wavenumber. If the antiferromagnetic order exists at low temperature of the Mott insulating phase, the antiferromagnetic fluctuations around $`(\pi ,\pi )`$ coexist with the density fluctuations around small wavenumber arising from the Mott criticality. It should also be noted that this fluctuation has a completely different origin from the ordinary charge-order fluctuations at a commensurate wave vector, although the present instability may also trigger the charge ordering.
It is naturally expected that the density fluctuations at small wavenumber inherent to the Mott transition may play a novel role in stabilizing the superconducting phase. In fact, the density fluctuations may be the origin of instabilities to various symmetry breakings including not only superconductivity but also charge and magnetic orderings, since the diverging density fluctuations are directly connected with the flattening of the quasiparticle dispersions at the Fermi level leading to the enhanced density of states. The diverging density of states widely enhances the instability for various orders.
The pairing mechanism arising from this density fluctuation was recently studied by a perturbative scheme of the mode-coupling theory Imada2005 . We here discuss superconductivity assuming the proximity to the Mott quantum critical point in detail. In the present scope, the density fluctuations have the largest energy scale over spin and orbital fluctuations and are the primary origin of the unusual properties. This proximity of the Mott transition indeed mediates the Cooper pairing through the enhanced density fluctuations. When we follow the perturbative scheme, the effective interaction between two electrons is obtained from the density fluctuations as
$$\mathrm{\Lambda }(q,i\omega _n)=UU^2\chi _X(q,i\omega _n)/2$$
(59)
up to the second order in $`U`$ with $`\chi _X`$ obtained from (51). Of course other fluctuations as spin fluctuations also affect the effective interaction. However, to extract the role of density fluctuations clearly, we ignore the contribution from spin fluctuations in this article as the first step. In fact, the density fluctuation plays the dominant role because of its large energy scale. When spin fluctuations are also considered, we expect that it reinforces the pairing, since both of the fluctuations enhance the same type of pairing symmetry while they do not interfere each other because of their fluctuations at very different wave numbers as we see below. We obtain the linearized Eliashberg equation for the superconducting gap $`\mathrm{\Delta }`$ as
$`\mathrm{\Delta }(q,\omega _n)`$ $`=`$ $`{\displaystyle \frac{T}{N}}{\displaystyle \underset{k,m}{}}G(k,i\omega _m)G(k,i\omega _m)`$ (60)
$`\times \mathrm{\Lambda }(qk,i(\omega _n\omega _m))\mathrm{\Delta }(k,i\omega _m),`$
where $`N`$ is the number of sites. This Eliashbrg equation is solved selfconsistently. Considering the level of perturbative approximations here, the first nontrivial way to solve this problem is to take the bare Green’s function for $`G`$, and we ignore the normal self-energy corrections to $`G`$. In the calculation of Green’s function in Eq.(60), the standard Hubbard model on the square lattice is employed as an example. However, the nonperturbative effect is taken into account through Eq.(51) for $`\chi _X`$ with the parameter values introduced above. Then Eq.(59) is inserted to Eq.(60) and the eigenvalue $`\lambda `$ is calculated for the right-hand side of Eq.(60). Namely, the linearized Eliashberg equation is solved selfconsistently for the relevant parameter values for the cuprate superconductors as cited above.
### V.2 Unconventional pairing
The solution of the Eliashberg equation Eq.(60) for the parameter values above shows that the right hand side of the linearized Eliashberg equation has the largest eigenvalue for the $`d_{x^2y^2}`$ pairing symmetry, which leads to the highest superconducting transition temperature $`T_{sc}`$ for this pairing symmetry as we show in Fig. 6. Figure 7 shows that the $`d_{x^2y^2}`$ pairing symmetry indeed wins over the eigenvalues for the other symmetries including the symmetry of the extended $`s`$-wave-symmetry pairing. It is remarkable that even though the fluctuations are at a small wavenumber $`Q`$, it generates the anisotropic pairing. This is because the effective interaction is repulsive in the most part of the Brillouin zone because of the first term in Eq. (59), while it becomes attractive only in the small wavenumber region. This forces the pairing to have an anisotropy with nodes. The $`d_{x^2y^2}`$ symmetry and its node position are understood because the largest gap grows in the $`(\pi ,0)`$ and $`(0,\pi )`$ regions, which is stabilized by the flat dispersion. Then the only possibility is to make nodes in the diagonal direction in the Brillouin zone. We note that the gap amplitude may be substantially underestimated because we have underestimated the flatness of the dispersion in Green’s function by taking the bare Green’s function instead of the correct one. Figure 6 shows how the eigenvalue grows with lowering temperatures. The superconducting transition temperature within this approximation is estimated from the temperature where the eigenvalue exceeds unity in Fig. 6. The transition temperature has the order of $`0.01t`$ to $`0.05t`$ as we see in Fig. 8, which corresponds to the order of 100K for the copper oxides when we take $`t0.4`$eV. It should be noted that the large energy scale (namely the Mott gap scale $`2`$ eV) of the density fluctuation represented here by the parameter value $`D_s=30X`$ is crucial for achieving such a high transition temperature. In other words, the high-energy excitations substantially contribute to enhancing the transition temperature. It is remarkable that within this simple approximation, the quantum Mott criticality has a dramatic effect on the superconductivity, which is comparable or even larger than that by the magnetic mechanism in the same level of approximations.
If we properly consider the self-energy effects in Eq.(60), we expect that the superconducting transition temperature $`T_{sc}`$ becomes vanishing at the Mott transition point although the pairing interaction is most enhanced at the Mott transition point because of the critical enhancement of the density fluctuations. These two should cause the separation of $`T_{sc}`$ and the gap amplitude. It leads to the pseudogap behavior in the underdoped region. This is left for future studies. We note that one has to be careful in taking account of the self-energy effects because a part of it appears through the density fluctuation itself, which is already taken into account here but is beyond the presently available perturbative treatment in the literature.
The present analyses based on the perturbative treatment and mean-field-type Eliashberg equation do not take into account fluctuation effects particularly in two-dimensional systems. In fact, in purely two-dimensional systems, we expect Berezinskii-Kosterlitz-Thouless (BKT) type transition for the gauge symmetry breaking of the superconductivity and this aspect is not considered here. Nonetheless, the present analyses have significance in the following points: First, even when the BKT transition is expected in pure two-dimensional system, its transition temperature has a comparable value to the mean-field results as known in the analyses of the XY model universality. Therefore, the present mechanism of the superconducting transition may also work for the BKT transition at the similar temperature scale, which can be inferred from the present simple approximations. Second, the significance of the present approach is that the superconducting mechanism arises from a completely new origin of the proximity from the quantum Mott criticality with the enhanced density fluctuations at small wavenumber. As a first step, clarification of possible relevance in the experimental situation is desired even at the mean-field level. Third, it is useful to compare consequences of the present mechanism with the conventional ones including magnetic fluctuation mechanisms at the same level of approximations. For the mean-field analyses, it is possible because we have many available results for the conventional mechanisms in the literature. Fourth, in the cuprates, very weak interlayer coupling, which still does not destroy the dominance of two-dimensional Mott criticality in a wide region, may sensitively induce real superconducting transition as the three-dimensional one. This circumstance may show that the mean-field treatment for the superconducting transition offers a qualitatively correct way of understanding if the pairing mechanism is correctly picked up. More detailed analyses with consideration of the inherent two-dimensional fluctuations and the resultant BKT transition is left for future studies.
The present pairing mechanism may also work near the bandwidth-control transition point. In fact the basic mechanism can be straightforwardly applied to the region near the critical end point of the bandwidth-control transition. This may explain the superconducting phase observed near the Mott transition point of $`\kappa `$-ET compounds family Kanoda0 .
The superconductivity near the valence instability point was studied theoretically as a model for CeCu<sub>2</sub>Si<sub>2</sub>, CeCu<sub>2</sub>Ge<sub>2</sub>, and other heavy fermion compounds Miyake ; Monthoux . Since we expect a similar novel criticality to the present Mott criticality, it would be an intriguing issue to pursue the mechanism of the superconductivity along the same line.
Since electron density fluctuations of course strongly couple with phonon modes, the phonons also supplement this density fluctuation mechanism. In fact, the strong coupling to small wavenumber phonon necessarily occurs, if the enhanced electron density fluctuations already exist in the Mott critical region. This may reinforce the strong coupling to phonons even when the conventional electron-phonon coupling $`\lambda `$ is rather weak. In fact, the $`B_{1g}`$ out-of-plane phonon mode with a small momentum transfer Cuk may have a relevance and presumable resultant kink structure in the angle-resolved photoemission spectra should be considered under this circumstance.
## VI Summary and Discussion
In this paper, it has been shown that the Mott transition is successfully described by a new framework for quantum phase transitions. The natural order parameter for the Mott transition is the electron doping concentration for the filling-control transition. For the bandwidth-control transition, the natural order parameter is the doublon (or holon) density. At zero temperature, the phase boundaries of metals and Mott insulators always exist as the Mott transition, which occurs either as the first-order or continuous transitions.
If the Mott transition occurs as the continuous transition at zero temperature and other spontaneous symmetry breakings are not involved, the metals and insulators are adiabatically connected at finite temperatures. We call this regime the quantum regime surrounding the $`T_c=0`$ boundary. On the other hand, if the transition occurs as the first-order transition at $`T=0`$ and terminates as the critical point above the Fermi degeneracy temperature, we call it the classical regime. Sandwiched by these two regimes, the marginal quantum critical region appears, where the first-order boundary and the continuous $`T_c=0`$ boundary meet at $`T=0`$.
Although the density fluctuation is completely suppressed in the Mott insulator because of the Mott gap, the criticality of the continuous Mott transition in the metallic side can be described by the critical enhancement of the density fluctuation at small wavenumbers in contrast to the naive expectation. This critical enhancement indeed occurs at finite-temperature Mott critical line as well as at the marginal quantum critical point. The marginal quantum criticality shows nontrivial and novel features.
In the classical regime around the high critical temperature, the Mott criticality is described by the Ising universality class. However, in the quantum regime, the clarified Mott criticality indicates the breakdown of the Ginzburg-Landau-Wilson (GLW) scheme. The quantum dynamics is derived from the path-integral formalism with the one-to-one correspondence between single-particle dynamics and two-particle correlations. The divergence of time scale at the transition is described by the dynamical exponent $`z=2`$ for the quantum region and $`z=4`$ for the marginally quantum region. Remarkably, it is shown that the free energy has nonanalytic expansion with respect to the order parameter, with power depending on spatial dimensionality $`d`$ for the quantum as well as the marginally quantum criticality, in marked contrast with the GLW expansion, which is demensionality independent. This unusual expansion results in dimensionality-dependent critical exponents, which also indicates that the Mott transition occurs always at the upper critical dimension at any $`d`$. Then, at any $`d`$, the scaling relations and the hyperscaling are still satisfied while the mean-field description is basically justified for the critical exponents except for the possible logarithmic corrections. They are totally described by a new universality class. Particularly when the Mott critical temperature becomes lowered just to zero temperature, the marginal quantum critical point appears and the critical exponents are given by $`\gamma =2d/2,\beta =d/2`$ and $`\delta =4/d`$.
The present theoretical framework for the quantum Mott criticality has clarified many aspects which are consistent with the experimental results near the Mott insulator. These are summarized in the following:
First, the universality class of the quantum Mott transition well explains the otherwise puzzling critical exponents recently discovered in $`\kappa `$-(ET)<sub>2</sub>Cu\[N(CN)<sub>2</sub>\]Cl; $`\gamma =1,\beta =1`$ and $`\delta =2`$. This exponents are identified as those at the marginal quantum critical point in two-dimensional systems. The scaling description is completely consistent with both the classical Ising-type transition observed in V<sub>2</sub>O<sub>3</sub> and the quantum transition observed in $`\kappa `$-(ET)<sub>2</sub>Cu\[N(CN)<sub>2</sub>\]Cl as well as in $`\kappa `$-(ET)<sub>2</sub>Cu<sub>2</sub>(CN)<sub>3</sub>.
Second, such unusual exponents and the new universality class at the marginal quantum critical point inevitably cause the differentiation of electrons in the momentum space even when the large Fermi surface with the contained Luttinger volume is expected. This draws a concrete picture how the Fermi liquid breaks down to the Mott insulator. The differentiation along the Fermi surface is the driving mechanism of the emergence of the flat dispersion, and the arc structure observed by the angle resolved photoemission experiments in the cuprate superconductors. The differentiation generates particular points on the Fermi surface responsible for the criticality, which is the reason why the hyperscaling relation is satisfied in the present theory.
Third, approaching the marginal quantum critical point, the system becomes more and more sensitive to the external perturbations and is easily driven to the inhomogeneous state, which has been suggested in various types of surface probes. We have estimated the typical length scale of the inhomogeneity determined from the balance of the Mott criticality and the long-range Coulomb interaction in the filling-control transition.
At the same time, the universality class of the Mott criticality at finite temperatures in the classical region is protected from the randomness because the Ising universality at finite temperatures is not influenced by the small randomness. The first-order transition is only driven by the mechanism of the Mott transition equivalent to the Ising class. The universality class of the critical point at the termination point of the first-order transition may receive fewer effects of Anderson localization, because the Anderson transition does not drive the first-order transition at all. The insensitivity to the randomness is particularly true in three dimensions. In purely two dimensional systems, however, the situation is nontrivial because of the sensitive random field effects on the Ising transition, which will be discussed in a separate paper. This is also the origin of the tendency for the spatial inhomogeneity. At the marginal quantum critical point, we also expect that the universality is protected against randomness at least in three dimensions because of diverging density of states at the gap edge and efficient screening through divergent density susceptibility. Near the marginal quantum critical point, the enhanced density of states implies that the Thomas-Fermi screening length becomes short, which leads to a more efficient screening of the random potential by the fewer carriers. Then contrary to the naive expectation, the density fluctuation at a long length scale appears. Along the $`T_c=0`$ boundary, the metallic phase may be under a severe effects of randomness and the continuous metal-insulator transition is eventually triggered by the Anderson localization.
The filling-control transition requires a special care because of the long-range Coulomb interaction. For charged electrons, the real divergence of the density fluctuation at strictly $`q=0`$ is suppressed. However, the density fluctuations are still strongly enhanced at a small nonzero wavenumber, which may cause unusual properties. In addition, if dispersive light carriers coexist with the carriers near the Mott insulator, the screening and compensation by such dispersive and good metallic carriers allows the critical divergence of the density fluctuation for the correlated carriers. This appears to be realized indeed in the cuprate superconductors.
The fourth point for the experimental relevance of the marginal quantum criticality is the non-Fermi-liquid properties. The non-Fermi-liquid properties observed in the doped Mott insulators with two-dimensional anisotropy are accounted for by the density fluctuations arising from the quantum Mott criticality. This in fact explains the $`T`$-linear resistivity in two dimensions, when we employ the mode coupling theory. In particular, the fluctuations at the energy scale as large as 1 eV observed as the long tail in the optical conductivity universal in transition metal compounds and organic compounds are explained not by spin or orbital fluctuations but by this density fluctuation of the quantum Mott criticality.
We have also shown that the mechanism of the high-$`T_c`$ superconductivity can be ascribed to the density fluctuations originating from the marginal quantum Mott criticality. Although the density fluctuation occurs at small wave number, it causes the unconventional pairing with a nodal structure. For the realistic parameter values of the cuprate superconductors, the solution of the linearized Eliashberg equation has the highest transition temperature of the order of 100K for the $`d_{x^2y^2}`$ wave symmetry . It is remarkable that, within the present level of approximation, the density fluctuations inherent near the quantum Mott criticality overlooked in the literature cause comparable or even larger effects than the spin fluctuation mechanisms extensively studied for the cuprate superconductors. The large energy scale of the density (charge) fluctuations may help even stronger instability towards the superconductivity, if this instability could be more carefully tuned by the design of material parameters. This is a challenging future task.
From the experimental point of view, it is highly desired to develop a good experimental probe for studying the dynamical and short-range density (charge) correlations. In contrast to the magnetic correlations well studied by neutron scattering and NMR, experimental probes for the wavenumber and frequency dependent charge correlations of electrons are poor. We can study the optical conductivity and dielectric functions only at zero wavenumber while STM and several microscopes can detect only the static structure on the surface. Raman scattering is not powerful enough so far for the study of the systematic wavenumber dependence. Electron energy loss spectra and inelastic X-ray scattering in principle probe the dynamical density correlations while the present energy resolution is rather poor. To uncover the whole feature of physics near the Mott transition, it would be highly desired to develop experimental probes for the frequency and momentum dependence of electron density correlations and it will make a breakthrough in this field. Our prediction from the present work is that the extended charge (density) fluctuations at small wavenumbers may be observed in the frequency and wavenumber dependent spectra in the critical region of a certain class of the Mott transitions in transition metal compounds and organic conductors. The density fluctuation may also occur with a compensation between two different types of carriers, good metallic carriers and the strongly correlated carriers near the Mott insulator. This compensated density fluctuation will require a more refined probe to be detected. This will reveal the missing ring of various puzzling properties including the mechanism of the high-$`T_c`$ cuprate superconductors and the criticality of the Mott transition in the organic conductors.
The author thanks D. Basov, K. Kanoda, F. Kagawa, S. Tajima and Y. Tokura for illuminating discussions on their experimental results. This work is supported by the grant-in-aid from the Ministry of Education, Culture, Sports, Science and Technology. A part of the computation has been performed at Supercomputer Center, Institute for Solid State Physics, University of Tokyo. |
warning/0506/math0506152.html | ar5iv | text | # TWISTED GRADED HECKE ALGEBRAS
## 1. Introduction
Drinfel’d defined graded Hecke algebras for any finite subgroup $`G`$ of $`GL(V)`$ . This definition was shown by Ram and Shepler to generalize the definition of graded versions of affine Hecke algebras for real reflection groups given by Lusztig, who was motivated by questions in representation theory . In the same paper Ram and Shepler classified all graded Hecke algebras for complex reflection groups. They found many nontrivial graded Hecke algebras, but also showed that some groups (such as $`G(r,1,n)=/rS_n`$ when $`r>2,n>3`$) have no nontrivial graded Hecke algebras. Etingof and Ginzburg showed that in the case of a finite symplectic group, the symplectic form itself arises naturally in the structure of any associated graded Hecke algebra, always yielding nontrivial examples, which they called symplectic reflection algebras . Many authors have studied representations of graded Hecke algebras, their subalgebras generated by certain idempotents, and connections to the geometry of the corresponding orbifolds $`V/G`$. (See for example .)
In this note we generalize these constructions by incorporating a two-cocycle $`\alpha `$ that represents an element of the cohomology group $`\mathrm{H}^2(G,^\times )`$. Chmutova introduced such a two-cocycle for symplectic groups, showing that it appears naturally in symplectic reflection algebras arising from nonfaithful group representations to Sp$`(V)`$ . Here our motivation comes from numerous papers on orbifolds in which such a cocycle, called discrete torsion, appears ( are just a few), as well as our finding that for some groups there is a nontrivial twisted graded Hecke algebra even if there is no nontrivial graded Hecke algebra (such as Example 2.16 below). We adapt the direct linear-algebraic approach of Ram and Shepler to derive criteria for existence of such nontrivial twisted graded Hecke algebras (Theorem 2.10, Corollary 2.13). We give examples for which the parameter spaces for twisted graded Hecke algebras are larger than that for graded Hecke algebras: $`S_nGL(^n)`$ via the permutation representation (Example 2.17), as well as $`(/\mathrm{}Z)^{n1}SL(^n)`$ (Example 2.16). On the other hand, if $`G`$ is symplectic, this is not the case (Example 2.15). It would be interesting to do a more thorough analysis of the parameter spaces of twisted graded Hecke algebras for various types of groups.
In Section 3 we show that twisted graded Hecke algebras are precisely particular types of deformations of crossed product algebras $`S(V)\mathrm{\#}_\alpha G`$ (Theorem 3.2). We use this idea to show that Ram and Shepler’s isomorphism between the different definitions of (untwisted) graded Hecke algebra given by Drinfel’d and Lusztig arises as an equivalence of deformations whose infinitesimals are thus cohomologous. This also puts previous results into a larger context: There are in general many more infinitesimal deformations (that is, Hochschild two-cocycles) of $`S(V)\mathrm{\#}_\alpha G`$ than those that lift to deformations that are (twisted) graded Hecke algebras. In some cases these other infinitesimals also lift to deformations of $`S(V)\mathrm{\#}_\alpha G`$: See Remark 2.14 herein and . We are not aware of a general result regarding such deformations.
We thank A. Ram for several stimulating conversations. It was a question of his that ultimately led to this work as well as to insight into other related projects. We thank R.-O. Buchweitz for explaining a computation of the relevant Hochschild cohomology to us.
Throughout, we will work over the complex numbers $``$, so that $`=_{}`$ unless otherwise indicated.
## 2. Twisted graded Hecke algebras
Our approach to graded Hecke algebras will be through crossed products (generalizations of skew group algebras), so we begin by summarizing this construction. Our definition below of (twisted) graded Hecke algebras involves a parameter $`t`$, which may be taken to be an indeterminate or any complex number. In Section 3 we will assume $`t`$ is an indeterminate, as we will discuss deformations over the polynomial ring $`[t]`$. Specializing $`t`$ to any nonzero complex number results in a definition of graded Hecke algebra equivalent to those in the literature.
Let $`V=^n`$, $`G`$ a finite subgroup of $`GL(V)`$, and $`\alpha :G\times G^\times `$ a two-cocycle, that is
(2.1)
$$\alpha (g,h)\alpha (gh,k)=\alpha (h,k)\alpha (g,hk)$$
for all $`g,h,kG`$. The action of $`G`$ on $`V`$ induces actions by algebra automorphisms on the tensor algebra $`T(V)`$ and the symmetric algebra $`S(V)`$.
Suppose $`S`$ is any associative $``$-algebra with an action of $`G`$ by automorphisms, for example $`S=T(V)`$ or $`S=S(V)`$. Then we may form the crossed product algebra $`S\mathrm{\#}_\alpha G`$, which is $`SG`$ as a vector space, and has multiplication
$$(rg)(sh)=\alpha (g,h)r(gs)gh$$
for all $`r,sS`$ and $`g,hG`$. This makes $`S\mathrm{\#}_\alpha G`$ an associative algebra as $`\alpha `$ is a two-cocycle. Note that $`S`$ is a subalgebra of $`S\mathrm{\#}_\alpha G`$, via $`S\stackrel{}{}S1`$, and that the subalgebra $`1G`$ of $`S\mathrm{\#}_\alpha G`$ is also known as the twisted group algebra $`^\alpha G`$. We accordingly abbreviate $`sg`$ by $`s\overline{g}`$ ($`sS,gG`$). The action of $`G`$ on $`S`$ becomes an inner action on $`S\mathrm{\#}_\alpha G`$, namely $`s\overline{g}s(\overline{g})^1`$ for all $`sS,gG`$.
The two-cocycle $`\alpha `$ is a coboundary if there is some function $`\beta :G^\times `$ such that
(2.2)
$$\alpha (g,h)=\beta (g)\beta (h)\beta (gh)^1$$
for all $`g,hG`$. The set of two-cocycles modulo coboundaries forms an abelian group under pointwise multiplication, called the Schur multiplier of $`G`$, and denoted $`\mathrm{H}^2(G,^\times )`$. If two cocycles differ by a coboundary, that is if they are cohomologous, the corresponding crossed products are isomorphic.
Replacing $`\alpha `$ by a cohomologous cocycle if necessary, we may assume that $`\alpha `$ is normalized so that $`\alpha (1,g)=\alpha (g,1)=1`$ for all $`gG`$. It follows that $`\overline{1}`$ is the multiplicative identity for $`^\alpha G`$. Thus $`(\overline{g})^1=\alpha ^1(g,g^1)\overline{g^1}=\alpha ^1(g^1,g)\overline{g^1}`$ for all $`gG`$.
For each $`gG`$, choose a skew-symmetric bilinear form $`a_g:V\times V`$ (arbitrary for now, and possibly 0). Let $`t`$ be an indeterminate and extend scalars, $`S\mathrm{\#}_\alpha G[t]:=[t](S\mathrm{\#}_\alpha G)`$. Let $`A`$ be the quotient of $`T(V)\mathrm{\#}_\alpha G[t]`$ by the (two-sided) ideal $`I[t]`$ generated by all
(2.3)
$$[v,w]\underset{gG}{}a_g(v,w)t\overline{g}$$
where $`v,wV`$ and $`[v,w]=vwwv`$. (We have omitted tensor symbols in elements of $`T(V)`$ for brevity.) Note that $`A`$ is additively isomorphic to a quotient of $`S(V)\mathrm{\#}_\alpha G[t]`$ via any choice of linear section $`S(V)T(V)`$ of the canonical projection of $`T(V)`$ onto $`S(V)`$. This quotient is proper if and only if there is more than one way to rearrange factors in a word by applying the relations (2.3). We say that $`A`$ is a twisted graded Hecke algebra if the above map yields $`AS(V)G[t]`$ as vector spaces over $`[t]`$. Equivalently, if we assign degree 1 to elements of $`V`$ and degree 0 to elements of $`G`$, then $`A`$ is a filtered algebra over $`[t]`$ and is a twisted graded Hecke algebra in case the associated graded algebra $`\mathrm{gr}A`$ is isomorphic to $`S(V)\mathrm{\#}_\alpha G[t]`$.
###### Remark 2.4.
Let $`\alpha `$ be the trivial two-cocycle $`\alpha (g,h)=1`$ for all $`g,hG`$, and let $`t=1`$. Then $`A`$ becomes the graded Hecke algebra of .
We will use the techniques of Ram and Shepler to determine the conditions on the bilinear forms $`a_g`$ under which $`A`$ is a twisted graded Hecke algebra. Let $`hG`$, and use (2.3) to obtain two expressions involving $`[h^1v,h^1w]`$. By direct substitution: $`[h^1v,h^1w]=_{gG}a_g(h^1v,h^1w)t\overline{g}`$. By conjugating (2.3) by $`(\overline{h})^1`$:
$`[h^1v,h^1w]`$ $`=`$ $`{\displaystyle \underset{gG}{}}a_g(v,w)t(\overline{h})^1\overline{g}\overline{h}`$
$`=`$ $`{\displaystyle \underset{gG}{}}a_g(v,w)\alpha ^1(h,h^1)\alpha (g,h)\alpha (h^1,gh)t\overline{h^1gh}.`$
Both must be in the ideal $`I[t]`$, and the assumed additive isomorphism $`AS(V)G[t]`$ allows us to equate coefficients of each $`\overline{g}`$ ($`gG`$) in these two expressions. Thus we have $`a_{h^1gh}(h^1v,h^1w)=\alpha ^1(h,h^1)\alpha (g,h)\alpha (h^1,gh)a_g(v,w)`$, and replacing $`v,w`$ by $`hv,hw`$, we obtain the first equation below. Again the isomorphism $`AS(V)G[t]`$ implies that in any expression $`wvuV^3`$, three applications of (2.3) in any order to obtain $`uvw`$ plus an element in $`tA`$ must yield the same element of $`A`$. Comparison results in the second equation below. Thus necessary conditions for $`A`$ to be a twisted graded Hecke algebra are:
(2.5)
$$a_{h^1gh}(v,w)=\alpha ^1(h,h^1)\alpha (g,h)\alpha (h^1,gh)a_g(hv,hw)$$
(2.6)
$$a_g(u,v)(gww)+a_g(v,w)(guu)+a_g(w,u)(gvv)=0$$
for all $`g,hG`$ and $`u,v,wV`$. These equations will be used in the proofs of the following lemma and theorem.
Not only are (2.5) and (2.6) necessary for $`A`$ to be a twisted graded Hecke algebra, but they are also sufficient: The relations (2.3) allow for rearrangement of any expression in $`A`$ to a particular form identified with an element of $`S(V)G[t]`$, and the relations (2.5) and (2.6) imply that such a form is unique. That is, (2.5) is equivalent to uniqueness of the canonical form of $`\overline{h}vu`$, and (2.6) is equivalent to uniqueness of the canonical form of $`wvu`$ ($`u,v,wV`$, $`hG`$), as demonstrated in for the case $`\alpha =1`$ (in the text above Lemma 1.5). The uniqueness of the form of $`wvu`$ is equivalent to the Jacobi identity in $`A`$,
(2.7)
$$[u,[v,w]]+[v,[w,u]]+[w,[u,v]]=0,$$
which is another way to express the condition (2.6). For the uniqueness of the canonical form of a monomial having more than three factors, it is helpful to note that replacing any of $`u,v,w`$ in the left side of (2.7) by an element of degree larger than 1 results in an element in the ideal generated by the left side of (2.7). (For more details, see .)
Let $`(,)`$ be any $`G`$-invariant nondegenerate Hermitian form on $`V`$. Define orthogonal complements of subspaces of $`V`$ via this form. Let $`gG`$ and $`V^g=\{vVgv=v\}`$. We will need the observation that $`(V^g)^{}=\mathrm{im}(g1)`$, which follows from the standard facts $`V^g=\mathrm{ker}(g1)`$, $`\mathrm{im}(g1)(V^g)^{}`$, and $`dim(\mathrm{im}(g1))=dim(\mathrm{ker}(g1)^{})`$. For any $`hG`$, denote by $`h^{}:(V^g)^{}(V^g)^{}`$ the composition of the linear maps $`h:(V^g)^{}V`$, $`VV/V^g`$, and a choice of isomorphism $`V/V^g\stackrel{}{}(V^g)^{}`$.
###### Lemma 2.8.
Suppose that $`A`$ is a twisted graded Hecke algebra for $`G`$ defined by skew-symmetric bilinear forms $`a_g`$ ($`gG`$) on $`V`$. For each $`g1`$, if $`a_g0`$ then $`\mathrm{ker}a_g=V^g`$, $`\mathrm{codim}(V^g)=2`$, and
$$a_g(hv,hw)=det(h^{})a_g(v,w)$$
for all $`hG`$ and $`v,w(V^g)^{}`$.
This proof is essentially the same as in Ram and Shepler \[20, Lemma 1.8\], with appropriate adjustments made for the possibly nontrivial two-cocycle $`\alpha `$.
###### Proof.
Let $`g1`$ with $`a_g0`$. Let $`vV`$ and $`wV^g`$. If $`v`$ is also in $`V^g`$, then $`a_g(v,w)(guu)=0`$ for all $`uV`$ by (2.6). As $`g1`$, there is some $`uV`$ with $`guu`$, so $`a_g(v,w)=0`$. If $`vV^g`$, let $`u=_{k=1}^rg^kv`$, where $`r`$ is the order of $`g`$, so that $`uV^g`$ (possibly $`u=0`$). As before, this implies $`a_g(u,w)=0`$, and we may rewrite this equation using the definition of $`u`$ and (2.5) as
$`0=a_g(u,w)`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{r}{}}}a_g(g^kv,w)`$
$`=`$ $`{\displaystyle \underset{k=1}{\overset{r}{}}}\alpha ^1(g^k,g^k)\alpha (g,g^k)\alpha (g^k,g^{1k})a_g(v,w).`$
The restriction of $`\alpha `$ to the cyclic subgroup $`g`$ generated by $`g`$ is a coboundary since the Schur multiplier of a cyclic group is trivial \[14, Prop. 1.1\]. Therefore there is a function $`\beta :g\times g^\times `$ such that $`\alpha (h,l)=\beta (h)\beta (l)\beta ^1(hl)`$ for all $`h,lg`$. It follows that the above coefficients of $`a_g(v,w)`$ are
$`\alpha ^1(g^k,g^k)\alpha (g,g^k)\alpha (g^k,g^{1k})`$
$`=`$ $`\beta ^1(g^k)\beta ^1(g^k)\beta (1)\beta (g)\beta (g^k)\beta ^1(g^{1k})\beta (g^k)\beta (g^{1k})\beta ^1(g)`$
$`=`$ $`\beta (1),`$
that is they are independent of $`k`$. The above calculation reduces to $`0=a_g(u,w)=r\beta (1)a_g(v,w)`$. As $`r\beta (1)0`$, this forces $`a_g(v,w)=0`$. Consequently $`V^g\mathrm{ker}a_g`$, and so $`(\mathrm{ker}a_g)^{}(V^g)^{}`$.
As $`a_g`$ is nonzero and skew-symmetric, we have $`dim((\mathrm{ker}a_g)^{})=\mathrm{codim}(\mathrm{ker}a_g)2`$. Let $`u,v`$ be two linearly independent elements of $`(\mathrm{ker}a_g)^{}(V^g)^{}`$ with $`a_g(u,v)0`$, so that in particular $`guu0`$ and $`gvv0`$. Let $`w`$ be any element of $`(V^g)^{}=\mathrm{im}(g1)`$, and write $`w=gw^{}w^{}`$ for some $`w^{}V`$. By (2.6), $`w`$ is a linear combination of $`guu`$ and $`gvv`$, two elements of $`(V^g)^{}`$. This implies $`dim(\mathrm{ker}a_g)^{}=dim(V^g)^{}=2`$, so that $`V^g=\mathrm{ker}a_g`$ and $`\mathrm{codim}(V^g)=2`$.
Finally, let $`v,w(V^g)^{}`$, $`hG`$, and
$`hv`$ $`=`$ $`a_{11}v+a_{21}w+x`$
$`hw`$ $`=`$ $`a_{12}v+a_{22}w+x^{}`$
where $`x,x^{}V^g`$ and $`a_{ij}`$ are scalars. Applying (2.5), we have $`a_{h^1gh}(v,w)=\alpha ^1(h,h^1)\alpha (g,h)\alpha (h^1,gh)a_g(hv,hw)`$ on the one hand, while evaluating via the above equations yields
(2.9)
$$\begin{array}{ccc}\hfill a_{h^1gh}(v,w)& =& \alpha ^1(h,h^1)\alpha (g,h)\alpha (h^1,gh)a_g(a_{11}v+a_{21}w+x,a_{12}v+a_{22}w+x^{})\hfill \\ & =& \alpha ^1(h,h^1)\alpha (g,h)\alpha (h^1,gh)(a_{11}a_{22}a_{12}a_{21})a_g(v,w)\hfill \\ & =& \alpha ^1(h,h^1)\alpha (g,h)\alpha (h^1,gh)det(h^{})a_g(v,w).\hfill \end{array}$$
Equating the two expressions for $`a_{h^1gh}(v,w)`$ gives $`a_g(hv,hw)=det(h^{})a_g(v,w)`$, as desired. ∎
Now suppose $`g1`$, $`a_g0`$, and $`hC(g)=\{hGhg=gh\}`$. Then (2.9) is equivalent to $`det(h^{})=\alpha (h,h^1)\alpha ^1(g,h)\alpha ^1(h^1,gh)`$. Applying (2.1) to the triple $`h,h^1,gh`$, we find that $`\alpha (h,h^1)\alpha ^1(h^1,gh)=\alpha (h,h^1gh)=\alpha (h,g)`$ when $`hC(g)`$, so that $`det(h^{})=\alpha (h,g)\alpha ^1(g,h)`$ for all $`hC(g)`$. This condition is independent of the choice of $`\alpha `$ in a given coset modulo coboundaries since coboundaries are symmetric on commuting pairs as is evident from (2.2). This is as expected since the determinant function is independent of such choices. Note that as $`gC(g)`$, (2.9) implies $`det(g^{})=1`$, that is $`gSL(V)`$. Further, in case $`g`$ is $`\alpha `$-regular, that is $`\alpha (g,h)=\alpha (h,g)`$ for all $`hC(g)`$, this determinant condition is simply $`det(h^{})=1`$ for all $`hC(g)`$. For nonregular elements $`g`$, this condition is different from that in the case of the trivial cocycle, leading to new examples such as Examples 2.16 and 2.17 below.
###### Theorem 2.10.
Let $`G`$ be a finite subgroup of $`GL(V)`$, $`\alpha :G\times G^\times `$ a normalized two-cocycle, and $`gG\{1\}`$. There is a twisted graded Hecke algebra $`A`$ with $`a_g0`$ if, and only if, $`\mathrm{ker}a_g=V^g`$, $`\mathrm{codim}(V^g)=2`$, and
(2.11)
$$det(h^{})=\alpha (h,g)\alpha ^1(g,h)$$
for all $`hC(g)`$.
Again the proof is similar to that of Ram and Shepler \[20, Thm. 1.9\] in the untwisted case. See also the paper of Etingof and Ginzburg , who used a criterion of Braverman and Gaitsgory adapted to Koszul algebras over $`G`$ .
###### Proof.
If $`A`$ is a twisted graded Hecke algebra, Lemma 2.8 and subsequent comments show the given conditions hold.
Conversely, suppose the stated conditions hold for $`g`$. Up to a scalar multiple, there is a unique skew-symmetric form on $`V`$ that is nondegenerate on $`(V^g)^{}`$ and has kernel $`V^g`$. Fix such a form $`a_g`$, and let
(2.12)
$$a_k(v,w)=\{\begin{array}{cc}\alpha ^1(h,h^1)\alpha (g,h)\alpha (h^1,gh)a_g(hv,hw),& \text{ if }k=h^1gh\\ 0,& \text{ otherwise}\end{array}$$
for all $`v,wV`$. In order for (2.5) to hold for all pairs of group elements, we must check that the definition of $`a_k`$ is independent of the choice of representative from the conjugacy class of $`g`$. Suppose that $`k=l^1h^1ghl`$, so that there are two ways to define $`a_k`$. One way yields
$$a_k(v,w)=\alpha ^1(hl,l^1h^1)\alpha (g,hl)\alpha (l^1h^1,ghl)a_g(hlv,hlw),$$
and the other yields
$$a_k(v,w)$$
$$=\alpha ^1(l,l^1)\alpha (h^1gh,l)\alpha (l^1,h^1ghl)\alpha ^1(h,h^1)\alpha (g,h)\alpha (h^1,gh)a_g(hlv,hlw).$$
These two expressions for $`a_k(v,w)`$ are indeed equal, as follows from five applications of (2.1), to the triple $`h^1,gh,l`$, to $`l^1,h^1,ghl`$, to $`g,h,l`$, to $`hl,l^1,h^1`$, and to $`h,l,l^1`$. By the discussion before Theorem 2.10, if $`h^1gh=g`$, then (2.12) coincides with $`a_g(v,w)`$. Therefore the definition of $`a_k`$ is independent of the choice of representative from the conjugacy class of $`g`$, and (2.5) holds. The argument of \[20, Lemma 1.8(b)\] applies without change to show that (2.6) holds as well: Their key observation is that (2.6) holds trivially if any one of $`u,v,w`$ is in $`V^g`$, while $`dim(V^g)^{}=2`$ implies that any three elements $`u,v,w`$ of $`(V^g)^{}`$ must be linearly dependent. Substituting a linear dependence relation into the left side of (2.6) yields 0 after some manipulation. Therefore $`A`$, defined by the $`\{a_k\}`$ in (2.12), is a twisted graded Hecke algebra. ∎
In summary, under the conditions in the theorem, $`a_g`$ is determined by its value $`a_g(v,w)`$ on a basis $`v,w`$ of $`(V^g)^{}`$, and for each $`hG`$, $`a_{h^1gh}(v,w)`$ is given by (2.12).
###### Corollary 2.13.
Let $`d`$ be the number of conjugacy classes of $`gG`$ such that $`\mathrm{codim}V^g=2`$ and $`det(h^{})=\alpha (h,g)\alpha ^1(g,h)`$ for all $`hC(g)`$. The sets $`\{a_g\}_{gG}`$ corresponding to twisted graded Hecke algebras $`A`$ form a vector space of dimension $`d+dim(^2V)^G`$.
###### Proof.
By (2.5) and (2.6), the only condition on $`a_1`$ is $`a_1(v,w)=a_1(hv,hw)`$ for all $`hG`$ since $`\alpha (h,h^1)=\alpha (h^1,h)`$, that is $`a_1`$ is a $`G`$-invariant element of $`(^2V)^{}`$. For each conjugacy class of elements $`gG`$ with $`\mathrm{codim}(V^g)=2`$ and $`det(h^{})=\alpha (h,g)\alpha ^1(g,h)`$ for all $`hC(g)`$, there are skew-symmetric forms $`a_g`$ with kernel $`V^g`$, determined by a single form $`a_g`$ (unique up to scalar multiple) for a representative $`g`$, and given by (2.12). ∎
###### Remark 2.14.
The results of this section apply if we relax one of the conditions on the forms $`a_g`$ a little: Allow $`a_g`$ to take values in $`V^g`$, that is $`g`$-invariant polynomials of degree at most 1. Again let $`A`$ be the quotient of $`T(V)\mathrm{\#}_\alpha G[t]`$ by the ideal $`I[t]`$ generated by all expressions of the form (2.3). As the image of each form $`a_g`$ has degree at most 1, the techniques of Ram and Shepler still apply to yield
$$a_{h^1gh}(v,w)=\alpha ^1(h,h^1)\alpha (g,h)\alpha (h^1,gh)\left(h^1a_g(hv,hw)\right)$$
in place of (2.5), and (2.6) is unchanged (but now considered as a relation in $`S(V)`$). That is, these are necessary and sufficient conditions for $`A`$ to be additively isomorphic to $`S(V)G[t]`$, leading to a generalization of Theorem 2.10.
In the remainder of this section, we give several examples.
###### Example 2.15.
(Symplectic groups.) Let $`V`$ be a finite dimensional symplectic vector space over $``$, that is $`dimV`$ is even and there is a nondegenerate skew-symmetric form $`\omega :V\times V`$. Let $`G`$ be a finite subgroup of the symplectic group Sp$`(V)`$ of all invertible linear transformations preserving $`\omega `$. Let $`\alpha :G\times G^\times `$ be the trivial two-cocycle $`\alpha (g,h)=1`$ for all $`g,hG`$. For each element $`gG`$ such that $`\mathrm{codim}V^g=2`$ (the symplectic reflections), let $`a_g=c_g\omega _g`$, where $`\omega _g`$ is defined to be $`\omega |_{(V^g)^{}}`$ on $`(V^g)^{}=\mathrm{im}(1g)`$ and $`0`$ on $`V^g=\mathrm{ker}(1g)`$, and where $`c_g`$ are scalars such that $`c_{hgh^1}=c_g`$ for all $`hG`$. Let $`a_1=c\omega `$ for a scalar $`c`$ and $`a_g=0`$ for all other $`gG`$. By \[9, Thm. 1.3\], the collection $`\{a_ggG\}`$ determines a graded Hecke algebra, called a symplectic reflection algebra. By Corollary 2.13, the parameter space for the possible twisted graded Hecke algebras cannot be larger than that for graded Hecke algebras in the symplectic case (cf. ).
The next two examples, by contrast, involve groups for which the parameter space for twisted graded Hecke algebras is larger than that for graded Hecke algebras.
###### Example 2.16.
(Elementary abelian groups.) Let $`n3`$, $`\mathrm{}2`$, $`V=^n`$ and $`G(/\mathrm{})^{n1}`$ the multiplicative subgroup of $`M_{n\times n}()`$ generated by the diagonal matrices (where $`q`$ is a primitive $`\mathrm{}`$th root of 1):
$`g_1`$ $`=`$ $`\mathrm{diag}(q,q^1,1,\mathrm{},1)`$
$`g_2`$ $`=`$ $`\mathrm{diag}(1,q,q^1,1,\mathrm{},1)`$
$`\mathrm{}`$
$`g_{n1}`$ $`=`$ $`\mathrm{diag}(1,\mathrm{},1,q,q^1)`$
Define a function $`\alpha :G\times G^\times `$ by
$$\alpha (g_1^{i_1}\mathrm{}g_{n1}^{i_{n1}},g_1^{j_1}\mathrm{}g_{n1}^{j_{n1}})=q^{_{1kn2}i_kj_{k+1}}.$$
The relation (2.1) may be checked directly. For each generator $`g_i`$ ($`i=1,\mathrm{},n1`$), as well as for $`g_n=g_1^1\mathrm{}g_{n1}^1`$, we have $`det(h^{})=\alpha (h,g_i)\alpha ^1(g_i,h)`$ for all $`hG`$, as follows from a straightforward computation. Thus we have a twisted graded Hecke algebra in which the forms $`a_{g_1},\mathrm{},a_{g_n}`$ are all nonzero (and other $`a_g=0`$). (Note that some of the relevant determinants are not 1. There is no nontrivial (untwisted) graded Hecke algebra in this case.) If $`\mathrm{}2`$, it may be checked that this gives the full parameter space of twisted graded Hecke algebras corresponding to this cocycle $`\alpha `$. If $`\mathrm{}=2`$, this determinant condition is in fact satisfied for all $`g`$ with $`\mathrm{codim}(V^g)=2`$, and so there is a larger parameter space of twisted graded Hecke algebras in this case. This example is discussed in a different context in \[24, Example 4.1\], which generalizes \[6, Example 4.7\].
###### Example 2.17.
(Symmetric groups.) Let $`G=S_n`$, acting on $`V=^n`$ by the permutation representation. Let $`\alpha `$ be the unique nontrivial two-cocycle (up to coboundary). The Schur representation group for $`S_n`$ is
$$\mathrm{\Gamma }_n=t_1,\mathrm{},t_{n1},\tau \tau ^2=1,t_i^2=1,\tau t_i=t_i\tau ,(t_it_{i+1})^3=1,(t_it_j)^2=\tau (ij2).$$
(See for example \[14, p. 179\].) A two-cocycle $`S_n\times S_n^\times `$ is determined up to coboundary by an irreducible representation of the Schur representation group $`\mathrm{\Gamma }_n`$ for which the central element $`\tau `$ necessarily acts as multiplication by a scalar. As $`\tau ^2=1`$, this scalar must be $`1`$ in case of the nontrivial cocycle $`\alpha `$. Let $`g=(12)(34)`$. The presentation of $`\mathrm{\Gamma }_n`$ may be used to check that
$$det(h^{})=\alpha (h,g)\alpha ^1(g,h)$$
for all $`hC(g)`$. For example, letting $`h=(12)`$, we have $`det(h^{})=1`$. Choose a section $`T:S_n\mathrm{\Gamma }_n`$ of the projection from $`\mathrm{\Gamma }_n`$ to $`S_n`$ that sends $`\tau `$ to $`1`$ and $`t_i`$ to the transposition $`(i,i+1)`$ as follows: Let $`T(12)=t_1`$, $`T(34)=t_3`$, $`T((12)(34))=t_1t_3`$, and choose other images arbitrarily. Then $`\alpha (h,g)\alpha ^1(g,h)`$ is the scalar by which the following element acts on an irreducible representation of $`\mathrm{\Gamma }_n`$:
$$T(12)T((12)(34))T((12)(12)(34))^1T((12)(34)(12))T(12)^1T((12)(34))^1$$
$`=`$ $`T(12)T(12)T(34)T(34)^1T(34)T(12)^1T((12)(34))^1`$
$`=`$ $`T(34)T(12)T(34)T(12)`$
$`=`$ $`\tau ,`$
thus the scalar is $`1`$ as desired. Other elements $`h`$ of $`C(g)`$ may be checked similarly; as $`hdet(h^{})`$ is a group homomorphism, it suffices to check the condition (2.11) on generators of $`C(g)`$. Therefore there is a twisted graded Hecke algebra corresponding to $`\alpha `$ with $`a_g0`$. (Compare with trivial $`\alpha `$, where $`a_g`$ is necessarily 0 for this choice of $`g`$ \[20, Table 1\].) With the choice $`g=(123)`$, we find as in the case of the trivial cocycle $`\alpha `$, that $`1=det(h^{})=\alpha (h,g)\alpha ^1(g,h)`$ for all $`hC(g)`$. Therefore the parameter space of twisted graded Hecke algebras is larger than that of graded Hecke algebras, and involves the conjugacy classes of all $`gG`$ for which $`\mathrm{codim}V^g=2`$.
## 3. Deformations of crossed products
In this section, we prove that the twisted graded Hecke algebras of the previous section are precisely the deformations of $`S(V)\mathrm{\#}_\alpha G`$ of a particular type. We use this connection to deformations to put our results and examples in a larger context.
Let $`t`$ be an indeterminate. Given any associative algebra $`R`$ over $``$ (for example $`R=S(V)\mathrm{\#}_\alpha G`$), a deformation of $`R`$ over $`[t]`$ is an associative $`[t]`$-algebra with underlying vector space $`R[t]=[t]R`$ and multiplication of the form
$$rs=rs+\mu _1(r,s)t+\mu _2(r,s)t^2+\mathrm{}+\mu _p(r,s)t^p$$
for all $`r,sR`$, where $`rs`$ denotes the product of $`r`$ and $`s`$ in $`R`$, the $`\mu _i:R\times RR`$ are $``$-bilinear maps extended to be $`[t]`$-bilinear, and $`p`$ depends on $`r,s`$. Associativity of $``$ implies that $`\mu _1`$ is a Hochschild two-cocycle, that is
(3.1)
$$\mu _1(w,r)s+\mu _1(wr,s)=\mu _1(w,rs)+w\mu _1(r,s)$$
for all $`w,r,sR`$, as well as further conditions on the $`\mu _i`$, $`i1`$. The cocycle $`\mu _1`$ is called the infinitesimal of the deformation. (See or for the details from algebraic deformation theory.)
Recall that $`S(V)\mathrm{\#}_\alpha G`$ is a graded algebra where we assign degree 1 to elements of $`V`$ and degree 0 to elements of $`G`$.
###### Theorem 3.2.
Up to isomorphism, the twisted graded Hecke algebras are precisely the deformations of $`S(V)\mathrm{\#}_\alpha G`$ over $`[t]`$ for which $`\mathrm{deg}\mu _i=2i`$ ($`i1`$).
###### Proof.
Let $`A`$ be a twisted graded Hecke algebra, as defined in Section 2. Let $`v_1,\mathrm{},v_n`$ be a basis of $`V`$ and choose the section of the projection from $`T(V)`$ to $`S(V)`$ in which a word in $`T(V)`$ is written in the order $`v_1^{i_1}\mathrm{}v_n^{i_n}`$. Express all elements of $`A`$ in terms of this section, writing group elements on the right. Such expressions in $`A`$ exist and are unique due to the additive isomorphism $`AS(V)\mathrm{\#}_\alpha G[t]`$. Now let $`r=v_1^{i_1}\mathrm{}v_n^{i_n}\overline{g}`$ and $`s=v_1^{j_1}\mathrm{}v_n^{j_n}\overline{h}`$ be elements of $`A`$. Denoting the product in $`A`$ by $``$, we have
$$rs=\alpha (g,h)v_1^{i_1}\mathrm{}v_n^{i_n}(gv_1^{j_1}\mathrm{}v_n^{j_n})\overline{gh}.$$
The factor $`v_1^{i_1}\mathrm{}v_n^{i_n}(gv_1^{j_1}\mathrm{}v_n^{j_n})`$ may now be rearranged using (2.3) repeatedly until it is in the standard form for elements of $`A`$ discussed above. The result will be of the form
$$rs=rs+\mu _1(r,s)t+\mu _2(r,s)t^2+\mathrm{}+\mu _p(r,s)t^p,$$
where $`rs`$ is identified with the product of $`r`$ and $`s`$ in $`S(V)\mathrm{\#}_\alpha G`$, via the additive isomorphism $`AS(V)\mathrm{\#}_\alpha G[t]`$, and $`\mu _1(r,s),\mathrm{},\mu _p(r,s)`$ are elements of $`A`$ also identified with elements of $`S(V)\mathrm{\#}_\alpha G`$. This is a finite process as each time (2.3) is applied, the degree drops. As the multiplication $``$ in the twisted graded Hecke algebra $`A`$ is bilinear and associative, the maps $`\mu _i`$ are bilinear and thus $`A`$ is a deformation of $`S(V)\mathrm{\#}_\alpha G`$ over $`[t]`$. (Alternatively, the existence of the obvious algebra isomorphism $`A/tAS(V)\mathrm{\#}_\alpha G`$ implies that $`A`$ is a deformation of $`S(V)\mathrm{\#}_\alpha G`$ over $`[t]`$.) The conditions on the $`\mu _i`$ stated in the theorem are consequences of the relations (2.3), by induction on the degree $`_{k=1}^n(i_k+j_k)`$ of a product $`v_1^{i_1}\mathrm{}v_n^{i_n}v_1^{j_1}\mathrm{}v_n^{j_n}`$.
Conversely, suppose that $`A`$ is a deformation of $`S(V)\mathrm{\#}_\alpha G`$ over $`[t]`$ satisfying the given conditions. By definition, $`AS(V)\mathrm{\#}_\alpha G[t]`$ as a vector space over $`[t]`$. Define a $`[t]`$-linear map $`\varphi :T(V)\mathrm{\#}_\alpha G[t]A`$ by
$$\varphi (v_{i_1}\mathrm{}v_{i_m}\overline{g})=v_{i_1}\mathrm{}v_{i_m}\overline{g}$$
for all words $`v_{i_1}\mathrm{}v_{i_m}`$ in $`T(V)`$ and $`gG`$. It may be checked that $`\varphi `$ is an algebra homomorphism since $`T(V)`$ is free on $`v_1,\mathrm{},v_n`$, and since $`\mu _i(^\alpha G,^\alpha G)=\mu _i(^\alpha G,V)=\mu _i(V,^\alpha G)=0`$ for all $`i1`$ by the degree condition on $`\mu _i`$. We will show, by induction on degree, that $`\varphi `$ is surjective: First note that $`\varphi (\overline{g})=\overline{g}`$ and $`\varphi (v\overline{g})=v\overline{g}`$ for all $`vV`$, $`gG`$. Now we would like to show that an arbitrary basis monomial $`v_{i_1}\mathrm{}v_{i_m}\overline{g}`$ ($`i_1\mathrm{}i_m`$) of $`A`$ is in $`\mathrm{im}(\varphi )`$. Assume $`v_{i_2}\mathrm{}v_{i_m}\overline{g}=\varphi (X)`$ for some element $`X`$ of $`T(V)\mathrm{\#}_\alpha G[t]`$. Then
$`\varphi (v_{i_1}X)`$ $`=`$ $`\varphi (v_{i_1})\varphi (X)`$
$`=`$ $`v_{i_1}v_{i_2}\mathrm{}v_{i_m}\overline{g}`$
$`=`$ $`v_{i_1}\mathrm{}v_{i_m}\overline{g}+\mu _1(v_{i_1},v_{i_2}\mathrm{}v_{i_m}\overline{g})t+\mu _2(v_{i_1},v_{i_2}\mathrm{}v_{i_m}\overline{g})t^2+\mathrm{}`$
By induction on $`m`$, as $`\mu _i`$ is a map of degree $`2i`$, each $`\mu _j(v_{i_1},v_{i_2}\mathrm{}v_{i_m}\overline{g})`$ is in $`\mathrm{im}(\varphi )`$. This implies $`v_{i_1}\mathrm{}v_{i_m}\overline{g}\mathrm{im}(\varphi )`$, and thus $`\varphi `$ is surjective.
It remains to determine the kernel of $`\varphi `$. Letting $`v,wV`$, we find that
$`\varphi (vw)`$ $`=`$ $`vw=vw+\mu _1(v,w)t`$
$`\varphi (wv)`$ $`=`$ $`wv=wv+\mu _1(w,v)t,`$
since $`\mathrm{deg}\mu _i=2i`$. As $`vw=wv`$ in $`S(V)`$, we may subtract to obtain $`\varphi (vwwv)=(\mu _1(v,w)\mu _1(w,v))t`$. Since $`\mathrm{deg}\mu _1=2`$ and $`\varphi (\overline{g})=\overline{g}`$ for all $`gG`$, this implies that
(3.3)
$$vwwv(\mu _1(v,w)\mu _1(w,v))t$$
is in the kernel of $`\varphi `$ for all $`v,wV`$. It also follows, as $`\mathrm{deg}\mu _1=2`$, that
$$\mu _1(v,w)\mu _1(w,v)=\underset{gG}{}a_g(v,w)\overline{g}$$
for some functions $`a_g:V\times V`$. By definition, $`a_g`$ is bilinear and skew-symmetric for each $`gG`$. Let $`I[t]`$ be the ideal of $`T(V)\mathrm{\#}_\alpha G[t]`$ generated by all such expressions (3.3), so that $`I[t]\mathrm{ker}\varphi `$. We will use a dimension count to show that $`I[t]=\mathrm{ker}\varphi `$: By the arguments of the previous section, $`T(V)\mathrm{\#}_\alpha G[t]/I[t]`$ is a quotient of $`S(V)\mathrm{\#}_\alpha G[t]`$, and so has dimension in each degree no greater than that of $`S(V)\mathrm{\#}_\alpha G[t]`$. Since $`\varphi `$ induces a map from $`T(V)\mathrm{\#}_\alpha G[t]/I[t]`$ onto $`S(V)\mathrm{\#}_\alpha G[t]`$, this forces $`I[t]=\mathrm{ker}\varphi `$ and thus $`A`$ is a twisted graded Hecke algebra. ∎
The proof of Theorem 3.2 shows that the skew-symmetric forms $`a_g`$ appearing in the relations (2.3) of a twisted graded Hecke algebra arise as coefficients in the skew-symmetrization of a Hochschild two-cocycle $`\mu _1`$ on $`S(V)\mathrm{\#}_\alpha G`$. Hochschild cohomology thus provides an alternate approach to the study of twisted graded Hecke algebras. An advantage of this approach is that it puts the computation of the possible skew-symmetric forms $`a_g`$ into a larger context. A disadvantage is that, given a Hochschild two-cocycle $`\mu _1`$ on an algebra $`R`$, there is no general method for determining whether it lifts to a deformation of $`R`$, nor for finding the $`\mu _i`$ ($`i2`$) in case it does lift. In case $`\mathrm{deg}(\mu _1)=1`$, Remark 2.14 points to existence of deformations of $`S(V)\mathrm{\#}_\alpha G`$ defined by quotients of $`T(V)\mathrm{\#}_\alpha G[t]`$ analogous to twisted graded Hecke algebras. There are a few special cases where deformations of $`S(V)\mathrm{\#}_\alpha G`$ are known whose infinitesimals $`\mu _1`$ have arbitrarily high degree . These examples were discovered after examining the relevant Hochschild cohomology:
$$\mathrm{HH}^{\text{}}(S(V)\mathrm{\#}_\alpha G)\left(\underset{gG}{}^{\text{}\mathrm{codim}V^g}((V^g)^{})det(((V^g)^{})^{})S(V^g)\overline{g}\right)^G$$
additively, where the superscript $`G`$ denotes elements invariant under the action induced by conjugation in $`^\alpha G`$ and the action of $`G`$ on $`V`$, and $`det`$ denotes the top exterior power. This follows from , see (6.2) and the formula before (6.4) where $`S(V)`$ is replaced by $`[M]`$, or more directly from the techniques of , where $`S(V)`$ is replaced by a Weyl algebra.
We conclude by giving a broader context for the isomorphism found by Ram and Shepler between Lusztig’s and Drinfeld’s definitions of (untwisted) graded Hecke algebras for real reflection groups . Let $`WGL(V)`$ be a finite real reflection group. Specifically, $`W`$ is generated by simple reflections $`s_1,\mathrm{},s_n`$ corresponding to simple roots $`\alpha _1,\mathrm{},\alpha _nV`$ for a root system $`R`$ in $`V`$. If $`\alpha R`$ is any root, we write $`s_\alpha `$ for the reflection corresponding to $`\alpha `$, so that
$$s_\alpha v=vv,\alpha \stackrel{ˇ}{}\alpha $$
where $`\alpha \stackrel{ˇ}{}=2\alpha /\alpha ,\alpha `$, and $`,`$ is the inner product. For each $`\alpha R`$, choose $`k_\alpha `$ in such a way that $`k_{g\alpha }=k_\alpha `$ for all $`gW`$, so that the number of independent parameters $`k_\alpha `$ is simply the number of different lengths of roots. Lusztig’s graded version of an affine Hecke algebra, as defined over $`[t]`$, is the quotient $`A^{}`$ of $`T(V)\mathrm{\#}W[t]`$ by the ideal generated by all
$$[v,w]\text{ and }\overline{s}_iv(s_iv)\overline{s}_i+k_{\alpha _i}v,\alpha _i\stackrel{ˇ}{}t$$
where $`v,wV`$ and $`i=1,\mathrm{},n`$.
Ram and Shepler defined the following graded Hecke algebra and showed that it is isomorphic to Lusztig’s algebra $`A^{}`$ given above (and their isomorphism extends to one over $`[t]`$): For each $`gW`$, let
$$a_g(v,w)=\frac{1}{4}\underset{\stackrel{\alpha ,\beta >0}{g=s_\alpha s_\beta }}{}k_\alpha k_\beta \left(v,\beta \stackrel{ˇ}{}w,\alpha \stackrel{ˇ}{}v,\alpha \stackrel{ˇ}{}w,\beta \stackrel{ˇ}{}\right).$$
In order to make the connection to Lusztig’s version, replace $`t`$ by $`t^2`$ in (2.3) and take the quotient $`A`$ of $`T(V)\mathrm{\#}W[t]`$ by the ideal generated by all
$$[v,w]\underset{gW}{}a_g(v,w)t^2\overline{g}.$$
The resulting quotient algebra $`A`$ is in fact a graded Hecke algebra defined over $`[t^2]`$, however the isomorphism with Lusztig’s algebra $`A^{}`$ is not defined over $`[t^2]`$. This isomorphism $`\mathrm{\Phi }_t:A^{}A`$ is given by
$$\mathrm{\Phi }_t(v)=vtv,h\text{ and }\mathrm{\Phi }_t(\overline{g})=\overline{g}$$
for all $`vV,gW`$, where $`h={\displaystyle \frac{1}{2}}{\displaystyle \underset{\alpha >0}{}}k_\alpha \alpha \stackrel{ˇ}{}\overline{s}_\alpha `$, and $`v,h={\displaystyle \frac{1}{2}}{\displaystyle \underset{\alpha >0}{}}k_\alpha v,\alpha \stackrel{ˇ}{}\overline{s}_\alpha `$. This may be extended uniquely to an algebra isomorphism by the calculations of Ram and Shepler. The reason for replacing $`t`$ by $`t^2`$ in the definition of $`A`$ is that it is necessary in order to extend their isomorphism to one defined over $`[t]`$. The only calculation that changes significantly in this context occurs in the proof that
$$[\mathrm{\Phi }_t(v),\mathrm{\Phi }_t(w)]=[v,w]+t^2[v,h,w,h]t[v,w,h]+t[w,v,h]=0.$$
This is true as $`[v,w,h]=[w,v,h]`$ and $`[v,h,w,h]=_{gW}a_g(v,w)\overline{g}`$ by computations in the proof of \[20, Thm. 3.5\].
In particular, the isomorphism $`\mathrm{\Phi }_t`$ is an equivalence of deformations over $`[t]`$. This implies that the Hochschild two-cocycle corresponding to the coefficients of $`t`$ in the deformation $`A^{}`$ of $`S(V)\mathrm{\#}W`$ is a coboundary, however there is a Hochschild two-cocycle arising from the coefficients of $`t^2`$ in products in $`A^{}`$ that will be cohomologous to that for $`A`$. (Again see or for details from the general theory of algebraic deformations.) |
warning/0506/math0506009.html | ar5iv | text | # Beyond Rouquier partitions
## 1. Introduction
In the course of investigating the truth of Broué’s Abelian defect group conjecture for symmetric groups, Rouquier singled out a special class of blocks of symmetric groups which he believed to have good properties. These blocks, and their corresponding blocks of the Iwahori-Hecke algebras (of type A) and the $`q`$-Schur algebras, are now known as Rouquier blocks, and are well understood, by the works of several authors (see, for example, , , , , ). In particular, there exist closed formulas, in terms of Littlewood-Richardson coefficients, for the decomposition numbers of these blocks when they are of ‘Abelian defect’.
These formulas coincide with those for the $`v`$-decomposition numbers arising from the canonical basis of the Fock space representation of $`U_v(\widehat{𝔰𝔩}_e)`$ obtained by Leclerc and Miyachi upon evaluation at $`v=1`$. In fact, such formulas for the $`v`$-decomposition numbers hold not only for the Rouquier partitions — the partitions indexing the simple modules of Rouquier blocks. As shown by Chuang and the author , whenever $`\kappa `$ is an $`e`$-core partition of Rouquier type — one having an abacus display in which the number of beads on each runner is non-decreasing as we go from left to right — then the formulas hold for the canonical basis element labelled by any partition lying in the class $`𝒫_\kappa `$ of $`e`$-regular partitions having $`e`$-core $`\kappa `$ and ‘locally small’ $`e`$-quotients.
However, this larger class of partitions is unsatisfactory for the following reasons:
* It does not contain any $`e`$-singular partition.
* It is not closed under the Mullineux involution $`\mu \mu ^{}`$; i.e. $`\mu 𝒫_\kappa `$ does not imply that $`\mu ^{}𝒫_\kappa ^{}`$ (even though, as shown in Proposition 3.7, that if the formula holds for the canonical basis element labelled by $`\mu `$, then it also holds for that labelled by $`\mu ^{}`$).
* $`\kappa `$ is not arbitrary.
We address all the above shortcomings in the first part of this paper. Our class $`𝒫_\kappa ^{}`$, consisting of partitions having $`e`$-core $`\kappa `$ and ‘locally small’ $`e`$-quotients (though our definition of ‘locally small’ is different from that for $`𝒫_\kappa `$), is defined for all $`e`$-core partitions $`\kappa `$, is closed under the Mullineux involution, and consists of $`e`$-singular partitions as well. Furthermore, if $`\kappa `$ is of Rouquier type, then $`𝒫_\kappa `$ is a proper subclass of $`𝒫_\kappa ^{}`$.
In the second part of this paper, we show that, upon evaluation at $`v=1`$, our formula for the canonical basis element labelled by a partition $`\mu 𝒫_\kappa ^{}`$ gives the corresponding decomposition numbers of the $`q`$-Schur algebras in characteristic $`l`$ (where $`q`$ is a primitive $`e`$-th root of unity if $`le`$, and $`q=1`$ otherwise) whenever $`l`$ is larger than the sizes of each constituent of the $`e`$-quotient of $`\mu `$. This thus includes cases where the blocks, in which the simple modules labelled by such partitions lie, are not of ‘Abelian defect’.
The paper is organized as follows: we begin in section 2 with a short account of the background theory which we require. In section 3, we introduce the notations used in this paper, and state the main theorem. Sections 4 and 5 are devoted to the proof of the main theorem.
## 2. Preliminaries
In this section, we give a brief account of the background theory we require, and introduce some notations and conventions which will be used in this paper. From now on, we fix an integer $`e2`$.
### 2.1. Partitions
Let $`𝒫_n`$ be the set of partitions of a natural number $`n`$, and denote the lexicographical ordering and dominance ordering on $`𝒫_n`$ by $``$ and $`\mathrm{}`$ respectively. Write $`𝒫=_n𝒫_n`$ for the set of all partitions.
Let $`\lambda =(\lambda _1,\lambda _2,\mathrm{},\lambda _s)`$ be a partition, with $`\lambda _1\lambda _2\mathrm{}\lambda _s>0`$. We write $`|\lambda |=_{i=1}^s\lambda _i`$ and $`l(\lambda )=s`$. The conjugate partition of $`\lambda `$ will be denoted as $`\lambda ^{}`$.
The James $`e`$-abacus (see, for example, \[4, Section 2.7\]) has $`e`$ runners, numbered $`0`$ to $`e1`$ from left to right, and its positions are numbered from left to right and down the rows, starting from $`0`$. Any partition may be displayed on such an abacus; this display is only unique when the number of beads used is fixed. In this paper, if $`\lambda `$ is a partition with $`e`$-core $`\kappa `$, we shall always display $`\lambda `$ on an $`e`$-abacus with $`l(\kappa )+Ne`$ beads for sufficiently large $`N`$, and denote its $`e`$-quotient by $`(\lambda ^0,\lambda ^1,\mathrm{},\lambda ^{n1})`$, where $`\lambda ^i`$ is the partition read off from runner $`i`$ of the abacus display of $`\lambda `$. Define the $`e`$-weight of a bead in the abacus display of $`\lambda `$ to be the number of vacant positions above it in the same runner. The $`e`$-weight of a runner is the sum of $`e`$-weights of all the beads in the runner, and the $`e`$-weight of $`\lambda `$ is the sum of $`e`$-weights of all the beads in the abacus display of $`\lambda `$.
If $`s+t=r`$ and $`\alpha 𝒫_s`$ and $`\beta 𝒫_t`$, let $`c_{\alpha \beta }^\lambda `$ denote the multiplicity of the ordinary irreducible character $`\chi ^\lambda `$ of the symmetric group $`𝔖_r`$ in the induced character $`\mathrm{Ind}_{𝔖_s\times 𝔖_t}^{𝔖_r}(\chi ^\alpha \chi ^\beta )`$. We refer the reader to \[10, I.9\] for a combinatorial description of $`c_{\alpha \beta }^\lambda `$, which is known as a Littlewood-Richardson coefficient. By convention, we define $`c_{\alpha \beta }^\lambda `$ to be $`0`$ when $`|\lambda ||\alpha |+|\beta |`$.
### 2.2. The Fock space representation
The algebra $`U_v(\widehat{𝔰𝔩}_e)`$ is the associative algebra over $`(v)`$ with generators $`e_r`$, $`f_r`$, $`k_r`$, $`k_r^1`$ $`(0re1)`$, $`d`$, $`d^1`$ subject to some relations (see, for example, \[7, §4\]). An important $`U_v(\widehat{𝔰𝔩}_e)`$-module is the Fock space representation $``$, which as a $`(v)`$-vector space has $`𝒫`$ as a basis. For our purposes, an explicit description of the actions of $`e_r`$ and $`f_r`$ will suffice.
Display a partition $`\lambda `$ on the $`e`$-abacus with $`t`$ beads, where $`tl(\lambda )`$ and $`e(r+t)`$. Let $`i`$ be the residue class of $`(r+t)`$ modulo $`e`$. Suppose there is a bead on runner $`i1`$ whose succeeding position on runner $`i`$ is vacant; let $`\mu `$ be the partition obtained when this bead is moved to its succeeding position. Let $`N_>(\lambda ,\mu )`$ (resp. $`N_<(\lambda ,\mu )`$) be the number of beads on runner $`i1`$ below (resp. above) the bead moved to obtained $`\mu `$ minus the number of beads on runner $`i`$ below (resp. above) the vacant position that becomes occupied in obtaining $`\mu `$. We have
$`f_r(s(\lambda ))`$ $`={\displaystyle \underset{\mu }{}}v^{N_>(\lambda ,\mu )}s(\mu );`$
$`e_r(s(\mu ))`$ $`={\displaystyle \underset{\lambda }{}}v^{N_<(\lambda ,\mu )}s(\lambda ),`$
where $`\mu `$ in the first sum runs over all partitions that can be obtained from $`\lambda `$ by moving a bead on runner $`i1`$ to its vacant succeeding position on runner $`i`$, while $`\lambda `$ in the second sum runs over all partitions that can be obtained from $`\mu `$ by moving a bead on runner $`i`$ to its vacant preceding position on runner $`i1`$.
In , Leclerc and Thibon introduced an involution $`x\overline{x}`$ on $``$, having the following properties (among others):
$$\overline{a(v)x}=a(v^1)\overline{x},\overline{e_r(x)}=e_r(\overline{x}),\overline{f_r(x)}=f_r(\overline{x})(a(v),x).$$
For $`k^+`$, we write $`[k]!=_{i=1}^k\frac{v^iv^i}{vv^1}`$, and let $`e_r^{(k)}=e_r^k/[k]!`$ and $`f_r^{(k)}=f_r^k/[k]!`$. Note that $`\overline{e_r^{(k)}(x)}=e_r^{(k)}(\overline{x})`$ and $`\overline{f_r^{(k)}(x)}=f_r^{(k)}(\overline{x})`$. There is a distinguished basis $`\{G(\sigma )\sigma 𝒫\}`$ of $``$, called the canonical basis, having the following characterization (\[9, Theorem 4.1\]):
1. $`G(\sigma )\sigma vL`$, where $`L`$ is the $`[v]`$-lattice in $``$ generated by $`𝒫`$.
2. $`\overline{G(\sigma )}=G(\sigma )`$.
The $`v`$-decomposition number $`d_{\lambda \sigma }(v)`$ is defined the coefficient of $`\lambda `$ in $`G(\sigma )`$; these numbers enjoy the following property:
###### Theorem 2.1 (\[7, Theorem 9, Proposition 11, Corollary 14\]).
We have
$`d_{\sigma \sigma }(v)`$ $`=1,`$
$`d_{\lambda \sigma }(v)`$ $`v_0[v]\text{ for all }\lambda \sigma .`$
Furthermore, $`d_{\lambda \sigma }(v)0`$ only if $`\sigma \mathrm{}\lambda `$ and $`\lambda `$ and $`\sigma `$ have the same $`e`$-core.
There is an involution $`\mu \mu ^{}`$, known as the Mullineux involution, on the set of $`e`$-regular partitions (see, for example, \[11, page 120\]). The $`v`$-decomposition numbers $`d_{\lambda \mu }(v)`$ and $`d_{\lambda ^{}\mu ^{}}(v)`$ are related in the following way:
###### Theorem 2.2 (\[6, Theorem 7.2\]).
We have $`d_{\lambda ^{}\mu ^{}}(v)=v^wd_{\lambda \mu }(v^1)`$ where $`w`$ is the $`e`$-weight of $`\mu `$.
The partition $`\mu _{}^{}{}_{}{}^{}`$ can be characterized in the following way using the $`v`$-decomposition numbers.
###### Theorem 2.3 (\[6, Corollary 7.7\]).
We have $`d_{\mu _{}^{}{}_{}{}^{}\mu }(v)=v^w`$, where $`w`$ is the $`e`$-weight of $`\mu `$, and $`\mathrm{deg}d_{\lambda \mu }(v)<w`$ for all $`\lambda \mu _{}^{}{}_{}{}^{}`$.
### 2.3. $`q`$-Schur algebras
Let $`𝔽`$ be a field of characteristic $`l`$, and let $`q𝔽^{}`$ be such that $`e`$ is the least integer such that $`1+q+\mathrm{}+q^{e1}=0`$. The $`q`$-Schur algebra $`𝒮_{𝔽,q}(n)=𝒮_{𝔽,q}(n,n)`$ over $`𝔽`$ has a distinguished class $`\{\mathrm{\Delta }^\mu \mu 𝒫_n\}`$ of right modules called Weyl modules. Each $`\mathrm{\Delta }^\mu `$ has a simple head $`L^\mu `$, and the set $`\{L^\mu \mu 𝒫_n\}`$ is a complete set of non-isomorphic simple modules of $`𝒮_{𝔽,q}(n)`$. The projective cover $`P^\mu `$ of $`L^\mu `$ (or of $`\mathrm{\Delta }^\mu `$) has a filtration in which each factor is isomorphic to a Weyl module; the multiplicity of $`\mathrm{\Delta }^\lambda `$ in such a filtration is well-defined, and is equal to the multiplicity of $`L^\mu `$ as a composition factor of $`\mathrm{\Delta }^\lambda `$. We denote this multiplicity as $`d_{\lambda \mu }^l`$, which is a decomposition number of $`𝒮_{𝔽,q}(n)`$.
The decomposition numbers in characteristic $`l`$ and those in characteristic $`0`$ are related by an adjustment matrix $`A_l`$: let $`D_l=(d_{\lambda \mu }^l)_{\lambda ,\mu 𝒫_n}`$, then $`D_l=D_0A_l`$. Furthermore, the matrix $`A_l`$ is lower unitriangular with nonnegative entries when the partitions indexing its rows and columns are ordered by a total order extending the dominance order on $`𝒫_n`$ (such as the lexicographic order). As a consequence, we have
###### Lemma 2.4.
$`d_{\lambda \mu }^ld_{\lambda \mu }^0`$.
The link between the $`v`$-decomposition numbers of the Fock space and the decomposition numbers of $`q`$-Schur algebras is established by Varagnolo and Vasserot:
###### Theorem 2.5 ().
$`d_{\lambda \mu }(1)=d_{\lambda \mu }^0`$.
Thus the canonical basis vector $`G(\mu )`$ of $``$ corresponds to the projective cover $`P^\mu `$ of $`q`$-Schur algebras, while the standard basis element $`\lambda `$ of $``$ corresponds to the Weyl module $`\mathrm{\Delta }^\lambda `$. Under this correspondence, the action of $`e_r,f_rU_v(\widehat{𝔰𝔩}_e)`$ on $``$ corresponds to $`r`$-restriction and $`r`$-induction (see, for example, \[11, 6.4\]) of modules of $`q`$-Schur algebras.
### 2.4. Jantzen order
Let $`\lambda `$ be a partition, and consider its abacus display, with $`k`$ beads say. Suppose in moving a bead, say at position $`a`$, up its runner to some vacant position, say $`aie`$, we obtain (the abacus display of) a partition $`\mu `$. Write $`l_{\lambda \mu }`$ for the number of occupied positions between $`a`$ and $`aie`$, and let $`h_{\lambda \mu }=i`$. Also, write $`\lambda \stackrel{𝜇}{}\tau `$ if the abacus display of $`\mu `$ with $`k`$ beads is also obtained from that of $`\tau `$ by moving a bead at position $`b`$ to $`bie`$, and $`a<b`$. Thus if $`\lambda \stackrel{𝜇}{}\tau `$, then the abacus display of $`\tau `$ with $`k`$ beads may be obtained from $`\lambda `$ in two steps: first move the bead at position $`a`$ to position $`aie`$ (which yields the abacus display of $`\mu `$), and then move the bead at position $`bie`$ to position $`b`$.
The Jantzen sum formula (see, for example, \[11, 5.32\]) provides an upper bound for the decomposition numbers, and may be stated as follows:
###### Theorem 2.6.
Let $`\lambda `$ and $`\mu `$ be distinct partitions, and let
$$J_{\lambda \mu }=(1)^{l_{\lambda \sigma }+l_{\tau \sigma }+1}(1+\nu _l(h_{\lambda \sigma }))d_{\tau \mu }^l,$$
where the sum runs through all $`\tau `$ and $`\sigma `$ such that $`\lambda \stackrel{𝜎}{}\tau `$, and where $`\nu _l`$ denotes the standard $`l`$-valuation if $`l>0`$ and $`\nu _0(x)=0`$ for all $`x`$. Then $`d_{\lambda \mu }^lJ_{\lambda \mu }`$, and $`d_{\lambda \mu }^l=0`$ if and only if $`J_{\lambda \mu }=0`$.
We write $`\lambda \tau `$ if there exists some $`\mu `$ such that $`\lambda \stackrel{𝜇}{}\tau `$. We further write $`\lambda <_J\sigma `$ (or $`\sigma >_J\lambda `$) if there exist partitions $`\tau _0,\mathrm{},\tau _r`$ such that $`\tau _0=\lambda `$, $`\tau _r=\sigma `$ and $`\tau _{i1}\tau _i`$ for all $`i=1,\mathrm{},r`$. It is clear that $`_J`$ (which means $`>_J`$ or $`=`$) defines a partial order on the set $`𝒫`$ of all partitions, and that if $`\lambda _J\mu `$, then $`\lambda `$ and $`\mu `$ have the same $`e`$-core and $`e`$-weight. Furthermore, the dominance order extends $`_J`$.
The Jantzen sum formula motivates the study of $`_J`$, as it provides this easy consequence.
###### Lemma 2.7.
Suppose $`d_{\lambda \mu }^l0`$. Then $`\lambda _J\mu `$.
Thus, if $`\lambda _J\mu `$, then $`d_{\lambda \mu }^l=0`$. Unfortunately, it is not easy to determine by inspection if $`\lambda _J\mu `$. To this end, we introduce another partial order $`_p`$.
Let $`\lambda `$ be a partition, and suppose that when displayed on an $`e`$-abacus with $`N`$ beads, the beads having positive $`e`$-weights are at positions $`a_1,a_2,\mathrm{},a_r`$, with say the bead at position $`a_i`$ having $`e`$-weight $`w_i`$. The induced $`e`$-sequence of $`\lambda `$, denoted $`s(\lambda )=s(\lambda )_N`$, is defined as
$$\underset{i=1}{\overset{r}{}}(a_i,a_ie,\mathrm{},a_i(w_i1)e),$$
where $`(b_1,b_2,\mathrm{},b_s)(c_1,c_2,\mathrm{},c_t)`$ denotes the (unique) weakly decreasing sequence obtained by rearranging the terms in the sequence $`(b_1,b_2,\mathrm{},b_s,c_1,c_2,\mathrm{},c_t)`$.
###### Example.
Let $`\lambda =(6,6,5,4)`$. Its 3-abacus display with $`6`$ beads is
$$\begin{array}{ccc}& & \\ & & \\ & & \\ & & \end{array}.$$
Thus $`s(\lambda )_6=(11,8)(10,7)(8,5)(6)=(11,10,8,8,7,6,5)`$.
We record some basic properties of induced $`e`$-sequences.
###### Lemma 2.8.
Let $`\lambda `$ be a partition. We have
1. $`s(\lambda )_0^w`$ where $`w`$ is the $`e`$-weight of $`\lambda `$;
2. if $`M_j`$ and $`m_j`$ denote respectively the largest occupied position and least vacant position on runner $`j`$ of the abacus display of $`\lambda `$, and $`s(\lambda )=(l_1,l_2,\mathrm{},l_w)`$, then $`M_j=\mathrm{max}_i\{l_il_ij(mode)\}`$ and $`m_j=\mathrm{min}_i\{l_il_ij(mode)\}e`$ whenever runner $`j`$ has positive $`e`$-weight.
3. if $`\mu `$ is the partition obtained by moving a bead in the abacus display of $`\lambda `$ (with $`N`$ beads) from position $`x`$ to position $`xie`$, then $`s(\lambda )_N=s(\mu )_N(x,xe,\mathrm{},x(i1)e)`$.
The partial order $`_p`$ on the set of partitions is defined as: $`\lambda _p\mu `$ if and only if $`\lambda `$ and $`\mu `$ have the same $`e`$-core and the same $`e`$-weight, say $`w`$, and $`s(\lambda )_Ns(\mu )_N`$ (for sufficiently large $`N`$) in the standard product order on $`_0^w`$.
###### Lemma 2.9.
If $`\lambda _J\mu `$, then $`\lambda _p\mu `$.
###### Proof.
It suffices to prove that $`\lambda <_p\mu `$ when $`\lambda \stackrel{𝜏}{}\mu `$. But if $`\lambda \stackrel{𝜏}{}\mu `$, then there exists integers $`i,x,y`$ with $`i1`$ and $`0x<y`$ such that
$`s(\lambda )`$ $`=s(\tau )(x,xe,\mathrm{},x(i1)e),`$
$`s(\mu )`$ $`=s(\tau )(y,ye,\mathrm{},y(i1)e),`$
by Lemma 2.8(3). Thus, $`s(\lambda )<s(\mu )`$ in the product order. ∎
## 3. Setup
In this section, we set up the notations which shall henceforth be used in this paper and state the main theorem. We devote the next two sections to the proof of the main theorem.
We shall consider partitions with a fixed $`e`$-core, say $`\kappa `$. We display all these partitions on an abacus with $`l(\kappa )+Ne`$ beads, where $`N`$ is ‘sufficiently large’. For each $`0i<e`$, let $`n_i`$ be the number of beads on runner $`i`$ of such an abacus. Define $``$ ($`=_\kappa )`$ as follows: $`ij`$ if and only if any of the following holds:
* $`n_i<n_j`$, or
* $`n_i=n_j`$ and $`ij`$.
We note that $``$ is independent of $`N`$, and is a partial order on $`\{0,1,\mathrm{},e1\}`$. Furthermore, as the abacus has $`l(\kappa )+Ne`$ beads, we have $`n_0n_i`$ for all $`0<i<e`$, so that $`0`$ is the minimal element with respect to $``$; we denote the maximal element by $`M`$.
For $`iM`$ (which means $`iM`$ and $`iM`$), define $`i^+`$ as the least (with respect to $``$) $`k`$ such that $`ik`$, and for $`0j`$, define $`j^{}`$ as the largest (with respect to $``$) $`k`$ such that $`kj`$.
For $`0<i<e`$, define
$$d_i=\{\begin{array}{cc}n_in_i^{},\hfill & \text{if }i^{}<i;\hfill \\ n_in_i^{}1,\hfill & \text{if }i^{}>i.\hfill \end{array}$$
Let $`\pi `$ be defined recursively as follows: $`\pi (0)=0`$ and $`\pi (i^+)=\pi (i)+1`$ for all $`0iM`$. Thus, $`\pi `$ is permutation on the set $`\{0,1,\mathrm{},e1\}`$, and $`\pi (i)\pi (j)`$ if and only if $`ij`$.
When there is a need to indicate the $`e`$-core $`\kappa `$ for which $``$, $`M`$, $`i^\pm `$, $`d_i`$ and $`\pi `$ are defined, we will include $`\kappa `$ as a subscript or superscript; we thus have $`_\kappa `$, $`M_\kappa `$, $`i^{\pm _\kappa }`$, $`d_i^\kappa `$ and $`\pi _\kappa `$.
The following Lemma is clear.
###### Lemma 3.1.
Let $`\kappa `$ be an $`e`$-core partition. The following statements are equivalent:
1. $`\kappa `$ has an abacus display in which the number of beads in each runner is non-decreasing as we go from left to right.
2. $`ij`$ if and only if $`ij`$.
3. $`i^+=i+1`$ for all $`0ie2`$.
4. $`i^{}=i1`$ for all $`0<i<e`$.
5. $`\pi (i)=i`$ for all $`0i<e`$.
When any of these statements hold, we say $`\kappa `$ is of Rouquier type.
We note that $`_\kappa `$, $`M_\kappa `$, $`i^{\pm _\kappa }`$, $`d_i^\kappa `$ and $`\pi _\kappa `$ and $`_\kappa ^{}`$, $`M_\kappa ^{}`$, $`i^{\pm _\kappa ^{}}`$, $`d_i^\kappa ^{}`$ and $`\pi _\kappa ^{}`$ are related as follows:
###### Lemma 3.2.
Let $`\kappa `$ be an $`e`$-core partition. Let $`\mathrm{\Phi }`$ be the involution on $`\{0,1,\mathrm{},e1\}`$ defined by $`\mathrm{\Phi }(i)M_\kappa i(mode)`$. Then
1. $`i_\kappa j`$ if and only if $`\mathrm{\Phi }(j)_\kappa ^{}\mathrm{\Phi }(i)`$;
2. $`M_\kappa =M_\kappa ^{}`$;
3. $`\mathrm{\Phi }(i^{\pm _\kappa ^{}})=\mathrm{\Phi }(i)^_\kappa `$, $`\mathrm{\Phi }(i^{\pm _\kappa })=\mathrm{\Phi }(i)^_\kappa ^{}`$;
4. $`d_i^\kappa ^{}=d_{\mathrm{\Phi }(i)^{+_\kappa }}^\kappa `$;
5. $`\pi _\kappa ^{}(i)=e1\pi _\kappa (\mathrm{\Phi }(j))`$.
###### Proof.
This follows from the fact that an abacus display of $`\kappa ^{}`$ may be obtained from that of $`\kappa `$ by rotating it through an angle of $`\pi `$ and reading the occupied positions as empty and the empty positions as occupied. ∎
Given two partitions $`\lambda `$ and $`\mu `$ with the same $`e`$-core, say $`\kappa `$, and $`e`$-quotients $`(\lambda ^0,\mathrm{},\lambda ^{e1})`$ and $`(\mu ^0,\mathrm{},\mu ^{e1})`$ respectively, define
$`\delta (\lambda ,\mu )`$ $`={\displaystyle \underset{j=1}{\overset{e1}{}}}\pi (j)(|\mu ^j||\lambda ^j|);`$
$`C_{\lambda \mu }`$ $`={\displaystyle \underset{j=0}{\overset{e1}{}}c_{\alpha ^j\beta ^j}^{\lambda ^j}c_{(\beta ^j^{})^{}\alpha ^j}^{\mu ^j}},`$
where the second sum runs over all partitions $`\alpha ^0,\mathrm{},\alpha ^{e1}`$ and $`\beta ^0,\mathrm{},\beta ^{M1},\beta ^{M+1},\mathrm{},\beta ^{e1}`$, and where $`\beta ^M=\mathrm{}=\beta ^0^{}`$. For convenience, we define $`C_{\lambda \mu }`$ and $`\delta (\lambda ,\mu )`$ to be $`0`$ if $`\lambda `$ and $`\mu `$ do not have the same $`e`$-core.
We note the following:
###### Lemma 3.3.
Let $`\mu `$ and $`\lambda `$ be partitions with $`e`$-core $`\kappa `$, and $`e`$-quotients $`(\mu ^0,\mu ^1,\mathrm{},\mu ^{e1})`$ and $`(\lambda ^0,\lambda ^1,\mathrm{},\lambda ^{e1})`$ respectively, such that for some partitions $`\alpha ^0,\mathrm{},\alpha ^{e1}`$ and $`\beta ^0,\mathrm{},\beta ^{M1},\beta ^{M+1},\mathrm{},\beta ^{e1}`$,
$$\underset{i=0}{\overset{e1}{}}\left(c_{\alpha _i\beta _i}^{\lambda ^i}c_{(\beta ^i^{})^{}\alpha ^i}^{\mu ^i}\right)0,$$
where $`\beta ^0^{}=\beta ^M=\mathrm{}`$. Then
1. $`\delta (\lambda ,\mu )=_{0iM}|\beta ^i|`$;
2. whenever $`x`$ is the least vacant position on runner $`j`$, $`y`$ is the largest occupied position on runner $`i`$, and $`ij`$, we have
$$xy>\{\begin{array}{cc}(d_j+1|\mu ^j^{}||\mu ^j||\mu ^{j^+}|)e,\hfill & \text{if }i^+=j;\hfill \\ (_{irj}d_r+1|\mu ^i||\mu ^{i^+}||\mu ^j||\mu ^{j^+}|)e,\hfill & \text{if }i^+j.\hfill \end{array}$$
###### Proof.
1. Since $`_{i=0}^{e1}\left(c_{\alpha _i\beta _i}^{\lambda ^i}c_{(\beta ^i^{})^{}\alpha ^i}^{\mu ^i}\right)0`$, we have, for all $`i`$, $`|\lambda ^i|=|\alpha _i|+|\beta _i|`$ and $`|\mu ^i|=|\beta ^i^{}|+|\alpha ^i|`$, so that $`|\lambda ^i||\mu ^i|=|\beta ^i||\beta ^i^{}|`$. Thus,
$`\delta (\lambda ,\mu )`$ $`={\displaystyle \underset{j=1}{\overset{e1}{}}}\pi (j)(|\mu ^j||\lambda ^j|)`$
$`={\displaystyle \underset{j=1}{\overset{e1}{}}}\pi (j)(|\beta ^j^{}||\beta ^j|)`$
$`={\displaystyle \underset{j=1}{\overset{e1}{}}}\pi (j)|\beta ^j^{}|{\displaystyle \underset{j=1}{\overset{e1}{}}}\pi (j)|\beta ^j|`$
$`={\displaystyle \underset{0jM}{}}(\pi (j)+1)|\beta ^j|{\displaystyle \underset{0jM}{}}\pi (j)|\beta ^j|`$
$`={\displaystyle \underset{0jM}{}}|\beta ^j|.`$
2. Let $`n_k`$ be the number of beads on runner $`k`$ in the abacus display of $`\kappa `$. Then $`x=(n_jl(\lambda ^j))e+j`$, and $`y=(n_i+\lambda _1^i1)e+i`$. Suppose $`ij`$. Then
$`xy`$ $`=(n_jn_i+1l(\lambda ^j)\lambda _1^i)e+(ji)`$
$`=\left({\displaystyle \underset{irj}{}}(n_rn_r^{})+1l(\lambda ^j)\lambda _1^i\right)e+(ji)`$
If $`i<j`$, then clearly
$`xy`$ $`({\displaystyle \underset{irj}{}}d_r+1l(\lambda ^j)\lambda _1^i)e+(ji)`$
$`>({\displaystyle \underset{irj}{}}d_r+1l(\lambda ^j)\lambda _1^i)e.`$
On the other hand, if $`i>j`$, then there exists $`irj`$ such that $`r^{}>r`$, so that $`n_rn_r^{}=d_r+1`$; thus
$`xy`$ $`({\displaystyle \underset{irj}{}}d_r+2l(\lambda ^j)\lambda _1^i)e+(ji)`$
$`>({\displaystyle \underset{irj}{}}d_r+1l(\lambda ^j)\lambda _1^i)e.`$
Part (2) now follows since
$$l(\lambda ^j)+\lambda _1^i|\lambda ^j|+|\lambda ^i|\{\begin{array}{cc}|\mu ^j^{}|+|\mu ^j|+|\mu ^{j^+}|,\hfill & \text{if }i^+=j;\hfill \\ |\mu ^i|+|\mu ^{i^+}|+|\mu ^j|+|\mu ^{j^+}|,\hfill & \text{if }i^+j.\hfill \end{array}$$
###### Definition 3.4.
Let $`𝒫_\kappa ^{}`$ be the collection of partitions $`\mu `$ having $`e`$-core $`\kappa `$ and $`e`$-quotient $`(\mu ^0,\mathrm{},\mu ^{e1})`$ such that
1. $`|\mu ^i^{}|+|\mu ^i|+|\mu ^{i^+}|d_i+1`$ for all $`i=1,\mathrm{},e1`$,
2. whenever $`i`$ and $`j`$ ($`ij`$) satisfy
1. $`|\mu ^j|+|\mu ^{j^+}|=d_j+1`$,
2. $`|\mu ^i^{}|+|\mu ^i|=d_i+1`$,
there exists $`k`$ such that $`ikj`$ and $`d_k>0`$.
Here, $`|\mu ^i^{}|`$ (resp. $`|\mu ^{i^+}|`$) is to be read as 0 when $`i^{}`$ (resp. $`i^+`$) is undefined.
###### Example.
Let $`e=5`$ and $`\kappa =(3,3)`$. Then $`(d_1,d_2,d_3,d_4)=(0,0,1,0)`$. Let $`\mu =(8,3,2,1,1,1)`$. Then $`\mu `$ has $`5`$-core $`\kappa `$ and $`5`$-quotient $`(\mathrm{},(1),\mathrm{},\mathrm{},(1))`$, so that $`\mu 𝒫_\kappa ^{}`$ (but $`\mu `$ is not an element of $`𝒫_\kappa `$ defined in ).
###### Remark.
1. As a partition having $`e`$-core $`\kappa `$ is Rouquier if and only if its $`e`$-weight is not more than $`\mathrm{min}\{d_i:1i<e\}+1`$, we see that $`𝒫_\kappa ^{}`$ includes all Rouquier partitions with $`e`$-core $`\kappa `$.
2. In , a collection $`𝒫_\kappa `$ of partitions with $`e`$-core $`\kappa `$ is defined when $`\kappa `$ is of Rouquier type. This is a subcollection of $`𝒫_\kappa ^{}`$.
The partitions in $`𝒫_\kappa ^{}`$ have the following nice property.
###### Lemma 3.5.
If $`\mu 𝒫_\kappa ^{}`$, and $`C_{\lambda \mu }0`$, then any occupied position on runner $`i`$ of the abacus display of $`\lambda `$ is less than any vacant position on runner $`j`$ as long as $`ij`$.
###### Proof.
This follows directly from the definition of $`𝒫_\kappa ^{}`$ and Lemma 3.3. ∎
We now state the main theorem of this paper.
###### Theorem 3.6.
Let $`\mu 𝒫_\kappa ^{}`$ with $`e`$-quotient $`(\mu ^0,\mathrm{},\mu ^{e1})`$ . Then
1. $`d_{\lambda \mu }(v)=C_{\lambda \mu }v^{\delta (\lambda ,\mu )}`$ for all $`\lambda 𝒫`$;
2. $`d_{\lambda \mu }^l=C_{\lambda \mu }`$ for all $`\lambda 𝒫`$ and $`l>\mathrm{max}_i(|\mu ^i|)`$.
We shall devote the next two sections to the proof of Theorem 3.6.
It is easy to describe the image under the Mullineux involuation of an $`e`$-regular partition for which Theorem 3.6(1) holds; furthermore, Theorem 3.6(1) also holds for the image:
###### Proposition 3.7.
Suppose that an $`e`$-regular partition $`\mu `$ has $`e`$-core $`\kappa `$ and $`e`$-quotient $`(\mu ^0,\mathrm{},\mu ^{e1})`$, and that $`d_{\lambda \mu }(v)=C_{\lambda \mu }v^{\delta (\lambda ,\mu )}`$ for all $`\lambda 𝒫`$. Then
1. $`\mu ^0=\mathrm{}`$;
2. $`\mu ^{}`$ has $`e`$-core $`\kappa ^{}`$ and $`e`$-quotient $`(\mathrm{},\mu ^{\mathrm{\Phi }(1)^{+_\kappa }},\mu ^{\mathrm{\Phi }(2)^{+_\kappa }},\mathrm{},\mu ^{\mathrm{\Phi }(e1)^{+_\kappa }})`$, where $`\mathrm{\Phi }`$ is the involution on $`\{0,1,\mathrm{},e1\}`$ defined by $`\mathrm{\Phi }(i)M_\kappa i(mode)`$;
3. $`d_{\tau \mu ^{}}(v)=C_{\tau \mu ^{}}v^{\delta (\tau ,\mu ^{})}`$ for all $`\tau 𝒫`$.
###### Proof.
By Theorem 2.3, $`\mu _{}^{}{}_{}{}^{}`$ is the unique partition such that $`d_{\mu _{}^{}{}_{}{}^{}\mu }(v)`$ has degree equals the $`e`$-weight of $`\mu `$. When $`_{j=0}^{e1}c_{\alpha ^j\beta ^j}^{\lambda ^j}c_{(\beta ^{j^_\kappa })^{}\alpha ^j}^{\mu ^j}0`$ with $`\beta ^0^{}=\beta ^{M_\kappa }=\mathrm{}`$, we have $`\mathrm{deg}(d_{\lambda \mu }(v))=\delta (\lambda ,\mu )=_{0jM_\kappa }|\beta ^j|`$ by Lemma 3.3(1). Thus $`\mathrm{deg}(d_{\lambda \mu }(v))=_{i=0}^{e1}|\mu ^i|`$ only if $`(\beta ^{j^_\kappa })^{}=\mu ^j`$ and $`\alpha ^j=\mathrm{}`$ for all $`0jM_\kappa `$. It follows that $`\mu ^0=\beta ^0^{}=\mathrm{}`$, giving part (1).
Furthermore, the $`e`$-quotient of $`\mu _{}^{}{}_{}{}^{}`$ equals
$$(\beta ^0,\beta ^1,\mathrm{},\beta ^{e1})=((\mu ^{0^{+_\kappa }})^{},(\mu ^{1^{+_\kappa }})^{},\mathrm{},(\mu ^{(e1)^{+_\kappa }})^{}),$$
where $`\mu ^{M^{+_\kappa }}`$ is to be read as $`\mathrm{}`$.
Now, when a partition has $`e`$-core $`\kappa `$, and $`e`$-quotient $`(\lambda ^0,\mathrm{},\lambda ^{e1})`$ say, then its conjugate has $`e`$-core $`\kappa ^{}`$ and $`e`$-quotient $`((\lambda ^{\mathrm{\Phi }(0)})^{},(\lambda ^{\mathrm{\Phi }(1)})^{},\mathrm{},(\lambda ^{\mathrm{\Phi }(e1)})^{})`$, where $`\mathrm{\Phi }`$ is the involution on $`\{0,1,\mathrm{},e1\}`$ defined by $`\mathrm{\Phi }(i)M_\kappa i(mode)`$. Thus the $`e`$-quotient of $`\mu ^{}`$ is
$$(\mu ^{\mathrm{\Phi }(0)^{+_\kappa }},\mu ^{\mathrm{\Phi }(1)^{+_\kappa }},\mu ^{\mathrm{\Phi }(2)^{+_\kappa }},\mathrm{},\mu ^{\mathrm{\Phi }(e1)^{+_\kappa }})=(\mathrm{},\mu ^{\mathrm{\Phi }(1)^{+_\kappa }},\mu ^{\mathrm{\Phi }(2)^{+_\kappa }},\mathrm{},\mu ^{\mathrm{\Phi }(e1)^{+_\kappa }}),$$
giving part (2).
Using Theorem 2.2, we need to show that $`C_{\lambda ^{}\mu ^{}}=C_{\lambda \mu }`$ and $`\delta (\lambda ^{},\mu ^{})=_{j=1}^{e1}|\mu _j|\delta (\lambda ,\mu )`$ for part (3). Since $`\mu ^0=\mathrm{}`$, we can simplify the expression of $`C_{\lambda \mu }`$, where $`\lambda `$ has $`e`$-core $`\kappa `$ and $`e`$-quotient $`(\lambda ^0,\mathrm{},\lambda ^{e1})`$, to
$$C_{\lambda \mu }=\left(\underset{\begin{array}{c}0ie1\\ i0,M_\kappa \end{array}}{}c_{\alpha ^i\beta ^i}^{\lambda ^i}\underset{\begin{array}{c}0je1\\ j0,\mathrm{\hspace{0.25em}0}^{+_\kappa },M_\kappa \end{array}}{}c_{(\beta ^{j^_\kappa })^{}\alpha ^j}^{\mu ^j}\right)c_{(\lambda ^0)^{}\alpha ^{0^{+_\kappa }}}^{\mu ^{0^{+_\kappa }}}c_{(\beta ^{M_\kappa ^_\kappa })^{}\lambda ^{M_\kappa }}^{\mu ^{M_\kappa }},$$
where the sum runs over all partitions $`\alpha ^1,\mathrm{},\alpha ^{M_\kappa 1},\alpha ^{M_\kappa +1},\mathrm{},\alpha ^{e2}`$ and $`\beta ^1,\mathrm{},\beta ^{M_\kappa 1},\beta ^{M_\kappa +1},\mathrm{},\beta ^{e2}`$. Similarly, we have
$$C_{\lambda ^{}\mu ^{}}=\left(\underset{\begin{array}{c}0ie1\\ i0,M_\kappa ^{}\end{array}}{}c_{\gamma ^i\delta ^i}^{\nu ^i}\underset{\begin{array}{c}0je1\\ j0,\mathrm{\hspace{0.25em}0}^{+_\kappa ^{}},M_\kappa ^{}\end{array}}{}c_{(\delta ^{j^_\kappa ^{}})^{}\gamma ^j}^{\rho ^j}\right)c_{(\nu ^0)^{}\gamma ^{0^{+_\kappa ^{}}}}^{\rho ^{0^{+_\kappa ^{}}}}c_{(\delta ^{M_\kappa ^{}^_\kappa ^{}})^{}\nu ^{M_\kappa ^{}}}^{\rho ^{M_\kappa ^{}}},$$
where the sum runs over all partitions $`\gamma ^1,\mathrm{},\gamma ^{M_\kappa ^{}1},\gamma ^{M_\kappa ^{}+1},\mathrm{},\gamma ^{e2}`$ and $`\delta ^1,\mathrm{},\delta ^{M_\kappa ^{}1},\delta ^{M_\kappa ^{}+1},\mathrm{},\delta ^{e2}`$, and $`e`$-quotients of $`\lambda ^{}`$ and $`\mu ^{}`$ are denoted as $`(\nu ^0,\mathrm{},\nu ^{e1})`$ and $`(\mathrm{},\rho ^1,\mathrm{},\rho ^{e1})`$ respectively. We now simplify $`C_{\lambda ^{}\mu ^{}}`$ using the fact that $`\nu ^i=(\lambda ^{\mathrm{\Phi }(i)})^{}`$ and $`\rho ^i=\mu ^{\mathrm{\Phi }(i)^{+_\kappa }}`$ for all $`i`$, and Lemma 3.2; we have
$`C_{\lambda ^{}\mu ^{}}`$ $`={\displaystyle \left(\underset{\begin{array}{c}0ie1\\ i0,M_\kappa \end{array}}{}c_{\gamma ^i\delta ^i}^{(\lambda ^{\mathrm{\Phi }(i)})^{}}\underset{\begin{array}{c}0je1\\ j0,\mathrm{\hspace{0.25em}0}^{+_\kappa ^{}},M_\kappa \end{array}}{}c_{(\delta ^{j^_\kappa ^{}})^{}\gamma ^j}^{\mu ^{\mathrm{\Phi }(j^_\kappa ^{})}}\right)c_{\lambda ^{M_\kappa }\gamma ^{0^{+_\kappa ^{}}}}^{\mu ^{M_\kappa }}c_{(\delta ^{M_\kappa ^_\kappa ^{}})^{}(\lambda ^0)^{}}^{\mu ^{0^+}}}`$
$`={\displaystyle \left(\underset{\begin{array}{c}0ie1\\ i0,M_\kappa \end{array}}{}c_{\gamma ^i\delta ^i}^{(\lambda ^{\mathrm{\Phi }(i)})^{}}\underset{\begin{array}{c}0je1\\ j0,M_\kappa ^_\kappa ^{},M_\kappa \end{array}}{}c_{(\delta ^j)^{}\gamma ^{j^{+_\kappa ^{}}}}^{\mu ^{\mathrm{\Phi }(j)}}\right)c_{\lambda ^{M_\kappa }\gamma ^{0^{+_\kappa ^{}}}}^{\mu ^{M_\kappa }}c_{(\delta ^{M_\kappa ^_\kappa ^{}})^{}(\lambda ^0)^{}}^{\mu ^{0^{+_\kappa }}}}`$
$`={\displaystyle \left(\underset{\begin{array}{c}0ie1\\ i0,M_\kappa \end{array}}{}c_{\gamma ^{\mathrm{\Phi }(i)}\delta ^{\mathrm{\Phi }(i)}}^{(\lambda ^i)^{}}\underset{\begin{array}{c}0je1\\ j0,\mathrm{\hspace{0.25em}0}^{+_\kappa },M_\kappa \end{array}}{}c_{(\delta ^{\mathrm{\Phi }(j)})^{}\gamma ^{\mathrm{\Phi }(j^_\kappa )}}^{\mu ^j}\right)c_{\lambda ^{M_\kappa }\gamma ^{\mathrm{\Phi }(M_\kappa ^_\kappa )}}^{\mu ^{M_\kappa }}c_{(\delta ^{\mathrm{\Phi }(0^{+_\kappa })})^{}(\lambda ^0)^{}}^{\mu ^{0^{+_\kappa }}}}`$
$`=C_{\lambda \mu }.`$
When $`C_{\lambda \mu }(=C_{\lambda ^{}\mu ^{}})0`$, then $`d_{\lambda \mu }(v)0`$, so that $`\lambda `$ and $`\mu `$ must have the same $`e`$-weight, i.e. $`_{i=0}^{e1}|\lambda ^i|=_{i=1}^{e1}|\mu ^i|`$ (since $`|\mu ^0|=0`$). Using Lemma 3.2, we have
$`\delta (\lambda ^{},\mu ^{})`$ $`={\displaystyle \underset{j=1}{\overset{e1}{}}}\pi _\kappa ^{}(j)(|\mu ^{\mathrm{\Phi }(j)^{+_\kappa }}||(\lambda ^{\mathrm{\Phi }(j)})^{}|)`$
$`={\displaystyle \underset{j=1}{\overset{e1}{}}}(e1\pi _\kappa (\mathrm{\Phi }(j)))(|\mu ^{\mathrm{\Phi }(j)^{+_\kappa }}||\lambda ^{\mathrm{\Phi }(j)}|)`$
$`=(e1){\displaystyle \underset{j=1}{\overset{e1}{}}}|\mu ^{\mathrm{\Phi }(j)^{+_\kappa }}|(e1){\displaystyle \underset{j=1}{\overset{e1}{}}}|\lambda ^{\mathrm{\Phi }(j)}|{\displaystyle \underset{j=1}{\overset{e1}{}}}\pi (\mathrm{\Phi }(j))(|\mu ^{\mathrm{\Phi }(j)^{+_\kappa }}||\lambda ^{\mathrm{\Phi }(j)}|)`$
$`=(e1)|\lambda ^{M_\kappa }|{\displaystyle \underset{\begin{array}{c}0je1\\ jM_\kappa \end{array}}{}}\pi (j)(|\mu ^{j^{+_\kappa }}||\lambda ^j|)`$
$`=(e1)|\lambda ^{M_\kappa }|{\displaystyle \underset{\begin{array}{c}0je1\\ jM_\kappa \end{array}}{}}\pi (j)|\mu ^{j^{+_\kappa }}|+{\displaystyle \underset{\begin{array}{c}0je1\\ jM_\kappa \end{array}}{}}\pi (j)|\lambda ^j|`$
$`=(e1)|\lambda ^{M_\kappa }|{\displaystyle \underset{j=1}{\overset{e1}{}}}(\pi (j)1)|\mu ^j|+{\displaystyle \underset{\begin{array}{c}0je1\\ jM_\kappa \end{array}}{}}\pi (j)|\lambda ^j|`$
$`={\displaystyle \underset{j=1}{\overset{e1}{}}}|\mu ^j|{\displaystyle \underset{j=1}{\overset{e1}{}}}\pi (j)(|\mu ^j||\lambda ^j|)`$
$`={\displaystyle \underset{j=1}{\overset{e1}{}}}|\mu ^j|\delta (\lambda ,\mu ).`$
This completes the proof. ∎
We end this section by obtaining some corollaries on the $`e`$-regular partitions contained in $`𝒫_\kappa ^{}`$ while assuming Theorem 3.6.
Firstly, the $`e`$-regular partitions in $`𝒫_\kappa ^{}`$ can be easily described:
###### Corollary 3.8.
Let $`\mu 𝒫_\kappa ^{}`$, with $`e`$-quotient $`(\mu ^0,\mathrm{},\mu ^{e1})`$. Then $`\mu `$ is $`e`$-regular if and only if $`\mu ^0=\mathrm{}`$.
###### Proof.
Suppose $`\mu `$ is $`e`$-regular. Then $`\mu ^0=\mathrm{}`$ by Theorem 3.6(1) and Proposition 3.7(1). Conversely, suppose $`\mu ^0=\mathrm{}`$. Then the least $`N`$ positions on runner $`0`$ in the abacus display of $`\mu `$ (with $`l(\kappa )+Ne`$ beads) are occupied, while the rest are vacant. By Lemma 3.5 (with $`i=0`$), the least $`(N1)`$ positions on runner $`j`$ are all occupied for all $`j>0`$. It follows that $`\mu `$ is $`e`$-regular. ∎
Also, $`𝒫_\kappa ^{}`$ is closed under the Mullineux involution:
###### Corollary 3.9.
Suppose $`\mu 𝒫_\kappa ^{}`$ is $`e`$-regular. Then $`\mu ^{}𝒫_\kappa ^{}^{}`$.
###### Proof.
This follows directly from Proposition 3.7(2), Lemma 3.2 and the definitions of $`𝒫_\kappa ^{}`$ and $`𝒫_\kappa ^{}^{}`$. ∎
## 4. Canonical basis
In this section, we provide the proof of part (1) of Theorem 3.6.
We begin with a reformulation of the Theorem.
###### Definition 4.1.
For any partition $`\mu `$, let $`H(\mu )=_{\lambda 𝒫}C_{\lambda \mu }v^{\delta (\lambda ,\mu )}\lambda `$.
Theorem 3.6(1) can thus be reformulated as follows:
If $`\mu 𝒫_\kappa ^{}`$, then $`G(\mu )=H(\mu )`$.
###### Lemma 4.2.
$`G(\mu )=H(\mu )`$ if and only if $`\overline{H(\mu )}=H(\mu )`$.
###### Proof.
Note that $`\delta (\lambda ,\mu )`$ is always non-negative. Furthermore, $`C_{\lambda \mu }0`$ and $`\delta (\lambda ,\mu )=0`$ if and only if $`\lambda =\mu `$, and $`C_{\mu \mu }=1`$. Thus, $`H(\mu )\mu _\lambda v[v]\lambda `$ and the lemma follows. ∎
We now review the results of . For an $`e`$-core partition $`\kappa `$ of Rouquier type, and integers $`a`$ and $`k`$ with $`1ae1`$ and $`k1`$, let
$$F_{a,k}=𝔣_a^{(k)}𝔣_{a+1}^{(k)}\mathrm{}𝔣_{e1}^{(k)}𝔣_{a1}^{(k)}𝔣_{a2}^{(k)}\mathrm{}𝔣_0^{(k)},$$
where, if $`j`$ is the residue class of $`i+l(\kappa )`$ modulo $`e`$, then $`𝔣_j=f_iU_v(\widehat{𝔰𝔩}_e)`$. Note that $`\overline{F_{a,k}(x)}=F_{a,k}(\overline{x})`$ for all $`x`$.
###### Lemma 4.3 (\[2, Lemma 3.1\]).
Let $`\kappa `$ be an $`e`$-core partition of Rouquier type. Let $`\lambda `$ be a partition with $`\kappa `$ and $`e`$-quotient $`(\lambda ^0,\mathrm{},\lambda ^{e1})`$, and consider its abacus display. Let $`x_a`$ and $`y_a`$ (resp. $`x_{a1}`$ and $`y_{a1}`$) be respectively the least vacant position and largest occupied position on runner $`a`$ (resp. $`a1`$). Suppose that
1. $`x_{a1}u>e`$ for all occupied positions $`u`$ lying on a runner to the left of runner $`a1`$;
2. $`ty_a>e`$ for all vacant positions $`t`$ lying on a runner to the right of runner $`a`$.
If $`1k(x_ay_{a1}1)/e`$, then
$$F_{a,k}(\lambda )=\underset{j=0}{\overset{k}{}}v^j\underset{\alpha ,\beta }{}c_{\lambda ^{a1}(j)}^\alpha c_{\lambda ^a(1^{kj})}^\beta \lambda (\alpha ,\beta ),$$
where $`\lambda (\alpha ,\beta )`$ denotes the partition having $`e`$-core $`\kappa `$ and $`e`$-quotient
$$(\lambda ^0,\mathrm{},\lambda ^{a2},\alpha ,\beta ,\lambda ^{a+1},\mathrm{},\lambda ^{e1}).$$
###### Note.
If $`(_{j=0}^{e1}c_{\alpha ^j\beta ^j}^{\lambda ^j}c_{(\beta ^{j1})^{}\alpha ^j}^{\mu ^j})0`$ with $`\beta ^1=\beta ^{e1}=\mathrm{}`$, then $`\delta (\lambda ,\mu )=_{j=0}^{e2}|\beta ^j|`$.
The proof of Proposition 4.1 of can be easily adapted to prove the following proposition.
###### Proposition 4.4.
Let $`\kappa `$ be an $`e`$-core partition of Rouquier type. Let $`\mu `$ be a partition with $`e`$-core $`\kappa `$ and $`e`$-quotient $`(\mu ^0,\mathrm{},\mu ^{e1})`$, and suppose that for all $`\lambda `$ such that $`C_{\lambda \mu }0`$, we have
$$F_{a,k}(\lambda )=\underset{j=0}{\overset{k}{}}v^j\underset{\alpha ,\beta }{}c_{\lambda ^{a1}(j)}^\alpha c_{\lambda ^a(1^{kj})}^\beta \lambda (\alpha ,\beta ).$$
Then
$$F_{a,k}(H(\mu ))=\underset{\eta }{}c_{\mu ^a(1^k)}^\eta H(\mu _\eta ),$$
where $`\mu _\eta `$ denotes the partition having $`e`$-core $`\kappa `$ and $`e`$-quotient
$$(\mu ^0,\mathrm{},\mu ^{a1},\eta ,\mu ^{a+1},\mathrm{}\mu ^{e1}).$$
The following is a direct consequence arising from the definition of $`𝒫_\kappa ^{}`$ and Lemma 3.3.
###### Lemma 4.5.
Let $`\kappa `$ be an $`e`$-core partition of Rouquier type. Let $`\mu 𝒫_\kappa ^{}`$ with $`e`$-quotient $`(\mu ^0,\mathrm{},\mu ^{e1})`$. Suppose $`\mu ^a\mathrm{}`$ for some $`a1`$. Let $`\tau ^a`$ be any partition such that $`|\tau ^a|<|\mu ^a|`$, let $`\tau `$ be the partition having $`e`$-core $`\kappa `$ and $`e`$-quotient $`(\mu ^0,\mathrm{},\mu ^{a1},\tau ^a,\mu ^{a+1},\mathrm{},\mu ^{e1})`$. Then any $`\lambda `$ with $`C_{\lambda \tau }0`$, $`a`$ and $`k=|\mu ^a||\tau ^a|`$ satisfy the hypothesis of Lemma 4.3.
###### Proposition 4.6.
Let $`\kappa `$ be an $`e`$-core partition of Rouquier type. Let $`\mu 𝒫_\kappa ^{}`$ . Then $`G(\mu )=H(\mu )`$.
###### Proof.
We prove by induction. Clearly, $`G(\kappa )=\kappa =H(\kappa )`$. Let $`|\mu |>|\kappa |`$, and suppose $`G(\tau )=H(\tau )`$ holds for all $`\tau 𝒫_\kappa ^{}`$ satisfying either $`|\tau |<|\mu |`$ or $`\tau <\mu `$. Let $`(\mu ^0,\mathrm{},\mu ^{e1})`$ be the $`e`$-quotient of $`\mu `$.
Suppose first that there exists $`a1`$ such that $`\mu ^a\mathrm{}`$. Let $`\tau ^a`$ be the partition obtained from $`\mu ^a`$ by removing its first column, and let $`\tau `$ denote the partition having $`e`$-core $`\kappa `$ and $`e`$-quotient $`(\mu ^0,\mathrm{},\mu ^{a1},\tau ^a,\mu ^{a+1},\mathrm{},\mu ^{e1})`$. Let $`k=|\mu ^a||\tau ^a|`$. By Lemmas 4.3 and 4.5 and Proposition 4.4, we have $`F_{a,k}(H(\tau ))=_\eta c_{\tau ^a(1^k)}^\eta H(\tau _\eta )`$. Now, $`\tau ,\tau _\eta 𝒫_\kappa ^{}`$ for all $`\eta `$ such that $`c_{\tau ^a(1^k)}^\eta 0`$. Furthermore, $`H(\mu )`$ occurs exactly once as a summand of $`F_{a,k}(H(\tau ))`$ and all the other summands $`H(\tau _\eta )`$ satisfy $`\tau _\eta <\mu `$. Thus, $`F_{a,k}(H(\tau ))`$, and all its summands $`H(\tau _\eta )`$ with $`\tau _\eta \mu `$ are bar-invariant by induction hypothesis, and hence so is $`H(\mu )`$. Thus $`G(\mu )=H(\mu )`$ by Lemma 4.2.
It remains to consider the case where $`\mu ^i=\mathrm{}`$ for all $`i>0`$. In this case, $`H(\mu )=\mu `$. Let $`\lambda `$ be a partition having $`e`$-core $`\kappa `$ and $`e`$-quotient $`(\lambda ^0,\mathrm{},\lambda ^{e1})`$. It is not difficult to see that if $`\lambda ^i\mathrm{}`$ for some $`i>0`$, then $`\lambda >\mu `$, so that $`d_{\lambda \mu }(v)=0`$. If $`\lambda ^i=\mathrm{}`$ for all $`i>0`$, then we shall show later in Corollary 5.6 that $`d_{\lambda \mu }^0=0`$ if $`\lambda \mu `$. Since $`d_{\lambda \mu }(1)=d_{\lambda \mu }^0`$ and $`d_{\lambda \mu }(v)_0[v]`$, we thus have $`d_{\lambda \mu }(v)=0`$ if $`\lambda \mu `$. Hence $`G(\mu )=\mu =H(\mu )`$. ∎
###### Proposition 4.7.
Let $`\kappa `$ be an $`e`$-core partition, and for $`0j<e`$, let $`n_j`$ be the number of beads on runner $`j`$ of its abacus display. Suppose that $`n_{i1}>n_i`$ for some $`i2`$. Given a partition $`\lambda `$ with $`e`$-core $`\kappa `$, let $`\mathrm{\Psi }(\lambda )`$ be the partition obtained from $`\lambda `$ by interchanging the $`(i1)`$-th and $`i`$-th runners of its abacus display. Let $`\mu 𝒫_\kappa ^{}`$. Then
1. $`\mathrm{\Psi }(\mu )𝒫_{\mathrm{\Psi }(\kappa )}^{}`$;
2. $`f_r^{(k)}H(\mu )=H(\mathrm{\Psi }(\mu ))`$;
3. $`e_r^{(k)}H(\mathrm{\Psi }(\mu ))=H(\mu )`$,
where $`r`$ is the residue class of $`l(\kappa )i`$ modulo $`e`$, and $`k=n_{i1}n_i`$.
###### Proof.
Part (1) follows from the definition of $`𝒫_\kappa ^{}`$ (and $`𝒫_{\mathrm{\Psi }(\kappa )}^{}`$). For parts (2) and (3), it suffices to show the following:
* for all partitions $`\lambda `$ with $`e`$-core $`\kappa `$, $`C_{\lambda \mu }=C_{\mathrm{\Psi }(\lambda )\mathrm{\Psi }(\mu )}`$ and $`\delta (\lambda ,\mu )=\delta (\mathrm{\Psi }(\lambda ),\mathrm{\Psi }(\mu ))`$;
* whenever $`C_{\tau \mu }0`$, the largest occupied position on runner $`i`$ of the abacus display of $`\tau `$ is less than the least vacant position on runner $`i1`$.
The first assertion follows from definition, while the second follows from Lemma 3.5. ∎
We are now ready to prove Theorem 3.6(1).
###### Proof of Theorem 3.6(1).
We prove by induction on $`N_\kappa =|\{(i,j)i<j,j_\kappa i\}|`$. If $`N_\kappa =0`$, then $`\kappa `$ is of Rouquier type, so that $`G(\mu )=H(\mu )`$ by Proposition 4.6. Assume thus $`N_\kappa >0`$. Then there exists $`2i<e`$ such that $`n_{i1}>n_i`$ (here, and hereafter in this proof, we keep the notations of Proposition 4.7). By Proposition 4.7(1), $`\mathrm{\Psi }(\mu )𝒫_{\mathrm{\Psi }(\kappa )}^{}`$. Since $`N_{\mathrm{\Psi }(\kappa )}<N_\kappa `$, we have $`G(\mathrm{\Psi }(\mu ))=H(\mathrm{\Psi }(\mu ))`$ by induction hypothesis. By Proposition 4.7(3), we have $`H(\mu )=e_r^{(k)}H(\mathrm{\Psi }(\mu ))=e_r^{(k)}G(\mathrm{\Psi }(\mu ))`$, so that $`H(\mu )`$ is bar-invariant. Hence $`G(\mu )=H(\mu )`$ by Lemma 4.2. ∎
## 5. $`q`$-Schur algebras
In this section, we provide the proof of part (1) of Theorem 3.6, which may be reformulated as follows: If $`\mu 𝒫_\kappa ^{}`$ with $`e`$-quotient $`(\mu ^0,\mu ^1,\mathrm{},\mu ^{e1})`$, then $`d_{\lambda \mu }^l=d_{\lambda \mu }^0`$ for all $`l>\mathrm{max}_i(|\mu ^i|)`$.
###### Definition 5.1.
Let $``$ be the equivalence relation defined on $`𝒫`$ as follows: $`\lambda \mu `$ if and only if $`\lambda `$ and $`\mu `$ are partitions having the same $`e`$-core, and $`|\lambda ^i|=|\mu ^i|`$ for all $`i`$, where $`(\lambda ^0,\lambda ^1,\mathrm{},\lambda ^{e1})`$ and $`(\mu ^0,\mu ^1,\mathrm{},\mu ^{e1})`$ are the $`e`$-quotients of $`\lambda `$ and $`\mu `$ respectively.
###### Definition 5.2.
Let $`\kappa `$ be an $`e`$-core partition of Rouquier type. Let $`\overline{𝒫}_\kappa `$ be a collection of partitions with $`e`$-core $`\kappa `$, having the following two properties:
1. On the abacus display of any partition in $`\overline{𝒫}_\kappa `$, any pair $`(x,y)`$, where $`x`$ is an occupied position on runner $`i`$ while $`y`$ is a vacant position on runner $`j`$, with $`i<j`$, satisfies $`x<y`$;
2. If $`\lambda ,\mu \overline{𝒫}_\kappa `$, $`\lambda \mu `$, with respective $`e`$-quotients $`(\lambda ^0,\mathrm{},\lambda ^{e1})`$ and $`(\mu ^0,\mathrm{},\mu ^{e1})`$, then every partition having $`e`$-core $`\kappa `$, and $`e`$-quotient of the form $`(\mu ^0,\mathrm{},\mu ^{i1},\lambda ^i,\mathrm{},\lambda ^{e1})`$ ($`1ie1`$) lies in $`\overline{𝒫}_\kappa `$.
###### Lemma 5.3.
Let $`\kappa `$ be an $`e`$-core partition of Rouquier type. Then $`𝒫_\kappa ^{}`$ satisfies the two conditions of $`\overline{𝒫}_\kappa `$.
###### Proof.
If $`\mu 𝒫_\kappa ^{}`$, then $`\mu `$ satisfies condition (1) of $`\overline{𝒫}_\kappa `$ by Lemma 3.5. That condition (2) of $`\overline{𝒫}_\kappa `$ also holds follows directly from the definition of $`𝒫_\kappa ^{}`$ . ∎
The following is an immediate consequence.
###### Lemma 5.4.
Let $`\kappa `$ be an $`e`$-core partition of Rouquier type. Let $`\lambda \overline{𝒫}_\kappa `$, and let $`s(\lambda )=(c_1,c_2,\mathrm{},c_w)`$ be its induced $`e`$-sequence. Then $`\overline{c_i}\overline{c_{i+1}}`$ for all $`1iw1`$, where $`\overline{x}`$ denotes the residue class of $`x`$ modulo $`e`$.
###### Proof.
This follows from Lemma 2.8(2) and condition (1) of $`\overline{𝒫}_\kappa `$. ∎
###### Proposition 5.5.
Let $`\kappa `$ be an $`e`$-core partition of Rouquier type. Suppose $`\lambda ,\mu \overline{𝒫}_\kappa `$ with $`\lambda \mu `$. If $`\lambda \tau `$ with $`\lambda \tau `$, then $`\tau _p\mu `$.
###### Proof.
Suppose $`\lambda \stackrel{𝜎}{}\tau `$. By Lemma 2.8(3), there exist integers $`N,x,y`$ with $`N1`$ and $`0x<y`$ such that
$`s(\lambda )`$ $`=s(\sigma )(x,xe,\mathrm{},x(N1)e),`$
$`s(\tau )`$ $`=s(\sigma )(y,ye,\mathrm{},y(N1)e).`$
Suppose $`\overline{x}=i`$ and $`\overline{y}=j`$, where $`\overline{a}`$ denotes the residue class of $`a`$ modulo $`e`$. Condition (1) of $`\overline{𝒫}_\kappa `$ forces $`j>i`$. Let $`s(\lambda )=(l_1,l_2,\mathrm{},l_w)`$, and let $`r`$ be the least index such that $`\overline{l_r}=i`$. Then $`\overline{l_s}>i`$ for all $`s<r`$ by Lemma 5.4, so that $`l_1,l_2,\mathrm{},l_{r1}`$ are terms of $`s(\sigma )`$. Let $`s(\mu )=(m_1,m_2,\mathrm{},m_w)`$. Then since $`\lambda \mu `$, we have $`\overline{m_t}=\overline{l_t}`$ for all $`t`$ by Lemma 5.4; in particular, $`\overline{m_s}>i`$ for all $`s<r`$. Let $`\rho `$ denote the partition having $`e`$-core $`\kappa `$ and $`e`$-quotient $`(\mu ^0,\mathrm{},\mu ^i,\lambda ^{i+1},\mathrm{},\lambda ^{e1})`$. Then $`\rho \overline{𝒫}_\kappa `$, and $`s(\rho )=(l_1,\mathrm{},l_{r1},m_r,\mathrm{},m_w)`$ by Lemma 5.4. In the abacus display of $`\rho `$, $`m_r`$ is the largest occupied position on runner $`i`$, so it is less than all the vacant positions on runners to the right of runner $`i`$ by condition (1) of $`\overline{𝒫}`$; in particular, $`m_r<y`$ since $`y`$ is a vacant position on runner $`j`$ of $`\lambda `$ (and $`j>i`$), and hence of $`\rho `$. Thus, $`m_r`$ is less than $`r`$ terms of $`s(\tau )`$, namely, $`l_1,l_2,\mathrm{},l_{r1},y`$. This implies that $`\mu _p\tau `$. ∎
###### Corollary 5.6.
Let $`\kappa `$ be an $`e`$-core partition of Rouquier type. Let $`\lambda ,\mu \overline{𝒫}_\kappa `$ with $`\lambda \mu `$. If $`l=0`$ or $`l`$ is greater than the size of each constituent of the $`e`$-quotient of $`\lambda `$, then $`d_{\lambda \mu }^l=\delta _{\lambda \mu }`$.
###### Proof.
By Theorem 2.6, it suffices to prove that $`J_{\lambda \mu }=0`$ when $`\lambda \mu `$. For each $`\tau `$ such that $`\lambda \tau `$, let $`j_{\tau \mu }=_\sigma (1)^{l_{\lambda \sigma }+l_{\tau \sigma }+1}`$, where $`\sigma `$ runs over all partitions such that $`\lambda \stackrel{𝜎}{}\tau `$. By Proposition 5.5 and Lemmas 2.9 and 2.7, $`d_{\tau \mu }^l=0`$ for all $`\tau `$ such that $`\lambda \tau `$ and $`\lambda \tau `$. Furthermore, the condition on $`l`$ implies that we always have $`\nu _l(h_{\lambda \sigma })=0`$. Thus $`J_{\lambda \mu }=_\tau j_{\tau \mu }d_{\tau \mu }^l`$, where $`\tau `$ runs over all partitions such that $`\lambda \tau `$ and $`\lambda \tau `$. Fix such a $`\tau `$; suppose $`\tau `$ can be obtained from $`\lambda `$ by first moving a bead at position $`x`$ on runner $`j`$ to position $`xNe`$ and then moving a bead at position $`yNe`$ on runner $`j`$ to $`y`$, with $`x<y`$. Let $`\rho _1`$ (resp. $`\rho _2`$) be the partition obtained from $`\lambda `$ by moving the bead at position $`x`$ (resp. $`yNe`$) to $`xNe`$. Then $`\lambda \stackrel{𝜎}{}\tau `$ if and only if $`\sigma =\rho _1`$ or $`\rho _2`$. Furthermore, $`l_{\lambda \rho _1}+l_{\tau \rho _1}`$ and $`l_{\lambda \rho _2}+l_{\tau \rho _2}`$ are of different parity. Thus $`j_{\tau \mu }=0`$. Hence $`J_{\lambda \mu }=0`$. ∎
We need the following results, which are analogous to each other, and which we believe are well-known, but we are not able to find an appropriate reference in the existing literature.
###### Proposition 5.7.
Suppose $`P_{}^\mu =_{i=1}^na_iP_{}^{\sigma _i}`$ (where $`()`$ is a composition of (divided powers) of $`r`$-induction functors), $`d_{\lambda \mu }^0=d_{\lambda \mu }^l`$ for all $`\lambda `$, and $`d_{\sigma _i\sigma _j}^l=\delta _{ij}`$. Then $`P_𝔽^\mu =_{i=1}^na_iP_𝔽^{\sigma _i}`$, and $`d_{\tau \sigma _i}^0=d_{\tau \sigma _i}^l`$ for all $`\tau `$ and all $`i`$.
###### Proposition 5.8.
Suppose $`P_{}^\mu =_{i=1}^na_iP_{}^{\sigma _i}`$ (where $`()`$ is a composition of (divided powers) of $`r`$-restriction functors), $`d_{\lambda \mu }^0=d_{\lambda \mu }^l`$ for all $`\lambda `$, and $`d_{\sigma _i\sigma _j}^l=\delta _{ij}`$. Then $`P_𝔽^\mu =_{i=1}^na_iP_𝔽^{\sigma _i}`$, and $`d_{\tau \sigma _i}^0=d_{\tau \sigma _i}^l`$ for all $`\tau `$ and all $`i`$.
We provide a proof of Proposition 5.7 only; that of Proposition 5.8 is entirely analogous.
###### Proof of Proposition 5.7.
Note that, since $`\mathrm{\Delta }_𝔽^\mu `$ and $`\mathrm{\Delta }_{}^\mu `$ have the same branching rule and $`d_{\lambda \mu }^0=d_{\lambda \mu }^l`$ for all $`\lambda `$, we have $`[P_𝔽^\mu :\mathrm{\Delta }_𝔽^\tau ]=[P_{}^\mu :\mathrm{\Delta }_{}^\tau ]=_{i=1}^na_id_{\tau \sigma _i}^0_{i=1}^na_id_{\tau \sigma _i}^l`$. We may assume that $`\sigma _i>_J\sigma _j`$ implies that $`i<j`$. We prove by induction on $`r`$ that $`_{i=1}^ra_iP_𝔽^{\sigma _i}`$ is a direct summand of $`P_𝔽^\mu `$. This is clear for $`r=1`$. Suppose $`_{i=1}^{r1}a_iP_𝔽^{\sigma _i}`$ is a direct summand of $`P_𝔽^\mu `$, so that
($``$)
$$\underset{i=1}{\overset{r1}{}}a_id_{\tau \sigma _i}^l=[\underset{i=1}{\overset{r1}{}}a_iP_𝔽^{\sigma _i}:\mathrm{\Delta }_𝔽^\tau ][P_𝔽^\mu :\mathrm{\Delta }_𝔽^\tau ]\underset{i=1}{\overset{n}{}}a_id_{\tau \sigma _i}^l.$$
for all $`\tau `$. If $`\tau >_J\sigma _r`$, then $`\sigma _j_J\tau `$ for all $`jr`$, so that $`d_{\tau \sigma _j}^l=0`$, and hence we have equality in $`()`$. This shows that if $`\tau >_J\sigma _r`$, then $`P_𝔽^\tau `$ is not a summand of $`P_𝔽^\mu `$. Now, $`[_{i=1}^{r1}a_iP_𝔽^{\sigma _i}:\mathrm{\Delta }_𝔽^{\sigma _r}]=_{i=1}^{r1}a_id_{\sigma _r\sigma _i}^l=0`$, while $`[P_𝔽^\mu :\mathrm{\Delta }_𝔽^{\sigma _r}]=a_r`$. Thus, $`a_rP_𝔽^{\sigma _r}`$ is a direct summand of $`P_𝔽^\mu `$. Hence, by induction, $`_{i=1}^na_iP_𝔽^{\sigma _i}`$ is a direct summand of $`P_𝔽^\mu `$, and so $`()`$ holds when $`r1`$ is replaced by $`n`$, and in fact with equality throughout. Thus $`P_𝔽^\mu =_{i=1}^na_iP_𝔽^{\sigma _i}`$, and $`d_{\tau \sigma _i}^0=d_{\tau \sigma _i}^l`$ for all $`\tau `$ and all $`i`$. ∎
We are now ready to prove Theorem 3.6(2).
###### Proof of Theorem 3.6(2).
We only need to show the Theorem holds when $`l>\mathrm{max}_i(|\mu ^i|)`$ and $`\lambda `$ has $`e`$-core $`\kappa `$ and $`e`$-quotient $`(\lambda ^0,\mathrm{},\lambda ^{e1})`$ with $`_{i=0}^{e1}|\lambda ^i|=_{i=0}^{e1}|\mu ^i|`$.
Let $`N_\kappa =|\{(i,j)i<j,j_\kappa i\}|`$.
Suppose first that $`N_\kappa =0`$. Then $`\kappa `$ is of Rouquier type.
If $`\lambda ^i\mathrm{}`$ for some $`i1`$, then $`\lambda >\mu `$ so that $`d_{\lambda \mu }^l=0=C_{\lambda \mu }`$. If $`\lambda ^i=\mathrm{}`$ for all $`i1`$, then $`\lambda \mu `$, so that $`d_{\lambda \mu }^l=\delta _{\lambda \mu }=C_{\lambda \mu }`$ by Corollary 5.6.
From the proof of Proposition 4.6, we can find a $`\tau 𝒫_\kappa ^{}`$ and an integer $`k`$ such that $`G(\mu )`$ occurs exactly once as a summand of $`F_{a,k}G(\tau )`$. Furthermore, if $`G(\alpha )`$ and $`G(\beta )`$ are summands of $`F_{a,k}G(\tau )`$, then $`\alpha \beta `$, so that $`d_{\alpha ,\beta }^l=\delta _{\alpha \beta }`$ by Corollary 5.6. We may assume that $`d_{\sigma \tau }^l=d_{\sigma \tau }^0`$ for all $`\sigma 𝒫`$ by induction hypothesis. Thus, $`d_{\lambda \mu }^l=d_{\lambda \mu }^0`$ for all $`\lambda 𝒫`$ by Proposition 5.7.
Now, suppose $`N_\kappa >0`$. Then there exists $`2i<e`$ such that $`n_{i1}>n_i`$ (here, and hereafter in this proof, we keep the notations of Proposition 4.7). By Proposition 4.7, we have $`e_r^{(k)}G(\mathrm{\Psi }(\mu ))=G(\mu )`$. Since $`N_{\mathrm{\Psi }(\kappa )}<N_\kappa `$, we may assume that $`d_{\sigma \mathrm{\Psi }(\mu )}^0=d_{\sigma \mathrm{\Psi }(\mu )}^l`$ for all $`\sigma 𝒫`$ by induction hypothesis. Thus $`d_{\lambda \mu }^0=d_{\lambda \mu }^l`$ for all $`\lambda 𝒫`$ by Proposition 5.8. ∎ |
warning/0506/math0506413.html | ar5iv | text | # Bounding right-arm rotation distances
## 1. Introduction
Rotation distance quantifies the difference in shape between two rooted binary trees of the same size by counting the minimum number of elementary changes needed to transform one tree to the other. Search algorithms are most efficient when searching balanced trees, which have few levels relative to the number of nodes in the tree. Thus one is often interested in calculating, or at least bounding, the number of these changes necessary to alter a given tree into another with a more desirable shape, such as a balanced tree.
If we allow these elementary changes, called rotations, to take place at any node, we obtain ordinary rotation distance. This was analyzed by Sleator, Tarjan and Thurston , who proved an upper bound of $`2n6`$ rotations needed to transform one rooted binary tree with $`n`$ nodes into any other, for $`n11`$. Furthermore, they showed that the $`2n6`$ bound is achieved for all sufficiently large $`n`$ and thus is the best possible upper bound. No efficient algorithm is known to compute rotation distance exactly, though there are polynomial-time algorithms of Pallo and Rogers which estimate rotation distance efficiently.
Here we expand on the study of restricted rotation distance begun in and . Restricted rotation distance allows rotations only at the root node and the right child of the root node. Restricted rotation distance is related to the word length of elements of Thompson’s group $`F`$ with respect to its standard finite generating set. This is illustrated in and involves the interpretation of elements of $`F`$ as pairs of finite binary rooted trees and Fordham’s method for computing the word length of an element of $`F`$ with respect to that standard finite generating set directly from such trees. These methods not only give an effective algorithm to compute restricted rotation distance, but they also give an effective algorithm to find the appropriate rotations which realize this distance.
Right and left rotations at a node $`N`$ of a rooted binary tree $`T`$ are defined to be the permutations of the subtrees of $`T`$ described in Figure 1. Right rotation at a node $`N`$ transforms the original tree $`T_1`$, given on the left side of Figure 1, to the tree $`T_2`$ on the right side of Figure 1. Left rotation at a node is the inverse operation. In all that follows, $`T_1`$ and $`T_2`$ denote trees with the same number of nodes.
In this paper, we discuss generalizations and variations of restricted rotation distance, in which rotations are again only allowed at specified nodes of the tree. We relate these distances to distinct word metrics on Thompson’s group $`F`$. We use this interpretation to exhibit linear bounds on the number of allowable elementary rotations needed to transform one tree with $`n`$ nodes into another, and show that these bounds are asymptotically sharp in the sense that the coefficients of the linear terms of the bounds are the best possible. These alternate definitions all allow rotations at the root node and at nodes connected to the root node by a path consisting entirely of right edges; that is, nodes that lie on the right side or right arm of the tree. The root node is considered to lie on the right side of the tree.
One complication that arises is that while the original restricted rotation distance is always defined between any two trees with the same number of nodes, this is no longer necessarily the case when we allow rotations at other collections of nodes along the right side of the tree. Some transformations between trees cannot be accomplished with a specified set of rotations without the nonstandard technique of adding additional nodes to the trees. Such a transformation is not permitted when computing rotation distances of any type. Below, we describe how to determine when such transformations are possible with a prescribed set of permitted locations at which to rotate. When this restricted right-arm rotation distance is defined, we provide an upper bound on its magnitude.
The sharp upper bound on the restricted rotation distance between two trees, each with $`n`$ nodes, obtained in is $`4n8`$, for $`n3`$. Below, we consider allowing additional rotations along the right side of the tree and note that allowing rotations at any finite collection of nodes on the right side of the tree does not change the multiplicative constant of 4 in the upper bound. It is only when we allow an infinite set of rotations along the right arm of the tree that we obtain the multiplicative constant of 2 in the upper bound, analogous to ordinary rotation distance. These rotation distances and bounds, which hold for sufficiently large $`n`$, are summarized in Table 1, where $`n`$ is the number of nodes in each tree.
Culik and Wood , in the course of studying ordinary rotation distance, showed that the rotation distance is never more than $`2n2`$, and in fact use only rotations on the right arm to show this bound. Since the ordinary rotation distance between two such trees can only be as much as $`2n6`$ for $`n11`$, it is remarkable that restricting rotations to the right side of the tree adds only four rotations to the upper bound.
Pallo explicitly studied right-arm rotation distance in , allowing rotations at all nodes along the right side of the tree. He described an algorithm for computing right-arm rotation distance which we show below is equivalent to finding the word length in Thompson’s group $`F`$ with respect to the standard infinite generating set.
The trees we consider are composed of edges and vertices. The vertices fall into two types: those of valence one and those of higher valence. The vertices of valence one are called exterior nodes or leaves or exposed leaves. The vertices of higher valence are called interior nodes. We shall use the terminology node to refer to a vertex which is an interior node, and leaf to refer to a vertex which is an exterior node.
A caret in a tree is composed of a node together with two downward directed edges. We will only consider finite, rooted binary trees with $`n`$ carets, equivalently, with $`n`$ nodes. Such trees are called extended binary trees in Knuth or 0-2 trees. The nodes and carets in a tree have a natural infix ordering. The exposed leaves in a tree are numbered from left to right, beginning with zero. A tree with $`n`$ carets yields $`n+1`$ exposed leaves. A caret with two exposed leaves is called an exposed caret, its leaves are termed siblings and those leaves are said to form a sibling pair.
A caret $`N`$ which is attached to the right (respectively left) edge of a caret $`M`$ is called the right (respectively left) child of $`M`$. A caret which has one edge on the left side of the tree is called a left caret. A caret which has one edge on the right side of the tree and is not the root caret is called a right caret. Similarly, we have left and right nodes. Carets which are neither right nor left are called interior carets. The union of left and right carets in a tree is called the spine of the tree. A tree consisting of only the root caret and $`n1`$ right carets is called the all-right tree with $`n`$ carets. An ancestor of a caret (resp. node) is any caret (resp. node) which lies along the shortest path between it and the root caret (resp. root node).
The connection between Thompson’s group $`F`$ and restricted rotation distance is described below. Thompson’s group $`F`$ is studied combinatorially in two ways: via a finite presentation and an infinite presentation. Computing restricted rotation distance between two trees is related to computing the word length of the element of $`F`$ described by those trees with respect to the standard finite generating set for the group $`F`$. Analogously, right-arm rotation distance corresponds to computing the word length of the element with respect to the word metric induced by the standard infinite generating set for $`F`$. Restricted right-arm rotation distances and restricted spinal rotation distances, defined below, relate to the word metric on $`F`$ with respect to other finite generating sets.
## 2. Thompson’s Group $`F`$
The connection between Thompson’s group $`F`$ and rotations at nodes of trees is described in and , using the work of Fordham . Here, we briefly describe this connection, and refer the reader to Cannon, Floyd and Parry for a survey of the properties of Thompson’s group $`F`$, and the further connections between elements of $`F`$ and pairs of binary rooted trees.
### 2.1. The infinite presentation of Thompson’s group $`F`$
Thompson’s group $`F`$ has a presentation with an infinite number of generators and relations:
$$𝒫=x_0,x_1,\mathrm{}|x_i^1x_nx_i=x_{n+1},i<n.$$
In this presentation, there are normal forms for elements given by
$$x_{i_1}^{r_1}x_{i_2}^{r_2}\mathrm{}x_{i_k}^{r_k}x_{j_l}^{s_l}\mathrm{}x_{j_2}^{s_2}x_{j_1}^{s_1}$$
with $`r_i,s_i>0`$, where the indices satisfy $`0i_1<i_2<\mathrm{}<i_k`$ and $`0j_1<j_2<\mathrm{}<j_l`$. This normal form is unique for a given element if we further require the reduction condition that when both $`x_i`$ and $`x_i^1`$ occur, so does $`x_{i+1}`$ or $`x_{i+1}^1`$, as discussed by Brown and Geoghegan . The relators provide a quick and efficient method for rewriting words into normal form, and form a complete rewriting system, as described by Brown . There is a natural shift homomorphism $`\varphi :FF`$ where $`\varphi (x_i)=x_{i+1}`$ which respects the relators, and the reduction from normal form to unique normal form is accomplished with a sequence of operations replacing words of the form $`ux_i\varphi (v)x_i^1w`$ with $`uvw`$, where $`\varphi (v)`$ is a subword which contains only generators of index $`i+2`$ and higher.
We note that $`F`$ can be generated by just $`x_0`$ and $`x_1`$ in the above presentation; the relators show that $`x_0`$ conjugates $`x_1`$ to $`x_2`$. Similarly, all higher-index generators are conjugates of $`x_1`$ by higher powers of $`x_0`$, as $`x_n=x_0^{(n1)}x_1x_0^{n1}`$. This leads to a finite presentation for $`F`$ with generating set $`\{x_0,x_1\}`$. In fact, $`x_0`$ and any higher index generator are sufficient to generate the group. Any two generators $`x_i,x_j`$ with $`ij`$ will generate an subgroup of $`F`$ which is isomorphic to the entire group but which is the entire group only when one of $`i`$ or $`j`$ is $`0`$.
We begin by proving that in the word metric arising from this infinite generating set, the normal form expressions are geodesic representatives for elements of $`F`$.
###### Lemma 2.1.
Let $`w`$ be an element of $`F`$, and $`\alpha `$ a word in the infinite generating set which is the unique normal form for $`w`$, as described above. Then $`\alpha `$ is a geodesic representative for $`w`$ in the word metric arising from the infinite generating set $`\{x_i\}`$ of $`F`$.
###### Proof.
Suppose that $`\alpha =x_{i_1}^{r_1}x_{i_2}^{r_2}\mathrm{}x_{i_k}^{r_k}x_{j_l}^{s_l}\mathrm{}x_{j_2}^{s_2}x_{j_1}^{s_1}`$ was not a geodesic representative for $`w`$ in this word metric. Then there is a shorter expression $`\beta `$, not necessarily in normal form, representing $`w`$ in this infinite generating set. It is clear from the relations of $`𝒫`$ that the conversion of $`\beta `$ into unique normal form can only preserve or decrease the length of $`\beta `$. Thus, after converting $`\beta `$ into normal form we have obtained a second expression for $`w`$ in unique normal form shorter than the initial unique normal form for $`w`$ given by $`\alpha `$, a contradiction. ∎
### 2.2. Tree pair diagrams for elements of Thompson’s group $`F`$
The group $`F`$ has a geometric description in terms of equivalence classes of tree pair diagrams. A tree pair diagram is a pair of finite rooted binary trees with the same number of nodes (or carets), or equivalently with the same number of leaves. We write $`w=(T_1,T_2)`$ to denote the two trees comprising a pair representing $`w`$. The equivalence between the geometric and algebraic interpretations of $`F`$ is described in , and examples of this equivalence and its connection with rotations are given in .
Given two trees $`T_1`$ and $`T_2`$ with the same number of carets, the word in normal form associated to $`w=(T_1,T_2)`$ is found as follows. The leaves of each tree are numbered from left to right, beginning with zero. The leaf exponent of a leaf numbered $`k`$ is the integral length of the longest path starting at leaf $`k`$ consisting entirely of left edges which does not touch the right side of the tree. The tree pair diagram $`(T_1,T_2)`$ has an associated normal form $`x_0^{f_0}x_1^{f_1}\mathrm{}x_n^{f_n}x_n^{e_n}\mathrm{}x_1^{e_1}x_0^{e_0}`$ where $`e_i`$ is the leaf exponent of leaf $`i`$ in tree $`T_1`$ and $`f_i`$ is the leaf exponent of leaf $`i`$ in $`T_2`$. An example of a tree with leaf exponents computed is given in Figure 2.
An element of $`F`$ is represented uniquely by a tree pair diagram satisfying the following reduction condition. A tree pair diagram $`(T_1,T_2)`$ is unreduced if both $`T_1`$ and $`T_2`$ contain a caret with two exposed leaves numbered $`i`$ and $`i+1`$. A tree pair diagram which is not unreduced is reduced. Geometrically, any tree pair diagram has a unique reduced form that is obtained by successively deleting exposed carets with identical leaf numbers from both trees, renumbering the leaves, and repeating this process until no further such reductions are possible. Elements of $`F`$ are equivalence classes of tree pair diagrams, where the equivalence relation is that two tree pairs are equivalent if they have a common reduced form.
This tree pair reduction condition corresponds exactly to the combinatorial reduction condition given above to ensure uniqueness for words in normal form in the infinite presentation of $`F`$. That is, if leaves $`i`$ and $`i+1`$ form a sibling pair in both $`T_1`$ and $`T_2`$, then in both cases, the leaf exponents of leaf $`i`$ will be non-zero in both trees and those for leaf $`i+1`$ will be zero, as leaf $`i+1`$ is a right leaf in both trees. So the corresponding normal form will contain both $`x_i`$ and $`x_i^1`$ but not $`x_{i+1}^{\pm 1}`$, meaning that the normal form can be reduced.
To perform the group operation on the level of tree pair diagrams, it may be necessary to use unreduced representatives of elements. Namely, to multiply $`(T_1,T_2)`$ and $`(S_1,S_2)`$, we create unreduced representatives $`(T_1^{},T_2^{})`$ and $`(S_1^{},S_2^{})`$ in which $`T_2^{}=S_1^{}`$, and write the product as the (possibly unreduced) element $`(T_1^{},S_2^{})`$. See for examples of group multiplication using tree pair diagrams for elements of $`F`$.
The reduced tree pair diagrams associated to the generators $`x_0,x_1`$ and $`x_n`$ are pictured in Figure 3. As explained in Lemmas 2.6 and 2.7 of , the generators $`x_0`$ and $`x_1`$ can be viewed in terms of rotations of rooted binary trees as well. The generator $`x_0`$ can be interpreted as a left rotation at the root of the left tree in the pair, yielding the right tree in the pair. Similarly, the generator $`x_1`$ performs a left rotation at the right child of the root, transforming the left tree in the pair to the right one. The inverses $`x_0^1`$ and $`x_1^1`$ perform right rotations at the root node and right child of the root node, respectively.
Analogously, right multiplication of an element $`w`$ given by a possibly unreduced representative $`(T_1,T_2)`$ by the generator $`x_0^1`$ yields the tree pair diagram $`(T_1^{},T_2^{})`$ in which $`T_1^{}`$ differs from $`T_1`$ by a left rotation at the root node, and $`T_2^{}=T_2`$. We can similarly interpret right multiplication by $`x_0,x_1^{\pm 1}`$ and $`x_n^{\pm 1}`$.
One complication that may arise when using the geometry of the tree pair diagrams to understand rotation distance is the possibility of requiring unreduced representatives in order to perform the group multiplication. Since elements of Thompson’s group are equivalence classes of tree pair diagrams, we can always multiply any group element $`w`$ by any group generator $`g`$. It is possible that we may have to add carets to the reduced tree pair diagram for $`w`$ in order to carry out this multiplication. From the standpoint of group theory, the reduced and unreduced tree pair diagrams are interchangeable. When considering rotation distance, however, we are not allowed to change the number of carets in the starting tree. Thus certain rotations, corresponding to multiplication by specific generators, may not be permitted when calculating rotation distance.
For example, we cannot perform a right rotation at the right child of the root to either of the trees in the tree pair diagram for $`x_0`$ as shown in Figure 3 because neither tree contains a left subtree of the right child of the root. As an element of Thompson’s group, we can enlarge any pair of trees to another representative in the equivalence class, and thus are able to multiply any element by any generator. A typical such application is shown in Figure 4 where a caret is added to a tree to be able to perform the desired rotation. The tree pair diagram $`w`$ does not have a left child of the right child of the root, so performing a right rotation at the right child of the root is not possible. However, the word $`\overline{w}`$ which represents the same element of $`F`$ does have a left child of the right child of the root and it is possible to perform the right rotation at the right child of the root there. We obtain $`\overline{w}`$ by adding an additional caret (indicated by dashing) to leaf number 2 in both trees of the tree pair diagram.
To describe when it is necessary to add a caret to a tree to perform a particular rotation, we make the following definitions. We say a right rotation at the root can be applied to a tree $`T`$ if the left subtree of the root of $`T`$ is non-empty. Similarly, we say a left rotation at the root can be applied to a tree $`T`$ if the right subtree of the root of $`T`$ is non-empty and we also adopt this terminology when performing rotations at other nodes along the spine of the tree.
Understanding when rotations can be performed on trees helps us develop the connection between rotations of trees and right multiplication by generators of $`F`$. If, for example, we have a tree pair diagram $`(T_1,T_1)`$ representing the identity and we can perform a left rotation at the root to $`T_1`$ to obtain $`x_0T_1`$, then the new tree pair diagram $`(T_1,x_0T_1)`$ is the tree pair diagram representing the word $`x_0`$ in $`F`$, and similarly the new tree pair diagram $`(x_0T_1,T_1)`$ is tree pair diagram representing the word $`x_0^1`$ in $`F`$.
We note that in $`F`$, multiplication by a generator may result in an unreduced tree pair diagram. So during the course of a sequence of multiplications by generators of $`F`$, the number of carets in the reduced tree pair diagram representing the partial products may fluctuate – rising when it is necessary to add one or more carets to apply a generator, and falling when multiplication by a generator results in an unreduced tree pair diagram. To understand rotation distance, however, as we apply a sequence of rotations to a single tree, we do not allow the number of carets in the tree to change.
The link between restricted rotation distance and Thompson’s group $`F`$ is the word metric on $`F`$ with respect to the generators $`\{x_0,x_1\}`$. Given two rooted binary trees $`T_1`$ and $`T_2`$ with the same number of nodes, we consider a minimal length word in $`x_0^{\pm 1}`$ and $`x_1^{\pm 1}`$ representing the element $`w=(T_1,T_2)F`$. As described in , this word gives a minimal sequence of rotations at the root and right child of the root which transform the tree $`T_1`$ into the tree $`T_2`$. It follows from Fordham that these minimal words which transform one tree into the other maintain a constant number of carets at each stage in the sequence of rotations. The issue of certain rotations altering the number of nodes in the tree does not arise in the case of restricted rotation distance.
More precisely, suppose that $`wF`$ is given by the tree pair diagram $`(T_1,T_2)`$, and a minimal length representative for $`w`$ is $`g_1g_2\mathrm{}g_n`$, where each $`g_i\{x_0^{\pm 1},x_1^{\pm 1}\}`$. Then the tree pair diagram $`(T_1,g_n\mathrm{}g_2g_1T_1)`$ will represent $`w`$, and we can think of the sequence of generators $`g_n\mathrm{}g_2g_1`$ as a sequence of rotations which transforms $`T_1`$ to $`T_2`$. At each stage of this process, we will be able to perform the rotation corresponding to the generator $`g_{i+1}`$ to the tree $`g_i\mathrm{}g_2g_1T_1`$ without adding extra carets. There may be reductions possible to tree pair diagrams, or equivalently to the normal forms, during this process, but from the standpoint of rotation distance we do not want to take advantage of these reductions. Instead, we keep the number of carets constant at each stage.
Equivalently, we can think of $`(g_n\mathrm{}g_2g_1T_1,T_2)`$ as a representative of the identity and witness the transformation of $`T_1`$ to $`T_2`$ by considering the sequence of tree pair diagrams
$$(T_1,T_2),(g_1T_1,T_2),\mathrm{},(g_n\mathrm{}g_2g_1T_1,T_2).$$
Below, we consider other possible locations for rotations to occur, and again exploit the link to Thompson’s group $`F`$, but now considering other appropriate generating sets for $`F`$, where the generators are chosen to reflect the locations where rotations are permitted. We assign a level to each node or caret in the tree as follows. The root node is defined to have level zero. The level of a node $`N`$ is the number of edges in a minimal length path connecting $`N`$ to the root node. The level of a caret $`C`$ is defined to be the level of the node associated to that caret.
Writing the generators $`x_n`$ for $`n>1`$ via the relators $`x_n=x_0^{(n1)}x_1x_0^{n1}`$, we relate each generator to the following rotation of a tree $`T`$. We denote the all-right tree with the appropriate number of carets by $``$. Group multiplication must be between a pair of elements, and each element corresponds to a pair of trees, so we use the tree $``$ as the positive tree corresponding to $`T`$. The product of the generator $`x_n`$ and the tree pair diagram $`(T,)`$ performs a right rotation to $`T`$ at the caret at level $`n`$ along the right arm of $`T`$. In all that follows, when we describe a generator as inducing a rotation on a single tree $`T`$ rather than on a tree pair diagram, we are forming the product with the pair $`(T,)`$ as above.
## 3. Metrics on $`F`$ and rotation distances
### 3.1. Relation to the word metric.
In , the word length with respect to the finite generating set $`\{x_0,x_1\}`$ of $`F`$ is used to compute the restricted rotation distance between a pair of trees, using techniques of Fordham . Fordham developed a method for computing the exact length of an element of $`F`$ directly from the reduced tree pair diagram representing that element.
###### Definition 3.1.
If $`T_1`$ and $`T_2`$ are trees with the same number of nodes, we define the restricted rotation distance $`d_{RR}(T_1,T_2)`$ as the minimal number of rotations required to transform $`T_1`$ to $`T_2`$, where rotations are allowed at the root and the right child of the root.
Restricted rotation distance is well-defined for any two trees with the same number of leaves, as shown in . We then obtain the following sharp bound on restricted rotation distance.
###### Theorem 3.2 (, Theorems 2 and 3).
Given two rooted binary trees $`T_1`$ and $`T_2`$ each with $`n`$ nodes, for $`n3`$, the restricted rotation distance between them satisfies $`d_{RR}(T_1,T_2)4n8`$. Furthermore, for $`n3`$, there are trees $`T_1^{}`$ and $`T_2^{}`$ with $`n`$ nodes realizing this bound; that is, with $`d_{RR}(T_1^{},T_2^{})=4n8`$.
Intermediate between the two-element generating set $`\{x_0,x_1\}`$ and the infinite generating set $`\{x_0,x_1,\mathrm{}\}`$ are other finite generating sets of the form $`\{x_0,x_{i_1},\mathrm{},x_{i_L}\}`$, where we arrange the indices of the (distinct) generators in increasing order. Analyzing the infinite generating set corresponds to allowing all rotations along the right side of the tree. Finite generating sets correspond to allowing finite collections of rotations at the root node and other nodes along the right side of the tree.
###### Definition 3.3.
Let $`𝒮=\{x_0,x_{i_1},\mathrm{},x_{i_L}\}`$ be a finite subset of the infinite generating set for $`F`$ and $`T_1`$ and $`T_2`$ be trees with the same number of leaves. We define $`d_{RRA}^𝒮(T_1,T_2)`$, the restricted right-arm rotation distance with respect to $`𝒮`$, as the minimal number of rotations required to transform $`T_1`$ to $`T_2`$, where the rotations are only allowed at levels $`0,i_1,\mathrm{},i_{L1}`$ and $`i_L`$ along the right side of the tree.
We will see below that unlike restricted rotation distance, restricted right-arm rotation distance may not be defined between all pairs of trees with the same number of nodes. We use the notation $`||_𝒮`$ to denote the word length of an element of $`F`$ with respect to the generating set $`𝒮`$. We now relate the restricted right-arm rotation distance $`d_{RRA}^𝒮(T_1,T_2)`$ to $`|(T_1,T_2)|_𝒮`$.
We consider two trees $`T_1`$ and $`T_2`$ each with $`n`$ nodes. The word length of the element $`w=(T_1,T_2)F`$ with respect to a generating set $`𝒮`$ is the length of the shortest expression for $`w`$ in that generating set. However, when considering the corresponding rotations to the tree pair diagram for $`w`$, we have no analogue of Fordham’s proof that a minimal length representative in these generators can be constructed while maintaining a constant number of nodes in each tree. Thus, it may be possible that a minimal length representative for $`w=(T_1,T_2)F`$ with respect to $`𝒮`$ includes some rotations which would require the addition of carets to the trees and are thus not permitted. Therefore, we see that the word length $`|(T_1,T_2)|_𝒮`$ provides only a lower bound on the rotation distance $`d_{RRA}^𝒮(T_1,T_2)`$, when this rotation distance is defined. If this word length corresponds to a sequence of rotations in which the number of nodes remains constant at each intermediate step, then we have computed the actual restricted right-arm rotation distance between the two trees. These cases will be addressed below.
For example, we consider the trees shown in Figure 5. The desired transformation from the top left tree $`T_1`$ drawn in solid lines to the top right tree $`T_2`$ drawn in solid lines would be given by $`x_1`$, a single right rotation at the right child of the root. But if the permitted locations for rotation are only at the root (corresponding to the generator $`x_0^{\pm 1}`$) and the right child of the right child of the root (corresponding to the generator $`x_2^{\pm 1}`$), it will be impossible to accomplish the desired transformation without adding extra nodes, and the corresponding restricted right-arm rotation distance is not defined between those two trees.
If we are permitted to add a node to the leftmost leaf of each tree, as shown with the dashed carets, to obtain the related problem of transforming the new tree $`T_1^{}`$ into $`T_2^{}`$ (drawn including the dashed carets) then the transformation would be possible using only the allowed rotations. The unreduced form of the top tree pair diagram $`(T_1^{},T_2^{})`$ drawn including the dashed caret is $`x_0x_2x_0^1`$ which reduces to $`x_1`$ in the usual manner in $`F`$, if desired. If rotations are permitted at the root and right child of the root, the rotations that transform $`T_1`$ to $`T_2`$ are exactly the same as those to perform the transformation from $`T_1^{}`$ to $`T_2^{}`$ and the added dashed caret is simply carried along intact. However, if we are only permitted to rotate at the root and right child of the right child of the root, the added caret is essential in allowing that transformation, though it does take two additional steps. We cannot transform $`T_1`$ to $`T_2`$ but we can easily transform $`T_1^{}`$ to $`T_2^{}`$ by rotating rightwards at the root, rightwards at the right child of the right child of the root, and then leftwards at the root, as pictured.
We can describe exactly when a tree $`T_1`$ can be transformed into $`T_2`$ without adding nodes with respect to a specified set of allowed rotations along the right-arm of the tree; that is, exactly when the restricted right-arm rotation distance is defined. First, we consider the case when the word in normal form associated to $`(T_1,T_2)`$ is already reduced; that is, when $`(T_1,T_2)`$ is a reduced tree pair diagram.
###### Lemma 3.4.
Let $`𝒮=\{x_0,x_{i_1},\mathrm{},x_{i_L}\}`$ be a generating set for $`F`$ with $`0<i_1<i_2<\mathrm{}<i_L`$. We consider the corresponding restricted right-arm rotation distance $`d_{RRA}^𝒮`$, where rotations are allowed at nodes at levels $`0,i_1,\mathrm{},i_L`$ on the right side of the tree. Suppose $`T_1`$ and $`T_2`$ are finite rooted binary trees with the same number of nodes forming a reduced tree pair diagram $`w=(T_1,T_2)F`$ with unique normal form given by
$$x_{i_1}^{r_1}x_{i_2}^{r_2}\mathrm{}x_{i_k}^{r_k}x_{j_l}^{s_l}\mathrm{}x_{j_2}^{s_2}x_{j_1}^{s_1}.$$
If $`x_t^{\pm 1}`$ for $`1ti_11`$ appears in this normal form, then the restricted right-arm rotation distance $`d_{RRA}^𝒮(T_1,T_2)`$ is not defined. Conversely, if no $`x_t^{\pm 1}`$ for $`1ti_11`$ appears in the unique normal form, then the restricted right-arm rotation distance $`d_{RRA}^𝒮(T_1,T_2)`$ is defined.
When $`i_1=1`$, it follows from Lemma 3.4 that $`d_{RRA}^𝒮`$ will always be defined. This includes the special case of restricted rotation distance, when $`𝒮=\{x_0,x_1\}`$.
###### Proof.
We recall that the leaf exponent of the leaf numbered $`n`$ in a tree is the length of the maximal path of left edges from leaf $`n`$ which does not reach the right side of the tree. Observe that the leaf exponent that changes as a result of a rotation at the node at level $`h`$ on the right arm of the tree corresponds to the leftmost leaf in the left subtree of the node where the rotation occurs.
First, we suppose that $`x_t^{\pm 1}`$ for $`1ti_11`$ appears in the unique normal form for $`(T_1,T_2)`$. So $`t`$ appears as the leaf number of a left leaf of a caret in either in $`T_1`$ or $`T_2`$ or possibly both. If the restricted right-arm rotation distance $`d_{RRA}^𝒮(T_1,T_2)`$ is defined, then the sequence of rotations transforming $`T_1`$ into $`T_2`$ does not change the number of nodes in the tree at any intermediate step and thus no leaves are added or removed during this process. So there is no potential renumbering of leaves, as there may be when considering equivalence classes of tree pairs in Thompson’s group $`F`$. We consider the leaf numbers whose exponents can be affected by rotations at the permitted nodes. Rotations are permitted at the root node and at levels $`i_j`$ along the right side of the tree. Rotations at the root can affect only the exponent of leaf zero, as it will be the leftmost leaf in the left subtree attached at the root. Other rotations can affect the exponents of leaves which are the leftmost leaves of left subtrees of right nodes at levels $`i_1`$ and lower. The left subtree of the right node at level $`h`$ will have leaves numbered at least $`h`$, so if $`t<i_1`$, then no rotation at level $`i_j`$ can affect the exponent of leaf $`t`$. So if there is a left leaf in the range $`1ti_11`$ present in $`T_1`$, rotations at the root cannot affect its exponent, and rotations at levels $`i_1`$ and greater cannot affect its exponent.
If leaf $`t`$ has different exponents in $`T_1`$ and $`T_2`$, since the allowed rotations cannot change its exponent, $`T_1`$ cannot be transformed into tree $`T_2`$ by the allowed rotations. If leaf $`t`$ is present in both trees with the same exponent, then since $`w`$ is in unique normal form, the exponent of leaf $`t+1`$ must also be non-zero in at least one of the trees. Moreover, leaves numbered $`t`$ and $`t+1`$ belong to the left subtree of the same node on the spine. Thus none of the allowed rotations can affect the leaf exponent of leaf $`t+1`$ as well. We iterate this argument with leaves $`t+1`$ and $`t+2`$. Thus, we see that if any $`x_t^{\pm 1}`$ with $`1ti_11`$ appears, then the two trees cannot be connected by any sequence of the allowed rotations without the addition of extra nodes.
Conversely, if $`x_t^{\pm 1}`$ for $`1ti_11`$ do not appear in the normal form, then we can rotate $`T_1`$ rightwards at the root by application of an appropriate power of $`x_0^k`$ so that all of the nontrivial subtrees then hang from the right arm of the tree at levels $`i_1`$ and greater. In Proposition 3.7 we show that the right-arm rotation distance is always defined between two trees with the same number of nodes. This allows us to finish the proof with the following argument.
We now use $`x_0,x_{i_j}`$ and conjugates of $`x_{i_j}`$ by powers of $`x_0`$ to rotate the tree to an all-right tree, just as in the infinite generating set, without adding any extra nodes. So we can transform $`T_1`$ to the all-right tree, and then from the all-right tree, we can again use $`x_0,x_{i_1}`$ and conjugates of $`x_{i_1}`$ by powers of $`x_0`$ (and possibly other $`x_{i_j}`$, if desired) to transform the all-right tree to $`T_2`$ without adding extra nodes. Thus, $`d_{RRA}^𝒮(T_1,T_2)`$ is defined. There may be more efficient ways of accomplishing this transformation but it is clear that there is at least one way of doing it without adding extra nodes, so the restricted right-arm rotation distance is defined. ∎
To understand the case where $`(T_1,T_2)`$ is an unreduced tree pair diagram, and thus we do not obtain the unique normal for the element directly from the leaf exponents, we introduce the notion of partial reduction. Partial reduction is similar to ordinary reduction except that we do not want to remove left nodes common to both trees. Stated algebraically, it means that if the normal form for the element contains instances of $`x_0`$ and $`x_0^1`$ but not $`x_1^{\pm 1}`$, we do not simplify the expression, as we do when $`x_k`$ and $`x_k^1`$ appear but not $`x_{k+1}^{\pm 1}`$ for $`k>0`$. The presence of these additional left nodes may allow us to perform rotations which would not be permitted otherwise without increasing the number of carets in the trees. This phenomenon occurs in the tree pairs shown in Figure 5, where the restricted right-arm rotation distance where rotations are permitted at the root and the right child of the right child of the root is defined between $`T_1^{}`$ and $`T_2^{}`$ but not between $`T_1`$ and $`T_2`$.
###### Definition 3.5.
A word $`w`$ in $`F`$ in normal form is partially reduced if it is of the form $`x_{i_1}^{r_1}x_{i_2}^{r_2}\mathrm{}x_{i_k}^{r_k}x_{j_l}^{s_l}\mathrm{}x_{j_2}^{s_2}x_{j_1}^{s_1}`$ with $`0i_1<\mathrm{}<i_k`$ and $`0j_1<\mathrm{}<j_l`$, with $`r_n`$ and $`s_n`$ all positive, and if we further require the partial reduction condition that for $`i>0`$, when both $`x_i`$ and $`x_i^1`$ occur, so does at least one of $`x_{i+1}`$ or $`x_{i+1}^1`$.
For any word $`w`$ in (not necessarily unique) normal form, there will be a maximal length word $`w^{}`$ satisfying the partial reduction condition which we can easily obtain using the procedure described above.
The partial reduction allows us to prove the following lemma, which describes when one given tree can be transformed into another with respect to a specified set of rotations, when the initial tree pair diagram is unreduced. The proof is identical to that of Lemma 3.4.
###### Lemma 3.6.
Let $`𝒮=\{x_0,x_{i_1},\mathrm{},x_{i_L}\}`$ be a generating set for $`F`$ with $`0<i_1<i_2<\mathrm{}<i_L`$. We consider the corresponding restricted right-arm rotation distance $`d_{RRA}^𝒮`$ where rotations are allowed at the root node and at right nodes of levels $`i_1,\mathrm{},i_L`$. Suppose $`T_1`$ and $`T_2`$ are finite rooted binary trees with the same number of nodes forming a tree pair diagram $`w=(T_1,T_2)F`$ and that $`w`$ has the partially reduced normal form of maximum length given by the word
$$w^{}=x_{i_1}^{r_1}x_{i_2}^{r_2}\mathrm{}x_{i_k}^{r_k}x_{j_l}^{s_l}\mathrm{}x_{j_2}^{s_2}x_{j_1}^{s_1}.$$
If $`x_t^{\pm 1}`$ for $`1ti_11`$ appears in this partially reduced normal form, then the restricted right-arm rotation distance $`d_{RRA}^𝒮(T_1,T_2)`$ is not defined. Conversely, if no $`x_t^{\pm 1}`$ for $`1ti_11`$ appears in this partially reduced normal form, then the restricted right-arm rotation distance $`d_{RRA}^𝒮(T_1,T_2)`$ is defined.
When rotations at all nodes along the right side of the tree are allowed, we obtain the right-arm rotation distance $`d_{RA}`$, understood by the methods of Culik and Wood and Pallo . Culik and Wood considered general rotation distance but used only rotations on the right side of the tree to show their upper bound of $`2n2`$, while Pallo intentionally restricts to only allow rotations on the right hand side of the tree. Pallo’s situation is analogous to restricted rotation distance, which considers only the rotations corresponding to the generators $`x_0`$ and $`x_1`$, because the word length once again yields the exact rotation distance.
###### Proposition 3.7.
Let $``$ denote the standard infinite generating set for $`F`$, and let $`T_1`$ and $`T_2`$ be binary trees, each with $`n`$ nodes. Then
$$d_{RA}(T_1,T_2)=|(T_1,T_2)|_{}.$$
###### Proof.
We will assume that the tree pair diagram $`(T_1,T_2)`$ is reduced. If it is not, we form the tree pair diagram $`(T_1^{},T_2^{})`$ representing the same group element which is reduced. The rotations necessary to transform $`T_1^{}`$ into $`T_2^{}`$ will also transform $`T_1`$ into $`T_2`$, since no additional rotations are necessary to alter the nodes which cause $`T_1`$ and $`T_2`$ to be unreduced. The nodes which were removed during the reduction are identical in both trees and are carried along unchanged during the rotations which transform $`T_1^{}`$ to $`T_2^{}`$. The leaf exponent method of associating the unique normal form to the tree pair diagram described above shows that each tree provides one part of the normal form; in the pair $`(T_1,T_2)`$ the tree $`T_1`$ corresponds to the terms with negative exponents and $`T_2`$ to those with positive exponents. We thus write the normal form as the product $`PN`$, where $`N`$ contains the generators with negative exponents, and $`P`$ those with positive exponents.
We see that $`N`$ is a word which rotates the tree $`T_1`$ into the all-right tree without requiring the addition of any nodes, and the subword $`P`$ is a string of generators which rotates the all-right tree into the tree $`T_2`$.
Thus we see that a lower bound for right arm rotation distance is $`|(T_1,T_2)|_{}`$, and an upper bound is given by combining the length of the strings $`P`$ and $`N`$. It follows from Lemma 2.1 that $`|(T_1,T_2)|_{}=|(T_1,)|_{}+|(T_2,)|_{}`$ where $``$ is the all-right tree with $`n`$ nodes, proving the proposition. ∎
### 3.2. Bounds on restricted rotation distances
Now that we have described the relationship between the different rotation distances and word lengths in $`F`$, we obtain numerical bounds on these rotation distances as summarized in Table 1. We note that word length of an element of $`F`$ computed with respect to a generating set of the form $`𝒮`$ given above has the potential to be much shorter than the word length of the same element computed with respect to the generating set $`\{x_0,x_1\}`$. Thus we might expect significantly smaller asymptotic upper bounds on restricted right-arm rotation distance than on restricted rotation distance. In fact, this is not the case, and the difference between the upper bounds on the two rotation distances is at most a constant.
The goal of this section is to prove that the multiplicative constant of $`4`$ in the upper bound on restricted rotation distance cannot be improved upon when we allow additional rotations along the right arm of the tree. Both of these rotation distances between two trees with $`n`$ nodes each, when defined, are bounded above by $`4n`$ minus a constant. This constant depends upon the particular finite set of rotations permitted. These bounds are shown to be sharp for restricted rotation distance in . We show below that they are asymptotically sharp for restricted right-arm rotation distance as well. While allowing additional rotations may shorten the restricted right-arm rotation distance between certain pairs of trees, asymptotically the worst-case scenario differs from restricted rotation distance only by an additive constant. One way to improve the multiplicative constant of $`4`$ is to allow rotation at an infinite collection of nodes along the right side of the tree, in which case the multiplicative constant of the bound may decrease to $`2`$.
The necessity of the constant $`4`$ is shown in two steps. We first show that the restricted right-arm rotation distance, when defined, is always bounded above by $`4n8`$, where $`n`$ is the number of nodes in either tree. We then show that there are words which approach this bound to within an additive constant.
###### Proposition 3.8.
Let $`𝒮=\{x_0,x_{i_1},x_{i_2},\mathrm{},x_{i_L}\}`$ be a generating set for $`F`$ with $`0<i_1<i_2<\mathrm{}<i_L`$, and let $`d_{RRA}^𝒮`$ be the corresponding restricted right-arm rotation distance. Let $`T_1`$ and $`T_2`$ be binary trees, each with $`n`$ nodes with $`n3`$, for which $`d_{RRA}^𝒮(T_1,T_2)`$ is defined. Then
$$d_{RRA}^𝒮(T_1,T_2)4n8.$$
###### Proof.
The case where $`i_1=1`$ is already addressed by the analysis of ordinary rotation distance, described in . We consider the element $`w=(T_1,T_2)F`$, where $`T_1`$ and $`T_2`$ are trees for which the relevant restricted right-arm rotation distance $`d_{RRA}^𝒮`$ is defined, and assume that $`i_1>1`$.
Case 1: The tree pair diagram $`(T_1,T_2)`$ is reduced.
In this case, we know that the normal form of $`w`$ contains no generators $`x_t^{\pm 1}`$ for $`1ti_11`$. In addition, this normal form can contain $`x_0`$ or $`x_0^1`$ but not both. If both $`x_0`$ and $`x_0^1`$ were present in the normal form with no $`x_1^{\pm 1}`$ generator, then the normal form could be reduced. We can assume by symmetry that the normal form for $`w`$ contains $`x_0^k`$ but no factors of $`x_0`$.
Using the correspondence between the normal form and the leaf exponents in the trees $`T_1`$ and $`T_2`$, we see that the leaves of both trees numbered from 1 through $`i_11`$ are either exposed right leaves of left nodes or exposed left leaves of right nodes. In $`T_1`$, denote the (possibly empty) subtrees of the left and right nodes by $`A_1,A_2,\mathrm{},A_n`$, where the smallest leaf number in $`A_1`$ is $`i_1`$. If $`A_1`$ is empty, by ”smallest leaf number”, we mean the number of the leaf attached to the spine of the tree in that position. Similarly, in $`T_2`$ denote these subtrees by $`B_1,B_2,\mathrm{},B_m`$, where the smallest leaf number in $`B_1`$ is $`i_1`$.
Let $`w^{}=wx_0^k`$, so that the tree pair diagram $`(S_1,T_2)`$ of $`w^{}`$ has tree $`S_1`$ containing a single left node, namely the root node, and $`i_11`$ right nodes with exposed left leaves, followed by right nodes having $`A_1,\mathrm{},A_n`$ as their left subtrees. The pair $`(S_1,T_2)`$ has the form given in Figure 6.
We consider the element $`vF`$ which has tree pair diagram $`(R_1,R_2)`$, where $`R_1`$ has a single left node, namely the root node, and the left subtree of the right node at height $`i`$ is $`A_i`$. The tree $`R_2`$ is defined analogously, using the subtrees $`B_i`$ from the original tree $`T_2`$. Since restricted rotation distance is well defined for all trees with the same number of nodes, we apply Theorem 3.2 to obtain the bound $`d_{RR}(R_1,R_2)4(n(i_11))8`$. This restricted rotation distance is realized by a string $`\alpha `$ of the generators $`\{x_0^{\pm 1},x_1^{\pm 1}\}`$.
We define a string of generators $`\alpha ^{}`$ by replacing each instance of $`x_1^{\pm 1}`$ in $`\alpha `$ with $`x_{i_1}^{\pm 1}`$. Then this string of generators exactly produces the tree pair diagram $`(S_1,T_2)`$. Since the number of nodes in each tree remains constant as each generator from $`\alpha `$ is applied to create $`(R_1,R_2)`$, the same is true as we multiply the generators in $`\alpha ^{}`$ to create $`w^{}=(S_1,T_2)`$.
Thus the restricted right-arm rotation distance with respect to $`𝒯=\{x_0,x_{i_1}\}`$ is bounded as follows:
$$d_{RRA}^𝒯(S_1,T_2)4(n(i_11))8.$$
Now we note that $`w=w^{}x_0^k`$, and since there were initially $`k+1`$ left nodes in the tree $`T_1`$, the number of nodes in each tree remains constant during these successive multiplications by $`x_0^1`$. Thus the string $`\alpha ^{}x_0^k`$ realizes the restricted rotation distance between the trees $`T_1`$ and $`T_2`$.
If $`ki_11`$, then the left nodes which are changed to right nodes under multiplication by $`x_0^k`$ do not appear in $`R_1`$ and $`R_2`$, and so are not represented in the upper bound given above. Thus, when the rotation distance is increased by $`k`$, we trivially extend the bound to
$$d_{RRA}^𝒯(T_1,T_2)4n8.$$
Since adding extra generators to the generating set, or equivalently allowing rotations at additional nodes, can only decrease the rotation distance, the upper bound still holds when we consider the entire generating set $`𝒮`$.
If $`ki_1`$, then the left nodes which are changed to right nodes by these multiplications by $`x_0^1`$ are of two types: those with exposed left leaves numbered from $`1`$ to $`i_11`$, and those with left subtrees of the form $`A_i`$. The first type of right node is not counted in the upper bound given above, and thus we increase the number of nodes in the bound by $`i_11`$ to (more than) account for the additional generators.
The right nodes of the second type, with left subtrees of the form $`A_i`$, are already counted in the bound given above. However, we recall that the word $`\alpha ^{}`$ which realizes the restricted right-arm rotation distance between $`S_1`$ and $`T_2`$, came from the word $`\alpha `$ in $`\{x_0^{\pm 1},x_1^{\pm 1}\}`$. We know from Fordham’s method of calculating word length with respect to the generating set $`\{x_0,x_1\}`$ directly from the tree pair diagram that each pair of nodes with the same infix number in each tree contributes a certain number of generators to this word length. Fordham calls this the weight of the pair of nodes. We see from Fordham’s table of weights that any pair of nodes in which one node is a right node has a weight of at most three. So using an extra generator of the form $`x_0^1`$ to transform this right node into a left node means that these nodes contribute at most four generators each to the length of the word realizing the restricted right-arm rotation distance between $`T_1`$ and $`T_2`$. We have thus shown the existence of the upper bound
$$d_{RRA}^𝒯(T_1,T_2)4n8.$$
Since rotation distance can only decrease when additional rotations are permitted, this extends immediately to show
$$d_{RRA}^𝒮(T_1,T_2)4n8.$$
Case 2: The tree pair diagram $`(T_1,T_2)`$ is not reduced.
In this case, since $`d_{RRA}^𝒮(T_1,T_2)`$ is defined, we know from Lemma 3.6 that the unreduced tree pair diagram $`(T_1,T_2)`$ yields a word $`\alpha `$ which is partially reduced and represents $`w`$. We obtain $`\alpha `$ by considering the unreduced normal form arising from $`(T_1,T_2)`$ and applying the usual reduction rules but without reducing instances of $`x_0`$ and $`x_0^1`$ with no $`x_1^{\pm 1}`$. We may also be able to partially reduce $`(T_1,T_2)`$ to correspond to this partially reduced normal form for $`w`$; in this case the number of nodes in the tree pair diagram may reduce to $`n^{}<n`$. The proof of this case is now identical to that of Case 1. This produces an upper bound of $`4n^{}8<4n8`$ on the restricted right-arm rotation distance between the two trees. ∎
We now show that the multiplicative constant of 4 is necessary for the above inequality.
###### Theorem 3.9.
Let $`𝒮=\{x_0,x_{i_1},\mathrm{},x_{i_L}\}`$. Then there exist trees $`T_1`$ and $`T_2`$, each with $`n`$ nodes, so that $`d_{RRA}^𝒮(T_1,T_2)`$ is defined, and with
$$d_{RRA}^𝒮(T_1,T_2)4n4i_L4.$$
The generating set $`𝒮`$ used in Theorem 3.9 corresponds to a series of rotations along the right side of the tree from levels $`0`$ to $`i_L`$ but does not necessarily include all rotations at levels within this range. We now enlarge our generating set to correspond to all rotations at levels $`0`$ to $`i_L`$, and work with this set $`𝒮^{}`$ in Theorem 3.10. It will be enough to use this larger set of generators and show that $`d_{RRA}^𝒮^{}(T_1,T_2)4n4i_L4`$. Thus we prove the following theorem.
###### Theorem 3.10.
Let $`𝒮^{}=\{x_0,x_1,x_2,\mathrm{},x_{i_L}\}`$, with $`m=i_L`$. Then there exist trees $`T_1`$ and $`T_2`$, each with $`n`$ nodes, so that $`d_{RRA}^𝒮^{}(T_1,T_2)`$ is defined, and with
$$d_{RRA}^𝒮^{}(T_1,T_2)4n4m4.$$
The elements we will use to prove this theorem have normal form
$$x_{m+2}x_{m+3}\mathrm{}x_{n2}x_{n3}^1x_{n4}^1\mathrm{}x_{m+1}^1$$
and tree pair diagram which we denote $`(T_1,T_2)`$. These elements are shown in Figure 7.
Fordham’s method for computing exact word length is only valid for the generating set $`\{x_0,x_1\}`$, so we bound the lengths of these elements indirectly in a series of lemmas which analyze how many generators are needed to change a “deeply buried” part of the tree.
We write $`[i,i+1]`$ if leaves $`i`$ and $`i+1`$ form a sibling pair. Performing a rotation corresponding to the generator $`x_n`$ on a tree $`T`$ is equivalent to taking the product of $`x_n`$ with the tree pair diagram $`(T,)`$, where $``$ is the tree consisting only of the root node and a series of right nodes. By analyzing the effect of a rotation at a right node on a tree $`T`$, we see that there are only three configurations of $`T`$ which allow a sibling pair to be created or destroyed. These are presented in Figure 8, where capital letters refer to nonempty subtrees of $`T`$ and lower case letters denote leaf numbers. Right rotation at the appropriate node $`N`$ along the right side of the tree has the following effect on the sibling pairs.
* The pair $`[a,b]`$ is destroyed and the pair $`[b,c]`$ is created.
* The pair $`[a,b]`$ is destroyed.
* The pair $`[b,c]`$ is created.
We can similarly consider left rotation at the node $`N`$, in which case we refer to Figure 9. Left rotation at node $`N`$ along the right side of the tree has the following effect on the sibling pairs.
* The pair $`[a,b]`$ is created and the pair $`[b,c]`$ is destroyed.
* The pair $`[a,b]`$ is created.
* The pair $`[b,c]`$ is destroyed.
From these observations, we can see immediately at which nodes it is possible to create and destroy sibling pairs with a set of rotations.
###### Lemma 3.11.
Suppose a tree $`T^{}`$ is obtained from a tree $`T`$ by applying a right rotation at a node $`N`$ at level $`n`$ on the right side of $`T`$.
1. If leaves $`m`$ and $`m+1`$ are siblings in $`T`$ and are not siblings in $`T^{}`$, then leaves $`m`$ and $`m+1`$ are the leaves of an exposed node whose parent is node $`N`$.
2. If we have the sibling pair $`[m,m+1]`$ in $`T^{}`$ but not in $`T`$, then $`[m,m+1]`$ must be the rightmost node in $`T^{}`$ and $`m`$ must be a leaf in $`T`$ whose parent is the node at level $`n`$ in $`T`$.
Similarly, the opposite conditions hold for left rotations.
We note that when $`N`$ is the root node, the only sibling pairs that might be affected by rotation at $`N`$ consist of the first two and the last two leaves in the tree.
When we consider the trees $`T_1`$ and $`T_2`$ in Figure 7 corresponding to the reduced tree pair diagram representing the element $`x_{m+2}x_{m+3}\mathrm{}x_{n2}x_{n3}^1x_{n4}^1\mathrm{}x_{m+1}^1`$, we see that in $`T_1`$, leaves $`n3`$ and $`n2`$ are siblings and in $`T_2`$, leaves $`n2`$ and $`n1`$ are siblings.
Now we consider the minimal number of rotations needed to change the sibling pairings from $`[n3,n2]`$ to $`[n2,n1]`$, expressed as a word $`w=g_1g_2\mathrm{}g_l`$, where each $`g_i𝒮^{}`$. We will need to first destroy the sibling pair $`[n3,n2]`$ and subsequently create the sibling pair $`[n2,n1]`$. The exposed nodes with siblings $`[n3,n2]`$ and $`[n2,n1]`$ are “deeply buried” in the sense that many rotations are required to affect those nodes and thus those leaf pairings. We measure this depth more precisely with the following definition.
###### Definition 3.12.
Let $`c`$ be an exposed caret, and $`\alpha _c`$ the minimal path from the node of $`c`$ to the spine of the tree. The node which is the endpoint of $`\alpha _c`$ lying on the spine of the tree is called the spinal ancestor of $`c`$.
Define $`G(c)=(r,s)`$ where $`r`$ is the number of edges in the path $`\alpha _c`$ and $`s`$ is the level in the tree of the spinal ancestor of $`c`$.
Note that the spinal ancestor of $`c`$ can be either a left or right caret. For example, in the tree $`T_1`$ for $`w`$ given in Figure 7, we consider $`c`$ as the caret with exposed leaves $`[n3,n2]`$. The length of the path $`\alpha _c`$ is $`nm3`$, so $`G(c)=(nm3,m+1)`$. If $`c^{}`$ is the caret in $`T_2`$ with exposed leaves $`[n2,n1]`$, then $`G(c)=(nm3,m+2)`$.
In Table 2 below, we summarize the changes in $`G(c)=(r,s)`$ when a single rotation at level $`k>0`$ is performed on the right arm of the tree containing the exposed caret $`c`$. We label each non-spinal node along the path $`\alpha _c`$, beginning with the one closest to the spine, as follows. We give the node the label $`R`$ if it belongs to a caret which is the right child of its parent, and $`L`$ if that node belongs to a caret which is the left child of its parent. If the first spinal ancestor of $`c`$ is on the right side of the tree, then $`\alpha _c`$ must begin with the label $`L`$. Analogously, if the first spinal ancestor of $`c`$ is on the left side of the tree, then $`\alpha _c`$ must begin with the label $`R`$.
The change in $`G(c)`$ under a single rotation is governed by two factors:
1. the relative positions of the levels $`k`$ and $`s`$, and
2. the first two labels along the path $`\alpha _c`$.
The following tables summarize the changes in $`G(c)=(r,s)`$ when different rotations are performed at level $`k`$ along the right side of the tree, so we are assuming that the spinal ancestor of $`c`$ lies on the right arm of the tree. If the spinal ancestor of $`c`$ lies on the left arm of the tree, then $`G(c)`$ is unaffected by a rotation along the right arm of the tree.
To give a lower bound on the restricted right-arm distance between the trees $`T_1`$ and $`T_2`$ which form the tree pair diagram for $`w`$, we consider the sibling pairings involving leaf $`n2`$. Leaf $`n2`$ is paired with leaf $`n3`$ in $`T_1`$ and paired with leaf $`n1`$ in $`T_2`$. In the following lemmas we bound the minimal number of rotations necessary to split these sibling pairs. Combined, these estimates yield the desired lower bound. The main tool is the ordered pair $`G(c)`$, which allows us to track the position of the exposed caret containing leaf $`n2`$ relative to the right arm of the tree.
###### Lemma 3.13.
Let $`w=(T_1,T_2)F`$ have normal form
$$x_{m+2}x_{m+3}\mathrm{}x_{n2}x_{n3}^1x_{n4}^1\mathrm{}x_{m+1}^1$$
where $`n>m+4`$. The tree $`T^{}`$ resulting from the application of at most $`2n2m3`$ rotations at locations at levels $`0`$ to $`m`$ along the right side of the tree to $`T_1`$ will contain the sibling pair $`[n3,n2]`$.
###### Proof.
We note that by Lemma 3.11, sibling pairs can be destroyed by a single rotation only when they are connected by one left edge to a node on the right side of the tree at level $`m`$ or less, or are the leaves of the rightmost caret in the tree. The exposed caret $`c`$ in $`T_1`$ with leaves $`n3`$ and $`n2`$ cannot be moved to be the rightmost caret of the tree, as all rotations preserve the natural infix order on the carets.
Thus, until the caret $`c`$ is connected by a single left edge to the right side of the tree at level at most $`m`$, and the correct rotation is performed to separate them, leaves $`n3`$ and $`n2`$ will remain a sibling pair. We use the ordered pair $`G(c)`$ to monitor the position of $`c`$ relative to the right arm of the tree while performing a series of rotations. The leaves $`n3`$ and $`n2`$ will remain sibling pairs until $`G(c)=(1,l)`$ for some $`lm`$, when a single rotation can be performed to separate these leaves.
We consider the sequence of trees $`S_0=T_1`$, $`S_1,S_2,\mathrm{},S_t`$ resulting from performing a series of $`t`$ rotations corresponding to a sequence of $`t`$ generators $`g_1g_2\mathrm{}g_t`$. Each $`S_i`$ is the result of applying $`g_i`$ to $`S_{i1}`$. We trace the images of the caret $`c`$ through this sequence and denote its image in $`S_i`$ by $`c_i`$. While the exposed leaves of each $`c_i`$ have the same leaf numbers in $`S_i`$, the entries in $`G(c_i)`$ may change as a result of each rotation. The possible changes in $`G(c_i)`$ are summarized in Table 2.
We note that it is possible to move the caret $`c`$ so that its spinal ancestor is on the left arm of the tree. Since rotations are not allowed along the left arm of the tree, caret $`c`$ must be returned to a subtree of a right node before the sibling pair $`[n3,n2]`$ can be split. This will not happen in any minimal length transformation.
We know that initially, $`G(c)=(nm3,m+1)`$, and the sibling pair $`[n3,n2]`$ is not destroyed until after $`G(c_i)=(1,l)`$, for some appropriate $`lm`$. The path $`\alpha _c`$ has labels $`LRRR\mathrm{}R`$. These labels remain unchanged as rotations are performed along the right arm of the tree. As the length of the path $`\alpha _{c_i}`$ is decreased, labels are removed sequentially from the beginning of this list, but the remaining labels are never changed by rotations along the right arm of the tree.
We see from Table 2 that the rotations which reduce the first coordinate of $`G(c_i)`$ fall into two types.
1. Rotations which decrease the first coordinate and increase the second.
2. Rotations which decrease the first coordinate and leave the second unchanged. These can only happen when the initial two labels of $`\alpha _{c_i}`$ are $`LL`$.
Along the initial path $`\alpha _c`$, there are no adjacent left labels. To create such a pair of labels, in order to perform a reduction of the first coordinate of $`G(c_i)`$ but leave the second coordinate unchanged, requires the creation of at least one additional caret with label $`L`$. While this is easily accomplished, its creation increases the first coordinate of $`G(c_i)`$. The resulting rotation then decreases this coordinate with no net change in $`G(c_i)`$. So we see that there will never be any rotations of this second type in a minimal sequence of rotations that splits the sibling pair $`[n3,n2]`$ in $`T_1`$.
To reduce the first coordinate of $`G(c)=(nm3,m+1)`$ to $`1`$, we will need at least $`nm2`$ rotations, all of the first type listed above. Each reduction will increase the second coordinate of $`G(c)`$ by one. We will need at least $`nm2`$ additional rotations to reduce the second coordinate back to its starting value, without changing the first coordinate. We then must perform at least one additional rotation to decrease the second coordinate to $`m`$ before the sibling pair in question can be split. This gives a minimum of $`2n2m3`$ rotations before the sibling pair $`[n1,n]`$ can be destroyed. ∎
We make an analogous argument in the lemma below to bound the minimal number of rotations necessary to split the sibling pair $`[n2,n1]`$ in the tree $`T_2`$.
###### Lemma 3.14.
Let $`w=(T_1,T_2)F`$ have normal form
$$x_{m+2}x_{m+3}\mathrm{}x_{n2}x_{n3}^1x_{n2}^1\mathrm{}x_{m+1}^1$$
where $`n>m+4`$. The tree $`T^{}`$ resulting from the application of at most $`2n2m2`$ rotations at locations at levels $`0`$ to $`m`$ along the right side of the tree to $`T_2`$ will contain the sibling pair $`[n2,n1]`$.
###### Proof.
We note that in this case, when $`c`$ is the caret with exposed leaves numbered $`n`$ and $`n+1`$, we have $`G(c)=(nm3,m+2)`$ and to reduce $`G(c)`$ to $`(0,l)`$ with $`lm`$ will take at least $`(nm2)+(nm2)+2=2n2m2`$ rotations by the same analysis as in Lemma 3.13. ∎
We combine these lemmas to prove Theorem 3.10.
Proof of Theorem 3.10. We consider the reduced tree pair diagram $`(T_1,T_2)`$ corresponding to the element $`w=x_{m+2}x_{m+3}\mathrm{}x_{n2}x_{n3}^1x_{n2}^1\mathrm{}x_{m+1}^1F`$, as in the lemmas above. If the restricted right-arm rotation distance between these two trees is $`d`$, then we consider the sequence of trees $`S_0=T_1`$, $`S_1,S_2,\mathrm{},S_d=T_2`$ resulting from performing that series of $`d`$ rotations to $`T_1`$ to get $`T_2`$. Lemma 3.13 shows that any application of $`2n2m3`$ allowed rotations to $`T_1`$ will still result in a tree with leaves $`n3`$ and $`n2`$ still paired, so in trees $`S_i`$ with $`0i2n2m3`$ leaf $`n2`$ must be paired with $`n3`$. We consider the tail end of that sequence, and find that Lemma 3.14 shows that in trees $`S_i`$ with $`d(2n2m2)id`$ leaf $`n2`$ must be paired with $`n1`$. Since it will take at least one additional rotation to change the pairing of leaf $`n2`$ from $`n3`$ to $`n1`$, the restricted right arm rotation distance between the two trees is at least $`4n4m4`$. ∎
Theorem 3.10 gives a family of pairs of trees with $`n`$ nodes satisfying a lower bound on restricted right-arm rotation distance with respect to a generating set $`𝒮^{}`$ which includes all generators from $`x_0`$ to $`x_m`$. Restricting the generating set to a subset $`𝒮`$ of $`𝒮^{}`$ which includes $`x_0`$ can only increase the restricted right-arm rotation distance between two trees, or cause it to be undefined. In the case of the words used in the proof of Theorem 3.10, it follows from Lemma 3.6 that the restricted right-arm rotation distance will still be defined when the generating set is further restricted. This follows because the smallest index in the normal form of the words used exceeds the highest level along the right arm of the tree where rotation is allowed. Thus we have proven Theorem 3.9 as well.
## 4. Bounding right-arm rotation distance
The original arguments of Culik and Wood which give a bound on ordinary rotation distance apply to right-arm rotation distance as well. Their argument is that any binary tree $`T`$ with $`n`$ nodes can be transformed to or from the all-right tree with $`n`$ nodes by no more than $`n1`$ rotations, all of which can be chosen to lie on the right arm of the tree. Thus, the right-arm rotation distance between two trees $`T_1`$ and $`T_2`$ each with $`n`$ nodes is no more than $`2n2`$, as we can transform $`T_1`$ to the all-right tree and from there transform it to $`T_2`$. While this bound is not optimal for the original rotation distance, we show that it is optimal for right-arm rotation distance.
###### Theorem 4.1.
For each $`n3`$, there are rooted binary trees $`T_1`$ and $`T_2`$ each with $`n`$ nodes so that the right-arm rotation distance between them satisfies $`d_{RA}(T_1,T_2)=2n2`$.
###### Proof.
To prove this we consider the elements of $`F`$ with normal form
$$x_0x_1x_2x_3\mathrm{}x_{n2}x_{n3}^1x_{n4}^1\mathrm{}x_1^1x_0^2$$
pictured in Figure 10, which have $`n`$ nodes and have word length $`2n2`$ with respect to the infinite generating set for $`F`$. It follows from Proposition 3.7 that this is also the right-arm rotation distance between the two trees. ∎
## 5. Left-arm and spinal rotation distances
We now consider rotation distances which include rotations at nodes along the left side of the tree, instead of or in addition to, nodes along the right side of the tree. It is clear by symmetry that restricted left-arm rotation distance, which allows rotations only at a finite collection of nodes on the left side of the tree and the root node, will satisfy the same bounds as restricted right-arm rotation distance. Similarly, left-arm rotation distance, which allows rotations at any node along the left arm of the tree, will satisfy the sharp upper bound of $`2n2`$ on trees with $`n`$ nodes, for $`n3`$.
Finally, we consider a rotation distance which allows rotations at the root node, a finite nonempty collection of nodes on the right side of the tree, and at a finite nonempty collection of nodes on the left side of the tree. Since all nodes where rotations are permitted lie on the spine of the tree, we call such a rotation distance a restricted spinal rotation distance. In terms of Thompson’s group $`F`$, left rotation at level $`n`$ on the left arm of the tree can be expressed as $`y_n=x_0^nx_1x_0^{n1}`$. So there is an extended infinite generating set for $`F`$ consisting of all $`x_n`$ and $`y_n`$ which corresponds to allowing rotations at any location on the spine. Here, we consider finite subsets of this enlarged generating set. Again, if we do not include $`x_0`$ in the generating set we consider, we generate either a subgroup isomorphic to $`F`$ or its direct square, so we restrict to the case where $`x_0`$ is included in the generating set.
###### Definition 5.1.
Let $`𝒮=\{x_0,x_{i_1},\mathrm{},x_{i_L},y_{j_1},\mathrm{},y_{j_l}\}`$ with $`i_1<i_2\mathrm{}<i_L`$ and $`j_1<j_2\mathrm{}<j_l`$ be a finite subset of the extended infinite generating set for $`F`$ and $`T_1`$ and $`T_2`$ be trees with the same number of leaves. We define $`d_{RS}^𝒮(T_1,T_2)`$, the restricted spinal rotation distance with respect to $`𝒮`$, as the minimal number of rotations required to transform $`T_1`$ to $`T_2`$, where the rotations are only allowed at levels $`0,i_1,\mathrm{},i_{L1}`$ and $`i_L`$ along the right side of the tree and at levels $`j_1,\mathrm{},j_l`$ on the left side of the tree.
Again, though allowing rotations at finitely many locations on both the right and left arms of the tree may reduce the rotation distance between some pairs of trees, we prove that the multiplicative constant of 4 in the upper bound cannot be decreased.
We first show that spinal rotation distance satisfies the same upper bound as restricted right-arm rotation distance.
###### Theorem 5.2.
Let $`𝒮=\{x_0,x_{i_1},\mathrm{},x_{i_L},y_{j_1},\mathrm{},y_{j_l}\}`$ with $`i_1<i_2\mathrm{}<i_L`$ and $`j_1<j_2\mathrm{}<j_l`$, where $`x_i`$ is a generator of $`F`$ and $`y_n=x_0^nx_1x_0^{n1}`$, and let $`d_{RS}^𝒮`$ be the corresponding spinal rotation distance. Let $`T_1`$ and $`T_2`$ be binary trees, each with $`n`$ carets with $`n3`$, for which $`d_{RS}^𝒮(T_1,T_2)`$ is defined. Then
$$d_{RS}^𝒮(T_1,T_2)4n8.$$
###### Proof.
Let $`𝒮^{}=\{x_0,x_{i_1},\mathrm{},x_{i_L}\}`$, which is also a generating set for $`F`$ in which each generator corresponds to a rotation along the right arm of the tree. Since adding additional elements to a generating set can only decrease the corresponding restricted rotation distance, we see immediately that $`d_{RS}^𝒮(T_1,T_2)d_{RS}^𝒮^{}(T_1,T_2)`$. Since $`d_{RS}^𝒮^{}(T_1,T_2)=d_{RRA}^𝒮^{}(T_1,T_2)`$, and Proposition 3.8 proves that $`d_{RRA}^𝒮^{}(T_1,T_2)4n8`$, the theorem follows. ∎
Now we show that the multiplicative coefficient of 4 is optimal in the same sense as with restricted right-arm rotation distance. For these examples, to avoid possible repeated excessive reductions in $`G(c)`$, we take words which have an exposed caret connected to the right-hand side of the tree with a path which alternates between branching right and left.
###### Theorem 5.3.
Let $`𝒮=\{x_0,x_{i_1},\mathrm{},x_{i_L},y_{j_1},\mathrm{},y_{j_l}\}`$ with $`i_1<i_2\mathrm{}<i_L`$ and $`j_1<j_2\mathrm{}<j_l`$, where $`x_i`$ is a generator of $`F`$ and $`y_n=x_0^nx_1x_0^{n1}`$. Then there exist trees $`T_1`$ and $`T_2`$ with $`n`$ nodes where $`n>\mathrm{max}\{i_L,j_l\}`$ for which $`d_{RS}^𝒮(T_1,T_2)`$ is defined that satisfy
$$d_{RS}^𝒮(T_1,T_2)4n4\mathrm{max}\{i_L,j_l\}12.$$
We again introduce a particular family of elements $`wF`$, represented by reduced tree pair diagrams $`(T_1,T_2)`$, which requires this lower bound on the restricted spinal rotation distance between $`T_1`$ and $`T_2`$. As in the proof of Theorem 3.10, we use the ordered pair $`G(c)`$ as the main tool of the proof. Since rotations are now permitted along the left arm of the tree, we must note the changes in $`G(c)`$ caused by a rotation along the left arm of the tree. As before, $`c`$ is an exposed caret, $`\alpha _c`$ is the minimal path from the node of $`c`$ to the spinal ancestor of $`c`$, each non-spinal node along $`\alpha _c`$ is given a label of $`R`$ or $`L`$, and $`G(c)=(r,s)`$. If the spinal ancestor of $`c`$ lies on the left arm of the tree, then the initial label along $`\alpha _c`$ must be $`R`$.
The following table summarizes the changes in $`G(c)`$ when a single rotation is performed along the left arm of the tree. If the spinal ancestor of $`c`$ is on the right arm of the tree, and a rotation is performed on the left arm of the tree, at level at least one, then $`G(c)`$ remains unchanged.
Rotation at the root caret, corresponding to multiplication by the generator $`x_0^{\pm 1}`$ can affect $`G(c)`$ in one of two ways.
1. If $`G(c)=(r,1)`$, then this rotation may change the arm of the tree on which the spinal ancestor of $`c`$ lies. If this is the case, then the first label along the path $`\alpha _c`$ will change from $`L`$ to $`R`$ or vice versa and $`G(c)`$ will remain $`(r,1)`$.
2. If $`G(c)(r,1)`$ then the $`r`$ coordinate remains unchanged, and the $`s`$ coordinate is changed by $`\pm 1`$, depending on the direction of the rotation.
Other rotations which do not add carets to the path $`\alpha _c`$ will not change the labels along $`\alpha _c`$. Rotations which decrease the length of $`\alpha _c`$ can only remove the initial label along the path.
The proof of Theorem 5.3 uses methods analogous to the proofs of Lemmas 3.13 and 3.14 and Theorems 3.9 and 3.10.
Proof of Theorem 5.3. First, we can assume without loss of generality that $`i_L`$ is larger than $`j_l`$, since if not, we can interchange the left and right sets of generators by taking reflections of the trees considered. Given the set $`𝒮`$, we define the level furthest from the root on the right side at which a rotation can take place as $`I=i_L`$ for convenience and consider the elements with $`m>I`$ of the form
$$w=(T_1,T_2)=x_{I+2}x_{I+3}^2x_{I+4}^2\mathrm{}x_m^2x_{m1}^2x_{m2}^2\mathrm{}x_{I+2}^2x_{I+1}^1.$$
The reduced tree pair diagram $`(T_1,T_2)`$ for $`w`$ is pictured in Figure 11. Each $`T_i`$ contains $`n=2mI`$ carets, with a deeply buried exposed caret in each tree connected to the right-hand side of the tree via a zigzag path. In $`T_1`$, we label the caret with the sibling pair $`[m1,m]`$ as $`c`$ and in $`T_2`$, we label caret with the sibling pair $`[m,m+1]`$ as $`d`$. We see that the lengths of $`\alpha _c`$ and $`\alpha _d`$ are both $`2m2I3`$, and have labels $`LRLRLRL\mathrm{}RL`$. We begin with $`G(c)=(2m2I3,I+1)`$ and $`G(d)=(2m2I3,I+2)`$ and to split each sibling pair, we need to reduce both $`G(c)`$ and $`G(d)`$ to $`(1,l)`$ for some appropriate $`l`$. As in the right-arm case described above, when counting the needed rotations, we count from the beginning of the transformation to determine the number of rotations needed to separate the sibling pair in caret $`c`$ and from the end of the transformation to determine the number of rotations to separate the pair in caret $`d`$.
As in the proof of Theorem 3.10, there are two ways that a single rotation can decrease the $`r`$ coordinate of $`G(c)`$, which we extract from Tables 2 and 3.
1. Rotations which decrease the first coordinate and increase the second.
2. Rotations which decrease the first coordinate and leave the second unchanged. These can only happen when the spinal ancestor of $`c`$ is on the right arm of the tree and initial two labels of $`\alpha _{c_i}`$ are $`LL`$ or when the spinal ancestor of $`c`$ is on the left arm of the tree and initial two labels of $`\alpha _{c_i}`$ are $`RR`$. We will call these bonus rotations.
We note that a single rotation which does not increase the $`r`$ coordinate either leaves the labels along $`\alpha _c`$ unchanged, or removes the initial label. There is no way for a single rotation to change a label in the middle of the path. Rotation at the root may change the initial label from $`R`$ to $`L`$ or vice versa, or leave all labels unchanged.
We first enlarge our generating set to $`𝒮^{}=\{x_0,x_1,x_2,\mathrm{}x_{i_L},y_1,y_2,\mathrm{},y_{j_l}\}`$, so that rotations are permitted at all nodes at levels zero through $`i_L`$ along the right side of the tree, and at all levels one through $`j_l`$ along the left side of the tree. As with Theorems 3.9 and 3.10, if we can produce trees $`T_1`$ and $`T_2`$ which require the desired lower bound on the spinal rotation distance $`d_{RS}^𝒮^{}(T_1,T_2)`$, then the same bound holds with respect to $`𝒮`$, since the removal of elements from the generating set can only increase the spinal rotation distance or cause it to be undefined. Since the elements in our example do not have any generators in their normal forms of index less than $`i_L`$, Lemma 3.6 guarantees that the relevant restricted right-arm rotation distance using just the rotations on the right side of the tree is defined. Thus the restricted spinal rotation distance using the larger set of rotations corresponding to the entire generating set $`𝒮`$ will also be defined.
We now give a lower bound on the minimal number of rotations necessary to separate the sibling pair $`[m1,m]`$ in $`T_1`$, which we can equivalently view as multiplication by a minimal sequence of generators. We recall that we must reduce $`G(c)`$ from $`(2m2I3,I+1)`$, where $`I=i_L`$, to $`(1,l)`$ for some level $`l`$ at which rotation is allowed on the correct arm of the tree. Every rotation, whether at a node on the right or left arm of the tree, which decreases the $`r`$ coordinate of $`G(c)`$, with the exception of bonus rotations, also increases the $`s`$ coordinate. The $`s`$ coordinate of $`G(c)`$ begins larger than $`i_L`$ (and thus $`j_l`$ as well), and so each non-bonus rotation which reduces the $`r`$ coordinate will require an additional rotation to counteract the corresponding increase in the $`s`$ coordinate.
We now consider the role of bonus rotations that can be used in a minimal length sequence of rotations. We note that the labels for $`\alpha _c`$ begin with $`LRLRLRL\mathrm{}`$. There is no natural occurrence in the sequence of labels of $`RR`$ or $`LL`$. In order to create an $`RR`$ or $`LL`$, we can use an application of at least one $`x_0^{\pm 1}`$ which will change the initial label but will not change $`G(c)`$. Thus each such bonus rotation would need to be accompanied by an $`x_0^{\pm 1}`$ which does not change $`G(c)`$.
Let $`\beta `$ be a minimal string of rotations which reduces $`G(c)=(2m2I3,I+1)`$ to $`(1,l)`$. We divide the rotations in $`\beta `$ into three groups:
1. $`p_1`$ non-bonus rotations which change the coordinates of $`G(c)`$,
2. $`p_2`$ bonus rotations which change the coordinates of $`G(c)`$, and
3. $`p_3`$ rotations which do not change $`G(c)`$.
The argument above showing that there is at least one $`x_0^{\pm 1}`$ accompanying each bonus rotation shows that $`p_3p_2`$. Totalling the effects of these $`p`$ rotations, we find that to reduce $`G(c)=(2m2I3,I+1)`$ to $`G(c)=(1,l)`$ with $`lI`$ will require at least $`2(2m2I4)+1=2n2I7`$ total rotations, since $`n=2mI`$. Thus, as in Lemma 3.13, any sequence of at most $`2n2I7`$ rotations applied to tree $`T_1`$ will have leaves $`m1`$ and $`m`$ as sibling pairs.
Similarly, working backward and considering tree $`T_2`$, we need to change the $`G(d)=(2m2I3,I+2)`$ to $`(1,l)`$ with $`lI`$, to be able to affect the sibling pair $`[m,m+1]`$. A similar calculation shows that any sequence of at most $`2n2I6`$ rotations applied to the tree $`T_2`$ will have leaves $`m`$ and $`m+1`$ as sibling pairs.
So we see that in the sequence of trees $`S_0=T_1,S_1,S_2,\mathrm{},S_d=T_2`$ exhibiting the transformation of $`T_1`$ into $`T_2`$ via rotation, the sibling pair $`[m1,m]`$ must be present in the first $`2n2I7`$ trees and the sibling pair $`[m,m+1]`$ must be present in the last $`2n2I6`$ trees. Including the rotation to change the sibling pairing of leaf $`m`$, we see that the spinal rotation distance must be at least $`4n4I12`$.
Note that the number of nodes in these trees $`n`$ is constructed to be $`2mI`$, so depending upon whether or not $`I`$ is even or odd, the number of nodes in the trees constructed using these examples for increasing $`m`$ will either always be even or always be odd. To obtain examples for all parity $`n`$ larger than $`I`$, we can repeat the argument above on examples with one additional caret, with normal forms
$$w^{}=(T_1,T_2)=x_{I+2}x_{I+3}^2x_{I+4}^2\mathrm{}x_m^2x_{m+1}x_m^1x_{m1}^2x_{m2}^2\mathrm{}x_{I+2}^2x_{I+1}^1$$
and find that the sibling pair $`[m,m+1]`$ must be present in the first $`2n2I5`$ steps and that the sibling pair $`[m+1,m+2]`$ must be present in the last $`2n2I6`$ steps, so in those cases the distance between the two trees must be at least $`4n4I10`$. Thus, we have that the bound holds for all $`n`$ larger than $`I`$ and thus the therorem. ∎ |
warning/0506/cond-mat0506173.html | ar5iv | text | # Mesoscopic fluctuations and intermittency in aging dynamics
## 1 Introduction
After a rapid quench of an external parameter, e.g. the temperature, many complex materials *age*, i.e. their properties slowly change with the *waiting time*, $`t_w`$, elapsed from the quench. Ever since the initial observations in polymers , evidence has accumulated that spin-glasses , type II superconductors , glasses , and soft condensed matter , among others, age in similar ways, e.g. : For observation times $`tt_w`$ physical averages are nearly constant, and autocorrelations and their conjugate linear response functions are connected by an equilibrium-like fluctuation-dissipation theorem (FDT). Conversely, for $`tt_w`$ they visibly drift and the FDT is violated. As was recently discovered, the drift happens in an *intermittent* fashion , i.e. through rare, large, and spatially heterogeneous re-arrangements, which appear as non-Gaussian tails in the probability density function (PDF) of configurational probes such as colloidal particle displacement and correlation or voltage noise fluctuations in glasses .
As aging phenomena are similar for a broad class of interactions, we seek a mesoscopic description, and assume that intermittent events, for short *quakes*, are the main source of de-correlation in non-equilibrium aging. In the framework of record dynamics , quakes are irreversible and are triggered by (energy) fluctuations of record magnitude. We show how this leads to a description of the configurational autocorrelation function, more specifically, the dependence of the shape of its PDF on $`t`$, $`t_w`$, the temperature $`T`$ and the system size $`N`$, which resembles observations for colloidal gels spin-glasses and kinetically constrained models . The model PDF is closely approximated by the Gumbel distributions widely used in the literature . The average and variance are given in close form as a function of $`t/t_w`$, $`T`$ and $`N`$. The average and the PDF, standardized to zero mean and unit variance, are in excellent agreement with spin-glass simulations. The agreement is rather poor for the variance itself, mainly because pseudo-equilibrium fluctuations are neglected.
## 2 The configuration auto-correlation PDF
In the model, a set of $`N`$ binary variables defines the system configuration. Without further loss of generality, and with an eye to the simulations of the E-A spin-glass model , we refer to these variables as spins, and to their changes of state as ‘flips’.
Configuration changes are gauged by the number of spins with different orientations at times $`t_w`$ and $`t_w+t`$. This Hamming distance, $`H`$, is simply related to the autocorrelation $`C`$ by
$$C(t_w,t)=12H(t_w,t_w+t)/N.$$
(1)
Initially, we focus on the probability $`P_H(h,t_w+t0,t_w)`$ for $`H=h`$ at time $`t_w+t`$, given $`H=0`$ at $`t_w`$, which we write as the average
$$P_H(h,t_w+t0,t_w)=\underset{s=0}{\overset{\mathrm{}}{}}P_S(s)P_H(hs)$$
(2)
over the conditional probability $`P_H(hs)`$ for $`H=h`$ given $`s`$ flips $`(s=0,1,\mathrm{}\mathrm{})`$. The weight function $`P_S(s)`$ is the probability for exactly $`s`$ flips during $`[t_w,t_w+t)`$.
Assuming for simplicity that flips occur at any site with probability $`1/N`$ leads to the master equation:
$$P_H(hs+1)=(1\frac{h1}{N})P_H(h1s)+\frac{h+1}{N}P_H(h+1s)0hN,$$
(3)
with the initial condition
$$P_H(h0)=\delta _{h,0}.$$
(4)
The equation has the formal solution
$$𝐏_H(hs)=𝒯^s𝐏_H(h0),$$
(5)
where $`𝒯`$ is the (bi-diagonal) stochastic matrix implicitly given by Eq. 3 and where the vector $`𝐏_H`$ has elements $`P_H(0s),P_H(1s)\mathrm{}P_H(Ns)`$.
The $`s`$ dependence of the conditional average and variance of $`H`$, $`\mu _H(s)`$ and $`\sigma _H^2(s)`$ can be gleaned from the moment generating function $`_{h=0}^NP_H(hs)z^h,z1`$. Omitting the details, one finds
$$\mu _H(s)=\frac{N}{2}(1(12/N)^s)$$
(6)
and
$$\sigma _H^2(s)=\frac{N}{4}(1(14/N)^s)+\frac{N^2}{4}((14/N)^s(12/N)^{2s}).$$
(7)
Since $`\sigma _H(s)\mu _H(s)`$ for large $`N`$, the r.h.s. of Eq. 2 is dominated in this limit by the term with index $`s(h)`$ implicitly given by $`h=\mu _H(s)`$. As a consequence, $`P_H`$ and $`P_S`$ acquire very similar shapes when standardized to zero average and unit variance.
To calculate $`P_S(s)`$ we need the probability that $`i`$ quakes occur between $`t_w`$ and $`t_w+t`$ and the distribution of the number of flips, for short ‘size’, of each quake. According to refs. , $`i`$ has a Poisson distribution with average
$$n_I(t_w,t)=\alpha (N)\mathrm{ln}(1+t/t_w).$$
(8)
The property $`\alpha (N)N`$, which removes the $`N`$ dependence of the exponent $`\lambda `$, (see Eq. 14), arises when the intermittent signal results from independent intermittent processes, stemming e.g. from locally thermalized clusters . The $`T`$ independence of $`\alpha `$ reflects the noise insensitivity of record dynamics , and holds within the low temperature range for which the description applies.
For the quake size, simulations of vortex dynamics yield a near exponential distribution. The same form is consistent with the (asymptotically) exponential distribution of the energy released by intermittent events. Hence, glossing over the integer nature of the sizes, we treat them as independent stochastic variables $`X_k`$, $`k=1,2\mathrm{}i`$, with the PDF
$$P_X(x)=q(T)\mathrm{exp}(q(T)x).$$
(9)
A temperature dependence of the reciprocal average quake size, $`q(T)`$, is allowed (but not required) by the theory, and is directly observable through the exponent $`\lambda `$, see Fig. 2. For typographical clarity this dependence is left understood, together with the dependence of $`n_I`$ on $`\mathrm{ln}(1+t/t_w)`$.
Considering first the conditional probability for $`S_i`$ flips for a given number $`i`$ of quakes, we note that $`S_i=_{k=1}^iX_k`$ has the gamma density
$$P_{S_i}(x)=q\frac{(qx)^{i1}}{(i1)!}\mathrm{exp}(qx).i>0.$$
(10)
Averaging the above expression over the Poisson distribution of $`i`$, and taking into account that $`\delta (s)\mathrm{exp}(n_I)`$ is the probability of no flips $`(i=0)`$, one finds
$$P_S(s)=\underset{i=1}{\overset{\mathrm{}}{}}P_{S_i}(s)\frac{n_I^i}{i!}\mathrm{exp}(n_I)+\delta (s)\mathrm{exp}(n_I).$$
(11)
With the variable $`z=(4qsn_I)^{1/2}`$, this is rewritten as
$$P_S(s)=2n_Ia\mathrm{exp}(qsn_I)I_1(z)/z+\delta (s)\mathrm{exp}(n_I),$$
(12)
where $`I_1`$ is the modified Bessel function of order one (See e.g. Abramowitz & Stegun, 9.6.10).
The $`\delta (s)`$ term in Eq. 12 will be neglected, since $`n_I`$ is large except for $`tt_w`$. With the term discarded, the standardized $`P_S(s)`$ has no $`T`$ dependence. This is seen, in brief, as follows: Using $`\mu _S=n_I/q`$ and $`\sigma _S^2=2n_I/q^2`$ for the average and variance of $`S`$, the standardized PDF, $`\sigma _SP_S((s\mu _S)/\sigma _S)`$, has no $`q`$ dependence. However, as $`q`$ carries the model full $`T`$ dependence, the latter disappears as well. Furthermore, due to its similarity with $`P_S`$, the standardized $`P_H`$ is also independent of $`T`$, as confirmed by Fig. 3,
Averaging Eqs. 6 and 7 over $`P_S(s)`$ , reintroducing the time dependence of $`n_I`$ and making the approximation $`\mathrm{ln}(12/N)2/N`$, one finds
$$\mu _H(t_w,t)=\frac{N}{2}\left(1(1+t/t_w)^{\lambda (T)}\right),$$
(13)
where
$$\lambda (T)=2\frac{\alpha (N)}{N}\frac{1}{q(T)}.$$
(14)
The average (macroscopic) form of the correlation function $`C`$, is obtained from Eqs. 13 and 1 as
$$\mu _C(t_w,t)=(1+t/t_w)^{\lambda (T)}.$$
(15)
Similar steps lead from Eq. 7 to
$$N\sigma _C^2(t_w,t)=1(1+t/t_w)^{2\lambda (T)}\left(12\lambda (T)\mathrm{ln}(1+t/t_w)\right).$$
(16)
Panel *(a)* of Fig. 1 shows, for three different values of $`n_I=\alpha (N)\mathrm{ln}(1+t/t_w)`$, the model PDF given by Eqs. 125 and 2. The two inserts show the $`n_I`$ dependence of $`\mu _C`$ and $`N\sigma _C^2`$ from a numerical evaluation of the model equations (circles) and from Eqs. 15 and 16 with the $`n_I(t,t_w)`$ dependence reintroduced (lines).
In standard form (see e.g. ref ), the one-parameter family of Gumbel densities is given by $`\mathrm{\Phi }_g(y)=\frac{bg^g}{\mathrm{\Gamma }(g)}\mathrm{exp}(b(yy_0)e^{b(yy_0)})`$, with $`b=\sqrt{(\mathrm{\Psi }^{}(g))}`$ and $`by_0=\mathrm{ln}(g)\mathrm{\Psi }(g)`$, where $`\mathrm{\Psi }`$ denotes the digamma function, $`\mathrm{\Psi }^{}`$ its derivative and $`g`$ is a real number. Gumbel densities empirically describe fluctuations in complex systems . Fig. 1,*(b)* shows that, except for $`n_I<1`$, our model PDF is closely approximated by the Gumbel PDF whose $`g`$ value minimizes the $`L_1`$ distance between the two. The left insert of the figure shows that this optimal $`g`$ value is linearly related to $`n_I`$ as $`g=0.300n_I0.185`$.
## 3 Comparison with simulation data
The (average) autocorrelation function $`\mu _C`$ for spin-glasses is well investigated . For $`t>t_w`$, $`\mu _C`$ is nearly a function of $`t/t_w`$, and can be fitted by a power-law with a temperature dependent exponent. E.g. Picco et al. found an excellent scaling using the variable $`\mathrm{ln}(t_w+t)\mathrm{ln}(t_w)`$.
The autocorrelation PDF for the E-A spin glass model nearly follows $`t/t_w`$ scaling according to Castillo et al. Chamon et al. also consider a kinetically constrained model with trivial statics. In both cases, the autocorrelation PDF, shifted to zero mean and rescaled to unit variance, is empirically fitted to time evolving Gumbel distributions, which are numerically equivalent to our model results (see Fig. 1).
For a more detailed comparison, we simulated the E-A spin-glass on a cubic lattice with $`N=16^3`$ using an event driven simulation technique , whose ‘intrinsic’ time unit corresponds, for large systems, to one Monte Carlo sweep. The data are sampled at $`20`$ time points, which are separated by a multiplicative factor of $`1.5`$, with start at $`t=100`$ and end at $`t2.2\times 10^5`$. Among these points, any ordered pair can be chosen for $`t_w`$ and $`t_w+t`$. For each set of physical parameters, $`1000`$ runs are performed with independent noise and couplings realizations, producing e.g. $`1000`$ data points for $`t/t_w=2216`$, and $`20000`$ points for $`t/t_w=1.5`$.
For a range of temperatures, the average spin-glass autocorrelation is plotted versus $`t/t_w`$ (symbols). Deviations from $`t/t_w`$ scaling appear for approx. $`t/t_w<4`$ and $`T>0.5`$, as seen in the left panel of Fig. 2 from the poor data collapse. Away from this parameter region, the data are well described by Eq.15 (full line). The $`T`$ dependence of the exponent is shown in the insert (circles), where the fit $`\lambda (T)=0.25T/T_g`$ is also shown (full line). $`T_g=0.97`$ is the critical temperature of the model . The same linear form, and with a similar slope coefficient, is found numerically in Kisker et al. . They, however, introduce a small $`t_w`$ dependence of $`\lambda `$, which is beyond the present model.
Summarizing, for low $`T`$ and not too small $`t/t_w`$, the model is able to describe the time dependences of the average autocorrelation with no free parameters. Furthermore, Eq. 14 links $`\lambda (T)`$ to a linear temperature increase of the average quake size.
To improve the statistics of the variance and PDF data in the $`t/t_w`$ scaling region where a comparison with the model is most interesting, we estimate the variance and its error-bar as the mean value and standard deviation over the set of variances for all pairs $`t_w,t`$ having the same ratio $`t/t_w`$. Similarly, the empirical frequencies of the $`C`$ values are calculated based on all data with the same $`t/t_w`$.
The value of $`N`$ stands in the model for an (unknown) number of thermally active spins, and appears in the autocorrelation variance, which vanishes linearly with $`1/N`$. This leads to an undetermined $`T`$ dependent scale factor, $`f_1`$, between model and simulation variance. Secondly, and more importantly, the data cannot be fitted without a second, ad hoc, offset parameter $`f_2`$, likely because the de-correlating effect of the pseudo-equilibrium fluctuations is altogether neglected. Hence, with respect to the variance, the theory only provides a qualitative description, which is captured by the empirical formula $`N\sigma _{emp.}^2(t_w,t)=f_1(N\sigma _C^2(t_w,t)+f_2)`$. For completeness, the latter is plotted (lines) in the right panel of Fig. 2 together with the simulation data with error bars. The parameter $`f_1`$ increases with $`T`$ within the range $`120`$, and $`f_2`$ remains close to $`1/10`$,
By contrast, an excellent agreement between predictions and data for the standardized PDF is achieved by a simple adjustement of the vertical scale of the latter. This scale, which is undetermined from the outset, is fitted to the properly normalized model PDF. The centering and rescaling are done using the data average and standard deviation. The results are plotted with $`1\sigma `$ error-bars in Fig. 3.
The three panels of the figure correspond, from left to right, to $`t/t_w=2.3`$, $`7.6`$ and $`25.6`$. In each panel, the data shown are for $`T=0.15,0.25,0.35`$ and $`0.5`$. Their collapse confirms the anticipated $`T`$ independence of the standardized autocorrelation PDF. The model predictions (full lines) contain one parameter, $`\alpha (N)`$, whence it is possible to determine one value of $`n_I`$ by data fitting. This was done for $`(t/t_w=25.6)`$—in the rightmost panel—yielding $`n_I=78`$ . Eq. 8 then gives $`\alpha =24`$, whence $`n_I=20`$ and $`49`$ for $`t_w=2.3`$ and $`7.6`$ respectively.
## 4 Conclusion
Based on the record dynamics description of intermittency the model develops the aging properties of the configuration autocorrelation after a deep quench. Its predictions for the average autocorrelation and the standardized PDF are accurate at low temperatures and for $`t>t_w`$. Together with allied efforts , the present results support the view that record-sized fluctuations are important for aging in metastable glassy systems.
## 5 Acknowledgments
Support from the Danish Natural Sciences Research Council is gratefully acknowledged. |
warning/0506/astro-ph0506311.html | ar5iv | text | # Gl 86 B: a white dwarf orbits an exoplanet host star
## 1 Introduction
Gliese 86 (thereafter Gl 86) is a K1 dwarf located at a distance of 10.9$`\pm `$0.08 pc (Hipparcos, ESA 1997). Queloz et al. (2000) found a 15.8 day periodical variation of its radial velocity. Because Gl 86 does not show any chromospheric activity or photometric variability they concluded that the variation of the radial velocity is induced by an exoplanet with a minimum mass of 4 $`M_{\mathrm{Jup}}`$ on an almost circular orbit (e=0.046, a=0.11 AU). Furthermore they reported a long-term linear trend in the radial velocity data (0.5 m s<sup>-1</sup> d<sup>-1</sup>) observed over a time span of 20 years in the CORAVEL program and also confirmed by CORALIE measurements (0.36 m s<sup>-1</sup> d<sup>-1</sup>). This is a clear signature of a further stellar companion in the Gl 86 system with a period longer than 100 yr (a$``$20 AU).
Els et al. (2001) found a faint companion 2 arcsec east of Gl 86 which clearly shares common proper motion. They also obtained near infrared photometry (J=14.7$`\pm `$0.2, H=14.4$`\pm `$0.2, and K=13.7$`\pm `$0.2) and from the derived color (J-K$``$1) they concluded that Gl 86 B must be substellar with a spectral type between late L to early T. However a substellar companion (m$``$78 $`M_{\mathrm{Jup}}`$) cannot explain the detected long term trend in the radial velocity of the exoplanet host star.
At the begin of 2005 we started a search for faint substellar companions of all exoplanet host stars known to harbor several substellar companions, i.e. several planets or as in the case of Gl 86 one planet and a reported brown dwarf companion (see Els et al. 2001). Our goal is to detect additional companions. We carry out our observations with NACO/VLT and its new simultaneous differential imaging device SDI which is particular sensitive for faint cool substellar companions exhibiting strong methane absorption features (T dwarfs), yielding a much higher contrast than standard AO imaging with NACO alone. All our targets are main-sequence stars with ages in the range of 1 to 10 Gyrs. From theoretical models (Baraffe et al. 2003) we expected that most of the substellar companions of our targets are T dwarfs cooler 1400 K, which are detectable with NACO/SDI close to the primary stars.
In section 2 we describe our SDI observations, the obtained astrometry and photometry in detail. With the achieved high contrast SDI imaging we furthermore proof that their is no further companion in the system which could induce the reported long term trend in the radial velocity of Gl 86 A (Queloz et al. 2000). Spectra of Gl 86 B were taken with NACO and were retrieved from the ESO public archive, yielding the surprising result that this companion is a white dwarf companion. We presented the spectroscopy in section 3. Finally section 4 summarizes and discusses all results of this letter.
## 2 NACO/SDI Observations
Detection of faint objects close to a much brighter source is the main challenge in the direct imaging search for substellar companions (brown dwarfs or planets) of stars. Within 1 arcsec of the bright central source the field is filled with speckles which are residuals from the none perfect adaptive optics correction of the incoming disturbed wavefront. The achievable signal to noise in this speckle-noise limited region does not increases with integration time, hence only a subtraction of the speckle pattern can improve the detection limit close to a bright source.
At the ESO/VLT the simultaneous differential imager (SDI) is offered for the AO-system NACO (Lenzen et al. 2004). Therein a double Wollaston prism splits the beam, coming from the AO system, into 4 beams which pass then through three different narrow band filters with central wavelengths 1.575, 1.600 and 1.625 $`\mu `$m and a bandwidth of 25 nm bandwidth each. The resulting speckle pattern of the four images is almost identical. Since cool (T$``$1400 K) objects exhibit a strong methane absorption band at 1.62 $`\mu `$m they appear much fainter in the 1.625$`\mu `$m filter than in the 1.575$`\mu `$m filter while the bright star and therefore also the speckle pattern has roughly equal brightness in all images. Subtracting the 1.625 $`\mu `$m SDI image from the image taken through the 1.575 $`\mu `$m SDI filter will effectively cancel out the speckle pattern of the star while the signal from the cool companion remains (see e.g. Biller et al. 2004).
We observed Gl 86 on 12. Jan 2005 with SDI (see Fig. 1). Due to the small SDI field of view (5x5 arcsec, tilted by 45 ), jittering cannot be used for background substraction. However we apply a small jitter with a box width of only 0.2 arcsec (10 SDI pixel) to correct for bad pixels. Per jitter cycle 5 object frames are taken, each is the average of 60 exposures of 2 s each. We always adjust the individual integration time so that only the central 9 SDI pixel of the primary point spread function are saturated. This improves the detection limit for faint companions at larger separations to the primary. At the end of each jitter cycle a sky-frame (10 arcsec offset from the target in Ra and Dec) is taken which is then subtracted from the 5 object frames to cancel out the bright infrared sky background. The sky-frames are taken in the same way as the the target-frames, i.e. 60 times 2 s integrations. The jitter cycle is repeated 4 times which yields a total on source integration time of 40 min. To distinguish between faint companions and any residual speckles we observe each target at the detector position angles 0 and 33 , i.e. 2x40 min integration time in total.
For infrared data-reduction we use the ESO Eclipse package. After flatfielding, we extract the four SDI quadrants (left-top 1.6 $`\mu `$m, right-top 1.575 $`\mu `$m, bottom both quadrants 1.625 $`\mu `$m) and apply image-registration, shifting, and final averaging on all individual frames. Because the radial position of the speckles is proportional to the wavelength all SDI images must be spatially rescaled. Finally the images are aligned, and their flux is adjusted to eliminate any differences in the quantum efficiency. The 1.575 $`\mu `$m and the 1.625 $`\mu `$m images (top-right and bottom-left quadrant) follow the same path through the SDI instrument, i.e. they provide very similar speckle pattern and yield the best contrast for methane rich companions. We therefore subtract these two SDI images. The difference frame taken at position angle 0 is then subtracted from the difference frame taken at position angle 33 . To filter out low spatial frequencies all frames are unsharp masked (see Fig. 2). The resulting difference frame compared to the image taken at 1.575 $`\mu `$m is shown in Fig. 2. The speckle pattern is effectively subtracted and a detection limit of 12.8 mag is reached at a separation of 0.5 arcsec.
For astrometrical calibration of the SDI camera we observed the binary HIP 9487. This system is listed in the Hipparcos catalogue ($`\rho `$=1.876$`\pm `$0.001 arcsec and $`\theta `$=278.500$`\pm `$0.001 at epoch 1991.25) and accurate astrometry is available for both components. Therefore the binary separation and position angle can be computed for the given observing epoch (12 Jan. 2005). We derive a pixelscale of 17.210$`\pm `$0.087 mas per pixel. The true north is slightly rotated to the east by 0.33$`\pm `$0.24 .
Due to the large proper and parallactic motion of Gl 86 ($`\mu _\alpha cos(\delta )`$ = 2092.59$`\pm `$0.56 mas/yr, $`\mu _\delta `$ = 654.49$`\pm `$0.49 mas/yr, and $`\pi `$ = 91.63$`\pm `$0.61 mas, Hipparcos ESA 1997) Els et al. (2001) already proved that Gl 86 A and B form a common proper motion pair, a result which is significantly (263$`\sigma `$ in separation and 279$`\sigma `$ in position angle) confirmed with our SDI observations (see Fig. 3). Gl 86 B is located 1.93$`\pm `$0.01 arcsec at a position angle of 104.0$`\pm `$0.3 . Fig. 3 illustrates the astrometry from Els et al. (2001) and our astrometric results. The expected change of separation and position angle is calculated assuming that only Gl 86 A is moving and Gl 86 B is an unrelated background star (see dashed lines in Fig. 3).
By comparing our astrometry with data from Els et al. (2001) we find a significant change in position angle of $`15.5\pm `$0.5 and +0.196$`\pm `$0.024 arcsec in separation for the given epoch difference. This is a clear evidence for orbital motion. The expected orbital motion for a companion at a separation of 21 AU is shown with dotted lines in Fig. 3 (0.090 arcsec/yr and 4.2 /yr, see also section 4).
With the J,H and K infrared magnitudes of Gl 86 B given by Els et al. (2001) and the Hipparcos parallax of Gl 86 A we derive the absolute magnitudes of Gl 86 B: $`M_\mathrm{J}`$=14.5$`\pm `$0.2 mag, $`M_\mathrm{H}`$=14.2$`\pm `$0.2 mag, $`M_\mathrm{K}`$=13.5$`\pm `$0.2 mag. If we assume that Gl 86 B is a brown dwarf companion these absolute magnitudes are consistent with a spectral type L7 to T5 (see Vrba et al. 2004). In our SDI images Gl 86 B is comparable bright in all three filter. We measure flux-ratios $`F{}_{1.575\mu \mathrm{m}}{}^{}/\mathrm{F}_{1.625\mu \mathrm{m}}`$ = 1.11$`\pm `$0.04. This clearly rules out spectral types later than T3, because for theses spectral types $`F{}_{1.575\mu \mathrm{m}}{}^{}/\mathrm{F}{}_{1.625\mu \mathrm{m}}{}^{}>1.3`$ due to methane absorption.
## 3 NACO Spectroscopy
Infrared spectra of Gl 86 B were taken in ESO observing program 070.C-0173(A) (extracted by us from the public archive). 8 spectra (120 s each) were taken in spectroscopic mode S27-SK-3 using the 86 mas slit, which yields a resolving power $`\lambda /\mathrm{\Delta }\lambda `$ = 1400. For background subtraction, 8 arcsec nodding was applied along the slit. To avoid that the bright primary is located on or close to the slit (saturation) the slit was orientated perpendicular to the position vector of Gl 86 B relative to the primary. All images are flatfielded (lamp flats) and the spectra are extracted, wavelength calibrated (with Argon lines) and finally averaged, with standard routines in IRAF. The resulting spectrum is flux calibrated with the photometric standard HIP 020677 (G2V). Figure 4 shows the flux calibrated NACO spectrum of Gl 86 B, together with spectra of M1, L5, and T5 dwarfs from Cushing et al. (2005). For the given spectral range the achieved signal to noise is $``$40. Neither characteristic molecular nor atomic absorption features are visible and the continuum is clearly different to those of L and T dwarfs. The spectrum is even steeper than the M1V reference spectrum which points to an effective temperature hotter than 3700 K. However, in the K-band, the gradient of the continuum is only slightly varying for effective temperatures higher than 4000 K, so that we can only give a low temperature limit for Gl 86 B.
Therefore we conclude that Gl 86 B is a cool white dwarf companion (see Fig.5). Due to its high surface gravity ($`log(g)>`$7) all absorption line are strongly broadened and therefore hardly detectable (see e.g. Br $`\gamma `$ absorption line at 2.17 $`\mu `$m, in Dobbie et al. 2005). We compare the absolute infrared photometry of Gl 86 B with data from Bergeron et al. (2000), who carried out a detailed photometric and spectroscopic analysis of cool white dwarfs, and derive an effective temperature of 5000$`\pm `$500 K.
## 4 Conclusions
We confirmed common proper motion of Gl 86 B to its primary and detected its orbital motion. This is a clear evidence that Gl 86 B is a bound companion of the exoplanet host star with a projected separation of 21 AU (1.93$`\pm `$0.01 arcsec). With the achieved high contrast SDI detection limit we can rule out any further stellar companions beyond 0.1 arcsec up to 2.1 arcsec, i.e. 1 AU to 23 AU in projected separation. T dwarf companions (T$`<`$1400 K) can be detected beyond 0.2. arcsecs (2 AU) and we are sensitive for faint substellar companions down to 35 $`M_{Jup}`$. The NACO spectrum of Gl 86 B is clearly different to the expected spectral type between L5 and T5 derived from infrared photometry. Gl 86 B is faint but its spectrum implies that it is even hotter than 3700 K (M1). Therefore we conclude that Gl 86 B is a cool white dwarf companion.
Queloz et al. (2000) report a long-term trend in the radial-velocity data of Gl 86 A (0.5 m s<sup>-1</sup> d<sup>-1</sup>). This is a clear hint on a further stellar companion in the Gl 86 system. We compute that with a separation of 21 AU, the mass of Gl 86 B must be 0.55 $`M_{\mathrm{}}`$ to induces this trend in the radial velocity. The derived mass is well consistent with a white dwarf companion. Santos et al. (2004) determined the mass of Gl 86 A to be 0.7 $`M_{\mathrm{}}`$, hence the total mass of the system is 1.25 $`M_{\mathrm{}}`$ and with a binary semi-major axis of 21 AU, the expected orbital time is 86 yr. This yields an orbital motion of 0.090 arcsec/yr and 4.2 /yr, assuming a circular orbit, which is also consistent with the measured orbital motion (see Fig. 3).
If we assume that both components of the binary have the same age, the white dwarf progenitor must be more massive than Gl 86 A (0.7 $`M_{\mathrm{}}`$) to be observable today as a white dwarf companion. According to Weidemann (2000) a 0.55 $`M_{\mathrm{}}`$ white dwarf is the remnant of 1 $`M_{\mathrm{}}`$ star. With the white dwarf models presented by Richer et al. (2000) and the derived effective temperature of Gl 86 B (5000$`\pm `$500 K) we can approximate a cooling timescale between 3 and 6 Gyr, i.e. the binary system should be 13 to 16 Gyrs old. For more massive white dwarf progenitors (2-4 $`M_{\mathrm{}}`$) the system age ranges between 2 to 8 Gyr.
Gl 86 B is the first confirmed white dwarf companion to an exoplanet host star. Theoretically, planets may or may not survive the red giant and asymptotic giant branch phases of stellar evolution. Planets which are located outside the red giant’s envelope, which reaches about a few hundred solar radii will survive. Closer companions will be either destroyed or migrate inward and become a close companion to the white dwarf remnant. According to Burleigh et al. (2002) it seems likely that distant planets (a$`>`$5 AU) survive the late stages of stellar evolution of main-sequence stars with masses in the range between 1 and 8 $`M_{\mathrm{}}`$, i.e. all white dwarf progenitors. In particular in the Gl 86 system the separation between the white dwarf and the exoplanet (21 AU) is large enough that it seems very well possible that the planet can survive the post main sequence phase of a F or G dwarf.
Furthermore we should mention that Gl 86 is one of the closest binaries known today to harbor an exoplanet. Only two other systems $`\gamma `$ Ceph (a$``$19 AU, e$``$0.36, see Hatzes et al. 2003) and HD 41004 (23 AU see Zucker et al. 2004) have comparable separations. In such close binary systems the dynamical stability of planets is limited to a small region around the planet host star. According to Holman & Wiegert (1999) the critical semi-major axis a<sub>c</sub> for planets around Gl 86 A is only 6.2 AU assuming a circular binary orbit (m = 0.55 $`M_{\mathrm{}}`$ and a = 21 AU). Due to mass loss during the post main sequence phase of the white dwarf progenitor, the binary separation was even smaller before (M<sub>tot</sub>a = const $``$ a<sub>old</sub> = 15.4 AU), i.e. a<sub>c</sub> = 3.7 AU. We calculate critical semi-major axis also for the two other close binaries and get a<sub>c</sub> = 7.5 AU for HD 41004 and a<sub>c</sub> = 4.0 AU for $`\gamma `$ Ceph. We should mention that all exoplanets detected in these close binary systems actually reside within the proposed long-time stable regions. However it would be of particular interest to search for further substellar companions in theses close binaries to verify with observational results the published theoretical constraints of planet stability in binary systems. |
warning/0506/astro-ph0506263.html | ar5iv | text | # 1 Profiles of resonance lines of K I (𝜆𝜆 766.6, 770.1 nm) and Na I (𝜆𝜆 589.1, 589.7 nm) computed in the frame of collisional broadening theory (van der Waals broadening) for 1200/5.0 C-model atmosphere of Tsuji [], see [] for more details.
ULTRACOOL DWARFS
Yakiv V. Pavlenko
<sup>1</sup>Main Astronomical Observatory, NAS of Ukraine
27 Akademika Zabolotnoho Str., 03680 Kyiv, Ukraine
e-mail: yp@mao.kiev.ua
We present results of modeling of spectra of M-, L-, T-dwarfs. Theoretical spectra are fitted to observed spectra to study the main parameters of the low-mass objects beyond the bottom of Main Sequence. Application of “lithium” and “deuterium” tests for assessment of ultra-cool dwarfs are discussed.
<sup> </sup><sup> </sup>footnotetext: © Yakiv V.Pavlenko 2004
INTRODUCTION
Population of ultracool (UC) dwarfs occupies the right-right-bottom quadrant below the bottom of the Main sequence. A lot of UC dwarfs was discovered after 1995 (see and for reviews). Basically, we can define at least 3 different populations of ultracool dwarfs:
— Low mass stars (LMS). Hydrogen burns in their core.
— Brown dwarfs (BD). Hydrogen cannot burns in their core. Their existence were predicted by Kumar , . Later inverstigations show that lithium burns inside the brown dwarfs of $`55M_j<M<75M_j`$ (see for more details). Here $`M_j`$ is mass of Jupiter: 1$`M_j=0.001M_{}`$. First brown dwarfs Teide1 and Gl 229B were discovered by groups of Rebolo and Nakajima , respectively. Deuterium should be depleted in atmospheres of brown dwarfs.
— Planets ($`M<13M_J`$) preserve deuterium (and lithium) during their evolution .
First spectral classifications of UC dwarfs were provided by Kirkpatrick et al. and Martín et al.. Today we can asses their spectra (see libraries of spectra on or ):
— M-dwarfs (GJ406, VB10, VB8, etc). TiO dominates in their spectra.
— L-dwarfs (GD169B, Kelu1,2MASS 0920+35, etc.). K and Na lines are the main features there (,), Ti and V atoms are bound into dust particles.
— T-dwarfs (Gl 229B, SDSS 0151, SDSS 1110, etc) – infrared spectra show CH<sub>4</sub> lines.
— planets (see list of discovered planets on web , and references therein). First confirmed discovery of planetary system 51 Peg was carried out by Mayor & Queloz (see Marcy et al. ).
M-, -L, -T dwarfs are of different effective temperatures and masses. Still, “the Main sequence” for brown dwarfs and L-, T-dwarfs forms the approximately horizontal line (Jupiter is on the left side radii-masses plot, see ) – the dependence of radii of UC dwarfs on mass is extremely weak due to degeneracy of the gas in their cores. As result, sizes of old brown dwarfs, L-dwarfs and Jupiter are comparatible.
As was noted by Zapatero Osorio (private communication) depending on age, T-dwarfs can be brown dwarfs (if they are old) or ”planetary objects” (masses below the deuterium burning limit, if they are young). Hence, very young T-dwarfs do not burn deuterium. Then, giant planets have been found by indirect techniques around stars. Young objects a few times more massive than Jupiter have been identified using direct imaging techniques. They are characterized by ultracool atmospheres (L and T types). These objects are free-floating in star-forming regions and very young clusters. This poses challenge to current theories of stellar and planetary formation (see Proc. of IAUS 211 ).
Different UC dwarfs are of different structure as well:
— inside the LMS we have core with hydrogen burning zone,
— Brown dwarfs burn deuterium, the most massive BDs ($`55M_J<M<75M_J`$) burn lithium within short time scales (see refs in ).
— planets are only objects without any nuclear burning processes. They preserve deuterium and lithium from times of their formation.
Models of formation of spectra of ultracool dwarfs
To model spectra and spectral energy distributions (SEDs) of ultracool dwarfs we should account a few complicate processes which govern physical state of their atmospheres:
* Dust formation processes. Due to the low temperatures and high pressure regime some molecular (and atomic) species are bound in different grain particles (see ). Indeed, molecular bands of VO and TiO are weaker in L-dwarf spectra in comparison with M-dwarfs.
* Damping of K and Na lines. Resonance doublets of K and Na form the most impressive features in spectra of L-dwarfs. Formally computed equivalent widths of there lines can be of order a few kÅ(see Fig. 1 and for more details).
* Dust opacities. Importance of account of dust opacities for a procedure of numerical modelling of spectra of L- and T- dwarfs was shown by Pavlenko et al.. Basically, the problem of the dust opacities in L-dwarf atmosphere is rather complicate – we should account absorption/scattering by particles of different composition, sizes, orientations. Moreover, recent recearches provide some evidences of cloudy structure of dust layers of L-dwarfs atmospheres (see materials of IAUS 211 ).
Optical spectra: K and Na lines
Resonance lines of Na I($`\lambda \lambda `$ 589.1, 589.7 nm) and K I ($`\lambda \lambda `$ 766.6, 770.1 nm) are very strong in spectra of UC dwarfs , because majority of alkali atoms exists there as neutral atoms. Na I resonance lines are stronger — in atmospheres of majority stars log N(Na) $`>`$ log N(K).
Lines of alkali metals observed in UC dwarfs spectra are pressure broadened. Extremely strong broadening of K and Na resonance lines provides an serious problem for their modelling. We can use for their wings modelling the traditional approach based on collisional interactions between atoms K and Na and H, He and molecule only for qualitative analysis .
More sophisticated approaches based on quantum-chemical consideration of the impact of potential fields provided by different species on levels K and Na were proposed recently by different groups (see and ).
On the other hand, in atmospheres of L-dwarfs the dust absorbs/scatters photons in wide spectral spectral regions. Dust opacity affects the overall spectral distributions (see for more details). Perhaps, for core and near wings of resonance lines K I and Na I we can still use the collisional approach .
Infrared spectra: H<sub>2</sub>O bands
Water bands cover the wide spectral regions in the infrared spectra of UC dwarfs (see and the poster by Lyubchik et al. on this session). For a long time the computation of the most complete lists of H<sub>2</sub>O is the real challenge for theoretical physics (see a review in ). In general, incompleteness of water line lists used for the numerical analysis of infrared spectra of UC dwarfs can increase our problems of stellar spectra computations in different ways:
– outer layers of model atmospheres computed with incomplete line lists of H<sub>2</sub>O are “too hot”.
– results of spectral synthesis can be affected by incompleteness of H<sub>2</sub>O lists.
Water bands in the IR are of interest for different topics. Infrared CO band at 2.3 and 4.5 micron can be used for determination of basic parameters of UC dwarfs: abundances, effective temperatures, rotational velocities (see , ). For their theoretical modelling the use of reliable list of H<sub>2</sub>O lines is of crucial importance (see for more details).
Lithium test
“Lithium test” was proposed Rebolo et al. to identify brown dwarfs from the population of LMS. Before 1995 L- and T-dwarfs were not known, and main attention was paid for the low-gravity M-dwarfs. They suggested that at least part of low-mass dwarfs in young open clusters should preserve their lithium. Observation of lithium lines in spectra of late M-dwarfs provides the direct evidence of their substellar nature. Pavlenko et al. showed that lithium lines can be detected in spectra of brown dwarfs despite of severe blending of the atomic lines by molecular bands. Later lithium lines were really found in spectra of some brown dwarfs (Teide1 , Kelu1, etc. see ).
On the other hand, observation of lithium lines in spectra of late-type low gravity dwarfs of open clusters provide the information about their age. Due to theoretical predictions (see refs. in ) the smallest objects should be cooled very quickly, i.e within time scales of a few Myrs. Still young, i.e low gravity dwarfs of ages 3-5 Myrs preserve their lithium as well. In Fig. 2 results of determination of litium abundances in atmospheres of the low-mass dwarfs of open cluster $`\sigma `$ Ori are showed. Note, these results are based on analysis of pseudequivalent widths of lithium lines (see and for more details) – measurements of the pseudoequivalent widthes are provided in respect to the local pseudocontinuum formed by molecular lines.
Perhaps, determination of masses of brown dwarfs is the main problem. Fortunately, often brown dwarfs form binary systems. The study of the low-mass objects is of special interest. First observations of GJ569B provide some evidences about its substellar nature (see refs in . Still, later observations of Martín et al. on Keck Telescope show that GJ569B is double system – GJ569Ba and GJ569Bb are orbiting with period 892 $`\pm `$ 25 days . Lithium test for this system is of crutial importance. However, in this case we should manage a combine spectrum formed in atmospheres of both components of different masses.
Later the application of the “lithium test” was discussed for L-dwarfs and even T-dwarfs (see ). Indeed, lithium lines were observed in spectra of some UC dwarfs.
Deuterim test
In the cores of the ultracool dwarfs the correlations effects between ions dominate lowering Coloumb barrier between particles (see for more details).Still temperatures in the interiors of UC dwarfs of masses M $`<`$ 13M<sub>J</sub> cannot be high enough (T $`<`$ 0.5 MK) to initiate there a nuclear burning of deuterium.
Béjar et al. propose to use observations of lines of deuterium contained species to determine the ages/masses of the smallest UC dwarfs. The task is very difficult in both theoretical and observational aspects. The simplest case would be proposed consists of the analysis of HDO/H<sub>2</sub>O spectra in the IR spectra of UC dwarfs . Still, HDO lines are very blended by H<sub>2</sub>O lines . From one side, we should have very accurate line lists both H<sub>2</sub>O and HDO. Observed intensities of HDO lines cannot exceed a few percent (see ibid and ). Moreover, IR spectrum of UC dwarfs should contain lines of other polyatomic species (CH<sub>4</sub> and others). These factors increase the demands for a capacity of observational facilities and the quality of theoretical data to identify and to carry the analysis of HDO lines in spectra of UC dwarfs.
ACKNOWLEDGEMENTS I thank Maria Rosa Zapatero Osorio (LAEFF, Spain) for her highlighted and helpful comments. I am grateful my collaborators and coauthors Hugh R.A. Jones (Univ. of Hertfordshire, UK), Rafael Rebolo, MartínEduardo, Víctor J. S. Béjar (IAC, Spain) for their contribution.
My inverstigations were supported by Small Recearch Grant of AAS and Royal Society travel grant and travel grants from Liverpool University (UK). |
warning/0506/cond-mat0506073.html | ar5iv | text | # Electron Mobility and Magneto Transport Study of Ultra-Thin Channel Double-Gate Si MOSFETs
## I Introduction
Advances in silicon-on-insulator (SOI) technology have made possible high quality thin channel single-gate and double-gate (DG) Si metal-oxide-field-effect-transistor (MOSFET) device structures celler:2003 , which are intensively explored at the momentshoji:1999 ; gamiz:2001b ; esseni:2003 ; prunnila:2004b . These both devices have many advantages over the standard bulk Si MOSFETs. However, DG MOSFETs are usually regarded as the most promising solution to the problems faced when the device/gate length is down-scaled into sub-50 nm regime (short channel effects) due to superior electrostatic gate control of the transistor channel charge.celler:2003 In addition of boosting the gate control the DG devices also provide other benefits in the form of enhanced electron mobility.gamiz:2001b ; esseni:2003 ; prunnila:2004b
Apart from the relevance to the microelectronics industry the SOI material has also enabled investigation of fundamental phenomena in SiO<sub>2</sub>-Si-SiO<sub>2</sub> quantum well structures. Unfortunately, the electron mobility at the Si-SiO<sub>2</sub> interface is intrinsically limited by the effects well known from standard bulk Si-MOSFETs. Despite this fact, the strong electronic and optical confinement provided by the SiO<sub>2</sub> barrier, many valley Si conduction band and indirect gap makes this quantum well system a unique tool to study several interesting effects in, for example, electron-hole liquids pauc:2004 and bi-layer electron systemstakashina:2004b ; prunnila:2005 .
In this work, we report on the fabrication and detailed room temperature and low temperature electronic properties of DG Si MOSFETs with Si well thickness in the range $`717`$ nm. Mobility, electron density and high magnetic field diagonal resistivity are mapped in large double-gate bias windows, enabling detailed investigation of the transport properties at different gate bias (electron distribution) symmetries.
## II Experimental
The DG MOSFETs were fabricated on commercially available 100 mm unibond (100) SOI wafers with n<sup>-</sup> (batch B) and p<sup>-</sup> (batch F) Si layer . The nominal Si film thickness was 400 nm and the buried oxide (BOX) was 400 nm thick. First, we exchanged the insulating handle wafer to a n<sup>+</sup> wafer to enable efficient metallic back gating at all temperatures. This procedure began by a growth of a 80 nm-thick dry oxide at 1000 C. Then the SOI wafer was vacuum bonded to a n<sup>+</sup> (111) Si wafer with $`2\times 10^{19}`$ $`\text{cm}^3`$ arsenic concentration. The bonded interface was annealed at 1100 C and the original insulating $``$500 $`\mu `$m-thick handle wafer was removed by etching in 25% tetramethyl ammonium hydroxide solution at 80 C. Finally the ”old” BOX layer was stripped in a 10% HF and as a result we had a SOI wafer with heavily doped handle and 360 nm thick SOI film and 80 nm thick BOX (back gate oxide).
The actual device fabrication for batch F begun by locally thinning the Si layer in some parts of the wafer in order to fabricate devices with different Si well thickness, $`t_\text{W}`$. This was done by utilizing standard nitride masking and thermal oxidation, i.e., local oxidation of silicon (LOCOS). Another two LOCOS steps were then used to define the thin channels and active areas of the devices. In the channel thinning step the Si thickness was reduced in the channel regions to create a recessed source-drain MOSFET structures. To define the active areas the parts of the Si layer that were not protected by the nitride mask were fully converted to SiO<sub>2</sub>. After oxide stripping and wafer cleaning a 40 nm-thick gate oxide was grown at 1000 C in oxygen - DCE (dichloroethylene) ambient. The applied DCE flow into the oxidation furnace corresponded to $``$2% of HCl. A 250 nm-thick polysilicon gate was deposited by CVD, implanted with As and then patterned with UV-lithography and plasma etching. The contact areas were implanted with As while the gate electrode protected the Si channel. A $``$500 nm thick CVD oxide was deposited and the implanted doses were activated at 950 C. Finally, after contact window opening and Al metallization, the samples were annealed in H<sub>2</sub>/N<sub>2</sub> ambient at 425 C for 30 min. Devices in batch B were fabricated in a similar fashion. The major difference was that in this batch the active areas of the devices were defined by etching through the Si layer instead of LOCOS process.
Figure 1 shows a schematic cross-section of our DG MOSFET device structure together with HRTEM image of a device from batch B. Properties of the devices reported here are listed in Table 1. The gate oxide thickness is $`t_{\text{OX}}=`$ 40 nm (43 nm) and the buried oxide thickness is $`t_{\text{BOX}}=`$ 83 nm (80 nm) for batch F (B). The cited well type is the Si layer type given by the SOI wafer manufacturer. No intentional doping is introduced into the channel in order to maximize the low temperature mobility. Further details about samples B-E721 and B-E742 can be found from prunnila:2005 and prunnila:2004b , respectively.
All electrical characteristics reported here were obtained from Hall bar structures with 100 $`\times `$ 1900 $`\mu `$m<sup>2</sup> channel dimensions and a 400 $`\mu `$m voltage probe distance. We used two different methods to determine the electron density $`N`$: ”split” capacitance-voltage (SCV) method sodini:1982 at room temperature and Shubnikov-de Haas (SdH) oscillations at low temperature. The mobility (or effective mobility) was determined from $`\mu =\sigma /eN`$ where $`\sigma `$ is the conductivity measured by a four point method and $`e`$ is the electronic charge. In the room temperature SCV measurements we used Agilent 4294A precision impedance analyzer at frequency of 691 Hz. This frequency was low enough to provide results that were independent of the channel resistance when all the voltage probes and the source and the drain were connected to the virtual ground of the analyzer. The channel conductivity was determined by Agilent 4156 C precision semiconductor parameter analyzer. DC offsets were systematically removed from this data by applying at least four different source-drain bias values at each gate voltage point. In the low temperature characterization the samples were mounted to a sample holder of a He-3 cryostat and the electrical measurements were performed utilizing standard low frequency lock-in techniques. A combination of voltage and current excitation was utilized in order to keep the source-drain bias from heating the electrons above the lattice temperature at sub-1 K temperatures.
## III Results and Discussion
### III.1 Room Temperature Mobility
Figure 2 shows experimental constant mobility and electron density contours of two DG MOSFETs with $`t_\text{W}=`$ 6.8 nm (F-E741) and $`t_\text{W}`$ = 17.3 nm (B-E742) as a function of top gate voltage $`V_{\text{ TG}}`$ and back gate voltage $`V_{\text{BG}}`$ at 300 K. The family of curves explicitly shows how the electron mobility behaves as a function of carrier density and gate biases. We can see that for both devices on any of the constant $`N`$ contours the mobility maximum occurs in the vicinity of the symmetric gate bias line $`V_{\text{TG}}/`$ $`t_{\text{OX}}=`$ $`V_{\text{BG}}/t_{\text{BOX}}`$ \[dashed black lines in Fig. 2\]. When gate bias asymmetry is increased the mobility decreases monotonically. The mobility enhancement towards the symmetric bias line can be related to the volume inversion/accumulation effect, where the electron gas spreads through out the whole Si well. This is illustrated by self-consistent Hartree electron distributions in Fig. 3.
The mobility modulation along the constant $`N`$ contours follows mainly from the modulation of phonon scattering, surface roughness scattering, and conduction effective mass.gamiz:2001b ; gamiz:2003 The first two have minimum and the latter maximum at symmetric bias. However, the phonon and the surface roughness scattering have the strongest influence on the mobility gamiz:2003 and this leads to the experimentally observed behavior of Fig. 2. The large difference between the magnitude of the mobilities of the two devices follows from the Si well thickness dependency of the effective mobility (not from sample quality) and it is consistent with Monte Carlo simulations gamiz:2003 and experiments of other groups esseni:2003 .
It is evident from Fig. 2(a) that the mobility maxima do not occur precisely on the symmetric bias line for the thinner device F-E741 at higher electron densities; the maxima are shifted towards larger (smaller) $`V_{\text{TG}}`$ ($`V_{\text{BG}}`$). More careful inspection reveals that the behavior is actually similar for the thicker device B-E742. The observed effect is brought out more clearly in Fig.4, which shows the top interface ($`V_{\text{BG}}=0`$), back interface ($`V_{\text{TG}}=0`$), and symmetric bias mobilities as a function of electron density . We can observe that the back interface mobility falls below the top interface mobility in both devices at high electron densities. As this top/back interface mobility difference increases as a function of electron density we interpret that the Si-BOX interface has larger surface roughness than the Si-gate oxide interface.
### III.2 Magneto Transport
Figure 5 shows the diagonal resistivity $`\rho _{xx}`$ measured as a function of top and back gate voltages from two DG MOSFETs F-E73 ($`t_\text{W}=7.0`$ nm) and F-E42 ($`t_\text{W}=14.2`$ nm) at $`B`$ = 9.0 T at 0.27 K. The numbers inside the axis indicate the Landau level (LL) filling factors $`\nu =Nh/eB`$. The bright and dark color correspond to low and high $`\rho _{xx}`$, respectively. In the white regions $`\rho _{xx}0`$ and Hall resistance $`\rho _{xy}`$ is equal to integer quantum Hall (QH) value $`h/e^2\nu `$. The obscured maximum slightly below the symmetric bias line and close to the threshold in both $`\rho _{xx}`$ data is an experimental artefact (present below $`1`$ K).note:arte
We first draw our attention to the LL filling factor behavior of the thin sample F-E73 \[Fig.5(a)\]. We can observe that the $`\rho _{xx}`$ minima are continuous trajectories, which suggest that only the ground 2D sub-band is populated. The constant LL trajectories $`\nu =4(k+1)`$, $`\nu =2k+3`$ and $`\nu =4k+2`$ ($`k=0,1,2\mathrm{}`$) can be related to cyclotron, valley and spin gaps, respectively. When the gate biases confine the electron gas closer to the back interface ($`V_{\text{BG}}>`$ $`V_{\text{TG}}t_{\text{BOX}}/`$ $`t_{\text{OX}}`$) the $`\rho _{xx}`$ minima corresponding to $`\nu =6`$ and $`\nu =10`$ become shallow. This weakening of the minima can be addressed to the stronger elastic scattering in the vicinity of the Si-BOX interface, which was already detected in the room temperature mobility (see also next sub-section), leading to disorder broadening of the LLs.
The thick sample F-E42 \[Fig.5(b)\] behaves similarly in comparison to F-E73 when electron density is low (small $`\nu `$) and also when $`V_{\text{BG}}0.5`$ V or $`V_{\text{TG}}0.25`$ V. This indicates single sub-band occupation and the $`\rho _{xx}`$ minimum trajectories can be addressed to cyclotron, valley and spin gaps as was done above. In the bias region where $`V_{\text{BG}}0.5`$ V, $`V_{\text{TG}}0.25`$ V and $`\nu >4`$ the trajectories are broken into a 2D pattern, which is a signature of two sub-band (bi-layer) transport muraki:1999 ; prunnila:2005 ; takashina:2004b . By analyzing the SdH oscillations of $`\rho _{xx}`$ as a function of inverse magnetic field (at 0.3 K) in the spirit of Ref. muraki:2000 along the symmetric bias line we find that when top gate is adjusted between $`V_{\text{TG}}=+0.9+2.78`$ V the energy spacing of the bonding - antibonding sub-bands, $`\mathrm{\Delta }_{\text{BAB}},`$ varies monotonically from $`23`$ to $`2`$ meV.
The threshold for the second sub-band is already at $`V_{\text{TG}}0.75`$ V on the symmetric bias line, which is roughly on the $`\nu =3`$ trajectory in Fig. 5(b). Therefore, Zeeman and valley splitting can push the lower valley spin down state of the second sub-band in to the valley gap $`\nu =3`$ of the first sub-band around the symmetric bias, where $`\mathrm{\Delta }_{\text{BAB}}`$ has a minimum value. This could explain the destruction of the $`\nu =3`$ QH state (appearance of finite $`\rho _{xx}`$) at symmetric gate bias. Note that the effect is complicated by the possible symmetry dependency of valley splittingtakashina:2004b and also by modulation of interface trap Coulomb scattering. These effects will be further explored elsewhere.
### III.3 Low Temperature Mobility
Figure 6 shows experimental constant mobility contours at 1.6 K and SdH electron density contours of device F-E73 as a function of $`V_{\text{TG}}`$ and $`V_{\text{BG}}`$. The electron density is determined from SdH oscillations at 0.27 K: $`N`$ is obtained from similar constant $`\nu `$ trajectories in $`\rho _{xx}`$ at constant magnetic field that were discussed above. Magnetic field value of $`B`$ = 2.5 T was chosen in order to keep all $`\rho _{xx}`$ minima clearly above zero, which enabled accurate determination of constant LL trajectories. In the conductivity measurement at 1.6 K a small magnetic field ($`0.2`$ T) was applied to suppress the quantum corrections of conductivity.
The electron density and gate bias symmetry dependency of low temperature mobility of F-E73 in Fig. 6 closely resembles that of room temperature effective mobility of F-E741 \[Fig. 2(a)\]. Due to absence of phonon scattering the asymmetry in the elastic scattering properties of the two interfaces is now brought out more clearly. It is evident that the maximum mobility is shifted from the symmetric gate bias position towards larger $`V_{\text{TG}}`$ and smaller $`V_{\text{BG}}`$ on any of the constant $`N`$ contours.
If we adjust the gate voltages along an arbitrarily chosen straight line (direction) in Fig. 6 in such a fashion that $`N`$ increases the mobility shows the typical MOSFET behavior, where the mobility first increases at low carrier density, reaches a maximum value and then decreases at high carrier density. This behavior can be addressed to well known elastic scattering mechanisms: Coulomb scattering that is most effective at low $`N`$ and interface roughness scattering that is most effective at high $`N`$ .ando:1982
When $`t_\text{W}10`$ nm the second sub-band becomes populated even when both gates have a modest positive bias, as was demonstrated in the previous sub-section. This alters the mobility behavior and in general the simple Coulomb/surface roughness-scattering picture is lost. Further effects arises from the fact that in some bias ranges the electrons in the second sub-band can be localized prunnila:2005 . The co-existence of localized and non-localized electrons complicates the scattering mechanisms popovic:1997 ; feng:1999 and discussion of such effects is beyond the scope of this paper. Therefore, full $`V_{\text{TG}}V_{\text{BG}}`$ mobility dependency of the devices with $`t_\text{W}10`$ nm will be reported elsewhere. Here we only cite the maximum mobilities, which can be found from Table 1. The maximum mobility $`\mu _{\mathrm{max}}`$ decreases with decreasing $`t_\text{W}`$ as expected. $`\mu _{\mathrm{max}}`$ scaling with $`t_{\text{ W}}`$ is consistent with recent experimental observations on single gate SOI devicesprunnila:2004 .
Finally, we note that as the maximum mobilities for all devices are relatively high it is unlikely that the mobility degradation at the back interface could originate from increased Coulomb scattering; the degradation is mainly due to interface roughness scattering. High mobility also indicates extremely low interface trap density, which justifies our assumption that in the room temperature measurements the SCV electron density corresponds to that of mobile electrons.
## IV Summary
In summary, we have reported on detailed room temperature and low temperature transport properties of ultrathin channel double-gate Si MOSFETs. The devices were fabricated on SOI wafers utilizing wafer bonding, which enabled us to use heavily doped metallic back gate. The devices showed mobility enhancement effects at symmetric gate bias at room temperature, which is the finger print of the volume inversion/accumulation effect. Small asymmetry in the mobility could be detected at 300 K between the top and back interfaces of the Si well. The effect could be enhanced at low temperatures and the mobility asymmetry was interpreted to arise from different interface roughness of the top and back interface. Low temperature peak mobilities of the reported devices scale monotonically with Si well thickness and the maximum low temperature mobility was 1.9 m<sup>2</sup>/Vs, which was measured from a 16.5 nm thick device. From the magneto transport data we observed single and two sub-band transport effects depending on the well thickness and gate biasing.
## V Acknowledgements
Technical assistance by M. Markkanen in the sample fabrication is gratefully acknowledged. K. Henttinen is thanked for the wafer bonding and S. Newcomb is thanked for the TEM analysis. This work has been partially funded by EU (IST-2001-38937 EXTRA), by the Academy of Finland (# 205467 CODE) and GETA graduate school. |
warning/0506/math0506405.html | ar5iv | text | # Auslander algebras and initial seeds for cluster algebras
## Introduction
Let $`\mathrm{k}`$ be a field and $`Q`$ be a Dynkin quiver. So the underlying graph $`|Q|`$ of $`Q`$ is a simply laced Dynkin diagram. We produce for the preprojective algebra $`\mathrm{\Lambda }`$ over $`\mathrm{k}`$ associated to $`Q`$ a module $`\mathrm{I}_Q`$ by pushing a minimal injective cogenerator over the Auslander algebra $`\mathrm{\Gamma }_Q`$ of $`\mathrm{k}Q`$ to $`\mathrm{\Lambda }\mathrm{mod}`$. It is easy to see that $`\mathrm{I}_Q`$ decomposes into $`r=|\mathrm{\Pi }|`$ pairwise non-isomorphic direct summands. We show that $`\mathrm{I}_Q`$ is a rigid module module, i.e. $`\mathrm{Ext}_\mathrm{\Lambda }^1(\mathrm{I}_Q,\mathrm{I}_Q)=0`$. Moreover, the Gabriel quiver $`\stackrel{ˇ}{A}_Q`$ of $`\mathrm{End}_\mathrm{\Lambda }(\mathrm{I}_Q)^{\text{op}}`$ is obtained from the Auslander-Reiten quiver $`A_Q`$ of $`\mathrm{k}Q`$ by inserting an extra arrow $`x\tau x`$ for each non-projective vertex $`x`$.
In we have shown that if $`M=_{i=1}^mM_i`$ for pairwise non-isomorphic indecomposable $`\mathrm{\Lambda }`$-modules $`M_i`$, then $`\mathrm{Ext}_\mathrm{\Lambda }^1(M,M)=0`$ implies $`mr`$. So our result shows that this maximum is assumed for each Dynkin quiver. By \[17, Theorem 2.2\] we conclude that $`\mathrm{I}_Q`$ is a maximal 1-orthogonal $`\mathrm{\Lambda }`$-module and thus $`\mathrm{End}_\mathrm{\Lambda }(\mathrm{I}_Q)`$ is a higher Auslander algebra in the sense of Iyama . It follows that each rigid module $`M`$ as above can be completed to a rigid module with $`r`$ pairwise non-isomorphic indecomposable direct summands. Note, that for the proof of the main result of it is essential that the quiver $`\stackrel{ˇ}{A}_Q`$ has no loops.
Let now $`G`$ be a complex simply connected simple Lie group of type $`|Q|`$ with $`NG`$ a maximal unipotent subgroup. Choose $`𝐢=(i_1,i_2,\mathrm{},i_r)`$ a reduced expression for the longest element $`w_0`$ of the Weyl group $`W`$ of $`G`$. It follows from that the coordinate ring $`[N]`$ is an (upper) cluster algebra. Associated to $`𝐢`$ one obtains an initial seed $`(\mathrm{\Delta }(j,𝐢)_{j=1,2,\mathrm{},r}^{},\stackrel{~}{B}(𝐢)^{})`$, where the $`\mathrm{\Delta }(j,𝐢)^{}`$ are certain generalized minors, and the exchange matrix $`\stackrel{~}{B}(𝐢)^{}`$ is obtained naturally from the quiver $`\stackrel{ˇ}{A}_Q`$ described above.
Next, $`𝐢`$ provides us with a convenient labelling of the indecomposable direct summands of $`\mathrm{I}_Q`$, that is, $`\mathrm{I}_Q=_{j=1}^r\mathrm{I}(j,𝐢)`$. Clearly the $`\mathrm{I}(j,𝐢)`$ are rigid. Thus, if we restrict to the special case $`\mathrm{k}=`$ the $`\mathrm{\Lambda }`$-modules $`\mathrm{I}(j,𝐢)`$ serve as natural labels for elements $`\rho _{\mathrm{I}(j,𝐢)}`$ of the dual of Lusztig’s semicanonical basis. This is a natural basis for $`[N]`$ and we show that $`\rho _{\mathrm{I}(j,𝐢)}=\mathrm{\Delta }(j,𝐢)^{}`$.
## 1. Main results
### 1.1.
We say that a quiver $`Q`$ is a Dynkin quiver if its underlying graph $`|Q|`$ is a Dynkin diagram of type $`𝖠,𝖣,𝖤`$. For a field $`\mathrm{k}`$ we consider the path category $`\mathrm{k}[Q]`$ (or $`\mathrm{k}Q`$ for short). The category $`\mathrm{k}Q\mathrm{mod}`$ of finitely presented $`\mathrm{k}`$-functors $`M:\mathrm{k}Q\mathrm{k}\mathrm{mod}`$ is equivalent to the category of finitely presented left modules over the corresponding path algebra which we denote by some abuse also by $`\mathrm{k}Q`$.
For a quiver $`Q`$ we consider the double $`\overline{Q}`$ which is obtained from $`Q`$ by adding a new arrow $`i\stackrel{a^{}}{}j`$ for each arrow $`i\stackrel{𝑎}{}j`$ in $`Q`$. The preprojective algebra $`\mathrm{\Lambda }`$ is the quotient of the path algebra $`\mathrm{k}\overline{Q}`$ by the ideal generated by the elements
$$\rho _q=\underset{\begin{array}{c}aQ_1\\ t(a)=q\end{array}}{}a^{}a\underset{\begin{array}{c}aQ_1\\ h(a)=q\end{array}}{}aa^{}\text{for}qQ_0,$$
see also 2.1.
In what follows, $`Q`$ will always be a connected Dynkin quiver. This implies that $`\mathrm{\Lambda }`$ is a finite-dimensional selfinjective algebra, which depends only on $`|Q|`$. Like $`\mathrm{k}Q`$, the algebra $`\mathrm{\Lambda }`$ can also be considered as a $`\mathrm{k}`$-category.
We have the universal covering $`F:\stackrel{~}{\mathrm{\Lambda }}\mathrm{\Lambda }`$ where $`\stackrel{~}{\mathrm{\Lambda }}`$ is the path category $`\mathrm{k}[Q]`$ modulo the usual mesh relations. The fundamental group $``$ of $`\mathrm{\Lambda }`$ acts on $`\stackrel{~}{\mathrm{\Lambda }}`$ via the translation $`\tau `$. Associated to $`F`$ we have the push-down functor $`F_\lambda :\stackrel{~}{\mathrm{\Lambda }}\mathrm{mod}\mathrm{\Lambda }\mathrm{mod}`$, see 2.3.
In $`Q`$ we find the Auslander-Reiten quiver $`A_Q`$ of $`\mathrm{k}Q`$ as a full convex subquiver. The Auslander category $`\mathrm{\Gamma }_Q`$ is the full subcategory of $`\stackrel{~}{\mathrm{\Lambda }}`$ which has the vertices of $`A_Q`$ as objects. Denote the inclusion of $`\mathrm{\Gamma }_Q`$ into $`\stackrel{~}{\mathrm{\Lambda }}`$ by $`J`$. There is a natural equivalence $`\mathrm{R}_Q`$ from $`\mathrm{\Gamma }_Q`$ to the category of indecomposable $`\mathrm{k}Q`$-modules, $`\mathrm{k}Q\mathrm{ind}`$. We say that an object $`x\mathrm{Obj}(\mathrm{\Gamma }_Q)`$ is projective if $`\tau x\mathrm{Obj}(\mathrm{\Gamma }_Q)`$, dually $`x`$ is injective if $`\tau ^1x\mathrm{Obj}(\mathrm{\Gamma }_Q)`$, see 2.4.
Associated to $`\mathrm{\Gamma }_Q`$ we consider the $`_0`$-graded category $`\stackrel{ˇ}{\mathrm{\Gamma }}_Q`$. It has the same objects as $`\mathrm{\Gamma }_Q`$ but the morphisms of degree $`i`$ are given by $`\stackrel{ˇ}{\mathrm{\Gamma }}_{Q,i}(x,y)=\mathrm{\Gamma }_Q(\tau ^ix,y)`$ if $`\tau ^ix\mathrm{Obj}(\mathrm{\Gamma }_Q)`$. We equip $`\stackrel{ˇ}{\mathrm{\Gamma }}_Q`$ with the natural composition. The Gabriel quiver $`\stackrel{ˇ}{A}_Q`$ of $`\stackrel{ˇ}{\mathrm{\Gamma }}_Q`$ is obtained from $`A_Q`$ by inserting an additional (degree $`1`$) arrow $`x\tau x`$ for each non-projective $`x\mathrm{Obj}(\mathrm{\Gamma }_Q)`$, see 3.3.
### 1.2. Start modules.
Let us write $`D`$ for the usual duality $`\mathrm{Hom}_\mathrm{k}(,\mathrm{k})`$. We denote by $`J^{}`$ the functor which considers a $`\mathrm{\Gamma }_Q`$-module (trivially) as a $`\stackrel{~}{\mathrm{\Lambda }}`$-module. Thus if apply the functor $`F_\lambda J^{}`$ to the injective $`\mathrm{\Gamma }_Q`$-module $`D\mathrm{\Gamma }_Q(,x)`$ we obtain a $`\mathrm{\Lambda }`$-module. We call
$$\mathrm{I}_Q:=\underset{x\mathrm{Obj}(\mathrm{\Gamma }_Q)}{}F_\lambda J^{}(D\mathrm{\Gamma }_Q(,x)).$$
the start module for $`\mathrm{\Lambda }`$ associated to $`Q`$. Note that $`F_\lambda J^{}(D\mathrm{\Gamma }_Q(,x))`$ is isomorphic to a submodule of $`F_\lambda J^{}(D\mathrm{\Gamma }_Q(,\tau ^1x))`$ if $`x\mathrm{Obj}(\mathrm{\Gamma }_Q)`$ is not injective, and $`F_\lambda J^{}(D\mathrm{\Gamma }_Q(,x))`$ is an injective $`\mathrm{\Lambda }`$-module if $`x\mathrm{Obj}(\mathrm{\Gamma }_Q)`$ is injective, see 2.4 and 3.1.
Consider $`=_{x,y\mathrm{Obj}(\mathrm{\Gamma }_Q)}\stackrel{ˇ}{\mathrm{\Gamma }}_Q(x,y)`$ as a (graded) associative $`\mathrm{k}`$-algebra with multiplication induced from the composition of morphisms.
###### Theorem 1.
Let $`\mathrm{\Lambda }`$ be the preprojective algebra associated to a Dynkin quiver $`Q`$. Then $``$ is isomorphic to $`\mathrm{End}_\mathrm{\Lambda }(\mathrm{I}_Q)^{\text{op}}`$. In particular, the Gabriel quiver of $`\mathrm{End}_\mathrm{\Lambda }(\mathrm{I}_Q)^{\text{op}}`$ is identified with $`\stackrel{ˇ}{A}_Q`$ as described above in 1.1. Moreover $`\mathrm{I}_Q`$ is rigid in the sense that $`\mathrm{Ext}_\mathrm{\Lambda }^1(\mathrm{I}_Q,\mathrm{I}_Q)=0`$.
The proof of this result is prepared in 3.4, 3.5, 3.6 and finished in 3.7.
### 1.3. Reduced expressions.
Let $`\pi :\mathrm{Obj}(\mathrm{\Gamma }_Q)Q_0`$ denote the map induced by the composition $`FJ`$. We call a total ordering $`x(1)<x(2)<\mathrm{}<x(r)`$ of the objects of $`\mathrm{\Gamma }_Q`$ adapted (to $`Q`$) if $`\mathrm{\Gamma }_Q(x(i),x(j))=0`$ for $`i<j`$. It is easy to find such orderings given that the quiver $`A_Q`$ is directed.
We call a vertex $`iQ_0`$ a source in $`Q`$ if no arrow ends at $`i`$. In this case we denote by $`s_i(Q)`$ the quiver which is obtained from $`Q`$ by reversing each arrow starting in $`i`$.
Let $`𝐢=(i_1,i_2,\mathrm{},i_r)`$ be a reduced expression for the longest element $`w_0W`$, that is $`w_0=s_{i_1}s_{i_2}\mathrm{}s_{i_r}`$ where $`r=|\mathrm{\Pi }|`$. We say that $`𝐢`$ is adapted to $`Q`$ if $`i_1`$ is a source in $`Q`$ and $`i_{k+1}`$ is a source in $`s_{i_k}\mathrm{}s_{i_1}(Q)`$ if $`1k<r`$. This is dual to the original definition in .
If $`x(1)<\mathrm{}<x(r)`$ is an adapted ordering then $`𝐢:=(\pi (x(1)),\mathrm{},\pi (x(r)))`$ is a reduced expression for the longest element $`w_0`$ of the Weyl group $`W`$ associated to $`|Q|`$, see 1.4 below. In fact, the adapted orderings of $`\mathrm{Obj}(\mathrm{\Gamma }_Q)`$ correspond in this way bijectively to the reduced expressions for $`w_0`$ which are adapted to $`Q`$, see \[3, Theorem 2.5\]. We set
$$\mathrm{I}(j,𝐢):=F_\lambda J^{}D\mathrm{\Gamma }_Q(,x(j))\text{ for }1jr.$$
Thus, $`𝐢`$ provides us with a convenient way of labelling the direct summands of $`\mathrm{I}_Q`$.
### 1.4.
Let now $`𝔤`$ be a complex simple Lie algebra of type $`|Q|`$ with the usual Serre generators $`e_i,h_i,f_i`$ for $`iQ_0`$. Thus the $`h_i`$ form a basis of the abelian subalgebra $`𝔥`$ and the $`e_i`$ resp. $`f_i`$ generate maximal nilpotent subalgebras $`𝔫`$ resp. $`𝔫_{}`$. The simple roots $`\alpha _i`$ form a basis of the dual space $`𝔥^{}`$ such that $`\alpha _i(h_j)=a_{i,j}`$ where $`(a_{i,j})_{i,jQ_0}`$ is the Cartan matrix of $`|Q|`$. The fundamental weights $`(\varpi _i)_{iQ_0}`$ are the basis of $`𝔥^{}`$ dual to the basis $`(h_i)_{iQ_0}`$ of $`𝔥`$.
The Weyl group $`W`$ is the subgroup of $`\mathrm{GL}(𝔥^{})`$ which is generated by the reflections $`s_i`$ for $`iQ_0`$ such that
$$s_i(\alpha )=\alpha \alpha (h_i)\alpha _i\text{ for }\alpha 𝔥^{}.$$
This is a finite reflection group.
Let now $`G`$ be a complex simply connected simple algebraic group with $`\mathrm{Lie}(G)=𝔤`$. It has maximal unipotent subgroups $`N`$ resp. $`N_{}`$ with $`\mathrm{Lie}(N)=𝔫`$ resp. $`\mathrm{Lie}(N_{})=𝔫_{}`$, and the maximal torus $`H`$ has $`\mathrm{Lie}(H)=𝔥`$. Moreover, we have standard embeddings $`\phi _i:\mathrm{SL}_2G`$ such that
$$\mathrm{exp}(tf_i)=\phi _i(\begin{array}{cc}1& 0\\ t& 1\end{array})\text{ and }\mathrm{exp}(te_i)=\phi _i(\begin{array}{cc}1& t\\ 0& 1\end{array})\text{ for }iQ_0.$$
Moreover set
$$\eta _i(t):=\phi _i(\begin{array}{cc}t& 0\\ 0& t^1\end{array})H\text{ for }iQ_0\text{ and }t^{}.$$
Next, recall that $`N_G(H)/H`$ is canonically isomorphic to the Weyl group $`W`$ defined above. In fact, it is possible to choose representatives $`\overline{w}N_G(H)`$ for the elements $`wW`$ such that
$`\overline{s_i}`$ $`=\mathrm{exp}(f_i)\mathrm{exp}(e_i)\mathrm{exp}(f_i),`$
$`\overline{uv}`$ $`=\overline{u}\overline{w}\text{ if }l(uv)=l(u)+l(v).`$
We identify the weight lattice $`P=_{iQ_0}\varpi _i`$ with the group of multiplicative characters of $`H`$ in such a way that $`\eta _i(t)^{\varpi _j}=t^{\delta _{i,j}}`$ for $`i,jQ_0`$. If we write $`\mathrm{?}^\lambda `$ for the character of $`H`$ corresponding to the weight $`\lambda `$ it follows that
$$h^{w(\lambda )}=(\overline{w}^1h\overline{w})^\lambda \text{ for }hH,\lambda P,wW.$$
### 1.5. Cluster algebras.
The coordinate ring of the affine base space $`[N_{}\backslash G]`$ consists of the functions $`f[G]`$ which are invariant under $`N_{}`$, i.e. $`f(g)=f(ng)`$ for all $`gG`$ and $`nN_{}`$. Now $`[N_{}\backslash G]`$ is naturally a $`G`$-module via $`gf(x)=f(xg)`$ for $`g,xG`$. It is well-known that each irreducible highest weight $`G`$-module $`L(\lambda )`$ can be realized as a direct summand of $`[N_{}\backslash G]`$ by taking
$$L(\lambda )=\{f[N_{}\backslash G]f(hg)=h^\lambda f(g)\text{ for }hH,gG\}.$$
For each $`L(\lambda )`$ we choose a highest weight vector $`u_\lambda `$ which we normalize by the condition $`u_\lambda (1_G)=1`$. Following we define for each fundamental weight $`\varpi _i`$ generalized minors
$$\mathrm{\Delta }_{\varpi _i,w(\varpi _i)}:=\overline{w}u_{\varpi _i}L(\varpi _i)$$
for any $`wW`$. In it is shown in particular that the coordinate ring of the double Bruhat cell $`G^{e,w_0}=B(B_{}\overline{w}_0B_{})`$ has the structure of an (upper) cluster algebra. Here, $`B`$ and $`B_{}`$ are opposite Borel subgroups of $`G`$ with $`BN`$ and $`B_{}N_{}`$.
For a reduced expression $`𝐢=(i_1,i_2,\mathrm{},i_r)`$ for $`w_0`$ which is adapted to $`Q`$ and $`k[n,1][1,r]`$ we set $`v_{>k}:=s_{i_r}s_{i_{r1}}\mathrm{}s_{i_{k+1}}`$ if $`k1`$ and $`v_{>k}:=w_0`$ if $`k1`$. Then, following set
$$\mathrm{\Delta }(k,𝐢):=\mathrm{\Delta }_{\varpi _{i_k},v_{>k}(\varpi _{i_k})}$$
where we take $`i_k=k`$ for $`k[n,1]`$. The $`\mathrm{\Delta }(k,𝐢)`$ for $`k[n,1][1,r]`$ form an initial cluster for $`[G^{e,w_0}]`$. There is also an easy algorithm to calculate from $`𝐢`$ the corresponding exchange matrix. Now set
$$e(𝐢):=\{i[1,r]x(i)\mathrm{Obj}(\mathrm{\Gamma }_Q)\text{ non-projective}\}.$$
There is a closely related upper cluster algebra structure on the coordinate ring $`[N]`$, whose initial cluster is given by the restrictions to $`N`$ of the functions $`\mathrm{\Delta }(k,𝐢)`$ with $`k[n,1]e(𝐢)`$ (see 4.4). The quiver associated to the corresponding exchange matrix is obtained from $`\stackrel{ˇ}{A}_Q`$ by removing the arrows between injective vertices. See Section 4 for a more detailed discussion.
### 1.6. Semicanonical basis.
In Lusztig introduces the semicanonical basis of the enveloping algebra $`U(𝔫)`$. Its elements are labelled naturally by the irreducible components of the preprojective varieties $`\mathrm{\Lambda }_𝐯`$ for $`𝐯_0^{Q_0}`$. In we started the study of the dual semicanonical basis which can be regarded as a basis of $`[N]`$. In particular, we found that to (the isoclass of) a rigid $`\mathrm{\Lambda }`$-module $`M`$ there corresponds naturally an element $`\rho _M`$ of this basis. If we set
$$\theta :[1,r][n,1]e(𝐢),j\{\begin{array}{cc}i\hfill & \text{ if }\mathrm{R}_Q(x(j))D\mathrm{k}Q(,i),\hfill \\ k\hfill & \text{ if }\tau ^1x(j)=x(k),\hfill \end{array}$$
we can state our second main result precisely:
###### Theorem 2.
For $`j[1,r]`$ we have $`\mathrm{\Delta }(j,𝐢)^{}:=\mathrm{\Delta }(\theta (j),𝐢)=\rho _{\mathrm{I}(j,𝐢)}`$.
The proof of this result is done after some preparation in Section 5.
### 1.7. Dualities.
Denote by $`\mathrm{?}^{}`$ the involution on $`\overline{Q}`$ with $`aa^{}`$ and $`a^{}a`$ for all $`aQ_1`$. This induces an anti-automorphism of $`\mathrm{\Lambda }`$ which we also denote by $`\mathrm{?}^{}`$. Thus we have a self-duality
$$S:\mathrm{\Lambda }\mathrm{mod}\mathrm{\Lambda }\mathrm{mod}\text{ with }SM(\mathrm{?})=DM(\mathrm{?}^{}).$$
Let us define for a reduced expression $`𝐢=(i_1,\mathrm{},i_r)`$ for $`w_0`$ which is adapted to $`Q`$
$$\mathrm{P}(j,𝐢):=F_\lambda J^{}\mathrm{\Gamma }_Q(x(j),)\text{ and }\mathrm{P}_Q:=\underset{j=1}{\overset{r}{}}\mathrm{P}(j,𝐢).$$
Then it is easy to see that $`\mathrm{P}_QS\mathrm{I}_{Q^{\text{op}}}`$. Thus, $`\mathrm{Ext}_\mathrm{\Lambda }^1(\mathrm{P}_Q,\mathrm{P}_Q)=0`$ and $`\mathrm{End}_\mathrm{\Lambda }(\mathrm{P}_Q)\mathrm{End}_\mathrm{\Lambda }(\mathrm{I}_{Q^{\text{op}}})^{\text{op}}`$. Now, $`\stackrel{ˇ}{A}_{Q^{\text{op}}}`$ is the Gabriel quiver of $`\mathrm{End}_\mathrm{\Lambda }(\mathrm{I}_{Q^{\text{op}}})^{\text{op}}`$, see Theorem 1, and it is not hard to see that $`\stackrel{ˇ}{A}_{Q^{\text{op}}}`$ may be identified with $`\stackrel{ˇ}{A}_Q^{\text{op}}`$ (recall that the same happens for Auslander-Reiten quivers: $`A_{Q^{\text{op}}}A_Q^{\text{op}}`$). So $`\mathrm{End}_\mathrm{\Lambda }(\mathrm{I}_{Q^{\text{op}}})\mathrm{End}_\mathrm{\Lambda }(\mathrm{P}_Q)^{\text{op}}`$ and $`\mathrm{End}_\mathrm{\Lambda }(\mathrm{I}_Q)^{\text{op}}`$ have the same Gabriel quiver $`\stackrel{ˇ}{A}_Q`$, but they are usually not isomorphic.
On the other hand, $`S`$ induces the canonical anti-automorphism $`\mathrm{?}^{}`$ on Lusztig’s algebra of constructible functions $`U(𝔫)`$ which commutes with the comultiplication, see \[29, Section 3.4\]. This yields by duality an automorphism $`\omega `$ of the coordinate ring $`[N]`$ which anti-commutes with the comultiplication. The corresponding anti-automorphism of $`N`$ leaves the one-parameter subgroups $`N,t\mathrm{exp}(te_i)`$ invariant.
Note that $`𝐢^{}:=(\mu (i_r),\mu (i_{r1}),\mathrm{},\mu (i_1))`$ is a reduced expression for $`w_0`$ which is adapted to $`Q^{\text{op}}`$, see 2.3 for the definition of $`\mu `$. We leave it to the reader to verify the following:
$$\rho _{\mathrm{P}(j,𝐢)}=\omega (\rho _{\mathrm{I}(r+1j,𝐢^{})})\text{for }1jr,$$
see also 1.6.
### 1.8. Triangulated structures.
In the appendix (Section 7) we point out that the category of injective $`\stackrel{~}{\mathrm{\Lambda }}`$-modules is triangulated. This implies that the category of injective $`\mathrm{\Lambda }`$-modules is also triangulated. This fact helps us to explain the unusual symmetries in the stable module categories over both categories by results of Freyd and Heller. Among other useful formulas we can recover the famous “6-periodicity” for modules over the preprojective algebra.
## 2. The universal cover of a preprojective algebra
### 2.1. Quiver categories.
Let $`Q=(Q_0,Q_1,t,h)`$ be a quiver with vertices $`Q_0`$, arrows $`Q_1`$ and $`t,h:Q_1Q_0`$ such that we have $`a:t(a)h(a)`$ for each arrow $`aQ_1`$. A path in $`Q`$ is a sequence of arrows $`a_na_{n1}\mathrm{}a_1`$ such that $`t(a_{i+1})=h(a_i)`$ for $`i=1,2,\mathrm{},n1`$.
On $`^{Q_0}`$ we have the Ringel bilinear form
$$𝐯,𝐰=\underset{iQ_0}{}𝐯(i)𝐰(i)\underset{aQ_1}{}𝐯(t(a))𝐰(h(a)).$$
Let $`\mathrm{k}`$ be a field. Since we need to consider infinite coverings of preprojective algebras we have to consider $`\mathrm{k}`$-categories rather than $`\mathrm{k}`$-algebras.
If $`Q`$ is a quiver we denote by $`\mathrm{k}[Q]=\mathrm{k}Q`$ the $`\mathrm{k}`$-category which has $`Q_0`$ as objects and with the morphism space $`\mathrm{k}Q(p,q)`$ having the paths from $`p`$ to $`q`$ as a basis. The composition is naturally induced from the concatenation of paths.
### 2.2. Conventions.
If $`𝒟`$ is a $`\mathrm{k}`$-category we denote by $`𝒟\mathrm{mod}`$ the category of finitely presented (covariant) $`\mathrm{k}`$-functors $`𝒟\mathrm{k}\mathrm{mod}`$. These functors are also called left modules, see for example \[16, Section 2.2\] for more details.
Let $`G`$ be a group of $`\mathrm{k}`$-automorphisms of $`𝒟`$ which acts from the left on $`𝒟`$. Then $`G`$ acts naturally from the right on $`𝒟\mathrm{mod}`$. If $`gG`$ and $`M`$ is a left $`𝒟`$-module, then we write $`M^g():=M(g^1)`$ for the twisted module. For example, if $`x𝒟`$ we get for the projective module $`𝒟(x,)`$ an isomorphism $`𝒟^g(x,)𝒟(gx,)`$.
### 2.3. Construction of $`\stackrel{~}{\mathrm{\Lambda }}`$.
Let $`Q=(Q_0,Q_1,t,h)`$ be a (connected) Dynkin quiver, so the underlying graph $`|Q|`$ is one of the following:
Define a (possibly trivial) involution $`\mu `$ on the vertices of $`Q`$ by
$$\mu (q)=\{\begin{array}{cc}n+1q\hfill & \text{ in case }𝖠_n,\hfill \\ 2n1q\hfill & \text{ in case }𝖣_n\text{ and }(n\text{ odd and }qn1),\hfill \\ 6q\hfill & \text{ in case }𝖤_6\text{ and }q5,\hfill \\ q\hfill & \text{ otherwise.}\hfill \end{array}$$
Following let $`Q`$ be the quiver with vertices $`Q_0=\times Q_0`$ and arrows $`Q_1=\times \{Q_1Q_1^{}\}`$ where
$`\stackrel{~}{t}(i,a)`$ $`=(i+1,t(a)),`$ $`\stackrel{~}{t}(i,a^{})`$ $`=(i,h(a)),`$
$`\stackrel{~}{h}(i,\alpha )`$ $`=(i,h(a)),`$ $`\stackrel{~}{h}(i,a^{})`$ $`=(i,t(a)).`$
We think of $`Q^{\text{op}}`$ as a subquiver of $`Q`$ via the embedding $`q(0,q)`$. The quiver $`Q`$ admits a translation automorphism $`\tau `$ induced by $`\tau (i,q)=(i+1,q)`$. Moreover we have a “Nakayama permutation” $`\widehat{\nu }`$ of $`Q`$. In order to define it we need the following auxiliary function: For any vertex $`q`$ denote by $`l(q)`$ the number of arrows pointing towards $`1`$ on the unique walk from $`1`$ to $`q`$. Now, $`\widehat{\nu }`$ is defined on the vertices by
$$\widehat{\nu }(p,q)=\{\begin{array}{cc}(p+q1+l(\mu (q))l(q),\mu (q))\hfill & \text{ in case }𝖠_n,\hfill \\ (p+n2+l(\mu (q))l(q),\mu (q))\hfill & \text{ in case }𝖣_n,\hfill \\ (p+q+2+l(\mu (q))l(q),\mu (q))\hfill & \text{ in case }𝖤_6\text{ and }q5,\hfill \\ (p+5,\mu (q))\hfill & \text{ in case }𝖤_6\text{ and }q=6,\hfill \\ (p+8,\mu (q))\hfill & \text{ in case }𝖤_7,\hfill \\ (p+14,\mu (q))\hfill & \text{ in case }𝖤_8.\hfill \end{array}$$
So, $`\widehat{\nu }`$ is a “translation reflection” stabilizing the “middle line” in case $`𝖠_n`$ and a translation in the cases $`𝖣_{2k}`$, $`𝖤_7`$ and $`𝖤_8`$. At the beginning of 6.1 we show as an example part of $`Q`$ for a quiver of type $`𝖣_5`$, together with the action of $`\widehat{\nu }`$.
We consider the mesh ideal $`I`$ in the category $`\mathrm{k}[Q]`$ which is generated by the elements
(2.1)
$$\underset{\begin{array}{c}aQ_1\\ t(a)=q\end{array}}{}(i,a^{})(i,a)\underset{\begin{array}{c}aQ_1\\ h(a)=q\end{array}}{}(i,a)(i+1,a^{})\text{ for all }(i,q)Q_0.$$
Note that $`\stackrel{~}{\mathrm{\Lambda }}:=\mathrm{k}[Q]/I`$ is independent of the orientation of $`Q`$ (up to relabelling of the vertices), and that $`I`$ is invariant under the quiver automorphisms $`\tau `$ and $`\widehat{\nu }`$. Thus we will use the same names for the induced automorphisms of $`\stackrel{~}{\mathrm{\Lambda }}`$. Note further that the commuting automorphisms $`\tau `$ and $`\widehat{\nu }`$ of $`\stackrel{~}{\mathrm{\Lambda }}`$ are determined up to isomorphism by their effect on objects.
The action of the group $``$ via $`\tau `$ on $`\stackrel{~}{\mathrm{\Lambda }}`$ induces self-equivalences $`\mathrm{?}^{(i)}`$ of $`\stackrel{~}{\mathrm{\Lambda }}\mathrm{mod}`$ with $`M^{(i)}():=M(\tau ^i)`$ for $`i`$. Moreover we obtain the covering functor $`F:\stackrel{~}{\mathrm{\Lambda }}\mathrm{\Lambda }`$ which sends $`(i,q)`$ to $`q`$. Associated to $`F`$ we have the push-down $`F_\lambda :\stackrel{~}{\mathrm{\Lambda }}\mathrm{mod}\mathrm{\Lambda }\mathrm{mod}`$ with
$$(F_\lambda M)(q)=_iM(i,q)$$
and the obvious effect on morphisms.
### 2.4. Auslander category.
We define a function $`N:Q_0_0`$ by the property
$$\tau ^{N(q)}(0,q)=\widehat{\nu }(0,\mu (q)).$$
This is well-defined by the construction of $`\widehat{\nu }`$ since $`\mu `$ is an involution on $`Q_0`$, see 2.3. The function $`N`$ depends on the orientation of $`Q`$ in case $`𝖠_n`$, $`𝖣_{2k+1}`$ and $`𝖤_6`$. In any case we have $`N(q)+N(\mu (q))=h(Q)2`$, where $`h(Q)`$ denotes the Coxeter number of $`|Q|`$.
Define now $`\mathrm{\Gamma }_Q`$ as the full subcategory of $`\stackrel{~}{\mathrm{\Lambda }}`$ which has the objects of the form $`(i,q)=\tau ^{iN(q)}\widehat{\nu }(0,\mu (q))`$ with $`qQ_0`$ and $`0iN(q)`$. In other words, we take the objects which lie between the two copies of $`Q^{\text{op}}`$ in $`Q`$ which are obtained via $`q(0,q)`$ resp. $`q\widehat{\nu }(0,q)`$. This is the Auslander category of $`\mathrm{k}Q`$. Note that $`\mathrm{\Gamma }_Q`$ depends on the orientation of $`Q`$.
By construction we have a full embedding $`\iota =\iota _Q:\mathrm{k}Q^{\text{op}}\mathrm{\Gamma }_Q`$ induced by $`q(0,q)`$. Moreover, $`\mathrm{\Gamma }_Q`$ is canonically equivalent to the category of indecomposable $`\mathrm{k}Q`$-modules via the functor
(2.2)
$$\mathrm{R}_Q:\mathrm{\Gamma }_Q\mathrm{k}Q\mathrm{ind},x\mathrm{\Gamma }_Q(\widehat{\nu }\iota ,x),$$
where $`\mathrm{\Gamma }_Q(\widehat{\nu }\iota ,x)=\mathrm{\Gamma }_Q(,x)\widehat{\nu }\iota :\mathrm{k}Q^{\text{op}}\mathrm{k}\mathrm{mod}`$ is a contravariant functor which we have to interpret as a left $`\mathrm{k}Q`$-module. For example, $`\mathrm{R}_Q(0,q)D\mathrm{k}Q(,q)`$ and $`\mathrm{R}_Q(\widehat{\nu }(0,q))\mathrm{k}Q(q,)`$.
Thus the Gabriel quiver of $`\mathrm{\Gamma }_Q`$ (which is the full subquiver of $`Q`$ with the vertices from $`\mathrm{\Gamma }_Q`$) is the Auslander-Reiten quiver $`A_Q`$ of $`\mathrm{k}Q`$.
Similarly, since we consider left modules, one obtains an equivalence
$$\mathrm{\Gamma }_Q\mathrm{inj}\mathrm{k}Q^{\text{op}}\mathrm{mod}(\mathrm{k}Q\mathrm{mod})^{\text{op}},II\widehat{\nu }\iota .$$
Here, $`\mathrm{\Gamma }_Q\mathrm{inj}`$ denotes the category of injective left $`\mathrm{\Gamma }_Q`$-modules.
Note that the $`qQ_0`$ parametrize the indecomposable projective-injective $`\mathrm{\Gamma }_Q`$-modules, namely we have
$$D\mathrm{\Gamma }_Q(,(0,q))\mathrm{\Gamma }_Q(\widehat{\nu }(0,q),).$$
Recall, that as an Auslander category $`\mathrm{\Gamma }_Q`$ has dominant dimension at least $`2`$ \[2, VI.5\]. This means by duality that each (indecomposable) injective $`\mathrm{\Gamma }_Q`$-module $`D\mathrm{\Gamma }_Q(,x)`$ has a projective presentation
$$P_{1,x}P_{0,x}D\mathrm{\Gamma }_Q(,x)0$$
with $`P_{0,x}`$ and $`P_{1,x}`$ also injective.
## 3. Start modules
### 3.1. Adjoint functors.
Let $`J:\mathrm{\Gamma }_Q\stackrel{~}{\mathrm{\Lambda }}`$ be the full embedding of locally bounded categories. Since $`\mathrm{\Gamma }_Q`$ is convex in $`\stackrel{~}{\mathrm{\Lambda }}`$ we get an exact functor “extension by $`0`$
$$J^{}:\mathrm{\Gamma }_Q\mathrm{mod}\stackrel{~}{\mathrm{\Lambda }}\mathrm{mod}\text{ with }(J^{}M)(x)=\{\begin{array}{cc}M(x)\hfill & \text{ if }x\mathrm{\Gamma }_Q,\hfill \\ 0\hfill & \text{ else.}\hfill \end{array}$$
For a $`\mathrm{\Gamma }_Q`$-module $`N`$ and $`x\mathrm{\Gamma }_Q`$ we have natural isomorphisms
$$\mathrm{Hom}_{\mathrm{\Gamma }_Q}(N,D\mathrm{\Gamma }_Q(,x))DN(x)\mathrm{Hom}_{\stackrel{~}{\mathrm{\Lambda }}}(J^{}N,D\stackrel{~}{\mathrm{\Lambda }}(,x)).$$
We conclude that $`J^{}`$ has a right adjoint $`J^\rho `$ which is defined by being left exact and
$$J^\rho D\stackrel{~}{\mathrm{\Lambda }}(,x)=\{\begin{array}{cc}D\mathrm{\Gamma }_Q(,x)\hfill & \text{ if }x\mathrm{Obj}(\mathrm{\Gamma }_Q),\hfill \\ 0\hfill & \text{ else. }\hfill \end{array}$$
The adjunction morphism $`\iota _M:J^{}J^\rho MM`$ is injective, its co-kernel is co-generated by $`_{y\mathrm{Obj}(\stackrel{~}{\mathrm{\Lambda }})\mathrm{Obj}(\mathrm{\Gamma }_Q)}M(y)`$.
###### 3.2 Lemma.
Let $`x,y\mathrm{Obj}(\mathrm{\Gamma }_Q)`$ and $`i`$. Then
$$\mathrm{Hom}_{\stackrel{~}{\mathrm{\Lambda }}}(J^{}D\mathrm{\Gamma }_Q(,x),J^{}D\mathrm{\Gamma }_Q^{(i)}(,y))\{\begin{array}{cc}\mathrm{\Gamma }_Q(\tau ^iy,x)\hfill & \text{ if }i0\text{ and }\tau ^iy\mathrm{Obj}(\mathrm{\Gamma }_Q),\hfill \\ 0\hfill & \text{ else.}\hfill \end{array}$$
In the first case we have more generally $`J^\rho (J^{}D\mathrm{\Gamma }_Q^{(i)}(,y))D\mathrm{\Gamma }_Q(,\tau ^iy)`$.
Note, that $`J^{}D\mathrm{\Gamma }(,y)`$ is the injective $`\mathrm{\Gamma }_Q`$-module with socle concentrated in $`y`$, but seen as $`\stackrel{~}{\mathrm{\Lambda }}`$-module, thus we may apply to it the translation functor $`\mathrm{?}^{(i)}`$, see 2.3.
Proof: The $`\stackrel{~}{\mathrm{\Lambda }}`$-module $`M=J^{}D\mathrm{\Gamma }_Q^{(i)}(,y)`$ has simple socle concentrated in the one-dimensional space $`M(\tau ^iy)`$. If $`\tau ^iy\mathrm{Obj}(\mathrm{\Gamma }_Q)`$ then this does not belong to the support of $`J^{}D\mathrm{\Gamma }_Q(,x)`$. On the other hand, if $`i<0`$, it is sufficient to show that there are no maps from an induced projective-injective module to $`M`$ since $`\mathrm{\Gamma }_Q`$ has dominant dimension $`2`$, see 2.4. Now,
$$0=M(\widehat{\nu }(0,q))=\mathrm{Hom}_{\stackrel{~}{\mathrm{\Lambda }}}(\stackrel{~}{\mathrm{\Lambda }}(\widehat{\nu }(0,q),),M)=\mathrm{Hom}_{\stackrel{~}{\mathrm{\Lambda }}}(J^{}\mathrm{\Gamma }_Q(\widehat{\nu }(0,q),),M)$$
with the first identity holding for $`i<0`$.
Thus, let $`i0`$ and $`\tau ^iy\mathrm{\Gamma }_Q`$. In this case we have in $`\stackrel{~}{\mathrm{\Lambda }}\mathrm{mod}`$ an injective presentation
$$0MD\stackrel{~}{\mathrm{\Lambda }}(,\tau ^ix)_jD\stackrel{~}{\mathrm{\Lambda }}(,y_j)^{m(j)}$$
for certain $`y_j\mathrm{Obj}(\stackrel{~}{\mathrm{\Lambda }})\mathrm{Obj}(\mathrm{\Gamma }_Q)`$. Our claim follows now from the construction of $`J^\rho `$. $`\mathrm{}`$
### 3.3. A graded category.
We construct the $`_0`$-graded category $`\stackrel{ˇ}{\mathrm{\Gamma }}_Q`$. It has the same objects as $`\mathrm{\Gamma }_Q`$, but the homogenous components are $`\stackrel{ˇ}{\mathrm{\Gamma }}_{Q,i}(x,y):=\mathrm{\Gamma }_Q(\tau ^ix,y)`$ if $`\tau ^ix\mathrm{Obj}(\mathrm{\Gamma }_Q)`$ and $`\stackrel{ˇ}{\mathrm{\Gamma }}_{Q,i}(x,)=0`$ otherwise. The natural composition is given by
$$\stackrel{ˇ}{\mathrm{\Gamma }}_{Q,j}(y,z)\stackrel{ˇ}{\mathrm{\Gamma }}_{Q,i}(x,y)\stackrel{ˇ}{\mathrm{\Gamma }}_{Q,i+j}(x,z),(\psi \varphi )\{\begin{array}{cc}\psi (\tau ^j\varphi )\hfill & \text{ if }\tau ^{i+j}x\mathrm{Obj}(\mathrm{\Gamma }_Q),\hfill \\ 0,\hfill & \text{ else.}\hfill \end{array}$$
By construction, each morphism in $`\stackrel{ˇ}{\mathrm{\Gamma }}_{Q,1}(x,)`$ factors through $`11_{\tau x}\stackrel{ˇ}{\mathrm{\Gamma }}_{Q,1}(x,\tau x)`$ if $`\tau x\mathrm{Obj}(\mathrm{\Gamma }_Q)`$, otherwise $`\stackrel{ˇ}{\mathrm{\Gamma }}_{Q,1}(x,)=0`$. Moreover we have
$$11_{\tau ^nx}\mathrm{}11_{\tau ^2x}11_{\tau x}=11_{\tau ^nx}\stackrel{ˇ}{\mathrm{\Gamma }}_{Q,n}(x,\tau ^nx)\text{ if }\tau ^nx\mathrm{Obj}(\mathrm{\Gamma }_Q),$$
where we consider on the left hand side $`11_{\tau ^ix}\stackrel{ˇ}{\mathrm{\Gamma }}_{Q,1}(\tau ^{i1}x,\tau ^ix)`$. Now, $`\stackrel{ˇ}{\mathrm{\Gamma }}_Q`$ can be described easily by a graded quiver $`\stackrel{ˇ}{A}_Q`$. It has the same vertices and degree 0 arrows as the Auslander-Reiten quiver $`A_Q`$ of $`\mathrm{k}Q`$ (i.e. the quiver of $`\mathrm{\Gamma }_Q`$) moreover there is a degree 1 arrow $`t_x:x\tau x`$ for each $`x\mathrm{Obj}(\mathrm{\Gamma }_Q)`$ with $`\tau x\mathrm{Obj}(\mathrm{\Gamma }_Q)`$. The degree $`0`$ relations are the mesh relations for $`\mathrm{\Gamma }_Q`$. Moreover, each degree $`0`$ arrow $`a:xy`$ with $`y`$ not projective gives rise to a degree $`1`$ relation
$$t_ya(\tau a)t_x.$$
This has to be interpreted as a zero-relation if $`x`$ is projective. A nice way to remember these relations is the following: For each arrow between to vertices which are not both injective there is a (generic) homogeneous length 2 relation in the opposite direction, see also 6.2.
### 3.4. Dynkin quivers.
Let us collect some basic facts about the representation theory of a Dynkin quiver $`Q`$. Define $`𝐢_q:=\underset{¯}{\mathrm{dim}}D\mathrm{k}Q(,q)`$ and $`𝐩_q:=\underset{¯}{\mathrm{dim}}\mathrm{k}Q(q,)`$ for $`qQ_0`$, the dimension vectors of the indecomposable injective and projective $`\mathrm{k}Q`$-modules, respectively. We have then for $`0iN(q)`$
(3.1)
$$\mathrm{\Phi }^i𝐢_q=\underset{¯}{\mathrm{dim}}(\tau _Q^iD\mathrm{k}Q(,q))=\underset{¯}{\mathrm{dim}}(\tau _Q^{iN(q)}\mathrm{k}Q(q,))=\mathrm{\Phi }^{iN(q)}𝐩_{\mu (q)},$$
where $`\mathrm{\Phi }`$ denotes the Coxeter transformation and $`\tau _Q`$ the Auslander-Reiten translate in $`\mathrm{k}Q\mathrm{mod}`$. Next, if $`,`$ denotes the Ringel bilinear form of $`\mathrm{k}Q`$ we have
* $`\underset{¯}{\mathrm{dim}}M,\underset{¯}{\mathrm{dim}}N=dim\mathrm{Hom}_Q(M,N)dim\mathrm{Ext}_Q^1(M,N)`$,
* $`𝐯,𝐰=\mathrm{\Phi }𝐯,\mathrm{\Phi }𝐰=𝐰,\mathrm{\Phi }𝐯`$,
* $`𝐯,𝐢_q=𝐯(q)`$ thus $`𝐯,_{qQ_0}𝐢_q=|𝐯|`$, where $`|𝐯|:=_{qQ_0}𝐯(q)`$,
see for example \[31, 2.4\].
###### 3.4.1 Lemma.
Let $`\tau ^i(0,p)=(i,p)`$ and $`\tau ^j\widehat{\nu }(0,q)`$ belong to $`\mathrm{Obj}(\mathrm{\Gamma }_Q)`$, then
$$dim\mathrm{\Gamma }_Q(\tau ^j\widehat{\nu }(0,q),\tau ^i(0,p))=\{\begin{array}{cc}(\mathrm{\Phi }^{i+j}𝐢_p)(q)\hfill & \text{ if }i+jN(p),\hfill \\ (\mathrm{\Phi }^{ij}𝐩_q)(p)\hfill & \text{ if }i+jN(\mu (q)),\hfill \\ 0\hfill & \text{ if }i+j>\mathrm{min}\{N(p),N(\mu (q))\}.\hfill \end{array}$$
Note that the three cases are not exclusive, however they cover obviously all possibilities for $`(i,j)[0,N(p)]\times [0,N(\mu (q))]`$.
Proof: In the first case we have
(3.2)
$$dim\mathrm{\Gamma }_Q(\tau ^j\widehat{\nu }(0,q),\tau ^i(0,p))=dim\mathrm{\Gamma }_Q(\widehat{\nu }(0,q),\tau ^{i+j}(0,q)).$$
On the other hand, the equivalence $`\mathrm{R}_Q`$ from $`\mathrm{\Gamma }_Q`$ to the category of indecomposable $`\mathrm{k}Q`$-modules commutes with translations
$$\mathrm{R}_Q(\tau ^i(0,q))\tau _Q^iD\mathrm{k}Q(,q)\text{for }0iN(q).$$
Thus (3.2) is equal to
$$dim\mathrm{Hom}_Q(\mathrm{k}Q(q,),\tau _Q^{i+j}D\mathrm{k}Q(,p))=dim\mathrm{Hom}_Q(\tau _Q^j\mathrm{k}Q(q,),\tau _Q^iD\mathrm{k}Q(,p)).$$
Our claim follows now from (3.1) since for a finite-dimensional $`\mathrm{k}Q`$-module $`M`$ we have that
$$\mathrm{Hom}_Q(\mathrm{k}Q(q,),M)M(q)D\mathrm{Hom}_Q(M,D\mathrm{k}Q(,q)).$$
The second case is treated similarly. Finally we have
$$\mathrm{\Gamma }_Q(\tau ^j\widehat{\nu }(0,q),\tau ^i(0,p))\stackrel{~}{\mathrm{\Lambda }}(\tau ^j\widehat{\nu }(0,q),\tau ^i(0,p))\stackrel{~}{\mathrm{\Lambda }}(\widehat{\nu }(0,q),\tau ^{i+jN(p)}\widehat{\nu }(0,\mu (p))).$$
The last term vanishes obviously for $`i+j>N(p)`$. A similar argument shows that $`\mathrm{\Gamma }_Q(\tau ^j\widehat{\nu }(0,q),\tau ^i(0,p))=0`$ for $`i+j>N(\mu (q))`$. $`\mathrm{}`$
###### 3.4.2 Lemma.
If $`N(p)N(q)>j0`$ holds for some $`p,qQ_0`$, then $`\mathrm{\Phi }^j𝐢_p,𝐢_q=0`$.
Proof: Since $`D\mathrm{k}Q(,q)`$ is injective we have
$$\mathrm{\Phi }^j𝐢_p,𝐢_q=dim\mathrm{Hom}_Q(\tau _Q^jD\mathrm{k}Q(,p),D\mathrm{k}Q(,q))=dim\stackrel{~}{\mathrm{\Lambda }}(\tau ^j(0,p),(0,q)).$$
On the other hand for $`N(p)j>N(q)`$ there is no path from $`\tau ^j(0,p)=(N(p)j,\mu (p))`$ to $`(0,q)`$ in $`Q`$. $`\mathrm{}`$
###### 3.5 Proposition.
With the notation of 1.2 and 3.4 we have
(3.3)
$$(\underset{¯}{\mathrm{dim}}\mathrm{I}_Q)(\mu (p))=\underset{qQ_0}{}\underset{i=0}{\overset{N(q)}{}}(i+1)(\mathrm{\Phi }^i𝐢_q)(p)$$
for $`pQ_0`$, and
(3.4)
$$dim=\underset{qQ_0}{}\underset{i=0}{\overset{N(q)}{}}\left(\left(\genfrac{}{}{0pt}{}{N(q)+2}{2}\right)\left(\genfrac{}{}{0pt}{}{i+1}{2}\right)\right)|\mathrm{\Phi }^i𝐢_q|.$$
Proof: For (3.3) we observe first that
$$\underset{¯}{\mathrm{dim}}F_\lambda J^{}D\mathrm{\Gamma }_Q(,(j,q))(\mu (p))=\underset{i=j}{\overset{N(q)}{}}(\mathrm{\Phi }^i𝐢_q)(p)$$
for $`(j,q)\mathrm{Obj}(\mathrm{\Gamma }_Q)`$ by 3.4.1 and the definition of the push-down $`F_\lambda `$. Now (3.3) follows from the definition of $`\mathrm{I}_Q`$.
For (3.4) we observe first that
$$|\underset{¯}{\mathrm{dim}}F_\lambda J^{}\mathrm{\Gamma }_Q(\tau ^j\widehat{\nu }(0,q),)|=\underset{k=j}{\overset{N(\mu (q))}{}}|\mathrm{\Phi }^k𝐩_q|$$
for $`\tau ^j\widehat{\nu }(0,q)\mathrm{Obj}(\mathrm{\Gamma }_Q)`$ again by 3.4.1 and the definition of the push-down. Now, by construction of $``$ we have
$`dim`$ $`={\displaystyle \underset{qQ_0}{}}{\displaystyle \underset{i=0}{\overset{N(\mu (q))}{}}}{\displaystyle \underset{j=0}{\overset{i}{}}}|\underset{¯}{\mathrm{dim}}F_\lambda J^{}\mathrm{\Gamma }_Q(\tau ^j\widehat{\nu }(0,q),)|`$
$`={\displaystyle \underset{qQ_0}{}}{\displaystyle \underset{i=0}{\overset{N(\mu (q))}{}}}{\displaystyle \underset{j=0}{\overset{i}{}}}{\displaystyle \underset{k=i}{\overset{N(\mu (q))}{}}}|\mathrm{\Phi }^k𝐩_q|`$
$`={\displaystyle \underset{qQ_0}{}}{\displaystyle \underset{i=0}{\overset{N(q)}{}}}{\displaystyle \underset{j=0}{\overset{i}{}}}{\displaystyle \underset{k=i}{\overset{N(q)j}{}}}|\mathrm{\Phi }^k𝐢_q|={\displaystyle \underset{qQ_0}{}}{\displaystyle \underset{i=0}{\overset{N(q)}{}}}{\displaystyle \underset{j=i}{\overset{N(q)+1}{}}}j|\mathrm{\Phi }^i𝐢_q|.`$
$`\mathrm{}`$
###### 3.6 Proposition.
Let $`Q`$ be a Dynkin quiver. Then for $`𝐯=_{qQ_0}_{d=0}^{N(q)}(i+1)\mathrm{\Phi }^i𝐢_q`$ holds
$$𝐯,𝐯=\underset{pQ_0}{}\underset{d=0}{\overset{N(p)}{}}\left(\left(\genfrac{}{}{0pt}{}{N(p)+2}{2}\right)\left(\genfrac{}{}{0pt}{}{d+1}{2}\right)\right)|\mathrm{\Phi }^d𝐢_p|.$$
Proof: For convenience let us write $`N(p,q):=\{0,1,\mathrm{},N(p)\}\times \{0,1,\mathrm{},N(q)\}`$ and $`N(p,q,):=\{(i,j)N(p,q)ij\}`$, similarly $`N(p,q,<):=N(p,q)N(p,q,)`$. Now we have
$`𝐯,𝐯=`$ $`{\displaystyle \underset{p,qQ_0}{}}{\displaystyle \underset{(i,j)N(p,q)}{}}(i+1)(j+1)\mathrm{\Phi }^i𝐢_p,\mathrm{\Phi }^j𝐢_q`$
$`=`$ $`{\displaystyle \underset{p,qQ_0}{}}{\displaystyle \underset{(i,j)N(p,q,)}{}}(i+1)(j+1)\mathrm{\Phi }^i𝐢_p,\mathrm{\Phi }^j𝐢_q`$
$`{\displaystyle \underset{p,qQ_0}{}}{\displaystyle \underset{(i,j)N(p,q,<)}{}}(i+1)(j+1)\mathrm{\Phi }^j𝐢_q,\mathrm{\Phi }^{i+1}𝐢_p`$
$`=`$ $`{\displaystyle \underset{p,qQ_0}{}}{\displaystyle \underset{(i,j)N(p,q,)}{}}(i+1)\mathrm{\Phi }^i𝐢_p,\mathrm{\Phi }^j𝐢_q`$
$`{\displaystyle \underset{p,qQ_0}{}}{\displaystyle \underset{i=N(q)+1}{\overset{N(p)}{}}}(i+1)(N(q)+1)\mathrm{\Phi }^i𝐢_p,\mathrm{\Phi }^{N(q)+1}𝐢_q`$
here, the second sum vanishes by Lemma 3.4.2, thus from 3.4
$`=`$ $`{\displaystyle \underset{p,qQ_0}{}}{\displaystyle \underset{(i,j)N(p,q,)}{}}(i+1)\mathrm{\Phi }^{ij}𝐢_p,𝐢_q`$
$`=`$ $`{\displaystyle \underset{p,qQ_0}{}}{\displaystyle \underset{d=0}{\overset{N(p)}{}}}\left({\displaystyle \underset{k=d}{\overset{\mathrm{min}\{N(p),N(q)+d\}}{}}}(k+1)\right)\mathrm{\Phi }^d𝐢_p,𝐢_q`$
$`=`$ $`{\displaystyle \underset{pQ_0}{}}{\displaystyle \underset{d=0}{\overset{N(p)}{}}}{\displaystyle \underset{k=d}{\overset{N(p)}{}}}(k+1)\mathrm{\Phi }^d𝐢_p,{\displaystyle \underset{qQ_0}{}}𝐢_q`$
$`{\displaystyle \underset{p,qQ_0}{}}{\displaystyle \underset{d=0}{\overset{N(p)}{}}}{\displaystyle \underset{k=N(p)+d}{\overset{N(p)}{}}}(k+1)\mathrm{\Phi }^d𝐢_p,𝐢_q.`$
Here again, the second sum vanishes by Lemma 3.4.2. Finally, $`\mathrm{\Phi }^d𝐢_p,_{qQ_0}𝐢_q=|\mathrm{\Phi }^d𝐢_p|`$ as observed at the beginning of 3.4, so our claim follows. $`\mathrm{}`$
### 3.7. Proof of Theorem 1.
The first claim follows directly from the construction of $`\stackrel{ˇ}{\mathrm{\Gamma }}_Q`$, Lemma 3.2 and the fact that
$$\mathrm{Hom}_\mathrm{\Lambda }(F_\lambda N,F_\lambda N)_i\mathrm{Hom}_{\stackrel{~}{\mathrm{\Lambda }}}(N,N^{(i)}).$$
For the second claim we have to show by \[8, Lemma 1\] that $`\underset{¯}{\mathrm{dim}}\mathrm{I}_Q,\underset{¯}{\mathrm{dim}}\mathrm{I}_Q=dim`$. This follows from the dimension formulas (3.3) and (3.4) of Proposition 3.5 together with Proposition 3.6.
## 4. The cluster algebras $`[G^{e,w_0}]`$ and $`[N]`$
In the next two sections we will use the setup from 1.4 and 1.5. We will need the following result, see for instance \[24, Section 4.4.3\]
###### 4.1 Lemma.
Let $`𝐢=(i_1,i_2,\mathrm{},i_m)`$ be a reduced expression for some element $`w^1W`$, and $`L(\lambda )`$ an irreducible representation of highest weight $`\lambda `$ for $`G`$. If $`u_\lambda `$ is a highest weight vector for $`L(\lambda )`$ then we have
$$\overline{w}u_\lambda =f_{i_m}^{(b_m)}\mathrm{}f_{i_2}^{(b_2)}f_{i_1}^{(b_1)}(u_\lambda )\text{ and }f_{i_m}(\overline{w}u_\lambda )=0.$$
Here $`b_1:=\lambda (h_{i_1})`$, $`b_k:=(s_{i_{k1}}\mathrm{}s_{i_2}s_{i_1}(\lambda ))(h_{i_k})`$ for $`2km`$ and $`f_{i_k}^{(b_k)}:=\frac{1}{b_k!}f_{i_k}^{b_k}U(𝔤)`$.
For $`vL(\lambda )`$, we write $`f_j^{\mathrm{max}}(v):=(1/m!)f_j^m(v)`$ where $`m=\mathrm{max}\{pf_j^p(v)0\}`$. With this notation we could restate the equality of the lemma as
$$\overline{w}u_\lambda =f_{i_m}^{\mathrm{max}}\mathrm{}f_{i_2}^{\mathrm{max}}f_{i_1}^{\mathrm{max}}(u_\lambda ).$$
### 4.2.
The group $`G`$ has Bruhat decompositions with respect to $`B`$ and $`B_{}`$, namely
$$G=\underset{uW}{}B\overline{u}B=\underset{vW}{}B_{}\overline{v}B_{}.$$
The intersection of two cells $`G^{u,v}=(B\overline{u}B)(B_{}\overline{v}B_{})`$ is called a double Bruhat cell.
In particular taking $`u=e`$, the unit in $`W`$, and $`v=w_0`$ we obtain $`G^{e,w_0}=B(B_{}\overline{w}_0B_{})`$, the intersection of $`B`$ with the big cell relative to $`B_{}`$. By \[4, Proposition 2.8\], $`G^{e,w_0}`$ consists of all elements $`x`$ of $`B`$ such that $`\mathrm{\Delta }_{\varpi _i,w_0(\varpi _i)}(x)0`$ for every $`i`$. This is a Zariski open subset of $`B`$, hence an algebraic variety of dimension $`n+r`$, where $`r=|\mathrm{\Pi }|`$ is the number of positive roots associated to the Dynkin type of $`G`$. Moreover, we see that the algebra of regular functions $`[G^{e,w_0}]`$ is obtained from $`[B]`$ by adjoining formal inverses to the functions $`\mathrm{\Delta }_{\varpi _i,w_0(\varpi _i)}`$.
On the other hand, $`N`$ can be described as the subvariety of $`B`$ given by the equations
$$\mathrm{\Delta }_{\varpi _i,\varpi _i}(x)=1,(1in).$$
Hence the algebra $`[N]`$ is the quotient of $`[B]`$ by the ideal generated by the elements $`(\mathrm{\Delta }_{\varpi _i,\varpi _i}1)_{i=1,2,\mathrm{},n}`$.
### 4.3.
In , Fomin and Zelevinsky have introduced a transcendence basis $`F(𝐢)`$ of the field of rational functions $`(G^{e,w_0})`$ consisting of certain generalized minors. In Berenstein, Fomin and Zelevinsky have shown that each $`F(𝐢)`$ can be taken as the initial cluster for a natural upper cluster algebra structure on the ring $`[G^{e,w_0}]`$. We are now going to recall their construction.
#### 4.3.1.
We add $`n`$ additional letters $`i_n,\mathrm{},i_1`$ at the beginning of $`𝐢`$, where $`i_j=j`$, and obtain an $`(r+n)`$-tuple
$$(i_n,\mathrm{},i_1,i_1,\mathrm{},i_r)=(n,\mathrm{},1,i_1,\mathrm{},i_r).$$
For $`k[n,1][1,r]`$ let
$$k^+=\{\begin{array}{cc}r+1\hfill & \text{if }|i_l||i_k|\text{ for all }l>k,\hfill \\ \mathrm{min}\{ll>k\text{ and }|i_l|=|i_k|\}\hfill & \text{otherwise}.\hfill \end{array}$$
Then $`k`$ is called $`𝐢`$-exchangeable if $`k`$ and $`k^+`$ are in $`[1,r]`$. Let $`e(𝐢)[1,r]`$ be the set of $`𝐢`$-exchangeable elements. One easily checks that $`e(𝐢)`$ contains $`rn`$ elements. More precisely, the set of indices $`i_k`$ for $`k[1,r]e(𝐢)`$ is exactly $`[1,n]`$.
#### 4.3.2.
Next, one defines a quiver $`\stackrel{~}{A}_𝐢`$ with set of vertices $`[n,1][1,r]`$. Assume that $`k`$ and $`l`$ are vertices such that the following hold:
* $`k<l`$;
* $`\{k,l\}e(𝐢)\mathrm{}`$.
There is an arrow $`kl`$ in $`\stackrel{~}{A}_𝐢`$ if and only if $`k^+=l`$, and there is an arrow $`lk`$ if and only if $`l<k^+<l^+`$ and $`a_{|i_k|,|i_l|}=1`$. Here, $`(a_{ij})_{1i,jn}`$ denotes the Cartan matrix of the root system of $`G`$. By definition these are all the arrows of $`\stackrel{~}{A}_𝐢`$.
###### 4.3.3 Remark.
If $`𝐢`$ is a reduced expression for $`w_0`$ which is adapted to a Dynkin quiver $`Q`$, then it is easy to obtain $`\stackrel{~}{A}_𝐢`$ from the Auslander-Reiten quiver $`A_Q`$. See the examples in 6.5 and 6.6.
#### 4.3.4.
Now define an $`(r+n)\times (rn)`$-matrix
$$\stackrel{~}{B}(𝐢)=(b_{kl})$$
as follows. The columns of $`\stackrel{~}{B}(𝐢)`$ are indexed by the elements in $`e(𝐢)`$, and the rows by $`[n,1][1,r]`$. Set
$$b_{kl}=\{\begin{array}{cc}1\hfill & \text{if there is an arrow }kl\text{ in }\stackrel{~}{A}_𝐢,\hfill \\ 1\hfill & \text{if there is an arrow }lk\text{ in }\stackrel{~}{A}_𝐢,\hfill \\ 0\hfill & \text{otherwise}.\hfill \end{array}$$
#### 4.3.5.
For $`k[n,1][1,r]`$ one defines a generalized minor $`\mathrm{\Delta }(k,𝐢)`$ as follows. For $`k[1,r]`$ set $`v_{>k}=s_{i_r}s_{i_{r1}}\mathrm{}s_{i_{k+1}}`$ and for $`k[n,1]`$ put $`v_{>k}=w_0`$. Then define
$$\mathrm{\Delta }(k,𝐢)=\mathrm{\Delta }_{\varpi _{|i_k|},v_{>k}(\varpi _{|i_k|})}.$$
Since $`s_j(\varpi _i)=\varpi _i`$ for $`ji`$, it is easy to see that if $`k[1,r]`$ is not exchangeable then $`\mathrm{\Delta }(k,𝐢)=\mathrm{\Delta }_{\varpi _{i_k},\varpi _{i_k}}`$. On the other hand for $`i[n,1]`$ we have $`\mathrm{\Delta }(i,𝐢)=\mathrm{\Delta }_{\varpi _i,w_0(\varpi _i)}`$.
It is known \[12, Theorem 1.12\] that this collection of $`n+r`$ minors is a transcendence basis of the field $`(G^{e,w_0})`$, for any reduced expression $`𝐢`$ of $`w_0`$. By 4.2, we see that if we remove from this collection the $`n`$ minors $`\mathrm{\Delta }_{\varpi _i,\varpi _i}`$ we obtain a transcendence basis of the field $`(N)`$.
#### 4.3.6.
Let $``$ be the field of rational functions over $``$ in $`n+r`$ independent variables $`\stackrel{~}{𝐱}=(x_n,\mathrm{},x_1,x_1,\mathrm{},x_r)`$. Let $`\overline{𝒜}(𝐢)_{}`$ denote the upper cluster algebra associated to the seed $`(\stackrel{~}{𝐱},\stackrel{~}{B}(𝐢))`$, a subalgebra of $``$ (see \[4, Definition 1.6\]). Here the non-exchangeable indices in $`[n,1][1,r]`$ label the generators of the coefficient group (see \[4, §2.2\]).
Berenstein, Fomin and Zelevinsky then show that the isomorphism of fields $`\phi _𝐢`$ from $``$ to $`(G^{e,w_0})`$ defined by
$$\phi _𝐢(x_k)=\mathrm{\Delta }(k,𝐢),(k[n,1][1,r]),$$
restricts to an algebra isomorphism $`\overline{𝒜}(𝐢)_{}[G^{e,w_0}]`$, see \[4, Theorem 2.10\].
Note that by varying the reduced expression $`𝐢`$ we obtain a priori several cluster algebra structures on $`[G^{e,w_0}]`$, but according to \[4, Remark 2.14\] all these structures coincide and give rise to the same cluster variables and clusters. Note also that in type $`𝖠_n`$, the upper cluster algebra $`\overline{𝒜}(𝐢)_{}`$ coincides with the cluster algebra $`𝒜(𝐢)_{}`$, see \[4, Remark 2.18\].
### 4.4.
Let $`\stackrel{~}{𝐱}^{}`$ be the subset of $`\stackrel{~}{𝐱}`$ obtained by removing the variables indexed by the $`n`$ non-exchangeable elements in $`[1,r]`$. Let $`^{}`$ be the field of rational functions over $``$ in the $`r`$ variables of $`\stackrel{~}{𝐱}^{}`$. Finally, let $`\stackrel{~}{B}(𝐢)^{}`$ be the matrix obtained from $`\stackrel{~}{B}(𝐢)`$ by removing the rows labelled by the $`n`$ non-exchangeable elements in $`[1,r]`$, and let $`\overline{𝒜}(𝐢)_{}^{}`$ denote the upper cluster algebra associated to the seed $`(\stackrel{~}{𝐱}^{},\stackrel{~}{B}(𝐢)^{})`$, a subalgebra of $`^{}`$. By 4.2, we see that the isomorphism of fields $`\phi _𝐢^{}:^{}(N)`$ defined by
$$\phi _𝐢^{}(x_k)=\mathrm{\Delta }(k,𝐢),(x_k\stackrel{~}{𝐱}^{}),$$
restricts to an algebra isomorphism $`\overline{𝒜}(𝐢)_{}^{}[N]`$.
Clearly the same remarks as in 4.3.6 apply to the cluster algebra structures of $`[N]`$.
### 4.5.
We shall now describe in representation theoretic terms the restriction to $`N`$ of the regular function $`\mathrm{\Delta }(k,𝐢)`$. Let $`L(\varpi _i)`$ be the fundamental irreducible $`𝔤`$-module with highest weight $`\varpi _i`$. Fix a highest weight vector $`u_{\varpi _i}`$.
It is well-known that $`L(\varpi _i)`$ can be realized in a canonical way as a subspace of the vector space $`[N]`$ by restricting the summand $`L(\varpi _i)`$ of $`[N_{}\backslash G]`$ to $`[N]`$. Thus, the highest weight vector $`u_{\varpi _i}`$ becomes identified with the constant regular function $`\mathrm{𝟏}`$. Using this identification we obtain the following lemma:
###### 4.5.1 Lemma.
For $`k[n,1][1,r]`$ we have
$$\mathrm{\Delta }(k,𝐢)=\{\begin{array}{cc}f_{i_r}^{\mathrm{max}}f_{i_{r1}}^{\mathrm{max}}\mathrm{}f_{i_{k+1}}^{\mathrm{max}}(u_{\varpi _{i_k}})\hfill & \text{ if }k[1,r],\hfill \\ f_{i_r}^{\mathrm{max}}f_{i_{r1}}^{\mathrm{max}}\mathrm{}f_{i_1}^{\mathrm{max}}(u_{\varpi _{|i_k|}})\hfill & \text{ if }k[n,1].\hfill \end{array}$$
In particular for $`k[n,1]`$, the minor $`\mathrm{\Delta }(k,𝐢)=\mathrm{\Delta }_{\varpi _k,w_0(\varpi _k)}`$ is a lowest weight vector of $`L(\varpi _k)`$.
Proof: This follows from \[5, p. 150–151\] and 4.1 by restricting regular functions on $`G`$ to the subgroup $`N`$, see also \[6, p. 113\]. $`\mathrm{}`$
## 5. Cluster variables and semicanonical basis
### 5.1.
Let $`Q`$ be a Dynkin quiver such that $`|Q|`$ is the diagram of $`G`$. In this section we prove that if $`𝐢`$ is a reduced expression for $`w_0`$ which is adapted to $`Q`$, the minors $`\mathrm{\Delta }(k,𝐢)[N]`$ coincide with certain dual semicanonical basis vectors coming from the injective modules of the Auslander algebra of $`Q`$. In particular, this shows that the set of minors $`\{\mathrm{\Delta }(k,𝐢)\}`$ depends only on $`Q`$, not on the choice of a particular expression $`𝐢`$ adapted to $`Q`$.
### 5.2.
Recall that by pushing down the injective modules of the Auslander algebra of $`Q`$ we obtain a set of indecomposable rigid modules $`\mathrm{I}(j,𝐢)`$ over the preprojective algebra $`\mathrm{\Lambda }`$, see 1.2 and 1.3. Since these modules are rigid, they have an open orbit in their module variety, and the closure of this orbit is an irreducible component. Therefore the module $`\mathrm{I}(j,𝐢)`$ can be used to label an element $`\rho _{\mathrm{I}(j,𝐢)}`$ of the dual semicanonical basis (see , \[19, Section 7.2\]).
### 5.3.
The proof of Theorem 2 will make use of certain results of that we shall now recall. Let $`I_i`$ denote the injective envelope of the simple $`\mathrm{\Lambda }`$-module $`S_i`$ with dimension vector $`𝐞_i`$, thus $`I_i=D\mathrm{\Lambda }(,i)`$. In the fundamental $`𝔤`$-module $`L(\varpi _i)`$ was realized in terms of the lattice of submodules of $`I_i`$. This goes as follows.
Let $``$ denote Lusztig’s algebra of constructible functions on the varieties of finite-dimensional $`\mathrm{\Lambda }`$-modules, and let $`^{}`$ be its graded Hopf dual, an algebra isomorphic to $`[N]`$. For a finite-dimensional $`\mathrm{\Lambda }`$-module $`X`$, let $`\delta _X`$ denote the linear form on $``$ obtained by evaluation at $`X`$. Then in the identification $`[N]^{}`$, the subspace $`L(\varpi _i)`$ gets identified to the subspace of $`^{}`$ spanned by the linear forms $`\delta _X`$ where $`X`$ runs over the lattice of submodules of $`I_i`$, and one has explicit formulas for the action of the Chevalley generators of $`𝔤`$ on each vector $`\delta _X`$ \[18, Theorem 3\]. In particular $`\delta _{I_i}=\rho _{I_i}`$ is a lowest weight vector of $`L(\varpi _i)`$, and $`\delta _0`$, where $`0`$ means the zero submodule of $`I_i`$, is a highest weight vector.
Let $`X`$ be a submodule of $`I_i`$. We have a short exact sequence of $`\mathrm{\Lambda }`$-modules
$$0XI_i\stackrel{𝑝}{}Y0,$$
where $`Y`$ is determined up to isomorphism by the isomorphism class of $`X`$, see \[18, Lemma 1\]. For $`j[1,n]`$ let $`m_j`$ denote the multiplicity of $`S_j`$ in the socle of $`Y`$, and let $`X_j`$ be the unique submodule of $`I_i`$ such that $`XX_jI_i`$ and $`X_j/X`$ is isomorphic to $`S_j^{m_j}`$. Thus $`X_j`$ is the pullback of $`p`$ and the inclusion of $`S_j^{m_j}`$ into $`Y`$.
###### 5.4 Lemma.
With the above notation, we have $`f_j^{\mathrm{max}}(\delta _X)=f_j^{(m_j)}(\delta _X)=\delta _{X_j}`$.
Proof: By \[18, Theorem 3 (ii)\], we have that
$$f_j^k(\delta _X)=_{𝔣=(X=X(1)\mathrm{}X(k))}\delta _{X(k)}$$
where the integral is over the variety of flags $`𝔣`$ of submodules of $`I_i`$ such that $`X(s)/X(s1)`$ is isomorphic to $`S_j`$ for all $`1<sk`$. For $`k>m_j`$ this variety is empty by definition of $`m_j`$, hence $`f_j^k(\delta _X)=0`$. For $`k=m_j`$, all flags $`𝔣`$ have their last step equal to $`X_j`$. Moreover since $`X_j/XS_j^{m_j}`$ this variety is isomorphic to the variety of complete flags in $`^{m_j}`$, whose Euler characteristic is $`m_j!`$. Hence $`f_j^{m_j}(\delta _X)=m_j!\delta _{X_j}`$, as claimed. $`\mathrm{}`$
###### 5.5 Lemma.
Let $`k[n,1]e(𝐢)`$ and $`j=\theta ^1(k)`$. Then, if $`i=|i_k|`$, the module $`\mathrm{I}(j,𝐢)`$ is a submodule of $`I_i`$ and in $`L(\varpi _i)`$ there holds
$$\delta _{\mathrm{I}(j,𝐢)}=\{\begin{array}{cc}f_{i_r}^{\mathrm{max}}f_{i_{r1}}^{\mathrm{max}}\mathrm{}f_{i_{k+1}}^{\mathrm{max}}(\delta _0)\hfill & \text{ if }ke(𝐢),\hfill \\ f_{i_r}^{\mathrm{max}}f_{i_{r1}}^{\mathrm{max}}\mathrm{}f_{i_1}^{\mathrm{max}}(\delta _0)\hfill & \text{ if }k[n,1].\hfill \end{array}$$
Proof: Recall that we defined the function $`\theta `$ in 1.5. If $`k=i[n,1]`$, then $`\mathrm{I}(j,𝐢)=I_i`$. On the other hand in this case the product $`f_{i_r}^{\mathrm{max}}f_{i_{r1}}^{\mathrm{max}}\mathrm{}f_{i_1}^{\mathrm{max}}`$ maps $`\delta _0`$ to the lowest weight vector of $`L(\varpi _i)`$, that is, to $`\delta _{I_i}`$, as required.
If $`ke(𝐢)`$, then $`\mathrm{R}_Q(x(j))`$ belongs to the $`\tau `$-orbit of $`\mathrm{R}_Q(x(k))`$ by definition of $`\theta `$, therefore to the $`\tau `$-orbit of $`D\mathrm{k}Q(,i)`$ see 1.3. It follows that $`\mathrm{I}(j,𝐢)`$ is a submodule of $`I_i`$, see 1.2.
More precisely, $`j=\theta ^1(k)=\mathrm{min}\{l[k+1,r]i_l=i\}`$. We conclude that in $`L(\varpi _i)`$ holds $`f_{i_{j1}}^{\mathrm{max}}\mathrm{}f_{i_{k+1}}^{\mathrm{max}}(\delta _0)=\delta _0`$ by 5.4, since the socle of $`I_i`$ is $`S_i`$.
Now consider for $`l[j1,r]`$ the $`\stackrel{~}{\mathrm{\Lambda }}`$-submodule $`\stackrel{~}{I}_{x(j)}(l,𝐢)`$ of $`D\stackrel{~}{\mathrm{\Lambda }}(,x(j))`$ with
$$\stackrel{~}{I}_{x(j)}(l,𝐢)(x(m)):=\{\begin{array}{cc}D\stackrel{~}{\mathrm{\Lambda }}(x(m),x(j))\hfill & \text{ if }ml,\hfill \\ 0\hfill & \text{ else.}\hfill \end{array}$$
With $`I_{x(j)}(l,𝐢):=F_\lambda \stackrel{~}{I}_{x(j)}`$ we see that $`I_{x(j)}(r,𝐢)=\mathrm{I}(j,𝐢)`$ is a submodule of $`I_i`$ since $`F_\lambda D\stackrel{~}{\mathrm{\Lambda }}(,x(j))I_i`$.
Using Lemma 5.4 we conclude that in $`L(\varpi _i)`$ we have
$$f_{i_l}^{\mathrm{max}}(\delta _{I_{x(j)}(l1,𝐢)})=\delta _{I_{x(j)}(l,𝐢)}\text{ for }l[j,r].$$
This yields that $`\delta _{\mathrm{I}(j,𝐢)}`$ is the extremal vector of $`L(\varpi _i)`$ with weight $`s_{i_r}s_{i_{r1}}\mathrm{}s_{i_{k+1}}(\varpi _i)`$, and the lemma follows. $`\mathrm{}`$
We can now finish the proof of Theorem 2. Using Lemma 4.5.1 and Lemma 5.5 we obtain that $`\mathrm{\Delta }(k,𝐢)=\delta _{\mathrm{I}(j,𝐢)}`$. But since $`\mathrm{I}(j,𝐢)`$ is rigid, its orbit is open and we have $`\delta _{\mathrm{I}(j,𝐢)}=\rho _{\mathrm{I}(j,𝐢)}`$.
## 6. Examples
Our running example will be the Dynkin quiver $`Q`$ of type $`𝖣_5`$ with the following orientation:
### 6.1. $`\mathrm{\Gamma }_Q`$ for $`𝖣_5`$.
We show the quiver $`A_Q`$ of $`\mathrm{\Gamma }_Q`$:
The isomorphism classes of the indecomposable representations of $`Q`$ correspond to the vertices of $`A_Q`$. Each such representation is uniquely determined by its dimension vector:
The dimension vector of the injective $`\mathrm{\Gamma }_Q`$-module $`D\mathrm{\Gamma }_Q(,(0,q))`$ is given by the $`q`$-components of the corresponding dimension vectors. The dimension vector of $`D\mathrm{\Gamma }_Q(,\tau ^j(0,q))`$ is obtained from this by “translation and cut-off”. Here we show for example $`\underset{¯}{\mathrm{dim}}D\mathrm{\Gamma }_Q(,(0,3))`$ and $`\underset{¯}{\mathrm{dim}}D\mathrm{\Gamma }_Q(,\tau (0,3))`$:
### 6.2. The category $`\stackrel{ˇ}{\mathrm{\Gamma }}_Q`$.
We display here the quiver of the (graded) category $`\stackrel{ˇ}{\mathrm{\Gamma }}_Q`$ associated to a quiver of type $`𝖣_5`$ with the same orientation as above. The arrows of degree one are dotted.
Recall that also the relations are easy to read off: For each arrow $`\alpha :xy`$ we have the corresponding mesh relation (2.1) from $`y`$ to $`x`$ if $`\alpha `$ is of degree $`1`$. Otherwise, if $`x`$ is not an “injective” vertex, (i.e. here if $`x\{(0,1),(0,3),(0,5)\}`$) there are one or two paths of length $`2`$ (and degree $`1`$) from $`y`$ to $`x`$. The first case occurs when $`y`$ is a “projective” vertex (i.e. here if $`y\{(3,2),(3,3),(2,5)\}`$) and the unique path of length $`2`$ from $`y`$ to $`x`$ is a zero-relation, otherwise the two paths form a commutativity relation.
### 6.3. A projective-injective $`\stackrel{ˇ}{\mathrm{\Gamma }}_Q`$-module.
We display for $`\stackrel{ˇ}{\mathrm{\Gamma }}_Q`$ as in 6.2 the projective module $`\stackrel{ˇ}{\mathrm{\Gamma }}_Q((0,5),)`$. Since $`(0,5)`$ is an injective vertex of $`A_Q`$ we have $`\stackrel{ˇ}{\mathrm{\Gamma }}_Q((0,5),)D\stackrel{ˇ}{\mathrm{\Gamma }}_Q(,(0,\mu (5))`$. Moreover, the projective modules $`\stackrel{ˇ}{\mathrm{\Gamma }}_Q((j,5),)`$ for $`0j2`$ are easily found as submodules of $`\stackrel{ˇ}{\mathrm{\Gamma }}_Q((0,5),)`$.
Each entry $`(n,q)`$ represents a basis vector which corresponds to a (graded) simple composition factor of this type. The arrows indicate as usual the action of $`\stackrel{ˇ}{\mathrm{\Gamma }}_Q`$.
### 6.4. Dimensions.
We include the result of some calculations of
$$d(Q):=dim\mathrm{End}_\mathrm{\Lambda }(\mathrm{I}_Q)$$
for specific orientations. This can be done quite easily on a computer using the formula (3.4).
$`A_n`$ $`:\text{}`$ $`d(A_n)`$ $`=2\left({\displaystyle \genfrac{}{}{0pt}{}{n}{5}}\right)+7\left({\displaystyle \genfrac{}{}{0pt}{}{n}{4}}\right)+9\left({\displaystyle \genfrac{}{}{0pt}{}{n}{3}}\right)+5\left({\displaystyle \genfrac{}{}{0pt}{}{n}{2}}\right)+n,`$
$`D_n`$ $`:\text{}`$ $`d(D_n)`$ $`=27\left({\displaystyle \genfrac{}{}{0pt}{}{n}{5}}\right)+43\left({\displaystyle \genfrac{}{}{0pt}{}{n}{4}}\right)+19\left({\displaystyle \genfrac{}{}{0pt}{}{n}{3}}\right)+2\left({\displaystyle \genfrac{}{}{0pt}{}{n}{2}}\right),`$
$`E_n`$ $`:\text{}`$ $`d(E_n)`$ $`=\{\begin{array}{cc}2444\hfill & \text{ if }n=6,\hfill \\ 13130\hfill & \text{ if }n=7,\hfill \\ 107114\hfill & \text{ if }n=8.\hfill \end{array}`$
### 6.5. An adapted ordering on $`\mathrm{Obj}(\mathrm{\Gamma }_Q)`$:
According to 1.3 we obtain for $`w_0`$ the following reduced expression which is adapted to $`Q`$:
$$𝐢=(4,2,1,3,5,4,2,1,3,5,4,2,1,3,5,4,2,3,4,1)$$
### 6.6. The quiver $`\stackrel{~}{A}_𝐢`$.
For the adapted expression $`𝐢`$ from 6.5 we obtain
$$\stackrel{~}{A}_𝐢=\text{}$$
The exchangeable vertices are $`\{1,2,3,4,5,6,7,8,9,10,11,12,13,14,16\}`$
## 7. Appendix: Using the triangulated structure
###### 7.1 Proposition (Freyd/Heller).
Let $`𝒟`$ be a triangulated $`\mathrm{k}`$-category with suspension functor $`\mathrm{\Sigma }`$.
* The category $`𝒟\mathrm{mod}`$ is a Frobenius category. Thus the stable category $`𝒟\text{-}\underset{¯}{\mathrm{mod}}`$ is triangulated with suspension functor $`\mathrm{\Omega }_𝒟^1`$, the inverse of Heller’s loop functor.
* In $`𝒟\text{-}\underset{¯}{\mathrm{mod}}`$ we have a functorial isomorphism $`M^\mathrm{\Sigma }\mathrm{\Omega }_𝒟^3M`$.
Part (a) is from \[13, Section 3\], see also \[30, Chapter 5\] for a modern treatment. Part (b) is a special case of \[22, §16\], see also \[26, Proposition B.2\].
###### 7.2 Remarks.
(a) We may consider $`𝒟`$ as a $`𝒟\text{-}𝒟`$-bimodule, i.e. a functor $`𝒟^{\text{op}}\times 𝒟\mathrm{k}\mathrm{mod}`$. Similarly, $`D𝒟`$ with $`D𝒟(a,b):=\mathrm{Hom}_\mathrm{k}(𝒟(b,a),\mathrm{k})`$ is a $`𝒟\text{-}𝒟`$-bimodule.
(b) Suppose that $`𝒟`$ admits Auslander-Reiten triangles with translate $`\tau :𝒟𝒟`$. In this case we set $`\nu :=\mathrm{\Sigma }\tau `$. Then the Auslander-Reiten formula $`𝒟(x,\mathrm{\Sigma }y)\mathrm{Hom}_k(𝒟(y,\tau x),\mathrm{k})`$ may be interpreted as an isomorphism of bimodules
(7.1)
$$D𝒟𝒟^{\nu ^1}\text{where }𝒟^{\nu ^1}(a,b):=𝒟(a,\nu b).$$
We conclude that
$$𝒩M:=D𝒟_𝒟M𝒟^{\nu ^1}_𝒟MM^{\nu ^1}$$
is a Nakayama functor for $`𝒟\mathrm{mod}`$ and $`\nu ^1`$ the corresponding Nakayama automorphism for $`𝒟`$, see for example \[14, Section 2\]. In this situation we will write
$$M^{(i)}:=M^{\tau ^i}.$$
(c) If moreover $`𝒟`$ is locally bounded, then $`𝒟\text{-}\underset{¯}{\mathrm{mod}}`$ admits Auslander-Reiten triangles with translation $`\tau _𝒟=\mathrm{\Omega }_𝒟^2𝒩`$. With Proposition 7.1 we obtain the functorial isomorphisms
(7.2)
$$M^{(1)}\tau _𝒟\mathrm{\Omega }_𝒟M\text{ and }\tau _𝒟^3MM^{\mathrm{\Sigma }\tau ^3}\text{ (in }𝒟\text{-}\underset{¯}{\mathrm{mod}}\text{)}.$$
In fact, we have $`\tau _𝒟\mathrm{\Omega }_D\mathrm{\Omega }^3𝒩\mathrm{?}^{\mathrm{\Sigma }\nu ^1}=\mathrm{?}^{\tau ^1}`$ and $`\tau _𝒟^3=\mathrm{\Omega }_𝒟^6𝒩^3\mathrm{?}^{\mathrm{\Sigma }^2\nu ^3}\mathrm{?}^{\mathrm{\Sigma }\tau ^3}`$. (d) If we consider in this context right modules (i.e. contravariant functors), we get in $`\underset{¯}{\mathrm{mod}}\text{-}𝒟`$ a functorial isomorphism $`M^\mathrm{\Sigma }\mathrm{\Omega }_𝒟^3M`$ and consequently
$$M^{(1)}\tau _𝒟\mathrm{\Omega }_𝒟M\text{ and }\tau _𝒟^3MM^{\mathrm{\Sigma }\tau ^3}\text{ (in }\underset{¯}{\mathrm{mod}}\text{-}𝒟\text{)}.$$
### 7.3. Derived categories.
It follows from Happel’s description \[21, I.5.6\] of the derived category $`𝒟^b(\mathrm{k}Q^{\text{op}}):=𝒟^b(\mathrm{k}Q^{\text{op}}\mathrm{mod})`$ that we have a natural equivalence
$$\stackrel{~}{\mathrm{\Lambda }}\mathrm{inj}𝒟^b(\mathrm{k}Q^{\text{op}}).$$
In particular, $`\stackrel{~}{\mathrm{\Lambda }}\mathrm{inj}`$ is a triangulated category which admits Auslander-Reiten triangles. The suspension functor resp. the Auslander-Reiten translate are
$$\mathrm{\Sigma }I=I^{\widehat{\nu }\tau }\text{ resp. }\tau I=I^{\tau ^1},$$
see 2.2. In our situation, these functors are up to isomorphism determined by their effect on objects. We conclude that we have an isomorphism of bimodules
(7.3)
$$\stackrel{~}{\mathrm{\Lambda }}^{\widehat{\nu }}D\stackrel{~}{\mathrm{\Lambda }},$$
see 7.2 (b).
In order to state our next result we introduce $`\underset{¯}{\mathrm{\Gamma }}_Q`$, the full subcategory of the Auslander category $`\mathrm{\Gamma }_Q`$ which contains all objects except those of the form $`\widehat{\nu }(0,q)`$ for $`qQ_0`$. We call $`\underset{¯}{\mathrm{\Gamma }}_Q`$ the stable Auslander category of $`\mathrm{k}Q`$.
###### 7.4 Proposition.
The category $`\stackrel{~}{\mathrm{\Lambda }}`$ is isomorphic to the repetitive category $`\underset{¯}{\overset{^}{\mathrm{\Gamma }}}_Q`$ of the stable Auslander category of $`\mathrm{k}Q`$.
Proof: Recall that by definition the objects of $`\underset{¯}{\overset{^}{\mathrm{\Gamma }}}_Q`$ are of the form $`(z,x)`$ with $`z`$ and $`x\mathrm{Obj}(\underset{¯}{\mathrm{\Gamma }}_Q)`$ and we have
$$\underset{¯}{\overset{^}{\mathrm{\Gamma }}}_Q((z,x),(z^{},x^{}))=\{\begin{array}{cc}\underset{¯}{\mathrm{\Gamma }}_Q(x,x^{})\hfill & \text{ if }z=z^{},\hfill \\ D\underset{¯}{\mathrm{\Gamma }}_Q(x,x^{})\hfill & \text{ if }z=z^{}1,\hfill \\ 0\hfill & \text{ else.}\hfill \end{array}$$
Here we define the dual $`\underset{¯}{\mathrm{\Gamma }}_Q\text{-}\underset{¯}{\mathrm{\Gamma }}_Q`$-bimodule $`D\underset{¯}{\mathrm{\Gamma }}_Q`$ by $`D\underset{¯}{\mathrm{\Gamma }}_Q(x,y):=\mathrm{Hom}_\mathrm{k}(\underset{¯}{\mathrm{\Gamma }}_Q(y,x),k)`$, compare 7.2 (a). Now, by the bimodule isomorphism (7.3) we see that the assignation
$$(z,(i,q))\tau ^i\widehat{\nu }^z(0,q)$$
induces an isomorphism $`\underset{¯}{\overset{^}{\mathrm{\Gamma }}}_Q\stackrel{~}{\mathrm{\Lambda }}`$. $`\mathrm{}`$
### 7.5. Conclusions.
(a) In our situation we note that in $`𝒟^b(\mathrm{k}Q)𝒟^b(\mathrm{k}Q^{\text{op}})`$ we have an isomorphism of functors
(7.4)
$$\mathrm{\Sigma }^2\tau ^{h(Q)}$$
where $`h(Q)`$ is the Coxeter number of $`|Q|`$. In fact, both functors coincide on objects as one easily verifies in $`\stackrel{~}{\mathrm{\Lambda }}`$. In our quiver situation this is sufficient. From (7.2) we obtain immediately the remarkable functorial isomorphism
(7.5)
$$\tau _{\stackrel{~}{\mathrm{\Lambda }}}^6M^{(h(Q)6)}$$
in $`\stackrel{~}{\mathrm{\Lambda }}\text{-}\underset{¯}{\mathrm{mod}}𝒟^b(\mathrm{k}Q)\text{-}\underset{¯}{\mathrm{mod}}`$.
(b) The action of the infinite cyclic group $`\tau `$ on $`\stackrel{~}{\mathrm{\Lambda }}`$ provides us with a Galois covering
$$F:\stackrel{~}{\mathrm{\Lambda }}\mathrm{\Lambda },$$
see for example . We conclude that $`\mathrm{\Lambda }\mathrm{inj}`$ is the orbit category of $`\stackrel{~}{\mathrm{\Lambda }}\mathrm{inj}`$ modulo the induced action of $`\tau `$. Now, $`\stackrel{~}{\mathrm{\Lambda }}\mathrm{inj}𝒟^b(\mathrm{k}Q)`$ is a triangulated category, and the hypothesis of are obviously fulfilled. Thus $`\mathrm{\Lambda }\mathrm{inj}`$ is also a triangulated category with Auslander-Reiten triangles and the corresponding translation $`\overline{\tau }`$ is the identity. Note moreover that the induced suspension is isomorphic to the corresponding Nakayama automorphism, i.e. $`\overline{\mathrm{\Sigma }}\overline{\nu }`$. By applying (7.2) to $`𝒟=(\mathrm{\Lambda }\mathrm{inj})^{\text{op}}`$ we conclude that
$$\tau _\mathrm{\Lambda }^3MM^{\overline{\nu }}\text{and}\tau _\mathrm{\Lambda }^6MM$$
holds functorially in $`\mathrm{\Lambda }\text{-}\underset{¯}{\mathrm{mod}}`$. In case $`\nu `$ is just a translation, see 2.3, we even have $`\tau _\mathrm{\Lambda }^3MM`$. This is our interpretation of the proof for the $`6`$-periodicity of $`\tau _\mathrm{\Lambda }`$ in . Even more directly by (7.2) we conclude that
$$\tau _\mathrm{\Lambda }\mathrm{\Omega }_\mathrm{\Lambda }MM$$
functorially. This means that the triangulated category $`\mathrm{\Lambda }\text{-}\underset{¯}{\mathrm{mod}}`$ is of Calabi-Yau dimension 2.
(c) Since $`\stackrel{~}{\mathrm{\Lambda }}\underset{¯}{\overset{^}{\mathrm{\Gamma }}}_Q`$ we have by Happel’s Theorem \[21, II.4\] $`\stackrel{~}{\mathrm{\Lambda }}\text{-}\underset{¯}{\mathrm{mod}}𝒟^b(\underset{¯}{\mathrm{\Gamma }}_Q\mathrm{mod})`$ as triangulated categories. If we consider the push-down functor
$$F_\lambda :\stackrel{~}{\mathrm{\Lambda }}\text{-}\underset{¯}{\mathrm{mod}}\mathrm{\Lambda }\text{-}\underset{¯}{\mathrm{mod}}$$
associated to the Galois covering $`F:\stackrel{~}{\mathrm{\Lambda }}\mathrm{\Lambda }`$ the subcategory of $`\mathrm{\Lambda }`$-modules of the first kind (i.e. the subcategory of objects which are isomorphic to a push-down) is equivalent to the orbit category $`𝒟^b(\underset{¯}{\mathrm{\Gamma }}_Q\mathrm{mod})/\tau \mathrm{\Sigma }^1`$ via the above identifications and (7.2), see . Here, $`\mathrm{\Sigma }`$ resp. $`\tau `$ are the suspension resp. the Auslander-Reiten translation in $`𝒟^b(\underset{¯}{\mathrm{\Gamma }}_Q\mathrm{mod})`$.
Now, for a Dynkin quiver $`Q`$ with Coxeter number $`h(Q)6`$ the algebras $`\underset{¯}{\mathrm{\Gamma }}_Q`$ are (quasi-) tilted. Thus, in these cases we find $`\mathrm{\Lambda }\text{-}\underset{¯}{\mathrm{mod}}𝒟^b(\underset{¯}{\mathrm{\Gamma }}_Q)/\mathrm{\Sigma }\tau ^1`$ is a cluster category in the sense of .
### 7.6. Remark
Let $`H`$ be a (basic, connected) finite-dimensional hereditary $`\mathrm{k}`$-algebra of finite representation type. Then $`H`$ is a species of type $`𝖠`$$`𝖦`$ in the sense of Dlab and Ringel, see . In this case we may also study the stable Auslander category $`\underset{¯}{\mathrm{\Gamma }}_H`$. The same argument as in 7.3 and 7.4 shows that $`\underset{¯}{\overset{^}{\mathrm{\Gamma }}}_H\mathrm{inj}𝒟^b(H^{^{\text{op}}})`$. The Auslander-Reiten translate in $`𝒟^b(H^{^{\text{op}}})`$ induces an automorphism $`\tau `$ of $`\underset{¯}{\overset{^}{\mathrm{\Gamma }}}_H`$. Thus we may consider the Galois covering $`\underset{¯}{\overset{^}{\mathrm{\Gamma }}}_H\underset{¯}{\overset{^}{\mathrm{\Gamma }}}_H/\tau `$. So we are tempted to consider $`\underset{¯}{\overset{^}{\mathrm{\Gamma }}}_H/\tau `$ as the preprojective algebra of $`H`$. However, $`\tau `$ is now in general not determined by its effect on objects since $`\mathrm{Out}_\mathrm{k}(H)=\mathrm{Aut}_k(H)/\mathrm{Inn}(H)`$ is possibly non-trivial. Thus we are (for the moment) unable to compare $`\underset{¯}{\overset{^}{\mathrm{\Gamma }}}_H/\tau `$ with the possible choices for the preprojective algebra of $`H`$ in the sense of Dlab and Ringel .
Anyway, if we denote by $`|H|`$ the (unoriented) diagram of $`H`$ then we find the list which we present in Figure 1 of interest due to its similarity with the cluster types of $`[N]`$. In the case of Coxeter number $`c(|H|)=6`$ we display the diagram of a canonical tubular algebra following Lenzing and the corresponding extended affine root system in the sense of Saito . In the case of $`𝖦_2`$ there are two different diagrams of canonical algebras which produce derived equivalent algebras (for adequate choices of bimodules), the corresponding two root systems are isomorphic up to the marking.
Note moreover, that in this case our previous calculations (7.5) predict that in the stable module category of the repetitive algebra $`\underset{¯}{\overset{^}{\mathrm{\Gamma }}}_H\text{-}\underset{¯}{\mathrm{mod}}`$ the Auslander-Reiten translate should be $`6`$-periodic. This is in fact the case for all (tubular) algebras which are derived equivalent to a canonical algebra with a diagram from our list.
### Acknowledgement
We like to thank Henning Krause for drawing our attention to the result on functor categories over triangulated categories. We also thank Markus Reineke for pointing out the reference for the material in 1.3. |
warning/0506/hep-ph0506139.html | ar5iv | text | # Monte-Carlo simulation of lepton pair production in 𝑝̄𝑝→𝑙⁺𝑙⁻+𝑋 events at 𝐸_{𝑏𝑒𝑎𝑚} = 14 GeV
## 1 Introduction
The measurements of lepton pair production in hadron-hadron interactions (in the following, MMTDY process, see and ) have already demonstrated their great potential for studying the properties of elementary particles. As for illustration, it is enough to mention the facts of discovery of charmed J(J/$`\mathrm{\Psi }`$) - meson as well as of beauty $`\mathrm{{\rm Y}}`$ \- meson which were done first in hadron-hadron collisions and confirmed later in $`e^+e^{}`$ experiments. Dilepton events may serve as a powerful tool to get out the information about the parton distribution functions (PDFs) in hadrons as it was already shown in a number of high energy experiments and theoretical papers, devoted to the data analysis in the framework of QCD . The plans to study this process are included into the LoI , TPR and of PANDA experiment at HESR <sup>1</sup><sup>1</sup>1Analogous arguments in a favor of studying of this process may be found also in a number of recent proposals for experimental program at HESR (see , ). which may provide an interesting information about quark dynamics inside the nucleon.
This intermediate energy experiment (in the following we shall consider the case of antiproton beam energy $`E_{beam}=14`$ GeV which corresponds to the center-of-mass energy of the $`p\overline{p}`$ system $`E_{cm}=5.3`$ GeV) may play an important role because it allows one to study the energy range where the perturbative methods of QCD (pQCD) come into interplay with a rich physics of bound states and resonances. The physics of hadron resonances formation and decay is strongly connected with the confinement problem, i.e. with the parton dynamics at large distances. A detailed and high-precision experimental study at PANDA may allow one to discriminate between a large variety of existing nonperturbative approaches and models that already exist or are under development now.
To reach the goals declared in the , , in connection with lepton pair production process , , one needs to know the possible energy, momentum and angle distributions of the produced individual leptons as well as the analogous distributions for the lepton pair as a whole. So, a detailed Monte-Carlo simulation of $`\overline{p}pl^+l^{}+X`$ ($`l=\mu ,e`$) process (see Fig.1) is needed. It is also clear that such a sort of simulation is also needed for a proper design of the muon system and electromagnetic calorimeter (see , ).
For this aim we utilized here, as for the first step, the well known event generator PYTHIA , which is based on the ideas of the quark parton model and is well tested and widely used for the simulation of hadron-hadron interactions. PYTHIA simulation is based on the use of the amplitudes of the relativistic quantum field theory. This allows a proper account of the relativistic kinematics during simulation of different physical variables distributions specific for $`\mu ^+\mu ^{}`$\- or $`e^+e^{}`$\- pair produced in event. In our case we use the perturbative QCD/QED parton level amplitude of the lepton pair production process $`q\overline{q}\gamma ^{}l^+l^{}`$ with the continuous spectrum of the invariant mass of lepton pair and the amplitude of J/$`\mathrm{\Psi }`$ \- resonance production process $`p+\overline{p}J/\mathrm{\Psi }+Xl^+l^{}+X`$, where J/$`\mathrm{\Psi }`$ decays through the leptonic channel ($`J/\mathrm{\Psi }l^+l^{}`$, $`l=\mu ,e`$ ), which are implemented into the PYTHIA package.
Let us underline that the results obtained here on the basis of PYTHIA simulation in some sense may allow to fix the predictions for lepton kinematical distributions which may be obtained in the framework of perturbative theory approach. Therefore, they may be useful at the analysis stage for defining the boundary between the predictions of perturbative and nonperturbative theoretical approaches.
In Section 2 we present the kinematical distributions for individual leptons. The set of plots with energy, transverse momentum and angle distributions are given together with some plots which show different kinds of Energy-Energy, Energy-Angle and Angle-Angle correlations between the physical variables of the leptons produced via the leading order quark level subprocess $`q\overline{q}l^+l^{}`$. The estimations of events loss due to a possible different choice of geometrical parameters of muon system and electromagnetic calorimeter are given.
Section 3 is devoted to another channel of lepton pair production. The process of $`J/\mathrm{\Psi }`$ resonance production (with its decay into a lepton pair) $`p+\overline{p}J/\mathrm{\Psi }+Xl^+l^{}+X`$ ($`l=\mu ,e`$) is simulated by use of PYTHIA6.4 which includes a sizeable set of parton level subprocesses that can give a contribution to this process. This process was chosen to be a benchmark process for PANDA experiment (see TPR ). Thus, the obtained here kinematical distributions of the final state leptons shall have a practical application. In the following we shall call this process as resonance production to distinguish it from the process considered in Section 2, where the invariant mass of two leptons produced in quark level subprocess $`q\overline{q}l^+l^{}`$ has a continuum spectrum.
Section 4 includes the distributions of the invariant mass and some other physical variables which are characteristic for the signal lepton pair as a whole system. The most interesting among them is the total transverse momentum of a lepton pair which is connected with the intrinsic transverse velocity of a quark inside the proton.
In Section 5 we estimate the size of the kinematical region in $`xQ^2`$ plane which can be available for measuring the quark distributions in PANDA experiment.
The problems connected with the background from the fake leptons, which may appear together with the signal lepton pair in one and the same event due to meson decays, as well as with the background from other than $`q\overline{q}l^+l^{}`$ subprocesses, are discussed in Sections 6 and 7, correspondingly. Also we present a set of cuts which allows to separate the background events from the signal ones. The efficiencies of the proposed cuts are also given.
In Section 8 we outline some important physical measurements which can be done by studying the lepton pair production at the energies available for PANDA experiment.
It is worth mentioning that the results presented here can also be useful for the future physical analysis of hadron decays at PANDA because the contribution from $`q\overline{q}l^+l^{}`$ events may be one of the main background sources in this kind of a study.
## 2 Distributions of leptons produced in $`p\overline{p}`$ collisions
We use PYTHIA6 to generate two samples (separately for muons and electrons) of 100 000 ”$`\overline{p}pl^+l^{}+X`$” events which include the $`22`$ quark level $`q\overline{q}\gamma ^{}l^+l^{}`$ subprocess. In the following, these events will be called as ”signal events”, and the muons/electrons produced in this subprocess will be called as ”signal” leptons. The fake leptons which are produced in hadron (mainly mesons) decays in the same signal event will be called ”decay” leptons. The simulation is done starting from the assumption of the ideal muon system and electromagnetic calorimeter covering $`360^o`$. No any cuts are used in the following Sections. They will be discussed in the subsection 6.5 and the Section 8.
We consider the case where both initial-state radiation (ISR) and final-state radiation (FSR) are switched on simultaneously by choosing the corresponding values of PYTHIA parameters. Also we have used the CTEQ3L parametrization of parton distributions and the default value of the parameter, which allows one to take into account the primordial $`k_T`$-effect (the effect of quark Fermi motion inside the hadron).
First we consider the distributions of some physical variables which describe the kinematics of individual leptons belonging to the $`l^+l^{}`$ pair. Simulation shows that there is no difference between $`e^+e^{}`$ and $`\mu ^+\mu ^{}`$ distributions. In all the following figures the vertical axis shows the number of events (per bin) that may be expected per year ($`10^7sec`$) for the luminosity $`L=210^{32}cm^2s^1`$ (=2 $`10^5mb^1s^1`$). The total number of expected events per year are shown as ”Integral” values in the figures. Let us underline that this number can be treated only as an estimation because it strongly depends on the models used in PYTHIA which is basically designed for much higher energies.
The distributions of the number of generated signal events ($`N_{ev}`$) versus the the energy $`E^{l^{(+/)}}`$ of the signal leptons, as well as versus the modulus of the transverse momentum $`P_T^{l^{(+/)}}`$ and of the polar (zenith) angles $`\theta ^{l^{(+/)}}`$, measured from the z-axis directed along the beam line, are given (top to bottom) in Fig.2. The left column of Fig.2 is for $`l^{}`$ distributions and the right one is for $`l^+`$. One can see from the top row of Fig.2 (plots a and b) that the most part of leptons energy is contained in the interval $`0<E^l<10`$ GeV. Its spectrum has a mean value $`<E^l>=2.6`$ GeV and a peak at $`E_{peak}^l0.5`$ GeV. The $`P_T^l`$ spectrum (see the middle row plots c and d of Fig.2) has an analogous peak at $`P_{T_{peak}}^l0.4`$ GeV. Behind their peaks both $`E^l`$ and $`P_T^l`$ spectra fall rather steeply. The main part of $`P_T^l`$ spectrum is confined within a rather narrow interval $`0<P_T^l<2`$ GeV.
The number of events spectrum versus the polar angle $`\theta ^l`$ (see bottom row plots e and f of Fig.2) has a peak around $`\theta ^l10^o`$ and the mean value $`<\theta ^l>=27.3^o`$. One sees that while the most of signal leptons fly in the forward direction ($`\theta ^l<90^o`$) there is still a small number of them which fly into the back hemisphere ($`\theta ^l>90^o`$).
Fig.3 includes the set of plots done separately for the signal leptons having the largest energy $`E_{fast}^l`$ (right column) in the lepton pair and for the leptons having a smaller energy $`E_{slow}^l`$ in the pair (left column). We shall call them, correspondingly, as ”fast” and ”slow” leptons. One can see that the energy spectrum of fast signal leptons (plot b) increases rather fast from the point $`E_{fast}^l0.5`$ GeV (more than $`90\%`$ of fast leptons have $`E_{fast}^l>1`$ GeV) up to the peak position at the point $`E_{fast}^l2.5`$ GeV ($`<E_{fast}^l>`$ = 3.85 GeV). Then it smoothly vanishes at $`E_{fast}^l=10`$ GeV.
In contrast to this picture, the analogous spectrum of the less energetic signal leptons (plot a) starts sharply from zero and reaches a peak at $`E_{slow}^l0.4`$ GeV (where the spectrum of the fast leptons only starts). Then it goes down and practically vanishes at the point $`E_{slow}^l5`$ GeV. One may see that the energy spectrum of slow leptons plot a in a pair is more than by half shorter than that one of fast leptons plot b and the mean value of slow leptons energy $`<E_{slow}^l>=1.36`$ GeV is about 3 times less than the mean energy of fast leptons $`<E_{fast}^l>=3.85`$ GeV.
The difference between the $`PT^l`$ spectra of fast and slow leptons is not so large (see, correspondingly, plots d and c in the middle row of Fig.3). They differ only by about 400 MeV shift to the left of the peak position of the slow leptons spectrum and about 340 MeV analogous shift of their mean transverse momentum value. Both of these spectra demonstrate that the main part of slow and fast leptons has $`PT^l>0.2`$ GeV.
The bottom row of the plots demonstrates that the polar (zenith) angle spectrum of less energetic leptons $`\theta _{slow}^l`$ (Fig.3 e) is shifted to the higher values as compared to the spectrum of fast leptons $`\theta _{fast}^l`$ (Fig.3 f). Its mean value $`<\theta _{slow}^l>=38.2^o`$ is more than twice as large as the analogous mean value of fast leptons: $`<\theta _{fast}^l>=16.5^o`$. Thus, we can conclude that almost all fast leptons fly in the forward direction ($`\theta _{fast}^l<80^o`$) and their spectrum practically finishes at $`\theta _{fast}^l60^o`$ (plot f) while about $`17\%`$ of slow leptons (plot e) have $`\theta _{slow}^l>60^o`$. It is worth noting that about $`5\%`$ of slow leptons may scatter into the back hemisphere.
Fig.4 contains two 3-Dimensional ”Angle-Energy” correlation plots for slow $`\theta _{slow}^l/E_{slow}^l`$ (plot $`𝐚`$) and fast $`\theta _{fast}^l/E_{fast}^l`$ (plot $`𝐛`$) leptons in signal pairs. Their vertical axes show the distribution of number of signal events ($`N_{ev}`$) multiplied by $`10^3`$. 2-dimensional plots $`𝐜`$ and $`𝐝`$ are obtained by projection of plots $`𝐚`$ and $`𝐛`$ onto the $`\theta E`$ planes. This allows to show the boundary contours of regions with different density of number of events. The right-hand color vertical strip in each plot $`𝐜`$ and $`𝐝`$ shows the correspondence between the contour color and the density of events. This color strip plays the role of the vertical z axis with the number of events $`N_{ev}`$ shown in plots a and b. Plots c and d of Fig.4 show the kinematical regions in $`\theta E`$ plane which are covered by slow and fast leptons, respectively. Plot c of Fig.4 demonstrates that the value of the polar angle of slow leptons $`\theta _{slow}^l`$ drops much steeply with the growth of their energy as compared to the behavior of the polar angle of fast lepton $`\theta _{fast}^l`$ shown at the plot $`𝐝`$.
After discussion of individual lepton distributions let us turn to the distributions that characterize the produced pair of leptons as a whole system. Fig.5 shows the Energy-Energy $`E_{slow}^l/E_{fast}^l`$ (plots a and c) and Angle-Angle $`\theta _{slow}^l/\theta _{fast}^l`$ (plots b and d) correlations. Plots c and d are the projections of 3D plots a and b on $`E_{slow}^lE_{fast}^l`$ and $`\theta ^{slow}\theta ^{fast}`$ planes, correspondingly.
Data taking of searched signal events is strongly influenced by the cuts on lepton energies which can be imposed from below to suppress the electronic noise and some other background which can be provided by detector effects. The analysis of the plots a and c of Fig.5 is summarized in the Table 1. It demonstrates (in $`\%`$) the loss of signal events due to application of the kinematical cut $`E_{fast}^l,E_{slow}^lE_{cut}`$ which sets the lower limit $`E_{cut}`$ on the value of lepton energy.
The efficiency of collection of signal events which contain leptonic $`l^+l^{}`$ pairs also depends on the angle coverage by the muon system and the Electromagnetic Calorimeter (ECAL). Plots $`𝐛`$ and $`𝐝`$ of Fig.5 show the Angle-Angle $`\theta _{slow}^l/\theta _{fast}^l`$ lepton correlation. The results of its analysis are given in the Table 2 which demonstrates which part (in $`\%`$) of signal events would be lost due to imposing the upper limit $`\theta _{cut}`$ (i.e. $`\theta _{slow}^l,\theta _{fast}^l\theta _{cut}`$ cut) on the size of muon system or ECAL.
The last line of Table 2 shows that even in the case when the muon system or the ECAL would cover the angle region $`\theta ^l90^o`$, about $`5\%`$ of the events containing the $`l^+l^{}`$ signal pairs would be lost. Nevertheless, such geometrical boundary allows to keep about 95$`\%`$ of signal events with electron or muon pairs. Therefore, we consider this choise of polar angle upper limit as a preferable one for the study of MMTDY process of lepton pair production (with the continuous mass spectrum of this pair).
## 3 Leptons from J/$`\mathrm{\Psi }`$ decay
The process of J/$`\mathrm{\Psi }`$ resonance production with its further decay into a lepton ($`l=\mu ,e`$) pair $`p+\overline{p}J/\mathrm{\Psi }+Xl^+l^{}+X`$ (see one of the possible reaction on Fig.6) was considered in the PANDA TDR as one of the benchmark processes. Therefore, modeling of the kinematical (energy, transverse momentum and angle) distributions of final state leptons is of practical interest.
Like in previous Section, we use the same event generator PYTHIA6.4 which includes the following set of subprocesses with J/$`\mathrm{\Psi }`$ production:
1) $`q_i\overline{q}_i\gamma ^{}c\overline{c}J/\mathrm{\Psi }l^+l^{}+X`$ 86) $`ggJ/\mathrm{\Psi }+gl^+l^{}+X`$
106) $`ggJ/\mathrm{\Psi }+\gamma l^+l^{}+X`$ 421) $`ggc\overline{c}[^3S_1^{(1)}]gl^+l^{}+X`$
422) $`ggc\overline{c}[^3S_1^{(8)}]gl^+l^{}+X`$ 423) $`ggc\overline{c}[^3S_0^{(8)}]gl^+l^{}+X`$
424) $`ggc\overline{c}[^3P_J^{(8)}]gl^+l^{}+X`$ 425) $`gqc\overline{c}[^3S_1^{(8)}]ql^+l^{}+X`$
426) $`gqc\overline{c}[^3P_J^{(8)}]ql^+l^{}+X`$ 427) $`ggc\overline{c}[^3S_1^{(1)}]ql^+l^{}+X`$
428) $`q\overline{q}c\overline{c}[^3S_1^{(8)}]gl^+l^{}+X`$ 429) $`q\overline{q}c\overline{c}[^1S_0^{(8)}]gl^+l^{}+X`$
430) $`q\overline{q}c\overline{c}[^3P_J^{(8)}]gl^+l^{}+X`$ 431) $`ggc\overline{c}[^3P_0^{(1)}]gl^+l^{}+X`$
432) $`ggc\overline{c}[^3P_1^{(1)}]gl^+l^{}+X`$ 433) $`ggc\overline{c}[^3P_2^{(1)}]gl^+l^{}+X`$
434) $`gqc\overline{c}[^3P_0^{(1)}]ql^+l^{}+X`$ 435) $`gqc\overline{c}[^3P_1^{(1)}]ql^+l^{}+X`$
436) $`gqc\overline{c}[^3P_2^{(1)}]ql^+l^{}+X`$ 437) $`qqc\overline{c}[^3P_0^{(1)}]gl^+l^{}+X`$
438) $`q\overline{q}c\overline{c}[^3P_1^{(1)}]gl^+l^{}+X`$ 439) $`q\overline{q}c\overline{c}[^3P_2^{(1)}]gl^+l^{}+X`$
The main contribution to the total cross section of J/$`\mathrm{\Psi }`$ production <sup>2</sup><sup>2</sup>2 Let us note that different theoretical models predict different values of J/$`\mathrm{\Psi }`$ production cross sections. comes from the following three subprocesses:
1) $`q_i\overline{q}_i\gamma ^{}c\overline{c}J/\mathrm{\Psi }l^+l^{}+X`$
428) $`q\overline{q}c\overline{c}[^3S_1^{(8)}]gl^+l^{}+X`$
430) $`q\overline{q}c\overline{c}[^3P_J^{(8)}]gl^+l^{}+X`$
Distributions of the final state leptons, produced in J$`/\mathrm{\Psi }`$ decay, are shown in Fig.7. They are obtained without a use of any cuts and cover the same ranges as the leptons produced in the continuum case (see Fig.2). From comparison of these two Figures one can see that the energy and momentum distributions presented in Fig.7 look very different to those of Fig.2. Thus, plot a of Fig.7 shows a rather flat distribution of the number of events versus the energy of leptons from J$`/\mathrm{\Psi }`$ decay, while the analogous plots a and b of Fig.2 demonstrate that the most part of leptons produced in the continuum case have small energies $`E^l<1`$ GeV. Analogously, as seen from the plot b of Fig.7, the peak of $`P_T^l`$ distribution in a case of J$`/\mathrm{\Psi }`$ production appears at the point $`P_T^l1.4`$ which is more than three times higher than the maximum value of $`P_T^l`$ in the plot b of Fig.2 ($`0.4`$ GeV). As one can see from the plot d of Fig.7, the maximum of cos $`(l^+,l^{})`$ is at $``$ 0.75. This value corresponds to the opening angle between the leptons of $`42^o`$.
Fig.8 is the analog of Fig.3. It presents the distributions of the same leptons produced in J$`/\mathrm{\Psi }`$ decay, but separately for ”slow” and ”fast” ones. One can see that the plots a and b of Fig.8 look much more different from the plots a and b of Fig.3. The main difference is that the energy spectra of slow ($`E_{slow}`$) and fast ($`E_{fast}`$) leptons, shown in the plots a and b of Fig.8, cover very different energy intervals which have very small area of their overlapping in the region of $`E^l45`$ GeV. It is also seen from the plots c and d of Fig.8 that transverse momenta of both slow and fast leptons have a rather close mean values about $`P_T^l1.4`$. Plots e and f of Fig.8 demonstrate a good separation of polar angle distributions of slow and fast leptons. It is seen that the most part of the spectrum of slow leptons lay behind the 20<sup>o</sup>, while the spectrum of fast leptons polar angles ranges in the interval from 0<sup>o</sup> to 25<sup>o</sup> and has the mean value $`\theta _{fast}^l=`$12.12 GeV.
Fig.9 includes the plots which show Angle-Energy correlations of the leptons produced in $`J/\mathrm{\Psi }`$ decay. Plots a and c are for slow leptons, plots b and d are for fast leptons. They are also very different from those shown in Fig.4 for the case of continuum lepton pair production via the process $`p\overline{p}+p\gamma ^{}l^+l^{}+X`$. It is seen that fast and slow leptons produced in $`J/\mathrm{\Psi }`$ decay cover very different and well separated regions (see plots c and d of Fig.9) while in the continuum case (see plots c and d of Fig.4) the distributions of fast and slow leptons have a wide area of their overlapping. These regions are rather narrow and have a strip form, that looks different from the case of continuum lepton pair production (see plots c and d of Fig.4).
It is of interest to compare the Energy-Energy and Angle-Angle correlation plots which are presented in Fig.10 for the case of lepton pair production from $`J/\mathrm{\Psi }`$ resonance decay and those shown in Fig.5 for leptons production in continuum case. It is seen from the plot c of Fig.10 that the energy region covered in $`E_{slow}E_{fast}`$-plane in the case of $`J/\mathrm{\Psi }`$-resonance production process fits well into the right corner of the region covered covered in the $`E_{slow}E_{fast}`$-plane shown in the plot c of Fig.5 for the case of continuum mass MMTDY process of lepton-antilepton pair production.
Moreovere, the analogous comparison of the plot d of Fig.10 with the plot d of Fig.5 allows to make an important observation that the area covered in the $`\theta _{slow}\theta _{fast}`$ plane fits well into the analogous aria of Angle-Angle correlation plot in Fig.5. From here one can conclude that the choice of polar angle boundary for muon system geometry, which was proposed in Section 2 basing on the analysis of MMTDY process, would be quite suitable for the study of both benchmark $`J/\psi `$ and MMTDY channels of lepton-antilepton pair production.
## 4 Distributions of the invariant mass, energy <br>and transverse momenta of lepton pairs
We consider here a set of physical variables which characterize a produced lepton pair as a whole system. These variables are constructed from the components of the total 4-momentum of initial state quark-antiquark system $`P_\alpha ^{q\overline{q}}=P_\alpha ^q+P_\alpha ^{\overline{q}}`$, ($`\alpha =0,1,2,3`$) and its analog $`P_\alpha ^{l^+l^{}}=P_\alpha ^{l^+}+P_\alpha ^l^{}`$ for a lepton pair ($`P^{l^\pm }`$ is the 4-momentum of lepton l) <sup>3</sup><sup>3</sup>3 $`P^l=(P_o^l,𝐏^l)`$, where $`𝐏=(P_x,P_y,P_z)`$ and $`P_o=\sqrt{M^2+𝐏^\mathrm{𝟐}}`$.. Figs.11a, 11b show, correspondingly, the distributions of the invariant masses of initial-state quark-antiquark pair
$$M_{inv}^{q\overline{q}}=\sqrt{(}P^{q\overline{q}})^2,$$
(1)
and the invariant mass of the final-state lepton-antilepton pair
$$M_{inv}^{l^+l^{}}=\sqrt{(}P^{l^+l^{}})^2=Q,Q^2=q^2=(P_\alpha ^{l^+}+P_\alpha ^l^{})^2,$$
(2)
been produced in the signal process $`\overline{p}pl^+l^{}+X`$ which goes through the quark level subprocess $`q\overline{q}\gamma ^{}l^+l^{}`$. Both invariant mass distributions look rather similar. They are rather short and drop steeply with the growth of invariant mass. The distribution of the invariant mass $`M_{inv}^{q\overline{q}}`$ of the initial-state $`q\overline{q}`$-system sharply starts at the point $`M_{inv}^{q\overline{q}}=1`$ GeV which is the left boundary point due to the internal PYTHIA restriction on the lowest value of the invariant mass of the initial state two-body system of any fundamental quark-parton $`22`$ subprocess. The spectrum finishes at $`M_{inv}^{q\overline{q}}=M_{inv}^{l^+l^{}}2.5`$ GeV.
Different to the spectrum of the invariant mass of the initial-state $`q\overline{q}`$-system, the invariant mass $`M_{inv}^{l^+l^{}}`$ of the final-state $`l^+l^{}`$ system has a very small tail at smaller than 1 GeV values $`M_{inv}^{l^+l^{}}<1`$ GeV (see Fig.11a) <sup>4</sup><sup>4</sup>4This left tail of $`M_{inv}^{l^+l^{}}`$ may appear due to the final state radiation (FSR) of photons by the produced leptons..
The spectrum of the total lepton pair energy $`E^{l^+l^{}}=E^{l^+}+E^l^{}`$ is shown in Fig.12a. It is seen that the lepton pair total energy distribution is by about 2 GeV longer than the spectrum of the fast leptons energy $`E_{fast}^l`$ (see Fig.3b).
The mean value of the total energy of the lepton pair, as seen from Fig.12$`𝐚`$, is about 5 GeV. From the same plot of Fig.12 it is clearly seen that in more than a half of events the lepton pair energy lays in the interval $`4E^{l^+l^{}}12`$ GeV. So, one can conclude that the produced lepton pairs are rather energetic and they can carry away quite a noticeable part of the total energy of the colliding $`\overline{p}p`$-system.
The distribution of the longitudinal component $`P_z^{l^+l^{}}=P_z^{l^+}+P_z^l^{}`$ of the total 4-momentum of lepton pair system (not shown here) has a shape which is very similar to the energy $`E^{l^+l^{}}`$ spectrum. The explanation of this fact follows from the shape of the distribution of the modulus of lepton pairs total transverse momentum
$$P_T^{l^+l^{}}=|\stackrel{}{P}_T^{l^+l^{}}|=|\stackrel{}{P}_T^{l^+}+\stackrel{}{P}_T^l^{}|,$$
(3)
which is presented as the “$`P_T^{l^+l^{}}`$ distribution” in Fig.12b. One may see that the distribution of the transverse component $`P_T^{l^+l^{}}`$ is much more narrow than that one of the energy $`E^{l^+l^{}}`$. It covers the region $`0<P_T^{l^+l^{}}<2`$ GeV, like it was in a case of a single lepton distribution (see plots c and d of Fig.2). $`P_T^{l^+l^{}}`$ has a peak position at about 1 GeV, that is twice as large as the analogous peak position of a single lepton $`P_T^l`$ shown in the same plots c and d of Fig.2. Thus, we see that the main contribution to the value of lepton pair energy comes from the longitudinal component $`P_z^{l^+l^{}}`$.
It is worth noting that according to Fig.1 the square of the invariant mass $`(M_{inv}^{l^+l^{}})^2=Q^2=q^2=(P^{l^+}+P^l^{})^2`$ has the meaning of the square of the momentum transferred from the quark-antiquark pair to the lepton pair. Therefore, it plays the same role as the $`Q^2`$ in the processes of deep-inelastic scattering (DIS) of lepton over the proton. From this point of view the diagram shown in Fig.1 looks like the cross chanel analog of the diagram of DIS scattering of a lepton over the proton with the inclusive production of a proton in the final state. The essential difference is that in DIS case the momentum transferred is defined by the relation $`q_{dis}=P_{in}^lP_{out}^l`$ ($`P_{in}^l`$ and $`P_{out}^l`$ are, respectively, the 4-momenta of incoming and outgoing leptons) <sup>5</sup><sup>5</sup>5 For DIS the value of $`Q^2`$ is defined as $`Q^2=q_{dis}^2=(P_{in}^lP_{out}^l)^2`$., which differs from the definition of $`Q^2`$ given above for the case of MMTDY process. Therefore, its square has the negative value $`q_{dis}^2`$ 0, while in MMTDY process its square is positive: $`q^2(P^{l^+}+P^l^{})^20`$.
To conclude this Section let us mention that the measurement of the invariant mass of lepton pair would allow to separate the background to $`J/\psi `$ production events up to a good accuracy.
## 5 Estimation of the size of $`xQ^2`$ region available for measurement of proton structure function
The distributions of Bjorken x-variables are shown in Fig.13 for up- (plot a) and down- (plot b) quarks <sup>6</sup><sup>6</sup>6 The distributions of antiquarks look similar to quark distributions for $`\overline{p}p`$ collisions. They represent the corresponding quark components of the proton structure function which was used in the present simulation with CTEQ3L PDF.
The most ineresting for us is the information about the size of the x-region which will be available at PANDA energies. We see from both plots of Fig.13 that x-varaiable spans the interval $`0.05<x<0.7`$. Recall that in $`\overline{p}p`$ collisions the transverse momentum $`P_T`$ plays the role of the transferred momentum q.
So, combining the results from Figs.11 and 13 we can conclude that, according to the results of simulation with PYTHIA, we can hope to get the information about valence quark distributions in the kinematical region defined by the following boundaries: $`0.05x0.7`$ and $`Q^26.2`$ GeV ($`Q^2q^2=(P_T^{l^+l^{}})^2`$). The important point that should be stressed here is that different to lepton-hadron scattering, which provides the information about the structure functions in the region of negative, i.e. ”space-like” values of the square of transversed momentum $`q_{dis}^2=(P_{in}^lP_{out}^l)^2`$, the ”annihilation” process $`\overline{p}pl^+l^{}+X`$ allows to get the information about the structure functions in a region of positive, i.e. ”time-like” values of $`q^20`$. Such a measurement of quark distributions will be a good supplement to the planned measurement of proton elastic formfactor in the region of ”time-like” values of $`q^20`$ (see ) and the planned measurements of deep-inelastic process in a region of small values of $`q^20`$ at JLab and DESY.
## 6 Fake leptons in signal events
The signal events, defined by the $`q\overline{q}l^+l^{}`$ subprocess, also contain some hadrons in the final state. Fortunately, their number is essentialy restricted by the upper limit on the beam energy that may be available at PANDA experiment. This circumstance may simplify greatly the identification of final state particles and the physical analysis due to reduction of the phase space and therefore to the reduction of the number of hadrons and other particles which may be produced in event directly or in the decays cascades of other hadrons. These hadrons may decay within the detector volume <sup>7</sup><sup>7</sup>7For the pion the $`c\tau `$ factor ($`\tau `$ is the mean life time of a particle, $`\tau _\pi =2.6E8`$ sec) is equal to $`c\tau =7.8`$ meter according to PDG. and thus produce the background leptons which may fake the signal leptons ($`\mu `$, e) produced in a signal annihilation subprocess. In signal events with electron-positron pair producton the fake electrons/positrons may appear also from muon decays.
In this Section we shall consider separatly the signal events with muon pair production (subsections 6.1 and 6.3) and the signal events with electron pair production (subsection 6.4). The subsection 6.2 includes the discussion of different kinematical distributions for charged pions which povide the main contribution to fake muons. Our analysis of fake leptons is based on two different samples. One of them is the sample with signal muon pair events production while the second one with electron pair production events. This samples were used in Section 2.
### 6.1 Fake muons
We shall first discuss the case of background muons which may be produced additionally to a signal ”$`\mu ^+\mu ^{}`$”- pair in the signal events due to hadron decays. For this reason, we shall call these fake muons also as “decay muons” (in the following we do not distinguish $`\mu ^+`$ and $`\mu ^{}`$). The distribution which shows the number of the main hadronic “parents” of muons in signal events <sup>8</sup><sup>8</sup>8 This kind of information can be extracted from PYTHIA event listings. is presented in the plot a of Fig.14, while Fig.14 b shows the distribution of muon “grandparents”, i.e. the “parents of muon parents”.
The correspondence between the bin number on the x-axis and the name of the related grandparent of the muon can be found from the right column of Table 7 of the Section 11 (Appendix: Tables). <sup>9</sup><sup>9</sup>9Sometimes, when the sea strange quarks take part in the fundamental $`q\overline{q}l^+l^{}`$ interaction, a parent virtual K-meson appears in PYTHIA event listing together with the “$`K^{+/}`$ -like string” (i.e. this string includes a strange quark) which origin flag points onto the colliding proton in the PYTHIA output listing. This case do corresponds to the line which is named as “$`K^{+/}`$-like string ” in the Table 7 of the Section 11 (Appendix: Tables). It is seen from the plot Fig.14 a that the charged pions (bin 2) deliver the main decay muons background while the contribution of $`K^\pm `$-mesons (bin 4) is of about one order less. From Fig.14 b and the right column of Table 7 of the Section 11 (Appendix: Tables) one can conclude that the strings (bin 2 in Fig.14 b), $`\omega `$ (bin 6) and $`\mathrm{\Lambda }^0`$ (bin 20) are the main grandparents of muons. It is of interest to consider in more detail the kinematical distributions of charged pions as the main source of fake muons and to compare them with the distributions of produced muons. We shall do it in the following subsection.
### 6.2 Kinematical distributions of parent pions
The left-hand side of Figure 15 includes (like Figs.2 and 3) three plots (a, c, e) containing the distributions (top to bottom) of the number of events versus the energy $`E_\pi `$ (plot a), the transverse momentum $`PT_\pi `$ (plot c) and the polar angle $`\theta _\pi `$ (plot e) of charged pions which appear in the sample of generated by PYTHIA signal muon events (100 000) based on the quark level subprocess $`q+\overline{q}\mu ^+\mu ^{}`$ (this sample was discussed in the Section 2 and in the begining of this Section).
The plot b in the same Fig.15 shows the distribution of number $`N_\pi `$ of charged $`\pi `$-mesons produced per event in generated signal events. The first left bin in this plot shows the number of events without charged pions. One may see that there is a huge number of signal events (about $`42\%`$) which do not contain at all any charged pions (the final states in the most part of these events, as it can be seen from the PYTHIA event listings, do include mostly nucleon-antinucleon pairs). Thus, with a good accuracy, we may expect that about $`42\%`$ (taking in account charged kaons decays) of events, which include the signal muon pairs, will not contain any additional fake decay muons. Let us underline that PYTHIA provides a good (at least one of the best if not the best) but still a model approximation <sup>10</sup><sup>10</sup>10Due to the fact that there is no complete physical and theoretical understanding of parton-to-hadron fragmentation processes, so far. to the hadronization processes.
From the second and the third bins of $`N_\pi `$ of the same plot b one sees that about $`24\%`$ of signal events may have only one charged pion and about $`27\%`$ of events may have two charged pions in the final states. The other bins of the same plot b demonstrate that about $`5\%`$ of events may have three charged pions, about $`1.5\%`$ of events do contain 4 final-state charged pions and there is a very small fraction of events containing from 5 to 6 final-state pions <sup>11</sup><sup>11</sup>11there is only one event which includes 7 pions is found among 100 000 of generated events..
The origin of all produced pions may be seen from the plot d (“$`\pi `$’s parents”) of Fig.15, where each bin on the x-axis corresponds to the parent particle of a pion (one can find the correspondence of the number of the bin to the name of a parent particle using the left-hand column of Table 8 of Section 11 (Appendix: Tables). From this plot d and Table 8 one can see that the dominant source (bin 2) for pion production are the strings (about $`22\%`$) which are one of the main objects in LUND fragmentation model. The bins 3-6 do correspond, respectively, to the $`\rho `$-, $`\eta `$\- and $`\omega `$-mesons. The decays of these light vector mesons give quite a noticeable contribution (about $`39\%`$ of all entries). The contribution of parent K-mesons (bins 7-10) is close to $`4\%`$ of all entries. The next sizable portion of pions (about $`21\%`$ of entries) comes from the family of $`\mathrm{\Delta }`$-resonances (bins 13-16) and from decays of $`\mathrm{\Lambda }^o`$ (see the corresponding 19-th bin which contains about $`10\%`$ of entries).
In the same way, the corresponding distributions of pion’s grandparents are shown in the plot f of Fig.15 (see also the right-hand column of Table 8 of Section 11, (Appendix: Tables). One can see that among the grandparents the denominating position belongs to the strings (their 4-th bin contains more than $`74\%`$ of entries). Then follow the bins which include K-mesons (bins 7-9) and also $`\eta ^{^{}}`$\- and $`\varphi `$\- mesons (bins 10 and 11) which all together include about $`2.5\%`$ of events. A new object in this grandparents plot, comparing to that one of the parents, is a group of diquarks (bins 12-14), which total contribution is about $`18.5\%`$ of entries. The group of $`\mathrm{\Sigma }`$\- and $`\mathrm{\Xi }`$-resonances (bins 16-23) gives a bit more than $`3.5\%`$ of pions parents.
### 6.3 Kinematical and vertex distributions of fake decay muons in signal events
Fig.16 includes a set of plots with the distributions of background fake decay muons, which are contained in the generated sample of 100 000 signal muon events described in the Section 2. This decay muons come from all possible decay channels of produced hadrons including the pion decays which were discussed above. It is clear that not all hadrons will decay within the detector volume. Therefore, in this subsection we have used the existing PYTHIA option which allows to take into account the restricted decay volume. As for the first approximation for a real PANDA detector volume we have chosen the cylinder with R=2.5m and L=8m. Because of this the number of entries in all the plots of Fig.16 (exept the plot b, which shows the number of all generated events) is equal to 16601. It means that the fraction of signal processes, which include fake muons, reduces to about 16.6 $`\%`$ after taking into account the detector size.
The distribution of a number of muon signal events versus a number of fake decay muons $`N_{bkg}^\mu `$, contained in each event, is shown in the plot b. It is seen that there can be up to 4 muons in the final state. From the first bin of this plot one may see that about $`83\%`$ of events have no background fake muons at all. This value agrees well with the number of entries shown in the plot a which was discussed above. The amount of events free of fake muons does not corresponds to the previously discussed amount of the events without charged pions because of applied restriction of the decay volume. It is also seen from the plot b that the numbers of events with one fake muon, with two and three fake muons are about by, respectively, one, two and three orders less than the number of events without any fake muon.
The left column of Fig.16 includes (from top to bottom) the energy $`E_{dec}^\mu `$, transverse momentum $`PT_{dec}^\mu `$ and polar angle $`\theta _{dec}^\mu `$ distributions of muons which appear from hadron decays. From comparison of these plots with their analogs from Fig.2, i.e. of those which are done for the signal muons, one may see that the fake muons are less energetic than the signal ones. One may also find that the mean value of the signal muons energy (see plots a, b of Fig.2) $`<E^\mu >`$ = 2.6 GeV corresponds to such a point in the energy spectrum of fake muons which decay in the restricted volume, where the contribution of fake muons (in the same signal events) is very low. Analogously, the mean value of the PT-distribution of signal muons $`<PT_{signal}^\mu >=0.7`$ GeV (see Fig.2) corresponds to that point where the spectrum of $`PT_{dec}^\mu `$ practically vanishes.
Therefore, one has to look for the set of some reasonable cuts (chosen with an account of the real detector effects) on a muon energy $`E^{mu}`$ as well as on its $`PT^{mu}`$ value, which may lead to an essential reduction of decay muons contribution and to keep at the same time the main part of signal events. For example, the comparison of the plots a and c of the Fig.2 with the analogous plots in Fig.16, leads to a conclusion that the cuts $`E^\mu >0.2`$ GeV, $`PT^\mu >0.2`$ GeV may allow to get rid of about $`66\%`$ of decay muons in the signal events and to save the most of signal (i.e. belonging to signal $`\mu ^+\mu ^{}`$\- pair) slow and fast muons. We shall return to this problem a bit later.
Another way which may help to discriminate the signal muons from the ”decay” ones is to use the information about the position of fake muon production vertex and the reconstruction of the invariant mass of the parent and grandparent hadrons. The plots d and f of Fig.16 contain the distributions of Vx- and Vz- (z-axis is chosen along the beam direction) components of the 3-vector $`V=(Vx,Vy,Vz)`$, which gives the position of a fake muon production vertex in millimeters (mm) Let us note that these plots are obtained within the PYTHIA level of simulation, i.e. without an account of details of detector construction and the effects caused by the magnetic field. Withing this approximation the distributions of Vx- and Vy- components have to be similar. For this reason, we shall show in the following only the distribution of Vx- component. Seen from the Fig.16 the Vz component of about $`35\%`$ of events may be rather close to zero, i.e. to the interaction point. The corresponding fake muons, produced near the interaction point, may give rise to the most difficult background. The contribution of the background muons from the other type of events (based mainly on minimum bias and QCD partonic subprocesses) will be discussed in the following Section 7.
### 6.4 Fake electrons
Now let us consider the situation with fake decay electrons and positrons background in the case of signal processes based on the quark level subprocess of electron-positron pair production <sup>12</sup><sup>12</sup>12 In the following we shall use “electron” as a common name for the background electrons and positrons.. Fake electrons produced in these decays will be also called in the following as “decay” ones. Recall, as it was already mentioned in the Section 2 and in the begining of the present one, that we shall use the sample of 100000 generated signal events with $`e^+e^{}`$ production.
Fig.17a presents the contribution of different parent particles into the process of creation of the fake electrons. It clearly demonstrates that among all shown sources of fake electrons the contribution of neutral pion $`\pi ^0`$ decay (bin 2) is a dominat one. It provides a much higher (about of one order) contribution than all of other decay channels. The electrons/positrons may appear in neutral pion decay only through the channel of Dalitz decay into a photon and an electron-positron pair: $`\pi ^0\gamma +e^+e^{}`$. The next contribution, which is by about one order less, give, in decending order, muon (bin 1), $`\eta `$\- and K-mesons decays as well as the decay of $`\mathrm{\Lambda }^0`$. Neutral pions in their turn may arise from decays of $`\eta `$\- and $`\omega `$-mesons or heavier mesons and baryons, produced as resonance states according to the LUND fragmentation model.
Fig.17 b shows the distribution of a number of background electrons versus the type of their grandparents. The correspondence between the bin numbers shown on the x-axis and the name of the grandparent particle can be found from the left column of the Table 7 of Appendix: (Tables which are presented in the Section 11). From comparison with this Table one can see that the strings (bin 2) are the main source of electron parents. Then follow a group of $`\rho `$-, $`\eta `$\- and $`\omega `$-mesons (bins 5-7) as well as the group of $`\mathrm{\Delta }`$-resonances (bins 15-16) and $`\mathrm{\Lambda }^o`$ (bin 20). Two orders less contribution is provided by the group of kaons (bins 8-11) which is followed by $`\eta ^{^{}}`$ resonance (bin 12) and the group of $`\mathrm{\Sigma }`$-resonances (bins 22-24).
The plots, containing the distributions of the fake decay electrons which appear in signal events (based on the $`q\overline{q}e^+e^{}`$ subprocess) are presented in Fig.18. They are done basing on the sample of 100 000 generated by PYTHIA signal events which were discussed in the Section 2 and contain signal $`e^+e^{}`$ pairs. Like in previouse subsection we use here the same restriction on the detector volume.
From the statistics frame in the upper part of plots one may see that in a case of decay electrons the number of entries is 1926 (i.e. the fraction of the signal processes which includes fake electrons is about $`2\%`$).
The comparison of the energy $`E_{dec}^e`$ and transverse momentum $`PT_{dec}^e`$ distributions of fake electrons in signal events, which are given in the left-hand column of Fig.18, with the analogous plots for fake muons from the left-hand column of Fig.16, allows to conclude that in electron case the distributions fall at least twice steeply than in muon one.
Analogously with the previous subsection, the comparison of the same $`E_{dec}^e`$ and transverse $`PT_{dec}^e`$ distributions of fake electrons in signal events, shown in plots a and c of Fig.18, with the plots a and c of Fig.2 for signal electrons leads to the conclusion that the soft cuts like $`E^e>0.2`$ GeV and $`PT^e>0.2`$ GeV may allow to eliminate the most of fake electrons at the cost of about $`10\%`$ loss of signal events.
The right-hand column of Fig.18 includes plots b and d. They contain the values of the Vx- and Vz- components of the 3-vector V which points the position of electron production vertex. In contrast to the form of vertex distribution for muon production case (see the plot d in Fig.16), the concentration of Vz-component of electron production (see plot d in Fig.18) near the interaction point poses the essential difference of a $`e^+e^{}`$ channel as compared to a case of $`\mu ^+\mu ^{}`$ channel. In the last case, the most of background muons are produced by light charged pions which may decay in flight at a rather large distance from the interaction point. It is also seen from plot f in Fig.18 and the analogous plot a in Fig.16 that the number of fake electrons or fake muons in signal events may be up to four in both cases.
### 6.5 Cuts for fake leptons reduction in signal events
The analysis of distributions discussed above leads to the conclusion that the following cuts:
$``$ 1) we select the events with the only two leptons with $`E_l`$ 0.2 GeV, $`PT_l`$ 0.2 GeV;
$``$ 2) the charges of these two leptons must be of the opposite sign;
$``$ 3) the vertex of lepton origin lies within the range $`R_{vtx}`$ 15 mm from the interaction point;
which, being applied to the sample of signal events, can allow one to select a subsample which include a strongly reduced fraction (fr) of events containing fake leptons ($`fr_e=0.008\%`$ for the case of electron pair production and $`fr_\mu =0.001\%`$ for the muon pair production). The loss of the signal events due to application of cuts 1)–3) is shown in Table 3.
One can see from this Table that it is possible to select the signal events which are almost free of background fake leptons at the cost of diminishing of the signal events sample by $``$ 17 $`\%`$ for $`\mu ^+\mu ^{}`$ and $``$ 14$`\%`$ for $`e^+e^{}`$ production.
## 7 QCD and minimum-bias background events
The other source of the background is the leptons produced in the minimum-bias (low - $`P_T`$ and diffractive scattering) events and QCD background these are mainly $`q+gq+g`$, $`g+gg+g`$ and $`q+q^{}q+q^{}`$) processes, where the possibility of appearance of two (and more) leptons in the final state is very high. For analysis of these processes $`10^6`$ events of antiproton diffraction over proton target with $`E_{beam}=`$14 GeV were generated with PYTHIA 6.4. These events include mentioned above processes, including the signal one $`\overline{q}+ql^+l^{}`$. According to PYTHIA the total cross section of these processes $`\sigma _{tot}^{bkg}=50.17`$mb is about $`10^7`$ times higher than the cross section of the signal MMTDY subprocess $`q+\overline{q}l^++l^{}`$: $`\sigma ^{\overline{q}ql^+l^{}}=5.5710^6`$mb. In the following subsections 7.1-7.2 we shall present the distributions obtained without use of any cuts.
### 7.1 Muon background.
The distribution of the parents of muons produced in background minimum-bias and QCD events is presented in Fig.19 a. It is seen that, like in a case of fake muons in signal events (see Fig. 14 a), the main contribution comes from $`\pi ^\pm `$\- and $`K^\pm `$-meson decays. Fig.19 b <sup>13</sup><sup>13</sup>13 which is slightly different from its analog in Fig.14 for the signal events. shows the distribution of muons grandparents. One can see (using the right-hand column of Table 7 of the Section 11 (Appendix: Tables) that the main grandparents of muons in QCD and minimum-bias events are the clusters and strings (bins 1, 2) as well as the $`\rho `$-, $`\eta `$\- and $`\omega `$-mesons (bins 3-6). Then follow the $`\mathrm{\Delta }`$-resonances (bins 16, 17) and $`\mathrm{\Lambda }^0`$ (bin 20).
The kinematical and other distributions of muons produced in the above mentioned generated background minimum-bias and QCD events are shown in Fig.20 It is seen that the kinematical distributions (plots b, c, e) do not differ so much (that is natural) from those of fake “decay” muons produced in the signal $`p\overline{p}l^+l^{}+X`$ processes (see Fig.16).
The distribution of the number of generated background events versus the amount of fake muons produced per event, i.e. $`N_{bkg}^\mu `$, is shown in plot b) of Fig.20. It differs noticeably from its analog shown in Fig.16 which contains only the distribution of fake ”decay” muons in signal lepton pair production events. The number of muons in the final state, contained in one and the same background event, can be up to 7. It means that the probability of production of the pair of fake muons with the charges of the opposite signs (like in signal events) is rather high in these background events. Such pairs may fake quite well the signal events.
The distribution plots of the production vertexes of background muons are shown in the right column of Fig.20. It is seen from the plot f of Fig.20 that the most of the muon production vertexes are spread over detector volume while for some of events these vertexes are rather close to the interaction point. So the information about the vertex position can be useful for background separation.
### 7.2 Electron background
Let us consider now the case of electrons produced in the background minimum-bias and QCD events. The distribution of the parents of background electron in the discussed above sample of minimum-bias and QCD events is presented in Fig.21 a. It is seen that, like in a case of fake electrons in signal events (see Fig. 17 a), the main contribution comes from $`\pi ^0`$-, $`\eta `$-, charged $`K`$-mesons (bins 2, 5 and 7, respectively) and also from muon decays (bin 1). The main source of fake electrons are decays of neutral pions ($`\pi ^0\gamma +e^+e^{}`$). They can appear directly or from decays of $`\rho `$-, $`\eta `$-, $`\omega `$\- , K- mesons, as well as from $`\mathrm{\Delta }`$\- resonances decays.
From the Fig.21 b one may see that in the minimum-bias and QCD events sample the structure of the distributions of electron grandparents is rather
different from its analog presented in Fig.17 b for a case of fake electrons in the sample of signal events with the signal electron-positron pairs. Namely, different to the plot Fig.17 b <sup>14</sup><sup>14</sup>14after normalization to an equal number of entries. the charged $`\rho `$-mesons (bin 5) takes the dominant position in the plot Fig.21 b. Then follows the noticeably increased contribution of clusters (bin 1) which height reaches the height of strings contribution (bin 2). It is also seen that on total the contribution of light vector mesons ($`\rho `$, $`\eta `$ and $`\omega `$ (bins 5-7)), as well as the contribution of K- (bins 8-11) and $`\eta ^{^{}}`$-, $`\varphi `$\- mesons (bins 12,13), have grown up as comparing to the higher bins (15, 16 and 20), which correspond to $`\mathrm{\Delta }`$\- and $`\mathrm{\Sigma }`$\- barions, respectively.
The distributions obtained from the sample of mentioned above generated minimum-bias and QCD background events, are shown in Fig.22. The kinematical distributions (plots a, c and e) are rather similar to the distributions of fake decay electrons in the signal events which were discussed in the subsection 6.4 and presented at Fig.18. From Fig.22 one may see that the total number of electrons in the sample of generated 1000000 minimum-bias and QCD events is equal to 37885. This number is of order less than the number of fake muons produced in the same sample of $`10^6`$ minimum-bias and QCD events (see previous subsection 7.1). Plot b) of Fig.22 shows the distribution of the number of generated minimum-bias and QCD events versus the number of decay electrons per event. The third bin in this plot, as well as the other bins to the right from it, show how many events may contain two and more electrons. In these events there may appear the $`e^+e^{}`$-pairs, which potentially may fake the signal events. It is clearly seen that the probability of appearance of 2 and more electrons in the final state reduces to a value of about few percents of the total number of generated events.
The plots d and f of Fig.22 show the distributions of the position of the electron production vertex in the background sample in the transversal ($`Vx_{dec}^e`$) and the longitudinal ($`Vz_{dec}^e`$) directions. It is seen that the most of background electrons, originating from hadron decays, are produced near the interaction point ($`V_x=0`$ and $`V_z=0`$), like it take place in Fig.18.
## 8 Background separation
To reduce the background contribution from the minimum bias and QCD events we added two new cuts to the previously used cuts 1)-3) (see Section 6.5). Finaly, we use the following selection cuts:
$``$ 1.) the events with the only two leptons with $`E_l`$ 0.2 GeV, $`PT_l`$ 0.2 GeV;
$``$ 2.) the charges of these two leptons must be of the opposite sign;
$``$ 3.) the vertex of lepton origin lies within the radius $`R_{vtx}`$ 15mm
from the interaction point;
$``$ 4.) $`M_{inv}(l^+,l^{})`$ 0.9 GeV;
$``$ 5.) lepton isolation criteria: the summed energy E<sub>sum</sub> of all the particles around
the lepton within the cone of the radius R=$`\sqrt{\mathrm{\Delta }_\phi ^2+\mathrm{\Delta }_\eta ^2}`$ = 0.2 in the $`\eta \phi `$ space
is not higher than E$`{}_{sum}{}^{max}=0.5`$ GeV <sup>15</sup><sup>15</sup>15the azimuth angle $`\phi `$ and the polar (zenith) angle $`\theta `$ are used to determine the direction of the 3-momentum of any particle. $`\eta `$ is the particle pseudorapidity, defined by the formula $`\eta =\mathrm{ln}tg(\theta /2)`$, where $`\theta `$ is the polar angle of the particle 3-momentum counted from the beam direction..
Here $`\mathrm{\Delta }_\phi =\phi _l\phi _p`$ is the difference of the lepton’s (l) azimuth angle $`\phi _l`$ and the azimuth angle $`\phi _p`$ of the particle (p), contained within the cone of the radius R around the lepton. Analogously, $`\mathrm{\Delta }_\eta =\eta _l\eta _p`$ is the difference of the lepton and the particle pseudorapidities.
Few words are in order now about the choice of the last two cuts. Tables 4 and 5 show, for muon pair and electron pair production cases, respectively, the influence of the variation of the cut on dilepton invariant mass $`M_{inv}(l^+,l^{})`$ on the loss of signal events and the value of signal to background ratio S/B (after application of the cuts 1.)-3.). It is seen from Table 4 that in $`\mu ^+\mu ^{}`$ case the growth of the maximal value of $`M_{inv}(l^+,l^{})`$ up to $`M_{inv}(l^+,l^{})=1.2`$ GeV allows to get rid completely of the background at cost of loosing of 35 $`\%`$ of signal events, while in $`e^+e^{}`$ case, see Table 5, the same upper limit leads only to S/B$`=`$ 2.3.
The results of all the five cuts sequent application to the sample of inelastic $`\overline{p}pX`$ events which contains the minimum-bias and QCD events (incuding the signal events based on the parton level annihilation subprocess $`q\overline{q}\gamma ^{}l^+l^{}`$ ) are collected in the Table 6. It is seen that the first three cuts allow to enlarge the the S/B ratio by about one order in a case of muon pair production. At the same time, the third cut is inefficient for the case of $`e^+e^{}`$-pair production. The forth cut $`M_{inv}(l^+,l^{})`$ 0.9 GeV allows to increase the S/B ratio in $`e^+e^{}`$ case by more than two orders and by about three ordes the S/B ratio for $`\mu ^+\mu ^{}`$ case.
The plots presented in the Figs.23, 24 are done to illustrate the action of the lepton isolation criterion used in the definition of the fifth cut. They show the distributions of the total energy of the particles which are contained within the cones of the radius R around the leptons. Figures 23 and 24 present these distributions for the electron and muon cases, correspondingly. By comparing the plots a (for leptons from the signal events) with the plots b (for background leptons from minimum-bias and QCD events) one can easily see that the signal events have much smaller summarized energy content within the cone of $`R0.2`$ than the energy content in background events. This observation is used in the cut 5.
From the Table 6 one can see that the last cut on the lepton isolation, i.e. the choice of only those final state leptons which have the restricted value of the summarized energy (not greater than $`E_{sum}`$=0.5 GeV) of other particles contained within the cone of some fixed radius R$`=\sqrt{\phi ^2+\eta ^2}=`$NR (NR$`=0.1,0.2,0.3\mathrm{}`$) around the direction of the lepton 3-momentum, allows one to achieve (choosing NR$`=0.2`$) the value of the signal to background ratio equal to S/B = 3.8 for electron production case and completely to get rid of background in muon production case. In both cases the application of the fifth cut leads to additional 8$`\%`$ loss of the signal events left after application of the first four cuts. Let us note that the same criteria, but with the use of a more restricted form of the forth cut $`M_{inv}(e^+,e^{})`$ 1.0 GeV, allows to increase the signal to background ratio up to S/B = 9 in $`e^+e^{}`$ case.
## 9 Remarks on the possibility of measuring the multi-parton interactons and the intrinsic quark transverse momentum in the proton.
In addition to the opportunity to get the information about parton distribution functions, which was already discussed in Section 5, let us also mention here three other ones which may be useful for studying quark dynamics in proton and its PDFs.
The first two are connected with the processes of two
$$\overline{p}pl_1^+l_1^{}+l_2^+l_2^{}+X$$
(4)
or even three lepton pairs
$$\overline{p}pl_1^+l_1^{}+l_2^+l_2^{}+l_3^+l_3^{}+X$$
(5)
production in one event.
Both of these processes can include two and three quark annihilation $`q\overline{q}l^+l^{}`$ subprocesses, respectively. The total cross sections of such processes can be smaller as compared to the $`\overline{p}pl^+l^{}+X`$ process which includes a single $`q\overline{q}l^+l^{}`$ subprocess. Nevertheless, they may contain interesting physical information which can be more easily extracted at the intermediate energies than at the higher ones.
First, the measurement of the characteristics of the system of other than lepton pairs particles, produced in the process (4), will give us the opportunity to get the information about the so-called ”underlying” event. The study of the analogous distribution in the process (5), in which all valence quarks (and antiquarks) in proton (and antiproton) will annihilate into lepton-antilepton pairs, may provide the information about gluon content in the proton. The understanding of the physics of the ”underlying” event, i.e. the interaction of partons which do not participate in the hard subprocess $`q\overline{q}l^+l^{}`$, is very important for the interpretation of the results of the present Tevatron and future LHC experiments.
The second opportunity, also provided by the processes (4) and (5), is the study of the so-called multiple parton hard interaction processes in pp- and $`p\overline{p}`$ interactions which are widely discussed in connection with the problem of a proper account of background contribution to the processes which are planned for seaches of the New Physics signals at Tevatron and LHC. It is worth mentioning that the measurements of the proceses (4) and (5) can be done for the case when final state lepton pairs would have different flavor, like $`\overline{p}pe^+e^{}+\mu ^+\mu ^{}+X`$. In such case we shall have the situation when four hadronic jets events, used in previous measurements of multiple interactions -, are substituted by events withfour leptons. This substitution shall increase the precision of measurement of the parameters of multiple interections as it was shown in recent measurements done with ”3jet + photon” final states , .
Besides getting the information about the fraction of multiple interactions their study opens the possibility to get the information about the spatial distribution of quarks within the proton. It is obvious that in a case of uniform distribution of quarks within the proton volume the occurrence of the first parton-parton interaction would not influence on the probability of happening of the second interaction, while in a case when the quarks are concentrated in small region the probability of happening of the second interaction becomes higher if one of the quarks has taken part in the first interaction.
The third opportunity is connected with the possibility to measure the characteristics of internal quark motion in the proton. This possibility is based on the fact that the shape of the distribution of the modulus of the vector sum of quark and antiquark transverse momentum vectors
$`P_T^{q\overline{q}}=|\stackrel{}{P}_T^{q\overline{q}}|=|\stackrel{}{P}_T^q+\stackrel{}{P}_T^{\overline{q}}|`$ (6)
practically coincides (due to the transverse momentum conservation law) with the shape of the above-considered distribution of the modulus of the lepton pair transverse momentum $`P_T^{l^+l^{}}`$, shown in the plot b of Fig.12. <sup>16</sup><sup>16</sup>16 Recall (see Section 4) that all plots in the Fig.12 are done for the case when both “Fermi motion” (or ”$`k_T`$-effect”) and ISR are switched “on”.
The variable $`P_T^{q\overline{q}}`$ is of special interest because it contains the information about two important physical features of quark dynamics inside the hadron. Indeed, in our case when a beam antiproton is directed along the z axis and it scatters over the proton fixed target, there may be only two sources of transverse motion of quarks in the initial state <sup>17</sup><sup>17</sup>17The values of transverse momenta of constituents in a target proton (which is at rest) as well as of those inside a beam antiproton (which moves along the z axis) are invariant under Lorentz boost along the z axis.:
A) internal Fermi-motion of quarks (with some transverse velocity) inside a proton, i.e. the so-called ”$`k_T`$-effect”;
B) initial-state radiation (ISR) of gluons or photons from quarks before hard quark-antiquark annihilation;
The importance of these two effects was recently discussed in connection with the interpretation of prompt photon production study in the experiment E706 at Fermilab and also with the study of ”$`\gamma /Z+jet`$” events, which are sensitive to the shape of gluon distribution, at the LHC and Tevatron , .
## 10 Conclusion.
The modeling of dilepton production in antiproton scattering over proton target $`\overline{p}pl^+l^{}+X`$ is done for the intermediate energy $`E_{beam}=14`$ GeV on the basis of PYTHIA event generator and the parton level subprocess of quark-antiquark annihilation $`\overline{q}ql^+l^{}`$.
The distributions of most essential kinematical variables of individual leptons, are are presented in Section 2. They show that the energy and angle spectra of the fast (most energetic) leptons in a pair are very different from those of slow leptons: the mean value $`<E_{fast}^l>=3.85`$ GeV is about three times higher than that one of slow leptons $`<E_{slow}^l>=1.36`$ GeV. The simulation has also shown a tendency which may be a rather general one: fast leptons fly predominantly at smaller angles $`<\theta _{fast}^l>=16.5^o`$ as compared to the angles of slow ones $`<\theta _{slow}^l>`$= $`38.2^o`$. It is worth noting that about $`6\%`$ of events may have slow leptons, that may scatter into the back hemisphere, i.e. $`\theta _{slow}^l>90^o`$. The angle-energy, energy-energy and angle-angle correlations among a slow and a fast lepton in the same lepton pair in event are also described in Section 2 and are presented in Figs.4, 5 together with the corresponding distributions of the number of events versus the corresponding lepton energies and angles.
These distributions allow one to estimate the energy, transverse momentum and angle ranges that may be covered by leptons produced in quark-antiquark annihilation process. They were useful for proper design of muon system and may be also used for the electromagnetic calorimeter. Tables 1 and 2 show the estimation of the loss of signal events depending on the choise of cuts on the lower values of lepton energy and, respectively, on the upper limit of the angle size of the muon system and the electromagnetical calorimeter. From these Tables one can see that, for instance, the choise of cuts $`E_{cut}^{min}=0.5`$ GeV and $`\theta _{max}^l=90^o`$ results in about 30$`\%`$ loss of signal events. The simulation PYTHIA has shown that one may expect to gain about $`710^7`$ MMTDY events per year for the luminosity $`L=210^5mb^1s^1`$.
The analogous study was done on the basis of PYTHIA in the Section 3 for the leptons which may appear in decay of J/$`\mathrm{\Psi }`$ mesons, produced in the benchmark process $`\overline{p}pJ/\mathrm{\Psi }+X`$. It was shown that the leptons, produced in J/$`\mathrm{\Psi }l^+l^{}`$ decay, fit well into the same angle regions as the leptons produced in MMTDY process $`\overline{p}pl^+l^{}+X`$. The reconstruction of the lepton pair invariant mass can allow to get rid of background without a sizable loss of signal events.
In Section 4 the study of kinematical characteristics of lepton pair as a whole system was done. It is shown that the spectrum of the invariant mass of the lepton pair decreases rather fast and vanishes at $`M_{inv}^{l+l}=2.5`$ GeV. At the same time the lepton pairs total energy $`E^{l^+l^{}}`$ spectrum starts at around 1 GeV and extends up to the value of 12 GeV. It is also demostraited that about a half of the events have the lepton pair energy higher than 5 GeV. It is shown that one can expect that in about 50$`\%`$ of events the lepton pairs would be rather energetic and they can take away from 33$`\%`$ up to 80$`\%`$ of the total energy of the final state system. The square of the invariant mass $`(M_{inv}^{l^+l^{}})^2=Q^2=q^2=(P^{l^+}+P^l^{})^2`$ has the meaning of the square of the momentum transfered from the hadronic system of quark-antiquark pair to the electromagnetic system of final state dilepton pair.
In Section 5 the analysis of distributions, obtained by PYTHIA, allowed to determine the region in x-Q<sup>2</sup>-plane which can be available for measuring the proton structure function at PANDA. This region is defined by the following boundaries: $`0.05x0.7`$ and $`Q^26.25`$ GeV. Let us emphasize that the measurements in this region of positive (”time-like”) $`q^2(P^{l^+}+P^l^{})^2=Q^20`$ would be a good extension of studies planned to be done at JLab in the region of negative (”space-like”) values of $`q_{dis}^2=(P_{in}^lP_{out}^l)^20`$.
An important problem of background suppresion is considered in Sections 6, 7 and 8. In Section 6 we have concentrated on the fake leptons which can appear from hadron decays in the same signal process. We studied signal processes with dimuon production separately from the processes with electron pair production using for this two separate event samples with signal muon pair and, respectively, electron pair production. First we have considered the case of $`\mu ^+\mu ^{}`$ production. The histograms which demonstrate the relative contribution of different parents and grandparents of produced muons are presented in the subsection 6.1. It is shown that charged pion decays produce the main part of fake muons. The next dominant source is the decays of charged K-mesons (its contribution is by more than two ordes less than that one from pions).
Some details about pion source are presented in the subsection 6.2 where it was shown that about 42$`\%`$ of signal events do not contain at all any charged pions, while 24$`\%`$ of them contain only one pion and 27$`\%`$ include two charged pion in the final state. About 5$`\%`$ of signal events include tree charged pions and 1.5$`\%`$ have four charged pions. This prediction of PYTHIA indicate that the reconstraction of invariant masses of parent and grandparent hadrons can be a quite reliable way for fixing the origin of produced fake muons.
The kinematical plots, shown in subsection 6.3 (and in subsection 6.4 for fake electrons) are done by applying a geometrical restriction on the detector volume. This restriction produced a strong reduction of the fraction of fake leptons. In result the energy and transverse momentum spectra of fake muons became essentialy shorter than their analogs shown in the Section 2. The fraction of signal events which include fake muons has reduced down to about 16.6$`\%`$ (in comparison with pion distributions). It means that, according to PYTHIA, about 83$`\%`$ of signal events would not contain fake muons at all due to the used restriction of the decay volume. The analogous parent, grandparent and kinematical distributions were obtained in subsection 6.4 for a case of background electrons. It is found that the application of geometrical restriction on the detector volume provides the reduction of signal events fraction (containing fake electrons) down to 2$`\%`$.
The set of three cuts are proposed in the subsection 6.5 which, being applied together with the mentioned above restriction on the detector volume, allows a further reduction of the fraction of the signal events containig fake decay leptons. Namely, they are: $`fr_\mu =0001\%`$ in a case of $`\mu ^+\mu ^{}`$ production and $`fr_e=0.008\%`$ in $`e^+e^{}`$. Let us underline that this strong reduction of the fraction the signal events including fake leptons was achieved at the cost of the loss of a noticeable number of selected signal events. These losses are: $`17\%`$ for $`\mu ^+\mu ^{}`$ and $`14\%`$ for $`e^+e^{}`$ production.
Much more dangerous background, provided by minimum-bias and QCD events, was studied in the Sections 7 and 8. The former includes two subsections in which the figures with parent and grandparent relative contributions as well as the with kinematical distributions for $`\mu ^+\mu ^{}`$ and $`e^+e^{}`$ cases are presented, respectively. Subsection 8 contains the set of five cuts which include three cuts which were previously used to reduce the number of events containing fake decay leptons. The fourth cut reduces the spectrum of the invariant mass of dilepton system by the condition $`M_{inv}(l^+,l^{})0.9`$ GeV, while the fifth cut uses the isolation criteria for the lepton. It is shown that dispite the fact that the cross section of minimum-bias process $`\sigma _{tot}^{minbs}`$ is by about 7 oders larger than the cross section of the signal process $`\sigma _{tot}^{\overline{q}ql^+l^{}}`$ the application of these five cuts allows to get rid completely of minimum-bias and QCD background contribution in $`\mu ^+\mu ^{}`$case and to reach the value of S/B = 3.8 for $`e^+e^{}`$ case. It is worth mentioning that the application of the fourth and the fifth cuts leades to an additional loss of the number of selected signal events by 8$`\%`$.
The Section 9 contains three important remarks about the physical potential of studing events having several lepton-antilepton pairs in the final state. First, it is stressed that the study of events with two (and, maybe, even three) lepton pairs would allow to enlarge the precision of the parameters of multiple quark interactions. It will also essentialy extend the region of QCD studies because up to now there were done only five dedicated measurements of such events in proton-proton and antiproton-proton collisions. The first three processes considered the case when four jets were produced in the final state, while the last two measurements have used the events in which the final state was including tree jets plus one direct photon. Therefore, the measurements in which the jets would be substituted by leptons would allow to reach a higher level of precision.
It was also noted that at the same time such events, based on the processes of valence quarks annihilation in the colliding protons, will provide a clean information about the dynamics of spectator quarks interactions, i.e. about the so-called ”underlaying events”. It is worth mentioning that the most interesting would be the measurements with three leptons pairs production, because in this case the underlaying processes would be defined mostly by soft gluon interactions.
The third opportunity, which was discussed in the Section 8, can be based on the measurement of the transverse momentum of lepton pair which is directly connected with the transverse momentum of the system of two annihilating quarks. The latter can be caused by the so-called ”$`k_Teffect`$” which is connected with the so called ”Fermi-motion” of quarks inside the proton or the radiation of gluons by initial-state quarks. This information is of a big interest for the interpretation of QCD effects observed at high energy hadron colliders.
It should be underlined that the present PYTHIA simulation does not take into account the detector effects (like the magnetic field and the material of the apparatus, for example). Hence, the obtained plots, which describe the unbiased distributions of the produced free particles, may be mainly useful for preliminary estimations and working out the criteria (cuts) for selection of experimental events for further physical analysis. The detailed GEANT simulation with account of detector design (and based on the simulated PYTHIA event sample described here) will be a subject of our following publication. Nevertheless, one can expect that there is a high probability that the main features of the real process, corrected to detector effects, would be rather similar to those shown in the plots presened in this article.
## 11 Acknowledgements.
One of the authors (A.N.S.) acknowledge support from Russian-German ”FAIR Russia Research Centre“ (which is also supported from Russian Federal Agency for Atomic Energy ”Rosatom“ and German Helmholtz Association).
## 12 Appendix: Tables. |
warning/0506/math0506459.html | ar5iv | text | # An Extension of Barbashin-Krasovski-LaSalle Theorem to a Class of Nonautonomous Systems
## 1 Introduction and Main Result
Let us consider the following time-varying dynamical system:
$$\dot{x}=f(t,x),xD,t𝐑$$
(1)
where $`D`$ is a domain in $`𝐑^n`$ containing the origin ($`0D𝐑^n`$). About $`f`$ we suppose the following:
1) $`f(t,0)=0`$, for any $`t𝐑`$;
2) Uniformly continuous in $`t`$, uniformly in $`xD`$, i.e. $`\epsilon >0\delta _\epsilon >0`$ s.t. $`t_1,t_2𝐑,|t_1t_2|<\delta _\epsilon `$ and $`xD,f(t_1,x)f(t_2,x)<\epsilon `$;
3) Uniformly local Lipschitz continuous in $`x`$ for any $`t𝐑`$ , i.e. for any compact set $`KD`$, there exists a positive constant $`L_K>0`$ such that:
$$f(t,x)f(t,y)L_Kxy,\mathrm{for}\mathrm{any}x,yK\mathrm{and}t𝐑$$
4) bounded in time, that means there exists a continuous function $`M:D𝐑`$ such that:
$$f(t,x)M(x),\mathrm{for}\mathrm{any}t𝐑$$
With these hypotheses we know that for any $`(t_0,x_0)𝐑\times D`$ there exists a unique solution of the Cauchy problem:
$$\begin{array}{c}\dot{x}=f(t,x)\\ x(t_0)=x_0\end{array}$$
(2)
with the initial data $`(t_0,x_0)`$; we denote by $`x(t;t_0,x_0)`$ this solution. One can define this solution for $`t(t_0T,t_0+T)`$ where $`T=sup_{r>0,B_r(x_0)D}\frac{r}{f_{B_r(x_0)}}`$, the supremum is taken over all positive radius such that the ball centered around $`x_0`$, $`B_r(x_0)=\{x𝐑^n|xx_0<r\}`$, is completely included in $`D`$ and $`f_{B_r(x_0)}=sup_{(t,x)𝐑\times \overline{B}_r(x_0)}f(t,x)`$ is a supremum norm of $`f`$ with respect to $`B_r(x_0)`$ (where is no confusion we denote $`B_r=B_r(0)`$). The function $`\gamma _{t,t_0}(x_0)=x(t;t_0,x_0)`$ is well defined for some bounded open set $`S`$, $`\gamma _{t,t_0}:SUD`$ (with $`U`$ open and bounded) and it is Lipschitz continuous with a Lipschitz constant given by $`L=exp(L_U|tt_0|)`$ ($`L_U`$ being the Lipschitz constant associated to $`f`$, as above, on the compact set $`\overline{U}`$). All these results can be found in any textbook of differential equations (for instance see ).
Our concern regards the stability behaviour of the equilibrium point $`\overline{x}=0`$. First we recall some definitions about stability (in Liapunov sense).
Definition We say the equilibrium point $`\overline{x}=0`$ for (1) is uniformly stable, if for any $`\epsilon >0`$ there exists $`\delta _\epsilon >0`$ such that for any $`t_0𝐑`$ and $`x_0𝐑`$ with $`x_0<\delta _\epsilon `$ the solution $`x(t;t_0,x_0)`$ is defined for all $`tt_0`$ and furthermore $`x(t;t_0,x_0)<\epsilon `$, for every $`t>t_0`$.
Definition We say that the equilibrium point $`\overline{x}=0`$ for (1) is uniformly asymptotic stable, if it is uniformly stable and there exists a $`\delta >0`$ such that for any $`t_0𝐑`$ and $`x_0D`$ with $`x_0<\delta `$ the solution $`x(t;t_0,x_0)`$ is defined for every $`tt_0`$ and $`lim_t\mathrm{}x(t;t_0,x_0)=0`$.
If in the definition of uniform stability we interchange ’there exists $`\delta _\epsilon >0`$’ with ’for any $`t_0𝐑`$ ’ (thus $`\delta `$ will depend on $`\epsilon `$ and $`t_0`$, $`\delta _{\epsilon ,t_0}`$) then the equilibrium is said (just) stable. If we proceed the same in the second definition we obtain that the equilibrium is asymptotic stable. For time-invariant systems there is no distinction between uniform stability and stability, or uniform asymptotic stability and asymptotic stability. In general case, the uniform (asymptotic) stability implies (asymptotic) stability, but the converse is not true (see for instance ).
We say that the dynamics (1) has a positive invariant set $`N`$ if for any $`t_0𝐑`$ and $`x_0N`$ the solution $`x(t;t_0,x_0)N`$ for all $`tt_0`$ for which it is well-defined. Then it makes sense to consider the dynamics restricted to $`N`$, i.e. the function:
$$X:𝐑^+\times 𝐑\times NN,X(\tau ;t_0,x_0)=x(\tau +t_0;t_0,x_0)$$
where $`\tau `$ runs up to a maximal value depending on $`(t_0,x_0)`$. Moreover, by considering the case of $`f`$ from (1) we obtain that $`X(\tau ;t_0,0)=0`$, for any $`\tau >0`$, $`t_0𝐑`$. Therefore we may define the corresponding stability properties of the restricted dynamics as above, where we replace $`D`$ by $`N`$.
The main result of this paper is given by the following theorem:
###### THEOREM 1
Consider the time-varying dynamical system $`(\text{1})`$ for which $`f`$ has the properties $`1)4)`$. Suppose there exists a function $`V:D𝐑`$ of class $`𝒞^1`$ such that:
$`H1)`$ $`V(x)0`$ for every $`xD`$ and $`V(0)=0`$;
$`H2)`$ There exists a continuous function $`W:D𝐑`$ such that:
$$\frac{dV}{dt}(t,x)=V(x)f(t,x)W(x)0$$
$`H3)`$ Let $`E=\{xD|W(x)=0\}`$ denote the zero-set (or kernel) of $`W`$; suppose that $`f`$ resticted to $`E`$ is time-invariant (i.e. $`f(t,x)=f(t_0,x)`$, for every $`t𝐑`$ and $`xE`$). Let us denote by $`N`$ the maximal positive invariant set in $`E`$ , i.e. for any $`x_0N`$ and $`t_0𝐑`$, $`x(t;t_0,x_0)N`$, for every $`t[t_0,t_0+T_{x_0})`$ in the maximal interval of definition of the solution.
Then the dynamics $`(\text{1})`$ has at $`\overline{x}=0`$ an uniformly asymptotic stable equilibrium point if and only if the dynamics restricted to $`N`$ has an asymptotic stable equilibrium at $`\overline{x}=0`$. $`\mathrm{}`$
Even if it has appeared in the literature in a more general setting (I refer to ), it is worth mentioning the form the Invariance Principle takes in this context:
###### THEOREM 2 (Invariance Principle)
Consider the time-varying dynamical system $`(\text{1})`$ for which $`f`$ has the properties $`1)4)`$. Suppose there exists a function $`V:D𝐑`$ of class $`𝒞^1`$ such that:
$`H1)`$ It is bounded below, i.e. $`V(x)V_0`$ for any $`xD`$ for some $`V_0𝐑`$;
$`H2)`$ There exists a continuous function $`W:D𝐑`$ such that:
$$\frac{dV}{dt}(t,x)=V(x)f(t,x)W(x)0$$
$`H3)`$ Let $`E=\{xD|W(x)=0\}`$ denote the zero-set (or kernel) of $`W`$; suppose that $`f`$ resticted to $`E`$ is time-invariant (i.e. $`f(t,x)=f(t_0,x)`$, for any $`t𝐑`$ and $`xE`$). Let us denote by $`N`$ the maximal positive invariant set included in $`E`$, i.e. for any $`x_0N`$ and $`t_0𝐑`$, $`x(t;t_0,x_0)N`$, for any $`t[t_0,t_0+T_{x_0})`$ in the maximal interval of definition of the solution.
Then any bounded trajectory of $`(\text{1})`$ tends to $`N`$, i.e. if $`(t_0,x_0)`$ is the initial data for a bounded solution included in $`D`$ then:
$$\underset{t\mathrm{}}{lim}d(x(t;t_0,x_0),N)=0$$
(3)
Remarks 1) There are two directions in which Theorem 1 generalizes the well-known Barbashin-Krasovskii-LaSalle’s Theorem (see , or ); firstly we require $`V`$ to be only nonnegative and not strictly positive, secondly we consider the time-varying dynamical systems. There exists in literature two earlier results in the first direction that I wish to comment. The first result that I am referring to is Lemma 5 from . In this lemma only autononous systems are considered and the restricted dynamics is required to be attractive in the sense that all trajectories should tend to the origin. I point out that only the requirement of attractivity is not enough; this can be seen in a trivial case, namely the 2 dimensional system given by Vinograd (conform ), for which the origin is an attractive equilibrium but not stable, and take $`V0`$. I stress out that for the purposes of their paper () Lemma 5 can be replaced by Theorem 7 of this paper without affecting the other results of that paper.
A second result has appeared in but not in a general and explicit form as here. In fact in the author is concerned with the stability of the large-scale systems which are already decomposed in triangular form. Thus, this result solves the problem only in the case when we can perform the observability decomposition of the dynamics (1) with respect to the output $`W(x)`$. This case requires a supplementary condition, namely the codistribution span by $`dW,dL_fW,\mathrm{},dL_f^nW`$ to be of constant rank on $`D`$ (see ). Among other requirements, this geometric condition implies also that $`N`$ is a manifold, whereas we do not assume here this rather strong assumption.
I acknowledge the existence of a recently published paper that deals with a similar extension of the Liapunov theorem, yet only for autonomous systems (). However, I was unaware of this result at the time I was working in this field (i.e. 1993–1995).
2) Some other papers deal with extensions of the invariance principle for nonautonomous systems. In two special cases, when the system is either asymptotically autonomous (in ) or asymptotically almost periodic (in ), the bounded solution tends to the largest pseudo-invariant set in $`E`$. However they use the classical Liapunov theorems to obtain the uniform boundedness of the solutions. Thus they require the existence of a strictly positive definite function playing the rôle of Liapunov function, while here we require only nonnegativeness of the Liapunov-like function. In other approaches an additional auxiliary function is assumed and by means of extra conditions the time in $`E`$ is controlled (see the results of Salvadori or Matrosov, e.g. in ). In a third approach an extra condition on $`\dot{V}`$ is considered without any additional condition on the vector field; such an approach is considered in .
3) The condition that the restricted dynamics to be uniformly asymptoticaly stable is necessary and sufficient . Thus it is a center-manifold-type result where a knowledge about a restricted dynamics to some invariant set implies the same property of the whole dynamics. We point out here that the set $`N`$ does not need to be a manifold.
4) One could expect that simple stability of the restricted dynamics would imply uniform stability of the restricted dynamics. But this is not true as we can see from the following example:
###### EXAMPLE 1
Consider the following autonomous planar system:
$$\{\begin{array}{c}\dot{x}=y^2\\ \dot{y}=y^3\end{array},(x,y)𝐑^2$$
(4)
The solution of the system is given by $`(x,y)(x+ln(1+y^2t),\frac{y}{\sqrt{1+y^2t}})`$. It is obvious that the equilibrium is not stable but if we take $`V=y^2`$ we have $`\frac{dV}{dt}=2y^4`$ and on the set $`E=N=\{(x,0),x𝐑\}`$ the dynamics is trivial stable $`\dot{x}=0`$.
The problem is not the nonisolation of the equilibrium, but the existence of some invariant sets in any neighborhood of the equilibrium;
5) Theorem 2 is the natural generalization of the Invariance Principle to the class of systems considered in this paper. The conclusion of this theorem applies only to bounded trajectories. Thus we have to know apriori which solutions are bounded. Since they are bounded we can extend them indefinitely in positive time. Thus it makes sense to take the limit $`t\mathrm{}`$ in $`(\text{3})`$. We mention that a more general Invariance Principle can be obtained even under weaker conditions than those from here (see ).
The organization of the paper is the following: in the next section we give the proof of these results. In the third section we consider the autonomous case and we present the systemic consequences related to the nonlinear Liapunov equation and a special type of zero-state detectability. In the fourth section we consider a nonlinear Riccati equation (or Hamilton-Jacobi equation) and we present a result of robust stabilizability by output feedback. The last section contains the conclusions and is followed by the bibliography.
## 2 Proof of the Main Results
We prove by contradiction the uniform stability of the equilibrium. For this, we construct a $`𝒞^1`$-convergent sequence of solutions that are going away from the origin and whose limit is a trajectory, thus contradicting the hypothesis.
For the uniform asymptotic stability, we prove first that the $`\omega `$-limit set of bounded trajectories is included in $`N`$ (implicitely proving the Invariance Principle - Theorem 2) and then we addapt a classical trick (used for instance in Theorem 34.2 from ) that the convergence of trajectories in $`\omega `$-limit set will attract the convergence of the bounded trajectory itself. In both steps we use essentially the time-invariant property of $`f`$ restricted to $`E`$. In proving the uniform stability we also obtain that the solution can be defined on the whole positive real set (can be completely extended in future).
Theorem 2 (the Invariance Principle) will follow simply from a lemma that we state during the proof of uniform attractivity.
First we need a lemma:
###### LEMMA 3
Let $`f`$ be a vector field defined on a domain $`D`$ and having the properties 1-4 as above. Let $`(t_i)_i`$ be a sequence of real numbers and $`(w_i)_i`$, $`w_i:[a,b]D`$ be a sequence of trajectories for the time-translated vector field $`f`$ with $`t_i`$, i.e. $`\dot{w}_i(t)=f(t+t_i,w_i(t))`$ .
If the trajectories are uniformly bounded, i.e. there exists $`M>0`$ such that $`w_i_{\mathrm{}}<M`$, for any $`i`$, then we can extract a subsequence, denoted also by $`(w_i)_i`$, uniformly convergent to a function $`w`$ in $`𝒞^1([a,b];D)`$, i.e. $`w_iw`$ and $`\dot{w}_i\dot{w}`$ both uniformly in $`𝒞^0([a,b];D)`$
Proof We apply the Ascoli-Arzelà Lemma twice: first to extract a subsequence such that $`(w_i)_i`$ is uniformly convergent and second to extract further another subsequence such that $`(\dot{w}_i)_i`$ is uniformly convergent. Then we obtain that $`lim_i\frac{d}{dt}w_i=\frac{d}{dt}lim_iw_i`$.
1. We verify that $`(w_i)_i`$ are uniformly bounded and equicontinuous. The uniformly boundedness comes from $`w_i_{\mathrm{}}<M`$. The equicontinuity comes from the uniformly boundedness of the first derivative. Indeed, since $`w_iM`$, the closed ball $`\overline{B}_M`$ is compact and $`f(t,)`$ is continuous on $`\overline{B}_M`$, there exists a constant $`A`$ such that $`f(t,x)A`$, for any $`(t,x)𝐑\times \overline{B}_M`$. Then:
$$\dot{w}_i(t)=f(t+t_i,w_i(t))A,\mathrm{for}\mathrm{any}i\mathrm{and}t[a,b]$$
Thus $`(w_i)_i`$ is relatively compact and we can extract a subsequence, that we denote also by $`(w_i)_i`$, which is uniformly convergent to a function $`w𝒞^0([a,b];D)`$.
2. We prove that $`(\dot{w}_i)_i`$ is relatively compact. We have already proved the uniform boundedness $`\dot{w}_i_{\mathrm{}}A`$. For the equicontinuity we use both the uniform continuity in $`t`$ and uniform local Lipschitz continuity in $`x`$, of $`f`$. Let $`L_M`$ be the uniform Lipschitz constant corresponding to the compact set $`\overline{B}_M`$. Then:
$$\dot{w}_i(t_1)\dot{w}_i(t_2)=f(t_i+t_1,w_i(t_1))f(t_i+t_2,w_i(t_2))$$
$$f(t_i+t_1,w_i(t_1))f(t_i+t_2,w_i(t_1))$$
$$+f(t_i+t_2,w_i(t_1))f(t_i+t_2,w_i(t_2))$$
Let $`\epsilon >0`$ be arbitrarily. Then we choose $`\delta _1`$ such that $`f(s_1,x)f(s_2,x)<\frac{\epsilon }{2},\mathrm{for}\mathrm{any}|s_1s_2|<\delta _1\mathrm{and}x\overline{B}_M`$ On the other hand: $`f(t_i+t_2,w_i(t_1))t(t_i+t_2,w_i(t_2))L_Mw_i(t_1)w_i(t_2)L_MA|t_1t_2|`$. Then we choose $`\delta =\mathrm{min}(\delta _1,\frac{\epsilon }{2L_MA})`$. Then the left-hand side from the above inequality is also bounded by $`\frac{\epsilon }{2}`$ for any $`t_1,t_2`$ with $`|t_1t_2|<\delta `$. Thus $`\dot{w}_i(t_1)\dot{w}_i(t_2)<\frac{\epsilon }{2}+\frac{\epsilon }{2}=\epsilon `$, for any $`i`$ and $`t_1,t_2[a,b]`$, $`|t_1t_2|<\delta `$.
We can now extract a second subsequence from $`(w_i)_i`$ such that $`(\dot{w}_i)_i`$ is also uniformly convergent and this ends the proof of lemma. $`\mathrm{}`$
Proof of Uniform Stability
Let us assume that the equilibrium is not uniformly stable. Then there exists $`\epsilon _0>0`$ such that for any $`\delta `$, $`0<\delta <\epsilon _0`$ there are $`x_0,t`$ and $`\mathrm{\Delta }>0`$ such that $`x_0<\delta `$ and $`x(t+\mathrm{\Delta };t,x_0)=\epsilon _0`$, $`x(t+\tau ;t,x_0)<\epsilon _0`$, for $`0\tau <\mathrm{\Delta }`$. We choose $`\epsilon _0`$ (eventually by shrinking it) such that $`\overline{B}_{\epsilon _0}N`$ is included in the attraction domain of the origin (for the restricted dynamics).
By choosing a sequence $`(\delta _i)_i`$ converging to zero we obtain sequences $`(x_{0i})_i`$, $`(t_i)_i`$ and $`(\mathrm{\Delta }_i)_i`$ such that: $`x_{0i}0`$ and $`x(t_i+\mathrm{\Delta }_i;t_i,x_{0i})=\epsilon _0`$
Let $`\delta <\epsilon _0`$ be such that for any $`z_0B_\delta N`$ we have $`x(t;0,z_0)<\frac{\epsilon _0}{2}`$ for any $`t>0`$ (such a choice for $`\delta `$ is possible since the dynamics restricted to $`N`$ is stable). Let $`i_0`$ be such that $`\delta _i<\delta `$, for $`i>i_0`$. We denote by $`(u_i)_{i>i_0}`$ the time moments such that $`x(t_i+u_i;t_i,x_{0i})=\delta `$ and $`x(t;t_i,x_{0i})>\delta `$ for $`t>t_i+u_i`$. Since the spheres $`\overline{S}_{\epsilon _0}`$ and $`\overline{S}_\delta `$ are compact we can extract a subsequence (indexed also by $`i`$) such that both $`x_i=x(t_i+\mathrm{\Delta }_i;t_i,x_{0i})`$ and $`y_i=x(t_i+u_i;t_i,x_{0i})`$ are convergent to $`x^{}`$, respectively to $`y^{}`$; $`x_ix^{}`$, $`y_iy^{}`$, $`x^{}=\epsilon _0`$, $`y^{}=\delta `$. Since $`V`$ is continuously nonincreasing on trajectories and $`lim_iV(x_{0i})=0`$, we get $`V(x^{})=V(y^{})=0`$. Therfore $`x^{},y^{}N`$.
Suppose $`f(x,t)A`$ on $`\overline{B}_{\epsilon _0}`$, for some $`A>0`$. Then one can easily prove that $`\mathrm{\Delta }_iu_i\frac{\epsilon _0\delta }{A}=T_1`$, for any $`i>i_0`$ (i.e. the flight time between two spheres of radius $`\delta `$ and $`\epsilon _0`$ has a lower bound).
Define now the time-translated vector fields $`f_i(t,x)=f(t+t_i+u_i,x)`$ and denote by $`w_i:[0,T_1]\overline{B}_{\epsilon _0}`$ the time-translated solutions $`w_i(t)=x(t+t_i+u_i;t_i,x_{0i})`$. Then: $`\dot{w}_i(t)=f_i(t,w_i(t))`$, $`0tT_1`$. By applying Lemma 3 we get a subsequence uniformly convergent to a trajectory $`w^1:[0,T_1]\overline{B}_{\epsilon _0}N`$, such that $`w^1(0)=lim_iw_i(0)=y^{}`$ and $`w^1(t)>\delta `$, for $`0<tT_1`$. If $`w^1(T_1)<\epsilon _0`$ we obtain that $`\mathrm{\Delta }_iu_iT_1>\frac{\epsilon _0w^1(T_1)}{A}`$, for some $`ii_1>i_0`$. Then, we denote $`T_2=T_1+\frac{\epsilon _0w^1(T_1)}{A}`$ and we repeat the scheme. We obtain another sequence which is uniformly convergent to a trajectory $`w^2:[0,T_2]\overline{B}_{\epsilon _0}N`$ such that $`w^2(0)=y^{}`$, $`w^2(t)>\delta `$, $`0<tT_2`$ and $`w^2(t)=w^1(t)`$, for $`0tT_1`$.
Thus we extend each trajectory $`w^k:[0,T_k]\overline{B}_{\epsilon _0}N`$ to a trajectory $`w^{k+1}:[0,T_{k+1}]\overline{B}_{\epsilon _0}N`$ such that $`T_{k+1}T_k`$, $`w^{k+1}(t)=w^k(t)`$ for $`0tT_k`$ and $`w^{k+1}(t)>\delta `$, for $`0<tT_{k+1}`$.
We end this sequence of extensions in two cases:
1) $`lim_kT_k=T^{}<+\mathrm{}`$ (the limit may be reached in a finite number of steps), in which case we have$`lim_kw^k(T_k)=\epsilon _0`$ and thus $`lim_kw^k(T_k)=x^{}`$; or:
2) $`lim_kT_k=+\mathrm{}`$.
In the first case we obtain a trajectory $`w^{}:[0,T^{}]\overline{B}_{\epsilon _0}N`$ such that $`w^{}(0)=y^{}`$, $`w^{}(T^{})=x^{}`$ with $`w^{}(0)=\delta `$ and $`w^{}(T^{})=\epsilon _0`$. But this is a contradiction with the choice of $`\delta `$ (and of stability of the restricted dynamics).
In the second case we obtain a trajectory $`w^{}:[0,\mathrm{})\overline{B}_{\epsilon _0}N`$ such that $`w^{}(0)=\delta <\epsilon _0`$ and $`w^{}(t)>\delta `$ for $`t>0`$. Thus $`lim_t\mathrm{}w^{}(t)0`$ contradicting the assumption that $`\overline{B}_{\epsilon _0}N`$ is included in the attraction domain of the origin. Now the proof is complete. $`\mathrm{}`$.
For the proof of uniformly attractivity we recall a few definitions and results:
Definition A point $`x^{}`$ is called $`\omega `$-limit point for the trajectrory $`x(t;t_0,x_0)`$ if there exists a sequence $`(t_k)_k`$ such that $`lim_k\mathrm{}t_k=\mathrm{}`$, $`x(t;t_0,x_0)`$ is defined for all $`t>t_0`$ and $`lim_kx(t_k;t_0,x_0)=x^{}`$.
The set of all $`\omega `$-limit points is called the $`\omega `$-limit set and is denoted by $`\mathrm{\Omega }(t_0,x_0)`$. It characterizes the trajectory $`x(t;t_0,x_0)`$ and it depends on the initial data $`(t_0,x_0)`$.
###### THEOREM 4 (Birkoff’s Limit Set Theorem, see )
A bounded trajectory approaches its $`\omega `$-limit set, i.e. $`lim_t\mathrm{}d(x(t;t_0,x_0),\mathrm{\Omega }(t_0,x_0))=0`$, where $`d(p,S)=inf_{xS}px`$ is the distance between the point $`p`$ and the set $`S`$. $`\mathrm{}`$
There is also a very useful result about uniformly continuous functions:
###### LEMMA 5 (Barbălat’s Lemma, see )
If $`g:[t_0,\mathrm{})\mathrm{}`$ is a uniformly continuous function such that the following limit exists and is finite, $`lim_t\mathrm{}_{t_0}^tg(\tau )𝑑\tau `$, then $`lim_t\mathrm{}g(t)=0`$. $`\mathrm{}`$
Proof of Uniform Attractivity
We already know that $`\overline{x}=0`$ is uniformly stable. What we have to prove is the uniform attractivity.
Let $`\epsilon _0>0`$ be chosen with the following properties: for any $`t_0`$ and $`x_0D\overline{B}_{\epsilon _0}`$ the positive trajectory $`x(t;t_0,x_0)`$ is bounded by $`\epsilon _1`$ (i.e. $`x(t;t_0,x_0)B_{\epsilon _1}`$); for any $`t_1`$ and $`x_1DB_{\epsilon _1}`$ the trajectory $`x(t;t_1,x_1)`$, $`t>t_1`$, is bounded by some $`M`$; and for any $`x_2NB_{\epsilon _1}`$ the trajectory $`x(t;t_0,x_2)`$ tends to the origin $`lim_t\mathrm{}x(t;t_0,x_2)=0`$. We are going to prove that $`lim_t\mathrm{}x(t;t_0,x_0)=0`$.
Let us consider the $`\omega `$-limit set $`\mathrm{\Omega }(t_0,x_0)`$. It is enough to prove that $`\mathrm{\Omega }(t_0,x_0)=\{0\}`$, because of Birkoff’s Limit Set Theorem.
Let $`x^{}\mathrm{\Omega }(t_0,x_0)`$ and suppose $`x^{}0`$. Let us denote by $`x(t)=x(t;t_0,x_0)`$ and $`g(t)=V(x(t))f(t,x(t))`$. Since the solution is continuous and bounded, so is $`g(t)`$. On the other hand
$$V(x(t))=V(x_0)+_{t_0}^tg(\tau )𝑑\tau $$
Since $`\dot{x}(t)=f(t,x(t))`$ and $`x(t)`$ is bounded we obtain that it is also uniformly continuous. Thus $`g(t)`$ is also uniformly continuous (recall we have assumed $`f(,x)`$ is uniformly continuous in $`t`$). Let $`(t_k)_k`$ be a sequence that renders $`x^{}`$ a $`\omega `$-limit point. Then $`lim_kV(x(t_k))=V(lim_kx(t_k))=V(x^{})`$. Since $`V(x(t))`$ is a decreasing function bounded below, there exists the limit: $`lim_t\mathrm{}V(x(t))=V(x^{})`$. Now, applying Barbălat’s Lemma we obtain $`lim_t\mathrm{}g(t)=0`$ or $`W(x^{})=0`$. Thus $`\mathrm{\Omega }(t_0,x_0)E`$, the kernel of $`W`$.
In this point we need a result about the behaviour of solutions starting at $`x^{}`$. We mention that the following lemma is a consequence of Theorem 3 from . But, since we are under stronger conditions, we have found a simpler proof that we are going to present here (our conditions are stronger because we need to obtain uniform stability and consequently boundedness of the solutions when Liapunov function is only positive semidefinite, which overall means a weaker condition).
###### LEMMA 6
The positive trajectory starting at $`x^{}`$ is included in $`E`$ and thus the $`\mathrm{\Omega }`$-limit set is a positive invariant set included in $`N`$.
Proof
Let $`\tau >0`$ be an arbitrary time interval. Let $`(t_k)_k`$ be the sequence that renders $`x^{}`$ a $`\omega `$-limit point for the trajectory $`x(t)=x(t;t_0,x_0)`$. Then, if we denote by $`x_k=x(t_k)`$ we have $`lim_kx_k=x^{}`$. Consider the following sequence of functions: $`w_k:[0,\tau ]D,w_k(t)=x(t+t_k;t_k,x^{})`$. We have chosen $`x_0,t_0`$ such that all these functions are bounded by $`M`$, i.e. $`w_k_{\mathrm{}}<M`$. We have $`w_k(0)=x^{}`$ and $`V(w_k(t))V(x^{})`$. Let us denote by $`y_k^t=x(t+t_k)`$, for any $`0t\tau `$, and let $`L_M`$ be the Lipschitz constant of $`f`$ on the compact $`\overline{B}_M`$. Then: $`y_k^tw_k(t)e^{L_mt}x_kx^{}`$ and, since $`lim_kx_k=x^{}`$ we get $`lim_ky_k^tw_k(t)=0`$. On a hand, since $`V(x^{})=lim_t\mathrm{}V(x(t))`$ and $`V`$ is nonincreasing on trajectories we have $`V(y_k^t)>V(x^{})`$ and also $`lim_kV(y_k^t)=V(x^{})=lim_kV(w_k(t))`$. On the other hand, since $`(w_k)_k`$ are uniformly bounded we apply Lemma 3 and we obtain a subsequence uniformly convergent to a function $`w𝒞^1([0,\tau ];D\overline{B}_M)`$. Obviously $`V(w(t))=V(x^{})`$ for any $`0t\tau `$. Thus $`W(w(t))=0`$ and $`w(t)E`$. On the other hand, since $`f`$ is continuous in $`(t,x)`$ we obtain that $`w`$ is an integral curve of $`f`$, i.e. $`\dot{w}(t)=f(t_{},w(t))`$, for $`0t\tau `$ and any $`t_{}`$. In particular, for $`t_{}=t_k`$ we get $`w(t)`$ is a solution of the same equation as $`w_k(t)`$ and $`w(0)=w_k(0)=x^{}`$. By the uniqueness of the solution they must coincide. Then $`x(t+t_k;t_k,x^{})E`$ for $`0t\tau `$. But $`\tau `$ was arbitrarily; thus $`x(t;t_0,x^{})E`$ for any $`t`$ and then $`x^{}N`$. $`\mathrm{}`$
Since the trajectory starting at $`x^{}`$ is included in $`N`$, it should converge to the origin (the equilibrium point). Let us denote by $`\epsilon =\frac{x^{}}{2}`$. From uniform stability there exists a $`\delta >0`$ such that for any $`\stackrel{~}{x}D`$, $`\stackrel{~}{x}<\delta `$ implies $`x(t_2;t_1,\stackrel{~}{x})<\epsilon `$, for any $`t_2>t_1`$. Let $`\mathrm{\Delta }t`$ be a time interval such that $`x(t;0,x^{})<\frac{\delta }{2}`$ for any $`t>\mathrm{\Delta }t`$. We consider the compact set $`C`$, the $`\frac{\delta }{2}`$-neighborhood of the compact curve $`\mathrm{\Gamma }=\{x(t;0,x^{})|0t\mathrm{\Delta }t\}`$:
$$C=\{xD|d(x,\mathrm{\Gamma })\frac{\delta }{2}\}=\underset{t[0,\mathrm{\Delta }t]}{}\overline{B_{\delta /2}(x(t;0,x^{}))}$$
which is the union of the closed balls centered at $`x(t;0,x^{})`$ and of radius $`\frac{\delta }{2}`$. We set $`\delta _1=\frac{\delta }{2}exp(L_C\mathrm{\Delta }t)`$ where $`L_C`$ is the uniform Lipschitz constant of $`f`$ on the compact set $`C`$. Since the solution is uniformly Lipschitz with respect to the initial point $`x_0`$ we have that for any $`t_1𝐑`$ and $`x_1`$ such that $`x_1x^{}<\delta _1`$ we get: $`x(t_1+\mathrm{\Delta }t;t_1,x_1)x(\mathrm{\Delta }t;0,x^{})<\frac{\delta }{2}`$ and then $`x(t_1+\mathrm{\Delta }t;t_1,x_1)<\delta `$. Furthermore, from the choice of $`\delta `$ we obtain that $`x(t_1+\tau ;t_1,x_1)<\epsilon `$, for any $`\tau >\mathrm{\Delta }t`$ or $`x(t_1+\tau ;t_1,x_1)x^{}>\epsilon `$, for any $`\tau >\mathrm{\Delta }t`$.
Now we pick a $`t_n`$ such that $`x(t_n;t_0,x_0)x^{}<\delta _1`$. Then, from the previous discussion $`x(t_n+\tau ;t_0,x_0)x^{}>\epsilon `$, for any $`\tau >\mathrm{\Delta }t`$ which contradicts the limit $`lim_kx(t_k;t_0,x_0)=x^{}`$. This contradiction comes from the hypothesis that $`x^{}0`$. Thus $`\mathrm{\Omega }(t_0,x_0)=\{0\}`$ and now the proof is complete. $`\mathrm{}`$
Proof of Theorem 2 (The Invariance Principle)
If $`x(t;t_0,x_0)`$ is a bounded trajectory then, from Birkoff’s Limit Set Theorem it approaches its $`\omega `$-limit set. On the one hand we can use Barbălat’s Lemma and prove that $`W`$ vanishes on $`\omega `$-limit set of bounded trajectories. On the other hand, as we have proved in Lemma 6, the $`\omega `$-limit set is invariant and included in $`N`$. Thus the bounded trajectory approaches the set $`N`$. $`\mathrm{}`$
## 3 The Autonomous Case: Consequences in Nonlinear Control Theory
Consider the following inputless nonlinear control system:
$$S\{\begin{array}{c}\dot{x}=f(x)\\ y=h(x)\end{array},xD𝐑^ny𝐑^p$$
(5)
such that $`f(0)=0`$, $`h(0)=0`$ and $`D`$ a neighborhood of the origin. Suppose $`f`$ is local Lipschitz continuous and $`h`$ continuous on $`D`$. Then denote by $`x(t,x_0)`$ the flow generated by $`f`$ on $`D`$ (i.e. the solution of $`\dot{x}=f(x)`$, $`x(0)=x_0`$), by $`E=kerh=\{xD|h(x)=0\}`$, the kernel of $`h`$ and by $`N`$ the maximal positive invariant set included in $`E`$, i.e. the set $`N=\{\stackrel{~}{x}D|h(x(t,\stackrel{~}{x}))=0\mathrm{for}\mathrm{any}t0\mathrm{such}\mathrm{that}x(t,\stackrel{~}{x})\mathrm{has}\mathrm{sense}\}`$.
We present two concepts of detectability for $`(\text{5})`$. The first one has been used by many authors lately (see for instance ).
Definition The pair $`(h,f)`$ is called zero-state detectable (or z.s.d.) if $`\overline{x}=0`$ is an attractive point for the dynamics restricted to $`N`$, i.e. there exists an $`\epsilon _0>0`$ such that for any $`x_0N`$, $`x_0<\epsilon _0`$, $`lim_t\mathrm{}x(t,x_0)=0`$.
Definition The pair $`(h,f)`$ is called strong zero-state detectable (or strong z.s.d.) if $`\overline{x}=0`$ is an asymptotical stable equilibrium point for the dynamics restricted to $`N`$, i.e. it is zero-state detectable and for some $`\epsilon _0`$ and for any $`x_0N`$ with $`x_0<\epsilon _0`$, $`lim_t\mathrm{}x(t,x_0)=0`$.
We see that strong z.s.d. implies z.s.d., but obviously the converse is not true.
In this framework, as a consequence of the main result we can state the following theorem:
###### THEOREM 7
For the inputless nonlinear control system $`(\text{5})`$ with $`f`$ local Lipschitz continuous and $`h`$ continuous, consider the following nonlinear Liapunov equation:
$$Vf+h^q=0$$
(6)
or the following nonlinear Liapunov inequality:
$$Vf+h^q0$$
(7)
for some $`q>0`$. Suppose there exists a positive semidefinite solution of $`(\text{6})`$ or $`(\text{7})`$ of class $`𝒞^1`$ defined on $`D`$ such that $`V(0)=0`$.
Then the pair $`(h,f)`$ is strong zero-state detectable if and only if $`\overline{x}=0`$ is an asymptoticaly stable equilibrium for the dynamics $`(\text{5})`$. $`\mathrm{}`$
Below we give an example:
###### EXAMPLE 2
Consider the dynamics:
$$\begin{array}{c}\dot{x}_1=x_1^3+\mathrm{\Psi }(x_2)\hfill \\ \dot{x}_2=x_2^3\hfill \end{array},(x_1,x_2)𝐑^2$$
(8)
where $`\mathrm{\Psi }:𝐑𝐑`$ is local Lipschitz continuous, $`\mathrm{\Psi }(0)=0`$ and there exist constants $`a>0`$, $`b1`$ such that:
$$|\mathrm{\Psi }(x)|a|x|^b,x_2$$
If we choose as output function $`h(x)=x_2^2`$ we see that the pair $`(h,f)`$ is strong zero-state detectable; indeed, the set $`E=\{x𝐑^2|h(x)=0\}=\{(x_1,0)|x_1𝐑\}`$ and the dynamics restricted to $`E`$ is $`\dot{x}_1=x_1^3`$ which is asymptoticaly stable.
Now, if we choose $`V(x)=\frac{x_2^2}{2}`$ we have $`\dot{V}=x_2^4`$ and thus $`V`$ is a solution of the Liapunov equation $`(\text{6})`$ with $`q=2`$. Then, the equilibrium is asymptoticaly stable, as a consequence of the theorem 7.
On the other hand we can explicitely solve for $`x_2`$: $`x_2(t)=\frac{x_{20}}{\sqrt{2(1+x_{20}^2t)}}`$ and then we have: $`|\mathrm{\Psi }(x_2(t))|C(1+Bt)^{1/2}`$ for some $`B,C>0`$ and any $`t0`$. Now the asymptotic stability follows as a consequence of Theorem 68.2 from (stability under perturbation).
## 4 An Application in Robust Stabilizability
We present here, as an application, a robust stabilizability result for a nonlinear affine control system. In fact it is an absolute stability result about a particular situation. More general results about absolute stability for nonlinear affine control system will appear in a forthcoming paper. We base our approach on the existance of a positive semidefinite solution of some Hamilton-Jacobi equation or inequality. Discussions about solutions of this type of equation may be found in .
Consider the following Single Input - Single Output control system:
$$\{\begin{array}{ccc}\hfill \dot{x}& =& f(x)+g(x)u\hfill \\ \hfill y& =& h(x)\hfill \end{array},xD𝐑^n,u,y𝐑$$
(9)
where $`f`$ and $`g`$ are local Lipschitz continuous vector fields on a domain $`D`$ including the origin, $`h`$ is a local Lipschitz real-valued function on $`D`$, and $`f(0)=0`$, $`h(0)=0`$. Consider also a local Lipschitz output feedback:
$$\phi :𝐑𝐑,\phi (0)=0$$
(10)
We define now two classes of perturbations associated to this feedback. Let $`a>0`$ be a positive real number. The first class contains time-invariant perturbations:
$$P_1=\{p:𝐑𝐑,p\mathrm{is}\mathrm{local}\mathrm{Lipschitz},p(0)=0\mathrm{and}|p(y)|<a|\phi (y)|,y0\}$$
while the second class is composed by time-varying perturbations:
$$P_2=\{p:𝐑\times 𝐑𝐑,p(y,t)\mathrm{is}\mathrm{local}\mathrm{Lipschitz}\mathrm{in}y\mathrm{for}t\mathrm{fixed}\mathrm{and}\mathrm{uniformly}\mathrm{continuous}\mathrm{in}t$$
$$\mathrm{for}\mathrm{any}y\mathrm{fixed},p(0,t)0\mathrm{and}\mathrm{there}\mathrm{exists}\epsilon >0\mathrm{such}\mathrm{that}|p(y,t)|<(a\epsilon )|\phi (y)|,y0,t\}$$
Now we can define more precisely the concept of robust stability:
Definition We say the feedback $`(\text{10})`$ robustly stabilizes the system $`(\text{9})`$ with respect to the class $`P_1P_2`$ if for any perturbation $`pP_1P_2`$ the closed-loop with the perturbed feedback $`\phi +p`$ has an asymptoticaly stable equilibrium at the origin.
In other words, we require that the origin to be asymptoticaly stable for the dynamics:
$$\dot{x}=f(x)+g(x)(\phi (h(x))+p(h(x),t))$$
(11)
for any $`pP_1P_2`$. Since the null function belongs to $`P_1`$, the feedback $`\phi `$ itself must stabilize the closed-loop too.
With these preparations we can state the result:
###### THEOREM 8
Consider the nonlinear affine control system $`(\text{9})`$ and the feedback $`(\text{10})`$. Suppose the pair $`(h,f)`$ is strong zero-state detectable and suppose the following Hamilton-Jacobi equation:
$$Vf+(\frac{1}{2}Vg+\phi h)^2(1a^2)(\phi h)^2=0,V(0)=0$$
(12)
or inequality:
$$Vf+(\frac{1}{2}Vg+\phi h)^2(1a^2)(\phi h)^20,V(0)=0$$
(13)
has a positive semidefinite solution $`V`$ of class $`𝒞^1`$ on $`D`$.
Then the feedback $`\phi `$ robustly stabilizes the system $`(\text{9})`$ with respect to the class $`P_1P_2`$.
Proof
Let us consider a perturbation $`pP_1P_2`$. Then, the closed-loop dynamics is given by $`(\text{11})`$. We compute the time derivative of the solution $`V`$ of $`(\text{12})`$ with respect to this dynamics:
$$\frac{dV}{dt}=Vf(x)+Vg(x)(\phi (h(x))+p(h(x),t))$$
After a few algebraic manipulations we get:
$$\frac{dV}{dt}(\frac{1}{2}Vgph)^2+(ph)^2a^2(\phi h)^2$$
Now, for $`pP_1`$, $`\frac{dV}{dt}`$ is time-independent and we may take for instance:
$$W(x)=(p(h(x)))^2a^2(\phi (h(x)))^20$$
For $`pP_2`$, $`\frac{dV}{dt}`$ is time-dependent and we define:
$$W(x)=(2a\epsilon \epsilon ^2)(\phi (h(x)))^20$$
Either a case or the other, we obtain (recall the definitions of $`P_1`$ and $`P_2`$):
$$\frac{dV}{dt}W(x)0$$
The kernel-set of $`W`$ is given by:
$$E=\{xD|W(x)=0\}=\{xD|h(x)=0\}$$
We see that the closed-loop dynamics $`(\text{11})`$ restricted to $`E`$ is simply given by $`\dot{x}=f(x)`$ and is time-independent. Moreover, since we have supposed $`(h,f)`$ is strong zero-state detectable, it follows that the restricted dynamics to the maximal positive invariant set in $`E`$ has an asymptoticaly stable equilibrium at the origin. Now, applying Theorem 1, the result follows. $`\mathrm{}`$
Let us consider now an example:
###### EXAMPLE 3
Consider the following planar nonlinear control system:
$$\{\begin{array}{c}\dot{x}_1=x_1^3+u\hfill \\ \dot{x}_2=x_2^3\hfill \\ y=x_2^3\hfill \end{array}$$
(14)
We are interested to find how robust the feedback $`\phi (y)=y`$ is, i.e. how large we can choose $`a`$ such that $`\phi `$ robustly stabilizes the system $`(\text{14})`$ with respect to the class $`P_1P_2`$.
The Hamilton-Jacobi equation $`(\text{12})`$ takes the form:
$$x_1^3\frac{V}{x_1}x_2^3\frac{V}{x_2}+(\frac{1}{2}\frac{V}{x_1}+x_2^3)^2(1a^2)x_2^6=0$$
or:
$$x_1^3\frac{V}{x_1}x_2^3\frac{V}{x_2}+\frac{1}{4}(\frac{V}{x_1})^2+x_2\frac{V}{x_1}+a^2x_2^6=0$$
A solution of this equation is:
$$V(x_1,x_2)=\frac{a^2}{4}x_2^4$$
For any $`a>0`$ it is positive semidefinite and the system $`(\text{14})`$ is strong zero-state detectable. Thus, as a consequence of theorem 8, we can choose $`a`$ arbitrary large such that $`\phi `$ robustly stabilizes the system $`(\text{14})`$ with respect to the class $`P_1P_2`$.
On the other hand, for any feedback $`\mathrm{\Phi }`$, local Lipschitz and:
$$|\mathrm{\Phi }(y)|a|y|,\mathrm{for}\mathrm{some}a>0$$
we have seen in the previous example that the closed-loop has an asymptoticaly stable equilibrium at the origin.
## 5 Conclusions
In this paper we study an extension of Barbashin-Krasovski-LaSalle and Invariance Principle to a class of time-varying dynamical systems. We impose two type of conditions on the vector field: one is regularity (we require uniformly continuity with respect to $`t`$ and uniformly local Lipschitz continuoity and boundedness with respect to $`x`$); the other condition requires the vector field to be time-invariant on the zero-set $`E`$ of an auxiliary function. In this setting we find that the asymptotic behaviour of the dynamics restricted to the largest positive invariant set in $`E`$ determines the asymptotic stability character of the full dynamics.
Then we study two applications in control theory. The first application concerns the notion of detectability. We give another definition for this notion, called strong zero-state detectability and we show how the existence of a positive semidefinite solution of the Liapunov equation or inequation is related to the asymptotic stability of the equilibrium. We obtain a nonlinear equivalent of the linear well-known result: if the pair $`(C,A)`$ is detectable and there exists a positive solution $`P0`$ of the Liapunov algebraic equation $`A^TP+PA+C^TC=0`$, then the matrix $`A`$ has all eigenvalues with negative real part.
The second application is on the problem of robust stabilizability. We give sufficient conditions such that a given feedback robustly stabilizes the closed-loop with respect to two sector classes of perturbations (time-invariant and time-varying). The condition is formulated in term of the existence of a positive solution of some Hamilton-Jacobi equation or inequality.
This last application opens the problem of absolute stability for nonlinear affine control systems, that will be considered in a forthcoming paper. |
warning/0506/math-ph0506040.html | ar5iv | text | # Proposition 2.1.
Large Parameter Behavior of Equilibrium Measures
Tamara Grava
Scuola Internazionale Superiore di Studi Avanzati, via Beirut 4, 34014 Trieste, Italy
grava@sissa.it
and
Fei-Ran Tian
Department of Mathematics, Ohio State University, Columbus, OH 43210
tian@math.ohio-state.edu
Abstract: We study the equilibrium measure for a logarithmic potential in the presence of an external field $`V_{}(\xi )+tp(\xi )`$, where $`t`$ is a parameter, $`V_{}(\xi )`$ is a smooth function and $`p(\xi )`$ a monic polynomial. When $`p(\xi )`$ is of an odd degree, the equilibrium measure is shown to be supported on a single interval as $`|t|`$ is sufficiently large. When $`p(\xi )`$ is of an even degree, the equilibrium measure is supported on two disjoint intervals as $`t`$ is negatively large; it is supported on a single interval for convex $`p(x)`$ as $`t`$ is positively large and is likely to be supported on multiple disjoint intervals for non-convex $`p(x)`$.
The support of the equilibrium measure shrinks to isolated points as $`|t|+\mathrm{}`$ in all the cases that we consider. For sufficiently large $`|t|`$, each topological component of the support contains a local minimizing point of the external field $`V_{}(\xi )+tp(\xi )`$; a “potential well” phenomenon.
$`\mathrm{\S }`$ 1 Introduction
In this paper, we study the following minimization problem with constraints
$$\underset{\{\psi 0,\psi 𝑑\xi =1\}}{\mathrm{Minimize}}\left[\frac{1}{2\pi }_{\mathrm{}}^+\mathrm{}_{\mathrm{}}^+\mathrm{}\mathrm{log}|\xi \eta |\psi (\xi )\psi (\eta )d\xi d\eta +_{\mathrm{}}^+\mathrm{}\mathrm{V}(\xi )\psi (\xi )d\xi \right].$$
(1.1)
The external field $`V(\xi )`$ is a $`C^{\mathrm{}}`$ function that satisfies
$$\underset{\xi \pm \mathrm{}}{lim}\frac{V(\xi )}{log(1+\xi ^2)}=+\mathrm{}.$$
(1.2)
Under this condition, the existence and uniqueness of the minimizer for (1.1) has been established . The measure $`\psi (\xi )d\xi `$, where $`\psi (\xi )`$ is the minimizer of (1.1), is called the equilibrium measure under the external field $`V(\xi )`$.
Equilibrium measures find applications in many branches of mathematical sciences. It is used to describe the partition function of the Hermitian one-matrix model in random matrix theory of statistical physics . It is also intrinsically connected to the free energy of the Yang-Mills theory . Finally, it plays an important role in orthogonal polynomials and approximation theory .
Although its importance to physics and approximation theory () is well known, the minimization problem (1.1) is not well understood. The minimizer is explicitly known only for a few cases where the external fields $`V(\xi )`$ are the simplest polynomials . It is therefore desirable to study the minimization problem for much more general $`C^{\mathrm{}}`$ external fields.
Our method to solve the minimization problem (1.1) is the same as one we used to solve a similar minimization problem for the zero dispersion limit of the KdV equation
$$u_t+6uu_x+ϵ^2u_{xxx}=0\text{with }u(x,0;ϵ)=u_0(x),$$
where the initial data $`u_0(x)`$ is a bounded decreasing function. The weak limit of the KdV solution $`u(x,t;ϵ)`$ as $`ϵ0`$ is determined by a minimization problem
$$\underset{\{\psi 0,\psi L^1\}}{\mathrm{Minimize}}\left[\frac{1}{2\pi }\mathrm{log}\left|\frac{\xi \eta }{\xi +\eta }\right|\psi (\xi )\psi (\eta )d\xi d\eta +\mathrm{V}(\xi ,\mathrm{x},\mathrm{t})\psi (\xi )d\xi \right],$$
(1.3)
where the space-time dependence is through the function $`\stackrel{~}{V}`$
$$\stackrel{~}{V}(\xi ,x,t)=x\xi 4\xi ^3t\theta (\xi )$$
(1.4)
and $`\theta (\xi )`$ is encoded with the initial information $`u_0(x)`$. We found that the problem (1.3) is intrinsically connected to the Euler-Poisson-Darboux equations . We utilized the solution of the equations to study the minimizer and hence the weak limit of the KdV solution.
The minimizer of (1.1) has been known to have a compact support if the external field $`V(\xi )`$ satisfies condition (1.2). The support is usually a union of a finite or infinite number of disjoint closed intervals. In this paper, we will study how the number of gaps in the support of the minimizer varies with respect to the external field. We would like to know whether there is any universality on the number of gaps.
We shall consider parameter-dependent external field $`V(\xi )=V_{}(\xi )+tp(\xi )`$, where $`t`$ is a parameter, $`V_{}(\xi )`$ is a smooth function of $`\xi `$ and $`p(\xi )`$ a monic polynomial. This is motivated by the space-time dependence of $`\stackrel{~}{V}`$ of (1.4) in the KdV minimization problem (1.3). We shall show that when $`p(\xi )`$ is of an odd degree, the minimizer of (1.1) has no gap in its compact support as $`|t|`$ is sufficiently large. When $`p(\xi )`$ is of an even degree, the minimizer has a single gap in its support as $`t`$ is negatively large; it has no gap for convex $`p(\xi )`$ as $`t`$ is positively large and is likely to have multi-gaps for non-convex $`p(\xi )`$. We note that the same minimizer may have an arbitrary number of gaps when $`|t|`$ is not large. A similar result has also been discovered in the case of the KdV weak limit . Namely, the KdV weak limit is of zero or single phase for all $`x`$ when $`t`$ is sufficiently large, i.e., the minimizer of (1.1) has either zero or one gap in its support.
The support of the minimizer shrinks to isolated points as $`|t|+\mathrm{}`$ in all the cases that we consider. Each topological component of the support contains a local minimizing point of the external field $`V(\xi )`$; a “potential well” phenomenon.
The organization of the paper is as follows.
In Section 2, we will use function theoretical methods to solve the minimization problem. We will formulate the minimizer in terms of solutions of the Euler-Poisson-Darboux equations.
In Section 3, we will use the Euler-Poisson-Darboux solutions to study the behavior of the minimizer when the parameter in the external field is sufficiently large.
$`\mathrm{\S }`$ 2 Solution of the Minimization Problem
In this section, we will use the method that we have developed for the KdV zero dispersion limit to solve the minimization problem (1.1).
Introducing a linear operator
$$L\psi (\xi )=\frac{1}{\pi }_{\mathrm{}}^+\mathrm{}\mathrm{log}|\xi \mu |\psi (\mu )𝑑\mu ,$$
(2.1)
we rewrite the quadratic functional of (1.1) as $`{\displaystyle \frac{1}{2}}<L\psi ,\psi >+<V,\psi >`$. Here $`<>`$ is the standard $`L^2`$ inner product. The Euler-Lagrange equations take the form
$`L\psi (\xi )V(\xi )`$ $`=`$ $`l\text{where }\psi >0,`$ (2.2)
$`L\psi (\xi )V(\xi )`$ $``$ $`l\text{where }\psi =0,`$ (2.3)
where $`l`$ is the Lagrange multiplier. It can be shown that $`\psi `$ is the minimizer iff $`\psi `$ is a nonnegative function that satisfies variational conditions (2.2-2.3) and the constraint
$$_{\mathrm{}}^+\mathrm{}\psi (\xi )𝑑\xi =1.$$
(2.4)
We make the ansatz that the support of $`\psi `$ consists of a finite union of disjoint intervals. One denotes $`I=\{\xi ;\psi >0\}`$ and writes
$$I=_{k=1}^{g+1}(u_{2k},u_{2k1}),$$
(2.5)
where $`u_{2g+2}<\mathrm{}<u_2<u_1`$. Hence, the support is the closure of $`I`$.
We now consider a slightly stronger version of (2.2) and (2.3),
$`L\psi (\xi )V(\xi )`$ $`=`$ $`l\text{on }I,`$ (2.6)
$`L\psi (\xi )V(\xi )`$ $`<`$ $`l\text{on }\backslash \overline{I},`$ (2.7)
where $`\overline{I}`$ denotes the closure of $`I`$. Since $`\psi 0`$, we must also have
$`\psi `$ $`>`$ $`0\text{on }I,`$ (2.8)
$`\psi `$ $`=`$ $`0\text{off }I.`$ (2.9)
Our strategy to construct the minimizer is to first find the solution $`\psi `$ of equations (2.4), (2.6) and (2.9) and then impose inequalities (2.7) and (2.8) on $`\psi `$. Since it is a non-negative function and since it satisfies (2.2), (2.3) and (2.4), $`\psi `$ will then be the minimizer.
To solve (2.4), (2.6) and (2.9), we use the fact that the operator $`L`$ of (2.1) is connected to the Hilbert transform $`H`$ on the real line.
$$L\psi (\xi )=_0^\xi H\psi (\tau )𝑑\tau ,$$
(2.10)
where
$$H\psi (\xi )=\frac{1}{\pi }P.V._{\mathrm{}}^{\mathrm{}}\frac{\psi (\mu )d\mu }{\xi \mu }.$$
This makes equations (2.6) and (2.9) amenable to the Riemann-Hilbert technique in function theory.
Differentiating (2.6) with respect to $`\xi `$ and using (2.10), one obtains
$$H\psi (\xi )=V^{}(\xi )\text{on I}.$$
(2.11)
To recover equation (2.6) from (2.11) by integration, one must have for $`k`$ $`=`$ $`1`$, $`2`$, $`\mathrm{}`$, $`g`$,
$$_{u_{2k+1}}^{u_{2k}}[H\psi (\xi )V^{}(\xi )]𝑑\xi =0.$$
(2.12)
Recalling the relation of the Hilbert transform to analytic function, one can write for real $`\xi `$
$$𝒢^+(\xi )=\psi (\xi )+\sqrt{1}H\psi (\xi ),$$
where $`𝒢^+(\xi )`$ is the boundary value on the real axis of a function
$$𝒢(z)=\frac{1}{\pi \sqrt{1}}_{\mathrm{}}^+\mathrm{}\frac{\psi (\mu )d\mu }{\mu z},$$
which is analytic in the upper half complex plane. In view of the constraint (2.4), we expand the Cauchy integral to obtain $`𝒢(z)=\frac{1}{\sqrt{1}\pi z}+O(1/z^2)`$ for large $`|z|`$.
Conditions (2.12) now take a new form
$$_{u_{2k+1}}^{u_{2k}}[Im𝒢^+(\xi )V^{}(\xi )]𝑑\xi =0$$
(2.13)
for $`k=1,2,\mathrm{},g`$.
Equations (2.9) and (2.11) then become a Riemann-Hilbert problem in function theory
$`Im𝒢^+(\xi )`$ $`=`$ $`V^{}(\xi )\text{on }I,`$ (2.14)
$`Re𝒢^+(\xi )`$ $`=`$ $`0\text{off }I,`$ (2.15)
where $`𝒢(z)`$ is analytic in the upper half complex plane. It follows from the Plemelj formula that
$$\stackrel{~}{𝒢}(z)=\frac{R(z,\stackrel{}{u})}{\pi }_I\frac{V^{}(\mu )}{R(\mu ,\stackrel{}{u})(\mu z)}𝑑\mu ,$$
where $`\stackrel{}{u}`$ denotes $`(u_1,u_2,\mathrm{},u_{2g+2})`$ and $`R(\xi ,\stackrel{}{u})=\sqrt{(\xi u_1)(\xi u_2)\mathrm{}(\xi u_{2g+2})}`$, is a solution to this Riemann-Hilbert problem. Here $`R(\xi ,\stackrel{}{u})`$ is set to be positive for $`\xi >u_1`$. It defines a Riemann surface with branch cuts along the set $`I`$ of (2.5).
To derive the equations governing the endpoints $`u_1`$, $`u_2`$, $`\mathrm{}`$, $`u_{2g+2}`$, one usually imposes condition (2.13) and the asymptotics $`𝒢(z)=\frac{1}{\sqrt{1}\pi z}+O(1/z^2)`$ for large $`|z|`$ . The former gives $`g`$ equations and the latter results in $`g+2`$ moment conditions. Consequently, one obtains exactly $`2g+2`$ algebraic equations on $`u_1`$, $`u_2`$, $`\mathrm{}`$, $`u_{2g+2}`$.
In this paper, we take a slightly different approach. We observe that the Riemann-Hilbert problem (2.14-2.15) has many other solutions. Indeed, it is obvious that
$$𝒢(z)=\frac{\frac{R^2(z,\stackrel{}{u})}{\pi }_I\frac{V^{}(\mu )}{R(\mu ,\stackrel{}{u})(\mu z)}𝑑\mu +\sqrt{1}Q(z)}{R(z,\stackrel{}{u})},$$
(2.16)
where $`Q(z)`$ is an arbitrary polynomial with real coefficients, is also a solution. We choose polynomial $`Q(z)`$ such that
1. it is a polynomial of degree $`2g+1`$.
2. $`𝒢(z)=\frac{1}{\sqrt{1}\pi z}+O(1/z^2)`$ for large $`|z|`$.
3. conditions (2.13) are satisfied.
Hence, constraint (2.4) and conditions (2.6) and (2.9) are built in the construction of $`𝒢`$.
It is quite obvious that such a polynomial $`Q`$ is unique.
We first analyze the boundary value of $`𝒢(z)`$ on the real axis. Necessarily, the Cauchy integral in the numerator of (2.16) will become a singular integral. Our key observation is that the latter is, more or less, the solution of a boundary value problem for the Euler-Poisson-Darboux equations.
###### Proposition 2.1.
The boundary value of $`𝒢(z)`$ at $`\xi u_i`$, $`i=1,2,\mathrm{},2g+2`$
$$𝒢^+(\xi )=\frac{2\sqrt{1}R^2(\xi ,\stackrel{}{u})\mathrm{\Phi }_g(\xi ,u_1,u_2,\mathrm{},u_{2g+1})+\sqrt{1}Q(\xi )}{R(\xi ,\stackrel{}{u})}+\sqrt{1}V^{}(\xi ),$$
(2.17)
where $`\mathrm{\Phi }_g(\xi ,u_1,\mathrm{},u_{2g+2})`$ satisfies the Euler-Poisson-Darboux equations
$`2(u_iu_j){\displaystyle \frac{^2\mathrm{\Phi }_g}{u_iu_j}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Phi }_g}{u_i}}{\displaystyle \frac{\mathrm{\Phi }_g}{u_j}},`$ (2.18)
$`2(\xi u_i){\displaystyle \frac{^2\mathrm{\Phi }_g}{\xi u_i}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Phi }_g}{\xi }}2{\displaystyle \frac{\mathrm{\Phi }_g}{u_i}},`$ (2.19)
$`\mathrm{\Phi }_g(u,u,\mathrm{},u)`$ $`=`$ $`{\displaystyle \frac{1}{2(g+1)!}}{\displaystyle \frac{d^{g+2}V(u)}{du^{g+2}}}.`$ (2.20)
We will omit the proof here, since it is similar to the proof of an analogous result on the KdV zero dispersion limit .
A formula for $`Q`$ can be derived from the loop conditions (2.13) and asymptotics $`𝒢(z)=\frac{1}{\sqrt{1}\pi z}+O(1/z^2)`$ for large $`|z|`$. To achieve this, we introduce a sequence of polynomials
$$P_{g,n}(\xi ,\stackrel{}{u})=\xi ^{g+n}+a_{g,1}\xi ^{g+n1}+\mathrm{}+a_{g,g+n}$$
(2.21)
whose coefficients are uniquely determined by
$$\frac{P_{g,n}(\xi ,\stackrel{}{u})}{R(\xi ,\stackrel{}{u})}=\xi ^{n1}+O(\frac{1}{\xi ^2})\text{for large }|\xi |,$$
(2.22)
and
$$_{u_{2k+1}}^{u_{2k}}\frac{P_{g,n}(\xi ,\stackrel{}{u})}{R(\xi ,\stackrel{}{u})}𝑑\xi =0k=1,2,\mathrm{},g.$$
(2.23)
###### Proposition 2.2.
$`Q(\xi ,\stackrel{}{u})`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{2g+2}{}}}[{\displaystyle \underset{l=1,li}{\overset{2g+2}{}}}(\xi u_l)]\mathrm{\Psi }_g(u_i,\stackrel{}{u})+{\displaystyle \frac{1}{\pi }}P_{g,0}(\xi ,\stackrel{}{u})+c_1P_{g,1}(\xi ,\stackrel{}{u})`$
$`+\mathrm{}+c_gP_{g,g}(\xi ,\stackrel{}{u}).`$
Here $`\mathrm{\Psi }_g(\xi ,u_1,u_2,\mathrm{},u_{2g+2})`$ satisfies the Euler-Poisson-Darboux equations (2.18) and (2.19) with the diagonal boundary value
$$\mathrm{\Psi }_g(u,u,\mathrm{},u)=\frac{1}{2(g+1)!}\frac{d^{g+1}V(u)}{du^{g+1}}.$$
(2.24)
The coefficients
$$c_k=2k\underset{l=0}{\overset{gk}{}}\mathrm{\Gamma }_l(\stackrel{}{u})q_{g,k+l}(\stackrel{}{u})k=1,2,\mathrm{},g,$$
where $`\mathrm{\Gamma }_l(\stackrel{}{u})`$’s come from the expansion
$$R(\mu ,\stackrel{}{u})=\mu ^{g+1}[\mathrm{\Gamma }_0(\stackrel{}{u})+\frac{\mathrm{\Gamma }_1(\stackrel{}{u})}{\mu }+\frac{\mathrm{\Gamma }_2(\stackrel{}{u})}{\mu ^2}+\mathrm{}].$$
(2.25)
The function $`q_{g,k}`$ is
$$q_{g,k}(\stackrel{}{u})=\frac{1}{2\pi \sqrt{1}}_I\frac{V(\mu )\mu ^{gk}}{R(\mu ,\stackrel{}{u})}𝑑\mu .$$
(2.26)
The proof is also similar to that of the KdV case .
The function $`q_{g,k}(\stackrel{}{u})`$ also satisfies equations of Euler-Poisson-Darboux type .
Finally, we postulate the boundedness of the minimizer, which is equal to the real part of $`𝒢^+(\xi )`$. The numerator in (2.17) must then vanish at the end points of set $`I`$
$$P(u_i,\stackrel{}{u})=0i=1,2,\mathrm{},2g+2,$$
(2.27)
where
$$P(\xi ,\stackrel{}{u})=2R^2(\xi ,\stackrel{}{u})\mathrm{\Phi }_g(\xi ,\stackrel{}{u})Q(\xi ,\stackrel{}{u}).$$
(2.28)
Since it is of degree $`2g+1`$ in $`\xi `$ and has $`2g+2`$ zeros because of (2.27), the polynomial $`Q(\xi ,\stackrel{}{u})`$ must be identically zero. The real part of $`𝒢^+`$ of (2.17), on the solution $`\stackrel{}{u}`$ of equations (2.27), becomes
$$\psi (\xi )=2Re\{\sqrt{1}R(\xi ,\stackrel{}{u})\}\mathrm{\Phi }_g(\xi ,\stackrel{}{u}).$$
(2.29)
Hence, $`\psi (\xi )`$ of (2.29) satisfies conditions (2.4), (2.6) and (2.9).
Inequality (2.8) is equivalent to
$$Re\{\sqrt{1}R(\xi ,\stackrel{}{u})\}\mathrm{\Phi }_g(\xi ,\stackrel{}{u})<0\text{for }\xi \text{ on I},$$
(2.30)
and (2.7) is equivalent to
$`{\displaystyle _{u_{2k+1}}^\xi }R(\mu ,\stackrel{}{u})\mathrm{\Phi }_g(\mu ,\stackrel{}{u})𝑑\mu `$ $`>`$ $`0\text{for }u_{2k+1}<\xi <u_{2k}\text{ and }1k2g+2,`$
$`{\displaystyle _{u_1}^\xi }R(\mu ,\stackrel{}{u})\mathrm{\Phi }_g(\mu ,\stackrel{}{u})𝑑\mu `$ $`>`$ $`0\text{for }\xi >u_1,`$ (2.31)
$`{\displaystyle _\xi ^{u_{2g+2}}}R(\mu ,\stackrel{}{u})\mathrm{\Phi }_g(\mu ,\stackrel{}{u})𝑑\mu `$ $`<`$ $`0\text{for }\xi <u_{2g+2}.`$
We summarize the above results in the following theorem
###### Theorem 2.3.
If $`(u_1,u_2,\mathrm{},u_{2g+2})`$ satisfies equations (2.27) and if inequalities (2.30-2.31) are satisfied, then $`\psi (\xi )`$ of (2.29) is a non-negative function that satisfies variational conditions (2.2) and (2.3) and constraint (2.4); so $`\psi (\xi )`$ must be the minimizer of the minimization problem (1.1).
$`\mathrm{\S }`$ 3 Large Parameter Results
In this section, we shall consider a one-parameter family of external field
$$V(\xi )=V_{}(\xi )+tp(\xi ).$$
(3.1)
Here, $`t`$ is a parameter, $`p(\xi )`$ is a monic polynomial of degree $`n`$ and $`V_{}(\xi )`$ is a $`C^{\mathrm{}}(\mathrm{},+\mathrm{})`$ function. We are interested in the behavior of the equilibrium measure when $`|t|`$ is large.
Our strategy for studying the large parameter behavior of the minimizer of (1.1) is as follows. For simplicity, we will assume $`V_{}(\xi )`$ to have a power function growth at $`\xi =\mathrm{}`$ or $`\xi =+\mathrm{}`$. This allows us to use the scaling technique to study the equilibrium measure when the parameter $`t`$ is sufficiently large.
$`\mathrm{\S }`$ 3.1 The degree of polynomial $`p(\xi )`$ is odd
We shall show that the equilibrium measure is supported on a single interval $`[u_2,u_1]`$ when $`|t|`$ is large. This corresponds to (2.5) for $`g=0`$.
We will use Theorem 2.3 to construct the minimizer of (1.1). Hence, we need to solve equations (2.27) for $`g=0`$ and verify inequalities (2.30-2.31).
We now consider the algebraic equations (2.27) for $`g=0`$. The function $`P(\xi ,u_1,u_2)`$ of (2.28) is
$$2(\xi u_1)(\xi u_2)\mathrm{\Phi }_0(\xi ,u_1,u_2)+(\xi u_2)\mathrm{\Psi }_0(u_1,u_1,u_2)+(\xi u_1)\mathrm{\Psi }_0(u_2,u_1,u_2)\frac{1}{\pi }.$$
Equations (2.27) exactly become
$`(u_1u_2)\mathrm{\Psi }_0(u_1,u_1,u_2){\displaystyle \frac{1}{\pi }}`$ $`=`$ $`0,`$ (3.2)
$`(u_2u_1)\mathrm{\Psi }_0(u_2,u_1,u_2){\displaystyle \frac{1}{\pi }}`$ $`=`$ $`0,`$ (3.3)
which are equivalent to
$`(u_1u_2)\sqrt{{\displaystyle \frac{}{u_2}}\mathrm{\Psi }_0(u_1,u_1,u_2)}\sqrt{{\displaystyle \frac{1}{\pi }}}`$ $`=`$ $`0,`$ (3.4)
$`\mathrm{\Psi }_0(u_1,u_1,u_2)+\mathrm{\Psi }_0(u_2,u_1,u_2)`$ $`=`$ $`0.`$ (3.5)
In the derivation of (3.4), we have used an identity
$$\mathrm{\Psi }_0(u_1,u_1,u_2)\mathrm{\Psi }_0(u_2,u_1,u_2)=2(u_1u_2)\frac{}{u_2}\mathrm{\Psi }_0(u_1,u_1,u_2).$$
(3.6)
This is easily verified by calculating
$`{\displaystyle \frac{}{\xi }}\left[2(\xi u_2){\displaystyle \frac{}{u_2}}\mathrm{\Psi }_0(\xi ,u_1,u_2)\mathrm{\Psi }_0(\xi ,u_1,u_2)\right]`$
$`=`$ $`2(\xi u_2){\displaystyle \frac{^2}{u_2\xi }}\mathrm{\Psi }_0+2{\displaystyle \frac{}{u_2}}\mathrm{\Psi }_0{\displaystyle \frac{}{\xi }}\mathrm{\Psi }_0.`$
The right hand side vanishes because $`\mathrm{\Psi }_g`$ satisfies the Euler-Poisson-Darboux equation (2.19). The function in the parenthesis is then independent of $`\xi `$; so we obtain
$$2(\xi u_2)\frac{}{u_2}\mathrm{\Psi }_0(\xi ,u_1,u_2)\mathrm{\Psi }_0(\xi ,u_1,u_2)=\mathrm{\Psi }_0(u_2,u_1,u_2).$$
Letting $`\xi =u_1`$ yields the identity (3.6).
We will rewrite equation (3.5) in another useful form. Its left hand side is a function of $`u_1`$ and $`u_2`$. Denote this function by $`H(u_1,u_2)`$. Since $`\mathrm{\Psi }_0(\xi ,u_1,u_2)`$ satisfies equations (2.18-2.19), we derive an equation for $`H(u_1,u_2)`$
$$2(u_1u_2)\frac{^2H}{u_1u_2}=\frac{H}{u_1}\frac{H}{u_2}.$$
The boundary condition (2.24) for $`\mathrm{\Psi }_0(\xi ,u_1,u_2)`$ implies $`H(u,u)=V^{}(u)`$. We then use formula (B.3) to obtain
$$H(u_1,u_2)=\frac{1}{\pi }_{u_2}^{u_1}\frac{V^{}(\mu )d\mu }{\sqrt{(u_1\mu )(\mu u_2)}}.$$
Hence, equation (3.5) is equivalent to
$$_{u_2}^{u_1}\frac{V^{}(\mu )d\mu }{\sqrt{(u_1\mu )(\mu u_2)}}=0.$$
(3.7)
We shall study the case when $`t`$ is negatively large in detail. The other case when $`t`$ is positively large can be handled in the same way.
We now make assumptions on $`V_{}(\xi )`$. Since $`V(\xi )=V_{}(\xi )+tp(\xi )`$ satisfies the growth condition (1.2), $`V_{}(\xi )`$ must grow faster than $`p(\xi )`$ as $`\xi +\mathrm{}`$ because of $`t<0`$. In contrast, we will make a very mild assumption on $`V_{}(\xi )`$ as $`\xi 1`$
$`V_{}(\xi )`$ $`=`$ $`C_+\xi ^{M_+}+h_+(\xi )\text{ for }\xi 1,`$ (3.8)
$`V_{}^{}(\xi )`$ $``$ $`0\text{for }\xi 1.`$ (3.9)
Here $`C_+`$ is a positive constant and $`M_+`$ is a positive constant that is bigger than $`n`$ and $`2`$. The number $`n`$ is the degree of the polynomial $`p(\xi )`$. Function $`h_+(\xi )`$ has an order less than $`n`$. More generally, we assume
$$\underset{\xi +\mathrm{}}{lim}\frac{h_+^{\prime \prime \prime }(\xi )}{\xi ^{M_+3}}=0.$$
(3.10)
This immediately implies
$$\underset{\xi +\mathrm{}}{lim}\frac{h_+(\xi )}{\xi ^{M_+}}=0,\underset{\xi +\mathrm{}}{lim}\frac{h_+^{}(\xi )}{\xi ^{M_+1}}=0,\underset{\xi +\mathrm{}}{lim}\frac{h_+^{\prime \prime }(\xi )}{\xi ^{M_+2}}=0.$$
(3.11)
We will solve equations (3.4) and (3.5) for large $`u_1`$ and $`u_2`$. This is motivated by the fact that $`V(\xi )`$ has a critical point between $`u_2`$ and $`u_1`$ in view of equation (3.7) and that the minimizing point of $`V(\xi )=V_{}(\xi )+tp(\xi )`$ moves to $`+\mathrm{}`$ as $`t\mathrm{}`$.
We will first consider the case $`p(\xi )=\xi ^n`$.
We will split $`\mathrm{\Psi }_0`$ of (3.4-3.5) into simpler terms. In view of its boundary data (2.24) for $`g=0`$, $`\mathrm{\Psi }_0`$ depends linearly on $`V`$. The decomposition of $`V=t\xi ^n+C_+\xi ^{M_+1}+h_+(\xi )`$ allows us to write
$$\mathrm{\Psi }_0=t\mathrm{\Psi }_{\xi ^n}+C_+\mathrm{\Psi }_{\xi ^{M_+}}+\mathrm{\Psi }_{h_+(\xi )}.$$
(3.12)
In view of the integral formula (B.4), $`\mathrm{\Psi }_{\xi ^n}`$ and $`\mathrm{\Psi }_{\xi ^{M_+}}`$ are homogeneous functions of $`(\xi ,u_1,u_2)`$ of orders $`n1`$ and $`M_+1`$, respectively.
Substituting (3.12) into equation (3.4) and then dividing it by $`|t|^{\frac{M_+}{2(M_+n)}}`$, we use the homogeneity of $`\mathrm{\Psi }_{\xi ^n}`$ and $`\mathrm{\Psi }_{\xi ^{M_+}}`$ to obtain
$`(U_1U_2)\{{\displaystyle \frac{}{U_2}}\mathrm{\Psi }_{\xi ^n}(U_1,U_1,U_2)+C_+{\displaystyle \frac{}{U_2}}\mathrm{\Psi }_{\xi ^{M_+}}(U_1,U_1,U_2)`$ (3.13)
$`+{\displaystyle \frac{1}{|t|^{\frac{M_+2}{M_+n}}}}{\displaystyle \frac{}{u_2}}\mathrm{\Psi }_{h_+(\xi )}(u_1,u_1,u_2)\}^{\frac{1}{2}}{\displaystyle \frac{1}{\sqrt{\pi }|t|^{\frac{M_+}{2(M_+n)}}}}=0.`$
Here $`U_1`$ and $`U_2`$ are
$$U_1=\frac{u_1}{|t|^{\frac{1}{M_+n}}},U_2=\frac{u_2}{|t|^{\frac{1}{M_+n}}}.$$
(3.14)
Similarly, substituting (3.12) into equation (3.5) and then dividing it by $`|t|^{\frac{M_+1}{M_+n}}`$, we obtain
$`\mathrm{\Psi }_{\xi ^n}(U_1,U_1,U_2)+C_+\mathrm{\Psi }_{\xi ^{M_+}}(U_1,U_1,U_2)+{\displaystyle \frac{1}{|t|^{\frac{M_+1}{M_+n}}}}\mathrm{\Psi }_{h_+(\xi )}(u_1,u_1,u_2)`$
$`\mathrm{\Psi }_{\xi ^n}(U_2,U_1,U_2)+C_+\mathrm{\Psi }_{\xi ^{M_+}}(U_2,U_1,U_2)+{\displaystyle \frac{1}{|t|^{\frac{M_+1}{M_+n}}}}\mathrm{\Psi }_{h_+(\xi )}(u_2,u_1,u_2)=0.`$
It follows from the behavior (3.10-3.11) of $`h_+(\xi )`$ that the three terms involving $`\mathrm{\Psi }_{h_+(\xi )}`$ in (3.13-S3.Ex8), together with their first derivatives with respect to $`U_1`$ and $`U_2`$, decay to zero as $`t\mathrm{}`$ if $`U_1`$ and $`U_2`$ of (3.14) are kept bounded.
We now denote $`1/t`$ by $`T`$ and solve equations (3.13) and (S3.Ex8) for $`U_1`$ and $`U_2`$ as functions of $`T`$ in the neighborhood of $`T=0`$.
First, equations (3.13) and (S3.Ex8) at $`T=0`$ become
$`(U_1U_2)\sqrt{{\displaystyle \frac{}{U_2}}\mathrm{\Psi }_{\xi ^n}(U_1,U_1,U_2)+C_+{\displaystyle \frac{}{U_2}}\mathrm{\Psi }_{\xi ^{M_+}}(U_1,U_1,U_2)}`$ $`=`$ $`0,`$ (3.16)
$`\mathrm{\Psi }_{\xi ^n}(U_1,U_1,U_2)+C_+\mathrm{\Psi }_{\xi ^{M_+}}(U_1,U_1,U_2)`$
$`\mathrm{\Psi }_{\xi ^n}(U_2,U_1,U_2)+C_+\mathrm{\Psi }_{\xi ^{M_+}}(U_2,U_1,U_2)`$ $`=`$ $`0.`$ (3.17)
These two equations have a solution $`U_1=U_2=U^{}`$, where $`U^{}`$ is determined by
$$\mathrm{\Psi }_{\xi ^n}(U^{},U^{},U^{})+C_+\mathrm{\Psi }_{\xi ^{M_+}}(U^{},U^{},U^{})=\frac{1}{2}\frac{d}{d\xi }[\xi ^n+C_+\xi ^{M_+}]|_{\xi =U^{}}=0,$$
where we have used the boundary condition (2.24) for $`\mathrm{\Psi }_{\xi ^n}`$ and $`\mathrm{\Psi }_{\xi ^{M_+}}`$. We hence obtain
$$U_1=U_2=U^{}=[\frac{n}{C_+M_+}]^{\frac{1}{M_+n}}$$
as a solution of (3.16-3.17).
Second, we calculate the Jacobian of equations (3.13) and (S3.Ex8) at $`U_1=U_2=U^{}`$ and $`T=0`$. Denote the left hand side of (3.13) by $`F_1(U_1,U_2,T)`$ and that of (S3.Ex8) by $`F_2(U_1,U_2,T)`$. Since $`\mathrm{\Psi }_g`$ satisfies the Euler-Poisson-Darboux equations (2.18), (2.19) and (2.24), we obtain
$`{\displaystyle \frac{F_1}{U_1}}={\displaystyle \frac{F_1}{U_2}}`$ $`=`$ $`\sqrt{{\displaystyle \frac{}{U_2}}\mathrm{\Psi }_{\xi ^n}(U^{},U^{},U^{})+C_+{\displaystyle \frac{}{U_2}}\mathrm{\Psi }_{\xi ^{M_+}}(U^{},U^{},U^{})}`$
$`=`$ $`\sqrt{{\displaystyle \frac{1}{8}}{\displaystyle \frac{d^2}{d\xi ^2}}[\xi ^n+C_+\xi ^{M_+}]}|_{\xi =U^{}}>0,`$
$`{\displaystyle \frac{F_2}{U_1}}={\displaystyle \frac{F_2}{U_2}}`$ $`=`$ $`4[{\displaystyle \frac{}{U_2}}\mathrm{\Psi }_{\xi ^n}(U^{},U^{},U^{})+C_+{\displaystyle \frac{}{U_2}}\mathrm{\Psi }_{\xi ^{M_+}}(U^{},U^{},U^{})]`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{d^2}{d\xi ^2}}[\xi ^n+C_+\xi ^{M_+}]|_{\xi =U^{}}>0.`$
Hence, the Jacobian of (3.13) and (S3.Ex8) is nonzero. Equations (3.13) and (S3.Ex8) then give $`U_1`$ and $`U_2`$ as functions of $`T`$ near $`T=0`$.
Therefore, equations (3.2-3.3) have a solution $`u_1(t)`$, $`u_2(t)`$ for negatively large $`t`$. The solution has the following asymptotics
$$u_1(t)U^{}|t|^{\frac{1}{M_+n}},u_2(t)U^{}|t|^{\frac{1}{M_+n}}\text{for }t1.$$
Moreover, it follows from (3.13) and (3.14) that the length of the interval $`[u_1(t),u_1(t)]`$, the support of the minimizer, shrinks to zero as $`t\mathrm{}`$.
To make sure that $`\psi (\xi )`$ of (2.29) is the minimizer, we also need to verify that inequalities (2.30-2.31) are satisfied.
We split $`\mathrm{\Phi }_0`$ in the fashion of (3.12)
$$\mathrm{\Phi }_0=t\mathrm{\Phi }_{\xi ^n}+C_+\mathrm{\Phi }_{\xi ^{M_+}}+\mathrm{\Phi }_{h_+(\xi )}.$$
Using the scaling (3.14), we write $`\mathrm{\Phi }_0(\xi ,u_1(t),u_2(t))`$ as
$$|t|^{\frac{M_+1}{M_+n}}[\mathrm{\Phi }_{\xi ^n}(\mathrm{\Xi },U_1,U_2)+C_+\mathrm{\Phi }_{\xi ^{M_+}}(\mathrm{\Xi },U_1,U_2)+\frac{1}{|t|^{\frac{M_+1}{M_+n}}}\mathrm{\Phi }_{h_+(\xi )}(\xi ,u_1,u_2)],$$
(3.18)
where $`\mathrm{\Xi }=\xi /|t|^{\frac{1}{M_+n}}`$. We use (3.11) to deduct that if $`\mathrm{\Xi }`$ is kept bounded, say, $`U^{}/3\mathrm{\Xi }3U^{}`$, the last term in the parenthesis goes to zero uniformly as $`t\mathrm{}`$. The sum in the parenthesis then has the limit
$`\mathrm{\Phi }_{\xi ^n}(\mathrm{\Xi },U^{},U^{})+C_+\mathrm{\Phi }_{\xi ^{M_+}}(\mathrm{\Xi },U^{},U)`$
$`=`$ $`{\displaystyle \frac{\frac{d}{d\mathrm{\Xi }}[\mathrm{\Xi }^n+C_+\mathrm{\Xi }^{M_+}]\frac{d}{d\xi }[\xi ^n+C_+\xi ^{M_+}]|_{\xi =U^{}}}{2(\mathrm{\Xi }U^{})}}`$
$`=`$ $`{\displaystyle \frac{\frac{d}{d\mathrm{\Xi }}[\mathrm{\Xi }^n+C_+\mathrm{\Xi }^{M_+}]}{2(\mathrm{\Xi }U^{})}}`$
$`=`$ $`{\displaystyle \frac{C_+M_+\mathrm{\Xi }^{n1}[\mathrm{\Xi }^{M_+n}U^{^{M_+n}}]}{2(\mathrm{\Xi }U^{})}}>0\text{for }\frac{U^{}}{3}\mathrm{\Xi }3U^{}.`$
Here we have used formula (B.5) for $`\mathrm{\Phi }_{\xi ^n}`$ and $`\mathrm{\Phi }_{\xi ^{M_+}}`$ in the first equality and $`U^{}=[n/C_+M_+]^{1/(M_+n)}`$ in the last two equalities. In view of (3.18), we have the inequality $`\mathrm{\Phi }_0(\xi ,u_1(t),u_2(t))>0`$ for $`\frac{1}{2}u_2(t)\xi 2u_1(t)`$ when $`t`$ is negatively large. This proves the inequality (2.30).
To prove inequality (2.31), we first show that $`\mathrm{\Phi }_0(\xi ,u_1(t),u_2(t))>0`$ uniformly for $`\xi <1`$ and $`\xi >2u_1(t)`$ when $`t`$ is negatively large.
We shall prove the uniform positivity for $`\xi <1`$ first. We use formula (B.6) to write $`\mathrm{\Phi }_0(\xi ,u_1(t),u_2(t))`$ as
$`{\displaystyle \frac{1}{2\sqrt{(\xi u_1)(\xi u_2)}}}[V^{}(\xi )`$ (3.19)
$`{\displaystyle \frac{1}{\pi }}{\displaystyle _{u_2}^{u_1}}{\displaystyle \frac{\sqrt{(\xi u_1)(\xi u_2)}}{\xi \mu }}{\displaystyle \frac{V^{}(\mu )d\mu }{\sqrt{(u_1\mu )(\mu u_2)}}}]`$
$`=`$ $`{\displaystyle \frac{1}{2\sqrt{(\xi u_1)(\xi u_2)}}}[V^{}(\xi )`$
$`{\displaystyle \frac{1}{\pi }}{\displaystyle _{u_2}^{u_1}}({\displaystyle \frac{\sqrt{(\xi u_1)(\xi u_2)}}{\xi \mu }}1){\displaystyle \frac{V^{}(\mu )d\mu }{\sqrt{(u_1\mu )(\mu u_2)}}}],`$
where in the last equality we have used equation (3.7).
We now show that the integral of (3.19) tends to zero uniformly for $`\xi <0`$ as $`t\mathrm{}`$. Because of the inequality
$$|\frac{\sqrt{(\xi u_1)(\xi u_2)}}{\xi \mu }1|\frac{u_1u_2}{u_2\xi }\frac{u_1u_2}{u_2}\text{for }\xi <0<u_2<u_1,$$
the integral is bounded by
$$\frac{u_1u_2}{u_2}_{u_2}^{u_1}\frac{|V^{}(\mu )|}{\sqrt{(u_1\mu )(\mu u_2)}}𝑑\mu $$
for $`\xi <0`$. This, in view of the scaling (3.14), is further bounded by a constant times
$$\frac{[U_1(t)U_2(t)]^2}{u_2(t)}|t|^{\frac{M_+}{M_+n}},$$
where we have used $`V^{}(\mu )=O((u_1(t)u_2(t))t^{\frac{M_+2}{M_+n}})`$ for $`u_2(t)\mu u_1(t)`$. This asymptotics can be verified by observing that $`V^{}(\mu )`$ has a zero between $`u_2(t)`$ and $`u_1(t)`$ because of (3.7) and by expanding $`V^{}(\mu )`$ around the zero. Since $`U_1(t)U_2(t)=O(|t|^{\frac{M_+}{2(M_+n)}})`$ because of (3.13), the integral of (3.19) tends to zero uniformly for $`\xi <0`$ as $`t\mathrm{}`$.
The first term $`V^{}(\xi )=tn\xi ^{n1}+V_{}^{}(\xi )`$ in the parenthesis of (3.19) is bounded from above by a negative constant for $`\xi 1`$ when $`t`$ is negatively large. This immediately follows from condition (3.9).
Therefore, $`\mathrm{\Phi }_0(\xi ,u_1(t),u_2(t))>0`$ uniformly for $`\xi <1`$ when $`t`$ is negatively large.
In the same way, we can prove that $`\mathrm{\Phi }_0(\xi ,u_1(t),u_2(t))>0`$ uniformly for $`\xi >2u_1(t)`$ when $`t`$ is negatively large. This together with $`\mathrm{\Phi }_0(\xi ,u_1(t),u_2(t))>0`$ for $`\frac{1}{2}u_2(t)\xi 2u_1(t)`$ proves the first half of (2.31).
It remains to prove the rest of (2.31), i.e.,
$$_\xi ^{u_2(t)}\sqrt{(\mu u_1(t))(\mu u_2(t))}\mathrm{\Phi }_0(\mu ,u_1(t),u_2(t))𝑑\mu <0$$
(3.20)
for $`\xi <u_2(t)`$. Since $`\mathrm{\Phi }_0(\mu ,u_1(t),u_2(t))>0`$ uniformly for $`\xi <1`$ and $`\xi u_2(t)/2`$ when $`t`$ is negatively large, it suffices to prove the inequality for $`1\xi u_2(t)/2`$.
Using formula (B.6) for $`\mathrm{\Phi }_0(\xi ,u_1(t),u_2(t))`$, we write the left hand side of the inequality of (3.20) as half of
$`V(\xi )+V(u_2(t))`$
$`{\displaystyle \frac{1}{\pi }}{\displaystyle _\xi ^{u_2(t)}}\sqrt{(\xi u_1)(\xi u_2)}[{\displaystyle _{u_2}^{u_1}}{\displaystyle \frac{V^{}(\mu )d\mu }{(\xi \mu )\sqrt{(u_1\mu )(\mu u_2)}}}]𝑑\xi .`$
The first two terms $`V(\xi )+V(u_2(t))V(u_2(t)/2)+V(u_2(t)<0`$ for $`1\xi u_2(t)/2`$. The whole third term, in view of the scaling (3.14), has an order in $`t`$ lower than that of $`V(u_2(t)/2)+V(u_2(t))`$. These prove that (S3.Ex28) is negative for $`1\xi u_2(t)/2`$ when $`t`$ is negatively large. We have therefore proved (3.20).
Inequalities (2.30-2.31) have thus been verified. By Theorem 2.3, the equilibrium measure for the external field $`V(\xi )=t\xi ^n+V_{}(\xi )`$ has no gap in its support when $`t`$ is negatively large enough.
The scaling nature of the above approach allows us to extend the result from $`p(\xi )=\xi ^n`$ to an arbitrary monic polynomial of degree $`n`$.
###### Theorem 3.1.
Under conditions (3.8-3.10) on $`V_{}(\xi )`$, the equilibrium measure for the external field $`V(\xi ,t)=V_{}(\xi )+tp(\xi )`$, where $`p(\xi )`$ is a monic polynomial of an odd degree $`n`$, has no gap in its support when $`t`$ is negatively large enough.
The case when $`t`$ is positively large can be treated in the same way. Here, we replace conditions (3.8-3.11) on $`V_{}(\xi )`$ by analogous ones
$`V_{}(\xi )`$ $`=`$ $`C_{}|\xi |^M_{}+h_{}(\xi )\text{ for }\xi 1,`$ (3.22)
$`V_{}^{}(\xi )`$ $``$ $`0\text{ for }\xi 1,`$ (3.23)
where $`C_{}>0`$ and $`M_{}>max\{2,n\}`$. Function $`h_+(\xi )`$ has the following behavior
$$\underset{\xi \mathrm{}}{lim}\frac{h_{}^{\prime \prime \prime }(\xi )}{|\xi |^{M3}}=0.$$
(3.24)
###### Theorem 3.2.
Under conditions (3.22-3.24) on $`V_{}(\xi )`$, the equilibrium measure for the external field $`V(\xi ,t)=V_{}(\xi )+tp(\xi )`$, where $`p(\xi )`$ is a monic polynomial of an odd degree $`n`$, has no gap in its support when $`t`$ is positively large enough.
$`\mathrm{\S }`$ 3.2 The degree of polynomial $`p(\xi )`$ is even
We shall show that the equilibrium measure for the $`V_{}(\xi )+tp(\xi )`$ has a single gap in its support when $`t`$ is negatively large. When $`t`$ is positively large, there are many possibilities. The equilibrium measure has no gap when $`p(\xi )`$ is convex and may have multi-gaps when $`p(\xi )`$ is non-convex.
We will first study the case when $`t`$ is negatively large. The minimizer will be supported on $`[u_4,u_3][u_2,u_1]`$. This corresponds to $`g=1`$ in (2.5).
For simplicity, we assume that $`V_{}(\xi )`$ of (3.1) is an even function. The evenness of $`V(\xi )=t\xi ^n+V_{}(\xi )`$ implies the evenness of the equilibrium measure $`\psi (\xi )`$. Consequently, its support must be symmetric about the origin. This means that $`u_3=u_2`$ and $`u_4=u_1`$. We therefore obtain from (2.28) that
$`P(\xi ,\stackrel{}{u})`$ $`=`$ $`2(\xi ^2u_1^2)(\xi ^2u_2^2)\mathrm{\Phi }_1(\xi ,\stackrel{}{u})+(\xi +u_1)(\xi ^2u_2^2)\mathrm{\Psi }_1(u_1,\stackrel{}{u})`$
$`+`$ $`(\xi ^2u_1^2)(\xi +u_2)\mathrm{\Psi }_1(u_2,\stackrel{}{u})+(\xi ^2u_1^2)(\xi u_2)\mathrm{\Psi }_1(u_2,\stackrel{}{u})`$
$`+(\xi u_1)(\xi ^2u_2^2)\mathrm{\Psi }_1(u_1,\stackrel{}{u}){\displaystyle \frac{1}{\pi }}\xi .`$
Equations (2.27) exactly become
$`2u_1(u_1^2u_2^2)\mathrm{\Psi }_1(u_1,\stackrel{}{u}){\displaystyle \frac{1}{\pi }}u_1`$ $`=`$ $`0,`$
$`2u_2(u_2^2u_1^2)\mathrm{\Psi }_1(u_2,\stackrel{}{u}){\displaystyle \frac{1}{\pi }}u_2`$ $`=`$ $`0,`$
which are equivalent to
$`(u_1u_2)\sqrt{2(u_1+u_2)}\sqrt{{\displaystyle \frac{\mathrm{\Psi }_1(u_1,\stackrel{}{u})}{u_2}}}\sqrt{{\displaystyle \frac{1}{\pi }}}`$ $`=`$ $`0,`$ (3.25)
$`\mathrm{\Psi }_1(u_1,\stackrel{}{u})+\mathrm{\Psi }_1(u_2,\stackrel{}{u})`$ $`=`$ $`0.`$ (3.26)
We have used the following identity when deriving (3.25)
$$\mathrm{\Psi }_1(u_1,\stackrel{}{u})\mathrm{\Psi }_1(u_2,\stackrel{}{u})=2(u_1u_2)\frac{}{u_2}\mathrm{\Psi }_1(u_1,\stackrel{}{u}).$$
Its proof is the same as the one for a similar identity (3.6).
In view of the integral formula (B.10) for its left hand side, equation (3.26) is equivalent to
$$_{u_2}^{u_1}\frac{V^{}(\mu )d\mu }{\sqrt{(u_1^2\mu ^2)(\mu ^2u_2^2)}}=0.$$
(3.27)
We further assume $`V_{}`$ satisfies the condition
$$V_{}(\xi )=C|\xi |^M+h(\xi )\text{ for }|\xi |1,$$
(3.28)
where $`C>0`$ and $`M>max\{3,n\}`$. Here $`h(\xi )`$ satisfies
$$\underset{\xi \pm \mathrm{}}{lim}\frac{h^{\prime \prime \prime \prime }(\xi )}{\xi ^{M4}}=0.$$
(3.29)
This implies
$$\underset{\xi \pm \mathrm{}}{lim}\frac{h(\xi )}{\xi ^M}=0,\underset{\xi \pm \mathrm{}}{lim}\frac{h^{}(\xi )}{\xi ^{M1}}=0,\underset{\xi \pm \mathrm{}}{lim}\frac{h^{\prime \prime }(\xi )}{\xi ^{M2}}=0,\underset{\xi \pm \mathrm{}}{lim}\frac{h^{\prime \prime \prime }(\xi )}{\xi ^{M3}}=0.$$
(3.30)
We will solve equations (3.25-3.26) for large $`u_1`$ and $`u_2`$.
Again, we will first consider $`p(\xi )=\xi ^n`$.
We split $`\mathrm{\Psi }_1`$ of (3.25-3.26) into simpler terms
$$\mathrm{\Psi }_1=t\mathrm{\Psi }_{\xi ^n}+C\mathrm{\Psi }_{\xi ^M}+h(\xi ),$$
where $`\mathrm{\Psi }_{\xi ^n}`$ and $`\mathrm{\Psi }_{\xi ^M}`$ are homogeneous functions of $`(\xi ,u_1,u_2)`$ of orders $`n2`$ and $`M2`$, respectively.
Introducing
$$U_1=\frac{u_1}{|t|^{\frac{1}{Mn}}},U_2=\frac{u_2}{|t|^{\frac{1}{Mn}}},$$
(3.31)
we scale equations (3.25-3.26) as
$`(U_1U_2)\sqrt{2(U_1+U_2)}\{{\displaystyle \frac{}{U_2}}\mathrm{\Psi }_{\xi ^n}(U_1,U_1,U_2,U_2,U_1)`$
$`+C{\displaystyle \frac{}{U_2}}\mathrm{\Psi }_{\xi ^M}(U_1,U_1,U_2,U_2,U_1)`$
$`+{\displaystyle \frac{1}{|t|^{\frac{M3}{Mn}}}}{\displaystyle \frac{}{u_2}}\mathrm{\Psi }_{h(\xi )}(u_1,u_1,u_2,u_2,u_1)\}^{\frac{1}{2}}{\displaystyle \frac{1}{\sqrt{\pi }|t|^{\frac{M}{2(Mn)}}}}=0,`$
$`\mathrm{\Psi }_{\xi ^n}(U_1,U_1,U_2,U_2,U_1)+C\mathrm{\Psi }_{\xi ^M}(U_1,U_1,U_2,U_2,U_1)`$
$`+{\displaystyle \frac{1}{|t|^{\frac{M2}{Mn}}}}\mathrm{\Psi }_{h(\xi )}(u_1,u_1,u_2,u_2,u_1)\mathrm{\Psi }_{\xi ^n}(U_2,U_1,U_2,U_2,U_1)`$
$`+C\mathrm{\Psi }_{\xi ^M}(U_2,U_1,U_2,U_2,U_1)+{\displaystyle \frac{1}{|t|^{\frac{M2}{Mn}}}}\mathrm{\Psi }_{h(\xi )}(u_2,u_1,u_2,u_2,u_1)=0.`$
At $`t=\mathrm{}`$, equations (S3.Ex37-S3.Ex39) become
$`(U_1U_2)\sqrt{2(U_1+U_2)}\{{\displaystyle \frac{}{U_2}}\mathrm{\Psi }_{\xi ^n}(U_1,U_1,U_2,U_2,U_1)`$ (3.34)
$`+C{\displaystyle \frac{}{U_2}}\mathrm{\Psi }_{\xi ^M}(U_1,U_1,U_2,U_2,U_1)\}^{\frac{1}{2}}=0,`$
$`\mathrm{\Psi }_{\xi ^n}(U_1,U_1,U_2,U_2,U_1)+C\mathrm{\Psi }_{\xi ^M}(U_1,U_1,U_2,U_2,U_1)`$ (3.35)
$`\mathrm{\Psi }_{\xi ^n}(U_2,U_1,U_2,U_2,U_1)+C\mathrm{\Psi }_{\xi ^M}(U_2,U_1,U_2,U_2,U_1)=0.`$
They have a solution $`U_1=U_2=\widehat{U}`$, where $`\widehat{U}`$ is defined by
$$\mathrm{\Psi }_{\xi ^n}(\widehat{U},\widehat{U},\widehat{U},\widehat{U},\widehat{U})+C\mathrm{\Psi }_{\xi ^M}(\widehat{U},\widehat{U},\widehat{U},\widehat{U},\widehat{U})=\frac{1}{4\widehat{U}}\frac{d}{d\xi }[\xi ^n+C\xi ^M]|_{\xi =\widehat{U}}=0,$$
where we have used formula (B.9) for $`\mathrm{\Psi }_{\xi ^n}`$ and $`\mathrm{\Psi }_{\xi ^M}`$ in the first equality. We hence obtain
$$U_1=U_2=\widehat{U}=[\frac{n}{CM}]^{\frac{1}{Mn}}$$
as a solution of (3.34-3.35).
It is straight forward to use the Implicit Function Theory to determine the solution of equations (S3.Ex37-S3.Ex39) in the neighborhood of $`U_1=U_2=\widehat{U}`$, $`t=\mathrm{}`$. This in turn gives solution $`u_1(t)`$, $`u_2(t)`$ of (3.25-3.26) for $`t1`$. The solution has the asymptotics
$$u_1(t)\widehat{U}|t|^{\frac{1}{Mn}},u_1(t)\widehat{U}|t|^{\frac{1}{Mn}}\text{for }t1.$$
We now verify inequalities (2.30-2.31). We will first prove (2.30) and the last two inequalities of (2.31). It suffices to show that $`\mathrm{\Phi }_1(\xi ,u_1(t),u_2(t),u_2(t),u_1(t))`$ is negative for $`\xi u_2(t)/2`$ and positive for $`\xi u_2(t)/2`$ when $`t1`$. Since $`\mathrm{\Phi }_1(\xi ,u_1(t),u_2(t),u_2(t),u_1(t))`$ is odd in $`\xi `$ in view of formula (B.8), we only need to prove its positivity for $`\xi u_2(t)/2`$.
We split $`\mathrm{\Phi }_1`$ into simpler terms and scale these terms
$`\mathrm{\Phi }_1(\xi ,u_1,u_2,u_2,u_1)`$ (3.36)
$`=`$ $`|t|^{\frac{M3}{Mn}}[\mathrm{\Phi }_{\xi ^n}(\mathrm{\Xi },U_1,U_2,U_2,U_1)+C\mathrm{\Phi }_{\xi ^M}(\mathrm{\Xi },U_1,U_2,U_2,U_1)`$
$`+{\displaystyle \frac{1}{t^{\frac{M3}{Mn}}}}\mathrm{\Phi }_{h(\xi )}(\xi ,u_1,u_2,u_2,u_1)],`$
where $`\mathrm{\Xi }=\xi /|t|^{\frac{1}{Mn}}`$. In view of condition (3.30) and the scaling (3.31), the last term in the parenthesis goes to zero uniformly for $`\widehat{U}/3\xi /|t|^{\frac{1}{Mn}}3\widehat{U}`$ as $`t\mathrm{}`$. The sum in the parenthesis then has the limit
$`\mathrm{\Phi }_{\xi ^n}(\mathrm{\Xi },\widehat{U},\widehat{U},\widehat{U},\widehat{U})+C\mathrm{\Phi }_{\xi ^M}(\mathrm{\Xi },\widehat{U},\widehat{U},\widehat{U},\widehat{U})`$
$`=`$ $`{\displaystyle \frac{\widehat{U}\frac{d}{d\mathrm{\Xi }}[\mathrm{\Xi }^n+C\mathrm{\Xi }^M]\mathrm{\Xi }\frac{d}{d\xi }[\xi ^n+C\xi ^M]|_{\xi =\widehat{U}}}{2\widehat{U}(\mathrm{\Xi }^2\widehat{U}^2)}}`$
$`=`$ $`{\displaystyle \frac{CM\mathrm{\Xi }^{n1}[\mathrm{\Xi }^{Mn}\widehat{U}^{Mn}]}{\mathrm{\Xi }^2\widehat{U}^2}},`$
where we have used (B.7) in the first equality and $`\widehat{U}=[n/CM]^{1/Mn}`$ in the last one. The limit is positive for $`\widehat{U}/3\mathrm{\Xi }3\widehat{U}`$. It then follows from equation (3.36) that $`\mathrm{\Phi }_1(\xi ,u_1(t),u_2(t),u_2(t),u_1(t))`$ is positive for $`u_2(t)/2\xi 2u_1(t)`$ when $`t1`$.
We now show that $`\mathrm{\Phi }_1(\xi ,u_1(t),u_2(t),u_2(t),u_1(t))>0`$ uniformly for $`\xi 2u_1(t)`$. We use formula (B.8) to write $`\mathrm{\Phi }_1(\xi ,u_1,u_2,u_2,u_1)`$ as
$`{\displaystyle \frac{\xi }{\sqrt{(\xi ^2u_1^2)(\xi ^2u_2^2)}}}[{\displaystyle \frac{V^{}(\xi )}{2\xi }}`$ (3.37)
$`+{\displaystyle \frac{1}{\pi }}{\displaystyle _{u_2}^{u_1}}{\displaystyle \frac{\sqrt{(\mu ^2u_1^2)(\mu ^2u_2^2)}}{(\mu ^2\xi ^2)}}{\displaystyle \frac{V^{}(\mu )d\mu }{\sqrt{(u_1^2\mu ^2)(\mu ^2u_2^2)}}}]`$
$`=`$ $`{\displaystyle \frac{\xi }{\sqrt{(\xi ^2u_1^2)(\xi ^2u_2^2)}}}[{\displaystyle \frac{V^{}(\xi )}{2\xi }}`$
$`+{\displaystyle \frac{1}{\pi }}{\displaystyle _{u_2}^{u_1}}({\displaystyle \frac{\sqrt{(\mu ^2u_1^2)(\mu ^2u_2^2)}}{\mu ^2\xi ^2}}1){\displaystyle \frac{V^{}(\mu )d\mu }{\sqrt{(u_1^2\mu ^2)(\mu ^2u_2^2)}}}],`$
where we have used (3.27) in the equality.
It is then straight forward to use the argument below (3.19) to show that the integral of (3.37) tends to zero uniformly for $`\xi 2u_1(t)`$ as $`t\mathrm{}`$.
The first term $`V^{}(\xi )=tn\xi ^{n1}+V_{}^{}(\xi )`$ in the parenthesis is bounded from below by a positive constant uniformly for $`\xi 2u_1(t)`$ when $`t1`$. This follows from the conditions (3.28-3.30) on $`V_{}(\xi )`$.
Function $`\mathrm{\Phi }_1(\xi ,u_1(t),u_2(t),u_2(t),u_1(t))`$ is therefore positive in view of (3.37). This together with the similar result for $`u_2(t)/2\xi 2u_1(t)`$ proves that it is positive uniformly for $`\xi u_2(t)/2`$ when $`t1`$.
We have therefore proved (2.30) and the last two inequalities of (2.31).
It remains to prove the rest of (2.31), i.e.,
$$_{u_2(t)}^\xi \sqrt{(\mu ^2u_1(t)^2)(\mu ^2u_2(t)^2)}\mathrm{\Phi }_1(\mu ,u_1(t),u_2(t),u_2(t),u_1(t))𝑑\mu >0$$
(3.38)
for $`|\xi |<u_2(t)`$ when $`t1`$.
First, the inequality (3.38) is valid for $`u_2(t)<\xi u_2(t)/2`$ since $`\mathrm{\Phi }_1`$ is negative for $`2u_1(t)\xi u_2(t)/2`$. It suffices to prove the inequality for $`u_2(t)/2\xi u_2(t)`$.
We use formula (B.8) to write the integral of (3.38) as
$`{\displaystyle \frac{V(\xi )V(u_2)}{2}}`$
$`+{\displaystyle \frac{1}{\pi }}{\displaystyle _{u_2}^\xi }\sqrt{(\nu ^2u_1^2)(\nu ^2u_2^2)}[{\displaystyle _{u_2}^{u_1}}{\displaystyle \frac{V^{}(\mu )d\mu }{(\mu ^2\nu ^2)\sqrt{(u_1^2\mu ^2)(\mu ^2u_2^2)}}}]\nu 𝑑\nu .`$
It is easy to use the scaling (3.31) to show that the above is positive for $`u_2(t)/2<\xi 0`$; hence, inequality (3.38) is verified for $`u_2(t)<\xi 0`$. Since $`\mathrm{\Phi }_1(\xi ,u_1,u_2,u_2,u_1)`$ is odd in $`\xi `$ in view of (B.8), the inequality (3.38) can be extended from $`u_2(t)<\xi 0`$ to $`u_2(t)<\xi u_2(t)`$. This completes the verification of inequalities (2.30-2.31).
By Theorem 2.3, the equilibrium measure for the external field $`V(\xi )=t\xi ^n+V_{}(\xi )`$ has a single gap in its support when $`t`$ is negatively large.
It is also easy to extend the result from $`p(\xi )=\xi ^n`$ to any even monic polynomial of degree $`n`$.
###### Theorem 3.3.
Suppose $`V_{}(\xi )`$ is an even function that satisfies condition (3.28-3.29). The equilibrium measure for the external field $`V(\xi ,t)=V_{}(\xi )+tp(\xi )`$, where $`p(\xi )`$ is an even monic polynomial of degree $`n`$, has a single gap in its support when $`t`$ is negatively large.
For the case that $`t`$ is positively large, we assume
$$V_{}^{\prime \prime }(\xi )0\text{for }|\xi |1.$$
(3.39)
Again, we first consider the case $`p(\xi )=\xi ^n`$
Since $`V_{}^{\prime \prime }(\xi )`$ is uniformly bounded in $`|\xi |A`$ for some large constant $`A`$, we may choose $`t`$ positively large enough so that $`V(\xi )=t\xi ^n+V_{}(\xi )`$ is a convex function for $`|\xi |ϵ_0`$, where $`ϵ_0`$ is a tiny positive number. It is well known that an everywhere convex external field has an equilibrium measure whose support is a connected finite interval . Hence, if $`V(\xi )=t\xi ^n+V_{}(\xi )`$ is also convex in $`|\xi |<ϵ_0`$ for large $`t>0`$, the corresponding equilibrium measure has no gap in its support.
It is therefore interesting to note that, even if $`V(\xi )=t\xi ^n+V_{}(\xi )`$ is never convex in a small neighborhood of $`\xi =0`$, the equilibrium measure can still be shown to have no gap in its support for large $`t>0`$.
The approach is similar to one we use in $`\mathrm{\S }`$ 3.1. We will present the proof briefly.
We still need to solve equations (3.2-3.3). Using
$$U_1=u_1t^{\frac{1}{n}},U_2=u_2t^{\frac{1}{n}},$$
we scale the equations as
$$(U_1U_2)\sqrt{\frac{}{U_2}\mathrm{\Psi }_{\xi ^n}(U_1,U_1,U_2)+\frac{1}{t^{\frac{1}{n}}}\frac{}{u_2}\mathrm{\Psi }_V_{}(u_1,u_1,u_2)}\sqrt{\frac{1}{\pi }}=0,$$
(3.40)
$`\mathrm{\Psi }_{\xi ^n}(U_1,U_1,U_2)`$ (3.41)
$`+{\displaystyle \frac{1}{t^{\frac{1}{n}}}}\mathrm{\Psi }_V_{}(u_1,u_1,u_2)+\mathrm{\Psi }_{\xi ^n}(U_2,U_1,U_2)+{\displaystyle \frac{1}{t^{\frac{1}{n}}}}\mathrm{\Psi }_V_{}(u_2,u_1,u_2)=0.`$
At $`t=+\mathrm{}`$, these equations become
$`(U_1U_2)\sqrt{{\displaystyle \frac{}{U_2}}\mathrm{\Psi }_{\xi ^n}(U_1,U_1,U_2)}\sqrt{{\displaystyle \frac{1}{\pi }}}`$ $`=`$ $`0,`$
$`\mathrm{\Psi }_{\xi ^n}(U_1,U_1,U_2)+\mathrm{\Psi }_{\xi ^n}(U_2,U_1,U_2)`$ $`=`$ $`0.`$
They have a solution $`U_1=\stackrel{~}{U},U_2=\stackrel{~}{U}`$, where $`\stackrel{~}{U}`$ is defined by
$$2\stackrel{~}{U}\sqrt{\frac{}{U_2}\mathrm{\Psi }_{\xi ^n}(\stackrel{~}{U},\stackrel{~}{U},\stackrel{~}{U})}\sqrt{\frac{1}{\pi }}=0.$$
It is not hard to get from a formula of type (B.11) for $`\mathrm{\Psi }_{\xi ^n}`$ that
$$\frac{}{U_2}\mathrm{\Psi }_{\xi ^n}(\stackrel{~}{U},\stackrel{~}{U},\stackrel{~}{U})=\frac{1}{4}\frac{(n1)!!}{(n2)!!}U^{n2}.$$
We hence obtain
$$U_1=\stackrel{~}{U},U_2=\stackrel{~}{U},\stackrel{~}{U}=[\frac{(n2)!!}{\pi (n1)!!}]^{\frac{1}{n}}.$$
as a solution of equations (3.40-3.41) at $`t=+\mathrm{}`$.
One can then use the Implicit Function Theory to show that (3.40-3.41) give $`U_1`$ and $`U_2`$ as functions of $`t`$ near $`t=+\mathrm{}`$ and
$$lim_{t+\mathrm{}}U_1(t)=\stackrel{~}{U},lim_{t+\mathrm{}}U_2(t)=\stackrel{~}{U}.$$
Therefore, equations (3.2-3.3) can be inverted to give $`u_1(t)`$ and $`u_2(t)`$, which have the asymptotics
$$u_1(t)\frac{\stackrel{~}{U}}{t^{\frac{1}{n}}},u_2(t)\frac{\stackrel{~}{U}}{t^{\frac{1}{n}}}\text{as }t1.$$
To verify the inequalities (2.30-2.31), it suffices to prove that $`\mathrm{\Phi }_0(\xi ,u_1(t),u_2(t))>0`$ uniformly for all $`\xi `$ when $`t`$ is sufficiently large.
We first decompose $`\mathrm{\Phi }_0`$
$`\mathrm{\Phi }_0(\xi ,u_1(t),u_2(t))`$ $`=`$ $`t\mathrm{\Phi }_{\xi ^n}(\xi ,u_1(t),u_2(t))+\mathrm{\Phi }_{V_{}(\xi )}(\xi ,u_1(t),u_2(t))`$
$`=`$ $`t^{\frac{2}{n}}[\mathrm{\Phi }_{\xi ^n}({\displaystyle \frac{\xi }{t^{\frac{1}{n}}}},U_1(t),U_2(t))+{\displaystyle \frac{1}{t^{\frac{2}{n}}}}\mathrm{\Phi }_{V_{}(\xi )}(\xi ,u_1(t),u_2(t))].`$
The first term in the parenthesis is bigger than a positive constant for all $`\xi `$ and large $`t`$. To see this, function $`\mathrm{\Phi }_{\xi ^n}`$, in view of formula (B.4), can be written as a multiple integral of the second derivative of $`\xi ^n`$. Hence, when $`\stackrel{~}{U}/2U_12\stackrel{~}{U}`$ and $`2\stackrel{~}{U}U_2\stackrel{~}{U}/2`$, the first term is bounded from below by a positive constant for all $`\xi /t^{1/n}`$.
To show $`\mathrm{\Phi }_0(\xi ,u_1(t),u_2(t))>0`$ for all $`\xi `$, it is enough to show that the second term $`\mathrm{\Phi }_{V_{}(\xi )}`$ is bounded from below for all $`\xi `$ and large $`t`$. To accomplish this, we deduct from condition (3.39) that $`V_{}^{\prime \prime }(\xi )`$ is bounded from below for all $`\xi `$. Function $`\mathrm{\Phi }_{V_{}(\xi )}`$, which can be written as a multiple integral of $`V_{}^{\prime \prime }`$ in view of formula (B.4), is therefore bounded from below for all $`\xi `$, $`u_1`$ and $`u_2`$.
We have therefore verified inequalities (2.30-2.31). By Theorem 2.3, the equilibrium measure for the external field $`V(\xi ,t)=V_{}(\xi )+t\xi ^n`$ is supported on a single interval when $`t1`$.
To generalize $`p(\xi )`$ from $`\xi ^n`$ to any convex polynomial $`p(\xi )`$, we notice that the lowest order term $`\xi ^m`$ in $`p(\xi )`$ must be of an even order and that its coefficient must also be positive. The previous analysis centered around $`\xi ^n`$ can be applied to $`\xi ^m`$.
###### Theorem 3.4.
Under condition (3.39), the equilibrium measure for the external field $`V(\xi ,t)=V_{}(\xi )+tp(\xi )`$, where $`p(\xi )`$ is a convex polynomial, has no gap in its support when $`t1`$.
We conclude this section by making an observation on the case of non-convex $`p(\xi )`$. Such a polynomial $`p(\xi )`$ has multiple “wells”. They will be amplified and become the “wells” of the external field $`V(\xi ,t)=V_{}(\xi )+tp(\xi )`$ as $`t`$ is positively large if $`V_{}(\xi )`$ satisfies condition (3.39). Its equilibrium measure is therefore likely to be supported on multiple disjoint intervals.
Appendix A. Algebro-Geometric Solution of the Riemann-Hilbert Problem
In Section 2, we use function theoretical methods to solve the Riemann-Hilbert problem (2.14) and (2.15). Our solution is given by formula (2.16); it is the cornerstone of Propositions 2.1 and 2.2. In this appendix, we will present yet another approach to the Riemann-Hilbert problem. We will solve it using an algebro-geometric method. More precisely, we will give another expression for the solution (2.16) and hence re-derive the formulae of Propositions 2.1 and 2.2.
The Riemann-Hilbert problem (2.14) and (2.15) has an intrinsic algebro-geometric structure. All its solutions are connected to the Riemann surface defined by the equation $`w^2=(\mu u_1)(\mu u_2)\mathrm{}(\mu u_{2g+2})`$. We choose the branch cuts along the set of $`I`$ of (2.5). It is remarkable that our algebro-geometric approach does not require the external field $`V(\xi )`$ to be an analytic function.
Our starting point is to generalize the Cauchy kernel. On a Riemann surface, there are many analogues of the Cauchy kernel. The most convenient one for the Riemann-Hilbert problem (2.14) and (2.15) is an Abelian differential of the third kind, denoted by $`K(\mu ,z)d\mu `$, with two simple poles at the points $`(z,\pm R(z,\stackrel{}{u}))`$ with residues $`\pm 1`$, respectively. Here, $`\pm R(z,\stackrel{}{u})`$ are the upper and lower sheets of the Riemann surface. This means that $`K(\mu ,z)d\mu `$ takes the form
$$K(\mu ,z)d\mu =\frac{d\mu }{R(\mu ,\stackrel{}{u})}\frac{R(z,\stackrel{}{u})}{\mu z}+\text{holomorphic terms}$$
and that its behavior, as $`\mu `$ is near the poles, is
$$K(\mu ,z)d\mu =\pm \frac{d\mu }{\mu z}+\text{regular terms}.$$
We may further require that
$$_{u_{2k+1}}^{u_{2k}}K(\mu ,z)𝑑\mu =0k=1,2,\mathrm{},g,$$
(A.1)
which is analogous to condition (2.13).
The differential $`K(\mu ,z)d\mu `$ then takes the form
$$K(\mu ,z)d\mu =\frac{d\mu }{R(\mu ,\stackrel{}{u})}\frac{R(z,\stackrel{}{u})}{\mu z}+\underset{k=1}{\overset{g}{}}\omega _k(\mu )d\mu _{u_{2k+1}}^{u_{2k}}\frac{d\xi }{R(\xi ,\stackrel{}{u})}\frac{R(z,\stackrel{}{u})}{\xi z},$$
where $`\omega _k(\mu )d\mu `$ is the basis of holomorphic differentials normalized along the intervals $`[u_{2k+1},u_{2k}]`$, $`k=1,\mathrm{},g`$.
Using the Riemann bilinear relations between $`K(\mu ,z)d\mu `$ and the differentials $`{\displaystyle \frac{P_{g,n}(\eta ,\stackrel{}{u})}{R(\eta ,\stackrel{}{u})}}d\eta `$ defined in (2.22), it is possible to reduce the above formula to the form
$$K(\mu ,z)d\mu =\frac{d\mu }{R(\mu ,\stackrel{}{u})}\frac{R(z,\stackrel{}{u})}{\mu z}+\frac{1}{2}\underset{m=1}{\overset{g}{}}\frac{\mu ^{gm}d\mu }{R(\mu ,\stackrel{}{u})}\underset{k=1}{\overset{m}{}}k\mathrm{\Gamma }_{mk}(\stackrel{}{u})\underset{p_{}}{\overset{p_+}{}}\frac{P_{g,k}(\eta ,\stackrel{}{u})}{R(\eta ,\stackrel{}{u})}𝑑\eta ,$$
(A.2)
where $`p_\pm =(z,\pm R(z,\stackrel{}{u}))`$ and $`\mathrm{\Gamma }_k(\stackrel{}{u})`$ are the coefficients of the expansion (2.25). Formula (A.2) can also be derived from the explicit form of the Bergmann kernel on the Riemann surface $`w^2=R^2(\xi ,\stackrel{}{u})`$ .
We next point out a remarkable symmetry property
$$\frac{}{z}K(\mu ,z)=\frac{}{\mu }K(z,\mu ).$$
It also follows from the Riemann bilinear relations .
We now use $`K(\mu ,z)d\mu `$ to construct a solution of the Riemann-Hilbert problem (2.14) and (2.15)
$`𝒢(z)`$ $`={\displaystyle \frac{1}{\pi }}{\displaystyle _I}V^{}(\mu )K(z,\mu )𝑑\mu +{\displaystyle \frac{\sqrt{1}}{\pi }}{\displaystyle \frac{P_{g,0}(z,\stackrel{}{u})}{R(z,\stackrel{}{u})}}`$ (A.3)
$`={\displaystyle \frac{}{z}}\left[{\displaystyle \frac{1}{\pi }}{\displaystyle _I}V(\mu )K(\mu ,z)𝑑\mu \right]+{\displaystyle \frac{\sqrt{1}}{\pi }}{\displaystyle \frac{P_{g,0}(z,\stackrel{}{u})}{R(z,\stackrel{}{u})}}.`$ (A.4)
To see this, we derive from (A.3) the boundary value of $`𝒢`$ at real $`\xi `$
$$𝒢^+(\xi )=\sqrt{1}V^{}(\xi )\frac{1}{\pi }P.V._IV^{}(\mu )K(\xi ,\mu )𝑑\mu +\frac{\sqrt{1}}{\pi }\frac{P_{g,0}(\xi )}{R(\xi ,\stackrel{}{u})}.$$
When $`\xi I`$ and $`\mu I`$, the kernel $`K(\xi ,\mu )`$ is real and $`R(\xi ,\stackrel{}{u})`$ is pure imaginary. We then derive $`Im𝒢^+(\xi )=V^{}(\xi )`$, which is (2.14). When $`\xi \backslash \overline{I}`$ and $`\mu I`$, the kernel $`K(\xi ,\mu )`$ is pure imaginary and $`R(\xi ,\stackrel{}{u})`$ is real. We instead have $`Re𝒢^+(\xi )=0`$, which is exactly (2.15).
We further claim that $`𝒢`$ of (A.3) equals the function defined in (2.16). To see this, it suffices to prove that $`𝒢`$ of (A.3) satisfies both the loop conditions (2.13) and the asymptotics $`𝒢(z)=\frac{1}{\sqrt{1}\pi z}+O(1/z^2)`$ for large $`|z|`$. Conditions (2.13) are easily verified, in view of similar conditions (A.1) on $`K(\mu ,z)d\mu `$ and (2.23) on $`P_{g,0}`$. For large $`|z|`$, the kernel $`K(z,\mu )`$ behaves like $`O(\frac{1}{z^2})`$. Indeed the first term of $`K(z,\mu )`$ is proportional to $`\frac{1}{R(z,\stackrel{}{u})(z\mu )}`$, which clearly decays at least as fast as $`\frac{1}{z^2}`$ as $`|z|\mathrm{}`$. The other terms of $`K(z,\mu )`$ are of the form $`\frac{z^{gm}d\mu }{R(z,\stackrel{}{u})}=O(\frac{1}{z^{m+1}})`$ for $`|z|`$ large and $`m=1,\mathrm{},g`$. This, together with the behavior (2.22) of $`P_{g,0}`$, justifies the asymptotics $`𝒢(z)=\frac{1}{\sqrt{1}\pi z}+O(1/z^2)`$ for large $`|z|`$.
Inserting the explicit form (A.2) into the expression (A.4), we obtain
$`𝒢(z)`$ $`=`$ $`{\displaystyle \frac{}{z}}\left[{\displaystyle \frac{R(z,\stackrel{}{u})}{\pi }}{\displaystyle _I}{\displaystyle \frac{V(\mu )d\mu }{R(\mu ,\stackrel{}{u})(\mu z)}}\right]`$
$`+2\sqrt{1}{\displaystyle \underset{m=1}{\overset{g}{}}}q_{g,m}(\stackrel{}{u}){\displaystyle \underset{k=1}{\overset{m}{}}}k\mathrm{\Gamma }_{mk}(\stackrel{}{u}){\displaystyle \frac{P_{g,k}(z,\stackrel{}{u})}{R(z,\stackrel{}{u})}}+{\displaystyle \frac{\sqrt{1}}{\pi }}{\displaystyle \frac{P_{g,0}(z,\stackrel{}{u})}{R(z,\stackrel{}{u})}},`$
where $`q_{g,m}(\stackrel{}{u})`$ is given in (2.26).
The boundary value of $`𝒢(z)`$ of (S3.Ex71) on the real axis can be obtained by observing that the first term has the boundary value at real $`\xi `$
$$\frac{}{\xi }\left[2\sqrt{1}R(\xi ,\stackrel{}{u})\mathrm{\Psi }_g(\xi ,\stackrel{}{u})+\sqrt{1}V(\xi )\right],$$
where $`\mathrm{\Psi }_g`$ is as given in Proposition 2.2. The function $`\mathrm{\Psi }_g(\xi ,\stackrel{}{u})`$ is related to $`\mathrm{\Phi }_g(z,\stackrel{}{u})`$ of Proposition 2.1 by an identity
$$\mathrm{\Phi }_g(\xi ,\stackrel{}{u})=\frac{}{\xi }\mathrm{\Psi }_g(\xi ,\stackrel{}{u})+\frac{1}{2}\underset{i=1}{\overset{2g+2}{}}\frac{\mathrm{\Psi }_g(\xi ,\stackrel{}{u})\mathrm{\Psi }_g(u_i,\stackrel{}{u})}{\xi u_i}.$$
(A.6)
We hence arrive at the boundary value of $`𝒢`$ at the real $`\xi `$
$$𝒢^+(\xi )=\frac{2\sqrt{1}R^2(\xi ,\stackrel{}{u})\mathrm{\Phi }_g(\xi ,\stackrel{}{u})+\sqrt{1}Q(\xi )}{R(\xi ,\stackrel{}{u})}+\sqrt{1}V^{}(\xi ),$$
where the polynomial $`Q(\xi ,\stackrel{}{u})`$ is
$$R^2(\xi ,\stackrel{}{u})\underset{i=1}{\overset{2g+2}{}}\frac{\mathrm{\Psi }_g(u_i,\stackrel{}{u})}{\xi u_i}+2\underset{m=1}{\overset{g}{}}q_{k,m}(\stackrel{}{u})\underset{k=1}{\overset{m}{}}k\mathrm{\Gamma }_{mk}(\stackrel{}{u})P_{g,k}(\xi ,\stackrel{}{u})+\frac{1}{\pi }P_{g,0}(\xi ,\stackrel{}{u}).$$
These are equivalent to the formulae of Propositions 2.1 and 2.2.
Appendix B. Euler-Poisson-Darboux Equations
The boundary value problem for the Euler-Poisson-Darboux equations (2.18-2.20) has one and only one solution. Its solution can be constructed using those of the following simpler problem as building blocks
$`2(x_1x_2){\displaystyle \frac{^2q}{x_1x_2}}`$ $`=`$ $`{\displaystyle \frac{q}{x_1}}\rho {\displaystyle \frac{q}{x_2}},\rho >0isaconstant,`$ (B.1)
$`q(x_1,x_1)`$ $`=`$ $`g(x_1).`$ (B.2)
A simple calculation shows that the solution of (B.1-B.2) is given by the formula
$$q(x_1,x_2)=C_0_1^1\frac{g(\frac{1+\mu }{2}x_1+\frac{1\mu }{2}x_2)}{\sqrt{1\mu ^2}}(1+\mu )^{\frac{\rho 1}{2}}𝑑\mu ,$$
where
$$C_0=\frac{1}{_1^1\frac{(1+\mu )^{\frac{\rho 1}{2}}}{\sqrt{1\mu ^2}}𝑑\mu }.$$
A change of integration variable gives another formula for the solution
$$q(x_1,x_2)=C_0\left[\frac{2}{x_1x_2}\right]^{\frac{\rho 1}{2}}_{x_2}^{x_1}g(x)\frac{(xx_2)^{\frac{\rho 2}{2}}}{(x_1x)^{\frac{1}{2}}}𝑑x,$$
(B.3)
where the square root is set to be positive for $`x`$ between $`x_1`$ and $`x_2`$.
Using a solution method of , one is able to construct the solution of (2.18-2.20) using a multiple integral
$$\mathrm{\Phi }_g(\xi ,\stackrel{}{u})=M_0_1^1\mathrm{}_1^1\stackrel{~}{f}\frac{\underset{k=2}{\overset{2g+2}{}}(1+\mu _k)^{\frac{k1}{2}}}{\sqrt{_{k=1}^{2g+2}(1\mu _k)}}𝑑\mu _1\mathrm{}𝑑\mu _{2g+2},$$
(B.4)
where $`\stackrel{~}{f}`$ denotes
$$f^{(g+2)}(\frac{1+\mu _{2g+2}}{2}(\mathrm{}(\frac{1+\mu _2}{2}(\frac{1+\mu _1}{2}\xi +\frac{1\mu _1}{2}u_1)+\mathrm{})+\frac{1\mu _{2g+2}}{2}u_{2g+2})$$
and the constant $`M_0`$ is chosen so that the boundary condition (2.20) is satisfied. Here $`f^{(g+2)}`$ denotes the $`(g+2)th`$ derivative of $`f`$.
The function $`\mathrm{\Psi }_g(\xi ,\stackrel{}{u})`$ satisfies the same equations (2.18-2.19) as $`\mathrm{\Phi }_g(\xi ,\stackrel{}{u})`$; but, instead of the boundary condition (2.20), it satisfies (2.24) with one derivative lower.
We now list some of the properties concerning $`\mathrm{\Phi }_0`$, $`\mathrm{\Phi }_1`$ and $`\mathrm{\Psi }_1`$. They are useful in the calculations in Section 3.
###### Property B.1.
$`\mathrm{\Phi }_0(\xi ,u,u)`$ $`=`$ $`{\displaystyle \frac{V^{}(\xi )V^{}(u)}{2(\xi u)}},`$ (B.5)
$`\mathrm{\Phi }_0(\xi ,u_1,u_2)`$ $`=`$ $`{\displaystyle \frac{V^{}(\xi )}{2\sqrt{(\xi u_1)(\xi u_2)}}}`$ (B.6)
$`{\displaystyle \frac{1}{2\pi }}{\displaystyle _{u_2}^{u_1}}{\displaystyle \frac{V^{}(\mu )d\mu }{(\xi u)\sqrt{(u_1u)(uu_2)}}},`$
$`\mathrm{\Phi }_1(\xi ,u,u,u,u)`$ $`=`$ $`{\displaystyle \frac{uV^{}(\xi )\xi V^{}(u)}{2(\xi ^2u^2)u}},`$ (B.7)
$`\mathrm{\Phi }_1(\xi ,u_1,u_2,u_2,u_1)`$ $`=`$ $`{\displaystyle \frac{V^{}(\xi )}{2\sqrt{(\xi ^2u_1^2)(\xi ^2u_2^2)}}}`$ (B.8)
$`+{\displaystyle \frac{\xi }{\pi }}{\displaystyle _{u_2}^{u_1}}{\displaystyle \frac{V^{}(\mu )d\mu }{(\mu ^2\xi ^2)\sqrt{(u_1^2\mu ^2)(\mu ^2u_2^2)}}},`$
$`\mathrm{\Psi }_1(u,u,u,u,u)`$ $`=`$ $`{\displaystyle \frac{V^{}(u)}{4u}},`$ (B.9)
$`\mathrm{\Psi }_1(u_1,u_1,u_2,u_2,u_1)`$ $`+`$ $`\mathrm{\Psi }_1(u_2,u_1,u_2,u_2,u_1)`$ (B.10)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _{u_2}^{u_1}}{\displaystyle \frac{V^{}(\mu )d\mu }{\sqrt{(u_1^2\mu ^2)(\mu ^2u_2^2)}}}.`$
Function $`V(\xi )`$ is assumed to be an even function in (B.7), (B.8), (B.9) and (B.10).
###### Proof.
In formula (B.4), $`\mathrm{\Phi }_0`$, $`\mathrm{\Phi }_1`$ and $`\mathrm{\Psi }_1`$ are written as multiple integrals of $`V^{\prime \prime }(\xi )`$, $`V^{\prime \prime \prime }(\xi )`$ and $`V^{\prime \prime }(\xi )`$, respectively. Since the smooth function $`V(\xi )`$ can be approximated by polynomials in $`C^3(S)`$ on every compact set $`S`$ as close as possible, it suffices to prove Property B.1 when $`V(\xi )`$ is simply a polynomial.
We will rely on the contour integral formulation of $`\mathrm{\Phi }_g`$
$$\mathrm{\Phi }_g(\xi ,\stackrel{}{u})=\frac{1}{4\pi \sqrt{1}}_\gamma \frac{V^{}(\mu )d\mu }{(\mu \xi )R(\mu ,\stackrel{}{u})},$$
(B.11)
where $`\gamma `$ is a contour enclosing the point $`\xi `$ and all the cuts along $`[u_{2k},u_{2k1}]`$, $`k=1,2,\mathrm{},g+1`$. To see this , it is easy to check that the right hand side, as a function of $`\xi `$ and $`\stackrel{}{u}`$, satisfies (2.18) and (2.19). Boundary condition (2.20) is also easily verified using the Cauchy Integral Formula.
Identity (B.5) is an easy consequence of (B.11) when $`u_1=u_2=u`$.
To prove (B.6), we replace the contour $`\gamma `$ by $`\gamma ^{}`$, which still encloses the cuts $`[u_2,u_1]`$, $`k=1,2,\mathrm{},g+1`$, but excludes the point $`\xi `$, and write $`\mathrm{\Phi }_0(\xi ,u_1,u_2)`$ as
$$\frac{1}{2}Res_{\mu =\xi }\left[\frac{V^{}(\mu )}{(\mu \xi )\sqrt{(\mu u_1)(\mu u_2)}}\right]+\frac{1}{4\pi \sqrt{1}}_\gamma ^{}\frac{V^{}(\mu )d\mu }{(\mu \xi )R(\mu ,\stackrel{}{u})}.$$
The first term gives the first term on the right of (B.6). Rewriting the second term as an integral along the cut $`[u_2,u_1]`$ gives the second term of (B.6).
Identities (B.7) and (B.8) can be proved in the same way.
To prove (B.9), we use an analogous formula of (B.11) for $`\mathrm{\Psi }_1(\xi ,u_1,u_2,u_2,u_1)`$
$$\frac{1}{4\pi \sqrt{1}}_\gamma \frac{V(\mu )d\mu }{(\mu \xi )\sqrt{(\mu ^2u_1^2)(\mu ^2u_2^2)}}.$$
(B.12)
Identity (B.9) is an immediate consequence of this formula and the Cauchy Integral Formula.
To prove (B.10), we notice from the formula of type (B.4) that $`\mathrm{\Psi }_1`$ can be written as a multiple integral of $`V^{\prime \prime }(\xi )`$. Since $`V(\xi )`$ is an even function of $`\xi `$, so is $`\mathrm{\Psi }_1(\xi ,u_1,u_2,u_2,u_1)`$. This evenness allows us to modify formula (B.12) to get a new contour integral formulation for $`\mathrm{\Psi }(\xi ,u_1,u_2,u_2,u_1)`$
$$\frac{1}{4\pi \sqrt{1}}_\gamma \frac{V(\mu )\mu d\mu }{(\mu ^2\xi ^2)\sqrt{(\mu ^2u_1^2)(\mu ^2u_2^2)}}.$$
The left hand side of (B.10) then equals
$`{\displaystyle \frac{1}{4\pi \sqrt{1}}}{\displaystyle _\gamma }\left({\displaystyle \frac{1}{\mu ^2u_1^2}}+{\displaystyle \frac{1}{\mu ^2u_2^2}}\right){\displaystyle \frac{V(\mu )\mu d\mu }{\sqrt{(\mu ^2u_1^2)(\mu ^2u_2^2)}}}`$
$`=`$ $`{\displaystyle \frac{1}{4\pi \sqrt{1}}}{\displaystyle _\gamma }{\displaystyle \frac{}{\mu }}\left({\displaystyle \frac{1}{\sqrt{(\mu ^2u_1^2)(\mu ^2u_2^2)}}}\right)V(\mu )𝑑\mu `$
$`=`$ $`{\displaystyle \frac{1}{4\pi \sqrt{1}}}{\displaystyle _\gamma }{\displaystyle \frac{V^{}(\mu )d\mu }{\sqrt{(\mu ^2u_1^2)(\mu ^2u_2^2)}}}.`$
We obtain the right hand side of (B.10) by writing the last integral along the cuts $`[u_1,u_2]`$ and $`[u_2,u_1]`$ and using the oddness of $`V^{}(\mu )`$.
The proof of Property B.1 is completed.
Acknowledgments. T. G. was supported in part by MISGAM ESF Programme and ENIGMA MRTN-CT-2004-5652. F.-R. T. did part of his research while he was visiting the Courant Institute in 2003, IMS of CUHK and SISSA in 2004. F.-R. T. is grateful to these institutions for their support. F.-R. T. was also supported in part by NSF Grant DMS-0103849 and Grant DMS-0404931 and by a John Simon Guggenheim Fellowship. |
warning/0506/hep-ex0506011.html | ar5iv | text | # Decays of Baryon Resonances into 𝚲𝐊⁺, 𝚺^𝟎𝐊⁺ and 𝚺⁺𝐊^𝟎
## 1 Introduction
A quantitative approach to strong interactions at low energies can not be constructed without detailed information about the properties of strongly interacting particles. In recent years, substantial progress had been achieved in understanding the spectrum and the properties of low mass mesons; the observations of new mesons consisting of heavy quarks is one of the hottest topics in hadron physics. However, an understanding of the interaction between quarks can hardly be reached on the basis of knowing only quark–antiquark systems. Only baryons can provide information if pairs of quarks like to cluster into diquarks or if their excitation spectrum unravels the full richness of three–particle dynamics. But at present, only the low–mass baryon resonances are reasonably well established experimentally Krusche:2003ik , even though the nature of some of these states is still under discussion. In particlar there is no consensus if the Roper $`\mathrm{N}(1440)\mathrm{P}_{11}`$, the $`\mathrm{N}(1535)\mathrm{S}_{11}`$, and the $`\mathrm{\Lambda }(1405)\mathrm{S}_{01}`$ resonances should be interpreted as excited three–quark states or if they are created by conventional or chiral meson–baryon interactions Meissner:1989nn ; Krehl:1999km ; Jaffe:2003sg ; Kaiser:1995cy ; Glozman:1995tb ; CaroRamon:1999jf ; Steininger:1996xw ; Jido:2003cb . Above 1.8 GeV, data become sparse, and even the density of states is unclear Isgur:1995ei .
Photoproduction of nucleon resonances in their decay to strange particles offers attractive possibilities. In the reaction $`\gamma \mathrm{p}\mathrm{\Lambda }\mathrm{K}^+`$, only $`\mathrm{N}^{}`$ and no $`\mathrm{\Delta }^{}`$ resonances can contribute. The reactions $`\gamma \mathrm{p}\mathrm{\Sigma }^0\mathrm{K}^+`$ and $`\gamma \mathrm{p}\mathrm{\Sigma }^+\mathrm{K}^0`$ contribute to both the $`\mathrm{N}^{}`$ and $`\mathrm{\Delta }^{}`$ series, but with different couplings (due to different Clebsch–Gordan coefficients). Since baryon resonances have large widths and are often overlapping, a reduction in the number of partial waves or/and rigid constraints between different data sets facilitates greatly the task to identify the leading contributions.
Some of the ‘missing’ resonances are predicted to couple strongly to $`\mathrm{\Lambda }\mathrm{K}`$ and $`\mathrm{\Sigma }\mathrm{K}`$ Capstick:1998uh hence they may contribute significantly to these particular channels. Furthermore, $`\mathrm{\Lambda }`$’s reveal their polarisation in their decay, so that polarisation variables are accessible without use of a polarised photon beam or a polarised target. The same is true for the $`\mathrm{\Sigma }\mathrm{K}`$ reaction through the $`\mathrm{\Sigma }^0\mathrm{\Lambda }\gamma `$ decay. And, last not least, high statistics photoproduction data are now available from SAPHIR Glander:2003jw , CLAS McNabb:2003nf , and LEPS Zegers:2003ux . In this letter we present results of a partial wave analysis of these data, point out differences between the data and common features. The SAPHIR data on reaction $`\gamma \mathrm{p}\mathrm{K}^0\mathrm{\Sigma }^+`$ Lawall:2005np were finalized only after completion of most of the fits described here. They are not included in the systematic evaluation of errors but only in the final fit. The changes in pole position and helicity couplings induced by the new SAPHIR data are marginal only.
Data on meson production off nucleons have been interpreted using different approaches. Feuster and Mosel developed a unitary effective Lagrangian model and described both meson– and photon–induced reactions on the nucleon. Data involving known baryon resonances of spin $`J3/2`$ and later $`J5/2`$ were fitted and $`\gamma \mathrm{N}`$, $`\pi \mathrm{N}`$, $`\pi \pi \mathrm{N}`$, $`\eta \mathrm{N}`$ and $`\mathrm{K}\mathrm{\Lambda }`$ partial widths were extracted Feuster:1997pq ; Feuster:1998cj ; Penner:2002ma ; Penner:2002md ; Shklyar:2004dy . The data have also been interpreted by Regge–model calculations Guidal:2003qs using only $`\mathrm{K}`$ and $`\mathrm{K}^{}`$ exchanges and no $`s`$–channel resonances. Only the gross features of the data were reproduced. Two models based on similar effective Lagrangian approaches maid ; Mart:1999ed ; Janssen:2001pe gave the correct order of magnitude of the total cross section but failed to reproduce differential cross sections. A structure near 1.9 GeV was interpreted Mart:1999ed as evidence for a ‘missing’ resonance at this mass. Quantum numbers $`\mathrm{D}_{13}`$ were tentatively assigned to the structure which seemed consistent with the angular distribution and quark model prediction Capstick:1998uh . Including data on $`\mathrm{K}^0\mathrm{\Sigma }^+`$ required adding the $`\mathrm{N}(1720)\mathrm{P}_{13}`$ resonance Mart:2000jv . In a more recent analysis Mart:2003ty , the structure finds a more complex interpretation, and the $`\mathrm{D}_{13}`$ partial wave is found to be resonant at 1740 MeV and at a higher ill–defined mass. Janssen et al. Janssen:2002jg do not see the need for introducing a $`\mathrm{N}(1895)\mathrm{D}_{13}`$ resonance. Other groups Bennhold:2000id ; Julia-Diaz:2005qj find evidence for a third $`\mathrm{S}_{11}`$ resonance which has been suggested to explain an anomaly in the $`\eta `$ photoproduction cross section Rebreyend:2000se . Usov and Scholten Usov:2005wy use a coupled channel frame derived from an effective Lagrangian to fit CLAS and SAPHIR data on $`\mathrm{\Lambda }\mathrm{K}^+`$ and $`\mathrm{\Sigma }\mathrm{K}`$. Eight $`\mathrm{N}^{}`$ and 3 $`\mathrm{\Delta }^{}`$ resonances are introduced, among them two $`\mathrm{P}_{11}`$ resonances at 1520 and 1850 MeV and a $`\mathrm{\Delta }(1855)\mathrm{P}_{33}`$. The CLAS collaboration concluded that interference between several resonant states must be important in this mass range, rather than a single well–separated resonance.
The partial wave analysis presented here is based on the operator expansion method described in detail in Anisovich:2004zz . The method is very convenient to describe $`s`$–channel resonances, to calculate contributions from triangle and box diagrams and to project $`t`$– and $`u`$–channel exchange amplitudes into $`s`$–channel partial waves. The data on $`\mathrm{\Lambda }\mathrm{K}^+`$ and $`\mathrm{\Sigma }\mathrm{K}`$ were fitted jointly with data on $`\pi ^0`$ and $`\eta `$ photoproduction (see Table 1). Included were the differential cross section for $`\pi ^0`$ and $`\eta `$ production from CB–ELSA Bartholomy:04 ; Crede:04 , the Mainz–TAPS data Krusche:nv on $`\eta `$ photoproduction, cross sections for $`\pi ^0`$ and $`\eta `$ photoproduction from GRAAL and beam asymmetry measurements GRAAL1 ; SAID1 ; GRAAL2 , and data on $`\gamma \mathrm{p}\mathrm{n}\pi ^+`$ SAID2 . The different data sets enter the fits with weights which are listed in the fifth column of Table 1. The fits minimise a pseudo–chisquare function
$`\chi _{\mathrm{tot}}^2={\displaystyle \frac{w_i\chi _i^2}{w_iN_i}}{\displaystyle N_i}.`$ (1)
where the $`N_i`$ are given as $`N_{\mathrm{data}}`$ (per channel) in the second and the weights in the last column of Table 1.
The partial wave solutions based on $`\pi ^0`$ and $`\eta `$ photoproduction data only were presented earlier in the two letter publications Bartholomy:04 ; Crede:04 . Aspects of the new fits related to the $`\pi ^0`$ and $`\eta `$ photoproduction are presented in the preceding paper anis . In this paper, those results are discussed which pertain to final states with open strangeness.
## 2 Photoproduction of open strangeness
### 2.1 $`𝚲𝐊^\mathbf{+}`$photoproduction
The CLAS data cover the mass ($`\sqrt{s}`$) range from the $`\mathrm{K}\mathrm{\Lambda }`$ threshold to 2.4 GeV. The differential angular distributions are given in 56 bins about 13 MeV wide at low energies and 10 MeV wide at high energies. The SAPHIR collaboration showed 36 angular distributions from threshold to 2.4 GeV in about 20 MeV wide mass bins. The CLAS and SAPHIR data are complemented by the coincident observation of the $`\mathrm{\Lambda }`$ recoil polarisation. The LEPS collaboration at SPring-8 published beam asymmetry measurements in 9 mass bins covering the same mass range.
The experimental total cross section was determined by summation of the experimental differential cross sections and, in regions where no data exist, predicted values determined by the fit. For some very forward SAPHIR data points, fit values were taken as well.
The SAPHIR and CLAS $`\gamma \mathrm{p}\mathrm{\Lambda }\mathrm{K}^+`$ total cross sections show a very narrow enhancement in the 1700 MeV region. It can easily be fitted by giving the $`\mathrm{N}(1650)\mathrm{S}_{11}`$ resonance a narrow ($`30`$ MeV) width or by introducing a new $`\mathrm{S}_{11}`$ or $`\mathrm{P}_{11}`$ resonance (where the former resonance gives a slightly better description). If such a state exists, it couples strongly only to the $`\mathrm{\Lambda }\mathrm{K}^+`$ channel and must have a very exotic nature. While the narrow peaks seem to be consistent in the SAPHIR and CLAS total cross sections, their origin is very different. In the SAPHIR data, the peak is connected with a larger differential cross section (compared to the fit) over a broad angular range; in the CLAS data, the peak originates from just two points in the very forward region. If these points are excluded, the CLAS total cross section is even smaller than given by the fit. Of course, a new narrow state should not be claimed on this basis; new data are needed to resolve this discrepancy and the narrow structure at 1700 MeV is disregarded here.
The Figs. 3 and 4 show the differential cross sections obtained by SAPHIR and CLAS and the results of the best fit described below. The agreement is rather good in both cases: most discrepancies between the two experiments can obviously be ascribed to an overall energy–dependent normalisation error.
The beam polarisation asymmetries are compared to the fit in Fig. 6. For these data, the calculations were made for narrow bins (with ten times smaller widths) and then averaged to fit the correspondent experimental data. The recoil polarisation, obtained from the weak–decay asymmetry of hyperons, provides further constraints to the solution. Data and fit, divided into 24 mass bins, are shown in Fig. 6.
### 2.2 $`\mathrm{\Sigma }\mathrm{K}`$ photoproduction
The $`\gamma \mathrm{p}\mathrm{\Sigma }\mathrm{K}`$ data obtained by CLAS and SAPHIR cover the same energy range as the $`\gamma \mathrm{p}\mathrm{\Lambda }\mathrm{K}^+`$ data. The data are also complemented by the coincident observation of the $`\mathrm{\Sigma }`$ recoil polarisation. The beam asymmetry was measured by LEPS and given in 9 energy intervals covering the same mass range as the $`\mathrm{\Lambda }\mathrm{K}^+`$ measurements. The CLAS data for $`\mathrm{\Sigma }^+\mathrm{K}^0`$ photoproduction carnahan are binned into 6 energy intervals of about 100 MeV width. This data is also used in the fits. The calculations were made for 10 MeV mass spacings and then averaged to fit the experimental bins.
As in the case of the $`\mathrm{\Lambda }\mathrm{K}^+`$ reactions, the total cross section measurements from SAPHIR and CLAS are not fully compatible. In Fig. 7 the different height of the total cross sections as obtained by the two experiments can be seen. Both data agree only after renormalisation with an energy dependent function as given in Fig. 2.
The SAPHIR and CLAS differential cross sections $`\mathrm{\Sigma }^0\mathrm{K}^+`$ and the result of the best fit are shown in Figs. 9 and 9. The agreement is rather good in both cases: as for $`\mathrm{\Lambda }\mathrm{K}^+`$ production, the discrepancy can be ascribed to an overall energy–dependent normalisation error. The beam polarisation asymmetries from LEPS are given in Fig. 11 and the $`\mathrm{\Sigma }`$ recoil polarisation data divided into 23 energy bins are shown in Fig. 11.
The total and differential cross sections $`\mathrm{\Sigma }^+\mathrm{K}^0`$ versus our fit are given in Figs. 13. In this case, no normalisation factor was applied.
### 2.3 The normalisation function
Finally, a comment seems appropriate concerning the inconsistency between total cross section obtained by the SAPHIR and CLAS collaborations. If the problem is connected with an error in the photon flux normalisation, the normalisation factor should not depend on energy. A fit with an energy independent factor increased $`\chi ^2`$ for $`\gamma \mathrm{p}\mathrm{\Lambda }\mathrm{K}^+`$ by 16 for the SAPHIR differential cross sections, and by 40 for the CLAS differential cross sections. The recoil polarisation data were described by a curve with an increase of $`\chi ^2`$ by 20. Such small differences are not seen in the pictures, and we conclude that the energy dependence does not play any critical role. The total $`\chi ^2`$ for $`\mathrm{\Sigma }\mathrm{K}`$ final states only changed by 20. A constant normalisation factor was determined to be 0.80 for $`\gamma \mathrm{p}\mathrm{\Lambda }\mathrm{K}^+`$ and 0.89 for $`\gamma \mathrm{p}\mathrm{\Sigma }\mathrm{K}`$. If all reactions were fitted with the same factor, it optimised at $`0.82\pm 0.03`$. This factor is very close to that obtained for the $`\mathrm{\Lambda }\mathrm{K}^+`$ channel, and the $`\chi ^2`$ changed only by another $`20`$. However, the calculated $`\mathrm{\Sigma }\mathrm{K}`$ total cross section was then systematically lower than the SAPHIR data points. In any case, including or excluding the energy dependence of the normalisation factor does not change any of the conclusions concerning masses or widths of baryon resonances.
## 3 Fit results
The fitting weights of the various data sets and their $`\chi ^2`$ contributions are given in Table 1. The weights of the $`\mathrm{\Lambda }\mathrm{K}^+`$ and $`\mathrm{\Sigma }\mathrm{K}`$ data were selected to obtain a good description of this data without noticeable deterioration of the fit to $`\pi `$ and $`\eta `$. The solution was carefully checked for stability. The weights of some channels were changed significantly, the resulting changes in the fit parameters were included in the final systematic errors. For example, increasing the weights from 3 to 35 for the data on beam asymmetry improved the description of this data, increased slightly the $`\chi ^2`$ for other data, but lead to very small (maximum 5 MeV) shifts in resonance positions and/or widths.
### 3.1 First fits
In first fits, all resonances seen in the analysis of $`\pi ^0`$ and $`\eta `$ photoproduction Bartholomy:04 ; Crede:04 were introduced in the fit. Unphysical solutions were obtained in some of the fits, and couplings of resonances had to be restricted. The $`\mathrm{\Sigma }\mathrm{K}`$ data have a rather pronounced peak in the 1800 MeV mass region; some solutions described the peak by a single resonance having a huge coupling to this channel. This amplitude created large interferences with other contributions. To avoid this class of solutions, we demanded that the $`\mathrm{\Lambda }\mathrm{K}^+`$ and $`\mathrm{\Sigma }\mathrm{K}`$ couplings should not exceed the couplings to the $`\pi \mathrm{p}`$ or $`\eta \mathrm{p}`$ channels by more than a factor 2. In this way, solutions were found producing an acceptable overall description of the data with $`\chi ^2`$ values very close to the best fit not having such restrictions. In the new fits, all resonances had couplings well within the boundaries. When a coupling constant fell onto a boundary value, the fit was not sensitive to this coupling.
In some particular mass regions, the fit deviated visibly from the data. These regions are now discussed in some detail. The discussion will lead to the final fit and to the results gathered in Table 2.
To estimate systematic errors, fits were performed excluding the SAPHIR or CLAS cross sections. Masses and widths of the resonances of the main solution are defined more precisely by the CLAS data; masses and widths changed by less then 4 MeV in the fit with the SAPHIR data excluded. If the CLAS cross sections were excluded, a few more significant changes of the resonance positions were found. All these uncertainties are included in the errors given in Table 2.
In Table 3 we give ratios of partial widths derived from the couplings given in Table 2. The relation between coupling constants and partial widths is given e.g. in eq. (11) of the preceding paper anis .
### 3.2 The low–mass region (1500–1750 MeV)
The cross section for the reaction $`\gamma \mathrm{p}\mathrm{\Lambda }\mathrm{K}^+`$ rises steeply above threshold and reaches a maximum value at about 1720 MeV, just 100 MeV above the $`\mathrm{\Lambda }\mathrm{K}^+`$ threshold. The peak cross section for $`\gamma \mathrm{p}\mathrm{\Sigma }\mathrm{K}`$ is reached only at about 1900 MeV. This behaviour suggests that near–threshold resonances should have strong couplings to $`\mathrm{\Lambda }\mathrm{K}^+`$. The fit gives, however, stronger subthreshold couplings to the $`\mathrm{\Sigma }\mathrm{K}`$ channel. These strong subthreshold $`\mathrm{\Sigma }\mathrm{K}`$ couplings parameterise a background which interferes near threshold with $`t`$– and $`u`$–channel exchanges.
The low–mass part is strongly influenced by the $`\mathrm{S}_{11}`$ partial wave, in particular by $`\mathrm{N}(1650)\mathrm{S}_{11}`$. As mentioned, the $`\mathrm{S}_{11}`$ partial wave is described by a two–pole four–channel $`K`$–matrix. The lower $`K`$–matrix pole can have its position in a very wide mass interval, in some fits it moved down to 1100 MeV. However, there is no visible change in the fit as long as the $`K`$–matrix pole is situated between 1200 and 1460 MeV. For $`K`$–matrix masses above 1400 MeV, the fit started to exceed the SAPHIR total cross section for $`\gamma \mathrm{p}\mathrm{\Lambda }\mathrm{K}^+`$ at $``$ 1850 MeV, but the CLAS data were described more precisely. The $`K`$–matrix pole is obviously defined only with an appreciable error, we quote $`(1440_{180}^{+40})`$ MeV. There is a strong dependence of the low mass $`K`$–matrix pole couplings on the pole position. For example when the mass of lowest $`K`$–matrix pole goes to lower values the $`\mathrm{\Sigma }\mathrm{K}`$ coupling increased very significantly providing the same overall $`\mathrm{S}_{11}`$ contribution to the $`\mathrm{\Sigma }\mathrm{K}`$ cross section.
The $`T`$–matrix poles are studied in the $`\sqrt{s}`$ complex plane. Due to the four–channel nature, the complex plane is split into 8 Riemann sheets, 4 of them are relevant for the discussion. The pole positions of the $`T`$–matrix amplitude are given in Table 2 together with squared $`K`$–matrix couplings. The analytical continuation of
$`\rho _a(s)={\displaystyle \frac{\sqrt{(s(m_\mu +m_B)^2))(s(m_\mu m_B)^2)}}{s}},`$
$`a=\pi \mathrm{N},\eta \mathrm{N},\mathrm{K}\mathrm{\Lambda },\mathrm{K}\mathrm{\Sigma }\mu =\pi ,\eta ,\mathrm{K}B=\mathrm{N},\mathrm{\Lambda },\mathrm{\Sigma }.`$ (2)
to the lower complex $`s`$–plane defines the sheet closest to the physical region above threshold, for $`Re(s)>(m_\mu +m_B)^2`$. For points on this sheet with $`Re(s)<(m_\mu +m_B)^2`$, the closest physical region is at the threshold. Let us denote this sheet as $`H`$. The sheet defined by the analytical continuation of the expression
$$\rho _a(s)=i\frac{\sqrt{((m_\mu +m_B)^2)s)(s(m_\mu m_B)^2)}}{s}$$
is closest to the physical region for $`Re(s)<(m_\mu +m_B)^2`$; this sheet is denoted as $`L`$. The first pole situated on the sheet $`HHLL`$ with respect to the $`\pi \mathrm{N}`$, $`\eta \mathrm{N}`$, $`\mathrm{K}\mathrm{\Lambda }`$, and $`\mathrm{K}\mathrm{\Sigma }`$ thresholds has a mass $`(1534i107)`$ MeV. The correspondent pole situated on the sheet $`HHHL`$ is very close in mass, $`(1518i121)`$ MeV. The second pole has a mass $`1710i105`$ MeV and is situated on the sheet $`HHHH`$. The correspondent pole on the sheet $`HHHL`$ has a mass of $`1703i115MeV`$. The thresholds situated near poles have only a weak influence on pole positions.
We found that $`T`$–matrix poles close to the physical region are very stable when the masses of the $`K`$–matrix poles are scanned in a large interval. Only a few minor changes occurred compared to the previous analysis Bartholomy:04 ; Crede:04 in which a two channel $`K`$–matrix was used. The errors given in the Table 2 are defined from a large set of solutions made under different assumptions; changes of the pole positions could even be larger than the errors quoted in the Table, though, when more channels are included.
To estimate unseen contributions from three–body final states, fits were performed with a five–channel $`K`$–matrix where the fifth channel provided an unknown inelasticity. It was parametrised as $`\pi \pi \mathrm{N}`$ phase volume. We found negligible contributions from this channel and poor convergency of the fits.
### 3.3 The intermediate mass range (1700–2200 MeV)
Small discrepancies were also seen in different distributions in the 1800–1900 MeV mass region. To resolve these discrepancies, resonances with different quantum numbers were added one by one. The most significant improvement came from a $`\mathrm{P}_{11}`$ state with mass $`(1840_{40}^{+15})`$ MeV and width $`(140_{15}^{+35})`$ MeV.
With this state included a very satisfactory description of all data sets (except for $`\mathrm{\Sigma }\mathrm{K}`$ polarisation, see below) was obtained up to 2100 MeV. The $`\chi ^2`$ was improved almost for all reactions. The fit of the recoil polarisations without this state (where the effect is clearly visible) is shown as dashed curves in Figs. 6,11. The SAPHIR data on $`\mathrm{K}^0\mathrm{\Sigma }^+`$ production were added at the end of this analysis. They are described with a $`\chi ^2=109`$ for 120 data points. The comparison of this solution with resonances having other quantum numbers but a similar mass is made in Table 4.
Fig. 15 a shows a mass scan of the resonance ($`\chi ^2`$ as a function of the assumed $`\mathrm{P}_{11}`$ mass). In the scan, the mass of the $`\mathrm{P}_{11}`$ mass was fixed to preset values while all other fit parameters were allowed to adjust newly. There is a clear $`\chi ^2`$ minimum at 1810–1850 MeV. In Fig. 15 a, the $`\chi ^2`$ changes are shown for the fit of $`\pi N`$ and $`\eta N`$ differential cross sections and beam asymmetry data. The $`\mathrm{P}_{11}`$ contribution leads to a minimum in $`\chi ^2`$ for all individual distributions, although the optimum is slightly lower for the latest GRAAL beam asymmetry data. The changes for the sum of $`\chi ^2`$ for these reactions is shown in Fig. 15 b. The correspondent pictures for the reactions with $`\mathrm{\Lambda }\mathrm{K}^+`$ and $`\mathrm{\Sigma }\mathrm{K}`$ final states are shown in Figs. 15 c,d. Here the $`\chi ^2`$ distribution as well has a minimum at 1840 MeV almost for all reactions. This fact provides strong evidence that the $`\mathrm{N}(1840)\mathrm{P}_{11}`$ is a genuine resonance and not an artefact of some data. The width is determined to be $`(140_{15}^{+35})`$ MeV. This is the minimum value when the width is determined as a function of the fitted mass: if the mass is shifted by 60 MeV from the central position the width increases almost by factor 2.
The $`\mathrm{N}(1840)\mathrm{P}_{11}`$ mass is considerably larger than that of the PDG $`\mathrm{N}(1710)\mathrm{P}_{11}`$. This is a discrepancy but, possibly, the $`\mathrm{N}(1710)\mathrm{P}_{11}`$ could be split into two resonances, an $`\mathrm{N}(1670)\mathrm{P}_{11}`$ and an $`\mathrm{N}(1840)\mathrm{P}_{11}`$. An $`\mathrm{N}(1670)\mathrm{P}_{11}`$ is predicted as member of an antidecuplet Diakonov:2003jj and evidence was reported for a narrow resonance at 1670 MeV which might have $`\mathrm{P}_{11}`$ quantum numbers Kuznetsov:2004gy . A mass scan was therefore performed searching for an additional $`\mathrm{P}_{11}`$ state, fixing the mass of $`\mathrm{P}_{11}(1840)`$ at the optimum value. The description of the CLAS differential cross sections in both final channels was slightly improved if this additional resonance had a large mass and a very wide width, but there was no clear minimum in the $`\chi ^2`$ distributions. In the analysis of $`\pi \mathrm{N}\mathrm{N}\pi `$ and $`\mathrm{N}\pi \pi `$ the $`\mathrm{P}_{11}`$ mass and width were determined to have values fully compatible with our findings Manley .
The recoil polarisation for $`\mathrm{\Sigma }\mathrm{K}`$ production still had systematic deviations between data and fit, and $`\mathrm{\Lambda }\mathrm{K}^+`$ production showed a discrepancy in the 2.2 GeV mass region. Also the recoil polarisation for $`\gamma \mathrm{p}\mathrm{\Lambda }\mathrm{K}^+`$ was not well described in the 2.2 GeV mass region. As before resonances with different quantum numbers were introduced one by one. Only adding contributions from the $`\mathrm{D}_{13}`$ wave improved the picture. Two $`\mathrm{D}_{13}`$ states needed to be introduced; the mass of the lower–mass state optimised for $`(1875\pm 25)`$ MeV and a width of $`(80\pm 20)`$ MeV. The result of the mass scan for this state is shown in Fig. 16. There is a clear minimum for the GRAAL $`\pi \mathrm{N}`$ differential cross section and for the $`\mathrm{\Sigma }\mathrm{K}`$ (SAPHIR, CLAS) differential cross sections. A rather shallow minimum is seen for the $`\mathrm{\Lambda }`$ recoil polarisation. The distribution of sum of $`\chi ^2`$ for all reactions with $`\mathrm{\Lambda }\mathrm{K}^+`$ and $`\mathrm{\Sigma }\mathrm{K}`$ final states is shown in Fig. 16 b. The other state was found at $`(2166\pm 35)`$ MeV, its width at $`(280\pm 65)`$ MeV. With this resonance, the recoil polarisation in $`\gamma \mathrm{p}\mathrm{\Lambda }\mathrm{K}^+`$ is described as well. Omitting this contribution from the fit yields the fit which is shown in Figs. 6 and 11 as dashed line. The mass scan of individual channels did not show a clear minimum for this state. The $`\chi ^2`$ for the recoil polarisation (differential cross section) show asymmetric minima, shallow on the high–mass (low–mass) side and steeper on the low–mass (high–mass) side. Only their sum gives a pronounced minimum in $`\chi ^2`$.
### 3.4 The high mass range ($`>`$ 2200 MeV)
At large photon energies, kaons are produced preferentially in forward direction. The forward peaks are well reproduced by the fit which assigns forward meson production to meson exchanges in the $`t`$–channel. In the mass range covered by CLAS and SAPHIR, the fractional contributions of the $`\mathrm{K}`$ and $`\mathrm{K}^{}`$ exchanges are 8% and 22% to the total cross section for $`\mathrm{\Lambda }\mathrm{K}^+`$ and 25% and 47% for $`\mathrm{\Sigma }\mathrm{K}`$, respectively. In the $`\mathrm{K}^0\mathrm{\Sigma }^+`$ differential cross section, there is no forward peak. The $`\mathrm{K}`$ reggeized exchange is supposed to be significantly suppressed, but the fit does not find a significant contribution from $`\mathrm{K}^{}`$ exchange neither.
In the $`u`$–channel, $`\mathrm{\Lambda }`$ and $`\mathrm{\Sigma }`$ exchanges both contribute about 10% to $`\mathrm{\Lambda }\mathrm{K}^+`$ and to $`\mathrm{\Sigma }\mathrm{K}`$.
### 3.5 Discussion
#### 3.5.1 Resonances with strong $`\mathrm{\Lambda }\mathrm{K}^+`$, $`\mathrm{\Sigma }\mathrm{K}`$ coupling
For $`\mathrm{N}(2000)\mathrm{F}_{15}`$ and $`\mathrm{N}(1870)\mathrm{D}_{13}`$ the masses were found to be about 70 MeV lower then in the previous analysis of data on $`\mathrm{p}\pi `$ and $`\mathrm{p}\eta `$ final states only Bartholomy:04 ; Crede:04 . As was pointed out in those publications, these resonances only weekly contributed to $`\mathrm{p}\pi `$ and $`\mathrm{p}\eta `$ cross sections and were helpful to describe the polarisation functions. In the present analysis we found that these states have significant couplings to $`\mathrm{\Sigma }\mathrm{K}`$ and adding the $`\mathrm{P}_{11}(1840)`$ state provided a flexibility to describe the polarisation data. The $`\mathrm{F}_{15}`$ state has now mass and width which are close to the values obtained by Höhler hoehler . Moreover, the mass and width of his solution can be used as fixed input values without changing the quality of the description.
For $`\mathrm{\Lambda }\mathrm{K}^+`$ production, the $`\mathrm{S}_{11}`$ partial wave provides $`(48\pm 5)`$% of the total cross section. The $`\mathrm{N}(1720)\mathrm{P}_{13}`$ contribution is the next strongest one with $`(19\pm 4)`$%. The newly observed states $`\mathrm{N}(1840)\mathrm{P}_{11}`$ and $`\mathrm{N}(2170)\mathrm{D}_{13}`$ contribute $`3\pm 1`$% and $`(2\pm 0.8)`$%, repectively. For $`\mathrm{\Sigma }\mathrm{K}`$ production, the $`\mathrm{S}_{11}`$ wave provides the strongest resonance contributions, $`(22\pm 6)`$%. The second strongest resonance is $`\mathrm{D}_{33}(1700)`$ with $`12\pm 4`$% while $`\mathrm{N}(1720)\mathrm{P}_{13}`$ contributes only about 1% to the cross section. The $`\mathrm{N}(1840)\mathrm{P}_{11}`$ and $`\mathrm{N}(2170)`$ $`\mathrm{D}_{13}`$ resonances contribute on the level of 7% and 1% to the $`\mathrm{\Sigma }\mathrm{K}`$ total cross section.
Although quite a large number of states contribute only weakly to the total cross sections, their amplitudes are important to describe the polarisation information. Here, the interference between large and small amplitudes can significantly change the polarisation function. The typical example is the beam asymmetry in the $`\eta `$ photoproduction reaction where a small contribution from $`\mathrm{N}(1520)\mathrm{D}_{13}`$ changes the behaviour of the polarisation function dramatically.
#### 3.5.2 Four $`\mathrm{D}_{13}`$ resonances
There are four nucleon resonances with quantum numbers $`I(J^P)=1/2(3/2)^{}`$ ($`\mathrm{D}_{13}`$). These quantum numbers can be formed with intrinsic orbital angular momentum $`L=1`$ and intrinsic spin $`S=1/2`$, or with $`S=3/2`$ and $`L=1`$ or $`L=3`$. The lowest mass state at 1520 MeV seems to be dominantly in a ($`J=3/2;L=1,S=1/2`$) state. Apart from $`\mathrm{\Delta }(1232)\mathrm{P}_{33}`$, it provides the largest contribution to the $`\mathrm{p}\pi ^0`$ photoproduction cross section. Its companion at 1700 MeV, dominantly ($`J=3/2;L=1,S=3/2`$), makes a much smaller contribution. It seems plausible that the two further states, $`\mathrm{N}(1870)\mathrm{D}_{13}`$ and $`\mathrm{N}(2170)\mathrm{D}_{13}`$, should have a dominantly ($`J=3/2;L=1,S=1/2`$) configuration. We note that the spacings in mass square are
$`\mathrm{M}^2(2166)\mathrm{M}^2(1875)=(1.19\pm 0.18)\mathrm{GeV}^2,`$ (3)
$`\mathrm{M}^2(1875)\mathrm{M}^2(1520)=(1.20\pm 0.10)\mathrm{GeV}^2,`$ (4)
which agrees with the N–Roper mass splitting
$`\mathrm{M}^2(1440)\mathrm{M}^2(938)=(1.19\pm 0.06)\mathrm{GeV}^2.`$ (5)
This leads to the conjecture that $`\mathrm{N}(1870)\mathrm{D}_{13}`$ and $`\mathrm{N}(2170)`$ $`\mathrm{D}_{13}`$ could be the first and second radial excitation of the $`\mathrm{N}(1520)\mathrm{D}_{13}`$. It was shown in Klempt:2002vp that all sequential baryon resonances in a given partial wave have mass square spacings of about $`1.142\mathrm{GeV}^2`$.
#### 3.5.3 Quark–model predictions for $`\mathrm{D}_{13}`$ resonances.
Quark–model calculations predict a very large number of states. This is exemplified here using the $`\mathrm{N}\mathrm{D}_{13}`$ resonance series and a comparison with the Bonn constituent–quark model Loring:2001kx . The lowest mass states with negative parity have $`L=1`$. The spatial wave function for $`L=1`$ has mixed symmetry. Hence, the spin–flavour wave function has to have mixed symmetry, it has to belong to a 70–plet. A mixed symmetry spin–flavour wave function can be realised for $`S=1/2`$ and $`S=3/2`$. These two states can mix, thus forming the two states belonging to the $`1\mathrm{}\omega `$ band, $`\mathrm{N}(1520)\mathrm{D}_{13}`$ and $`\mathrm{N}(1700)\mathrm{D}_{13}`$.
In the third excitation band, the total orbital angular momentum can be $`L=1`$ (with $`S=1/2,\mathrm{\hspace{0.17em}3}/2`$) or $`L=3`$ (with $`S=3/2`$). Spatial wave functions can be constructed now which are symmetric, antisymmetric, or of mixed symmetry, resulting in a total of eight different states. They do not only mix among themselves; mixing with higher configurations belonging to the fifth excitation band has to be considered as well. As a result, the Bonn model predicts $`\mathrm{N}\mathrm{D}_{13}`$ resonances at 1472 and 1622 MeV in the first excitation band and at 1918, 1988, 2146, 2170, 2190, 2223, 2231, and 2271 MeV in the third excitation band. (The quark model in Capstick:1993kb gives very similar results.) The underlined masses indicate those resonances which are dominantly spin 1/2 resonances within a 70–plet. It is conceivable that the photon couplings of these resonances are larger than those of other states. This conjecture could provide a natural explanation why $`\mathrm{N}(1870)\mathrm{D}_{13}`$ and $`\mathrm{N}(2170)\mathrm{D}_{13}`$ are observed and the other states not. However, model calculations of baryon decays do not reproduce this pattern. If the other resonances exist in addition, their discovery will require data of much higher statistics and additional polarisation data. New data from pion–induced reactions will likely be mandatory as well.
#### 3.5.4 Do four $`\mathrm{S}_{11}`$ resonances exist ?
There are claims for four nucleon resonances with quantum numbers $`\mathrm{S}_{11}`$. A survey of these results can be found in Saghai:2004pt . The observation of the two high–mass states in photoproduction seems questionable since differential cross sections in the higher mass range are not well reproduced; in pion–induced reactions, the introduction of a third and fourth $`\mathrm{S}_{11}`$ resonance at $`(1846\pm 47)`$ and $`(2113\pm 70)`$ MeV improves data description considerably Chen:2002mn . In this analysis, the two lowest mass states at 1535 and 1650 MeV are observed but we do not find any need for introducing additional $`\mathrm{S}_{11}`$ states. The pairs of resonances $`\mathrm{N}(1535)\mathrm{S}_{11}`$ and $`\mathrm{N}(1520)\mathrm{D}_{13}`$, $`\mathrm{N}(1846)\mathrm{S}_{11}`$ and $`\mathrm{N}(1870)\mathrm{D}_{13}`$, and $`\mathrm{N}(2133)\mathrm{S}_{11}`$ and $`\mathrm{N}(2170)\mathrm{D}_{13}`$, may provide reasonable spin doublets, with small spin–orbit splittings. Resonances belonging to spin triplets (or degenerate quartets) like $`\mathrm{N}(1650)\mathrm{S}_{11}`$ and $`\mathrm{N}(1700)\mathrm{D}_{13}`$ are only weakly excited and their radial excitations are not observed in the data discussed here.
## 4 Summary
Results of an analysis of hyperon photoproduction in the reactions $`\gamma \mathrm{p}\mathrm{\Lambda }\mathrm{K}^+`$ and $`\gamma \mathrm{p}\mathrm{\Sigma }\mathrm{K}`$ using data from the CLAS, LEPS and SAPHIR collaborations are presented. The data are analysed in a combined fit with data on $`\pi ^0`$ and $`\eta `$ photoproduction. The SAPHIR and the CLAS data are compatible only if a normalisation factor in the order of 0.85 is introduced. The combined fit yields results which are compatible with the results on $`\pi ^0`$ and $`\eta `$ photoproduction reported earlier by the CB–ELSA collaboration but requires to introduce new baryon resonances. In particular, a $`\mathrm{P}_{11}`$ state was observed in the region of 1840 MeV which contributes to almost all reactions. The analysis highlights the existence of four $`\mathrm{D}_{13}`$ resonances, $`\mathrm{N}(1520)\mathrm{D}_{13}`$, $`\mathrm{N}(1700)\mathrm{D}_{13}`$, $`\mathrm{N}(1870)\mathrm{D}_{13}`$, and $`\mathrm{N}(2170)\mathrm{D}_{13}`$. A comparison with the Bonn quark model suggests that the main component of their flavour wave functions all belong to 70–plets; the weakly excited $`\mathrm{N}(1700)\mathrm{D}_{13}`$ has dominantly intrinsic spin 3/2 while the other 3 resonances have mostly spin 1/2.
### Acknowledgment
We would like to thank the CB–ELSA/TAPS collaboration for numerous discussions on topics related to this work. We acknowledge financial support from the Deutsche Forschungsgemeinschaft within the SFB/TR16. The St. Petersburg group received funds from the Russian Foundation for Basic Research (grant 04-02-17091). U. Thoma thanks for an Emmy Noether grant from the DFG. A. Anisovich and A. Sarantsev acknowledge support from the Alexander von Humboldt Foundation. |
warning/0506/nlin0506061.html | ar5iv | text | # Transmitting a signal by amplitude modulation in a chaotic network.
## I General setting and purposes.
In this section we recall the main results established in CS and state the questions addressed in the present paper. Consider a set of $`N`$ nodes (or relays) connected on a graph. The link from $`j`$ to $`i`$ is denoted by $`J_{ij}`$. Links are oriented ($`J_{ij}J_{ji}`$) and signed. The sign mimics excitation/inhibition effects. These effects are obviously present in biological network but they can also exist in communication networks. For example, regulation systems exist, designed to optimize the bandwidth capacity. These systems can balance the activity of a given relay with another one, resulting in an effective excitation/inhibition. Note that, in our model, the $`J_{ij}`$’s do not depend on the state of the node, but the linear response theory accommodates this generalisation. A zero link means that there is no connection from $`j`$ to $`i`$. In the sequel the $`J_{ij}`$’s are fixed and do not evolve in time. Moreover, as argued in the introduction, we are interested in the behavior of a network having a fixed set of $`J_{ij}`$’s and we do not consider statistical results relying on averages over some probability distribution for the $`J_{ij}`$’s.
The activity of a node $`i`$ is characterized by a continuous variable $`x_i`$. It is determined by the set of signals coming from the nodes connected to $`i`$, the signal coming from $`j`$ being weighted by $`J_{ij}`$. Denote by $`u_i(t)`$ the total input received by $`i`$ at time $`t`$. This is a function of the $`x_j`$’s and of the $`J_{ij}`$’s. We assume that the activity of $`i`$ evolves according to $`x_i(t+1)=f(u_i(t))`$, where $`f(x)`$ is a sigmoidal transfer function with a slope $`g`$ (e.g. $`f(x)=tanh(gx)`$). Moreover, we consider the case where $`g>>1`$, namely the sigmoid is strongly nonlinear. Note that sigmoid transfer functions are encountered in neural networks (when the neuron activity is described in terms of frequency rates), in genetic networks (Hille function), and they may also be suitably represent amplification/saturation effects in communication protocols such as TCP/IP. This amplification/saturation effect is actually the most important characteristics in what follows.
Assume indeed that we superimpose to the “background” input $`u_i(t)`$ of the node $`i`$ a small signal $`\xi _i(t)`$. How does this signal propagates inside the network ? Because of the sigmoidal shape of the transfer functions the answer depends crucially on the activity of the nodes. Assume, for the moment and for simplicity, that the time-dependent signal $`\xi _i(t)`$ has variations substantially faster than the variations of $`u_i`$. Consider then the cases depicted in Fig. 1a,b. In the first case the signal $`\xi _i(t)`$ is amplified by $`f`$, without distortion if $`\xi _i(t)`$ is weak enough. On the other hand, it is damped and distorted by wild nonlinear effects due to saturation, in Fig. 1b. This example shows that the signal propagation in such a network must take into account the topological structure of the graph as well as the nonlinear effects. This simple remark leads one to reconsider the notion of “hub”. A hub is a relay with a strong connectivity. From a topological point of view, this is certainly a very important node. But, when considering signal propagation in networks with saturating relays, the role of a hub might be temporarily weakened if this hub is maintained in a saturated state by the global activity.
In some situations, it is possible to analyze the combined effects of topology and nonlinearity on the propagation of weak amplitude signals. This is the goal of the method developed in CS . On technical grounds one first needs to make the assumption that the global spontaneous activity is chaotic. The relevance of chaos for realistic situations may be debated (note however that chaos generically occurs in the model presented below) and we deal with this point in the discussion. But, at the present stage, our point is slightly different. On one hand we provide an example where the combined effects of topology and nonlinearity can be handled, and, on the other hand, we establish that one can study the propagation and the effects of a signal superimposed upon the chaotic background in spite of (and in fact thanks to) chaos.
To be more specific consider the following model. The input signal $`u_i(t)`$ is a function of the activity $`x_j(t)`$ of the units $`j`$ connected to $`i`$ and it is given by $`u_i(t)=_jJ_{ij}x_j(t)`$. Then the global dynamics writes:
$$𝐮(t+1)=𝐆\left[𝐮(t)\right]=𝒥.f(𝐮(t)),$$
(1)
where $`𝐮(t)=\left\{u_i(t)\right\}_{i=1}^N`$ and where we used the notation $`f(𝐮(t))=\left\{f(u_i(t))\right\}_{i=1}^N`$. $`𝒥`$ is the matrix of coupling coefficients.
The saturation of the sigmoid discussed above (fig. 1b) has the dynamical effect of producing volume contraction in the phase space. Indeed, the Jacobian matrix is given by $`DG_𝐮=𝒥\mathrm{\Lambda }(𝐮)`$, where $`\mathrm{\Lambda }`$ is a diagonal matrix with $`\mathrm{\Lambda }_{ii}(𝐮)=f^{}(u_i)`$, and its determinant is given by $`detDG_𝐮=det(𝒥)\times _{i=1}^Nf^{}(u_i)`$. The determinant has an absolute value strictly lower than $`1`$ provided that some $`u_i`$’s are strong enough (corresponding to a saturation of the corresponding unit).
Consequently, the asymptotic dynamics settle onto an attractor. This attractor is generically unique (provided that one breaks the symmetry $`𝐮𝐮`$ of the transfer function $`f(x)=tanh(gx)`$, with, e.g. a small time independent threshold added to the local field $`𝐮`$).
The dynamics (1) generically exhibits a transition to chaos by quasi-periodicity as $`g`$ increases, for suitable choices of the $`J_{ij}`$’s (e.g. independent, identically distributed random variables with scaled mean and variance IJBC ; PD ; EPL ; JP ). \[Recall however that the $`J_{ij}`$’s are fixed during the evolution and that we consider a specific realization of the $`J_{ij}`$’s (we do not average over the disorder)\]. Thus the dynamics asymptotically settles onto a chaotic attractor provided that $`g`$ is sufficiently large. The statistical properties of the dynamics on its attractor are characterized by the Sinai-Ruelle-Bowen measure $`\rho `$ (SRB) which is obtained as the weak limit of the Lebesgue measure $`\mu `$ under the dynamical evolution:
$$\rho =\underset{n+\mathrm{}}{lim}𝐆^n\mu .$$
(2)
In the following we will assume that all Lyapunov exponents are bounded away from zero (weak hyperbolicity). Then for each $`𝐮supp\rho `$, where $`supp\rho `$ is the support of $`\rho `$, there exists a splitting $`E^s(𝐮)E^u(𝐮)`$ such that $`E^u(𝐮)`$, the unstable space, is locally tangent to the attractor (the local unstable manifold) and $`E^s(𝐮)`$, the stable space, is transverse to the attractor (locally tangent to the local stable manifold). Let us emphasize that the stable and unstable spaces depend on $`𝐮`$ (while the Lyapunov exponents are $`\mu `$ almost surely constant). Let us consider a point $`𝐮`$ on the attractor and make a small perturbation $`\delta _𝐮`$. This perturbation can be decomposed as $`\delta _𝐮=\delta _𝐮^u+\delta _𝐮^s`$ where $`\delta _𝐮^uE^u(𝐮)`$ and $`\delta _𝐮^sE^s(𝐮)`$. $`\delta _𝐮^u`$ is locally amplified with an exponential rate (given by the largest positive Lyapunov exponent). On the other hand $`\delta _𝐮^s`$ is damped with an exponential speed (given by the largest negative Lyapunov exponent).
Assume now that we superimpose a weak signal upon the (chaotic) activity. For simplicity, we shall assume that the signal does not depend on the state of the system (linear response still applies in this case, but the equations (4,5) do not hold anymore). Denote by $`𝝃`$ the vector $`\left\{\xi _i\right\}_{i=1}^N`$. The new dynamical system is:
$$\stackrel{~}{𝐮}(t+1)=𝐆\left[\stackrel{~}{𝐮}(t)\right]+𝝃(t)=\stackrel{~}{𝐆}\left[\stackrel{~}{𝐮}(t)\right]$$
(3)
The weak signal $`𝝃(t)`$ may be viewed as small perturbation of the trajectories of the unperturbed system (1). Consequently it has a decomposition on the local stable and unstable space. The stable component is exponentially damped. The unstable one is amplified by the dynamics and nonlinear terms rapidly scramble and mix the signal. Consequently, it becomes soon impossible to distinguish the signal from the chaotic background.
This is the effect observed on individual trajectories. However, the situation is substantially different if one considers the average effect of the signal, the average being performed with respect to the SRB measure $`\rho `$ of the unperturbed system. It has been established in CS that the average variation of the local field $`u_i`$ under the influence of the signal is given, to the linear order, by:
$$\rho \left[\delta _{u_i}(t)\right]\stackrel{\mathrm{def}}{=}\stackrel{~}{u}_i(t)u_i(t)=\underset{\sigma =\mathrm{}}{\overset{t}{}}\chi (\sigma )𝝃(t\sigma 1),$$
(4)
where $`\chi (\sigma )`$ is the matrix :
$$\chi (\sigma )=\rho (d𝐮)D𝐆_𝐮^\sigma $$
(5)
that writes in explicit form:
$$\chi _{ij}(\sigma )=\underset{\gamma _{ij}(\sigma )}{}\underset{l=1}{\overset{\sigma }{}}J_{k_lk_{l1}}\underset{l=1}{\overset{\sigma }{}}f^{}(u_{k_{l1}}(l1)).$$
(6)
We used the shortened notation $`<>`$ for the average with respect to $`\rho `$. The sum holds on each possible path $`\gamma _{ij}(\sigma )`$, of length $`\sigma `$, connecting the unit $`k_0=j`$ to the unit $`k_\sigma =i`$, in $`\sigma `$ steps. One remarks that each path is weighted by the product of a topological contribution depending only on the weight $`J_{ij}`$ and of a dynamical contribution. Since $`f`$ is a sigmoid the weight of a path $`\gamma _{ij}(\sigma )`$ depends crucially on the state of saturation of the units $`k_0,\mathrm{},k_{\sigma 1}`$ at times $`0,\mathrm{},\sigma 1`$. In particular, if $`f^{}(u_{k_{l1}}(l1))>1`$ a signal is amplified while it is damped if $`f^{}(u_{k_{l1}}(l1))<1`$ (see Fig. 1). Consequently, though a signal has many choices for going from $`j`$ to $`i`$ in $`\sigma `$ time steps, some paths may be “better” than some others, in the sense that their contribution to $`\chi _{ij}(\sigma )`$ is higher. In particular, this contribution depends strongly on the time correlations between the levels of saturation $`f^{}(u_j),f^{}(u_{k_1}),\mathrm{},f^{}(u_i)`$ of the units $`j,k_1,\mathrm{},i`$ composing the path.
This observation leads us to several remarks.
1. The paths $`ji`$ are a priori not equivalent. They have a different weight that depends on one hand on the topological contribution and on the other hand on the nonlinearity of the transfer function (this last effect indeed does not exists for linear transfer functions).
2. The average effect of the signal, measured at time $`t`$, is a sum of a large number of contributions, resulting from the various possible paths, with time delayed versions of the signal. Assuming for example that $`𝝃`$ is periodic, the observed effect is that of a sum of waves with different amplitudes and delays. The global effect can be weak if the waves interfere in a destructive way, or strong if they interfere in a constructive way. This suggests that resonances may occur. If one computes the Fourier transform of $`\chi _{ij}(t)`$, denoted by $`\widehat{\chi }_{ij}(\omega )`$ and called complex susceptibility, one observes resonance peaks corresponding to poles in the complex plane (see CS ). An example is given in the section II.2.
3. A natural (physicists) reflex would be to seek these resonance in the Fourier transform of the correlation function $`C_{ij}(t)=u_i(t)u_j(0)u_iu_j`$ of the pair $`ij`$. Indeed, the fluctuation-dissipation theorem basically tells us that a susceptibility is (the Fourier transform of) a correlation function. This is true for dynamical systems encountered in physics, where the (microscopic) dynamics preserves the volume in the phase space (Liouville theorem). But this is no longer true in our case where the dynamics contracts the phase space volume. Actually, the complex susceptibility contains more information than the power spectrum.
Indeed, the Jacobian matrix, as we saw, can be split into 2 parts, corresponding to the action of $`𝐆`$ in the local stable and unstable space, respectively. This means that the response function (7) (the corresponding susceptibility) decomposes in a stable and an unstable part. Each part has its resonances and they can be drastically different. On the one hand it has been shown by Ruelle Ruelle that the unstable contribution is actually a correlation function (this is a generalized version of the fluctuation-dissipation theorem). Henceforth the resonances of the unstable part are contained in the power spectrum. They are called “Ruelle-Pollicott resonances” RP and they do not depend on the observable (provided the observables belong to the same suitable functional space). Practically, in our case, this means that these resonances do not depend on the pair $`ij`$ <sup>2</sup><sup>2</sup>2More precisely the pole, whose real part is the frequency value of the maximum, and whose imaginary part is the resonance width, does not depend on the pair, but the residue, corresponding of the value of maximum, depends on it (see e.g. Fig. 4). Furthermore, since $`|\widehat{C}_{ij}(\omega )|=|\widehat{C}_{ji}(\omega )|`$ the analysis of these resonances does not tell us which units excites and which unit responds (see section II.3). In this sense the analysis of the correlation does not display causal information (except the trivial property $`C_{ij}(t)=C_{ji}(t)`$). But the main drawback of correlations functions is that they do not contain all possible resonances.
Indeed, the stable part displays additional resonances, called in the following stable resonances. They are not Ruelle-Pollicott resonances. They may also depend on the pair $`ij`$. Indeed, the susceptibility of the pair $`ij`$ is in general distinct from the susceptibility of the pair $`ji`$ and they may have distinct resonances (see Fig. 5a). Note also that the corresponding response function are causal since the stable directions introduce an arrow of time (see Fig. 5b). Finally, on numerical grounds, the computation of complex susceptibilities affords a better resolution frequency than the computation of correlation functions (see section II.3).
Therefore the complex susceptibilities give us essential information about the average effect of a periodic signal, applied by some unit onto some other unit. This is investigated in some details in sections II.2,II.6.
4. The equation (6) opens in principle the possibility of inducing a response of $`i`$, by exciting $`j`$ with a suitable frequency, even if there is no direct link between the 2 units. On the contrary, there may exist a direct link between $`j`$ and $`i`$ and, in spite of this, there may not be a measurable effect if the frequency of the excitation applied to $`j`$ does not correspond to a high response of $`i`$. This enhances the effect of nonlinearities in the effective capacities of the network. This is discussed in the section II.2. The notion of “hub” is in particular revisited in the section II.4.
5. The existence of the stable part may lead to a violation of the standard wisdom stating that the characteristic time of return to equilibrium is equal to the characteristic time for mixing. Actually, stable resonances introduce additional time scales that can be relatively longer than the mixing time. (The mixing time is given by the Ruelle-Pollicott resonance that is the closest to the real axis in the complex plane). As a matter of fact, it is in principle possible to observe the average effect of a signal corresponding to a kick, over a time substantially larger than the correlation time. An example is shown in section II.3.
6. Finally, the existence of resonances opens the possibility of using amplitude modulation to transmit a signal from a specific unit to another specific one, in spite of (thanks to) chaos. The original signal is then recovered by a suitable averaging procedure. This is discussed in section II.6. Let us emphasize that the procedure suggested here is not control of chaos. We do not stabilize the dynamics on a periodic orbit by a suitable perturbation. The perturbed dynamics stays chaotic and we use some natural properties of chaos, such as mixing, to reconstruct any weak signal by a suitable average.
In the next section, we present an example, based on the model (1), supporting these claims. For this we select a specific set of $`J_{ij}`$’s, randomly drawn, and we focus on the characteristics of this particular network. We do not perform statistical averages on the distribution of the $`J_{ij}`$’s. As argued in the introduction we want indeed to provide analysis tools allowing an user to extract the characteristics of the network he is currently using. One may however ask about the genericity of this example. What happens for a different set of $`J_{ij}`$’s ? What happens if one increases the size ? The statistical behavior of this model has been widely studied in IJBC ; PD ; JP . It has been shown that chaos generically occurs provided $`g`$ is sufficiently large. The average critical value for the transition to chaos has been analytically computed in EPL ; JP . The thermodynamic limit was also fully characterized for a mean field version. From these studies and from the theoretical arguments developed above we claim that the behavior described is generic in this model. Another example of resonance curves has been produced in CS for a fully connected version of the coupling matrix.
Obviously, some features such as the resonance frequencies are specific to the choice of the $`J_{ij}`$’s. But this is precisely what interests us. Starting from a chaotic network with specific resonances, we are able to compute numerically the susceptibilities by suitably exciting the nodes. This procedure does not require an a priori knowledge about the dynamics. From the susceptibilities curve we extract the resonances of this network and we then use them for applications.
One may also ask about the genericity of the model itself. This point is examined in the discussion. At this stage we simply want to remark that most of the effects exhibited in the next section are predicted from the general theory presented above. These effects are non-intuitive and depart widely from the conventional wisdom about chaotic systems. They also open new perspectives in the study of networks. This example is therefore designed to check that these theoretical predictions can be realized in at least one example. Also, the analysis performed here can be easily reproduced in other examples or models, possibly more realistic (see the discussion).
## II Signal propagation
### II.1 Model example
The numerical simulations presented here have been performed on the following example. The number of units was fixed to $`N=9`$. The network is sparse. Each unit receives connection from exactly $`K=4`$ other units. The $`J_{ij}`$’s have been drawn at random according to a Gaussian distribution with mean zero and a variance $`\frac{J^2}{K}`$. This ensures the correct normalization of the local fields $`u_i`$ IJBC . The version of the $`𝒥`$ for which all simulations have been performed is :
$$\left[\begin{array}{ccccccccccccccccc}0& 0& 0.213& 0& 0.469& 0& 0& 0.69& 0.318& & & & & & & & \\ 1.131& 0.822& 0& 0& 0& 0& 0& 0.007& 0.301& & & & & & & & \\ 0& 0.234& 0& 0& 0.51& 0.283& 0.177& 0& 0& & & & & & & & \\ 0& 0.644& 0& 0& 0.033& 0& 1.187& 0.722& 0& & & & & & & & \\ 0& 0& 0& 0& 0.511& 0.579& 0.495& 0.269& 0& & & & & & & & \\ 0& 1.015& 0& 0& 0& 0& 1.312& 0.684& 0.365& & & & & & & & \\ 0& 0& 0.852& 0.342& 0.389& 0& 0& 0.041& 0& & & & & & & & \\ 0& 0.416& 0& 0& 0.084& 0& 0.287& 0.208& 0& & & & & & & & \\ 0& 0& 0& 0.649& 0& 0& 0.331& 0.140& 1.023& & & & & & & & \end{array}\right]$$
Coupling matrix $`𝒥`$ used in the simulations below.
(Note that the corresponding graph is not decomposable). The corresponding network is drawn in Fig. 2. Blue stars correspond to inhibitory links and red crosses to excitatory links. It is for example easy to see that the unit $`7`$ is a “hub” in the sense that it sends links to almost every units, while $`0`$, $`2`$, $`3`$ or $`5`$ send at most two links.
A small constant $`\theta _i`$ has been added to each $`u_i`$ to break down the symmetry $`𝐮𝐮`$ (i.e. $`u_i(t)=_jJ_{ij}x_j(t)+\theta _i`$).
The corresponding dynamics exhibits a transition to chaos by quasi-periodicity. For $`g=3`$ the dynamics has a strange attractor. There is one positive Lyapunov exponent ($`\lambda _1=0.153`$) and $`8`$ negative Lyapunov exponents (with $`\lambda _2=0.427`$). Hence the system is weakly hyperbolic (all Lyapunov exponent bounded away from zero). The spectrum is stable to small variations of $`g`$. The Kaplan-Yorke dimension is $`1.64`$.
### II.2 Computation of susceptibilities.
In a nutshell (see CS for more details) the computation of susceptibilities consists in perturbing the trajectories of (1) by two perturbations $`𝝃^{(1)}(t)=ϵ𝐞_j\mathrm{cos}(\omega t)`$ and $`𝝃^{(2)}(t)=ϵ𝐞_j\mathrm{sin}(\omega t)`$. If $`\stackrel{~}{𝐮}^{(1)},\stackrel{~}{𝐮}^{(2)}`$ denote the variables of the corresponding perturbed systems then one may write (for $`\omega 0`$):
$$\widehat{\chi }_{ij}(\omega )=\underset{T\mathrm{}}{lim}\frac{1}{Tϵ}\underset{t=0}{\overset{T}{}}e^{i\omega (t1)}[\stackrel{~}{u}_i^{(1)}(t)+i\stackrel{~}{u}_i^{(2)}(t)].$$
(7)
This provides a straightforward way to compute the susceptibility where most of the computing time goes into computing the orbits $`\stackrel{~}{𝐮}^{(k)}(t)`$. From a numerical point of view the precision of (7) can be improved by performing an additional average over several trajectories. This allows one to compute error bars.
Note that the average is directly performed on the perturbed trajectories, and not on the difference between the perturbed and unperturbed trajectories. This has two consequences. On one hand this avoids to iterate simultaneously the perturbed and unperturbed system, compute the difference, and renormalize it when it becomes too large, to keep only linear effects. Instead, the computation (7) includes all nonlinear effects and these are precisely these effects that permit to compute the average with respect to $`\rho `$. Indeed, the SRB measure (2) is exactly an average over a typical trajectory including all nonlinear effects such as mixing and folding.
The perturbation $`𝝃^{(i)}(t)`$ decomposes into a stable and unstable part. The unstable part is rapidly amplified by the initial condition sensitivity, then nonlinear contributions arise, leading to mixing and to an effective average on the attractor, when $`T\mathrm{}`$. This average is the Fourier transform of a correlation function at the frequency $`\omega `$. Since correlation functions decay exponentially in chaotic system <sup>3</sup><sup>3</sup>3The exponential decay can be proved if the system is uniformly hyperbolic but uniform hyperbolicity can not be checked numerically. Henceforth, one assumes that the system behaves as if it “were” uniformly hyperbolic. Note however that uniform hyperbolicity is a sufficient but not a necessary condition for exponential decay. Without exponential decay the sum (7) may diverge leaving us without linear response theory. Consequently, on practical grounds, one has to check that the sum (7) does not diverge. the sum (7) converges and gives a finite contribution to the unstable part. The stable part is damped by contraction, but since it is applied in a continuous way, one obtains an effective summation of the effects of a sinusoidal perturbation transverse to the attractor. One finally obtains a finite quantity giving the (average) response of the system to a perturbation having projections in the stable and unstable directions. When $`ϵ`$ is small this is the linear response. Note however that there is a priori no condition on $`ϵ`$ in eq. (7) <sup>4</sup><sup>4</sup>4There is in fact an hidden one: $`ϵ`$ must not be too large to ensure that the perturbed system is still chaotic. Indeed, clearly a too big $`ϵ`$ will irremediably kill the chaos and give rise to a periodic regime.. This means that $`\widehat{\chi }_{ij}(\omega )`$ corresponds to the response of the system to the perturbation, possibly including nonlinear contributions in $`ϵ`$, whenever $`ϵ`$ is too large. One has therefore to check that $`ϵ`$ is small enough to ensure that $`\widehat{\chi }_{ij}`$ does not vary when $`ϵ`$ varies on a small interval. An example is given in section II.5.
Some examples of susceptibilities computed in this way are depicted in Fig. 3 with $`ϵ=10^3`$. The computation has been done with $`T=262144`$ and $`100`$ trajectories corresponding to have $`\mathrm{26.214.400}`$ points for each $`\omega `$. The frequency sampling is $`\delta \omega =\frac{2\pi }{4096}`$.
Several remarks can be made. First, as expected, there are resonance peaks common to all pairs. For example, there is a common peak located at $`\omega =0.57`$ with the same width (corresponding to the imaginary part of the corresponding pole) Note however that the height can be different (it corresponds to the value of the residue). It is also clear from inspection of Figure 3 that there are resonance peaks common to the susceptibilities corresponding to the same emitting unit (e.g. $`0.45;0.84;1.01`$ in Fig.3a; $`0.74;1.37`$ if Fig. 3b). There are also peaks that exists only for some pairs (e.g. $`0.125;0.2;2.52`$ in Fig. 3a; $`0.12,0.74,2.83`$ in Fig. 3b; $`0.125,0.70;1.26`$ in Fig. 3 c). Also, a simple glance at Fig. 3c shows that characteristics of some of the resonance peaks (namely the frequency corresponding to the maximal response and/or the width) of a receiving unit depend on the exciting unit. For example, applying a signal to the unit $`6`$ with the frequency $`\omega =1.26`$ will induce a strong response of the unit $`5`$, while applying the same signal with the same frequency to the units $`0`$ or $`5`$ will induce a weak response (see section II.2 and Fig. 10 for more details). This suggests that a suitable filtering of the global signal arriving at $`5`$ will produce a good signal to noise ratio in the case $`65`$ while it will be poor in the case $`05`$ or $`55`$ (in this last case the unit does not “feel” the signal even if it is applied to itself, because this signal is hidden into the chaotic background). Note finally that some resonance peaks are relatively high ($`20`$) corresponding to an efficient amplification of a signal with suitable frequency.
It is also clear from these figures that the intensity of the resonance has no direct connection with the intensity or the sign of the coupling and is mainly due to nonlinear effects. For example, there is no direct connection from $`0`$ to $`3`$ or $`5`$ but nevertheless these units react strongly to a suitable signal injected at unit $`0`$. On the other hand, there certainly exists a link between the resonance curves and the topological connectivity of the node: $`7`$ is a topological “hub” that sends links to almost every units, and the response curves of the units are rather similar. On the contrary, $`0`$ sends a unique link to $`1`$. Consequently, the resonance curves for the other units correspond to “indirect” paths and they look different.
### II.3 Susceptibilities versus correlations.
As discussed above, Ruelle’s theory states that the susceptibility $`\widehat{\chi }_{ij}(\omega )=\widehat{\chi }_{ij}^s(\omega )+\widehat{\chi }_{ij}^u(\omega )`$ where $`\widehat{\chi }_{ij}^s(\omega )`$ ($`\widehat{\chi }_{ij}^u(\omega )`$) is the stable (unstable) part. Consequently, $`\widehat{\chi }_{ij}(\omega )`$ contains stable and unstable resonances. Since unstable resonances are present in the Fourier transform of the correlation function $`\widehat{C}_{ij}(\omega )`$ it is natural to compare $`\widehat{\chi }_{ij}(\omega )`$ and $`\widehat{C}_{ij}(\omega )`$. An example is given in Fig. 4. As expected, one observes common peaks but there are additional peaks in the susceptibility. Note also that the correlation curves are all similar and have the same peaks (only the value of the maxima change).
This figure calls however for an important remark. While the numerical method used for the computation of the susceptibilities allows one to have a rather high resolution in frequency ($`\delta \omega _{min}=\frac{2\pi }{4096}`$) and to detect narrow resonance peaks, the computation of the correlation function is submitted to much more stringent limitations. Indeed, it is well known that, due to the initial condition sensitivity, the maximum time resolution is (assuming an attractor with a diameter of order $`1`$) $`t_{max}=\frac{1}{\lambda _1}\mathrm{ln}(\eta )`$, where $`\lambda _1`$ is the maximum Lyapunov exponent ($`\lambda _1=0.153`$ in our case) and $`\eta `$ is the round off error ($`10^{16}`$ on a Pentium, in double precision) ER . This gives a $`t_{max}`$ of order $`240`$ corresponding to a frequency resolution $`\delta \omega _{min}\frac{2\pi }{240}=0.026`$. This is the narrowest width of the resonance peaks that one can measure. Using specific libraries in quadruple precision ($`\eta 10^{32}`$) will only divide by 2 the frequency resolution. Since this effect is due to initial condition sensitivity in the unstable directions, the stable part of the susceptibility is not subjected to these limitations. A consequence of this remark is however that by glancing at the example presented here one cannot say whether a resonance peak corresponds to a stable or an unstable resonance (except for some striking cases such as $`\omega =0.57`$). This would require a more careful investigation of the corresponding poles, but this is not necessary for the scope of the present work (see CS for a computation of the poles). Indeed, a simple glance at Fig. 3a,b,c reveals a large number of peaks and many of them are not present in the correlation curves. For the applications discussed in the following, all what we need to know is that the susceptibility contains all resonances (stable and unstable) while the correlation only contains unstable resonances.
Another observation further underlines the difference between susceptibilities and correlation functions. The fluctuation-dissipation theorem of non-equilibrium statistical physics asserts that the response function to a non-equilibrium perturbation can be expressed in terms of a correlation function. Consequently, the information about non-equilibrium relaxation is included in the equilibrium fluctuations. It follows in particular that the relaxation time towards equilibrium is equal to the decorrelation time (mixing time).
Consider now figure 5. The Fourier transform of the susceptibility $`\widehat{\chi }_{ij}(\omega )`$ is the average response $`R_{ij}(t)`$ of $`i`$ to an instantaneous kick applied to $`j`$ at $`t=0`$. The first row of Fig. 5 shows these responses for the pairs $`13`$, $`31`$ and $`55`$ as well as the corresponding time correlations. Clearly, the coherence time observed in the response is substantially longer than the correlation time. Henceforth, the time for returning to equilibrium is different, and in our case longer, than the mixing time. This shows that the (average) effect of a kick can be observed on very long time scales in spite of chaos.
One also notes that the response $`13`$ is drastically different from the response $`31`$ while correlation functions are identical (up to the symmetry $`C_{13}(t)=C_{31}(t)`$). This difference is even more striking when observing the Fourier transforms. Since $`C_{ij}(t)=C_{ji}(t)`$ the graph of $`|\widehat{C}_{13}(\omega )|`$ and $`|\widehat{C}_{31}(\omega )|`$ are identical. Consequently, observing a resonance peak in the Fourier transform of the correlation function for a pair $`ij`$ does not tell us “who excites whom”. On the other hand, the graph of $`|\widehat{\chi }_{13}(\omega )|`$ and $`|\widehat{\chi }_{31}(\omega )|`$ displays clearly different peaks. At a frequency $`\omega =2.52`$, $`3`$ excites $`1`$, but $`1`$ does not excite $`3`$. Consequently, the susceptibility provides causal informations contrary to correlation functions. The difference comes from the fact that correlation functions deal with the dynamics “on” the attractor, while susceptibilities consider perturbations on the attractor as well as transverse to the attractor. As a matter of fact, the presence of stable directions introduces an explicit arrow of time and causality.
### II.4 Revisiting the notion of “hubs”.
The resonance curves leads us to seriously revisiting the notion of hub. As indicated in Fig. 2, the node $`7`$ is a topological hub. However, its ability to propagate a weak periodic signal with frequency $`\omega `$ depends on $`\omega `$. The previous analysis leads then us to propose a notion of “effective” connectivity based on susceptibility curves. For a given frequency $`\omega `$, we plot the modulus of the susceptibility $`|\chi _{ij(\omega )}|`$ with a representation assigning to each pair $`i,j`$ a circle whose size is proportional to the modulus. Some examples are represented in Fig. 6. We clearly see in this figure that changing the frequency changes the effective network.
For example, with a frequency $`\omega =0.125`$ (Fig. 6a), the node $`1`$ has a strong ability to transmit signals towards the node $`5`$ (namely the response of this unit is high). On the contrary, nodes $`5,6`$ and …$`7`$ (the topological hub) have weak performances in signal transmission at this frequency. Moreover, one sees that $`7`$ is a bad sender and a bad receiver: in this sense it is not a hub at this frequency. With a frequency $`0.57`$ (unstable resonances) the effective network has a rather symmetric structure and basically all units respond to this excitation (however with a different amplitude). Also, some units present a strong affinity with some others, at a specific frequency. This affinity is however not completely specific: the unit $`3`$ “likes” the frequency $`\omega =0.84`$ (Fig. 6c) whatever is the unit emitting it (but the best excitation is provided by unit $`7`$). Obviously, one also checks that for frequencies that do not correspond to resonances (such as $`\omega =2.33`$ in Fig. 6f) the response is essentially inexistent whatever the pair.
Finally, this figure shows that it is basically possible to excite any unit from any other one in such a way that this unit (and possibly a few other but not all the other units) have a maximal response. This can be observed in more details in Fig. II.4. This is a matrix where the entry $`i,j`$ (receiver/sender) contains the frequency where the modulus of the susceptibility is maximum (first value) and the value of this maximum (second value). Clearly, some units are more “excitable” than others (such as $`5`$).
All these effects are due to a combination of topology and dynamics and they cannot be read in the connectivity matrix $`𝒥`$.
$$\left[\begin{array}{ccccccccccccccccc}(3.1,5.84)& (0.572,8.91)& (0.756,5.97)& (3.14,10.6)& (3.05,6.8)& (3.03,6.3)& (0.724,6.71)& (3.14,7.85)& (3.14,8.7)& & & & & & & & \\ (0.563,9.77)& (0.569,22.2)& (0.795,11)& (0.479,10.7)& (0.577,12.4)& (0.373,5.55)& (0.68,13.4)& (0.569,14.5)& (0.577,6.44)& & & & & & & & \\ (1.35,4.45)& (0.569,8.14)& (1.27,7.13)& (0.719,4.42)& (1.28,8.73)& (1.39,5.05)& (1.41,7.83)& (0.844,7.07)& (1.3,4.5)& & & & & & & & \\ (0.563,15.4)& (0.569,35.3)& (0.795,21.9)& (0.463,20.9)& (0.799,19.8)& (0.816,11.5)& (0.71,19.8)& (0.844,22.1)& (0.482,11.3)& & & & & & & & \\ (1.35,4.6)& (0.569,5.32)& (1.31,7.67)& (0.742,4.04)& (1.28,10.5)& (1.39,5.22)& (1.26,8.3)& (0.844,6.88)& (1.3,5.53)& & & & & & & & \\ (0.36,21.6)& (0.127,28.2)& (0.385,23.5)& (0.437,27.8)& (0.138,17.6)& (0.364,17.6)& (1.26,15.5)& (0.569,19.1)& (0.129,15)& & & & & & & & \\ (0.141,6)& (0.569,10.1)& (1.31,7.58)& (0.437,8.37)& (1.28,9.33)& (1.39,4.53)& (1.26,9.09)& (0.854,6.79)& (1.3,5.23)& & & & & & & & \\ (0.563,6)& (0.569,14.1)& (0.795,7.86)& (0.479,8.9)& (0.563,7.53)& (0.399,4.37)& (0.71,7.91)& (0.569,10)& (0.5,4.61)& & & & & & & & \\ (3.13,19.4)& (3.12,14.9)& (3.02,16.9)& (3.14,32)& (3.04,21.1)& (3.03,20.2)& (3.02,13.4)& (3.14,26.5)& (3.14,25.8)& & & & & & & & \end{array}\right]$$
Matrix where the entry $`i,j`$ (receiver/sender) contains the frequency where the modulus of the susceptibility is maximum (first value) and the value of this maximum (second value)
\[Extrapolating further the possibilities suggested by these figures, one may imagine to apply at node $`3`$ a superposition of signals, with amplitude modulation, but with a different carrier frequency (e.g. $`\omega _1=0.125`$ and $`\omega _2=2.33`$), such that $`5`$ and $`8`$ respond simultaneously to their own resonance frequency. However, this operation requires to be strictly in the linear response regime (small $`ϵ`$).\]
### II.5 Effect of $`ϵ`$.
A linear response theory assumes that $`ϵ`$ is small enough, so that $`\widehat{\chi }_{ij}`$, corresponding to the first order term in the $`ϵ`$ expansion of $`<\delta 𝐮>`$, is independent of $`ϵ`$. Henceforth, to check that the value $`ϵ=10^3`$ chosen in our simulations is sufficiently small, we have to verify that multiplying or dividing $`ϵ`$ by some (small factor) does not change the susceptibility. Actually, later on we will also be interested in larger values of $`ϵ`$ where the nonlinear terms in the $`ϵ`$ expansion are non negligible (see section II.7). Although this regime brings the system out of the linear response setting, it provides interesting stability properties for amplitude modulation. However, it is well known that nonlinearities change the resonance structure. Consequently, we have investigated the influence of increasing $`ϵ`$ on the susceptibilities. Some examples are depicted in Fig. 7 for the pair $`31`$
One remarks that the susceptibility is stable in the range $`[5\times 10^4;2.\times 10^3]`$ (note however that the signal is more noisy when $`ϵ`$ is weaker, explaining the larger fluctuations for $`ϵ=5\times 10^4`$). Increasing $`ϵ`$ further leads to distortions in the resonance curve (see section II.7).
### II.6 Amplitude modulation.
The existence of strong amplitude specific resonances opens up the possibility for transmitting a signal carrying information from a node to a target node, in such a way that, with a suitable filtering of the chaotic background, the initial signal can be recovered. For this, one may use amplitude modulation where the characteristic time scale for the modulation is sufficiently long. More precisely, one performs the average (7) with a sliding time window whose width is sufficiently large to have small fluctuations but remains sufficiently small compared to the characteristic time for the modulation.
To illustrate this point we have superimposed a signal with periodic amplitude modulation, $`\xi (t)=ϵcos(\omega _Mt)sin(\omega _0t)`$ where $`\omega _0`$ is a resonance frequency and $`\omega _M`$ the frequency of the amplitude modulation. Note that, in the case $`ϵ=10^3`$ the variations of the signal amplitude stay within the limits where linear response theory applies. However, for such a weak signal amplitude, the signal/noise ratio is very large and we had to perform the average (7) over a time windows of width $`T=10^6`$. This imposes strong constraints on the modulation frequency, which has to be smaller than $`\frac{\pi }{T}`$ to have a correct sampling of the signal. The simulations have been done with $`\omega _M=\frac{2\pi }{2.1123T}\mathrm{2.97.10}^6`$. The (arbitrary) non integer factor $`2.1123`$ has been introduced to avoid commensurability between the frequency corresponding to the sliding window ($`\frac{2\pi }{T}`$) and $`\omega _M`$.
We have first considered (Fig. 8) the case $`73`$ with a resonance frequency $`\omega =0.57`$ (see Fig.3a). One clearly sees that the signal can be recovered by the averaging procedure (7). Note also that one gets an effective amplification by a factor $`20`$ in agreement with the resonance curve 3a. Consequently, and contrary to conventional intuition about chaotic systems, a weak signal superimposed upon a chaotic background can be recovered provided one performs a suitable average over the chaotic dynamics. In some sense, this idea is already contained in Boltzmann’s work where macroscopic observable values are obtained by averaging over the microscopic molecular chaos.
We have then investigated the possibility of sending a signal from a unit to some target by suitably selecting the frequency. In section II.2 we have given the example of exciting $`5`$ with a frequency $`0.125`$ or $`1.26`$. In the first case, it is expected that a signal emitted from $`0`$ or $`1`$ will be correctly received by $`5`$ while the same signal emitted from $`5`$ or $`6`$ will not be distinguished from the chaotic background. This is verified in Fig. 9. The first column represents the response after performing the average (7) on the perturbed system $`\stackrel{~}{u}`$. The right column represents the same average performed on the unperturbed system (without signal) $`u`$. It is clear that the signal emitted by the units $`0,1`$ is correctly recovered by unit $`5`$ (with however some distortions) while there is no clear difference between the perturbed and unperturbed cases when the units $`5,6`$ are emitting the same signal. Note that the average corresponding to each pair have been performed for different initial conditions. This explains why the figures in the right column are different.
The case $`\omega =1.26`$ is presented in the Fig. 10
### II.7 nonlinear regime.
The main drawback of the previous examples is the weakness of the signal. In order to have an efficient elimination of the chaotic background one needs to average over a long time window, limiting de facto the modulation frequencies, and, even so, the decoded signal is not completely satisfactory. It is then reasonable to increase the amplitude of the signal. But first one has then to check that the resulting dynamics remains chaotic. Indeed, too large a signal will irretrievably “kill” the background. Even when $`ϵ`$ is weak enough so that one may still consider the signal as a perturbation, increasing the signal/noise ratio can drive the system outside the linear response regime. Indeed, as suggested in section II.5 the susceptibility curves are modified when $`ϵ`$ is larger than $`5.10^3`$. In this section we investigate this effect more carefully.
First we note that an explicit formula for nonlinear corrections have been worked out by Ruelle in NLRuelle . However, it is hardly tractable, even in the case of model (1) where the linear response has a simple form. We used then the numerical computation (7) for $`ϵ=10^2`$. In Fig. 11a,b,c we have represented the susceptibility curves for the same cases as in Fig. 3a,b,c, section II.2. One observes sharper resonance peaks. This corresponds to having poles approaching the real axis when increasing the strength of the periodic forcing. Certainly, for sufficiently large $`ϵ`$ and for specific resonant frequencies, one expects the dynamics to “lock” on a periodic orbit with the effect of “killing” the chaotic regime. This is however not yet the situation for the value of $`ϵ`$ investigated here, as verified here (see e.g. Fig. 12).
As an example, we have excited the unit $`0`$ with a frequency $`\omega =0.32`$, corresponding to a sharp resonance and with an amplitude modulation frequency $`\omega _0=2.9710^5`$ (note that this frequency is ten times higher than the previous one. We were then able to shrink the sliding window by a factor $`10`$). The response of the $`9`$ units is plotted in Fig. 12. The signal/noise ratio is substantially better. We also observe that several units respond, but the more accurate response corresponds to the unit $`5`$.
How does the perturbed evolution of this unit look like? The unperturbed and perturbed trajectories are plotted in Fig. 13. One does not see any difference. In particular the perturbed dynamics is still chaotic.
Assume now that there is a user at node $`5`$, observing the dynamics. Without filtering, he does not notice any difference between the system with and without signal. But, if he knows the carrier frequency, he is able to recover a signal emitted from the unit $`0`$ out of the chaotic background. This suggests a way to encode hidden information in a chaotic signal. There is in fact a little bit more. The same user located at $`2,3`$ or $`8`$ will not be able to recover a sufficiently good signal. In this sense, the nonlinearity allows us to send the signal to specific targets by a suitable choice of the frequency modulation.
## III Discussion.
In this paper, we have exhibited some non-intuitive aspects of networks with nonlinear relays and chaotic dynamics. These effects were predicted on the basis of general theoretical arguments, analyzing the linear response of such systems to an excitation by a signal injected at some place in the network. We have in particular argued that saturation effects in the nonlinear transfer function of the node induces the presence of resonances, the stable ones, that are not present in the correlation functions. These resonances, corresponding to the response of the system to out-of-equilibrium perturbations, can be used to produce unexpected results in chaotic systems, such as the transmission of a signal with amplitude modulation from a node to some target. Though the signal is basically weak, it can be recovered by a suitable averaging procedure, in spite of chaos. Moreover, thanks to chaos, this recovery can only be performed if the receiving user knows the carrier frequency and if he is located at the right node. Note that this transmission is robust with respect to noise, as we checked.
Furthermore, it would be most interesting to realise an implementation of this scheme on “real” experimental chaotic networks. For example, the frequency dependent averaging of eq. (7) could be implemented with a “lock-in” amplifier Libbrecht .
We have presented an example supporting these conjectures. We have briefly discussed in section II the genericity of this example. But what about the genericity of the model itself? As discussed in the introduction this model contains some essential features such as the competition of excitation/inhibition, the asymmetry of the interactions, and the saturation of the transfer functions that are basically present in biological networks or in some communication networks. These features generically produces chaos in systems like (1), provided that the nonlinearity is sufficiently large (see DNnet for a recent review). However, one does not necessarily have chaos in “real” networks. It might indeed well be that biological networks, for example, are often closer to intermittency than to chaos and closer to bifurcation points than to structurally stable hyperbolic systems (see Bak ; DNnet ). In this case, the application of the methods developed here may lead to fundamental questions such as: do we still have a linear response theory in this case ? As discussed in the paper, when one approaches a structurally unstable point (bifurcation point) the susceptibility may diverge, as it does in physical systems at a second order phase transition. Note however, that the the way the susceptibility diverges is already a crucial information. Actually, in systems undergoing a second order phase transition, the divergence occurs if one takes the thermodynamic limit, but physical systems are finite. In the same way, the divergence of our susceptibility may occur when one takes the infinite time limit, but real networks are investigated on finite times. This question deserves therefore further investigations.
Nevertheless, the simple fact that we have been able to produce one example of the effects theoretically predicted for such networks, obliges, in our opinion, the community to review some a priori. Though it is helpful to investigate the topological properties of complex networks (small world, scale free graphs and so on) it is in no way sufficient for characterizing the ability of transmission of such networks with active nodes. Moreover, the dynamics of signal propagation is not a superposition of the graph properties and of the local input/output dynamics. There is a complex nonlinear feedback between the two (as revealed for example in eq. (6)) and one has to study it as a whole. This may require the development of new tools such as the linear response theory presented here. The development of such tools opens new perspectives and suggest new applications for communication networks, as well as a better understanding of biological networks. |
warning/0506/hep-th0506146.html | ar5iv | text | # Quantum field theory on manifolds with a boundary
## 1 Introduction
Quantum field theory (QFT) can be defined by a functional integral
$$d\mu (\varphi )=𝒟\varphi \mathrm{exp}(W(\varphi ))$$
(1)
where $`W`$ is the classical action. From the point of view of formal properties (translational invariance) of such a functional integral it should not matter whether we write in it $`\varphi `$ or $`\varphi +\varphi _0`$. However, if the action is defined on a manifold with a boundary then the dependence on the boundary value of $`\varphi `$ seems to be crucial -. This means that a formal invariance under translations in function space $`\varphi \varphi +\varphi _0`$ must be broken in the definition of the functional integral in refs.-. Then, the dependence on the boundary value breaks some symmetries present in the classical action $`W`$. Such an approach to QFT disagrees with the conventional one based on the mode summation or perturbation expansion in the number $`N`$ of components or in the coupling constant. In this paper we discuss a relation between the two approaches in the framework of the functional integral. In the AntiDeSitter models of refs.- the boundary appears at the spatial infinity and coincides with the (compactified) Minkowski space. In the Euclidean version of the AntiDeSitter space (in the Poincare coordinates) the boundary can be realized as the Euclidean subspace of the hyperbolic space.
We consider a Riemannian manifold $``$ with the boundary $``$. The metric on $``$ is denoted by $`G`$ and its restriction to $``$ by $`g`$. We shall denote the coordinates on $``$ by $`X`$ and their restriction to $``$ by $`x`$; close to the boundary we write $`X=(y,x)`$. The action for a minimally coupled massless free scalar field $`\varphi `$ reads
$$W_0(\varphi )=_{}𝑑X\sqrt{G}G^{AB}_A\varphi _B\varphi (\varphi ,𝒜\varphi )$$
(2)
The non-negative bilinear form (2) is defined on a certain domain $`D(𝒜)`$ of functions. Such a bilinear form determines a self-adjoint operator $`𝒜`$ (the definition of $`𝒜`$ depends on the choice of $`D(𝒜)`$) in the Hilbert space of square integrable functions. The free (Euclidean) quantum field $`\varphi `$ is defined by $`𝒜^1`$ in the sense that the kernel of $`𝒜^1`$ (the Green function) provides a definition of the two-point correlation function of $`\varphi `$. The Green function $`𝒢`$ is a solution of the equation
$$𝒜𝒢_AG^{AB}\sqrt{G}_B𝒢=\delta $$
(3)
where $`G=detG_{AB}`$. The solution of eq.(3) is not unique . If $`𝒢^{}`$ is another solution of eq.(3) then $`𝒢^{}=𝒢+𝒮`$ where $`𝒮`$ is a solution of the equation
$$𝒜𝒮=0$$
(4)
We can determine $`𝒢`$ unambiguously imposing some additional requirements, e.g., requiring that $`𝒢=0`$ on the boundary. The various definitions of $`𝒢`$ correspond to various choices of $`D(𝒜)`$ in the definition of the bilinear form (2).
## 2 The functional measure
To the free action (2) we add a local interaction $`V`$. Now, the total action reads
$$W=W_0+W_I=_{}𝑑X\sqrt{G}G^{AB}_A\varphi _B\varphi +_{}𝑑X\sqrt{G}V(\varphi )$$
(5)
We can give a mathematical definition of the formal functional measure (1)
$$d\mu _V(\varphi )=Z_0^1d\mu _0(\varphi )\mathrm{exp}(W_I)$$
(6)
where the Gaussian measure $`\mu _0`$ is a mathematical realization of the formal integral
$$d\mu _0(\varphi )=𝒟\varphi \mathrm{exp}(W_0)$$
The partition function $`Z_0`$ in eq.(6)
$$Z_0=𝑑\mu _0\mathrm{exp}(W_I)$$
(7)
determines a normalization factor.
We do not discuss in this paper some divergence problems which may arise if $`V`$ is a local function of $`\varphi `$ and $``$ has an infinite volume. We may assume that $`W_I`$ has been properly regularized. We have a suggestion how to construct a regular QFT at the end of this paper.
A functional measure $`\mu `$ defines a probability distribution of fields $`\varphi (X)`$; more precisely the ”smeared out ” fields
$$(\varphi ,f)=𝑑X\sqrt{G}\varphi (X)f(X).$$
In probability theory (see, e.g., ) it is convenient to treat the probability measure $`\mu `$ defined on some sets of random fields $`\varphi `$ as one of many possible realizations of the probability space $`(\mathrm{\Omega },\mathrm{\Sigma },P)`$. The random field $`\varphi :\mathrm{\Omega }R`$ (at fixed $`X`$) is a map from the set $`\mathrm{\Omega }`$ to the set of real numbers such that the two-point correlation function is an average over the ”sample paths” $`\omega \mathrm{\Omega }`$
$$\varphi (X)\varphi (Y)=_\mathrm{\Omega }𝑑P(\omega )\varphi _\omega (X)\varphi _\omega (Y)$$
where $`P`$ is a probability measure on the $`\sigma `$-algebra $`\mathrm{\Sigma }`$ of subsets of $`\mathrm{\Omega }`$. The Gaussian measure gives a realization of the Gaussian random field $`\varphi _\omega `$. It is defined by the mean
$$m(X)=𝑑\mu (\varphi )\varphi (X)\varphi (X)$$
and the covariance
$$\begin{array}{c}𝒢(X,Y)=𝑑\mu (\varphi )(\varphi (X)\varphi (X))(\varphi (Y)\varphi (Y))\hfill \\ =(\varphi (X)\varphi (X))(\varphi (Y)\varphi (Y))\hfill \end{array}$$
(8)
or by its characteristic function $`S`$
$$S[if]=𝑑\mu \mathrm{exp}(i(\varphi ,f))=\mathrm{exp}(i(m,f)\frac{1}{2}(f,𝒢f))$$
Note that if we make a shift in the function space and define $`\stackrel{~}{\varphi }=\varphi m`$ then $`\stackrel{~}{\varphi }`$ has zero mean. Hence, we could subtract the mean value defining a new Gaussian measure
$$d\stackrel{~}{\mu }(\stackrel{~}{\varphi })=d\mu (\stackrel{~}{\varphi }+m)$$
The Gaussian measure is quasiinvariant under a shift $`\chi `$ if there exists an integrable function $`\rho (\varphi ,\chi )`$ such that
$$d\mu (\varphi +\chi )=d\mu (\varphi )\rho (\varphi ,\chi )$$
(9)
It is easy to see by a calculation of the characteristic function of both sides of eq.(9) that the measure $`\mu `$ is quasiinvariant under the shift $`\chi `$ if
$$\rho (\varphi ,\chi )=\mathrm{exp}((\varphi ,B\chi )\frac{1}{2}(\chi ,C\chi ))$$
(10)
and the following equations are satisfied
$$\chi =𝒢B\chi $$
(11)
$$(B\chi ,𝒢B\chi )=(\chi ,C\chi )$$
(12)
Eqs.(9)-(12) express the formal invariance of the functional measure (1) under translations in the function space. If these conditions are not satisfied then it really does matter what is the shift $`\chi `$. In some papers on AdS-CFT correspondence -, the choice is made $`𝒢(X,Y)=𝒢_D(X,Y)`$ where $`𝒢_D`$ is the Dirichlet Green function (vanishing on the boundary) and $`\varphi (X)=\varphi _0(X)`$ where $`\varphi _0(X)`$ is a solution of the equation
$$𝒜\varphi _0=0$$
(13)
with a fixed boundary condition $`\mathrm{\Phi }`$. We can see that eqs.(11)-(12) cannot be satisfied if $`\chi =\varphi _0`$. Hence, the partition function $`Z[\mathrm{\Phi }]`$ may depend on the boundary value $`\mathrm{\Phi }`$.
In general, choosing in QFT the boundary field $`\varphi _00`$ we break some symmetries of the classical action (5). As an example we could consider the hyperbolic space with the metric
$$ds^2=y^2(dy^2+dx_1^2+\mathrm{}.+dx_d^2)$$
(14)
The hyperbolic space (14) has compactified $`R^d`$ as the boundary . The hyperbolic space can be considered as an Euclidean version of AntiDeSitter space ($`AdS_{d+1}`$ has compactified Minkowski space as a boundary at conformal infinity ). It is also a Euclidean version of DeSitter space. However, the Poincare coordinates (14) are inappropriate for an analytic continuation of quantum fields from the hyperbolic space to DeSitter space (there is also no boundary at conformal infinity of DeSitter space).
The action (5) in the hyperbolic space is invariant under $`R^d`$ rotations and translations whereas the quantum field theory with a fixed boundary value of $`\varphi _0`$ would not be invariant under these symmetries. The approach to QFT assuming a boundary condition $`\mathrm{\Phi }`$ for the field $`\varphi _0`$ and the Dirichlet boundary condition for the Green function leads to a different quantum field theory than the one developed in refs.. The latter is determined by the mean $`\varphi =0`$ and a choice of the Green function (the free propagator $`𝒢`$ solving eq.(3)) which does not vanish on the boundary. A possible way to determine the propagator is to construct it for a real time by a mode summation and subsequently to continue analytically the propagator to the imaginary time (for a class of models this is done in ; the mode summation is also not unique). It seems reasonable to choose the Green function $`𝒢`$ which has the symmetries of the action $`W_0`$ (1)as in . Then, the functional measure (6) will have the symmetries of the action (5).
## 3 An average over the boundary values
After the heuristic discussion in sec.2 of functional integration over fields with a fixed boundary value we prove in this section that the approach starting form the free propagator $`𝒢`$ is equivalent to a quantization around a classical solution $`\varphi _0`$ with a prescribed boundary value $`\mathrm{\Phi }`$ if subsequently an average over all such boundary values is performed. First, let us assume (in the sense that for the bilinear forms $`(f,𝒢f)(f,𝒢_Df)`$)
$$𝒢𝒢_D$$
(15)
If the operator $`𝒜`$ is an elliptic operator then the inequality (15) follows from the maximum principle for elliptic operators . We are interested also in operators $`𝒜`$ with singular or vanishing coefficients which need not be elliptic. It is not clear whether the inequality (15) can be satisfied for such operators. However, the inequality (15) still holds true for the Green functions of singular operators discussed in which are expressed by a path integral. The Dirichlet condition imposes a restriction on the class of paths. Hence, the integral over a restricted set of paths is bounded by $`𝒢`$ in eq.(15).
If the inequality (15) is satisfied then there exists a positive definite bilinear form $`𝒢_B`$ such that
$$𝒢(X,X^{})=G_D(X,X^{})+𝒢_B(X,X^{})$$
(16)
Clearly on the boundary
$$𝒢(0,x;0,x^{})=𝒢_B(0,x;0,x^{})𝒢_E(x,x^{})$$
(17)
$`𝒢_E`$ defines a non-negative bilinear form on the set of functions defined on the boundary $``$.
Theorem 1
Let $`\mu _0`$ be the Gaussian measure with the mean zero and the covariance $`𝒢`$. Assume that $`𝒢`$ and $`𝒢_D`$ are real positive definite bilinear forms satisfying the inequality (15). Then, there exist independent Gaussian random fields $`\varphi _D`$ and $`\varphi _B`$ with the mean equal zero and the covariance $`𝒢_D`$ and $`𝒢_B`$ resp. such that for any integrable function $`\mathrm{exp}(W_I)F`$
$$\begin{array}{c}𝑑\mu _0(\varphi )\mathrm{exp}(W_I(\varphi ))F(\varphi )\hfill \\ =𝑑\mu _D(\varphi _D)𝑑\mu _B(\varphi _B)\mathrm{exp}(W_I(\varphi _D+\varphi _B))F(\varphi _D+\varphi _B)\hfill \end{array}$$
(18)
In this sense
$$\varphi =\varphi _D+\varphi _B$$
(19)
The theorem and its proof can be found in . It is easy to check eq.(18) for the generating functional ( then $`\mathrm{exp}(W_I(\varphi ))F(\varphi )=\mathrm{exp}(\varphi ,J)`$). On a perturbative level the general formula (18) follows from the one for the generating functional. For the general theory of ”conditioning” (15) see . Eq.(18) is discussed in the lattice approximation in (sec.8.1). Another derivation and a discussion of its relevance to the AdS-CFT correspondence can be found in . The relevance of an average over the boundary values for the Hamiltonian formulation of the quantum field theory is discussed in .
Let us note that on a formal level
$$d\mu _D(\varphi _D)=𝒟\varphi _D\mathrm{exp}(\frac{1}{2}(\varphi _D,𝒜_D\varphi _D))$$
where $`𝒜_D`$ is the Laplace-Beltrami operator with the Dirichlet boundary conditions. On a formal level $`𝒜_D\varphi _0=0`$. Hence, $`d\mu _D(\varphi _D+\varphi _0)=d\mu _D(\varphi _D)`$ although strictly speaking the shift of $`\mu _D`$ by $`\varphi _0`$ does not make sense because $`\varphi _0`$ does not vanish on the boundary. We treat the r.h.s. of eq.(18) (before an integration over $`\varphi _B`$) as a rigorous version of the QFT shifted by a classical solution. This interpretation is suggested by
Theorem 2
Let $`𝒢_D`$ be the Dirichlet Green function of the operator $`𝒜`$ (eq.(3)). Let $`𝒢`$ be another real solution of eq.(3) satisfying the inequality (15). Then, there exists a Gaussian random field $`\mathrm{\Phi }`$ defined on the boundary $``$ with the mean zero and the covariance $`𝒢_E`$ such that ( in the sense of $`L^2(dP)`$ integrals )
$$\varphi _B(X)=_{}𝑑x_b\sqrt{g}𝒟(X,x_b)\mathrm{\Phi }(x_b)$$
(20)
where $`𝒟(X,x_b)`$ is the Green function solving the boundary value problem for eq.(13).
Proof: Let us note that $`𝒢_D`$ as well as $`𝒢`$ satisfy the same equation (3). Then, their difference $`𝒢_B=𝒢𝒢_D`$ satisfies the equations
$$𝒜(X)𝒢_B(X,X^{})=𝒜(X^{})𝒢_B(X,X^{})=0$$
(21)
and the boundary condition $`𝒢_B(0,x;0,x^{})=𝒢_E(x,x^{})`$. We can solve eq.(21) with the given boundary condition $`𝒢_E`$
$$𝒢_B(X,X^{})=_{}𝑑x_b\sqrt{g}𝒟(X,x_b)_{}𝑑x_b^{}\sqrt{g}𝒟(X^{},x_b^{})𝒢_E(x_b,x_b^{})$$
(22)
where $`𝒟`$ is the Green function solving the Dirichlet boundary problem for eq.(13). The bilinear form $`𝒢_E`$ (17) defines a Gaussian field $`\mathrm{\Phi }`$ on the probability space $`(\mathrm{\Omega },\mathrm{\Sigma },P)`$ (see sec.2) with the mean zero and the covariance
$$\mathrm{\Phi }(x)\mathrm{\Phi }(x^{})=𝒢_E(x,x^{})$$
(23)
We define $`\stackrel{~}{\varphi }_B`$ by the r.h.s. of eq.(20) where the integral can be understood in the sense of the $`L^2(dP)`$ convergence of the Riemann sums (see, e.g.,). For the proof of the theorem ($`\stackrel{~}{\varphi }=\varphi `$) it is sufficient to show that the covariance of $`\stackrel{~}{\varphi }_B`$ coincides (as a bilinear form) with $`𝒢_B`$. This is a consequence of eq.(22).
Let us note that there exists the Gaussian measure $`\nu _B`$ such that
$$𝑑\nu _B(\mathrm{\Phi })\mathrm{\Phi }(x)\mathrm{\Phi }(x^{})=𝒢_E(x,x^{})$$
$`\nu _B`$ can be defined by $`\mu _B`$ as $`\nu _B=\mu _BT`$ where $`T(\mathrm{\Phi })=\varphi _B`$ is the one to one map (20) expressing the solution of eq.(13) by its boundary value. We can see that if there is a QFT with a two-point function $`𝒢`$ non-vanishing on the boundary then there is the unique choice of $`\varphi _0`$ solving eq.(13) such that $`\varphi =\varphi _D+\varphi _0`$ is a realization of a random field with the boundary value $`\mathrm{\Phi }`$. An average over $`\mathrm{\Phi }`$ leads to the Green function $`𝒢`$.
In the example of the hyperbolic space (14) (with the Poincare coordinates) the solution (20) of the Dirichlet boundary problem (13) can be expressed by its boundary value $`\varphi _B(y=0,x)=\mathrm{\Phi }(x)`$
$$\varphi _B(X)=y^{\frac{d}{2}}𝑑p\mathrm{exp}(ipx)|p|^{\frac{d}{2}}K_{\frac{d}{2}}(|p|y)\stackrel{~}{\mathrm{\Phi }}(p)$$
(24)
where $`\stackrel{~}{\mathrm{\Phi }}`$ denotes the Fourier transform of $`\mathrm{\Phi }`$ and $`K_\nu `$ is the modified Bessel function of order $`\nu `$ . Comparing with eq.(20) we obtain ($`X=(y,x)`$)
$$𝒟(X,x^{})=(2\pi )^dy^{\frac{3}{2}d+1}𝑑p\mathrm{exp}(ip(xx^{}))|p|^{\frac{d}{2}}K_{\frac{d}{2}}(|p|y)$$
The two-point function $`𝒢_E`$ resulting from the QFT on the hyperbolic space constructed in refs. is $`𝒢_E=\mathrm{ln}|xx^{}|`$ . In the Fourier transforms (up to an inessential normalization) we have ( see the discussion in )
$$\stackrel{~}{\mathrm{\Phi }}(p)\stackrel{~}{\mathrm{\Phi }}^{}(p^{})=\delta (pp^{})|p|^d$$
(25)
We may apply eq.(25) to calculate $`\varphi _B(X)\varphi _B(X^{})`$. As a solution of eq.(21) after an analytic continuation to the real time it must coincide with the Hadamard two-point function (vacuum expectation value of an anticommutator of quantum scalar fields) which is usually denoted by $`G^{(1)}(X,X^{})`$ ( the formula for $`𝒢_B`$ in the hyperbolic space can be found in and for $`𝒢_D`$ in )
$$𝒢_B(X,X^{})=\varphi _B(X)\varphi _B(X^{})=(yy^{})^{\frac{d}{2}}𝑑p\mathrm{exp}(ip(xx^{}))K_{\frac{d}{2}}(|p|y)K_{\frac{d}{2}}(|p|y^{})$$
(26)
## 4 Non-linear boundary value problem
We can modify the formulation (6)-(13) of QFT on manifolds with a boundary so that the interaction $`V(\varphi )`$ is taken into account already at the classical level. Then, instead of eq.(13) we consider the equation
$$𝒜\psi =V^{}(\psi )$$
(27)
or in the integral form
$$\psi (X)=\varphi _B(X)+𝑑X^{}\sqrt{G}𝒢_D(X,X^{})V^{}(\psi (X^{}))$$
(28)
where $`\varphi _B`$ is defined in eq.(20). In order to express the functional integral (6) in terms of $`\psi `$ let us introduce a differential operator
$$𝒜_\psi =𝒜+V^{\prime \prime }(\psi )$$
(29)
Define $`𝒢_D^\psi `$ as the Dirichlet Green function of $`𝒜_\psi `$. Let $`\mu _\psi `$ be the Gaussian measure with the mean zero and the covariance $`𝒢_D^\psi `$. Then, the formula (18) reads (under the assumption that the function on the r.h.s. of eq.(30) is integrable)
$$\begin{array}{c}𝑑\mu _0(\varphi )\mathrm{exp}(W_I(\varphi ))F(\varphi )=𝑑\mu _B(\varphi _B)\mathrm{exp}(W_0(\varphi _B)W(\psi ))det(𝒜_\psi )^{\frac{1}{2}}det(𝒜)^{\frac{1}{2}}\hfill \\ d\mu _\psi (\varphi _D)\mathrm{exp}(dX\sqrt{G}V(\varphi _D+\psi )+dX\sqrt{G}V(\psi )+dX\sqrt{G}V^{}(\psi )\varphi _D\hfill \\ +\frac{1}{2}dX\sqrt{G}\varphi _DV^{\prime \prime }(\psi )\varphi _D)F(\varphi _D+\psi )\hfill \end{array}$$
(30)
For the proof let us shift variables in eq.(18) and apply eqs.(9)-(12). Then,
$$\begin{array}{c}d\mu _D(\varphi _D+\chi )\mathrm{exp}(𝑑X\sqrt{G}V(\varphi _D+\varphi _B+\chi ))F(\varphi _D+\varphi _B+\chi )\hfill \\ =d\mu _D(\varphi _D))\mathrm{exp}(\frac{1}{2}dX\sqrt{G}\chi 𝒜\chi )F(\varphi _D+\psi )\hfill \\ \mathrm{exp}\left(𝑑X\sqrt{G}\chi 𝒜\varphi _D𝑑X\sqrt{G}V(\varphi _D+\psi )\right)\hfill \end{array}$$
(31)
where in the second step we inserted $`\chi =\psi \varphi _B`$ ($`\chi =0`$ on the boundary, hence the shift is admissible). Next, we make use of $`𝒜\varphi _B=0`$ (then $`𝒜\psi =𝒜\chi `$), subtract the two first terms of the Taylor expansion of $`V(\varphi _D+\psi )`$ in $`\psi `$ and apply the formula for a Gaussian integral of an exponential of a quadratic form
$$d\mu _D(\varphi _D))\mathrm{exp}(\frac{1}{2}dX\sqrt{G}\varphi _DV^{\prime \prime }(\psi )\varphi _D)=(det𝒜_\psi )^{\frac{1}{2}}d\mu _\psi (\varphi _D)$$
(32)
The final result is expressed in eq.(30). In this equation $`\mathrm{exp}(W(\psi ))`$ is the effective action in the tree approximation (discussed by ) and $`det𝒜_\psi ^{\frac{1}{2}}`$ gives the one-loop approximation to the effective action in QFT with the boundary value $`\mathrm{\Phi }`$ . The remaining $`d\mu _\psi (\varphi _D)`$ integral in eq.(30) can be calculated in perturbation expansion. It starts with higher powers $`n`$ ($`n3`$) of $`\varphi _D`$ leading to corrections in higher loops to the effective action.
## 5 Conclusions
In this section we derive some relations between correlation functions with respect to various measures discussed in earlier sections. Let us define
$$Z[\mathrm{\Phi }]=\mathrm{exp}(W_0(\varphi _B))𝑑\mu _D(\varphi _D)\mathrm{exp}(W_I(\varphi _D+\varphi _B))$$
(33)
where $`\varphi _B`$ is defined in eq.(20) with $`\mathrm{\Phi }`$ as a fixed boundary value. The definition (33) is introduced in such a way that it agrees with the large $`N`$ formula of and the semiclassical calculations of and the ones in eq.(30) (see also a discussion in ).
If $`Z[\mathrm{\Phi }]`$ is the generating functional then there exists a field $`𝒪(𝐱)`$ such that
$$Z[\mathrm{\Phi }]=\mathrm{exp}(_{}𝒪(x)\mathrm{\Phi }(x)\sqrt{g}𝑑x)$$
(34)
Treating $`Z[\mathrm{\Phi }]`$ as the generating functional we can calculate
$$\begin{array}{c}\frac{\delta }{\delta \mathrm{\Phi }(𝐱_1)}\mathrm{}.\frac{\delta }{\delta \mathrm{\Phi }(𝐱_n)}Z[\mathrm{\Phi }]_{|\mathrm{\Phi }=0}=(𝒟\frac{\delta }{\delta \varphi _B})(𝐱_1)\mathrm{}.(𝒟\frac{\delta }{\delta \varphi _B})(𝐱_n)\hfill \\ \mathrm{exp}(W_0(\varphi _B))d\mu _D(\varphi _D)\mathrm{exp}(W_I(\varphi _D+\varphi _B))_{|\varphi _B=0}\hfill \end{array}$$
(35)
where
$$(𝒟\frac{\delta }{\delta \varphi _B})(𝐱)𝑑X\sqrt{G}𝒟(X,𝐱)\frac{\delta }{\delta \varphi _B(X)}$$
and
$$W_0(\varphi _B)=\frac{1}{2}(\varphi _B,𝒜\varphi _B)$$
We wish to compare these correlation functions with the ones of the bulk field $`\varphi `$ defined by the generating functional
$$S[J]=𝑑\mu _0(\varphi )\mathrm{exp}(W_I(\varphi )+(J,\varphi ))$$
(36)
Then, the correlation functions can be calculated from the formula
$$\begin{array}{c}\frac{\delta }{\delta J(X_1)}\mathrm{}.\frac{\delta }{\delta J(X_n)}Z[J]_{|J=0}=(𝒢\frac{\delta }{\delta \varphi _c})(X_1)\mathrm{}.(𝒢\frac{\delta }{\delta \varphi _c})(X_n)\hfill \\ \mathrm{exp}(\frac{1}{2}(\varphi _c,𝒜\varphi _c))d\mu _0(\varphi )\mathrm{exp}(W_I(\varphi +\varphi _c))_{|\varphi _c=0}\hfill \end{array}$$
(37)
where on the r.h.s. we have absorbed the linear term of the exponential (36) into a shift of the measure according to eqs.(9)-(10) with $`\varphi _c=𝒢J`$ and $`J=𝒜\varphi _c`$ . It is clear from eqs.(35) and (37) that a perturbative calculation of $`n`$-point correlation functions of $`𝒪`$ and $`\varphi `$ involves the same graphs and only the propagators are different. The relation between (35) and (37) has been discovered earlier by Duetsch and Rehren ( see also ). A Hamil
tonian derivation of the relation between differentiation with respect to boundary values and sources $`J`$ can be found in .
The field theory in the bulk (6)-(7) is an integral over $`Z[\mathrm{\Phi }]`$
$$Z_0=𝑑\nu _B(\mathrm{\Phi })\mathrm{exp}(W_0(\varphi _B))Z[\mathrm{\Phi }]$$
(38)
We can obtain a connection between some other correlation functions. Generalizing eq.(38) let us define the generating functional $`S_D[\varphi _B;J]`$ in the $`\varphi _D`$ theory shifted by a background field $`\varphi _B`$
$$S_D[\varphi _B;J]=𝑑\mu _D(\varphi _D)\mathrm{exp}(W_I(\varphi _D+\varphi _B))\mathrm{exp}(𝑑X\sqrt{G}J\varphi _D)$$
(39)
Then, from eq.(18) the generating functional for correlation functions of the fields $`\varphi `$ in the model (6) is
$$S[J]=𝑑\mu _B(\varphi _B)\mathrm{exp}(𝑑X\sqrt{G}J\varphi _B)S_D[\varphi _B;J]$$
(40)
It can be seen that $`\varphi _D`$ and $`\varphi _B`$ enter symmetrically in eq.(18). Hence, we may also write
$$S[J]=𝑑\mu _D(\varphi _D)\mathrm{exp}(𝑑X\sqrt{G}J\varphi _D)S_B[\varphi _D;J]$$
(41)
Differentiating both sides of eq.(40) and (41) we obtain a relation between correlation functions of the fields $`\varphi `$,$`\varphi _D`$ and $`\varphi _B`$. The form of the correlation functions in the model (6) at the boundary points $`𝐱_j`$ is a simple consequence of eq.(41)
$$\varphi (𝐱_1)\mathrm{}\varphi (𝐱_n)=𝑑\mu _D(\varphi _D)𝑑\nu _B(\mathrm{\Phi })\mathrm{exp}(W_I(\varphi _D+\varphi _B))\mathrm{\Phi }(𝐱_1)\mathrm{}\mathrm{\Phi }(𝐱_n)$$
(42)
In particular, if the interaction is concentrated only on the boundary
$$W_I(\varphi )=_{}𝑑𝐱\sqrt{g}V(\varphi (0,𝐱))$$
then $`\varphi _D=0`$ in $`W_I`$ in eq.(42) and the functional integral (42) is the same as in QFT on $``$ defined by the free field measure $`d\nu _B`$ with the covariance $`𝒢_E(𝐱,𝐱^{})`$. In the case of the hyperbolic space this covariance is logarithmic. Hence, ultraviolet problem is the same as for quantum fields in two dimensions.
We think that the QFT theory on a boundary of a curved manifold is interesting for itself because of its remarkable regularity expressed (for the hyperbolic space) in the strong decay (25) in the momentum space. However, the main result of this paper is formulated in eqs.(40)-(42 ). The formulas connecting the correlation functions of fields in various field theoretic models can shed some light on relations of the AdS-CFT type.
Acknowledgements
The author is grateful to an anonymous referee for pointing out ref.. The formulas at the beginning of sec.5 resulted from my effort to derive the results of Duetsch and Rehren in the continuum as the addendum to an earlier version of my paper. |
warning/0506/nucl-th0506026.html | ar5iv | text | # Δ resonance contribution to two-photon exchange in electron-proton scattering
## Abstract
We calculate the effects on the elastic electron-proton scattering cross section of the two-photon exchange contribution with an intermediate $`\mathrm{\Delta }`$ resonance. The $`\mathrm{\Delta }`$ two-photon exchange contribution is found to be smaller in magnitude than the previously evaluated nucleon contribution, with an opposite sign at backward scattering angles. The sum of the nucleon and $`\mathrm{\Delta }`$ two-photon exchange corrections has an angular dependence compatible with both the polarisation transfer and the Rosenbluth methods of measuring the nucleon electromagnetic form factors.
The electromagnetic form factors reflect the essentially non-local nature of the nucleon in its interactions with photons. As the basic observables parametrising nucleon compositeness, the form factors have long been studied both experimentally and theoretically. This interest has been renewed recently due to the increased precision of electron-proton scattering experiments and the availability of two alternative methods of extracting the form factors from the data: the Rosenbluth method – also known as the longitudinal-transverse (LT) separation technique Wal94 ; Qat05 – and the polarisation-transfer (PT) technique Jon00 . If one uses the traditional one-photon exchange calculation to extract the form factors, the two methods lead to apparently incompatible results: while the PT method yields a ratio of the electric to magnetic form factors which falls off linearly with the square of the momentum transfer $`Q^2`$, the LT separation experiments give an approximately constant ratio Jon00 ; Bra02 ; Arr03 . Finding an explanation of this discrepancy is important for the use of electron-proton scattering as a precise and reliable tool in hadronic physics.
Several theoretical studies Blu03 ; Gui03 have suggested that the problem could be at least partially resolved by including higher-order two-photon exchange corrections in the analysis of electron-proton scattering data, in addition to the lowest order one-photon exchange (Born) approximation. The recent explicit calculation Blu03 has shown that with the two-photon exchange taken into account in the analysis of electron-proton scattering, the ratio of the form factors extracted from the LT separation measurements becomes more compatible with the ratio from the PT experiments. However, the two-photon exchange diagrams calculated in Ref. Blu03 contained only nucleons in the intermediate state; the contribution of other hadrons has not been included until now. In view of the prominent role of the $`\mathrm{\Delta }`$ resonance (unlike other excited states) in many hadronic reactions, it is essential to evaluate its contribution to the two-photon exchange in electron-proton scattering. Without an explicit calculation the results with only the nucleon intermediate state can only be viewed as suggestive in resolving the discrepancy. Some aspects of the $`\mathrm{\Delta }`$ contribution were addressed before Dre57 ; Cam69 , using various approximate approaches. These earlier studies demonstrated the importance of treating the $`\mathrm{\Delta }`$ on a par with the nucleon in considering higher-order corrections to electron-proton scattering.
This letter presents a quantum field theoretical calculation of the two-photon exchange “box” and “crossed-box” diagrams with a $`\mathrm{\Delta }`$ resonance in the intermediate state. We will show that the $`\mathrm{\Delta }`$ two-photon exchange correction is somewhat smaller in magnitude than that of the nucleon. At backward scattering angles the $`\mathrm{\Delta }`$ and nucleon contributions tend to partially cancel each other, their sum nevertheless yielding a predominantly negative two-photon exchange correction. We will show that the modified cross section has an angular dependence consistent with both the LT separation and PT measurements of the form factors.
We consider scattering of electrons (mass $`m_e0.511\times 10^3`$ GeV) off protons (mass $`M_N0.938`$ GeV) with the four-momenta assigned as $`e(p_1)+p(p_2)e(p_3)+p(p_4)`$. The differential cross section for this process is written in the form $`d\sigma =d\sigma _B(1+\delta _N+\delta _\mathrm{\Delta })`$ where $`d\sigma _B`$ is the lowest-order Born contribution (i. e. the cross section obtained from the one-photon exchange tree diagram) and $`\delta _N`$ ($`\delta _\mathrm{\Delta }`$) is the higher-order correction obtained from two-photon exchange diagrams containing nucleons ($`\mathrm{\Delta }`$’s) in the intermediate state. (Other higher-order effects which should be included in the formula for $`d\sigma `$ – such as the vacuum polarisation and the electron-photon vertex corrections – have been extensively studied in the past and are known Max00 to be irrelevant to the differences between the PT and LT analyses; we therefore focus here on the two-photon exchange effects only.) It is convenient to divide $`d\sigma `$ by the well-known factor describing the scattering from a structureless “proton” (see, e. g., Bjo64 ) and thus use the reduced cross section
$$d\sigma _R=\left[G_M^2(Q^2)+\frac{ϵ}{\tau }G_E^2(Q^2)\right](1+\delta _N+\delta _\mathrm{\Delta }).$$
(1)
Here the Born contribution is written in terms of the electric and magnetic form factors of the proton, $`G_E(Q^2)`$ and $`G_M(Q^2)`$, which are functions of the momentum transfer squared $`Q^2q^24\tau M_N^2=(p_1p_3)^2`$. The kinematic variable $`ϵ`$ is related to the scattering angle $`\theta `$ through $`ϵ=[1+2(1+\tau )\mathrm{tan}^2(\theta /2)]^1`$, which is equal to the photon polarisation in the Born approximation.
We denote the Born scattering amplitude as $`_B`$ and the two-photon exchange amplitudes with the nucleon and $`\mathrm{\Delta }`$ intermediate states as $`_N^{\gamma \gamma }`$ and $`_\mathrm{\Delta }^{\gamma \gamma }`$, respectively. From the equation $`d\sigma =d\sigma _B(1+\delta _N+\delta _\mathrm{\Delta })=\left|_B+_N^{\gamma \gamma }+_\mathrm{\Delta }^{\gamma \gamma }\right|^2`$, where $`d\sigma _B=\left|_B\right|^2`$, we derive to first order in the electromagnetic coupling $`e^2/(4\pi )1/137`$:
$$\delta _{N,\mathrm{\Delta }}=2\frac{\text{Re}\left(_B^{}_{N,\mathrm{\Delta }}^{\gamma \gamma }\right)}{\left|_B\right|^2}.$$
(2)
The nucleon part $`\delta _N`$ of the two-photon exchange was analysed in Ref. Blu03 . Below we will evaluate the $`\mathrm{\Delta }`$ two-photon exchange contribution $`\delta _\mathrm{\Delta }`$. The scattering amplitude $`_\mathrm{\Delta }^{\gamma \gamma }`$ is given by the sum of the box and crossed-box loop diagrams depicted in Fig. 1.
We use the $`\gamma N\mathrm{\Delta }`$ vertex of the following form Kon01 :
$`\mathrm{\Gamma }_{\gamma \mathrm{\Delta }N}^{\nu \alpha }(p,q)iV_{\mathrm{\Delta }in}^{\nu \alpha }(p,q)=i{\displaystyle \frac{eF_\mathrm{\Delta }(q^2)}{2M_\mathrm{\Delta }^2}}\{g_1[g^{\nu \alpha }p/q/p^\nu \gamma ^\alpha q/\gamma ^\nu \gamma ^\alpha pq+\gamma ^\nu p/q^\alpha ]`$
$`+g_2[p^\nu q^\alpha g^{\nu \alpha }pq]+(g_3/M_\mathrm{\Delta })[q^2(p^\nu \gamma ^\alpha g^{\nu \alpha }p/)+q^\nu (q^\alpha p/\gamma ^\alpha pq)]\}\gamma _5T_3,`$ (3)
where $`M_\mathrm{\Delta }1.232`$ GeV is the $`\mathrm{\Delta }`$ mass, $`p_\alpha `$ and $`q_\nu `$ are the four-momenta of the incoming $`\mathrm{\Delta }`$ and photon, respectively, and $`g_1`$, $`g_2`$ and $`g_3`$ are the coupling constants.<sup>1</sup><sup>1</sup>1We use the notation and conventions of Ref. Bjo64 throughout. An analysis of Eq. (3) in the $`\mathrm{\Delta }`$ rest frame suggests that $`g_1`$, $`g_2g_1`$ and $`g_3`$ may be interpreted as magnetic, electric and Coulomb components, respectively, of the $`\gamma N\mathrm{\Delta }`$ vertex. The form factor in Eq. (3) is necessary for ultraviolet regularisation of the loop integrals evaluated below; we use the simple dipole form
$$F_\mathrm{\Delta }(q^2)=\frac{\mathrm{\Lambda }_\mathrm{\Delta }^4}{\left(\mathrm{\Lambda }_\mathrm{\Delta }^2q^2\right)^2},$$
(4)
where $`\mathrm{\Lambda }_\mathrm{\Delta }`$ is the cutoff. The form factor entails some model-dependence of our results, which is unavoidable in any dynamical hadronic calculation. The isospin transition operator $`T_3`$ is defined by the relations $`_{\alpha =1}^3T_\alpha ^{}T_\alpha =1`$ and $`T_\alpha T_\beta ^{}=\delta _{\alpha \beta }\tau _\alpha \tau _\beta /3`$, where $`\tau _{1,2,3}`$ are the usual Pauli matrices. The vertex with an outgoing $`\mathrm{\Delta }`$ is given by the Dirac conjugate of Eq. (3), $`\mathrm{\Gamma }_{\gamma N\mathrm{\Delta }}^{\alpha \nu }(p,q)iV_{\mathrm{\Delta }out}^{\alpha \nu }(p,q)=\gamma _0\left[\mathrm{\Gamma }_{\gamma \mathrm{\Delta }N}^{\nu \alpha }(p,q)\right]^{}\gamma _0`$, with $`p_\alpha `$ and $`q_\nu `$ the four-momenta of the outgoing $`\mathrm{\Delta }`$ and incoming photon, respectively. The $`\gamma N\mathrm{\Delta }`$ vertex is orthogonal to the four-momenta of both the photon and the $`\mathrm{\Delta }`$:
$$q_\nu \mathrm{\Gamma }_{\gamma \mathrm{\Delta }N}^{\nu \alpha }(p,q)=0,p_\alpha \mathrm{\Gamma }_{\gamma \mathrm{\Delta }N}^{\nu \alpha }(p,q)=0.$$
(5)
The first of these equations ensures the usual electromagnetic gauge invariance of the calculation while the second allows us to use only the physical spin $`3/2`$ component,
$$S_{\alpha \beta }^\mathrm{\Delta }(p)=\frac{i}{p/M_\mathrm{\Delta }+i0}𝒫_{\alpha \beta }^{3/2}(p),𝒫_{\alpha \beta }^{3/2}(p)=g_{\alpha \beta }\frac{1}{3}\gamma _\alpha \gamma _\beta \frac{1}{3p^2}(p/\gamma _\alpha p_\beta +p_\alpha \gamma _\beta p/),$$
(6)
of the Rarita-Schwinger propagator Rar41 , the background spin $`1/2`$ component vanishing when contracted with the adjacent $`\gamma N\mathrm{\Delta }`$ vertices Pas99 . At present we do not include a width in the $`\mathrm{\Delta }`$ propagator as its influence on the unpolarised cross section should be small.
The loop integrals corresponding to the box and crossed-box diagrams in Fig. 1 can be written as
$$_\mathrm{\Delta }^{\gamma \gamma }=e^4\frac{d^4k}{(2\pi )^4}\frac{N_{box}^\mathrm{\Delta }(k)}{D_{box}^\mathrm{\Delta }(k)}e^4\frac{d^4k}{(2\pi )^4}\frac{N_{xbox}^\mathrm{\Delta }(k)}{D_{xbox}^\mathrm{\Delta }(k)},$$
(7)
with the numerators and denominators given by
$`N_{box}^\mathrm{\Delta }(k)`$ $`=`$ $`\overline{U}(p_4)V_{\mathrm{\Delta }in}^{\mu \alpha }(p_2+k,qk)[p/_2+k/+M_\mathrm{\Delta }]𝒫_{\alpha \beta }^{3/2}(p_2+k)V_{\mathrm{\Delta }out}^{\beta \nu }(p_2+k,k)U(p_2)`$ (8)
$`\times `$ $`\overline{u}(p_3)\gamma _\mu [p/_1k/+m_e]\gamma _\nu u(p_1),`$
$`N_{xbox}^\mathrm{\Delta }(k)`$ $`=`$ $`\overline{U}(p_4)V_{\mathrm{\Delta }in}^{\mu \alpha }(p_2+k,qk)[p/_2+k/+M_\mathrm{\Delta }]𝒫_{\alpha \beta }^{3/2}(p_2+k)V_{\mathrm{\Delta }out}^{\beta \nu }(p_2+k,k)U(p_2)`$ (9)
$`\times `$ $`\overline{u}(p_3)\gamma _\nu [p/_3+k/+m_e]\gamma _\mu u(p_1),`$
$`D_{box}^\mathrm{\Delta }(k)`$ $`=`$ $`\left[k^2+i0\right]\left[(kq)^2+i0\right]\left[(p_1k)^2m_e^2+i0\right]\left[(p_2+k)^2M_\mathrm{\Delta }^2+i0\right],`$ (10)
$`D_{xbox}^\mathrm{\Delta }(k)`$ $`=`$ $`D_{box}^\mathrm{\Delta }(k)|_{p_1kp_3+k},`$ (11)
where $`U`$ and $`u`$ denote the proton and electron four-spinor wave functions, respectively. Compared to the case of the nucleon Blu03 , the presence of a $`\mathrm{\Delta }`$ in the intermediate state entails a more complicated structure of the numerator. Also the loop integrals with a $`\mathrm{\Delta }`$ are not infrared divergent, in contrast with the nucleon contribution where the infrared part is very important Tsa61 ; Max00 . The evaluation of Eq. (7) involves preliminary algebraic manipulations to effect cancellations between terms in the numerators and denominators and subsequent integration of the thus simplified expressions. The result is obtained analytically in terms of the standard Passarino-Veltman dilogarithm functions tHo79 . In the calculation we used the computer package “FeynCalc” Mer91 .
The first and second loop integrals in Eq. (7) must be mutually related by crossing symmetry, which can be formulated in terms of the numerator of Eq. (2) using the Mandelstam variables $`s=(p_1+p_2)^2`$, $`t=(p_1p_3)^2`$ and $`u=(p_2p_3)^2=2M_N^2+2m_e^2ts`$. Denoting $`f^{\gamma \gamma }(s,t)_B^{}_\mathrm{\Delta }^{\gamma \gamma }`$ and writing it as the sum $`f^{\gamma \gamma }(s,t)=f_{box}^{\gamma \gamma }(s,t)+f_{xbox}^{\gamma \gamma }(s,t)`$, where the first (second) term is calculated using only the first (second) integral in Eq. (7), the crossing symmetry requires that
$$f_{xbox}^{\gamma \gamma }(s,t)=f_{box}^{\gamma \gamma }(u,t)|_{u=2M_N^2+2m_e^2ts}f^{\gamma \gamma }(s,t)=f^{\gamma \gamma }(2M_N^2+2m_e^2ts,t).$$
(12)
We calculated the integrals in Eq. (7) explicitly and checked that our results obey the crossing symmetry constraint Eq. (12).
The $`\mathrm{\Delta }`$ two-photon exchange correction to the differential cross section can be expressed as a quadratic form in the $`\gamma N\mathrm{\Delta }`$ coupling constants $`g_M=g_1`$, $`g_E=g_2g_1`$ and $`g_C=g_3`$:
$$\delta _\mathrm{\Delta }=C_Mg_M^2+C_{ME}g_Mg_E+C_Eg_E^2+C_Cg_C^2+C_{EC}g_Eg_C+C_{MC}g_Mg_C,$$
(13)
with the coefficients depending on the kinematical variables. The relative contributions of the coupling constants $`g_M`$, $`g_E`$ and $`g_C`$ to $`\delta _\mathrm{\Delta }`$ can be assessed from Table 1, where the $`C_M`$, $`C_{ME}`$, etc. are given as functions of $`ϵ`$ at two fixed $`Q^2`$ values.
In this calculation we used the dipole $`\gamma N\mathrm{\Delta }`$ form factor Eq. (4) with the cutoff $`\mathrm{\Lambda }_\mathrm{\Delta }=0.84`$ GeV, which describes a $`\mathrm{\Delta }`$ resonance whose mean-square radius is comparable to that of the nucleon. This choice is consistent with various parametrisations from pion electroproduction Cam69 ; Sat01 .
In the following we will discuss the results obtained with the fixed coupling constants $`g_M=7`$ and $`g_E=2`$. These couplings were used in the Dressed K-matrix Model Kon01 (adjusted for a different normalisation of the vertex used in the present calculation), yielding a good coupled-channel description of pion-nucleon scattering, pion photoproduction and Compton scattering at low and intermediate energies. For example, the $`E2/M1`$ ratio obtained in Ref. Kon01 from the pion photoproduction multipoles at the position of the $`\mathrm{\Delta }`$ resonance, is $`R_{EM}=\text{Im}E_{1+}^{3/2}/\text{Im}M_{1+}^{3/2}\times 100\%3\%`$, in agreement with the PDG Eid04 value: $`(2.5\pm 0.5)\%`$. Recent analyses Sat01 of pion electroproduction suggest that the Coulomb coupling constant $`g_C`$ is small and negative. In our calculation we will vary $`g_C`$ in the range $`[2,0]`$. With these values of $`g_M`$, $`g_E`$ and $`g_C`$ one can see from Eq. (13) and Table 1 that the magnetic coupling dominates the $`\mathrm{\Delta }`$ two-photon exchange correction whereas the electric coupling has a much smaller effect. Since the contribution of the Coulomb component is strongly suppressed (not exceeding $`0.2\%`$) we will omit it from further discussion, setting $`g_C=0`$ in the rest of the paper.
The $`ϵ`$ dependence of the sum of the $`\mathrm{\Delta }`$ and nucleon two-photon exchange corrections is shown in Fig. 2, for two fixed values of $`Q^2`$. The dependence on the $`\gamma N\mathrm{\Delta }`$ form factor can be seen by comparing the results obtained with the cutoffs $`\mathrm{\Lambda }_\mathrm{\Delta }=0.84`$ GeV and $`\mathrm{\Lambda }_\mathrm{\Delta }=0.68`$ GeV (the latter choice corresponds to a $`\mathrm{\Delta }`$ which is spatially “bigger” than the nucleon).
The purely nucleon contribution, shown for comparison, was calculated as in Ref. Blu03 using the $`\gamma NN`$ form factors extracted from the PT experiments Jon00 ; Bra02 . The $`\mathrm{\Delta }`$ correction is more prominent at higher momentum transfers. The $`\mathrm{\Delta }`$ tends to reduce the effect of the nucleon two-photon exchange, making the modulus of the negative nucleon correction somewhat smaller at backward angles (i. e. at low $`ϵ`$). The combined effect of the nucleon and $`\mathrm{\Delta }`$ two-photon exchanges produces a negative correction to the cross section at small $`ϵ`$, decreasing in magnitude as $`ϵ`$ increases.<sup>2</sup><sup>2</sup>2The diminishing of the two-photon exchange correction at forward angles is consistent with the analysis of electron-proton and positron-proton scattering data Arr04 . The main features of the $`\mathrm{\Delta }`$ contribution – its smallness and its tendency to attenuate the nucleon contribution at backward angles – are insensitive to the $`\gamma N\mathrm{\Delta }`$ form factor, being to that extent model-independent. The detailed interplay between the $`\mathrm{\Delta }`$ and the nucleon contributions can be more complicated, especially at forward angles, as can be seen from Fig. 2.
The calculated differential cross section is shown by the solid lines in Fig. 3, including the Born term and the sum of the two-photon exchange corrections $`\delta _N+\delta _\mathrm{\Delta }`$ with the nucleon and the $`\mathrm{\Delta }`$ intermediate states.
The reduced cross section Eq. (1), scaled for convenience by the square of the standard dipole form factor $`G_D(Q^2)=1/(1+Q^2/0.84^2)^2`$, is compared in Fig. 3 with the LT separation measurements from SLAC Wal94 (at $`Q^2=4`$ and $`6`$ GeV<sup>2</sup>) and JLab Qat05 (at $`Q^2=2.64`$ GeV<sup>2</sup>). The dotted lines show the Born contribution alone, using the nucleon form factors $`G_{E,M}(Q^2)`$ taken from the analysis of the JLab PT experiment Jon00 ; Bra02 . One can see that including only the Born term is inadequate in the analysis of the data. The addition of the two-photon exchange correction increases the slope of the cross section, also exhibiting some nonlinearity in $`ϵ`$. Thus the results of the PT and LT separation experiments become essentially compatible by including the nucleon and $`\mathrm{\Delta }`$ two-photon exchange corrections.
To summarise, we calculated the correction to the electron-proton scattering cross section due to the two-photon exchange with a $`\mathrm{\Delta }`$ intermediate state, treated on the same footing as the intermediate nucleon contribution. For realistic choices of the $`\gamma N\mathrm{\Delta }`$ vertex we found that the $`\mathrm{\Delta }`$ contribution alters the cross section by an amount from $`1\%`$ to $`+2\%`$, and is largest at backward scattering angles. For the cross section obtained using the LT separation technique, the $`\mathrm{\Delta }`$ two-photon exchange contribution slightly reduces the magnitude of the (negative) nucleon correction. Generally, the cross section including the nucleon and $`\mathrm{\Delta }`$ two-photon exchange corrections has the angular dependence which can accommodate the results of both the LT separation and PT methods of measuring the nucleon form factors. This calculation therefore provides explicit and compelling evidence that the two-photon exchange contribution (with the lowest mass, $`N`$ and $`\mathrm{\Delta }`$ intermediate states) can resolve the form factor discrepancy. To reconcile these two methods completely, theoretical analyses of the data might need additional ingredients. For example, one may take into account the dependence of the $`\gamma NN`$ and $`\gamma N\mathrm{\Delta }`$ vertices on the nucleon and $`\mathrm{\Delta }`$ off-shell momenta (as was suggested in Kon05 ). Heavier hadron resonances or quark degrees of freedom should also become important at higher momentum transfers (see e. g. Afa05 ).
###### Acknowledgements.
This work of S. K. and P. G. B. was supported in part by NSERC (Canada). The Southeastern Universities Research Association (SURA) operates the Thomas Jefferson National Accelerator Facility (Jefferson Lab) for the DOE under contract DE-AC05-84ER-40150. We thank John Arrington for useful comments. |
warning/0506/cond-mat0506392.html | ar5iv | text | # Microcanonical Ensemble Extensive Thermodynamics of Tsallis Statistics
## I Introduction
The equilibrium statistical mechanics and thermodynamics are well defined theories in modern physics Gibbs ; Balescu . Applications of these theories are restricted by investigation of the so-called thermodynamic or statistical systems which are constrained by several rigid requirements Kvasn . One of the first attempts to construct the generalized equilibrium statistical mechanics based on the mathematical redefinition of the Boltzmann-Gibbs statistical entropy and the principles of the information theory belongs to C. Tsallis Tsal88 . Until recently, there has been a great deal of interest in studying nonextensive thermodynamics due to its relevance in many fields of physics Tsal99 ; Gudima . However, many fundamental features regarding the violation of the zero law of thermodynamics and the principle of additivity remain unclear Abe0 ; Parv1 . Note that these difficulties have resulted in the occurrence of a large number of variants of the Tsallis generalized statistical mechanics Tsal98 .
The statistical mechanics investigates thermodynamic systems which are defined solely by the specification of macroscopic variables on the basis of the theory of probability and the microscopic laws of the classical and quantum mechanics. The evolution of the macroscopic system with a large number of degrees of freedom is impossible to describe by only dynamic methods. Therefore, the Gibbs idea of statistical ensembles is usually used Zub . All information about the macrostate of the system is contained in the phase distribution function, which evolves according to the Liouville equation or in the statistical operator whose evolution with time is described by the von Neumann equation. To derive the phase distribution function and the statistical operator is the primary goal of the nonequilibrium statistical mechanics.
In particular, the equilibrium statistical mechanics implies that one uses the Gibbs equilibrium statistical ensembles. In the state of thermodynamic equilibrium of the system, the phase distribution function and the statistical operator do not depend on time. Therefore, they are functions only of the first integrals of motion of the dynamic system. In this case, the mechanical laws and the Liouville and von Neumann equations do not allow one to determine unequivocally the equilibrium distribution function and the statistical operator Balescu ; Zub . Therefore, an obvious dependence of the equilibrium distributions on the macroscopic variables of the state of the system is defined by introducing additional postulates. The traditional way is based on the Gibbs postulate of the equiprobability of the dynamic states of the isolated system Gibbs . The alternative way rests on the Jaynes principle of a maximum of the information entropy Jaynes . The statistical mechanics constructed on the Gibbs equilibrium distributions, which corresponds to the Boltzmann-Gibbs statistical entropy, completely satisfies all postulates of the equilibrium thermodynamics.
Standard treatments of the Tsallis statistics point out that the entropic index $`q`$ is an additional intensive parameter, which has a fixed value for different thermodynamic systems Tsal98 . This concept leads to shortcomings of thermodynamics and needs to be reconsidered. As shown further, these problems can be resolved by the assumption that the parameter $`\xi =1/(q1)`$ is the extensive argument of the statistical entropy.
The paper is organized as follows. In the second section, the microscopic foundation of the Tsallis generalized statistical mechanics is given. The microcanonical equilibrium distribution function and statistical operator are deduced in the third section. In the fourth section, the performance of the thermodynamic principles in the microcanonical ensemble of the Tsallis statistics are proved. The developed formalism is exemplified in the fifth section by treating the classical microcanonical ideal gas.
## II Microscopic Foundation of Tsallis Statistical Mechanics
A macrostate of a system with a large number of degrees of freedom is imperfectly known at the microscopic level. The system can be found in any dynamic state compatible with the external macroscopic conditions. Therefore, for the macroscopic system it is possible to maintain only the probabilistic description of dynamic processes. For this reason, in the statistical mechanics the Gibbs idea of statistical ensembles is straightforward. The macroscopic state of the system is represented as a set of a large number of copies of the dynamic system under identical macroscopic conditions. Each system of the ensemble is represented by a point in phase space. Any physical observable $`A`$ of the macroscopic system is represented as the expectation value $`A^t`$ of the dynamic variable $`A(x,p,t)`$ with the phase distribution function $`\varrho (x,p,t)`$. The evolution with time of a phase distribution function is governed by the Liouville equation. According to the Liouville theorem, the volume of a region in phase space remains constant in the process of movement of phase points. The phase distribution function is constant along the phase trajectories, $`\varrho (x,p,t)=\varrho (x^{},p^{},t^{})`$. To describe the quantum many-particle systems the mixed states are considered. A macrostate thus appears as a set of possible microstates, which are set up by state vectors $`|\mathrm{\Psi }_r(t)`$, $`r=1,2,\mathrm{}`$, each with its own probability $`w_r`$ for its occurrence. The statistical operator $`\varrho (t)`$ allows to determine the expectation value of a dynamic variable $`A`$ regardless of the choice of the set of quantum states $`\{|\mathrm{\Psi }_r(t)\}`$. The evolution with time of a statistical operator is governed by the von Neumann equation. In the state of thermal equilibrium of the macroscopic system, the phase distribution function $`\varrho _{eq}(x,p)`$ and the statistical operator $`\varrho _{eq}`$ should not depend on the time $`t`$. Therefore, from the Liouville and von Neumann equations it follows that they are the first integrals of motion which must depend only on the first constants of motion of the system. Moreover, if these quantities are unequivocal and additive, then there exist only four such integrals of motion: energy $`H`$, the total momentum vector $`𝐏`$, the total angular momentum vector $`𝐌`$, and the number of particles $`N`$. In Balescu ; Zub , the basic definition of the microscopic foundation of the statistical mechanics is explained in detail. The present investigation rests on this assumption.
The equilibrium distribution function and the equilibrium statistical operator are not determined unequivocally by the mechanical laws. To express the equilibrium distributions from the macroscopic variables of state, the introduction of additional postulates is required. A traditional way to construct the equilibrium distributions is based on the Gibbs postulate of the equiprobability of all accessible dynamic states of the isolated system Gibbs . An alternative way for this is based on the statistical definition of the entropy and the use of the Jaynes principle explored in the information theory Jaynes . In the present study, we suggest a new method based on the laws of the equilibrium thermodynamics.
For this reason, we briefly recall the general laws of the macroscopic equilibrium thermodynamics. The thermodynamic systems are the object of the equilibrium thermodynamics and they must satisfy some obligatory conditions Kvasn . First, these are the systems of a large number of particles interacting with each other and with external fields. Second, for every thermodynamic system the zero law of thermodynamics is fulfilled, i.e., for such a system there exists a state of thermal equilibrium, which eventually is reached by the system at the fixed external conditions. This principle guaranties the existence of the special thermal measure, the temperature $`T`$, which is a general characteristic of all thermodynamic systems in equilibrium contact and which does not depend on the place and the method of measurements.
Third, for thermodynamic systems the principle of additivity is valid: all variables belong to two classes of additivity, according to the reaction of a given physical one to the division of the equilibrium system into the equilibrium macroscopic parts, for example, into two parts. The extensive variables can be split into two parts, and they should be proportional to the actual amount of matter present, $`_{1+2}=_1+_2`$. On the other hand, the intensive variables have to keep its values and cannot depend on the size of the system, $`\varphi _{1+2}=\varphi _1=\varphi _2`$. As an example, we may consider the thermodynamic systems which may be fixed in terms of the macroscopic variables of state $`T,V,N`$. In this case, the thermodynamic principle of additivity is implemented if intensive quantities are functions of intensive arguments, and extensive variables are proportional to the number of particles of the system multiplied by the intensive quantity. Such dependence of extensive and intensive variables is provided by the thermodynamic limit Kvasn . In this respect, all expressions have to be exposed to a formal limiting procedure $`N\mathrm{},V\mathrm{},v=V/N=\mathrm{const}`$, and only main asymptotics on $`N`$ should be kept. Then the extensive variables $``$ can be written $`(\alpha >0)`$
$$(T,V,N)|{}_{\genfrac{}{}{0pt}{}{N\mathrm{}}{v=\mathrm{const}}}{}^{}=N(f(T,v)+O(N^\alpha ))\stackrel{as}{=}Nf(T,v),$$
(1)
whereas the intensive variables $`\varphi `$ take the following form:
$$\varphi (T,V,N)|{}_{\genfrac{}{}{0pt}{}{N\mathrm{}}{v=\mathrm{const}}}{}^{}=\varphi (T,v)+O(N^\alpha )\stackrel{as}{=}\varphi (T,v),$$
(2)
where $`v=V/N`$ is the specific volume and $`f=/N`$ is the specific $``$. Note that the thermodynamic limit is a one-limiting procedure. The transitions not coordinated among themselves $`N\mathrm{}`$ and $`V\mathrm{}`$ have no physical sense, as in this case we would get results for either the superdense system or the empty one.
Fourth, in relation to the thermodynamic systems the first, the second and the third principles of thermodynamics are fulfilled, being the mathematical basis of the macroscopic theory. The first principle postulates the energy $`E`$ conservation law. The second principle of thermodynamics in the axiomatic formulation of R.J. Clausius postulates the existence of a function of state $`S_\mathrm{T}`$, called entropy. The absolute value of entropy is determined from the third law of thermodynamics or the Nernst theorem. The first and the second principles of thermodynamics for the quasistatic reversible processes are combined to give the fundamental equation of thermodynamics:
$$TdS_\mathrm{T}=dE+pdV+Xdz\mu dN,$$
(3)
where $`z=(z_1,\mathrm{},z_k)`$ and $`V`$ are the ”thermodynamic coordinates”; $`X=(X_1,\mathrm{},X_k)`$ and $`p`$ play the role of the associated ”forces”; $`\mu =\{\mu _i\}`$ are the chemical potentials and $`N=\{N_i\}`$ are the number of particles for each kind $`i`$, respectively. The second law of thermodynamics for nonequilibrium states, also formulated by R.J. Clausius, refers to the irreversible processes. This principle gives the direction of a real process allowing one to investigate the properties of equilibrium states as extreme ones. The most complete account of the equilibrium and non-equilibrium processes and the role of the characteristic times in the macroscopic thermodynamics is found in Kvasn ; Prigogine .
The expectation values of a dynamic variables with the equilibrium distribution function and the equilibrium statistical operator must satisfy all the postulates of the equilibrium thermodynamics. The connection between the distribution function and the macroscopic thermodynamical variables of the state is provided by the statistical entropy. Usually it is determined on the base of the information entropy. Let us define the Tsallis information entropy, which recently has received wide popularity due to the property of nonextensivity and which is used for construction of the so-called generalized statistical mechanics Tsal88 ; Tsal98 . The Tsallis information entropy for the discrete distribution of probabilities $`\{p_i\}`$ for $`W`$ independent elementary events is defined in the following manner Tsal88 :
$$S_{\mathrm{inf}}=k\underset{i=1}{\overset{W}{}}\frac{p_ip_i^q}{1q},\underset{i=1}{\overset{W}{}}p_i=1,$$
(4)
where $`k`$ is the Boltzmann constant and $`q𝐑`$ is the real parameter accepting values $`0<q<\mathrm{}`$. In the limit $`q1`$, we come to the well-known expression for the Boltzmann-Gibbs-Shannon entropy, $`S_{\mathrm{inf}}^{(BGS)}=k_{i=1}^Wp_i\mathrm{ln}p_i`$. The information entropy (4) is known as Havrda-Charvat-Daróczy-Tsallis entropy (see Rag99 ). However, in this paper, we shall use the short name for it. The information entropy is considered to be a measure of uncertainty of information concerning the statistical distribution $`\{p_i\}`$. The main its properties can be found in Tsal99 .
Let us introduce for further convenience a new representation for the Tsallis information entropy with a new parameter $`\xi `$
$$S_{\mathrm{inf}}=k\xi \underset{i=1}{\overset{W}{}}p_i(1p_i^{1/\xi }),\xi =\frac{1}{q1},$$
(5)
The parameter $`\xi `$ takes the values $`\mathrm{}\xi 1`$ for $`0<q1`$ and $`0<\xi \mathrm{}`$ for $`1q<\mathrm{}`$. In particular, in the limiting case for the value of the parameter $`q=1`$, we have $`\xi =\pm \mathrm{}`$.
Now, the Tsallis statistical entropy in the classical and quantum mechanics can be defined as follows
$$S(t)=k\xi \varrho (x,p,t)[1\varrho ^{1/\xi }(x,p,t)]𝑑\mathrm{\Gamma },S(t)=k\xi \mathrm{Tr}\{\varrho (t)[1\varrho ^{1/\xi }(t)]\},$$
(6)
where $`\varrho (x,p,t)`$ is the phase distribution function and $`\varrho (t)`$ is the statistical operator. It is easy to show that the Tsallis statistical entropy is not additive for the fixed value of $`\xi `$ and it is constant along the phase trajectories of the dynamic system. As the total time derivative from the phase distribution function is equal to zero, $`d\varrho /dt=0`$, valid from the Liouville equation and the Liouville theorem, the total time derivative from the classical entropy immediately yields the equality
$$\frac{dS(t)}{dt}=k\xi \frac{d\varrho (x,p,t)}{dt}[1(1+\frac{1}{\xi })\varrho ^{1/\xi }(x,p,t)]𝑑\mathrm{\Gamma }=0.$$
(7)
For the quantum ensembles, the Tsallis statistical entropy does not depend on time. Note that for the Gibbs statistical entropy this problem is inherent as well Zub .
## III Microcanonical Ensemble
In this section, the microcanonical distribution function and the statistical operator will be expressed through the variables of state of the isolated system $`(E,V,z,N)`$. Let us consider the equilibrium statistical ensemble of the closed energetically isolated systems of $`N`$ particles at the constant volume $`V`$ and the thermodynamic coordinate $`z`$. It is supposed that all systems have identical energy $`E`$ within $`\mathrm{\Delta }EE`$.
To begin with, we turn to instances of the classical case. The Tsallis equilibrium statistical entropy (6) represents a function of the parameter $`\xi `$ and a functional of the equilibrium phase distribution function $`\varrho _{eq}(x,p)`$:
$$S(\xi ,\{\varrho _{eq}\})=k\xi \underset{D}{}\varrho _{eq}(x,p)(1\varrho _{eq}^{1/\xi }(x,p))𝑑\mathrm{\Gamma }_N,$$
(8)
where $`d\mathrm{\Gamma }_N=dxdp`$ is an infinitesimal element of phase space. Let the phase distribution function $`\varrho _{eq}(x,p)`$ be distinct from zero only in the region of phase space $`D`$, which is defined by inequalities $`EH(x,p)E+\mathrm{\Delta }E`$ and be normalized to unity:
$$\underset{D}{}\varrho _{eq}(x,p)𝑑\mathrm{\Gamma }_N=1.$$
(9)
The phase distribution function depends on the first additive integrals of motion of the system. In particular, it is a function of the Hamiltonian, $`\varrho _{eq}(x,p)=\varrho _{eq}(H(x,p))`$. Moreover, the Hamilton function $`H(x,p)`$ has the parametrical dependence upon the number of particles $`N`$, volume $`V`$ of the system and $`z`$.
For an isolated system, in the state of thermal equilibrium the thermodynamic entropy $`S_\mathrm{T}(E,V,z,N)`$ has its maximal value. Hence, the fundamental equation of thermodynamics (3) for the quasiequilibrium processes is implemented. Changes of the variables of state at transition from one equilibrium state to another nearby state are equal to zero, $`dE=0,dV=0,dz=0`$ and $`dN=0`$. Therefore, from the basic equation of thermodynamics (3) it follows immediately that the thermodynamic entropy at the fixed values of $`E,V,z,N`$ is constant:
$$(dS_\mathrm{T})_{EVzN}=0.$$
(10)
To express the phase distribution function $`\varrho _{eq}(x,p)`$ through the variables of state $`(E,V,z,N)`$, let us replace the equilibrium thermodynamic entropy $`S_\mathrm{T}`$ of the macroscopic system with the Tsallis statistical one (8), $`S_\mathrm{T}(E,V,z,N)S(\xi ,\{\varrho _{eq}\})`$, and substitute it in Eq. (10). Taking into account Eqs. (9) and (10), one finds
$`dS={\displaystyle \frac{S}{\xi }}d\xi +{\displaystyle \underset{D}{}}{\displaystyle \frac{\delta S}{\delta \varrho _{eq}}}𝑑\varrho _{eq}𝑑\mathrm{\Gamma }_N`$ $`=`$ $`0,`$ (11)
$`{\displaystyle \underset{D}{}}𝑑\varrho _{eq}𝑑\mathrm{\Gamma }_N`$ $`=`$ $`0,`$ (12)
where the symbol $`d`$ before the functions $`S,\xi `$ and $`\varrho _{eq}`$ is the total differential in variables $`(E,V,z,N)`$. One should note that the unequivocal conformity between statistical and thermodynamic entropies is satisfied for the case where the parameters $`\xi `$ and $`\{\varrho _{eq}\}`$ are the functions of the variables of state $`(E,V,z,N)`$ of the isolated system. Let us put
$$\xi =\frac{1}{q1}=z.$$
(13)
Since $`d\xi =0`$ and $`d\varrho _{eq}=0`$, we obtain from Eqs. (11) and (12),
$$\frac{\delta S(z,\{\varrho _{eq}\})}{\delta \varrho _{eq}}=k\alpha ,$$
(14)
where $`\alpha `$ is a certain constant, and $`k`$ is the Boltzmann constant, which was introduced for convenience. Substituting Eq. (8) into (14), we obtain
$$\varrho _{eq}^{1/z}(x,p;E,V,z,N)=\frac{z\alpha }{z+1}.$$
(15)
The parameter $`\alpha `$ has been eliminated by using Eqs. (8) and (9):
$$\varrho _{eq}(x,p;E,V,z,N)=\left[1\frac{S}{kz}\right]^z.$$
(16)
Equations (16) and (9) together give
$$\left[1\frac{S}{kz}\right]^z=\underset{D}{}𝑑\mathrm{\Gamma }_N=\mathrm{\Delta }(H(x,p)E)𝑑\mathrm{\Gamma }_NW(E,V,N),$$
(17)
where $`\mathrm{\Delta }(\epsilon )`$ is the function distinct from zero only in the interval $`0\epsilon \mathrm{\Delta }E`$, where it is equal to unit. The statistical weight $`W(E,V,N)`$ is meant as a dimensionless phase volume, i.e., the number of dynamic states inside a layer $`\mathrm{\Delta }E`$. Based on this, we get the equipartition probability from Eq. (16) as a function of the thermodynamic ensemble variables, energy $`E`$, volume $`V`$, number of particle $`N`$, and parameter $`z`$ Zub :
$$\varrho _{eq}(x,p;E,V,z,N)=W^1(E,V,N)\mathrm{\Delta }(H(x,p)E).$$
(18)
Thus, using Eq. (17), we can write the entropy as (cf. Tsal88 ; Gross )
$$S(E,V,z,N)=kz[1W^{1/z}(E,V,N)]=kz[1e^{S_\mathrm{G}(E,V,N)/kz}],$$
(19)
where $`S_\mathrm{G}`$ is the Gibbs entropy Gibbs ; Zub for the microcanonical ensemble $`(E,V,N)`$:
$$S_\mathrm{G}(E,V,N)=k\mathrm{ln}W(E,V,N).$$
(20)
The quantum microcanonical ensemble and the corresponding equilibrium distribution function are in some respects analogous to the familiar classical ones. Let the probability distribution for quantum states of the system be different from zero only in the layer $`EE_iE+\mathrm{\Delta }E`$ and be normalized to unity:
$$\underset{i}{}w_i=1,EE_iE+\mathrm{\Delta }E.$$
(21)
The Tsallis equilibrium statistical entropy is a function of the parameter $`\xi `$ and probabilities $`\{w_i\}`$:
$$S(\xi ,\{w_i\})=k\xi \underset{i}{}w_i(1w_i^{1/\xi }).$$
(22)
The repeated use of the above procedure will lead us to the formula
$$\left[1\frac{S}{kz}\right]^z=\underset{i}{}\mathrm{\Delta }(E_iE)W(E,V,N).$$
(23)
The statistical weight $`W(E,V,N)`$ is equal to the number of quantum states in the layer $`\mathrm{\Delta }E`$. The quantum microcanonical distribution becomes
$$w_i(E,V,z,N)=W^1(E,V,N)\mathrm{\Delta }(E_iE).$$
(24)
The statistical operator corresponding to the microcanonical distribution of probabilities of quantum states (24) can be written as Zub
$$\varrho _{eq}(E,V,z,N)=W^1(E,V,N)\mathrm{\Delta }(HE),$$
(25)
where the operator function $`\mathrm{\Delta }(HE)`$ is determined in the diagonal representation by the matrix elements $`k|\mathrm{\Delta }(HE)|k^{}=\mathrm{\Delta }(E_kE)\delta _{kk^{}}`$. The quantum statistical entropy is calculated similarly to the classical one (19) with statistical weight (23). Note that the classical and quantum microcanonical distributions (18) and (24) are extreme equilibrium ones which correspond to a maximum of the Tsallis statistical entropy Tsal88 . The distribution functions (18) and (24) obtained by the thermodynamic method described here are identical with ones obtained by the Jaynes principle. The index $`q`$ for the Jaynes principle is a fixed parameter and does not depend on the variables of state of the system. In this case, the Tsallis statistics does not satisfy the zero law of thermodynamics Parv1 .
## IV Thermodynamics of microcanonical ensemble
It is well-known from the conventional statistical mechanics that in the thermal equilibrium the Gibbs entropy of the microcanonical ensemble is an extensive variable, and it has all peculiarities of the thermodynamic entropy in the thermodynamic limit Kvasn ; Zub . Mathematically, this implies that the Gibbs entropy $`S_\mathrm{G}`$ is a homogeneous function of variables $`E,V`$ and $`N`$ of the first order, i.e., one has the following property Prigogine :
$$S_\mathrm{G}(\lambda E,\lambda V,\lambda N)=\lambda S_\mathrm{G}(E,V,N),$$
(26)
where $`\lambda `$ is a certain constant. After substitution of Eq. (20) into (26), it is easy to check up that the statistical weight $`W`$ must satisfy the following requirement:
$$W(\lambda E,\lambda V,\lambda N)=W^\lambda (E,V,N).$$
(27)
Taking into account Eqs. (19) and (26), one finds the following peculiarity of the Tsallis entropy
$$S(\lambda E,\lambda V,\lambda z,\lambda N)=\lambda S(E,V,z,N),$$
(28)
which shows that the Tsallis entropy in the microcanonical ensemble is a homogeneous function of variables $`E,V,z,N`$ of the first order. In other words, it is extensive. It is essential to make clear that the homogeneity property of quantities (26)-(28) is realized only in the thermodynamic limit.
Differentiating Eq. (28) with respect to $`\lambda `$, and putting $`\lambda =1`$, we obtain the well-known Euler theorem for the homogeneous functions:
$$E\left(\frac{S}{E}\right)_{V,z,N}+V\left(\frac{S}{V}\right)_{E,z,N}+z\left(\frac{S}{z}\right)_{E,V,N}+N\left(\frac{S}{N}\right)_{E,V,z}=S.$$
(29)
Using the thermodynamic relations following from the fundamental equation of thermodynamics (3) in case of the isolated thermodynamic system $`(E,V,z,N)`$
$$\left(\frac{S}{E}\right)_{V,z,N}=\frac{1}{T},\left(\frac{S}{V}\right)_{E,z,N}=\frac{p}{T},\left(\frac{S}{z}\right)_{E,V,N}=\frac{X}{T},\left(\frac{S}{N}\right)_{E,V,z}=\frac{\mu }{T},$$
(30)
we get the Euler theorem Prigogine :
$$TS=E+pV+Xz\mu N.$$
(31)
Applying the differential operator with respect to the ensemble variables $`(E,V,z,N)`$ on Eq. (31), we obtain the fundamental equation of thermodynamics
$$TdS=dE+pdV+Xdz\mu dN$$
(32)
and the Gibbs-Duhem relation Prigogine
$$SdT=Vdp+zdXNd\mu .$$
(33)
Equation (33) means that the variables $`T`$, $`\mu `$, $`X`$ and $`p`$ are not independent. The fundamental equation of thermodynamics (32) provides the first principle
$$\delta Q=dE+pdV+Xdz\mu dN$$
(34)
and the second law of thermodynamics
$$dS=\frac{\delta Q}{T}.$$
(35)
Here $`\delta Q`$ is a heat transfer by the system to the environment for quasistatic transition of the system from one equilibrium state to another nearby state.
Let us investigate the homogeneity properties of the variables $`T,p,\mu `$ and $`X`$. Substituting Eq. (19) into (30), we obtain the following expressions for the temperature $`T`$ Parv1 :
$$T(E,V,z,N)=T_\mathrm{G}(E,V,N)W^{1/z}(E,V,N)=T_\mathrm{G}(E,V,N)e^{S_\mathrm{G}(E,V,N)/kz}$$
(36)
and the variable $`X`$
$$X(E,V,z,N)=kT_\mathrm{G}(E,V,N)[e^{S_\mathrm{G}(E,V,N)/kz}1S_\mathrm{G}(E,V,N)/kz].$$
(37)
The pressure and the chemical potential of the system are equivalent with the pressure $`p_\mathrm{G}`$ and the chemical potential $`\mu _\mathrm{G}`$ of the Gibbs statistics, respectively, $`p(E,V,z,N)=p_\mathrm{G}(E,V,N)`$ and $`\mu (E,V,z,N)=\mu _\mathrm{G}(E,V,N)`$. These equations were derived by using the thermodynamic relations for the temperature $`T_\mathrm{G}`$, the pressure $`p_\mathrm{G}`$ and the chemical potential $`\mu _\mathrm{G}`$ of the Gibbs statistics, and taking into account Eq. (20):
$$\frac{1}{T_\mathrm{G}}=\left(\frac{S_\mathrm{G}}{E}\right)_{V,N},\frac{p_\mathrm{G}}{T_\mathrm{G}}=\left(\frac{S_\mathrm{G}}{V}\right)_{E,N},\frac{\mu _\mathrm{G}}{T_\mathrm{G}}=\left(\frac{S_\mathrm{G}}{N}\right)_{E,V}.$$
(38)
The Gibbs quantities $`T_\mathrm{G},p_\mathrm{G},\mu _\mathrm{G}`$ are the homogeneous functions of the variables of state $`(E,V,N)`$ of the zero order. This can be proved by using Eqs. (38) and (26). Then, a combination of Eqs. (36) and (26) allows us to write the relation for the temperature $`T`$:
$$T(\lambda E,\lambda V,\lambda z,\lambda N)=T(E,V,z,N).$$
(39)
Similarly to Eq. (39), the relations for the pressure $`p(\lambda E,\lambda V,\lambda z,\lambda N)`$, the chemical potential $`\mu (\lambda E,\lambda V,\lambda z,\lambda N)`$, and the variable $`X(\lambda E,\lambda V,\lambda z,\lambda N)`$ are fulfilled. Thus, the temperature $`T`$, the pressure $`p`$, the chemical potential $`\mu `$, and quantity $`X`$ are the homogeneous functions of the variables $`E,V,z,N`$ of the zero order. So they are intensive variables Prigogine .
Let us prove in more detail the thermodynamic principle of additivity Kvasn . For instance, we assume that $`\lambda =1/N`$ and introduce the following specific variables:
$$\epsilon =\frac{E}{N},v=\frac{V}{N},\stackrel{~}{z}=\frac{z}{N}=\frac{1}{(q1)N}.$$
(40)
Thus, Eqs. (28) and (39) for the entropy and the temperature of the system, by using (40) with respect to $`\lambda =1/N`$, can be rewritten as
$$s(\epsilon ,v,\stackrel{~}{z})=\frac{1}{N}S(E,V,z,N)$$
(41)
and
$$T(\epsilon ,v,\stackrel{~}{z})=T(E,V,z,N),$$
(42)
where $`s(\epsilon ,v,\stackrel{~}{z})`$ is the specific entropy, $`s=S/N`$, which depends only on the intensive variables $`\epsilon ,v`$ and $`\stackrel{~}{z}`$. For the pressure $`p`$, the chemical potential $`\mu `$, and $`X`$, we have equations similar to that for the temperature (42). So, comparing Eqs. (41) and (42) with the thermodynamic equations (1) and (2), we conclude that the entropy $`S`$ is an extensive variable, as it is proportional to the number of particles $`N`$ multiplied by an intensive variable $`s`$, but the temperature $`T`$, the pressure $`p`$, the chemical potential $`\mu `$, and $`X`$ are intensive variables.
Let us divide the system into two parts ($`1`$ and $`2`$) and require that the total number of particles of the system should be equal to the sum of the number of particles of each subsystem separately and the specific quantities (40) should be equal among themselves
$$N_{1+2}=N_1+N_2,\epsilon _{1+2}=\epsilon _1=\epsilon _2,v_{1+2}=v_1=v_2,\stackrel{~}{z}_{1+2}=\stackrel{~}{z}_1=\stackrel{~}{z}_2.$$
(43)
Then, the variables $`E,V`$ and $`z`$ are extensive. Taking into account Eq. (43), one finds
$$s_{1+2}(\epsilon _{1+2},v_{1+2},\stackrel{~}{z}_{1+2})=s_1(\epsilon _1,v_1,\stackrel{~}{z}_1)=s_2(\epsilon _2,v_2,\stackrel{~}{z}_2).$$
(44)
Multiplying it by the first equation from (43) and using (41), we get
$$S_{1+2}(E_{1+2},V_{1+2},z_{1+2},N_{1+2})=S_1(E_1,V_1,z_1,N_1)+S_2(E_2,V_2,z_2,N_2).$$
(45)
Thus, in the microcanonical ensemble the Tsallis entropy is an extensive variable. Furthermore, Eqs. (42) and (43) allow us to write
$$T_{1+2}(E_{1+2},V_{1+2},z_{1+2},N_{1+2})=T_1(E_1,V_1,z_1,N_1)=T_2(E_2,V_2,z_2,N_2).$$
(46)
So, in the thermodynamic limit, the zero law of thermodynamics and the thermodynamic principle of additivity (see. (1) and (2)) for the Tsallis statistics in the microcanonical ensemble are valid. Here, the thermodynamic limit denotes the limiting statistical procedure $`N\mathrm{}`$ at $`\epsilon =\mathrm{const},v=\mathrm{const}`$, and $`\stackrel{~}{z}=\mathrm{const}`$ with keeping the main asymptotics on $`N`$. This may be explicitly seen by making an expansion of the functions of the state in powers of the small parameter $`1/N`$ ($`N1`$$`,|z|1`$, i.e. $`q1`$) with the large finite values of the variables $`E,V,z`$. For the extensive functions the term proportional to $`N`$ is held and for the intensive variables only the term proportional to $`N^0`$ is kept (cf. Eqs. (1) and (2)). Note that the correct thermodynamic limit, $`(q1)N=\mathrm{const}`$, for the Tsallis statistics has already been discussed in Botet et al. Botet1 ; Botet . In Abe Abe1 , the thermodynamic limit for the Tsallis statistics is wrong because the limits $`N\mathrm{}`$ and $`|z|\mathrm{}`$ are not coordinated among themselves. This procedure destroys the connection between the variables $`N`$ and $`z`$ in the functions of the state of the system (see Section II). In the case of the Boltzmann-Gibbs limit we make an expansion of the functions of the state in powers of the small parameter $`1/\stackrel{~}{z}`$ ($`|\stackrel{~}{z}|1`$, $`N1`$, i.e. $`q1`$) and hold only the zero term of the power expansion.
It is important to note that some authors Abe0 ; Vives to make the connection with the equilibrium thermodynamics interpret the equations similar to our Eqs. (19) and (36) as the extensive representation of the Tsallis statistics in the terms of the physical temperature. But, as was shown in Parv1 , these equations are only the transformation formulas from the Tsallis statistics to the new extensive one.
## V The perfect gas
The thermodynamic principle of additivity can thoroughly be investigated in the framework of a classical nonrelativistic ideal gas. In the microcanonical ensemble $`(E,V,z,N)`$, the statistical weight (17) of the perfect gas of $`N`$ identical nucleons is given by Das
$$W(E,V,N)=\frac{V^N}{N!}\frac{d^3p_1\mathrm{}d^3p_N}{(2\pi \mathrm{})^{3N}}\delta \left(\underset{i=1}{\overset{N}{}}\frac{\stackrel{}{p}_i^2}{2m}E\right)=\frac{V^N}{N!}\left(\frac{m}{2\pi \mathrm{}^2}\right)^{\frac{3}{2}N}\frac{E^{\frac{3}{2}N1}}{\mathrm{\Gamma }(\frac{3}{2}N)},$$
(47)
where $`m`$ is the nucleon mass. In the thermodynamic limit ($`N1`$, $`\epsilon =E/N=\mathrm{const}`$, $`v=V/N=\mathrm{const}`$) from Eq. (47), it follows immediately that Kvasn
$$W^{1/N}(E,V,N)=v\left(\frac{m\epsilon e^{5/3}}{3\pi \mathrm{}^2}\right)^{3/2}w(\epsilon ,v).$$
(48)
So Eq. (48) proves relation (27) for the statistical weight $`W`$ with $`\lambda =1/N`$. Then, the Tsallis entropy (19) is reduced to
$$S(E,V,z,N)=Ns(\epsilon ,v,\stackrel{~}{z}),s(\epsilon ,v,\stackrel{~}{z})=k\stackrel{~}{z}\left[1w^{1/\stackrel{~}{z}}\right].$$
(49)
Comparing Eq. (49) with (1), we conclude that the Tsallis entropy is an extensive variable. Note that in the limit $`\stackrel{~}{z}\pm \mathrm{}`$, we obtain the formula for the Gibbs specific entropy:
$$s(\epsilon ,v,\stackrel{~}{z})|_{\stackrel{~}{z}\pm \mathrm{}}=k\mathrm{ln}ws_\mathrm{G}(\epsilon ,v).$$
(50)
Substituting (49) into (30), we get
$$T(E,V,z,N)=\frac{2}{3}\frac{\epsilon }{k}w^{1/\stackrel{~}{z}}T(\epsilon ,v,\stackrel{~}{z}).$$
(51)
The temperature (51) is a function of the specific variables $`\epsilon ,v`$ and $`\stackrel{~}{z}`$. Therefore, it is an intensive variable by virtue of Eq. (2). In the limit $`\stackrel{~}{z}\pm \mathrm{}`$, we obtain the well-known formula for the Gibbs statistics
$$T(\epsilon ,v,\stackrel{~}{z})|_{\stackrel{~}{z}\pm \mathrm{}}=\frac{2}{3}\frac{\epsilon }{k}=T_\mathrm{G}(\epsilon ,v).$$
(52)
In a similar way, the pressure $`p`$, the chemical potential $`\mu `$, and $`X`$ become
$`p(E,V,z,N)`$ $`=`$ $`{\displaystyle \frac{2}{3}}{\displaystyle \frac{\epsilon }{v}}p(\epsilon ,v,\stackrel{~}{z}),`$ (53)
$`\mu (E,V,z,N)`$ $`=`$ $`{\displaystyle \frac{2}{3}}\epsilon \left[{\displaystyle \frac{5}{2}}\mathrm{ln}w\right]\mu (\epsilon ,v,\stackrel{~}{z}),`$ (54)
$`X(E,V,z,N)`$ $`=`$ $`{\displaystyle \frac{2}{3}}\epsilon \left[1+{\displaystyle \frac{1}{\stackrel{~}{z}}}\mathrm{ln}ww^{1/\stackrel{~}{z}}\right]X(\epsilon ,v,\stackrel{~}{z}).`$ (55)
Note that the pressure $`p`$ and the chemical potential $`\mu `$ for the classical ideal gas in the microcanonical ensemble do not depend on the parameter $`\stackrel{~}{z}`$, and they are equal to respective quantities of the Gibbs statistics. Then, Eqs. (49), (51), and (53)-(55) yield the Euler theorem (31) in terms of the specific variables
$$Ts=\epsilon +pv\mu +X\stackrel{~}{z}.$$
(56)
In the limit $`\stackrel{~}{z}\pm \mathrm{}`$, the pressure $`p`$ and the chemical potential $`\mu `$ remain unchanged but the variable $`X=0`$. So, by the example of the classical ideal gas the principle of additivity for the Tsallis statistics is proved. The Euler theorem (31) or (56) shows that in the Tsallis statistics the quantities $`z=1/(q1)`$ and $`X`$ should be the variable of the state and the associated ”force”, respectively.
At this point one important quantity must be noted, the heat capacity $`C_V=1/(T/E)_{V,z,N}`$. In the framework of the ideal gas of $`N`$ nucleons it can be written as
$$C_V=\frac{3}{2}kNw^{1/\stackrel{~}{z}}\left(1+\frac{3}{2}\frac{1}{\stackrel{~}{z}}\right)^1.$$
(57)
Figure 1 shows the specific heat (left) and the temperature (center) vs the parameter $`\stackrel{~}{z}`$. The calculations are done for the system of nucleons at the specific energy $`\epsilon =50`$ MeV and the specific volume $`v=3/\rho _0`$, where $`\rho _0=0.168`$ fm<sup>-3</sup>. It is of great interest that both the heat capacity and the temperature sharply change their shape in the region of small values of $`\stackrel{~}{z}`$ and considerably defer from their Gibbs limit, which in the figure is indicated by arrows. In the region of $`3/2<\stackrel{~}{z}<0`$, the heat capacity is negative. It is remarkable that such a behaviour has really been caused by the decrease of the temperature with $`\epsilon `$. This dependence can be seen even better in right panel of Fig. 1 which shows the temperature vs the specific energy of the system for different values of the parameter $`\stackrel{~}{z}`$. The negative heat capacity in the microcanonical ensemble is closely related to the phase transition of the first order where the entropy is a convex function Chomaz ; Gross1 . The Fig. 1 clearly shows that the variable $`\stackrel{~}{z}`$ is the order parameter and the system is physically unstable in the region $`3/2<\stackrel{~}{z}<0`$ at the fixed values of the variables $`(v,\stackrel{~}{z})`$, where the entropy is a convex function of $`\epsilon `$. For example, entropy $`s\epsilon ^{3/2}`$ at $`\stackrel{~}{z}=1`$. Note that this critical feature of the system is not related with negative values of the parameter $`q`$ Tsal98 as we take into account the condition $`1+1/N\stackrel{~}{z}>0`$ ($`N1`$). The crossing point of all curves $`(\epsilon _0,v_0)`$ in right panel of Fig. 1 is the point where the Tsallis and Gibbs entropies vanish, $`S=0`$ and $`S_G=0`$ or $`w(\epsilon _0,v_0)=1`$. Thus, the following values of the energy $`\epsilon >\epsilon _0`$ and the volume $`v>v_0`$ of the system have the physical sense. Note that the temperature $`T`$ do not depend on energy $`\epsilon `$ at $`\stackrel{~}{z}=3/2`$ and it is equal with temperature $`T_0=T(\epsilon _0,v_0)`$.
It should be remarked here that the microcanonical ensemble of the Tsallis statistics in the thermodynamic limit is equivalent with the canonical one. If we introduce the variable $`z`$ in the formulas for the perfect gas in the canonical ensemble of Abe et al. Abe0 ; Abe2 , then in the thermodynamic limit we recover the above functions of the state of the microcanonical ensemble. For instance, in the thermodynamic limit, Eq. (51) can be obtained from the energy $`U_q(T,V,N)`$ given in Abe0 ; Abe2 .
## VI Conclusions
In this paper, we have explored the microscopic foundation of the generalized equilibrium statistical mechanics based on the Tsallis statistical entropy. The viewpoint utilized here considers that the microcanonical ensemble is most convenient to analyze the fundamental questions of the statistical mechanics. We summarize our main principles.
Here, the Gibbs idea of the statistical ensembles defined within the framework of the quantum and classical mechanics was used. In this approach, the equilibrium phase distribution function and the statistical operator do not depend on time, and they are functions of the additive first integrals of motion of the system by virtue of performance of Liouville and von Neumann equations. Additionally, these main quantities are functions of the macroscopic variables of state of the system. To derive the distribution functions, in contrast with the Jaynes principle, the new thermodynamic method based on the fundamental equation of thermodynamics and statistical definition of the functions of the state of the system was given.
In this paper, we have made the following claim. The index $`\xi `$ of the Tsallis entropy should be an extensive variable of the state of the system. As a result of this assumption, we obtain that in the microcanonical ensemble the Tsallis entropy represents the homogeneous function of the variables $`E,V,z,N`$ of the first order. The temperature of the system is an intensive variable, and, consequently, the zero law of thermodynamics is satisfied. Other functions of state of the system are either extensive or intensive. Thus, in the thermodynamic limit, $`\stackrel{~}{z}=1/(q1)N=\mathrm{const}`$, in the Tsallis statistics the thermodynamic principle of additivity is carried out. Note that the Tsallis information entropy is nonextensive because the parameter $`\xi `$ is a certain intensive constant. Also it is necessary to note that the Tsallis statistical entropy as well as the Gibbs one has an essential lack. Both the entropies do not depend on time while the thermodynamic entropy grows up to achieve its maximal value in the state of thermal equilibrium. The extensive property of the Tsallis entropy in the microcanonical ensemble yields the Euler theorem which permits one to find the fundamental equation of thermodynamics and the Gibbs-Duhem relation. Thus, the first and the second principles of thermodynamics are fulfilled. Note that in the limit, $`\stackrel{~}{z}\pm \mathrm{}`$, all expressions of the Tsallis statistics take the form of the conventional Gibbs statistical mechanics. So the Tsallis statistical mechanics in the microcanonical ensemble satisfies all postulates of the equilibrium thermodynamics.
Finally, the classical nonrelativistic ideal gas of $`N`$ identical nucleons in the microcanonical ensemble was considered to illustrate the principles which were elucidated in the general theory. It has been shown that in the thermodynamic limit the statistical weight, the entropy, the temperature, and other quantities are the homogeneous functions of the first and zero order of the variables of state, respectively. Note that for ideal gas the Euler theorem was accomplished and in the limit, $`\stackrel{~}{z}\pm \mathrm{}`$, all expressions resembled ones of the Gibbs statistics.
Acknowledgments: This work has been supported by the Moldavian-US Bilateral Grants Program (CRDF project MP2-3045). We acknowledge valuable remarks and fruitful discussions with T.S. Biró, R. Botet, K.K. Gudima, M. Płoszajczak, and V.D. Toneev. |
warning/0506/astro-ph0506048.html | ar5iv | text | # VLA Imaging of the Intriguing HI Cloud HIJASS J1021+6842 in the M 81 Group
## 1 Introduction
The nearby M 81 group provides an ideal testbed for studies of dwarf galaxies and the effects of interactions between galaxies. Thanks to its proximity, it is one of the best studied groups, both regarding the interaction between its three luminous galaxies (M 81, M 82 and NGC 3077, e.g., Yun, Ho, & Lo, 1994), and the studies of individual group members, notably several of the dwarf galaxies (e.g., Puche et al., 1992; Walter & Brinks, 1999; Walter et al., 2002; Ott et al., 2001).
Boyce et al. (2001) recently conducted a blind HI survey of the M 81 group using a multibeam receiver on the 76–m Lovell telescope with a resultant beamsize of $``$ 12′. Interestingly, only one new object was discovered in their observations, HIJASS J1021+6842, which is the subject of this letter. This implies that the census of galaxies with HI gas in the M 81 group must be nearly complete down to HI masses of order 10<sup>7</sup> M.
HIJASS J1021+6842 lies at an angular distance of $`105^{}`$ (or a minimum separation of 110 kpc at the distance of M 81) from IC 2574. The proximity in sky position and radial velocity between IC 2574 and HIJASS J1021+6842 implies that HIJASS J1021+6842 is a probable member of the M 81 group at a distance of $`4`$ Mpc and a possible companion of IC 2574 at a distance of $`4`$ Mpc (Boyce et al. (2001); Karachentsev et al. (2002)). Boyce et al. (2001) noted the lack of an optical counterpart for HIJASS J1021+6842 in the second–generation red Digital Sky Survey, and speculated that it may be either a very low surface brightness companion of IC 2574, or the debris of a tidal encounter between IC 2574 and one of the galaxies around M 81.
Thus, HIJASS J1021+6842 is an intrinsically interesting object, and potentially a key object for understanding the process of galaxy formation from tidal debris. The present higher resolution HI observations provide a requisite first step to understanding the true nature of this object.
## 2 Observations and Analysis
We used the NRAO<sup>1</sup><sup>1</sup>1The National Radio Astronomy Observatory is a facility of the National Science Foundation operated under cooperative agreement by Associated Universities, Inc. Very Large Array (VLA) to obtain HI spectral line observations of HIJASS J1021+6842 in the D configuration (integration time: 106 minutes) and in the C configuration (148 minutes) in 2004. We observed 0841+708 as the complex gain calibrator and 0542+498 (3C147) and 1331+305 (3C286) as flux and bandpass calibrators. The spectra consisted of 128 channels with a width of 5.2 km s<sup>-1</sup> (after on–line Hanning smoothing). The data were analyzed using standard calibration/mapping tasks in the AIPS software package. Four channels near velocities of 0 km s<sup>-1</sup> were contaminated by Galactic foreground emission (see Fig. 1). To calculate the integrated HI map (moment 0) we carefully inspected and blanked consecutive channel maps to remove the extended Galactic emission from the data cube. In doing so we are helped by the fact that smoothly distributed emission, such as that due to foreground Galactic HI, is filtered out in interferometric observations.
Given the intrinsically faint HI emission of HIJASS J1021+6842, images were made with a uvtaper of 5.3 k$`\lambda `$ to emphasize weak, extended structures. This resulted in a beamsize of $`60^{\prime \prime }\times 52^{\prime \prime }`$ (60<sup>′′</sup>=1.17 kpc) and an rms of 0.8 mJy beam<sup>-1</sup> per channel. Given the considerable extent of the HI emission, it was important to correct for primary beam attenuation in all subsequent, quantitative analysis.
## 3 HI Distribution and Kinematics
### 3.1 HI Distribution
Figure 1 shows a mosaic of 12 channel maps (before primary beam correction) stepped at every other channel. This mosaic shows that HI is detected from $``$ $``$50 to $``$$`+`$70 km s<sup>-1</sup> in radial velocity, which is a larger range in velocity than the FWHM of $``$ 50 km s<sup>-1</sup> reported by Boyce et al. (2001). It is likely that the synthesis observations have allowed us to overcome confusion with Galactic HI emission, enabling us to detect HI at velocities near 0 km s<sup>-1</sup> and lower velocities. Additionally, the higher resolution imaging shows that the single source discovered by Boyce et al. (2001) breaks up into several regions of emission. There is a radial velocity gradient with positive velocities in the eastern components and negative velocities in the western components. What is particularly striking about these images is the fact that the HI is extended over more than 30 kpc which makes HIJASS J1021+6842 bigger than even the most massive dwarf galaxies in the M 81 group (e.g., IC 2574: Walter & Brinks (1999) Ho II: Puche et al. (1992)).
Figure 2 shows a map of the total HI column density of HIJASS J1021+6842. We identify seven regions of emission (labeled I–VII in Fig. 2). The maximum HI column density is found in region VI and corresponds to $`1.8\times 10^{20}`$ atoms cm<sup>-2</sup>. Region III is a marginal detection and needs follow up observations.
The total HI emission shown in Fig. 2 corresponds to a total HI mass of 1.5 $`\times `$ 10<sup>8</sup> M. This is considerably larger than the value reported by Boyce et al. (2001). We attribute the difference to the fact that most of the emission is situated in small clumps (compared to the Lovell beam, FWHM: $``$12), that the area over which the clumps are spread out is larger than the size of the Lovell beam, and that HIJASS J1021+6842 extends across velocities containing Galactic HI emission. The HI masses of the individual regions are listed in the caption of Fig. 2.
### 3.2 HI Kinematics and Dynamical Mass Estimate
The channel maps (Fig. 1) show that there is a velocity gradient across the HI extent of HIJASS J1021+6842 extending from region I (near 70 km s<sup>-1</sup>) to region VII ($``$50 km s<sup>-1</sup>). The velocities at which the individual regions appear brightest in the channel maps are given in Fig. 2 (right).
One interpretation of the velocity gradient is that the emission is gravitationally bound and that the velocity gradient can be interpreted as being due to a (broken) annulus in rotation. Under this assumption we find for a radius of 15 kpc and a V<sub>max</sub> of 40 km s<sup>-1</sup> a minimum dynamical mass of $`5.5\times 10^9`$ M. The resulting ratio of dynamical over detected mass (or M<sub>dyn</sub>/M<sub>HI</sub>) is $`>`$ 10 for the ensemble. This would be a large ratio which is higher than is typically observed in lower luminosity dwarf irregular galaxies (Skillman, 1996). We stress however, that a tidal origin may also result in the observed velocity gradient which would imply a significantly lower dynamical mass.
## 4 Constraints on the True Nature of HIJASS J1021+6842
### 4.1 Lack of an Optical Counterpart
As noted by Boyce et al. (2001), there is no obvious optical counterpart to HIJASS J1021+6842 which is detected on images of the DSS. In addition to this, the M 81 group has been the subject of many very thorough searches for small and low surface brightness members (e.g., van den Bergh, 1959; Karachentseva, 1968; Boerngen & Karachentseva, 1982; Caldwell et al., 1998; Karachentseva & Karachentsev, 1998). These searches reached faint surface brightness levels and lead to the discovery of dwarfs as faint as m$`{}_{\mathrm{V}}{}^{}18`$, corresponding to M$`{}_{\mathrm{V}}{}^{}10`$ at the distance of the M81 group (thus reaching down to the extreme low–mass end of the luminosity function). As yet, none of these searches have come up with candidate counterparts for HIJASS J1021+6842 (e.g., Caldwell et al., 1998). One possible concern is the presence of a bright star in the field (HD 89343, RA(2000.0) = 10 21 03.3, DEC(2000.0) = $`+`$68 44 52, V = 6 mag.) which may limit the local sensitivity to low surface brightness optical emission.
It should be noted, though, that the HI distribution is so different from the other M 81 group dwarfs that perhaps one would not necessarily expect stars to be associated with this system. Normally, star formation is associated with HI column densities above a certain threshold (e.g., Kennicutt, 1989). Empirically this threshold hovers around $`10^{21}`$ atoms cm<sup>-2</sup> for dwarf irregular galaxies. The fact that no optical counterpart has been found appears to be consistent with this threshold. Clearly, deep optical follow–up observations are needed to determine at higher confidence whether or not stars are associated with HIJASS J1021+6842.
Another possible explanation for the lack of stars in HIJASS J1021+6842 is the possibility that we are seeing the early stages of formation of a new galaxy. This could either be formation from a primodial cloud or formation from tidal debris as discussed, e.g., by Makarova et al. (2002) or Bournaud et al. (2004) (and references therein). Given the relatively dense environment of the central M 81 triplet, the position of HIJASS J1021+6842 seems to be an unlikely location for the contemporary precipitation of a primordial galaxy. If HIJASS J1021+6842 is tidal debris (perhaps from the outer parts of one of the larger galaxies in the M 81 group), then the small velocity differences between its components allow for the possibility that the entire system is gravitationally bound, which could potentially lead to future concentration and the formation of stars.
### 4.2 Group Membership
Since HIJASS J1021+6842 was discovered in a survey of the M 81 group, it is natural to assume that it is a member of that group. Nonetheless, HIJASS J1021+6842 is an unusual system and, before assuming group membership, it is important to consider evidence in favor or against this assumption. An overview of the relative locations of the M 81 group members in the plane of the sky is presented in (Karachentsev et al., 2002), their Figure 1. The systemic velocity of $``$ 30 km s<sup>-1</sup> is certainly consistent with membership in the M 81 group (Karachentsev et al., 2002). In fact, it is only offset from M 81 by 2.3° in angle (corresponding to a minimum distance of $``$ 140 kpc) and 40 km s<sup>-1</sup> in radial velocity. In addition, HIJASS J1021+6842 is situated very close to the M 81 group member IC 2574 both in position (projected distance: 110 kpc) and velocity ($`\mathrm{\Delta }`$v=30 km s<sup>-1</sup>). This may suggest that IC 2574 and HIJASS J1021+6842 are companions, a scenario which is tentatively supported by the HIJASS observations which indicate a low column density HI bridge connecting them.
What is the probability that this object is part of the Local Group rather than the M 81 group? In that case HIJASS J1021+6842 would likely be classified as a high velocity cloud (HVC). There are several problems with such a hypothesis, however. A strong argument against it is its fairly high column density, of order $`10^{20}`$ atoms cm<sup>-2</sup> which is at least an order of magnitude higher than what is found even in the Magellanic Stream (Putman et al., 2003). Given its angular size, HIJASS J1021+6842 would be classified as a Compact HVC. Again, typical column densities for CHVCs are at least an order of magnitude lower (de Heij, Braun, & Burton, 2002). An additional strong argument arguing against HIJASS J1021+6842 being an HVC is that typical HVCs in the direction of M 81 have velocities in the range $``$150 to $``$200 km s<sup>-1</sup> (Wakker & van Woerden, 1997), i.e. at much lower velocities. We therefore conclude that HIJASS J1021+6842 is indeed associated with the M 81 group.
## 5 Conclusions
We present VLA HI 21cm observations of HIJASS J1021+6842 discovered by Boyce et al. (2001) in the direction of the M 81 group. Given its location, column densities and systemic velocity, HIJASS J1021+6842 is very likely a member of the M 81 group. Perhaps the most striking result is that the HI emission is distributed over an area of 30 kpc in diameter which is much larger than the extent of the largest dwarf members of the M 81 group. This, and the lack of detectable stars makes the HIJASS J1021+6842 system absolutely unique amongst the M 81 group members. Other HI clouds known to date which have no detected optical counterparts are the SW clump of HI 1225+01 (Giovanelli & Haynes, 1989; Chengalur et al., 1995; Turner & MacFadyen, 1997), HIPASS J0731-69 (Ryder et al., 2001), and HI clouds in Virgo (Davies et al., 2004; Minchin et al., 2005). We detect 1.6 $`\times `$ 10<sup>8</sup> M of HI distributed over roughly 30 kpc in HIJASS J1021+6842. Peak HI column densities are of order $`1.8\times 10^{20}`$ atoms cm<sup>-2</sup>, which is well below the empirical threshold for star formation activity to commence. This may explain why to date no optical counterpart has been identified.
The individual clouds which make up the system are either self gravitating entities orbiting within their mutual gravitational potential, or the densest concentrations of a huge, 30 kpc diameter, gently rotating, very low surface density cloud. Assuming that the entire complex is gravitationally bound, we derive a minimum dynamical mass of $`5.5\times 10^9`$ M, which would be more than an order of magnitude more massive than the luminous (HI) mass. It should be stressed, though, that other scenarios, such as a past stripping event leading to the formation of tidal debris, may also result in the observed velocity structure (thus leading to a significantly lower dynamical mass of the system).
The present observations clearly demonstrate the uniqueness of HIJASS J1021+6842. Future high sensitivity HI synthesis observations of HIJASS J1021+6842 and its surroundings and deep optical and UV imaging will be necessary to elucidate the true nature of this enigmatic object.
EDS is grateful for partial support from a NASA LTSARP grant No. NAG5-9221 and the University of Minnesota. We thank our referee, Dr. Mike Disney, for useful comments which helped to improve the presentation of this paper. |
warning/0506/nucl-th0506085.html | ar5iv | text | # Bosonization of the Pairing Hamiltonian
## 1 Introduction
The problem of the bosonization of finite and infinite fermionic systems has been extensively explored in the past. All these investigations have been prompted by the recognition that in nature this phenomenon is widely occurring. Best examples of it are offered by the superconductivity in metals with the related BCS theory and by the superfluidity of certain atomic nuclei where it is signalled by the existence of a gap (about 1 MeV) in the energy spectrum, by the reduction of about a factor of two in the moment of inertia with respect to the shell model prediction and by the staggering behavior of the separation energies (odd-even effect in the mass number A). More generally in the nuclear case the bosonization phenomenon is of course epitomized by the outstandingly successful Interacting Boson Model (IBM) by Arima and Iachello . In this paper we first address the problem of the bosonization in the path integral formalism. Although this approach is general and hence keeps its validity in bosonizing any fermionic action with a quartic interaction, for sake of illustration we shall consider the well-known pairing hamiltonian which reads
$$\widehat{H}\widehat{H}_0+\widehat{H}_P=\underset{\nu =1}{\overset{L}{}}e_\nu \underset{m_\nu =j_\nu }{\overset{j_\nu }{}}\widehat{a}_{j_\nu m_\nu }^{}\widehat{a}_{j_\nu m_\nu }G\underset{\mu ,\nu =1}{\overset{L}{}}\widehat{A}_\mu ^{}\widehat{A}_\nu .$$
(1)
In (1) $`\widehat{a}_{j_\nu m_\nu }`$ is the destruction operator of a fermion in a single particle level (whose number is $`L`$) characterized by angular momentum $`j_\nu `$, third component $`m_\nu `$ and single particle energy $`e_\nu `$, whereas the operator
$$\widehat{A}_\nu =\underset{m_\nu =1/2}{\overset{j_\nu }{}}(1)^{j_\nu m_\nu }\widehat{a}_{j_\nu ,m_\nu }\widehat{a}_{j_\nu m_\nu }$$
(2)
destroys a pair of fermions with total angular momentum $`J=0`$ in the level $`j_\nu `$.
The Hamiltonian (1) is indeed:
1. relevant for both nuclear and condensed matter physics. In nuclear physics it is known to represent, together with the quadrupole force, a dominant part of the residual nucleon-nucleon interaction and a crucial microscopic ingredient of the IBM . If the interacting fermions are electrons, $`\widehat{H}`$ (in particular $`\widehat{H}_P`$) is the Hamiltonian underlying the BCS theory of superconductivity . Moreover, the same Hamiltonian governs collective phenomena in the physics of liquids and metal clusters;
2. a simple example to illustrate the bosonization of fermionic systems. The problem of relating the fermionic Hamiltonian (1) to a bosonic one, like the IBM, with the same spectrum, has been the object of several investigations and still represents a challenging question.
Furthermore:
3. for infinite systems, the BCS theory, of which the hamiltonian (1) offers a simple realization, is an example of spontaneous symmetry breaking of the gauge group $`U(1)`$, with the associated appearance of a Goldstone boson , set up with excitations consisting of the addition and removal of pairs of particles to and from the system. The question then arises whether this Goldstone mode survives in finite systems and, if so, how the Goldstone and Higgs excitations should be identified.
4. The Hamiltonian (1) has been shown to be exactly integrable , which means that the number of constants of motion equals the number of degrees of freedom. This does not imply, however, that explicit expressions for its eigenvalues and eigenstates are easily obtained, since the constants of motion are usually complicated operators which, in general, can only be diagonalized numerically. For this reason it is useful to provide approximate analytic expressions for the eigenvalues and eigenvectors of $`\widehat{H}_P`$.
This report is organized as follows. In Section 2 we illustrate how the path integrals scheme works and apply it to (1) in the degenerate case, where only one single particle level is considered, deducing the well-known zero seniority spectrum in the path integral formalism, discussing the appearance of a Goldstone boson and studying the seniority excitations within a matrix approach in Sec. 2.5. In Section 3 we study, by directly solving the Richardson equations, the problem of one and two pairs living in many non-degenerate single particle levels, deriving analytic expressions for the energies of both the collective and the trapped states. Finally, in Sec. 4 we briefly discuss the relation between the exact and BCS solution to the pairing problem in terms of Bogolioubov quasi-particles.
## 2 The degenerate case
The pairing Hamiltonian has been first diagonalized by Kerman, Lawson and MacFarlane in the simple case of one degenerate single particle level by using group theoretical methods. In fact $`\widehat{H}`$ can be cast in the following form
$$\widehat{H}=2\underset{\nu =1}{\overset{L}{}}e_\nu \widehat{S}_{z\nu }+\underset{\nu =1}{\overset{L}{}}e_\nu \mathrm{\Omega }_\nu G\widehat{S}_+\widehat{S}_{},$$
(3)
$`\mathrm{\Omega }_\nu =j_\nu +1/2`$ being the pair degeneracy of the level $`j_\nu `$, in terms of the quasi-spin operators
$`\widehat{S}_+={\displaystyle \underset{\nu =1}{\overset{L}{}}}\widehat{A}_\nu ^{},\widehat{S}_{}={\displaystyle \underset{\nu =1}{\overset{L}{}}}\widehat{A}_\nu ,`$ (4)
$`\widehat{S}_z={\displaystyle \frac{1}{2}}{\displaystyle \underset{\nu =1}{\overset{L}{}}}\widehat{S}_{z\nu }={\displaystyle \underset{\nu =1}{\overset{L}{}}}{\displaystyle \underset{m_\nu =1/2}{\overset{j_\nu }{}}}\left(\widehat{a}_{j_\nu m_\nu }^{}\widehat{a}_{j_\nu m_\nu }\widehat{a}_{j_\nu m_\nu }\widehat{a}_{j_\nu m_\nu }^{}\right)`$ (5)
which span a SU(2) algebra. Hence the pure pairing Hamiltonian $`\widehat{H}_P`$ can be immediately diagonalized, yielding
$$G<\widehat{S}_+\widehat{S}_{}>=G\left[N/2\left(\mathrm{\Omega }N/2+1\right)v/2\left(\mathrm{\Omega }v/2+1\right)\right],$$
(6)
where $`N`$ is the eigenvalue of the particle number operator (the number of fermions is assumed here to be an even integer)
$$\widehat{N}=\underset{\nu =1}{\overset{L}{}}\underset{m_\nu =j_\nu }{\overset{j_\nu }{}}\widehat{a}_{j_\nu m_\nu }^{}\widehat{a}_{j_\nu m_\nu }$$
(7)
and $`\mathrm{\Omega }=_\nu \mathrm{\Omega }_\nu `$ the total degeneracy. The seniority quantum number
$$v=\mathrm{\Omega }2S,$$
(8)
$`S(S+1)`$ being the eigenvalue of $`\widehat{S}^2`$, represents the number of unpaired particles.
Now, in the degenerate case $`(L=1)`$ $`e_\nu =e`$ and the full Hamiltonian becomes
$$\widehat{H}=2e(\widehat{S}_z+\mathrm{\Omega })G\widehat{S}_+\widehat{S}_{}.$$
(9)
Hence the spectrum generating algebra is still SU(2) and the energy turns out to read
$$E(n,s)=2enG\left[n\left(\mathrm{\Omega }n+1\right)s\left(\mathrm{\Omega }s+1\right)\right].$$
(10)
In (10), for future convenience, the pair number $`n=N/2\mathrm{\Omega }`$ and the quantum number
$$s=v/2\mathrm{min}\{n,\mathrm{\Omega }n\},$$
(11)
which counts the number of broken pairs, have been introduced.
The spectrum (10) is associated with two independent types of excitations: one is related to the addition or removal of one pair of fermions and the other to the breaking of a pair. The former, described by the quantum number $`n`$, is a Goldstone boson associated with the spontaneous breaking of the global gauge invariance reflecting the particle number conservation; the latter, described by the quantum number $`s`$, can be viewed as corresponding to the Higgs excitations .
In an infinite system the energy of the Goldstone bosons vanishes with the associated quantum number. This does not occur in a finite system, but, to the extent that the energy spectrum of the latter displays a pattern similar to that of an infinite system, it should exhibit two quite different energy scales. Actually the excitation energies associated with both the quantum numbers $`n`$ and $`s`$ appear to be of the same order, namely $`g\mathrm{\Omega }`$. However the Goldstone nature of the energy spectrum associated with $`n`$ is clearly apparent when one considers the excitations with respect to the minimum of (10). This occurs for
$$n=n_0=[\nu _0],$$
(12)
$`[\mathrm{}]`$ meaning integer part, with
$$\nu _0\frac{1}{2}(\mathrm{\Omega }+1)\frac{e}{G}.$$
(13)
Introducing then the shifted quantum number
$$\nu =nn_0,$$
(14)
(10) becomes
$$E(n_0+\nu ,s)=G\nu ^2+2G\nu (n_0\nu _0)Gn_0(2\nu _0n_0)+Gs(\mathrm{\Omega }s+1).$$
(15)
Therefore the addition (or removal) of one pair of nucleons with respect to the ground state requires an energy of order $`G`$: this is the energy of the Goldstone boson. The energy required to break a pair, namely the seniority energy, is instead of order $`G\mathrm{\Omega }`$: this is the energy of a Higgs boson.
### 2.1 The path integral approach
In this Section we show that the pairing spectrum (10) for $`s=0`$ can be obtained in the framework of the Feynman path integral and illustrate how an effective bosonic action can be set up in this case.
To do this it is necessary to introduce the odd Grassmann variables $`\lambda _m`$, $`\lambda _m^{}`$ associated to the fermionic destruction and creation operator $`\widehat{a}_m`$, $`\widehat{a}_m^{}`$ (we omit the index $`j`$ since we are considering only one single particle level) and the even elements
$$\phi _m(t)=(1)^{jm}\lambda _m(t)\lambda _m(t),\phi _m^{}(t)=(1)^{jm}\lambda _m^{}(t)\lambda _m^{}(t),$$
(16)
describing pairs of fermions with $`J_z=0`$.
Since we restrict ourselves to the zero seniority part of the spectrum (10), we define
$$\mathrm{\Phi }(t)=\underset{m>0}{}(1)^{jm}\lambda _m(t)\lambda _m(t),$$
(17)
namely the even elements of the Grassmann algebra associated to the $`J=0`$ composites (2), also referred to as “hard bosons”. The analysis of seniority excitations, which will be studied in 2.5 in a different framework, would clearly require the introduction of $`J0`$ composites.
Note that the index of nilpotency of the composites (17) is $`\mathrm{\Omega }`$ (i.e., $`\mathrm{\Phi }^n=0`$, $`n>\mathrm{\Omega }`$), reflecting the Pauli principle.
In terms of the above variables the euclidean action associated to the Hamiltonian (1) reads
$$S=\tau \underset{t=N_0/2}{\overset{N_0/21}{}}\left\{\underset{m}{}\lambda _m^{}(t)(_t^++e)\lambda _m(t1)G\mathrm{\Phi }^{}(t)\mathrm{\Phi }(t1)\right\},$$
(18)
where the time has been discretized ($`\tau `$ is the time spacing), $`N_0`$ is the number of sites of the lattice, $`_t^+`$ the discrete time derivative
$$_t^\pm f(t)=\pm \frac{1}{\tau }[f(t\pm 1)f(t)]$$
(19)
and $`e`$ the single particle fermion energy. The variables $`\lambda _m`$ at the time $`t=N_0/2`$ are related to the ones at $`t=N_0/21`$ by antiperiodic boundary conditions
$$\lambda _m(N_0/2)=\lambda _m(N_0/21).$$
(20)
The action (18) can be written in terms of the even variables as
$$S=\underset{t=N_0/2}{\overset{N_0/21}{}}\left\{\underset{m>0}{}\left[\phi _m^{}(t)\phi _m(t)x^2\phi _m^{}(t)\phi _m(t1)\right]G\tau \mathrm{\Phi }^{}(t)\mathrm{\Phi }(t1)\right\},$$
(21)
where
$$x=1\tau e.$$
(22)
Following the approach of Ref. for performing Berezin integrals over composite variables, the generating functional and the correlation functions of $`n`$ spin zero pairs of nucleons can then be expressed as integrals over the $`\phi `$ variables of the composites, according to
$$Z=[d\phi ^{}d\phi ]e^S$$
(23)
and
$$<\mathrm{\Phi }^n(t_2)\left[\mathrm{\Phi }^{}(t_1)\right]^n>=\frac{1}{Z}[d\phi ^{}d\phi ]\mathrm{\Phi }^n(t_2)\left[\mathrm{\Phi }^{}(t_1)\right]^ne^S,$$
(24)
respectively. In the continuum limit the following result is then obtained
$$<\mathrm{\Phi }^n(t_2)\left[\mathrm{\Phi }^{}(t_1)\right]^n>\underset{^{\tau 0}}{}e^{(\beta _2\beta _1)E_n}$$
(25)
for the correlation function, with $`\beta _{1,2}=\tau t_{1,2}`$ and
$$E_n=n[2eG(\mathrm{\Omega }n+1)]$$
(26)
for the energy spectrum. The latter exactly coincides with (10) when $`s=0`$: hence the well-known zero-seniority formula for the spectrum of the system is recovered.
Concerning the bosonization of the fermionic system, although in the pairing Hamiltonian the fermionic and bosonic degrees of freedom get mixed, a purely bosonic effective action for the field $`\mathrm{\Phi }`$ can be deduced using the variables above introduced. It reads
$$S_{eff}^{(\mathrm{\Phi })}=\underset{t=N_0/2}{\overset{N_0/21}{}}\left\{\underset{k=1}{\overset{\mathrm{\Omega }}{}}\alpha _k\left[\mathrm{\Phi }^{}(t)\right]^k\left[\mathrm{\Phi }^k(t)+x^{2k}\mathrm{\Phi }^k(t1)\right]G\tau \mathrm{\Phi }^{}(t)\mathrm{\Phi }(t1)\right\},$$
(27)
$`\alpha _k`$ being numerical coefficients whose asymptotic behavior is found to be
$$\alpha _k\underset{^\mathrm{\Omega }\mathrm{}}{}\mathrm{\Omega }^{(2k1)}.$$
(28)
The action $`S_{eff}^{(\mathrm{\Phi })}`$ is equivalent to the fermionic action (18) in the sense that it leads to the same spectrum (10). Worth noticing is that its structure is characterized by the absence of terms of the type $`[\mathrm{\Phi }^{}(t)]^k[\mathrm{\Phi }(t)]^{kp}[\mathrm{\Phi }(t1)]^p`$ ($`p0`$), due to the occurrence of nontrivial cancellations.
### 2.2 Goldstone bosons in the Hubbard-Stratonovitch approach
In what follows we again consider the part of the spectrum (10) associated with the addition and removal of pairs, showing its connection with a Goldstone boson.
It is well-known that the spontaneous breaking of the local electromagnetic gauge invariance plays a crucial role in the theory of superconductivity . Here the symmetry group U(1) of the global gauge transformations $`\psi e^{i\alpha }\psi `$ is broken by the non-vanishing expectation value of pair operators (the Cooper pairs), which carry charge $`2`$: the symmetry group is thus reduced to the unbroken subgroup $`Z_2`$, consisting of the two gauge transformations with $`\alpha =0`$ and $`\pi `$. The Goldstone field associated to this symmetry breaking lives in the coset space U(1)/$`Z_2`$ and displays only derivative interactions.
In a finite system, like the one considered here, a spontaneous symmetry breaking cannot rigorously occur, since tunnelling takes place between the various possibly degenerate states and the true ground state turns out to be a unique linear superposition of the degenerate states . However, if the two above mentioned features of the Goldstone fields in infinite systems, namely that they parameterize the coset space U(1)/$`Z_2`$ and have only derivative interactions, survive in finite systems, then the identification of a Goldstone boson becomes of relevance not only for a deeper understanding of the bosonization mechanism, but also for a convenient choice of the variables.
To explore this occurrence we perform a Hubbard-Stratonovitch linearization of the action (18), introducing an integration over auxiliary fields in the generating functional
$$Z=[d\lambda d\lambda ^{}d\eta d\eta ^{}]e^S,$$
(29)
where the new action is
$`S`$ $`=`$ $`\tau {\displaystyle \underset{t=N_0/2}{\overset{N_0/21}{}}}\{G[\eta ^{}(t)\eta (t)+\eta ^{}(t)\mathrm{\Phi }(t)+\eta (t)\mathrm{\Phi }^{}(t)]`$ (30)
$`+`$ $`{\displaystyle \underset{m=j}{\overset{j}{}}}\left[\lambda _m^{}(t)(_t^++e)\lambda _m(t)\right]\}`$
and the auxiliary bosonic fields $`\eta ^{}`$ and $`\eta `$ satisfy periodic boundary conditions.
Next we introduce the following polar representation for the auxiliary fields
$`\eta `$ $`=`$ $`\sqrt{\rho }e^{2i\theta },\eta ^{}=\sqrt{\rho }e^{2i\theta }`$ (31)
and explore whether the field $`\theta `$ can be identified with the Goldstone boson.
First, for the change of variable (31) to be one to one (with the only exception of the point $`\rho =0`$), $`\theta `$ must vary in the range $`0\theta <\pi `$: hence the field $`\theta `$ lives in the coset space of the broken symmetry group $`U(1)`$ of particle conservation with respect to the unbroken subgroup $`Z_2`$, as appropriate to a Goldstone field .
Second, the Goldstone field should display only derivative couplings in the action, which is not the case for the field $`\theta `$ after the transformation (31). However the non-derivative coupling can be eliminated introducing the following transformation on the nucleon fields:
$$\lambda _m=e^{i\theta }\psi _m,\lambda _m^{}=e^{i\theta }\psi _m^{}.$$
(32)
In fact, as a consequence of the above, the following operators
$`q^\pm `$ $`=`$ $`\mathrm{exp}(i\theta )_t^\pm \mathrm{exp}(\pm i\theta )\pm e,`$ (33)
whose time matrix elements are
$`\left(q^+\right)_{t_1t_2}`$ $`=`$ $`{\displaystyle \frac{1}{\tau }}\left[\mathrm{exp}\left\{i\tau \left(_t^+\theta \right)_{t_1}\right\}\delta _{t_2,t_1+1}\delta _{t_1t_2}\right]+e\delta _{t_1t_2}`$ (34)
$`\left(q^{}\right)_{t_1t_2}`$ $`=`$ $`{\displaystyle \frac{1}{\tau }}\left[\delta _{t_1t_2}\mathrm{exp}\left\{i\tau \left(_t^+\theta \right)_{t_2}\right\}\delta _{t_2,t_11}\right]e\delta _{t_1t_2},`$ (35)
will appear in the action. In the latter the $`\theta `$ field appears only in these operators and therefore under derivative, as appropriate to a Goldstone field.
### 2.3 The saddle point expansion
Having proved the Goldstone nature of the excitations associated to the addition or removal of fermionic pairs, we now proceed to deduce the energy spectrum through the saddle point expansion.
After integrating over the fermionic fields $`\psi _m`$ and $`\psi _m^{}`$ we get for the generating functional the expression
$$Z=_0^{\mathrm{}}\left[d\rho \right]_0^\pi \left[d\theta \right]\mathrm{exp}(S_{eff}),$$
(36)
with
$$S_{eff}=\tau \underset{t}{}G\rho \text{Tr}\mathrm{ln}\left(q^{}q^++G^2\rho \right),$$
(37)
where the trace is meant to be taken over the quantum number $`m>0`$ and the time. The $`U(1)`$ symmetry is now realized in the invariance of $`S_{eff}`$ under the substitution
$$\theta \theta +\alpha ,$$
(38)
with $`\alpha `$ time independent.
To perform the saddle point expansion one should first look for a minimum of the effective action (37) at constant fields. Denoting with $`\overline{\rho }`$ the time-independent component of the $`\rho `$ field and defining
$$M=\sqrt{e^2+G^2\overline{\rho }}$$
(39)
and
$$P^1=(1e\tau )_t^+_t^{}+M^2,$$
(40)
we can write $`S_{eff}`$ at constant fields as follows
$$\overline{S}_{eff}=\tau \underset{t}{}G\overline{\rho }\text{Tr}\mathrm{ln}P^1.$$
(41)
By performing the trace over $`m`$ and converting the sum into a time integral, taking first the $`N_0\mathrm{}`$ limit and then letting $`\tau 0`$ (with $`\tau N_0`$ constant), the above becomes
$$\overline{S}_{eff}=\tau N_0\left(G\overline{\rho }+\mathrm{\Omega }e\mathrm{\Omega }M\right),$$
(42)
whose minimum with respect to $`\overline{\rho }`$ occurs when
$$M=\frac{G\mathrm{\Omega }}{2},$$
(43)
or, equivalently, at
$$\overline{\rho }=\overline{\rho }_0\frac{1}{(2G)2}\left[(G\mathrm{\Omega })^24e^2\right]=\frac{\mathrm{\Delta }^2}{G^2}.$$
(44)
It is remarkable that the dimensionless $`\overline{\rho }_0`$, when multiplied by $`G^2`$, coincides with the well-known gap $`\mathrm{\Delta }`$ characterizing the BCS theory in the one single particle level case. The action $`\overline{S}_{eff}`$ at the minimum is then
$$S_0=\tau N_0\left(\frac{M\mathrm{\Omega }}{2}\frac{e^2}{G}+\mathrm{\Omega }e\right).$$
(45)
Next we perform an expansion around the saddle point. We start by defining the fluctuation of the static $`\rho `$-field according to
$$\rho =\overline{\rho }_0+r=\overline{\rho }_0\left(1+\frac{r}{\overline{\rho }_0}\right)$$
(46)
and by considering the generating functional (36), now written as
$$Z=_{\overline{\rho }_0}^{\mathrm{}}\left[dr\right]_0^\pi \left[d\theta \right]\mathrm{exp}(S_{eff}).$$
(47)
Obviously this expansion is justified only if the quantum fluctuations in (46) are small. We will show at the end of this Section that this is indeed the case.
To proceed further we rewrite $`S_{eff}`$ in the form
$$S_{eff}=\tau \underset{t}{}G\left(\overline{\rho }_0+r\right)+\text{Tr}\mathrm{ln}P\text{Tr}\mathrm{ln}\left[11+P\left(R_1+R_2\right)\right],$$
(48)
where
$$R_1=q^{}q^++(_t^++e)(_t^{}e)$$
(49)
and
$$R_2=G^2r.$$
(50)
We set then
$$S_{eff}=\underset{r=0}{\overset{\mathrm{}}{}}S_r,$$
(51)
the term $`S_0`$ being the saddle point contribution, given by (45). This grows like $`\mathrm{\Omega }^2`$, but it contains also a term of order $`\mathrm{\Omega }`$ and a term of order one, which should be kept if an expansion in powers of $`1/\mathrm{\Omega }`$ is sought for. However, we prefer to stick to the definition (40) for the operator $`P`$ and to compute the further contributions to the expansion (51) (the quantum fluctuations) by developing the last logarithm in the right-hand-side of (48): the terms thus obtained are naturally organized in powers of $`M^1`$. It is worth noticing that this expansion does not break the $`U(1)`$ invariance.
Actually we shall confine ourselves to consider the first and second order contributions in $`_t\theta `$ and $`r`$.
The first order action stems from the term linear in $`r`$ and from the first term in the expansion of the logarithm and reads
$$S_1=\tau G\underset{t}{}r_t\text{Tr}\left[P\left(R_1+R_2\right)\right].$$
(52)
The explicit computation shows that all the terms linear in $`r`$ cancel out: hence the $`r`$-integration remains undefined. However the contributions arising from the second term in the expansion of the logarithm make the integral over $`r`$ well defined. The second order contribution will not be reported here and the details can be found in Ref. . One is thus left with the result
$$S_1=\text{Tr}(PR_1)=\frac{1}{G}\left[1+(M+\frac{3}{2}e)\tau \right]\tau \underset{t=\mathrm{}}{\overset{\mathrm{}}{}}\theta (_t^+_t^{})\theta .$$
(53)
Note that these contributions are of order $`\mathrm{\Omega }`$ and 1, namely are $`O(1/\mathrm{\Omega })`$ with respect to the saddle point one.
In order to obtain the Goldstone boson energies we must find out how they depend upon the single particle energy $`e`$. For this purpose we have to perform in the integral expressing the generating functional $`Z_1`$ (associated with the action $`S_1`$) the $`\theta `$-integration, which is not gaussian, because $`\theta `$ is compact. Yet we can choose $`_t\theta `$ as a new integration variable, thus rendering the integral gaussian, getting
$$\frac{1}{N_0\tau }\mathrm{ln}Z_1=\frac{3}{4}e+\frac{M}{2}=\frac{3e+G\mathrm{\Omega }}{4}.$$
(54)
Finally, in order to deal with a specific system (e.g. a nucleus), the particle number must be fixed. This can be accomplished by replacing $`e`$ with
$$ϵ=e\mu ,$$
(55)
$`\mu `$ being the chemical potential and using
$$<\widehat{N}>=\frac{1}{\tau N_0}\frac{}{\mu }\mathrm{ln}Z=\frac{1}{\tau N_0}\frac{}{ϵ}\mathrm{ln}Z,$$
(56)
where $`\widehat{N}`$ is the particle number operator. Now, replacing $`<\widehat{N}>`$ with $`2n`$ and noticing that $`M`$ does not depend on $`\mu `$ (since it is independent of $`e`$, see Eq. (43)), Eq. (56) becomes
$$n=\frac{1}{N_0\tau }\frac{}{ϵ}(S_0\mathrm{ln}Z_1)=\frac{ϵ}{G}+\frac{\mathrm{\Omega }+3/4}{2},$$
(57)
which gives
$$\mu =G\left(n\mathrm{\Omega }/23/8\right)+e$$
(58)
for the chemical potential. Hence, in the presence of the chemical potential, the energy of the system becomes
$`E_{n,0}`$ $`=`$ $`{\displaystyle \frac{1}{\tau T}}(S_0\mathrm{ln}Z_1)+2\mu n`$ (59)
$`=`$ $`2enGn(\mathrm{\Omega }n+3/4)+{\displaystyle \frac{G}{8}}\left(5\mathrm{\Omega }+{\displaystyle \frac{9}{8}}\right),`$
which reproduces the excitation spectrum of the pairing hamiltonian with good accuracy, the relative difference between the exact ground state energy $`G[(\mathrm{\Omega }+1)/2e/G]^2`$ (namely Eq. (10) with $`s=0`$ and $`n=n_0`$) and the one corresponding to Eq. (59) being of order $`1/\mathrm{\Omega }`$.
Note that the same result can be obtained, rather than through the chemical potential, by introducing in the path integral the particle number projection operator
$$𝒫_n=_\pi ^{+\pi }\frac{d\alpha }{2\pi }e^{i(\widehat{N}2n)\alpha },$$
(60)
as shown in Ref. .
An important comment is in order on the validity of our expansion, which depends upon the size of $`\overline{\rho }_0`$. We assess the latter by replacing in (44) $`e`$ by $`ϵ`$ and using (58) (dropping the irrelevant term -3/8), thus obtaining
$$\overline{\rho }_0=n\left(\mathrm{\Omega }n\right),$$
(61)
which attains its maximum value $`\overline{\rho }_0=\mathrm{\Omega }^2/4`$ for $`n=\mathrm{\Omega }/2`$. This corresponds to the situation when the level where the pairs live is half-filled. When this situation is almost realized, namely when the shell is neither fully filled nor almost empty, the functional integral defining $`Z`$ becomes gaussian and an expansion in $`r/\overline{\rho }_0`$ can be performed. On the other hand $`\overline{\rho }_0`$ attains its lowest value when $`n=1`$ or $`n=\mathrm{\Omega }`$. In this case there is no shift in (46) and the $`\rho `$ field acts only through its fluctuations, which are small, thus supporting the validity of the expansion. It is then not surprising that our approach yields the (almost) correct excitation spectrum for any value of the pair number.
### 2.4 The hamiltonian of the s-bosons
We are now in a position of deriving the bosonic hamiltonian corresponding to the effective action previously obtained. The most general, particle conserving, quartic hamiltonian for a system of $`s`$-bosons, confined to live in one single particle level, reads in normal form (we omit for brevity the indices specifying the single particle levels on which the $`\widehat{b}`$ operators act)
$$H(\widehat{b}^{},\widehat{b})=h_1\widehat{b}^{}\widehat{b}+h_2\widehat{b}^{}\widehat{b}^{}\widehat{b}\widehat{b},$$
(62)
$`\widehat{b}^{},\widehat{b}`$ being bosonic creation-annihilation operators acting in a Fock space and satisfying canonical commutation relations.
The parameters $`h_1,h_2`$ should be fixed by requiring that the spectrum of (62) coincides with the pairing Hamiltonian one. This is indeed possible and it turns out that the choice
$$h_1=2eG\mathrm{\Omega },h_2=G$$
(63)
accomplishes the job. However the Hamiltonian thus obtained, being intrinsically bosonic, patently violates the Pauli principle and therefore the condition $`n<\mathrm{\Omega }`$ should be added a posteriori, when using (62) in dealing with a system of fermions.
On the other hand in our framework this condition naturally emerges. Indeed if we write the path integral associated to (62)
$$Z=[db^{}db]\mathrm{exp}(S),$$
(64)
where
$$S=\tau \underset{t=1}{\overset{N_0}{}}\left[b_{t+1}^{}_tb_t+H(b_{t+1}^{},b_t)\right],$$
(65)
$`b^{},b`$ being holomorphic variables satisfying periodic boundary conditions in time, and perform the same saddle point expansion previously illustrated, from the comparison with our effective action we get
$$h_1=2eG\mathrm{\Omega }G/4,h_2=G,n<\mathrm{\Omega }.$$
(66)
The inequality in the above equation, expressing the Pauli principle, is necessary for the two path integrals (36) and (64) to coincide and follows from the positivity of $`\overline{\rho }_0`$ as given in Eq. (61): thus in our approach this condition, far from being artificial, is necessarily implied by the formalism itself. Obviously the considerations following (51) hold valid here as well, hence $`h_1`$ will be affected by an error of order $`1/M`$. Accordingly the values (66) essentially coincide with (63), thus giving the exact spectrum.
Summarizing, we have constructed an effective bosonic Hamiltonian, namely
$$H(\widehat{b}^{},\widehat{b})=(2eG\mathrm{\Omega })\widehat{b}^{}\widehat{b}+G\widehat{b}^{}\widehat{b}^{}\widehat{b}\widehat{b},$$
(67)
which yields the exact zero-seniority pairing spectrum and naturally incorporates, when treated within the path integral approach, the Pauli condition $`n<\mathrm{\Omega }`$, keeping track of the fermionic nature of the composites.
### 2.5 Seniority excitations: a matrix approach
In this Section we address the problem of the whole spectrum of the degenerate pairing Hamiltonian, including the excited states with non zero seniority $`v`$. For convenience we set $`e=0`$ and consider only the interaction Hamiltonian $`H_P`$.
The path integral approach followed in the $`v=0`$ case has not been pursued till now for $`v0`$ since it becomes quite cumbersome. Hence in this case it is preferable to work in the hamiltonian formalism.
In the framework of the creation and annihilation operators the commutator
$$[\widehat{A},\widehat{A}^{}]=\mathrm{\Omega }\left(1\frac{\widehat{n}}{\mathrm{\Omega }}\right)$$
(68)
is non-canonical, thus rendering not trivial to find the eigenstates of $`H_P`$. An approach circumventing this difficulty consists in employing again even Grassmann variables: in so doing the difficulty associated with the non-canonical nature of the commutator (68) disappears, since the associated Grassmann variables do commute.
Within this formalism we study whether and how the Fock basis, set up to with determinants of single particle states, can be reduced to a minimal dimension without loosing the physical information we search for and whether this minimal basis can be expressed in terms of composite bosons .
In seeking for the reduction of the basis dimensions, we shall be guided by the two main features of $`H_P`$, namely that it is expressed solely in terms of the $`\phi `$ and is invariant for any permutation of the $`\phi `$.
These properties, indeed, urge us to express also the vector of the basis in terms of the variables (16). The action of $`H_P`$ on these states is then obtained following the lines illustrated in , the result being
$$H_P\psi (\phi ^{})=[d\phi ^{}d\phi ^{}]K_P(\phi ^{},\phi ^{})e^{{\scriptscriptstyle \phi ^{}\phi ^{}}}\psi (\phi ^{})=E\psi (\phi ^{}),$$
(69)
where the integration is over the even elements of the Grassmann algebra and the kernel reads
$$K_P(\phi ^{},\phi ^{})=H_P(\phi ^{},\phi ^{})e^{{\scriptscriptstyle \phi ^{}\phi ^{}}}.$$
(70)
We then attempt to diagonalize the $`H_P`$ associated with $`n`$ pairs in a basis set up with states represented as products of $`n`$ factors $`\phi ^{}`$’s. Since the number of these is $`\left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }}{n}\right)`$, the very large reduction of the basis dimension entailed by the choice of the variables (16) is fully apparent: indeed in terms of fermionic degrees of freedom, the corresponding basis would have had a dimension $`\left(\genfrac{}{}{0pt}{}{2\mathrm{\Omega }}{2n}\right)`$.
Moreover, since the variables $`\lambda `$ anticommute, each vector of this basis is antisymmetric in the exchange of any pair of fermions, thus fulfilling the Pauli principle.
We now actually explore whether, for a given $`\mathrm{\Omega }`$ and $`n`$, the eigenstates of $`H_P`$ can be cast into the form of a superposition of products of $`n`$ variables $`\phi ^{}`$, namely
$$\psi (\phi ^{})=\underset{m=1}{\overset{\left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }}{n}\right)}{}}\beta _m[\phi _{m_1}^{}\mathrm{}\phi _{m_n}^{}]_m,$$
(71)
the index $`m`$ identifying the set of quantum numbers $`\{m_1,m_2\mathrm{}m_n\}`$ and the $`\beta _m`$ being complex coefficients.
From the condition (11) it follows that the number of distinct eigenvalues of $`H_P`$ is given by $`[\mathrm{min}\{n,\mathrm{\Omega }n\}+1]`$, each of them having a degeneracy
$$\delta _s=\left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }}{s}\right)\left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }}{s1}\right),$$
(72)
$`s`$ being, we recall, the number of broken pairs (in our convention a binomial coefficient with a negative lower index vanishes).
We look for the eigenvalues of $`H_P`$ by diagonalizing the symmetric matrix $`_P(EH_P)/G`$, hence dimensionless, given by
$$\left(\begin{array}{ccc}+n& & 01\\ & \mathrm{}& \\ 01& & +n\end{array}\right).$$
(73)
The dimension of the above is $`\left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }}{n}\right)`$ and the symbol $`01`$ indicates that the upper (lower) triangle of the matrix is filled with zeros and ones. Indeed the matrix elements of $`_P`$ are one when the bra and the ket differ by the quantum state of one (out of $`n`$) pair, otherwise they vanish. The diagonal matrix elements simply count the number of pairs and $`=E/G`$.
We now look for a further dimensional reduction of the basis such that the resulting matrix Hamiltonian has the same eigenvalues of the original one, however all being non degenerate. For this purpose an elementary combinatorial analysis shows that the number of ones in each row (column) of the matrix (73) is given by $`n(\mathrm{\Omega }n)`$. Indeed a non-vanishing matrix element has the row specified by $`n`$ indices whereas, of the indices identifying the column, $`n1`$ should be extracted from those fixing the row in all the possible ways, which amounts to $`n`$ possibilities. The missing index should then be selected among the remaining $`\mathrm{\Omega }n`$ ones: hence the formula $`n(\mathrm{\Omega }n)`$ follows.
Note that the cases with $`n`$ and $`\mathrm{\Omega }n`$ pairs are equivalent, an occurrence which is also manifest in the exact spectrum (10). Hence in the following we shall confine ourselves to consider $`n\frac{\mathrm{\Omega }}{2}`$ only.
To write down explicitly the matrix (73) it is convenient to divide the set of the $`\mathrm{\Omega }`$ even Grassmann variables, whose quantum numbers identify the levels where the $`n`$ pairs are placed, into two subsets: one with $`\mathrm{\Omega }n`$ and the other with $`n`$ elements (to be referred to as I and II, respectively). This partition leads to a basis with dimension $`n+1`$, independently of the value of $`\mathrm{\Omega }`$. The associated eigenvalues correspond to the breaking of 0, 1, 2 $`\mathrm{}n`$ pairs.
Indeed in this instance the $`\left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }}{n}\right)`$ entries of each row and column of the matrix can then be grouped into $`n+1`$ sets, the first one corresponding to the $`n`$ pairs placed in the $`\mathrm{\Omega }n`$ levels of I, the remaining $`n`$ levels of II being empty (see Fig. 1a).
The number of configurations belonging to this first set is $`d_0=\left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }n}{n}\right)`$. To the second set are associated configurations with $`n1`$ pairs in I and one pair in II (see Fig. 1b), their number being $`d_1=n\left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }n}{n1}\right)`$. In general the $`(k+1)`$-th set embodies configurations with $`nk`$ pairs in I and $`k`$ pairs in II, their number being
$$d_k=\left(\genfrac{}{}{0pt}{}{n}{k}\right)\left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }n}{nk}\right)\text{with}0kn(k\text{integer}).$$
(74)
Clearly the total number of configurations is
$$\underset{k=0}{\overset{n}{}}d_k=\left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }}{n}\right).$$
(75)
With this organization of the levels the matrix (73) splits into $`(n+1)^2`$ rectangular blocks $`B_{kj}`$ (with $`0k,jn`$) of dimension $`d_k\times d_j`$ and, since the pairing Hamiltonian only connects states differing by the quantum number of one pair, the blocks with $`|kj|2`$ will have vanishing elements: the matrix thus becomes block-tridiagonal.
This suggests the introduction of the following orthonormal set of $`n+1`$ even commuting variables (the composites)
$$\mathrm{\Phi }_k^{}=\frac{1}{\sqrt{d_k}}\underset{m=1+s_{k1}}{\overset{s_k}{}}\left[\phi _{m_1}^{}\mathrm{}\phi _{m_n}^{}\right]_m,$$
(76)
where $`k`$ varies as in (74), $`s_j=_{l=0}^jd_l`$, being $`s_1=0`$ and $`m`$ again identifies the set $`\{m_1,m_2\mathrm{}m_n\}`$. Note that the variables (76) have in general an index of nilpotency higher than one.
The definition (76) also reflects our desire that the composite bosons keep as much as possible of the symmetry of $`H_P`$. And indeed the $`\mathrm{\Phi }_k^{}`$, while not fully symmetric with respect to the interchange of the $`\phi ^{}`$, turn out to be invariant with respect to the interchange of the $`\phi ^{}`$ belonging either to the set I or to the set II. This symmetry property enforces the maximum coherence among the components of $`\mathrm{\Phi }_k^{}`$. It is remarkable that composite variables corresponding to combinations of the $`\phi _{m_1}^{}\phi _{m_2}^{}\mathrm{}\phi _{m_n}^{}`$ different from (76) not only hold a lower symmetry than the one displayed by (76), but may also lead, as we have verified in some instances, to the wrong eigenvalues.
Now in the minimal basis (76) $`_P`$ is represented by a $`(n+1)\times (n+1)`$ matrix whose generic element $`\left(M_n\right)_{ki}<\mathrm{\Phi }_k^{}|_P|\mathrm{\Phi }_i^{}>`$ obtains by summing the elements of the block $`B_{ki}`$ of the $`\left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }}{n}\right)\times \left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }}{n}\right)`$ matrix, but for the normalization factor $`1/\sqrt{d_kd_i}`$. The sum is performed by recognizing that all the blocks have the same number of ones in each row. Specifically, in the diagonal block $`B_{kk}`$ the number of “ones” in each row is $`(nk)(\mathrm{\Omega }2n+2k)`$. In the upper diagonal block $`B_{k,k+1}`$ each row contains instead
$$c_k=(nk)^2$$
(77)
ones. Since the total number of ones in each row of the matrix (73) is $`n(\mathrm{\Omega }n)`$, the number of ones in the rows of the lower diagonal block $`B_{k,k1}`$ will be
$$b_k=k(\mathrm{\Omega }2n+k).$$
(78)
As a consequence the non vanishing elements of the matrix $`M_n`$ turn out to be
$`\left(M_n\right)_{k,k+1}`$ $`=`$ $`<\mathrm{\Phi }_k^{}|_P|\mathrm{\Phi }_{k+1}^{}>=\sqrt{{\displaystyle \frac{d_k}{d_{k+1}}}}c_k,`$ (79)
$`\left(M_n\right)_{k+1,k}`$ $`=`$ $`<\mathrm{\Phi }_{k+1}^{}|_P|\mathrm{\Phi }_k^{}>=\sqrt{{\displaystyle \frac{d_{k+1}}{d_k}}}b_{k+1}`$ (80)
and
$$\left(M_n\right)_{kk}=<\mathrm{\Phi }_k^{}|_P|\mathrm{\Phi }_k^{}>=a_k+n+n(\mathrm{\Omega }n)b_kc_k.$$
(81)
Clearly $`\left(M_n\right)_{k,k+1}=\left(M_n\right)_{k+1,k}`$, since the operator $`_P`$ is Hermitian and the basis (76) is orthonormal. The matrix thus becomes tridiagonal, reading
$`M_n=\left(\begin{array}{ccccccc}a_0& \sqrt{\frac{d_0}{d_1}}c_0& 0& & & & \\ \sqrt{\frac{d_0}{d_1}}c_0& a_1& \sqrt{\frac{d_1}{d_2}}c_1& 0& & & \\ 0& \sqrt{\frac{d_1}{d_2}}c_1& a_2& \sqrt{\frac{d_2}{d_3}}c_2& 0& & \\ & & & & & & \\ & 0& \sqrt{\frac{d_{k1}}{d_k}}c_{k1}& a_k& \sqrt{\frac{d_k}{d_{k+1}}}c_k& 0& \\ & & & & & & \\ & & & & 0& \sqrt{\frac{d_{n1}}{d_n}}c_{n1}& a_n\end{array}\right)`$ (89)
and the associated eigenfunctions should be expanded in terms of the composite variables (76), namely
$$\psi (\mathrm{\Phi }^{})=\underset{k=0}{\overset{n}{}}u_k\mathrm{\Phi }_k^{}.$$
(90)
It is convenient to cast the coefficients of the expansion (90), fixed by the eigenvalue equation
$$M_n\stackrel{}{u}=0,$$
(91)
into the form
$$u_k\sqrt{d_k}w_k.$$
(92)
Thus the $`\left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }}{n}\right)`$ equations for the $`\beta `$’s reduce to $`n+1`$ equations for the $`w`$’s and the associated eigenvalues $``$ obey the secular equation
$$D_n=det\left(M_n\right)=0,$$
(93)
where $``$ enters into the diagonal matrix elements $`a_k`$.
Now from the general theory of symmetric tridiagonal matrices one knows that Eq. (93) has $`n+1`$ distinct and real roots and these are found by applying the recursive relation
$$D_n=a_nD_{n1}\frac{d_{n1}}{d_n}(c_{n1})^2D_{n2}$$
(94)
for increasing values of $`n`$. Hence we have
$`D_0`$ $`=`$ $`y`$
$`D_1`$ $`=`$ $`y(y\mathrm{\Omega })`$
$`D_2`$ $`=`$ $`y(y\mathrm{\Omega })[y2(\mathrm{\Omega }1)]`$
$`\mathrm{}`$
$`D_n`$ $`=`$ $`y(y\mathrm{\Omega })[y2(\mathrm{\Omega }1)][y3(\mathrm{\Omega }2)]\mathrm{}[yn(\mathrm{\Omega }n+1)],`$ (95)
where
$`y`$ $`=`$ $`n(\mathrm{\Omega }n)++n.`$ (96)
We thus see that $`D_n`$ has all the zeros of $`D_{n1}`$ plus an extra one for $`y=n(\mathrm{\Omega }n+1)`$.
Moreover (95) allows us to write down for the general solution of (93) the expression
$$y=p(\mathrm{\Omega }p+1)\text{with}0pn.$$
(97)
Hence the well-known formula for the spectrum of the pairing Hamiltonian
$$=(np)(\mathrm{\Omega }np+1),$$
(98)
is recovered, the index $`p`$ coinciding with the pair seniority quantum number $`s`$.
Let us now consider the eigenfunctions of $`\widehat{H}_P`$. For a tridiagonal matrix, a recursive relation among the eigenvectors components, similar to (94), can also be established. In the specific case of the matrix $`M_n`$ it reads
$$d_kc_kw_{k+1}=d_{k1}c_{k1}w_{k1}d_ka_kw_k,$$
(99)
where again $`0kn`$ and quantities with negative indices are meant to be zero.
Thus for the lowest eigenvalue $`y=0`$ ($`p=0`$, zero seniority) we have
$$w_0^{(s=0)}=w_1^{(s=0)}=\mathrm{}.=w_n^{(s=0)},$$
(100)
namely the collective state
$$\psi _{s=0}=\left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }}{n}\right)^{1/2}\underset{k=0}{\overset{n}{}}\sqrt{d_k}\mathrm{\Phi }_k^{}=\left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }}{n}\right)^{1/2}\underset{m=1}{\overset{\left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }}{n}\right)}{}}\left[\phi _{m_1}^{}\mathrm{}\phi _{m_n}^{}\right]_m.$$
(101)
Indeed in (101) all the components of the wave-function, i.e. the monomials $`\left[\phi _{m_1}^{}\mathrm{}\phi _{m_n}^{}\right]_m`$, are coherently summed up. This state obtains for a specific partition of the levels defining the matrix $`M_n`$. However, any other partition would lead to the same result, being all the weights of the components equal. As a consequence the state (101) is non degenerate.
Concerning the second eigenvalue $`y=\mathrm{\Omega }`$ ($`p=1`$, pair seniority $`s=1`$), according to Eq. (99) the components of its eigenstate are
$$w_k^{(s=1)}=𝒩_1\left(k\mathrm{\Omega }n^2\right).$$
(102)
Finally the components of the state associated with a generic pair seniority $`s`$ turn out to read
$`w_k^{(s)}`$ $`=`$ $`𝒩_s{\displaystyle \underset{j=0}{\overset{s}{}}}(1)^j{\displaystyle \frac{\left(\genfrac{}{}{0pt}{}{k}{j}\right)\left(\genfrac{}{}{0pt}{}{s}{j}\right)\left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }s+1}{j}\right)}{\left(\genfrac{}{}{0pt}{}{n}{j}\right)^2}}=`$ (103)
$`=`$ $`𝒩_s{}_{3}{}^{}F_{2}^{}(k,\mathrm{\Omega }+s1,s;n,n;1)`$
$`{}_{3}{}^{}F_{2}^{}`$ being a generalized hypergeometric function. In the above the binomial $`\left(\genfrac{}{}{0pt}{}{k}{j}\right)`$ is meant to vanish when $`j>k`$.
In particular for the vector of the maximum seniority $`\stackrel{}{w}^{(s=n)}`$, corresponding to $`y=n(\mathrm{\Omega }n+1)`$ ($`s=n`$), one has
$$w_k^{(s=n)}=𝒩_n\left[(1)^k\frac{(\mathrm{\Omega }2n+k)!(nk)!}{(\mathrm{\Omega }2n)!}\right].$$
(104)
In (102), (103) and (104) $`𝒩_1`$, $`𝒩_s`$ and $`𝒩_n`$ are normalization constants.
It is interesting to reobtain in the present formalism the following important result : for values of $`\mathrm{\Omega }`$ large with respect to $`n`$, the number of dominant components in $`\stackrel{}{w}^{(s)}`$ decreases with $`s`$, reflecting the weakening of collectivity with increasing seniority. Indeed in the limit $`\mathrm{\Omega }>>n`$ the eigenvalues are
$$ys\mathrm{\Omega }$$
(105)
and moreover
$$c_k<<\mathrm{\Omega },\frac{d_{k1}}{d_k}c_{k1}k\mathrm{\Omega }\text{and}a_k\left(sk\right)\mathrm{\Omega },$$
(106)
being $`0sn`$.
By inserting the above limits in (99) the $`w_k^{(s)}`$ components corresponding to the eigenvalues (105) are found to be
$$w_0^{(s)}=w_1^{(s)}=\mathrm{}=w_{s1}^{(s)}=0$$
(107)
and
$$w_k^{(s)}=\frac{\left(\genfrac{}{}{0pt}{}{k}{s}\right)}{\left(\genfrac{}{}{0pt}{}{n}{s}\right)}w_n^{(s)},k=s,\mathrm{},n.$$
(108)
We thus see that the state with seniority $`s`$ has indeed, in the basis of the $`\sqrt{d_k}\mathrm{\Phi }_k^{}`$ and in the large $`\mathrm{\Omega }`$ limit, $`s`$ vanishing components, the remaining $`ns+1`$ ones being expressed, through (108), via the single component $`w_n^{(s)}`$. In other words, when $`\mathrm{\Omega }`$ is large, the collectivity of a state with seniority $`s`$ decreases as $`s`$ increases, because its components become fewer and fewer and, furthermore, the surviving components are more and more expressible through a single one.
## 3 The non degenerate case
We now turn to explore the case of the pairing Hamiltonian $`\widehat{H}`$ acting on a set of non degenerate single particle levels. The actual situation one faces in applying the pairing Hamiltonian to real systems corresponds indeed to having more than one energy $`e_\nu `$ in (1). In this respect nuclei and metals represent two extreme situations of the non-degenerate case: in the former a major shell is typically split into five or six single particle levels of different angular momenta, in the latter the number of non-degenerate levels entering into a band corresponds to a significant fraction of the Avogadro number. Moreover in a heavy nucleus the number of pairs living in a level may be as large as, say, eight while in a metal is one. Recently, a renewed and widespread interest for the pairing problem in the non-degenerate frame has flourished in connection with the physics of ultrasmall metallic grains, possibly superconducting , and of Bose-Einstein condensation .
As we shall illustrate, in the non-degenerate case one more type of excitations occurs with respect to the degenerate one, corresponding to promoting pairs above the Fermi sea: these are zero seniority excitations which, in the BCS language, are described as four quasi-particle states (more precisely, as two quasi-particles and two quasi-holes states). These are conveniently classified in terms of a new quantum number, the “like-seniority”, and reflect the occurrence of a quantum phase transition.
Furthermore, and more importantly, in the non-degenerate case critical values $`G_{\mathrm{cr}}`$ of the coupling arise which split the physics into a region governed by the mean field (when $`G<G_{\mathrm{cr}}`$) and a region governed by the pairing force (when $`G>G_{\mathrm{cr}}`$). This competition does not show up in the degenerate case.
### 3.1 The Richardson exact equations
In this Section we derive in an alternative way the well-known Richardson equations , using the Hamiltonian formalism expressed in the framework of even Grassmann variables previously introduced. As already mentioned, this is particularly useful for avoiding the non canonical commutators, which would naturally appear in the standard Hamiltonian treatment of the pairing problem.
In the following we consider the simpler case of zero-seniority states, although the approach can be generalized to include seniority excitations by introducing other composite variables, as done for one pair in Ref. .
Starting from the normal kernel of the Hamiltonian (1) written in terms of Grassmann variables
$$H=\underset{\nu =1}{\overset{L}{}}e_\nu \underset{m_\nu =j_\nu }{\overset{j_\nu }{}}\lambda _{j_\nu m_\nu }^{}\lambda _{j_\nu m_\nu }G\underset{\mu ,\nu =1}{\overset{L}{}}\mathrm{\Phi }_\mu ^{}\mathrm{\Phi }_\nu ,$$
(109)
we search for eigenstates of $`n`$ pairs of fermions in the $`s`$-quasibosons subspace as products of $`n`$ factors, namely
$$\psi _n(\mathrm{\Phi }^{})(m)=\underset{k=1}{\overset{n}{}}_k^{}(m),$$
(110)
where
$$_k^{}(m)=\underset{\nu =1}{\overset{L}{}}\beta _\nu ^{(k)}(m)\mathrm{\Phi }_\nu ^{}$$
(111)
is a superposition of $`s`$-quasibosons placed in all the available levels and the index $`m=(m_1,\mathrm{},m_n)`$ labels the unperturbed configuration from where the state develops as the pairing force is switched on. The set of values of $`m`$ corresponds to the possible states available to the system. When no confusion arises the index $`(m)`$ will be dropped.
It is convenient to start from the effective Hamiltonian
$$_{\mathrm{eff}}(\phi ^{},\phi )=\underset{\nu =1}{\overset{L}{}}2e_\nu \underset{m_\nu =1/2}{\overset{j_\nu }{}}\phi _{j_\nu m_\nu }^{}\phi _{j_\nu m_\nu }G\underset{\mu ,\nu =1}{\overset{L}{}}\mathrm{\Phi }_\mu ^{}\mathrm{\Phi }_\nu ,$$
(112)
coincident with (109) in the $`s`$-quasibosons subspace spanned by the states (110). Indeed while terms like $`\lambda ^{}\lambda `$ count the number of particles, $`\phi ^{}\phi \lambda ^{}\lambda ^{}\lambda \lambda `$ counts the number of pairs. The eigenvalue equation can then be written as
$$[d\phi ^{}d\phi ^{}]_{\mathrm{eff}}(\phi ^{},\phi ^{})\mathrm{exp}\left(\underset{\nu ,m_\nu }{}(\phi _{j_\nu m_\nu }^{}+\phi _{j_\nu m_\nu }^{})\phi _{j_\nu m_\nu }^{}\right)\psi _n(\mathrm{\Phi }^{})=E_n\psi _n(\mathrm{\Phi }^{}),$$
(113)
since in the expansion of the exponentials only the even powers, hence only the $`\phi `$ variables, survive. By performing the integrals over the $`\phi ^{}`$s one gets then
$$E_n(m)=\underset{k=1}{\overset{n}{}}E_k(m),\beta _\mu ^{(k)}(m)=\frac{C_k(m)}{2e_\mu E_k(m)},$$
(114)
where $`C_k(m)`$ are normalization factors (see for their expression) and $`E_k(m)`$ (henceforth referred to as “pair energies”) are the solutions of the non-linear system
$$\underset{\mu =1}{\overset{L}{}}\frac{\mathrm{\Omega }_\mu }{2e_\mu E_k}\underset{l=1,lk}{\overset{n}{}}\frac{1}{E_lE_k}=\frac{1}{G},k=1,\mathrm{}n.$$
(115)
We thus see that the Grassmann variables formalism enables us to recover the same equations obtained within, for instance, the quasi-spin framework and usually known as Richardson’s equation.
Since the equations (115) deal with pairs of fermions (e.g. nucleons) coupled to an angular momentum $`J=0`$, their eigenvalues, given by $`E=_kE_k`$, are those of the zero-seniority states ($`v=0`$). Importantly, these eigenvalues display different degrees of collectivity: hence they are conveniently classified in terms of the latter.
To clarify this point we introduce a number $`v_l`$, that counts in a given state the number of particles prevented to take part into the collectivity, not because they are blind to the pairing interaction (indeed they are coupled to $`J=0`$), but because they remain trapped in between the single particle levels, even in the strong coupling regime. This number is directly linked to the number $`N_G`$, first introduced by Gaudin in Ref. , through the relation $`v_l=2N_G`$ and might be considered as a sort of “like-seniority” (hence the notation $`v_l`$<sup>1</sup><sup>1</sup>1Note however that, whereas the seniority $`v`$ counts the fermions coupled to $`J0`$, the number $`v_l`$ refers to $`J=0`$ pairs. since it reduces to the standard seniority $`v`$ for large $`G`$: as pointed out in , $`N_G`$ has the significance of the number of pair energies which remain finite as $`G`$ goes to infinity.
Specifically we shall ascribe the value $`v_l`$= 0 to the fully collective state, $`v_l`$=2 to a state set up with a trapped pair energy while the others display a collective behavior, $`v_l`$=4 to the state with two trapped pair energies and so on.
Worth noticing is that in the degenerate case $`e_\mu =e`$ the result (10) for vanishing seniority is easily recovered by multiplying Eq. (115) by $`(2eE_k)`$ and summing over $`k`$. Indeed, exploiting the identity
$$\underset{k=1}{\overset{n}{}}\underset{l(k)=1}{\overset{n}{}}\frac{E_k}{E_lE_k}=\frac{1}{2}n(n1),$$
(116)
one immediately gets
$$E=2enGn\left(\underset{\mu =1}{\overset{L}{}}\mathrm{\Omega }_\mu n+1\right),$$
(117)
which coincides with (10) for $`s=0`$.
In the non degenerate case an exact analytic solution of the Richardson system cannot be obtained and one has to resort to numerical calculations. However, approximate analytic expressions for both the ground and the excited states can be obtained and they help in understanding the occurrence of critical phenomena. This will be illustrated in the next two Sections for the cases of 1 and 2 pairs.
### 3.2 One pair
We start by addressing the problem of one pair living in many single particle levels and show that, even in this case, a critical value of the coupling constant $`G`$ can occur which separates two different regimes .
For $`n=1`$ the Richardson system (115) reduces to the equation
$$\underset{\nu =1}{\overset{L}{}}\frac{\mathrm{\Omega }_\nu }{2e_\nu E}=\frac{1}{G},$$
(118)
which yields $`L`$ eigenvalues $`E(m)`$ ($`1mL`$), while the wave function
$$\psi (\mathrm{\Phi }^{})(m)=\underset{\nu =1}{\overset{L}{}}\beta _\nu (m)\mathrm{\Phi }_\nu ^{}$$
(119)
is characterized by the coefficients $`\beta _\nu (m)=C(m)/(2e_\nu E(m))`$, according to Eq. (114).
It is straightforward to solve equation (118) numerically: the solutions can be graphically displayed as the intersections of the left-hand-side of (118) with the straight line $`E=1/G`$ . Two classes of states appear: the first one embodies the lowest energy state, which lies below the lowest single particle level for an attractive interaction and represents a collective state when $`G`$ is large with respect to the single particle energies; the other contains the so-called “trapped” solutions, which lie in between the single particle levels.
Concerning the wave function, the one corresponding to the lowest eigenvalue, if $`G`$ is large enough to organize a collective motion, has all the $`\beta `$ coefficients sizable, reflecting the high degree of collectivity of the state. On the contrary, in the eigenstates of the trapped solutions only one $`\beta `$-coefficient dominates, so that their wave functions are close to the ones of an unperturbed state.
Let us now consider the collective state in the strong coupling regime. Here it is convenient to recast Eq. (118) as
$$\underset{\nu }{}\frac{\mathrm{\Omega }_\nu }{2e_\nu E}=\frac{1}{E2\overline{e}}\underset{\nu }{}\frac{\mathrm{\Omega }_\nu }{12{\displaystyle \frac{e_\nu \overline{e}}{E2\overline{e}}}}=\frac{1}{G}$$
(120)
and expand in the parameter $`2(e_\nu \overline{e})/(E2\overline{e})`$. Defining
$$\overline{e}=\frac{_\nu \mathrm{\Omega }_\nu e_\nu }{\mathrm{\Omega }}$$
(121)
in leading order we get
$$E_0=2\overline{e}\mathrm{\Omega }G,\mathrm{\Omega }=\underset{\nu }{}\mathrm{\Omega }_\nu ,$$
(122)
which coincides with the degenerate case value. Owing to the definition (121), the next-to-leading order correction vanishes. To proceed further we rewrite (120) as
$$\frac{\mathrm{\Omega }}{E2\overline{e}}\underset{n=0}{\overset{\mathrm{}}{}}\frac{2^nM^{(n)}}{(E2\overline{e})^n}=\frac{\mathrm{\Omega }}{E2\overline{e}}\underset{n=0}{\overset{\mathrm{}}{}}\left(\frac{G\mathrm{\Omega }}{E2\overline{e}}\right)^nm^{(n)}\alpha ^n=\frac{1}{G},$$
(123)
where the generalized moments of the distribution of single particle levels $`e_\nu `$, namely
$$M^{(n)}=\sigma ^nm^{(n)}=\frac{1}{\mathrm{\Omega }}\underset{\nu }{}\mathrm{\Omega }_\nu (e_\nu \overline{e})^n,$$
(124)
have been introduced together with the dimensionless expansion parameter
$$\alpha =\frac{2\sigma }{G\mathrm{\Omega }}$$
(125)
and the variance
$$\sigma =\sqrt{\frac{1}{\mathrm{\Omega }}\underset{\nu }{}\mathrm{\Omega }_\nu (e_\nu \overline{e})^2}.$$
(126)
The strong coupling regime then corresponds to $`\alpha 1`$. Note that the second moment is the square of the variance $`\sigma `$, whereas the third, $`M^{(3)}=\sigma ^3\gamma `$, and fourth, $`M^{(4)}=(c+3)\sigma ^4`$, moments are related to the skewness $`\gamma `$ and to the kurtosis $`c`$ of the distribution, respectively. Note also that $`\alpha =\sigma M^1`$, being $`M^1`$ the dimensionful expansion parameter introduced in Eq. (43).
In a perturbative frame, setting $`E^{(n)}=E^{(n1)}+\delta `$ and linearizing in $`\delta `$, we get for the lowest energy
$$E=2\overline{e}G\mathrm{\Omega }\left[1+\alpha ^2\gamma \alpha ^3+(1+c)\alpha ^4+𝒪(\alpha ^5)\right],$$
(127)
an expression valid for $`\alpha 1`$ and for any single particle energy distribution, but in particular when the fluctuation of the latter around $`\overline{e}`$ is small.
The evaluation of the trapped eigenvalues can be easily performed numerically. This approach however hides some interesting properties of the solutions, which correspond to a transition from the mean field to the pairing domain, as we shall illustrate in the following.
The number of trapped solutions of the secular equation (118) for one pair and $`L`$ single particle levels is $`L1`$. They can be identified by a quantum number $`\nu =2,\mathrm{}L`$, which labels the level where the pair lives when $`G=0`$ ($`\nu =1`$ corresponds to the lowest unbounded energy). Since they are trapped, namely
$$2e_{\nu 1}<E^{(\nu )}<2e_\nu ,$$
(128)
a new variable $`z^{(\nu )}(0,1)`$ can be defined according to
$$E^{(\nu )}=2e_{\nu 1}+2z^{(\nu )}(e_\nu e_{\nu 1}).$$
(129)
If we now isolate in the secular equation the terms associated with the poles in $`2e_{\nu 1}`$ and $`2e_\nu `$, we get
$$\frac{\mathrm{\Omega }_{\nu 1}}{z^{(\nu )}}+\frac{\mathrm{\Omega }_\nu }{z^{(\nu )}1}=\underset{\mu (\nu ,\nu 1)=1}{\overset{L}{}}\frac{\mathrm{\Omega }_\mu }{{\displaystyle \frac{e_\mu e_{\nu 1}}{e_\nu e_{\nu 1}}}z^{(\nu )}}\frac{2(e_\nu e_{\nu 1})}{G}.$$
(130)
Next we let the discrete variable $`\nu `$ become continuous by means of the Euler-McLaurin formula , which replaces the sum with an integral, i.e.
$$\underset{\nu =a}{\overset{b}{}}f(\nu )\frac{1}{2}f(a)+\frac{1}{2}f(b)+_a^bf(u)𝑑u.$$
(131)
The formula (131) holds valid if $`f(\nu )`$, $`\nu `$ to be viewed as a complex variable, is analytic in the strip $`a\text{Re}\nu b`$ ($`a`$ and $`b`$ being integers).
To proceed further one must specify the single particle energies and the associated degeneracies. Here we shall first consider the case of a harmonic oscillator well, and later generalize the results. The harmonic oscillator unperturbed energies are labeled by an index $`k=0,1,\mathrm{}𝒩1`$ and read
$$e_k\stackrel{~}{e}_k\mathrm{}\omega _0=(k+3/2)\mathrm{}\omega _0,(k=0,1,\mathrm{}𝒩1)$$
(132)
with associated pair degeneracies
$$\mathrm{\Omega }_k=(k+1)(k+2)/2.$$
(133)
The total degeneracy $`\mathrm{\Omega }`$, the average energy $`\overline{e}`$ and the variance $`\sigma `$ are found to be
$$\mathrm{\Omega }=\frac{1}{6}𝒩(𝒩+1)(𝒩+2),$$
(134)
$$\overline{e}=\frac{3}{4}(𝒩+1)\mathrm{}\omega _0$$
(135)
and
$$\sigma ^2=\frac{3}{80}(𝒩1)(𝒩+3)(\mathrm{}\omega _0)^2,$$
(136)
respectively, entailing the following expression for the expansion parameter:
$$\alpha ^{(h.o.)}=\frac{3\mathrm{}\omega _0}{G}\sqrt{\frac{3}{5}}\frac{\sqrt{(𝒩1)(𝒩+3)}}{𝒩(𝒩+1)(𝒩+2)}.$$
(137)
Using the dimensionless quantities $`\stackrel{~}{G}=G/\mathrm{}\omega _0`$ and $`\stackrel{~}{z}=z/\mathrm{}\omega _0`$ we then rewrite (130) in the Euler-McLaurin approximation obtaining ($`k`$ labels the single particle levels)
$$\frac{\mathrm{\Omega }(k1)}{\stackrel{~}{z}(k)}+\frac{\mathrm{\Omega }(k)}{\stackrel{~}{z}(k)1}=\phi ^{\mathrm{E}\mathrm{McL}}(k,\stackrel{~}{z}(k))$$
(138)
with
$`\phi ^{\mathrm{E}\mathrm{McL}}(k,\stackrel{~}{z})`$ $`=`$ $`{\displaystyle \frac{2}{\stackrel{~}{G}}}{\displaystyle \frac{3}{4}}(4k+2\stackrel{~}{z}+3)+{\displaystyle \frac{1}{4}}(𝒩1)(𝒩+2k+2\stackrel{~}{z}+3)`$ (139)
$`+`$ $`{\displaystyle \frac{k(1k)}{4(\stackrel{~}{z}+1)}}{\displaystyle \frac{(k+2)(k+3)}{4(\stackrel{~}{z}2)}}`$
$``$ $`{\displaystyle \frac{1}{2(k+\stackrel{~}{z}1)}}{\displaystyle \frac{𝒩(𝒩+1)}{4(k+\stackrel{~}{z}𝒩)}}`$
$`+`$ $`{\displaystyle \frac{1}{2}}(k+\stackrel{~}{z})(k+\stackrel{~}{z}+1)\mathrm{log}\left|{\displaystyle \frac{(\stackrel{~}{z}+1)(k+\stackrel{~}{z}𝒩)}{(\stackrel{~}{z}2)(k+\stackrel{~}{z}1)}}\right|.`$
The Euler-McLaurin approximation provides the key for studying analytically the limiting case $`𝒩\mathrm{}`$, which is helpful in shedding light on the properties of the trapped eigenvalues. For this purpose we define
$$\lambda =\frac{k}{𝒩},$$
(140)
which, when $`𝒩`$ is large, tends to become a continuous variable in the interval $`(0,1)`$.
Moreover we express the coupling constant $`\stackrel{~}{G}`$ in terms of $`\alpha `$ according to (137). Then Eq. (138) becomes
$$\varphi (\stackrel{~}{z})=\xi _{\mathrm{}}(\lambda )$$
(141)
with
$$\varphi (\stackrel{~}{z})\frac{1}{\stackrel{~}{z}}+\frac{1}{\stackrel{~}{z}1}+\frac{1}{2(\stackrel{~}{z}+1)}+\frac{1}{2(\stackrel{~}{z}2)}\mathrm{log}\left|\frac{\stackrel{~}{z}+1}{\stackrel{~}{z}2}\right|$$
(142)
and
$$\xi _{\mathrm{}}(\lambda )\frac{1}{\lambda }+\mathrm{log}\frac{1\lambda }{\lambda }+\frac{1\frac{8}{3}\sqrt{\frac{5}{3}}\alpha }{2\lambda ^2}.$$
(143)
The numerical solutions of Eq. (141) are displayed in Fig. 2 for different values of $`\alpha `$. It is interesting to follow the behavior with $`\alpha `$ of the curves, from the $`\alpha 0`$ case (strong coupling), which still carries the fingerprints of the harmonic oscillator, to the $`\alpha \mathrm{}`$ one (weak coupling), which corresponds to the straight line $`\stackrel{~}{z}=1`$.
The behavior of the curves is ruled both by the function $`\xi _{\mathrm{}}(\lambda )`$, which goes to $`\mathrm{}`$ when $`\lambda 1`$, whereas for $`\lambda 0`$
$$\xi _{\mathrm{}}(\lambda )\stackrel{\lambda 0}{}\{\begin{array}{cc}+\mathrm{}& \mathrm{for}\alpha <\alpha _{\mathrm{cr}}\\ \mathrm{}& \mathrm{for}\alpha >\alpha _{\mathrm{cr}}\end{array}\text{with}\alpha _{\mathrm{cr}}=\frac{3}{8}\sqrt{\frac{3}{5}}0.29,$$
(144)
and by the monotonic decrease of $`\varphi (\stackrel{~}{z})`$ in the interval $`(0,1)`$, at whose endpoints it assumes the values
$$\underset{\stackrel{~}{z}0}{lim}\varphi (\stackrel{~}{z})=+\mathrm{}\text{and}\underset{\stackrel{~}{z}1}{lim}\varphi (\stackrel{~}{z})=\mathrm{}.$$
(145)
Thus, when $`\lambda 1`$,
$$\stackrel{~}{z}(\lambda )1+\frac{1}{4\mathrm{log}(1\lambda )}$$
(146)
no longer depends upon $`\alpha `$ and all the curves in Fig. 2 coalesce to 1. Indeed since the pairing interaction is of finite range, a pair trapped in highly excited harmonic oscillator states has the two partners sufficiently de-localized to be little affected by the interaction. In this connection it is worth reminding that the classical limit is achieved by letting the degeneracy of the single particle levels become very large . Thus the coincidence of all the eigenvalues, for any $`G`$, in $`\lambda `$=1 also reflects the evolution from quantum to classical mechanics of our system.
For $`\lambda 0`$ two cases occur: if $`\alpha <\alpha _{\mathrm{cr}}`$ then $`\stackrel{~}{z}(\lambda )0`$, whereas if $`\alpha >\alpha _{\mathrm{cr}}`$ then $`\stackrel{~}{z}(\lambda )1`$. Thus a transition occurs at $`\alpha _{\mathrm{cr}}`$, as illustrated in Fig. 2: the almost parabolic behavior of $`\stackrel{~}{z}(\lambda )`$ for small $`\alpha `$ (which was first observed in Ref. ) is strongly distorted for $`\alpha \alpha _{\mathrm{cr}}`$ for small $`\lambda `$. For larger $`\alpha `$ a smoother behavior is recovered. In particular in Fig. 2 a marked minimum is seen to develop for $`\alpha `$ above, but close to, the critical value.
The convergence of all the curves of Fig. 2 to $`\stackrel{~}{z}(0)=0`$ for $`\alpha <\alpha _{\mathrm{cr}}`$ reflects the pressure exercised by the infinite number of the high-lying, large degeneracy, levels on the low-lying, low-degeneracy, ones. This occurrence might be understood on the basis of a sum rule the trapped solutions should fulfill and of the nature of the pairing force. Actually the sum of the eigenvalues $`\stackrel{~}{z}(\lambda )`$ can be exactly computed in the two limiting cases $`\alpha =0`$ and $`\alpha =\mathrm{}`$: from the Viète equations one indeed has
$$\mathrm{\Sigma }(\alpha =0)\frac{1}{𝒩1}\underset{k=1}{\overset{𝒩1}{}}\stackrel{~}{z}(k)=\frac{1}{4}$$
(147)
and
$$\mathrm{\Sigma }(\alpha =\mathrm{})=1,$$
(148)
respectively. These values set the limits for the area under the curves of Fig. 2. The existence of a critical point reflects the fact that, since the action of the pairing force is gauged by the product $`G\mathrm{\Omega }_k`$, above some critical value of $`\alpha `$ the system prefers to obey the sum rule by lifting the lowest eigenvalues (corresponding to the lowest degeneracies) to the unperturbed values. Indeed, as we shall illustrate, if the degeneracy of the single particle levels is not growing fast enough with $`k`$, then no transition occurs.
The case of finite $`𝒩`$ will not be discussed here. We just mention that, as proved in Ref. , the eigenvalues obtained in the $`𝒩\mathrm{}`$ are very robust with respect to variations of $`𝒩`$, when $`G`$ is large: indeed they keep their validity even for values of $`𝒩`$ as small as five.
It is now natural to ask whether the transition previously discussed is peculiar of the harmonic oscillator model or is more general. To answer this question, we take $`𝒩`$ large and consider two classes of models with equally spaced single particle energies, as in the harmonic oscillator case, but degeneracies growing (model $`a`$) or decreasing (model $`b`$) with $`k`$, namely
$`\mathrm{\Omega }_k^{(a)}`$ $`=`$ $`(k+1)^\gamma 𝒩^\gamma \lambda ^\gamma `$ (149)
$`\mathrm{\Omega }_k^{(b)}`$ $`=`$ $`\left(𝒩k\right)^\gamma =𝒩^\gamma (1\lambda )^\gamma ,`$ (150)
with $`\gamma 0`$. The associated expansion parameters are then
$$\alpha ^{(a)}=\alpha ^{(b)}=\sqrt{\frac{(\gamma +1)^3}{(\gamma +3)}}\frac{2}{(\gamma +2)\stackrel{~}{G}𝒩^\gamma }.$$
(151)
Note that for $`\gamma `$=2 and $`𝒩`$ large the harmonic oscillator expression for $`\alpha `$ is recovered, up to a factor of two.
By performing again the Euler-McLaurin approximation and studying the behavior of the trapped solution a transition is found to occur only for the models $`a`$, with
$$\alpha _{\mathrm{cr}}^{(a)}=\sqrt{\frac{(\gamma +1)^3}{\gamma +3}}\frac{1}{\gamma (\gamma +2)},$$
(152)
clearly behaving as $`1/\gamma `$ when $`\gamma 0`$. On the contrary, for the models $`b`$ no transition can take place. Thus a transition occurs only if the degeneracy grows with $`k`$, and the ‘strong coupling’ domain becomes wider ($`\alpha _{\mathrm{cr}}`$ increases) as $`\gamma `$ approaches zero: here the transition disappears. If the degeneracy decreases with $`k`$ (case b), no transition exists.
### 3.3 Two pairs
We now address the problem of two pairs living in many levels, which is already quite instructive on the general case of many pairs and on the evolution towards the BCS limit.
In the previous Section we have shown that even when the system consists of only one pair a transition between two different regimes is found, a precursor of the quantum phase transition occurring in infinite systems. When the system is made of several pairs, this phenomenon becomes more complicated, since many critical values of the coupling $`G`$ arise. These critical values correspond to a particular “escape” mechanism of the trapped solutions from the grid of the single particle levels, which is necessary for the occurrence of a fully collective state, as we shall now illustrate.
As we shall see, as for the one pair case, a transiton between two different regimes still occurs, but now the number of critical points can be, under suitable conditions, two, due to the larger number of possible configurations. If the number of single particle levels tends to infinity the two critical values of $`G`$ merge into one, which, notably, coincides with the $`G_{\mathrm{cr}}`$ of a one pair system. In correspondence of this $`G_{\mathrm{cr}}`$ the system undergoes a transition from a mean field to a pairing dominated regime.
The Richardson equations (115) reduce, for $`n=2`$, to the following system of two equations
$$\{\begin{array}{ccc}1G\underset{\mu =1}{\overset{L}{}}\frac{\mathrm{\Omega }_\mu }{2e_\mu E_1}+\frac{2G}{E_2E_1}& =& 0\\ 1G\underset{\mu =1}{\overset{L}{}}\frac{\mathrm{\Omega }_\mu }{2e_\mu E_2}+\frac{2G}{E_1E_2}& =& 0.\end{array}$$
(153)
As discussed in Section 3 the solutions are classified in terms of like-seniority: we shall ascribe the value $`v_l`$= 0 to the fully collective state, $`v_l`$=2 to a state set up with one trapped pair energy and one displaying a collective behavior and $`v_l`$=4 to the state with two trapped pair energies.
Let us first consider the weak coupling limit, where of course no collective mode develops (hence $`v_l`$ has no significance). Adopting a perturbative treatment for the pair energies $`E_1`$, $`E_2`$ we write
$$E_i=2e_{\mu _i}+Gx_i,$$
(154)
$`Gx_i`$ being a perturbation. Then
$$\underset{\mu }{}\frac{\mathrm{\Omega }_\mu }{2e_\mu E_i}=\frac{\mathrm{\Omega }_{\mu _i}}{Gx_i}+\underset{\mu \mu _i}{}\frac{\mathrm{\Omega }_\mu }{2(e_\mu e_{\mu _i})Gx_i}$$
(155)
and, expanding in $`G`$, the system (153) becomes
$$\{\begin{array}{ccc}1+\frac{\mathrm{\Omega }_{\mu _1}}{x_1}G\underset{\mu \mu _1}{}\frac{\mathrm{\Omega }_\mu }{2(e_\mu e_{\mu _1})}+\frac{2G}{2(e_{\mu _2}e_{\mu _1})}+O(G^2)& =& 0\\ 1+\frac{\mathrm{\Omega }_{\mu _2}}{x_2}G\underset{\mu \mu _2}{}\frac{\mathrm{\Omega }_\mu }{2(e_\mu e_{\mu _2})}+\frac{2G}{2(e_{\mu _1}e_{\mu _2})}+O(G^2)& =& 0,\end{array}$$
(156)
where the indices $`\mu _1`$ and $`\mu _2`$ select one configuration out of the unperturbed ones. At the lowest order in $`G`$ (weak coupling regime), if $`\mu _1\mu _2`$ (namely if the two pairs sit on different single particle levels when $`G=0`$), one has $`x_i=\mathrm{\Omega }_{\mu _i}`$ and the pair energies $`E_i=2e_{\mu _i}G\mathrm{\Omega }_{\mu _i}`$ are real. Thus the total energy $`E=E_1+E_2`$ becomes
$$E=2(e_{\mu _1}+e_{\mu _2})G(\mathrm{\Omega }_{\mu _1}+\mathrm{\Omega }_{\mu _2}).$$
(157)
If $`\mu _1=\mu _2`$, which is possible only if $`\mathrm{\Omega }_{\mu _1}>1`$, the Richardson system is
$$\{\begin{array}{ccc}1+\frac{\mathrm{\Omega }_{\mu _1}}{x_1}G_{\mu \mu _1}\frac{\mathrm{\Omega }_\mu }{2(e_\mu e_{\mu _1})}+\frac{2}{x_2x_1}+O(G^2)& =0& \\ 1+\frac{\mathrm{\Omega }_{\mu _2}}{x_2}G_{\mu \mu _2}\frac{\mathrm{\Omega }_\mu }{2(e_\mu e_{\mu _2})}+\frac{2}{x_1x_2}+O(G^2)& =0& \end{array}$$
(158)
and from its solution
$$x_{1,2}=(\mathrm{\Omega }_{\mu _1}1)\pm i\sqrt{\mathrm{\Omega }_{\mu _1}1}$$
(159)
the complex conjugate pair energies $`E_1=E_2^{}`$ are obtained. The total energy reads then
$$E=4e_{\mu _1}2G(\mathrm{\Omega }_{\mu _1}1)$$
(160)
and is, of course, real.
In comparing (157) and (160) with the energy of one pair system in the weak coupling regime, namely
$$E=2e_\mu G\mathrm{\Omega }_\mu ,$$
(161)
one sees that, while (157) corresponds to the sum of two contributions like (161), hence the two pairs ignore each other, in (160) a positive (repulsive) energy $`2G`$ associated to the Pauli blocking appears.
Let us now consider the strong coupling domain. We shall deal only with the states $`v_l=0`$ and $`2`$ since the $`v_l=4`$ states are of minor physical interest.
The $`v_l=0`$ state arises from an unperturbed configuration with the two pairs in the lowest single particle level (if $`\mathrm{\Omega }_1>1`$) or in the two lowest ones (if $`\mathrm{\Omega }_1=1`$).
Following the same procedure as in the one pair case, we introduce the variables
$$x_i=\frac{E_i2\overline{e}}{G\mathrm{\Omega }}$$
(162)
and solve the system (153) in terms of the expansion parameter $`\alpha `$ defined in (125) through a recursive linearization, getting
$`x_{1,2}`$ $`=`$ $`{\displaystyle \frac{(\mathrm{\Omega }1)}{\mathrm{\Omega }}}\alpha ^2{\displaystyle \frac{(\mathrm{\Omega }2)}{(\mathrm{\Omega }1)}}+\alpha ^3\gamma {\displaystyle \frac{(\mathrm{\Omega }4)}{(\mathrm{\Omega }1)}}`$ (163)
$`\pm `$ $`i{\displaystyle \frac{\sqrt{\mathrm{\Omega }1}}{\mathrm{\Omega }}}i{\displaystyle \frac{1}{2}}\alpha ^2{\displaystyle \frac{\mathrm{\Omega }}{(\mathrm{\Omega }1)^{3/2}}}\pm i\alpha ^3\gamma {\displaystyle \frac{\mathrm{\Omega }}{(\mathrm{\Omega }1)^{3/2}}}+𝒪(\alpha ^4).`$
Note that the two pair energies are always complex conjugate and the system’s total energy reads
$$E=4\overline{e}2G(\mathrm{\Omega }1)\left[1+\alpha ^2\frac{\mathrm{\Omega }(\mathrm{\Omega }2)}{(\mathrm{\Omega }1)^2}\alpha ^3\gamma \frac{\mathrm{\Omega }(\mathrm{\Omega }4)}{(\mathrm{\Omega }1)^2}+𝒪(\alpha ^4)\right].$$
(164)
Hence the leading order contribution coincides with the degenerate case collective eigenvalue (see Eq. 10), namely
$$E_{\mathrm{deg}}=4\overline{e}2G(\mathrm{\Omega }1),$$
(165)
which therefore represents a good estimate when the spreading of the mean field levels is small with respect to $`G\mathrm{\Omega }`$, as it is natural to expect.
We note that, when $`\mathrm{\Omega }1`$, (164) becomes just twice the value (127) of the collective energy of one pair of nucleons living in $`L`$ levels in the strong coupling regime. Moreover, the imaginary part of $`x_1`$ and $`x_2`$ goes to zero as $`1/\sqrt{\mathrm{\Omega }}`$. Thus, in this limit, the Pauli interaction between the two pairs vanishes, as expected: the two pairs behave as two free quasi-bosons condensed in a level whose energy is given by (127).
Let us now come to the like-seniority $`v_l=2`$ states.
If the coupling term (expressing the Pauli principle) in the Richardson equations (153) were absent, then the eigenvalues of the $`v_l=2`$ states could be simply obtained by adding the collective energy $`E_1`$ carried by one pair and the trapped energy $`E_2`$ carried by the other pair.
This situation is recovered in the very strong coupling limit, where all the single particle energies become essentially equal to $`\overline{e}`$ and both $`\overline{e}`$ and $`E_2`$ are very small with respect to $`E_1`$. The first equation of the system (153) then becomes
$$\frac{1}{G}+\frac{\mathrm{\Omega }}{E_1}\frac{2}{E_1}=0$$
(166)
yielding
$$E_1=G(\mathrm{\Omega }2),$$
(167)
namely the energy of the state with two pairs and $`v=2`$ in the one level problem. This result sets a correspondence between states with $`v_l=2`$ and $`v=2`$, thus connecting seniority and “like-seniority” (or the physics of a ‘broken’ and a ‘trapped’ pair).
In the strong coupling domain however the Pauli term in (153) cannot be neglected, but is well approximated by $`2/\mathrm{\Omega }`$: hence the system (153) decouples and can be recast as
$$\{\begin{array}{ccc}\frac{1}{G_{\mathrm{eff}}^{(1)}}_\mu \frac{\mathrm{\Omega }_\mu }{2e_\mu E_1}& =& 0\\ \frac{1}{G_{\mathrm{eff}}^{(2)}}_\mu \frac{\mathrm{\Omega }_\mu }{2e_\mu E_2}& =& 0,\end{array}$$
(168)
where
$`{\displaystyle \frac{1}{G_{\mathrm{eff}}^{(1)}}}{\displaystyle \frac{1}{G}}+{\displaystyle \frac{2}{E_2^{(\nu )}E_1^{(1)}}}{\displaystyle \frac{1+{\displaystyle \frac{2}{\mathrm{\Omega }}}}{G}}`$ (169)
$`{\displaystyle \frac{1}{G_{\mathrm{eff}}^{(2)}}}{\displaystyle \frac{1}{G}}{\displaystyle \frac{2}{E_2^{(\nu )}E_1^{(1)}}}{\displaystyle \frac{1{\displaystyle \frac{2}{\mathrm{\Omega }}}}{G}}.`$ (170)
Therefore the Pauli principle in the large $`G`$ regime just re-scales the coupling constant in such a way that the pairing interaction is quenched for the collective state and enhanced for the trapped ones.
The results obtained in the strong coupling domain are in agreement with the findings of Ref. , where the problem of the pairing Hamiltonian for small superconducting grains is studied. The main difference with the present approach is that there the single particle levels are non-degenerate: hence the results of Ref. for two pairs can be recovered from ours by setting $`\mathrm{\Omega }_\nu =1`$.
Let us now discuss the existence of one or more critical values for the coupling constant $`G`$.
In the weak coupling regime, when a state evolves from an unperturbed one with the two pairs living in the same level, then the pair energies $`E_1`$ and $`E_2`$ are always complex conjugate. On the other hand when the state evolves from an unperturbed one having the two pairs living in two different levels, then $`E_1`$ and $`E_2`$ are real.
By contrast, in the strong coupling regime the pair energies $`E_1`$ and $`E_2`$ of the $`v_l=0`$ state are always complex conjugate. It is then clear that, if the degeneracy $`\mathrm{\Omega }_1`$ of the lowest single particle level is greater than one, then the pair energies $`E_1`$ and $`E_2`$ are complex conjugate in both the weak and strong coupling regime and their behavior with $`G`$ is smooth.
On the other hand if $`\mathrm{\Omega }_1=1`$, since in the weak coupling limit the two pairs must live on different levels, $`E_1`$ and $`E_2`$ are necessarily real in a neighbourhood of the origin, but become complex in the strong coupling regime. Thus a singularity in their behavior as a function of the $`G`$ is bound to occur.
In this second case it appears natural to surmise that the singularity takes place when $`E_1`$ and $`E_2`$ coincide. In fact in this case the Pauli term of the Richardson equations diverges (reflecting the attempt of four particles to sit in one level that can host only two of them), and it must be compensated by a divergence in the sum entering into the system (153): this can only happen if $`E_1`$ or $`E_2`$ coincides with an unperturbed eigenvalue.
The situation is portrayed in fig. 3, where the behavior of the pair energies with $`\stackrel{~}{G}`$ is displayed for a harmonic oscillator well assuming $`L=3`$.
Clearly in this case two values of $`\stackrel{~}{G}_{\mathrm{cr}}`$ occur. The pair energies $`E_1`$ and $`E_2`$, real in the weak coupling limit, coincide at the critical point $`\stackrel{~}{G}_{\mathrm{cr}}^{(1)+}`$ (their common value being $`2e_1`$) and then become complex conjugate; the energy of the associated state evolves in the $`v_l=0`$ collective mode. By contrast, for the $`v_l=2`$ state, arising from the configuration with two pairs living on the second level (which is allowed for the harmonic oscillator well), the two pair energies $`E_1`$ and $`E_2`$ are complex conjugate in the weak coupling limit, coalesce into the energy $`2e_1`$ at the critical point $`\stackrel{~}{G}_{\mathrm{cr}}^{(1)}`$ and then become real. One of the two solutions remains trapped above $`2e_1`$, while the other evolves into a collective state: the sum of the two yields the energy of the $`v_l=2`$ state.
Analytic expression of the critical values of $`G`$ can be obtained starting from the Richardson equations and assuming $`\mathrm{\Omega }_\nu =1`$. Indeed in Ref. the expression
$$G_{\mathrm{cr}}^{(\nu )\pm }=\frac{1}{𝒫_{(1)\nu }\pm 𝒫_{(2)\nu }}=\left[\underset{\mu =1(\mu \nu )}{\overset{L}{}}\frac{\mathrm{\Omega }_\mu }{2e_\mu 2e_\nu }\pm \sqrt{\underset{\mu =1(\mu \nu )}{\overset{L}{}}\frac{\mathrm{\Omega }_\mu }{(2e_\mu 2e_\nu )^2}}\right]^1,$$
(171)
where $`G^+`$ is the critical value of the coupling related to the ground state and $`G^{}`$ the one related to the excited states, is derived. In the above the quantities
$$𝒫_{(k)\nu }=\left\{\underset{\mu =1(\mu \nu )}{\overset{L}{}}\frac{\mathrm{\Omega }_\mu }{(2e_\mu 2e_\nu )^k}\right\}^{\frac{1}{k}}$$
(172)
are the inverse moments of the level distribution. We thus recover, although through a somewhat different route, the results found long ago by Richardson .
It is important to stress that for the occurrence of a critical value of $`G`$ a single particle level with $`\mathrm{\Omega }=1`$ must exist. Such level is the lowest one in a harmonic oscillator: for this potential well both the ground and the first excited state of a $`n=2`$ system carry one critical value of the coupling constant, namely $`G_{\mathrm{cr}}^{(1)+}`$ and $`G_{\mathrm{cr}}^{(1)}`$, respectively. In correspondence of these values of $`G`$ the pair energies take on the value $`E_1=E_2=2e_1`$. Thus for a $`n=2`$ system two (at most) critical points exist on a $`\mathrm{\Omega }=1`$ level.
Finally, owing to the relevance of the $`\mathrm{\Omega }=1`$ degeneracy, we consider the model of $`L`$, for simplicity equally spaced, single particle levels all having $`\mathrm{\Omega }=1`$, a situation occurring in metals and in deformed nuclei. In this instance two positive $`G_{\mathrm{cr}}`$ always exist in the lowest level when $`L3`$ (in fact $`G_{\mathrm{cr}}^{(1)}\mathrm{}`$ for $`L=2`$). Moreover a positive $`G_{\mathrm{cr}}^{(1)}`$ implies complex $`E_1`$ and $`E_2`$ for $`G<G_{\mathrm{cr}}^{(1)}`$ and, since for small $`G`$ both the pair energies are real, they should evolve from an excited unperturbed configuration having the two pairs placed in the two lower-lying single particle levels, as illustrated in fig.4. Numerically, for this model, we have found that two $`G_{\mathrm{cr}}`$ appear on the second level when $`L9`$ and on the third level when $`L16`$. Thus in the $`\mathrm{\Omega }=1`$ model two pairs can form, so to speak, a quartet only if they live on adjacent single particle levels in the unperturbed configuration. Furthermore the more excited the unperturbed configuration is the more not only $`G`$, but also $`L`$, should be larger for the merging to occur, a fact clearly reflecting the competition between the mean field and the pairing force. Finally we observe that here, at variance with the finding of Ref. , $`G_{\mathrm{cr}}^{(\nu +1)+}>G_{\mathrm{cr}}^{(\nu )+}`$: this is simply because we measure the single particle energies from the bottom of the well rather than from the Fermi surface.
It is interesting to establish the relationship of the above discussed critical values of $`G`$ with the critical value (144) found for one pair, which corresponds to
$$G_{\mathrm{cr}}\frac{8}{L^2}$$
(173)
when $`L`$ is large. For the $`n=2`$ system we first notice that, at large $`L`$, in $`G_{\mathrm{cr}}^{(1)\pm }`$ the moment $`𝒫_1`$ dominates over $`𝒫_2`$: hence, in this condition,
$$G_{\mathrm{cr}}^{(1)+}G_{\mathrm{cr}}^{(1)}.$$
(174)
Moreover, again for $`L`$ large enough,
$$𝒫_1\frac{L^2}{8},$$
(175)
entailing the equality
$$G_{\mathrm{cr}}(n=1)=G_{\mathrm{cr}}(n=2)$$
(176)
in the asymptotic $`L`$ limit. From this outcome one sees that also when $`n=2`$ the relevant dynamical element for $`G<G_{\mathrm{cr}}^{(1)+}`$ is the mean field, whereas for $`G>G_{\mathrm{cr}}^{(1)}`$ it is the pairing force, as far as the system’s ground state is concerned.
### 3.4 Transition amplitudes
Further insight into the critical behavior of the 2-pair system can be gained by studying the pair transfer matrix elements from a one pair to a two pairs state as a function of $`G`$.
This can be achieved by constructing the wave functions according to Eqs. (110), (111), (114) with $`n=1`$ and $`n=2`$ and, with these, calculate the transition amplitudes induced by the operator $`\widehat{A}^{}`$ (see Eq. (2)), namely the matrix elements
$$<n=1|\widehat{A}^{}|n=0>\text{and}<n=2|\widehat{A}^{}|n=1>.$$
(177)
In the fully degenerate case the general expression for these transition amplitudes between states with any number $`n`$ of pairs and seniority $`v`$, can be obtained using the expression (6) for the excitation spectrum. Indeed, inserting into the pairing spectrum a complete set of states $`|n^{},v^{}>`$, one has
$$\underset{n^{},v^{}}{}<n+1,v|\widehat{A}^{}|n^{},v^{}><n^{},v^{}|\widehat{A}|n+1,v>=\left(n+1\frac{v}{2}\right)\left(\mathrm{\Omega }n1\frac{v}{2}\right),$$
(178)
thus getting for the transition amplitude
$$<n+1,v|\widehat{A}^{}|n,v>=\sqrt{\left(n+1\frac{v}{2}\right)\left(\mathrm{\Omega }n\frac{v}{2}\right)}.$$
(179)
Note the $`G`$-independence of (179), which, moreover, is diagonal in the seniority quantum number. From the previously illustrated analogy between seniority and like-seniority, we expect that for large $`G`$ the matrix elements (177) display an asymptotic behavior coinciding with (179), namely
$`<1,v_l=0|\widehat{A}^{}|0>`$ $`\underset{G\mathrm{}}{\overset{}{}}`$ $`\sqrt{\mathrm{\Omega }}`$ (180)
$`<2,v_l|\widehat{A}^{}|1,v_l>`$ $`\underset{G\mathrm{}}{\overset{}{}}`$ $`\{\begin{array}{cc}\sqrt{2(\mathrm{\Omega }1)}\hfill & \text{for }v_l=0\hfill \\ \sqrt{\mathrm{\Omega }2}\hfill & \text{for }v_l=2.\hfill \end{array}`$ (181)
Actually this happens only if $`v_l`$ is conserved. However, the number $`v_l`$ is not conserved at finite $`G`$: only in the strong coupling regime the vanishing of the matrix element between states of different Gaudin numbers occurs. Indeed for finite $`G`$ we get
$$\begin{array}{c}<(m_1,m_2)|\widehat{A}^{}|(m)>=\hfill \\ \hfill \frac{1}{𝒩}\underset{\mu \nu }{}[\stackrel{~}{\beta }_\mu ^{(1)}(m_1,m_2)\stackrel{~}{\beta }_\nu ^{(2)}(m_1,m_2)]^{}\left[\stackrel{~}{\beta }_\mu (m)\sqrt{\mathrm{\Omega }_\nu }+\stackrel{~}{\beta }_\nu (m)\sqrt{\mathrm{\Omega }_\mu }\frac{2}{\sqrt{\mathrm{\Omega }_\mu }}\beta _\mu (m)\delta _{\mu \nu }\right],\end{array}$$
(182)
$`𝒩`$ being a normalization constant.
Using the Richardson equations to get rid of the sums, we recast (182) as follows
$$<(m_1,m_2)|\widehat{A}^{}|(m)>=\frac{1}{𝒩}\frac{2C_1(m_1,m_2)C_2(m_1,m_2)C_1(m)}{G[(E_1(m_1,m_2)E(m)][(E_2(m_1,m_2)E(m)]},$$
(183)
which is real since $`E_1`$ and $`E_2`$ are either real or complex conjugate.
We display in fig. 5 the amplitudes (183) for the transfer from a $`n=2`$ to a $`n=1`$ system (divided, for obvious convenience, by $`\sqrt{\mathrm{\Omega }}`$). We consider the $`n=1`$ system either in the ground collective state ($`m=1`$) or in the first excited trapped state ($`m=2`$). On the other hand the $`n=2`$ system is either in the $`v_l=0`$ (namely $`(m_1,m_2)=(1,2)`$) or in the $`v_l=2`$ ($`(m_1,m_2)=(2,2)`$) state. The calculation is performed for a harmonic oscillator well with $`L=4,10,20`$ and $`40`$.
The matrix elements clearly exhibit a conspicuous enhancement in the proximity of $`G_{\mathrm{cr}}`$ (marked by an asterisk in the figure), which most clearly reflects the onset of an ODLRO (off-diagonal-long-range-order) in the system.
The behavior of the transition matrix element around the critical points is quite delicate since both the numerator and the denominator in (182) tend to vanish when $`E_1`$ and $`E_2`$ tend to $`2e_1`$. In this limit the transition amplituse (182) reads
$$<(m_1,m_2)|\widehat{A}^{}|(m)>\underset{GG_{\mathrm{cr}}^\pm }{\overset{}{}}\frac{2}{G_{\mathrm{cr}}^\pm }\frac{_{\mu 1}{\displaystyle \frac{\sqrt{\mathrm{\Omega }_\mu }\stackrel{~}{\beta }_\mu (m)}{2ϵ_\mu 2ϵ_1}}\pm \stackrel{~}{\beta }_1(m)𝒫_{(2)}}{\sqrt{6𝒫_{(2)}^4\pm 8𝒫_{(2)}𝒫_{(3)}^3+2𝒫_{(4)}^4}}.$$
(184)
Now, the inverse moments of the level distribution (172) vanish in the $`L\mathrm{}`$ limit: hence (184) diverges. As above mentioned, this occurrence, not appearing in fig. 5 because of the chosen normalization, relates to the ODLRO which sets in into a system close to a phase transition.
Specifically the diagonal amplitudes for large $`L`$ behave according to
$`<(m_1,m_2)|\widehat{A}^{}|(m)>\stackrel{L\mathrm{}}{}\{\begin{array}{cc}\sqrt{2\mathrm{\Omega }}\theta \left(\stackrel{~}{G}\stackrel{~}{G}_{\mathrm{cr}}\right)v_l=0\hfill & \\ \sqrt{\mathrm{\Omega }}\theta \left(\stackrel{~}{G}\stackrel{~}{G}_{\mathrm{cr}}\right)v_l=2.\hfill & \end{array}`$ (185)
The off-diagonal amplitudes behave instead as
$$<(m_1,m_2)|\widehat{A}^{}|(m)>\stackrel{L\mathrm{}}{}\delta \left(\stackrel{~}{G}\stackrel{~}{G}_{\mathrm{cr}}\right).$$
(186)
From the figures it is also clear that the critical value of $`G`$ increases with $`L`$ and $`\stackrel{~}{G}_{\mathrm{cr}}L^28`$ as $`L\mathrm{}`$ (see Eq. (173)).
## 4 BCS and the exact solution
In concluding this work we shortly comment on the relation between the exact solution of the Richardson equations and the BCS solution in terms of Bogoliubov quasi-particles.
To this scope we self-consistently solve the well-known BCS equations for $`N`$ fermions living in $`L`$ levels, namely
$`v_\nu ^2={\displaystyle \frac{1}{2}}\left[1{\displaystyle \frac{ϵ_\nu \mu }{\sqrt{(ϵ_\nu \mu )^2+\mathrm{\Delta }^2}}}\right]=1u_\nu ^2,`$ (187)
$`{\displaystyle \underset{\nu =1}{\overset{L}{}}}{\displaystyle \frac{\mathrm{\Omega }_\nu }{\sqrt{(ϵ_\nu \mu )^2+\mathrm{\Delta }^2}}}={\displaystyle \frac{2}{G}},`$ (188)
$`{\displaystyle \underset{\nu =1}{\overset{L}{}}}\mathrm{\Omega }_\nu \left[1{\displaystyle \frac{ϵ_\nu \mu }{\sqrt{(ϵ_\nu \mu )^2+\mathrm{\Delta }^2}}}\right]=N,`$ (189)
where $`\mathrm{\Delta }=G_\nu \mathrm{\Omega }_\nu u_\nu v_\nu `$ the gap, $`\mu `$ the chemical potential and $`ϵ_\nu =e_\nu Gv_\nu ^2`$, and we evaluate in this framework the excitation energy of a system with seniority $`v`$, given by the energy of $`v`$ quasi-particles, namely
$$E_{BCS}=\underset{\nu =1}{\overset{v}{}}\sqrt{(ϵ_\nu \mu )^2+\mathrm{\Delta }^2}.$$
(190)
In fig. 6 we display and compare the exact excitation energies $`E_{\mathrm{exc}}=E(v_l)E(g.s.)`$ for a $`L=3`$ harmonic oscillator well and the corresponding Bogoliubov’s quasi-particles predictions (190) for a $`v=2`$ and a $`v=4`$ state. It appears that for $`\stackrel{~}{G}`$ larger than the highest critical value ($`\stackrel{~}{G}_{\mathrm{cr}}^{}0.6`$) both the BCS and the exact excitation energy become linear functions of $`\stackrel{~}{G}`$ and, remarkably, are very close to each other, in particular for $`v_l=2`$.
Note that the excitations in the BCS picture amount to break a pair, whereas in the Richardson frame the like-seniority excitations correspond, microscopically, to two particles-two holes (the $`v_l=2`$) and to four particles-four holes (the $`v_l=4`$) excitations without the breaking of any pair. In other words they are associated with the promotion of one or two pairs to higher lying single particle levels.
The result in Fig. 6 thus strengthens the correspondence between seniority and “like-seniority”. Moreover it confirms that the Richardson exact solutions behave as (190) in the strong $`G`$ limit, as was proved by Gaudin in the large $`L`$ limit . It is remarkable that this appears to be approximately true already for $`L=3`$.
## 5 Conclusions
In this report we have presented two approaches to the bosonization of fermionic systems based on the introduction of composite variables associated with the relevant degrees of freedom. One relates to the change of variables in the Berezin integral expressing the generating functional, the other to the straight solution of the Richardson equations.
Our testing ground has been the simple pairing Hamiltonian, which acts between fermions coupled to total angular momentum $`J=0`$. Such system exhibits three kinds of excitations, related to the addition of removal of pairs, breaking of pairs (the so-called seniority excitations) and promotion of $`J=0`$ pairs to higher energy levels, respectively.
The first kind of excitation connects the ground state energies of different systems, in particular atomic nuclei, and is associated with a Goldstone mode in the spectrum. Although the system we consider is finite, and therefore cannot truly display a spontaneous symmetry breaking mechanism (as it happens instead in superconductivity), nevertheless the precursor signature of a Goldstone boson can be clearly identified. This has been demonstrated, at least in the degenerate case, by performing a Hubbard-Stratonovitch linearization of the action and showing that the associated bosonic auxiliary field indeed endows the properties of a Goldstone boson.
The seniority excitations have been studied, again in the degenerate case, in the framework of the Hamiltonian formalism cast in the language of Grassmann variables and a minimal basis of composite fields has been set up which reproduces the full pairing spectrum in the degenerate case together with the associated wave functions.
When many single particle levels are present, new facets of the pairing physics disclose themselves.
First a third type of excitations (promotion of pairs to higher levels) can occur. These have been classified in terms of a new quantum number, the “like-seniority”, closely related to the Gaudin number , which has been shown to coincide with the true seniority when the strength of the interaction is large. Second, critical values of the coupling constant $`G`$ appear, signalling the transition between two different regimes: one dominated by the mean field, the other by the pairing bosonizing interaction. Notably the number of the critical values of $`G`$, which is crucially linked to the exsistence of a lowest single particle level with pair degeneracy $`\mathrm{\Omega }=1`$, increases with the number of pairs: however all the $`G_{\mathrm{cr}}`$ collapse into a common value when the number of levels tends to infinity.
The next step in the study of the pairing Hamiltonian will be to consider the general case of $`n`$ pairs living in a set of single particle levels and explore whether the signatures of a Goldstone boson and of the Higgs excitation are still present, so shading further light on the phase transition to superconductivity occurring in finite and infinite Fermi systems. In this connection the escape mechanism of the pairs from the grid of the single particle levels should be carefully examined in relation to the pair degeneracy of the latter.
Finally, worth to be explored are new scheme recently proposed for the bosonization of finite Fermi systems, specifically of atomic nuclei: in particular the one outlined in Refs. which, always in the framework of the path integral, introduces the concept of boson dominance in computing the trace expressing the partition function of a nucleus.
Acknowledgements
Most of the work presented here is the fruit of a longstanding collaboration with R. Cenni, A. Molinari and F. Palumbo, whom we would like to thank for many useful discussions and careful reading of the manuscript. |
warning/0506/hep-ex0506021.html | ar5iv | text | # SPECTROSCOPIC MEASUREMENTS USING THE H1 AND ZEUS DETECTORS
## 1 Introduction
At the $`ep`$-collider HERA, interactions are studied at a centre-of-mass energy of 300–320 GeV by the H1 and ZEUS multi-purpose detectors. The virtuality, $`Q^2`$, of the exchanged photon allows to distinguish two kinematical regimes: photoproduction ($`Q^2<1`$GeV<sup>2</sup>) and deep-inelastic-scattering (DIS) ($`Q^2>1`$GeV<sup>2</sup>). As recent results from fixed-target experiments have evidence for a narrow baryon resonance decaying to $`K^+n`$ $`^\mathrm{?}`$ and $`K_s^0p`$ $`^\mathrm{?}`$, interpreted as a pentaquark, spectroscopic measurements covering this and related topics have been performed at HERA.
## 2 Evidence for a narrow baryonic state decaying to $`K_s^0p`$ and $`K_s^0\overline{p}`$ in DIS
A resonance search$`^\mathrm{?}`$ has been made in the $`K_s^0p`$ and $`K_s^0\overline{p}`$ invariant-mass spectrum measured with the ZEUS detector using an integrated luminosity of 121 pb<sup>-1</sup>. The search was performed in the central rapidity region of inclusive DIS for $`Q^2`$ above 1 GeV<sup>2</sup>. The results support the existence of a state like those observed by the fixed-target experiments, with a mass of $`1521.5\pm 1.5`$ (stat.) $`{}_{1.7}{}^{}{}_{}{}^{+2.8}`$ (syst.) MeV and a Gaussian width consistent with the experimental resolution of $`2`$MeV. The signal is visible at moderate $`Q^2`$ and, for $`Q^2>20`$GeV<sup>2</sup>, contains $`221\pm 48`$ events (Fig. 1 (a)). The probability of a similar signal anywhere in the range 1500–1560 MeV arising from fluctuations of the background is below $`6\times 10^5`$.
## 3 Search for pentaquarks decaying to $`\mathrm{\Xi }\pi `$ in DIS
The ZEUS collaboration has also performed a search$`^\mathrm{?}`$ for pentaquarks decaying to $`\mathrm{\Xi }^{}\pi ^{}`$ ($`\mathrm{\Xi }^{}\pi ^+`$) and corresponding antiparticles using the same data sample as in the previous analysis, to investigate the observation of the NA49 experiment$`^\mathrm{?}`$. A clear signal for $`\mathrm{\Xi }^0(1530)\mathrm{\Xi }^{}\pi ^+`$ was observed. However, no signal for any new baryonic state was observed at higher masses in either the $`\mathrm{\Xi }^{}\pi ^{}`$ or $`\mathrm{\Xi }^{}\pi ^+`$ channels (Fig. 1 (b) and (c)). The searches in the antiparticle channels were also negative. Upper limits on the ratio of a possible $`\mathrm{\Xi }_{3/2}^{}`$ ($`\mathrm{\Xi }_{3/2}^0`$) signal to the $`\mathrm{\Xi }^0(1530)`$ signal were set in the mass range 1650–2350 MeV.
## 4 Evidence for a narrow anti-charmed baryon state by H1
The H1 collaboration observed$`^\mathrm{?}`$ a narrow resonance in $`D^{}p`$ and $`D^+\overline{p}`$ invariant-mass combinations in a data sample corresponding to an integrated luminosity of $`75`$pb<sup>-1</sup> in DIS. The decay channel $`D^+D^0\pi _s^+(K^{}\pi ^+)\pi _s^+`$ (and the corresponding antiparticle decay) was used to identify the $`D^+`$ mesons, the (anti-)proton candidates were selected using particle identification based on the differential energy loss. The resonance has a mass of $`M(D^{}p)=3099\pm 3`$ (stat.) $`\pm 5`$ (syst.) MeV and a measured Gaussian width of $`12\pm 3`$ (stat.) MeV, compatible with the experimental resolution (Fig. 2 (a)). The probability for a background fluctuation to produce a signal as large as observed is less than $`4\times 10^8`$. The region of $`M(D^{}p)`$ in which the signal is observed contains a richer yield of $`D^{}`$ mesons and exhibits a harder proton candidate momentum distribution than is the case for the side bands in $`M(D^{}p)`$. The fraction of $`D^{}`$ mesons originating from the resonance decay is found to be $`1.46\pm 0.32\%`$. A signal with compatible mass and width is also observed in an independent photoproduction sample.
The resonance is interpreted as an anti-charmed baryon with a minimal constituent quark composition of $`uudd\overline{c}`$, together with its charge conjugate.
## 5 Search for a narrow anti-charmed baryonic state decaying to $`D^\pm p^{}`$ by ZEUS
A similar search$`^\mathrm{?}`$ has been performed by the ZEUS collaboration in the $`D^\pm p^{}`$ invariant-mass spectrum using an integrated luminosity of $`126`$pb<sup>-1</sup>. The decay channels $`D^+D^0\pi _s^+(K^{}\pi ^+)\pi _s^+`$ and $`D^+D^0\pi _s^+(K^{}\pi ^+\pi ^+\pi ^{})\pi _s^+`$ (and the corresponding antiparticle decays) were used to identify $`D^+`$ mesons in a combined DIS and photoproduction sample. No resonance structure was observed in the mass spectrum in DIS as well as in photoproduction. The upper limit on the fraction of $`D^{}`$ mesons originating from such a resonance decay is found to be 0.23% (95% C.L.); the upper limit for DIS with $`Q^2>1`$GeV<sup>2</sup> is 0.35% (95% C.L.). To guarantee the highest possible comparability between the measurements of the H1 and the ZEUS collaborations the analysis has been repeated with requirements as close as possible to those of the H1 collaboration (Fig. 2 (b) and (c)); no signal was observed. The results of the two collaborations are incompatible.
## 6 Observation of $`K_s^0K_s^0`$ resonances in DIS
Inclusive $`K_s^0K_s^0`$ production has been studied$`^\mathrm{?}`$ by the ZEUS collaboration in DIS using an integrated luminosity of $`120`$ pb<sup>-1</sup>. Two states are observed at masses of $`1537_8^{+9}`$ MeV and $`1726\pm 7`$ MeV, as well as an enhancement around $`1300`$ MeV (Fig. 3). The state at $`1537`$ MeV is consistent with the well established $`f^{}`$. The state at $`1726`$ MeV may be the glueball candidate $`f_0(1710)`$. However, its width of $`38_{14}^{+20}`$ MeV is narrower than the $`125\pm 10`$ MeV observed by previous experiments for the $`f_0(1710)`$.
## 7 Conclusions
Most of the recent spectroscopic measurements using the H1 and ZEUS detectors cover exotic resonances. The ZEUS collaboration has observed a resonance in $`K_s^0p(\overline{p})`$ which could be the $`\mathrm{\Theta }^+`$ which has been observed recently by several other experiments. The H1 collaboration has observed a narrow resonance in $`D^{}p`$ which could be a candidate for a charmed pentaquark. However, this observation is not been confirmed by the ZEUS collaboration.
## References |
warning/0506/quant-ph0506187.html | ar5iv | text | # Quantum feedback cooling of a single trapped ion in front of a mirror
## I Introduction
Laser cooling and trapping of single ions BlattRMP ; WinelandCooling ; EITCooling is one of the highlights in the development of quantum optics during the last two decades. Single trapped ions are a laboratory paradigm of a quantum system, which can be prepared and controlled on the single quantum level, and whose time evolution can be monitored continuously by observing the scattered light in photodetection or homodyne measurements BlattRMP . By continuous observation of a single quantum system CarmichaelBook ; GZ we learn the state of the system, as described by a conditional system density matrix $`\rho _c(t)`$, and based on this knowledge we can act back on the system, giving rise to quantum feedback control of the system of interest WM\_coll ; Feedback\_Squeezing ; FB\_CQED ; FB\_Ion ; FB\_Meso . In the present paper we present a theory of *quantum feedback cooling of a single trapped ion*: by extracting from the scattered light the position of the ion in the trap, we implement a feedback loop on the system in the form of a damping force with the purpose of cooling the ion motion in the trap. Development of this theory is not only of fundamental interest in quantum optics, but the particular setup studied is motivated by ongoing experimental efforts BushevThesis ; FeedbackCoolingExp to implement quantum feedback cooling of single trapped ions in laboratory. Indeed the present theoretical results provides a quantitative basis for the understanding of these experiments FeedbackCoolingExp .
The particular setup studied in the present paper is a single laser cooled trapped ion in front of a mirrorMirrorColl , as illustrated in Fig. 1, and motivated by present experiments BushevThesis ; FeedbackCoolingExp . A single ion is stored a distance $`L`$ from a mirror in a harmonic trapping potential. The ion is assumed to be a two-level system weakly driven near resonance by laser light. Light is scattered into both the mirror mode, as well as the other other “background” modes of the radiation field. By detuning the laser on the red side of the atomic transition, the ion is laser cooled to a temperature corresponding to Doppler limit, where the mean occupation of the trap levels is much larger than one (i.e. far away from the sideband cooling limit to the ground state of the trap). Motion of the ion adds sidebands of the light scattered into the mirror mode displaced by the trap frequency. Observing the scattered light of these motional sidebands allows us to infer the position of the ion in the sense of continuous measurement theory, and feed back a damping force proportional to the momentum to implement quantum feedback cooling. In this paper we will first formulate a continuous measurement theory to read the position of the trapped ion from the scattered light using the language of stochastic Schrödinger Equations GZ ; CarmichaelBook . Building on general quantum feedback theory formulated by Wiseman and Milburn WM\_coll ; Feedback\_Squeezing , we will then derive a quantum feedback master equation for the motion of the trapped ion. This will allow us to study the dynamics and limits of quantum feedback cooling.
For the setup studied in this paper the continuous readout of the ion position is based on light scattering into the mirror mode, with additional photons scattered into all other “background modes” of the radiation field. Spontaneous emission is intrinsically associated with a momentum recoil of the ion, which perturbs the ion motion, i.e. contributes a heating mechanism for the ion. In a parallel paper Rabl we study a quantum feedback scheme based on a *dispersive readout of the velocity* of the trapped ion to avoid this unwanted heating. It is based on the large variation of the index of refraction with the Doppler effect near a *dark state resonance in an atomic $`\mathrm{\Lambda }`$-system* (based on electromagnetically induced transparency).
The paper is structured as follows. Sec. II presents the basic dynamic equations for the motion of an ion in front of a mirror. Quantum feedback equations are formulated in Sec. III, while results are presented in Sec. IV.
## II Model and Basic Equations
In this section we will develop the basic equations for continuous measurement of the photons in the mirror mode of the electromagnetic field. We will start with a detailed description of our model in terms of a Schrödinger equation for the coupled atom-bath system and the exciting laser. Continuous measurement theory provides us with a quantum stochastic Schrödinger equation and hence a quantum stochastic master equation in the Lamb-Dicke limit, where we adiabatically eliminate the excited state from the two-level atom. We will then derive the photocurrent obtained by detecting mirror mode photons and the corresponding stochastic master equation for the conditional density operator in the white noise (diffusive) limit.
### II.1 Single trapped ion in front of a mirror
We consider a single trapped ion which is placed at a distance $`L`$ from a mirror as indicated in Fig. 1 BlattRMP ; CiracLaserCooling ; Dorner02 ; EschnerNature . For the harmonic motion we assume a 1D model in the $`z`$-direction (identical to the mirror axis). The harmonic trap has an oscillation frequency $`\nu _T`$, and we denote the destruction (creation) trap operator by $`a`$ ($`a^{}`$). The electronic degrees of the ion form a two-level atom with atomic transition frequency $`\omega _{eg}`$, with ground state $`|g`$ and excited state $`|e`$. We drive the two-level system with a laser with frequency $`\omega _L`$ which couples the ground to the excited state with the Rabi frequency $`\mathrm{\Omega }`$ and a detuning from the atomic resonance $`\mathrm{\Delta }_L=\omega _L\omega _{eg}`$. The atomic system Hamiltonian can thus be written as
$$H_{\mathrm{sys}}=\nu _Ta^{}a\mathrm{\Delta }_L|ee|\frac{1}{2}\mathrm{\Omega }\left(\mathrm{e}^{\mathrm{i}k_{\mathrm{eff}}\widehat{z}}\right|eg|+\mathrm{h}.\mathrm{c}.)$$
(1)
Note that in this paper we set $`\mathrm{}=1`$. In the interaction term we allow for a laser field incident at an angle $`\chi `$ with respect to an axis normal to the $`z`$-axis. The momentum recoil due to absorption of a laser photon is represented by $`k_{\mathrm{eff}}\widehat{z}=\eta \mathrm{sin}\chi \left(a+a^{}\right)\stackrel{~}{\eta }\left(a+a^{}\right)`$ where the Lamb-Dicke parameter $`\eta =2\pi a_0/\lambda `$ is the ratio of the size of the ground state and the laser wavelength. Due to the geometry of the system in consideration, the (quantized) electric field consists of two contributions, $`\stackrel{}{E}^{(+)}=\stackrel{}{E}_m^{(+)}+\stackrel{}{E}_b^{(+)}`$, where the $`\stackrel{}{E}_m^{(+)}`$ denotes the modes restricted by the boundary condition of the mirror and $`\stackrel{}{E}_b^{(+)}`$ the remaining background modes Dorner02 ; EschnerNature . We adopt a 1D model for the mirror mode and write for the electric field operator
$$\stackrel{}{E}_m^{(+)}\left(z\right)=\mathrm{i}_0^{\mathrm{}}𝑑\omega \alpha _\omega \stackrel{}{e}\mathrm{sin}\left(k(\omega )z\right)b_m(\omega )$$
(2)
with $`\alpha _\omega `$ a normalization factor for the mode function. The internal states of the atom couple to the vacuum field by an electric dipole transition. Denoting by $`\stackrel{}{d}`$ the dipole matrix element, and introducing Pauli operator notation for the two level system, $`\sigma _{}=|ge|`$, the system-bath coupling Hamiltonian is
$$H_{\mathrm{int}}=\stackrel{}{d}\left(\stackrel{}{E}_b^{(+)}(\widehat{z})+\stackrel{}{E}_m^{(+)}\left(\widehat{z}\right)\right)\sigma _{}+\mathrm{h}.\mathrm{c}.$$
(3)
The total Hamiltonian for the ion coupled to the radiation field is
$$H=H_{\mathrm{sys}}+H_{\mathrm{bath}}+H_{\mathrm{int}}.$$
(4)
Here $`H_{\mathrm{bath}}`$ is the free Hamiltonian for the radiation field. We write this Hamiltonian as the sum of a Hamiltonian for the mirror and the background modes $`H_{\mathrm{bath}}=H_m+H_b`$. In our 1D model the mirror mode Hamiltonian has the form $`H_m=𝑑\omega \omega b_m^{}(\omega )b_m(\omega )`$ with $`b_m(\omega )`$ photon destruction operators, satisfying bosonic commutation relations $`[b_m(\omega ),b_m^{}(\omega ^{})]=\delta (\omega \omega ^{})`$. Similar expression can be given for the background modes.
In analyzing this problem we are interested in the situation where the time delay $`\tau _M=2L/c`$ of the emitted light bouncing from mirror back to the atom is much shorter than the system time scales, in particular the spontaneous emission time from the excited state, $`\tau _M1/\mathrm{\Gamma }`$, and the timescales associated with the laser interactions $`\tau 1/\mathrm{\Omega }`$, $`1/|\mathrm{\Delta }_L|`$. This justifies the Markov approximation for the emission into the mirror modes, where we refer to Dorner02 for a complete analysis.
In the following we will denote the total spontaneous emission rate of the atom by $`\mathrm{\Gamma }=\mathrm{\Gamma }_m+\mathrm{\Gamma }_b.`$ Here $`\mathrm{\Gamma }_m=\epsilon \mathrm{\Gamma }`$ with $`\epsilon `$ the fraction of the solid angle covered by the lens is the emission rate into mirror mode, and $`\mathrm{\Gamma }_b=(1\epsilon )\mathrm{\Gamma }`$ the emission rate into the background modes.
### II.2 Quantum Stochastic Schrödinger Equation
The dynamics of our model is summarized in the Schrödinger Equation
$`|\dot{\mathrm{\Psi }}(t)=`$ $`[\mathrm{i}H_{\mathrm{sys}}+`$ (5)
$`+`$ $`\sqrt{\mathrm{\Gamma }_m}\sigma _{}\mathrm{sin}\left(k_{eg}(L+\widehat{z})\right)b_m^{}(t)+\mathrm{h}.\mathrm{c}.+`$
$`+`$ $`\sqrt{\mathrm{\Gamma }_b}{\displaystyle _1^{+1}}du\sqrt{N(u)}\sigma _{}\mathrm{e}^{\mathrm{i}uk_{eg}\widehat{z}}b_u^{}(t)+\mathrm{h}.\mathrm{c}.]|\mathrm{\Psi }(t)`$
We choose to formulate the problem in the language of a *Quantum Stochastic Schrödinger equation* (QSSE) GZ , which allows for a direct connection with continuous measurement of the scattered light, and provides a direct link to quantum feedback theory developed in the following subsections.
In Eq. (5) $`|\mathrm{\Psi }(t)`$ is the Schrödinger state vector of the combined atom-field system, i.e. the laser-driven trapped ion including the mirror and background modes of the radiation field. The first term on the RHS is the time evolution due to the system Hamiltonian (1).
The second and third line describe the interaction of the two-level atom with the mirror mode and the background modes, respectively. We assume that these radiation modes are initially in the vacuum state. In writing Eq. (5) we have transformed to an interaction picture with respect to the free Hamiltonian of the radiation field $`H_B`$ GZ . As a result, we have introduced bath operators for the mirror mode $`b_m(t)=1/\sqrt{2\pi }𝑑\omega b_m(\omega )e^{i\omega t}`$. In the Markovian limit these operators satisfy bosonic commutation relations
$$[b_m(t),b_m^{}(s)]=\delta (ts).$$
(6)
In a Quantum Langevin formulation GZ $`b_m(t)`$ represents a quantum noise operator. Thus the second line of (5) describes the emission of photons by the atom into the mirror mode, with the center of the ion trap displaced a distance $`L`$ from the mirror. We note that the motion of the ion couples to the light via the recoil, as seen by the appearance of $`\widehat{z}`$ in the mirror mode function. This coupling imparts information of the ion motion on the light emitted in the mirror mode. In the following subsections we will analyze this scattered light to continuously monitor the atomic motion, with the goal of implementing a feedback loop to cool the ion. The coupling strength to the mirror mode is proportional to the square root of the spontaneous emission probability into the mirror mode $`\mathrm{\Gamma }_m\epsilon \mathrm{\Gamma }`$ with $`\epsilon `$ the fraction of the solid angle (typically $`\epsilon `$ is much smaller than one).
The third line in Eq. (5) represents spontaneous emission of the ion into the background modes. This is a coupling term familiar from the theory of laser cooling of two-level atoms CiracLaserCooling . Spontaneous emission into the background mode is again associated with a recoil of the ion motion. In our 1D model for the motion of the trapped ion, it is the projection of this momentum on the $`z`$-axis which is the relevant momentum transfer. Denoting by $`\theta `$ the angle between the emitted photon and the $`z`$-axis, and $`u=\mathrm{cos}\theta `$, we associate the transition for the excited state to the ground state including the momentum transfer with the operator $`e^{iuk_{eg}\widehat{z}}\sigma _{}`$, where $`k_{eg}\omega _{eg}/ck_L`$. Spontaneous photons can be emitted in all directions into the background modes consistent with the dipole radiation pattern of the given electronic transition. We denote this (normalized) angular dependence by $`N(u)`$. Thus the integral over $`u`$ in the last line of Eq. (5) realizes photon emission into all of these possible directions. The operators $`b_u(t)`$ are again photon destruction (or noise) operators associated with these emission directions. They satisfy commutation relations
$$[b_u(t),b_u^{}^{}(s)]=\delta (uu^{})\delta (ts),$$
(7)
and commute with the mirror bath operators $`b_m(t)`$ introduced above. The coupling strength to the background modes is proportional to $`\sqrt{\mathrm{\Gamma }_b}\sqrt{(1\epsilon )\mathrm{\Gamma }}`$. For red laser detuning $`\mathrm{\Delta }_L<0`$ the cycle of laser excitation followed by spontaneous emission into the background mode leads to laser cooling.
### II.3 Ito form of the Quantum Stochastic Schrödinger Equation
To give a meaning to the white noise limit (compare Eqs. (6,7)), we must interpret the Schrödinger equation (5) as a quantum stochastic *Stratonovich* equation GZ . As usual, it is more convenient to work with an Ito form, where Wiener noise increments
$$dB_{m,u}(t)=_t^{t+dt}𝑑sb_{m,u}(s)$$
(8)
“point to the future”, i.e. are statistically independent of $`|\mathrm{\Psi }(t)`$. These Wiener noise increments satisfy the Ito table
$`dB_m(t)dB_m^{}(t)`$ $`=dt,`$ (9a)
$`dB_u(t)dB_u^{}^{}(t)`$ $`=\delta (uu^{})dt,`$ (9b)
which follow from Eqs. (6,7), the other entries of the Ito table being zero. The resulting Ito QSSE is
$`(\mathrm{I})d|\mathrm{\Psi }(t)=`$ $`[\mathrm{i}H_{\mathrm{eff}}dt+\sqrt{\mathrm{\Gamma }_m}C_m(\widehat{z})dB_m^{}\left(t\right)`$ (10)
$`+\sqrt{\mathrm{\Gamma }_b}{\displaystyle _1^{+1}}du\sqrt{N(u)}C_u(\widehat{z})dB_u^{}(t)]|\mathrm{\Psi }(t)`$
Here, we have introduced the “jump operators”
$`C_u(\widehat{z})`$ $`=\mathrm{e}^{\mathrm{i}uk_{eg}\widehat{z}}\sigma _{},`$ (11a)
$`C_m(\widehat{z})`$ $`=\mathrm{sin}\left(k_{eg}\left(L+\widehat{z}\right)\right)\sigma _{},`$ (11b)
which are associated with the emission of a photon in the background modes and the mirror modes, respectively. Furthermore, we have defined an *effective non-hermitian system Hamiltonian*
$$H_{\mathrm{eff}}=H_{\mathrm{sys}}\frac{\mathrm{i}}{2}\left[\mathrm{\Gamma }_b+\mathrm{\Gamma }_m\mathrm{sin}^2(k_{eg}(L+\widehat{z}))\right]|ee|.$$
(12)
The non-hermitian part of $`H_{\mathrm{eff}}`$ arises from the Ito correction in the conversion process. Physically, it corresponds to the radiation damping of the excited state due to the total radiation field. We also note that the photon absorption terms have disappeared in Eq. (10) due to $`dB_{m,u}(t)|\mathrm{\Psi }(t)=0`$. This follows from our assumption of an initial vacuum state.
### II.4 Quantum Stochastic Master Equation
We are interested in the time evolution of our system where the photons emitted in the mirror mode are detected by a photon counter, while the background modes remain unobserved. Therefore, we are only interested in the dynamics of the reduced density operator $`\widehat{W}\left(t\right)\mathrm{Tr}_b\left\{|\mathrm{\Psi }\left(t\right)\mathrm{\Psi }\left(t\right)|\right\}`$ where we trace over the background modes of the radiation field. We emphasize that $`\widehat{W}\left(t\right)`$ still contains all the degrees of freedom of the mirror modes, in addition to the internal and external atomic dynamics.
Using Ito calculus (see Appendix A) we obtain the quantum stochastic master equation (QSME)
$`(\mathrm{I})d\widehat{W}\left(t\right)=`$ $`\mathrm{i}\left(H_{\mathrm{eff}}\widehat{W}\left(t\right)\widehat{W}\left(t\right)H_{\mathrm{eff}}^{}\right)dt`$ (13)
$`+`$ $`\mathrm{\Gamma }_m𝒥\left[C_m(\widehat{z})\right]dB_m^{}\left(t\right)\widehat{W}\left(t\right)dB_m\left(t\right)`$
$`+`$ $`\sqrt{\mathrm{\Gamma }_m}\left(C_m(\widehat{z})dB_m^{}\widehat{W}(t)+\widehat{W}(t)C_m^{}(\widehat{z})dB_m\right)`$
$`+`$ $`\mathrm{\Gamma }_b{\displaystyle _1^{+1}}𝑑uN(u)𝒥\left[C_u(\widehat{z})\right]\widehat{W}\left(t\right)𝑑t`$
with $`H_{\mathrm{eff}}`$ defined in Eq. (12). For the “recycling terms” we use the notation
$$𝒥\left[c\right]\rho c\rho c^{}.$$
(14)
Before proceeding we note that for $`\epsilon =0`$, i.e. no coupling to the mirror modes, Eq. (13) reduces to the standard master equation for 1D laser cooling of a two-level atom GZ . In this case $`\widehat{W}`$ is only an atomic density operator containing the internal and motional dynamics. For $`\epsilon 0`$, we still have a stochastic equation with the mirror bath degrees of freedom included.
### II.5 Adiabatic elimination of the excited state and Lamb-Dicke limit
We will simplify the above QSSE (10) and QSME (13) with two assumptions. First, we assume weak laser excitation to the excited state, $`\mathrm{\Omega }\mathrm{max}(\mathrm{\Gamma },\left|\mathrm{\Delta }\right|)`$. Second, we assume a small Lamb-Dicke parameter $`\eta 2\pi a_0/\lambda 1`$ (tight trap): this allows us to expand the exponents $`\mathrm{e}^{\mathrm{i}k\widehat{z}}\mathrm{e}^{\mathrm{i}\eta (a+a^{})}=1+\mathrm{i}\eta (a+a^{})+𝒪(\eta ^2)`$. Both of these assumptions are well satisfied in present experiments BlattRMP
To eliminate the weakly populated excited level, we go back to Eq. (10) and expand the state vector $`|\mathrm{\Psi }(t)`$ into ground state and excited state components,
$$|\mathrm{\Psi }\left(t\right)|\psi _g\left(t\right)|g+|\psi _e\left(t\right)|e.$$
(15)
As shown in Appendix B we can eliminate $`|\psi _e(t)`$ in perturbation theory in the Ito QSSE (10) to obtain an effective equation for $`|\psi _g\left(t\right)`$. In a similar way as for Eq. (13) we obtain a QSME for the partially reduced density operator
$$\widehat{w}\left(t\right)\mathrm{Tr}_b\left\{|\psi _g\left(t\right)\psi _g\left(t\right)|\right\},$$
(16)
given by
$`(\mathrm{I})d\widehat{w}(t)=`$ $`\mathrm{i}\left[h_{\mathrm{eff}}\widehat{w}(t)\widehat{w}(t)h_{\mathrm{eff}}^{}\right]dt`$ (17)
$`+\gamma 𝒥\left[c_m(\widehat{z})\right]dB_m^{}(t)\widehat{w}(t)dB_m(t)`$
$`+\sqrt{\gamma }\left(c_m(\widehat{z})dB_m^{}(t)\widehat{w}(t)+\widehat{w}(t)c_m^{}(\widehat{z})dB_m(t)\right)`$
$`+_b\widehat{w}(t)dt.`$
The first three lines give the dynamics of the ion motion coupled to the mirror mode. The fourth line describes the traced-out action of the background mode on the ion motion, i.e. laser cooling of the ion.
In Eq. (17) we have defined an effective Hamiltonian acting only on the motional states of the ion,
$$h_{\mathrm{eff}}=H_T\frac{\mathrm{i}}{2}\gamma c_m^{}(\widehat{z})c_m(\widehat{z}).$$
(18)
where we expand the eliminated jump operators to second order in the Lamb-Dicke limit with the center of the trap at $`k_{eg}L=\pi /4`$:
$$c_m(\widehat{z})\frac{1}{\sqrt{2}}\left(1+\eta \left(a+a^{}\right)\frac{1}{2}\eta ^2\left(a+a^{}\right)^2\right).$$
(19)
The parameter
$$\gamma =\epsilon \mathrm{\Gamma }\frac{\mathrm{\Omega }^2}{4}\frac{1}{\mathrm{\Delta }_L^2+\frac{\mathrm{\Gamma }}{2}}$$
(20)
is the optical pumping rate into the mirror mode. The first three lines of Eq. (17) thus describe the motional state coupled via laser excitation followed by spontaneous emission to the mirror mode.
The Liouvillian $`_b`$ in the fourth line of Eq. (17) is the standard laser cooling Liouvillian for weak field excitation and in the Lamb-Dicke limit BlattRMP ; CiracLaserCooling ; WinelandCooling ,
$`_b\widehat{w}\left(t\right)`$ $`=A_{}𝒟\left[a\right]\widehat{w}\left(t\right)+A_+𝒟\left[a^{}\right]\widehat{w}\left(t\right)`$ (21)
$`\mathrm{\Gamma }_{\text{eff}}(N+1)𝒟\left[a\right]\widehat{w}\left(t\right)+\mathrm{\Gamma }_{\text{eff}}N𝒟\left[a^{}\right]\widehat{w}\left(t\right),`$
where we have used the notation
$$𝒟\left[c\right]\rho c\rho c^{}\frac{1}{2}\left(c^{}c\rho +\rho c^{}c\right).$$
(22)
The rates
$$A_\pm =\eta ^2\frac{\mathrm{\Omega }^2}{4}\mathrm{\Gamma }_b\left(\frac{\mathrm{sin}^2\chi }{\left(\mathrm{\Delta }_L\nu _T\right)^2+\frac{\mathrm{\Gamma }^2}{4}}+\frac{\alpha }{\mathrm{\Delta }_L^2+\frac{\mathrm{\Gamma }^2}{4}}\right).$$
(23)
have the meaning of cooling (heating) terms for red laser detuning $`\mathrm{\Delta }_L<0.`$ With $`\mathrm{\Gamma }_{\mathrm{eff}}=A_{}A_+>0`$ and for $`\mathrm{\Delta }_L<0`$ we have
$$N=\frac{A_+}{A_{}A_+},$$
(24)
which is the final mean trap occupation established by laser cooling (alone). We have also used the abbreviation $`\alpha =𝑑uu^2N(u)`$ for the dipole transition parameter and $`\chi `$ is the incident angle of the laser beam. With these definitions the mirror mode optical pumping rate (20) can be written as $`\gamma =\epsilon N\mathrm{\Gamma }_{\mathrm{eff}}/(1+\alpha )\eta ^2`$, and from $`\mathrm{\Gamma }_{\mathrm{eff}}\mathrm{sin}^2\chi `$ and $`N1/\mathrm{sin}^2\chi `$ we see that this pumping rate is independent from the angle of the incoming laser beam.
In the following we will study a scenario FeedbackCoolingExp where the laser cooling establishes a steady state with a mean trap occupation $`N1`$ (i.e. far from the ground state), as represented by the second line in Eq. (17). This is the limit of Doppler cooling, which is obtained if $`\mathrm{\Gamma }\nu _T`$. The minimally obtainable steady state energy in this limit is $`\mathrm{}\mathrm{\Gamma }(\alpha +1)/2`$. By observing the spontaneous emission into the mirror mode (see first two lines of Eq. (17)), we will infer the position of the atom to apply a feedback loop to cool the system (far) below the laser cooling limit.
For completeness we note that in the case where the mirror mode is not observed, the reduced system density operator $`\rho (t)\mathrm{Tr}_m\left\{\widehat{w}(t)\right\}`$ obeys the master equation
$`\dot{\rho }(t)`$ $`=\mathrm{i}[H_T,\rho (t)]+\gamma 𝒟\left[c_m(\widehat{z})\right]\rho (t)+_b\rho (t)`$ (25)
$`\mathrm{i}[H_T,\rho (t)]+_{\mathrm{LC}}\rho (t)_0\rho (t),`$
which contains the dynamics from the free ion motion, and the dissipative dynamics from the emission into the mirror mode and laser cooling. In a second order expansion in terms of $`\eta `$, we have
$$𝒟\left[c_m\left(\widehat{z}\right)\right]=\eta ^2\mathrm{cos}^2(k_{eg}L)𝒟\left[a+a^{}\right]+𝒪\left(\eta ^3\right)$$
(26)
which, multiplied by $`\gamma `$, is typically much smaller than $`\mathrm{\Gamma }_{\mathrm{eff}}N`$ and thus the corrections in the heating and cooling rates will be neglected here.
### II.6 Continuous observation of the mirror mode
We measure the photons emitted into the mirror modes by a photon counter as shown in Fig. 1. We denote by $`N_c(t)`$ the number of photon counts at time $`t`$. A particular count trajectory is characterized the photon detection times $`t_1,t_2,\mathrm{}`$. Our knowledge of the state of the system, given by the internal and external degrees of the ion, for a given count trajectory is represented by a conditional density matrix $`\rho _c\left(t\right)`$ GZ .
Given the state of the system at time $`t`$, $`\rho _c\left(t\right)`$, the detection of a mirror mode photon in a time interval $`(t,t+dt]`$ is associated with a quantum jump of the atom described by
$$\rho _{c,\mathrm{jump}}\left(t+dt\right)=\frac{𝒥\left[c_m\right]\rho _c\left(t\right)}{\mathrm{Tr}\left\{𝒥\left[c_m\right]\rho _c\left(t\right)\right\}}$$
(27)
where according to (11b) the atom returns to from the excited state to the ground state, and momentum is transferred to the ion motion in accordance with the mirror mode function. In the case of no observed photon, the system evolves with the effective non trace-preserving Liouvillian $`L_0`$
$$\rho _{c,\mathrm{no}\mathrm{jump}}\left(t+dt\right)=\left(1+L_0dt\right)\rho _c\left(t\right)$$
(28)
where
$$L_0\rho \mathrm{i}\left[h_{\mathrm{eff}}\rho \rho h_{\mathrm{eff}}^{}\right]+_b\rho $$
and $`h_{\mathrm{eff}}`$ is defined in Eq. (18). The expected number of counts in the interval $`(t,t+dt]`$ is with $`dN_c(t)=N_c\left(t+dt\right)N_c(t)`$
$$dN_c(t)=p_{\mathrm{emission}}^{(t,t+dt]}=\gamma \mathrm{Tr}_{\mathrm{sys}}\left\{𝒥\left[c_m\right]\rho _c\left(t\right)\right\}dt$$
(29)
In view of $`dN_c(t)=0`$ or $`1`$, for this point process we have the Ito table $`dN_c^2\left(t\right)=dN_c\left(t\right)`$ and $`dN_c\left(t\right)dt=0`$.
We can summarize the above *a posteriori* time evolution in an Ito stochastic Schrödinger equation (see. eg. GZ )
$`(\mathrm{I})\rho _c\left(t\right)=`$ $`_0\rho _c\left(t\right)dt+`$ (30)
$`+({\displaystyle \frac{𝒥\left[c_m\right]\rho _c\left(t\right)}{\mathrm{Tr}_{\mathrm{sys}}\left\{𝒥\left[c_m\right]\rho _c\left(t\right)\right\}}}\rho _c\left(t\right))\times `$
$`\times \left(dN_c\left(t\right)\gamma \mathrm{Tr}_{\mathrm{sys}}\left\{𝒥\left[c_m\right]\rho _c\left(t\right)\right\}dt\right)`$
where $`_0`$ is defined in (25). This equation gives the time evolution of the conditional density matrix of the ion $`\rho _c\left(t\right)`$ for a particular count trajectory. Not observing, i.e. tracing over the mirror mode, is equivalent to taking the ensemble average over all count trajectories in (30). In this case, we recover the master equation $`\dot{\rho }\left(t\right)=_0\rho \left(t\right)`$ for the *a priori* dynamics GZ .
### II.7 Diffusion approximation
In the previous subsection we considered photon counting of the light emitted in the mirror modes, and the associated time evolution of the system described by the condition density operator $`\rho _c\left(t\right)`$. We are interested in learning the motion (position) of the atom from the scattered light in the sense of continuous measurement. The goal is to use this information to control the motion of the atom, and eventually act back on the atom to cool it.
The scattered light of a weakly driven trapped atom ResFluorescence consists of (i) a strong elastic component at the frequency of the driving laser (see vertical transitions in Fig. 2), and (ii) weak motional sidebands at the trap frequency $`\nu _T`$ suppressed by the Lamb-Dicke parameter $`\eta `$. The information on the motion of the atom is encoded in the “motional sidebands”. We find it convenient to formulate the problem in a way, where we focus directly on the contributions of these sidebands to the photon count signal. The physical picture is that the elastic component acts like a “(strong) local oscillator” which beats with the “(weak) light emitted from the sidebands”. This situation is reminiscent of homodyne measurements in quantum optics GZ ; CarmichaelBook , and will lead in the following to a description in terms of a *diffusive stochastic process* rather than a point process associated with the photon counting described above. The formal expansion parameter is $`\eta 1`$ (Lamb-Dicke limit).
From the previous subsection we know that the mean number of photon counts in $`(t,t+dt]`$ is
$$dN_c(t)\stackrel{kL=\pi /4}{=}\frac{1}{2}\gamma dt+\gamma \eta a+a^{}_c(t)dt+O(\eta ^2).$$
(31)
The first term is elastic scattering. The second term, which is first order in $`\eta `$, is proportional to $`\stackrel{~}{z}a+a^{}`$, i.e. includes information on the ion motion. Here and in the following we take the center of the trap to be on the slope of the standing wave, i.e. $`k_{eg}L=\pi /4`$.
Following the analysis of homodyne detection GZ ; CarmichaelBook , we split the stochastic variable $`dN_c(t)`$ into a deterministic and a remaining stochastic part, thus defining $`dY_c\left(t\right)`$,
$$dN_c(t)\frac{1}{2}\gamma dt+\eta dY_c\left(t\right)$$
(32)
and we can show (cf. Appendix C) that $`dY_c(t)`$ is a Gaussian stochastic variable with non-zero mean, i.e.
$$dY_c(t)=\sqrt{\gamma /2}/\eta dW(t)+\gamma \stackrel{~}{z}_c(t)dt$$
with $`dW(t)`$ a Wiener increment satisfying $`dW^2(t)=dt`$.
This leads us to define a photocurrent where we subtract the large constant contribution from the elastic scattering process,
$`I_c(t)`$ $`=\eta {\displaystyle \frac{dY_c(t)}{dt}}`$ (33)
$`=\gamma \eta \stackrel{~}{z}_c(t)+\sqrt{{\displaystyle \frac{\gamma }{2}}}\xi \left(t\right).`$
with $`\xi \left(t\right)`$ Gaussian white noise $`\xi \left(t\right)\xi \left(t^{}\right)=\delta \left(tt^{}\right)`$ (shot noise). We see that $`I_c(t)`$ follows $`\stackrel{~}{z}_c(t)`$ and thus represents a continuous measurement of the position of the ion. The information on the motion is contained in the sidebands of the current, i.e. in the frequency components centered around $`\pm \nu _T`$.
In the diffusive approximation the conditional density matrix $`\rho _c(t)`$ GZ ; CarmichaelBook obeys
$$(\mathrm{I})d\rho _c(t)=\left[_0dt+\sqrt{\frac{\gamma }{2}}dW(t)_m\right]\rho _c(t)$$
(34)
where
$$_m\rho _c(t)=2\eta \left(\stackrel{~}{z}\rho _c(t)+\rho _c(t)\stackrel{~}{z}2\stackrel{~}{z}_c(t)\rho _c(t)\right)$$
(35)
and Eq. (34) is derived from (30) in the diffusive limit $`\eta 1`$ (cf. Appendix C).
## III Quantum Feedback Cooling
In the previous section we have reformulated the continuous observation of the ion motion through spontaneous light scattering into mirror modes in a form reminiscent of homodyne detection. This will allow us below to study feedback cooling of trapped ions building on the *Wiseman-Milburn theory of quantum feedback* WM\_coll ; Feedback\_Squeezing .
In Eq. (33) we have obtained a current which is proportional to the mean value of the *position* of the atom. We want to use this information to feed back an appropriate force proportional to the *momentum* to damp the motional state of the atom BushevThesis ; FeedbackCoolingExp . The information about the position is encoded in the motional sidebands of the current. In a harmonic trap of known frequency any combination of the average position and momentum can be obtained by shifting the sideband current by a phase of $`\varphi `$, if the trap frequency is much faster than any other (cooling) timescale in the problem (weak coupling limit). This phase $`\varphi `$ can be controlled electronically, and for $`\varphi =\pi /2`$, the shifted current follows the momentum. A force, which is proportional to this current, can damp the motion of the ion.
### III.1 Feedback current
We model the feedback circuit as shown in Fig. 3. First, the signal $`I_c(t)`$ given by Eq. (33) is mixed with a local oscillator of frequency $`\omega _0\nu _T`$ to shift the signal of the motional sideband to zero frequency. Then the current is sent through a band pass filter of width $`B`$ to cut off rapidly oscillating terms. The filter is described by a filter function $`Z(\omega )`$, centered around zero frequency. At the end the signal is mixed again with the local oscillator and amplified by a factor $`G`$. The feedback current can then be written as
$$I_{\mathrm{fb},c}\left(t\right)=G\mathrm{cos}(\omega _0t)_{\mathrm{}}^t𝑑\tau \stackrel{~}{Z}(t\tau )\mathrm{cos}\left(\omega _0\tau +\varphi \right)I_c(\tau ),$$
(36)
where $`\stackrel{~}{Z}\left(\tau \right)`$ is the Fourier transform of the band pass function $`Z(\omega )`$. The feedback Hamiltonian is specified in the next subsection.
To evaluate the expression for the current, it is convenient to change to a basis which is rotating with the frequency of the local oscillator $`\omega _0`$ by applying the unitary transformation $`U\mathrm{exp}(\mathrm{i}\omega _0a^{}at)`$. The evolution timescale of the density operator in this new frame, $`\stackrel{~}{\rho }_c(t)U\rho _c(t)U^{}`$ is determined by the detuning $`\delta =\omega _0\nu _T`$ and the cooling rates $`G\gamma ,\mathrm{\Gamma }_{\mathrm{eff}}`$. Under the assumption, that these frequencies are smaller than the filter bandwidth $`B`$, the feedback current is given by
$$I_{\mathrm{fb},c}\left(t\right)=G\left[\gamma \eta X_\varphi _c^I(t)+\sqrt{\frac{\gamma }{2}}\mathrm{\Xi }\left(t\right)\right]\mathrm{cos}\left(\omega _0t\right).$$
(37)
The first term in this expression, $`X_\varphi _c^I\mathrm{Tr}_{\mathrm{sys}}\{X_\varphi \stackrel{~}{\rho }_c(t)\}`$ is the slowly varying expectation value of the quadrature component
$$X_\varphi a\mathrm{e}^{\mathrm{i}\varphi }+a^{}\mathrm{e}^{\mathrm{i}\varphi }$$
(38)
(in the rotating frame). The second contribution in Eq. (37) is defined as
$$\mathrm{\Xi }(t)_{\mathrm{}}^t𝑑\tau \mathrm{cos}\left(\omega _0\tau +\varphi \right)\stackrel{~}{Z}(t\tau )\xi \left(\tau \right).$$
(39)
It describes the noise which passes through the feedback circuit. The stochastic mean of $`\mathrm{\Xi }(t)`$ is zero due to the vanishing mean of the white noise variable $`\xi (t)`$, and the correlation function is given by
$$\mathrm{\Xi }(t)\mathrm{\Xi }(t^{})\delta _B\left(tt^{}\right)+𝒪\left(\frac{B}{\omega _0}\right).$$
(40)
Here $`\delta _B(tt^{})`$ denotes a delta-function for functions which vary on a slow timescale much larger than $`B^1`$.
Thus for a clear separation of timescales,
$$G\gamma ,\delta ,\mathrm{\Gamma }_{\mathrm{eff}}B\omega _0,\nu _T,$$
(41)
the current given in Eq. (37) is proportional to the slowly varying expectation value of $`X_\varphi `$, and has a noise term which is delta-correlated on a timescale of the system evolution in the rotating frame.
### III.2 Quantum Feedback Dynamics
The feedback current ((37)) for $`\varphi =\pi /2`$ is proportional to the slowly varying momentum of the particle. For the cooling of the ion motion, we apply a linear force which is proportional to the the feedback current (37). For a trapped ion, this can be realized by applying a voltage on the trap electrodes, which leads to a displacement of the trap center. The effect of the feedback force is given by the interaction picture Hamiltonian
$$H_{\mathrm{fb}}=I_{\mathrm{fb},c}(t\tau )\stackrel{~}{z}_I(t).$$
(42)
In this equation, $`\stackrel{~}{z}_I(t)U^{}\stackrel{~}{z}U`$ is proportional to the position operator in the interaction picture, while $`\tau `$ denotes the finite time delay in the feedback loop, which we require to be much smaller than the trap frequency $`\tau 1/\nu _T`$. The master equation (34) has to be complemented with the feedback term,
$$(\mathrm{S})\left[d\stackrel{~}{\rho }_c(t)\right]_{\mathrm{fb}}=I_{\mathrm{fb},c}(t\tau )\left(\mathrm{i}\right)[\stackrel{~}{z}_I(t),\stackrel{~}{\rho }_c(t)]dt.$$
(43)
which has to be interpreted as a *Stratonovich* stochastic differential equation Feedback\_Squeezing . For the slow dynamics of the density matrix in the rotating frame, we can make a rotating wave approximation and neglect rapidly rotating terms $`\mathrm{exp}(\pm 2\mathrm{i}\omega _0t)`$. The filtered noise (39) is delta-correlated on timescales slower than $`B^1`$, thus we have the coarse grained evolution of the density matrix
$`(\mathrm{S})\left[d\stackrel{~}{\rho }_c(t)\right]_{\mathrm{fb}}`$ $`={\displaystyle \frac{G}{2}}\gamma \eta X_\varphi _c^I(t\tau )dt𝒦\stackrel{~}{\rho }_c(t)`$ (44)
$`+{\displaystyle \frac{G}{2}}\sqrt{{\displaystyle \frac{\gamma }{2}}}dW_\mathrm{\Xi }(t\tau )𝒦\stackrel{~}{\rho }_c(t),`$
with the feedback operator
$$𝒦\stackrel{~}{\rho }_c(t)\mathrm{i}[\stackrel{~}{z},\stackrel{~}{\rho }_c(t)]$$
(45)
and the “slow” Wiener increment $`dW_\mathrm{\Xi }\left(t\right)\mathrm{\Xi }\left(t\right)dt`$.
The total evolution of the system is determined by the conditioned master equation (34) plus the contribution from the feedback loop (44). To combine the two equations, we have to convert Eq. (44) from Stratonovich to Ito form. The total conditioned evolution is
$`(\mathrm{I})d\stackrel{~}{\rho }_c(t)`$ $`=\stackrel{~}{}_0\stackrel{~}{\rho }_c+\sqrt{{\displaystyle \frac{\gamma }{2}}}dW(t)\stackrel{~}{\rho }_c(t)`$
$`+({\displaystyle \frac{G}{2}}\gamma \eta X_\varphi _c^I(t\tau )dt+{\displaystyle \frac{G^2}{16}}\gamma 𝒦dt+`$
$`+{\displaystyle \frac{G}{2}}\sqrt{{\displaystyle \frac{\gamma }{2}}}dW_\mathrm{\Xi }(t\tau ))𝒦\stackrel{~}{\rho }_c(t),`$ (46)
where
$$\stackrel{~}{}\stackrel{~}{\rho }_c_{\mathrm{LC}}\stackrel{~}{\rho }_c\mathrm{i}[\delta a^{}a,\stackrel{~}{\rho }_c]$$
(47)
(cf. Eq. (25)) is the laser cooling Liouvillian in the rotating frame.
Because the exact photocurrent can not be kept track of in experiments, Eq. (46) is of limited use. The goal is to derive an equation for the ensemble averaged density operator. We follow the derivation given by Wiseman and Milburn in Feedback\_Squeezing , where the measured current is fed back directly, and adopt it for our model. Assuming that the state at time $`t\tau `$ and all previous times is known, we take the ensemble average $`E[]`$ of Eq. (46) over the trajectories in $`(t\tau ,t]`$. We then formally divide by $`dt`$ and for convenience redefine $`\rho (t)E[\stackrel{~}{\rho }_c(t)]`$:
$`(\mathrm{I})\dot{\rho }(t)=`$ $`\stackrel{~}{}\rho (t)+{\displaystyle \frac{G}{2}}\gamma \eta X_\varphi _c^I(t\tau )𝒦\rho (t)`$ (48)
$`+`$ $`{\displaystyle \frac{G}{2}}\sqrt{{\displaystyle \frac{\gamma }{2}}}𝒦E\left[\mathrm{\Xi }(t\tau )\stackrel{~}{\rho }_c(t)\right]+{\displaystyle \frac{G^2}{16}}\gamma 𝒦^2\rho (t).`$
The density matrix $`\rho (t)`$ is still conditioned on the evolution up to time $`t\tau `$, but not conditioned on trajectories in $`(t\tau ,t]`$. The ensemble average $`E[X_\varphi ^I_c(t\tau )\stackrel{~}{\rho }_c(t)]`$ factorizes because $`\rho _c(t\tau )`$ is assumed known. Under the Markov approximation, we let $`\tau `$ go to zero, while due to the coarse graining of the time evolution in Eq. (43), $`dt`$ will still be larger than this small delay. An expansion in $`\tau `$ yields
$`\stackrel{~}{\rho }_c(t)`$ $`=\left[1+𝒪\left(\tau \right)\right]\stackrel{~}{\rho }_c(t\tau +dt)=`$ (49)
$`=\left[1+𝒪\left(\tau \right)\right]\left[1+\sqrt{{\displaystyle \frac{\gamma }{2}}}dW(t\tau )\right]\stackrel{~}{\rho }_c(t\tau ).`$
We now can evaluate the remaining ensemble average in Eq. (48) because $`dW(t\tau )`$ is stochastically independent from $`\stackrel{~}{\rho }_c(t\tau )`$. We obtain
$`E\left[\mathrm{\Xi }(t\tau )\stackrel{~}{\rho }_c(t)\right]=`$ $`\sqrt{\gamma }E\left[\mathrm{\Xi }(t\tau )\xi (t\tau )\right]\rho (t)`$ (50)
$``$ $`\sqrt{{\displaystyle \frac{\gamma }{2}}}\eta (X_\varphi \rho (t)+\rho (t)X_\varphi `$
$`2X_\varphi _c^I(t\tau )\rho (t)),`$
and thus the term in the last line, a conditional expectation value, cancels with the second term on the right hand side of Eq. (48). In going from the first to the second line in Eq. (50) we have dropped terms $`\mathrm{exp}(\pm \mathrm{i}\omega _0t)`$.
With this last step, we can finally evaluate Eq. (48) and write down the *quantum feedback master equation* (compare for the motional degrees of freedom:
$$\dot{\rho }=\stackrel{~}{}\rho +\frac{G}{4}\gamma \eta 𝒦\left(X_\varphi \rho +\rho X_\varphi \right)+\frac{G^2}{16}\gamma 𝒦^2\rho .$$
(51)
The first term on the right hand side $`\stackrel{~}{}`$ is the laser cooling Liouvillian (47) in the rotating frame. The second term with $`𝒦`$ given in Eq. (45) in the master equation is the feedback term. It acts back on the system and is responsible for cooling if we choose the parameters $`\delta `$ and $`\varphi `$ appropriately. The last term in the master equation is a diffusive term of the form of a double commutator.
## IV Results
In the last section we have shown that for a separation of timescales $`\delta ,\mathrm{\Gamma }_{\mathrm{eff}}B\omega _0,\nu _T`$ we obtain an unconditioned (non-selective) master equation for the motional density matrix in the rotating frame. By inserting the definitions of $`\stackrel{~}{}`$ and $`𝒦`$ the master equation reads
$`\dot{\rho }`$ $`=\mathrm{i}\delta [a^{}a,\rho ]+A_{}𝒟[a]+A_+𝒟[a^{}]+`$ (52)
$`\mathrm{i}{\displaystyle \frac{G}{4}}\gamma \eta [\stackrel{~}{z},X_\varphi \rho +\rho X_\varphi ]{\displaystyle \frac{G^2}{16}}\gamma [\stackrel{~}{z},[\stackrel{~}{z},\rho ]].`$
We have used the previously introduced variables $`\stackrel{~}{z}=a+a^{}`$ and $`X_\varphi =a\mathrm{e}^{\mathrm{i}\varphi }+a^{}\mathrm{e}^{\mathrm{i}\varphi }`$. In the first line of Eq. (52) we recover the master equation for laser cooling, with the corresponding heating and cooling rates $`A_\pm `$ given in Eq. (23). The second line describes the effect of the feedback loop, where $`\gamma =\epsilon N\mathrm{\Gamma }_{\mathrm{eff}}/(1+\alpha )\eta ^2`$ is the emission rate in the mirror mode and $`G`$ is the gain parameter amplifying the feedback current. The first term in the second line depends on the phase shift $`\varphi `$ and as we will show below, leads to the expected damping for $`\varphi =\pi /2`$. The second term arises from the noise in the feedback current and leads to a momentum diffusion, i.e. heating.
We will derive solutions of the feedback master equation (51), which is bilinear in the position and momentum $`\widehat{z}`$ and $`\widehat{p}_z`$. It is convenient to use a Wigner function representation GZ of the density matrix. This gives rise to a Fokker-Planck equation Risken for the Wigner function $`W(\overline{z},\overline{p},t)`$ with dimensionless position and momentum variables $`x_1\overline{z}=z\sqrt{m\nu _T/2}`$ and $`x_2\overline{p}=p_z/\sqrt{2m\nu _T}`$,
$`{\displaystyle \frac{W(\overline{z},\overline{p},t)}{t}}=`$ $`{\displaystyle \underset{i,j}{}}\kappa _{ij}{\displaystyle \frac{}{x_i}}\left(x_jW(\overline{z},\overline{p},t)\right)+`$ (53)
$`+{\displaystyle \underset{i,j}{}}D_{ij}{\displaystyle \frac{^2W(\overline{z},\overline{p},t)}{x_ix_j}}.`$
The $`\kappa _{ij}`$ are independent of the phase space variables and $`D_{ij}`$ is diagonal, thus Eq. (53) describes an Ornstein-Uhlenbeck process Risken with drift matrix
$$\kappa =\frac{\mathrm{\Gamma }_{\mathrm{eff}}}{2}\left(\begin{array}{cc}1& 2\stackrel{~}{\delta }\\ 2G\eta \stackrel{~}{\gamma }\mathrm{cos}\varphi +2\stackrel{~}{\delta }& 12G\eta \stackrel{~}{\gamma }\mathrm{sin}\varphi \end{array}\right)$$
(54)
and the diagonal terms of the diffusion matrix
$$(D_{11},D_{22})=\frac{\mathrm{\Gamma }_{\mathrm{eff}}}{8}(2N+1,2N+1+\frac{1}{2}G^2\stackrel{~}{\gamma }).$$
(55)
Here we have introduced the dimensionless detuning $`\stackrel{~}{\delta }\delta /\mathrm{\Gamma }_{\mathrm{eff}}`$ and decay rate $`\stackrel{~}{\gamma }\gamma /\mathrm{\Gamma }_{\mathrm{eff}}`$ normalized with respect to the width of the sidebands. The Gaussian Wigner function is uniquely determined by it’s first and second position and momentum moments, and we will use the notation
$$\overline{z}^r\overline{p}^s_W𝑑\overline{z}𝑑\overline{p}\overline{z}^r\overline{p}^sW(\overline{z},\overline{p},t),$$
(56)
which equals the symmetric expectation value of the corresponding operators. The bilinearity of Eq. (51) with respect to position and momentum gives rise to a closed set of equations for the first and second moments of the Wigner function individually and are given in Appendix D.
We are interested in the motional energy of the ion, which is related to the expectation value of the phonon number by $`E=\mathrm{}\nu _T(a^{}a+1/2)`$. The expectation value for the number operator can be read off from the second moments of the Wigner function:
$$a^{}an=\overline{z}^2_W+\overline{p}^2_W\frac{1}{2}$$
(57)
We will calculate this quantity for different choices of parameters in the following subsections.
### IV.1 Cold damping
In this subsection we show results for $`\varphi =\pi /2`$ and $`\delta =0`$, i.e. the center of the band pass filter is set exactly to the trap frequency. As derived in Appendix D the number expectation value for the steady state in this case is given by
$$n_{ss}=\frac{N+\frac{1}{2}\eta \stackrel{~}{\gamma }\left(2N1\right)G+\frac{1}{8}\stackrel{~}{\gamma }G^2}{1+2\eta \stackrel{~}{\gamma }G}$$
(58)
Taking the gain $`G=0`$ yields $`n_{ss}=N`$, i.e. if we do not use the feedback current to influence the ion, the steady state occupation will be the one for standard laser cooling. We see that the slope of the occupation number is negative at $`G=0`$, i.e.
$$n_{ss}/G|_{G=0}=\stackrel{~}{\gamma }\eta (2N+1)/2<0,$$
(59)
and for $`G\mathrm{}`$ it diverges (note that in our model $`G\gamma `$ has to be smaller than $`B`$). Thus our theory yields a non-vanishing optimal gain $`G_{\mathrm{min}}`$ for which the occupation number has a minimum smaller than $`N`$,
$$G_{\mathrm{min}}=\frac{\sqrt{1+8(2N+1)\eta ^2\stackrel{~}{\gamma }}1}{2\eta \stackrel{~}{\gamma }}.$$
(60)
Inserting this into Eq. (58) yields an expression for the minimal occupation number:
$$n_{\mathrm{min}}=\frac{4(2N1)\eta ^2\stackrel{~}{\gamma }1+\sqrt{1+8(2N+1)\eta ^2\stackrel{~}{\gamma }}}{16\eta ^2\stackrel{~}{\gamma }}.$$
(61)
With increasing solid angle $`\epsilon `$ we collect more information about the motional state of the system and hence the minimum $`n_{ss}`$ is expected to decrease, which is shown in Fig. 4. With increasing $`\epsilon `$ the optimal gain is decreasing, because the feedback noise term is growing with $`G^2`$ while the damping term is linear in $`G`$.
We show in Fig. 5 the decrease in the steady state phonon number with the gain. The relative decrease is larger with a higher laser cooling steady state phonon number $`N`$. For lower $`N`$, the mirror decay rate $`\gamma N`$ is smaller and thus we get less information about the motional state of the atom, which limits the feedback cooling.
We will now expand $`n_{ss}`$ in the limit of large ($`N1`$) occupation numbers. For a series expansion of (61) the formal expansion parameter is $`N\sqrt{\epsilon }`$, thus an expansion in the (usually also small) $`\epsilon `$ is only possible for very low $`N`$. We make an expansion for large $`N`$ in the opposite limit (Doppler limit), while the condition $`N\sqrt{\epsilon }1`$ has to be satisfied. $`N`$ can be tuned with e.g. with the laser detuning $`\mathrm{\Delta }_L`$. Then the minimal occupation number approximately reads
$$n_{\mathrm{min}}=\frac{N}{2}+4\sqrt{\frac{1+\alpha }{\epsilon }}\frac{1+\alpha }{N\epsilon },$$
(62)
which implies that for a sufficiently large collection angle the minimal obtainable phonon number is above $`N/2`$ and thus feedback cooling alone cannot give a steady state. The reduction in the energy of the ion with time is due to the reduction in $`\overline{p}^2_W`$, while $`\overline{z}^2_W`$ is constant, as is shown in the time evolution in Fig. 6. Thus the Wigner function for the steady state will not be rotationally invariant, but “classically squeezed” in the momentum direction.
A phase space picture can demonstrate the action of the feedback on the system state (see Fig. 9(a)). By feeding back a linear force $`f`$ to the ion, we effectively apply a unitary operator of the form
$$U(t)\mathrm{exp}\left(\mathrm{i}fxt\right).$$
(63)
This operator acts as a momentum kick on a state with a magnitude proportional to the momentum, which we have chosen by setting $`\varphi =\pi /2`$. The points in the Wigner function will tend towards the x-axis, while the diffusion term will counteract the feedback term, leading to a steady state Wigner function.
The difference in the position and momentum variance can be quantified; we will give an expression for the amount of “squeezing”, i.e. the ratio between the two half-axis of the error-ellipse for the Wigner function in phase space is obtained by rotating the axes of the ellipse:
$`r_\sigma `$ $`{\displaystyle \frac{\text{semiminor axis}}{\text{semimajor axis}}}={\displaystyle \frac{1f}{1+f}}`$ (64a)
$`f`$ $`{\displaystyle \frac{\sqrt{\left(\sigma _{zz}\sigma _{pp}\right)^2+4\sigma _{zp}^2}}{\sigma _{zz}+\sigma _{pp}}}`$ (64b)
Here $`\sigma _{zz}=\overline{z}^2_W\overline{z}_W^2`$ and $`\sigma _{pp}=\overline{p}^2_W\overline{p}_W^2`$ are the variances of position and momentum, respectively, and $`\sigma _{zp}=\overline{z}\overline{p}_W\overline{z}_W\overline{p}_W`$. As mentioned, due to the affection of only the $`\sigma _{pp}`$ component, $`\sigma _{zp}=0`$ in the case $`\varphi =\pi /2`$. The range of the squeezing parameter is $`0<r_\sigma 1`$, where a small value corresponds to strong squeezing and for $`r_\sigma =1`$ the state is symmetric.
The time dependent Fokker-Planck equation is solvable analytically and the timescale of the cooling process is given by the eigenvalues of the drift matrix (54), which are in this case $`\mathrm{\Gamma }_{\mathrm{eff}}`$ and $`\mathrm{\Gamma }_{\mathrm{eff}}+2\eta \gamma G`$ corresponding to the usual Doppler cooling and the feedback cooling. This shows that the feedback cooling happens on a timescale faster than laser cooling alone.
### IV.2 Variable feedback phase
For a phase $`\varphi \pi /2`$, the magnitude of the feedback force is proportional to the projection of the momentum on an other rotated axis in phase space. We have pointed out in Eq. (63) that the action of the linear force (shifted trap) is always a momentum kick. Thus the particle will always be “kicked too hard” or not hard enough towards the phase space center. We will calculate the regions of stability where the feedback can still lead to a steady state. Such a steady state will only form if the both eigenvalues of the matrix $`\kappa `$ are positive. One eigenvalue of this matrix is $`\mathrm{\Gamma }_{\mathrm{eff}}`$ for arbitrary $`\varphi `$, giving again the usual Doppler cooling, and the other eigenvalue is $`\mathrm{\Gamma }_{\mathrm{eff}}2G\eta \gamma \mathrm{sin}\varphi `$, which is always positive for negative angles. For positive angles $`\varphi >0`$, the gain has to fulfill the condition $`G<\mathrm{\Gamma }_{\mathrm{eff}}/2\gamma \eta \mathrm{sin}\varphi `$. If this condition is satisfied, a steady state number expectation value exists and reads:
$`n_{ss}`$ $`=\left[(1\eta \stackrel{~}{\gamma }G\mathrm{sin}\varphi )(12\eta \stackrel{~}{\gamma }G\mathrm{sin}\varphi )\right]^1\times `$ (65)
$`\times [N+{\displaystyle \frac{1}{2}}(4N1)\eta \stackrel{~}{\gamma }G\mathrm{sin}\varphi +`$
$`+{\displaystyle \frac{1}{8}}\stackrel{~}{\gamma }G^2\left(1+4\stackrel{~}{\gamma }\eta ^2\left(2N+12\mathrm{sin}^2\varphi \right)\right)`$
$`\eta \stackrel{~}{\gamma }^2G^3\mathrm{sin}\varphi ].`$
From Eq. (65) we can see that an energy decrease via feedback cooling is only possible for angles $`\pi <\varphi <0`$ by calculating the slope $`n_{ss}/G|_{G=0}`$. Because Eq. (65) is of higher order in $`G`$ than the equation we had for $`\varphi =\pi /2`$, (58), we will not give an analytical solution for the minimal gain and number occupation here. We also find that for $`\varphi \pi /2`$ the optimal occupation number is higher than $`\varphi =\pi /2`$ (compare related studies in VitaliPhase ). The steady state occupation number for varying $`\varphi `$ as a function of the gain is plotted in Fig. 7, where we can see that for non-optimal phases the range of $`G`$ for $`n_{ss}<N`$ is shrinking.
For the special case of $`\varphi =\pi `$ or $`\varphi =0`$, no cooling can be observed any more and the number expectation value is quadratic in $`G`$. In principle, a steady state with $`n_{ss}>N`$ always exists with
$$n_{ss}=N+\frac{1}{8}\stackrel{~}{\gamma }G^2\left(1+4\eta ^2\stackrel{~}{\gamma }\left(2N+1\right)\right).$$
(66)
The more interesting feature of the $`\varphi =\pi `$ case is that in the master equation (51) the feedback term (second term) reduces in a rotating wave approximation to a Hamiltonian term of the form $`\mathrm{i}[\mathrm{\Delta }\nu a^{}a,\rho ]`$. For this case we observe a small shift $`\mathrm{\Delta }\nu `$ in the frequency of the trap linearly proportional to the gain. In this paper, we have not discussed the detailed experimental setup used to apply the force to the ion, which would be necessary for the knowledge of the exact forces acting on the ion. For $`\varphi =\pi `$ one can measure the frequency shift in the location of the sideband and determine the conversion factor from the gain parameter $`G`$ used in this paper and an experimental gain factor, which might be the real electronic gain in the feedback loop.
### IV.3 Rotation in Phase Space
We have shown that the phase $`\varphi =\pi /2`$ we chose leads to the lowest energy of the motional state of the ion. The variance for the position operator $`\overline{z}^2_W=\left(2N+1\right)/4`$ remains constant with time as shown e.g. in Fig. 6, thus posing a lower limit to the obtainable energy. The detuning $`\delta `$ of the local oscillator in the feedback loop from the trap frequency creates a tunable slow rotation of the (interaction picture) Wigner function in phase space. This results in “squeezing” of all quadrature components (see Fig. 9(b)), and the Wigner function can regain a symmetric shape. Of course the timescale for this rotation has to be much slower than the filter bandwidth $`B`$.
For the time evolution of the variances, the effect of the detuning is illustrated in Fig. 8. We see the time evolution of an initially thermal (symmetric) state with an occupation number of $`N`$. In contrast to Fig. 6 the width of the Wigner function in the momentum and the position space are alternately decreased until they reach the new feedback steady value. For a larger detunings the two variances decrease equally in time and energetically lower states can be reached.
For a rotation of the Wigner function with the frequency $`\delta `$, we have to compare this rotation timescale with the cooling timescale $`\gamma `$. For $`\gamma `$ comparable to $`\delta `$ the optimal phase is is shifted with respect to $`\pi /2`$ because the Wigner function is rotating in phase space during the cooling time. When the detuning is much larger than the cooling rate, the Wigner function ellipse direction will not be resolved during the cooling time and thus the optimal phase returns to $`\pi /2`$. By numerical optimization (Fig. 10) we find that the optimal phase is shifted from $`\pi /2`$ asymmetrically with respect to the detuning $`\delta `$. It reaches it’s maximum excursion for a value of $`\delta /\mathrm{\Gamma }_{\mathrm{eff}}1`$, for higher detunings the optimal phase approaches $`\pi /2`$ again. For these optimal values, we plot in Fig. 11 the squeezing parameter $`r_\sigma `$ (64a), which is one for a symmetric Gaussian state. We see that the state at no detuning is “classically squeezed” as we already mentioned in subsection IV.1 and the squeezing increases up to $`\stackrel{~}{\delta }1`$, then upon approaching $`\stackrel{~}{\delta }\mathrm{}`$ the squeezing parameter approaches one, and the state is thermal.
For increasing $`\delta `$, we also show that the number expectation value is decreasing. We will not give an analytic expression for $`n_{ss}`$ for an arbitrary $`\delta `$ here. We merely calculate the minimal number of phonons in the limit of large $`\delta `$. For this we require an additional separation of the timescales between the effective feedback cooling rate and the detuning, while the other timescale inequalities still hold:
$$\gamma \delta B.$$
(67)
With these new conditions we take the detuning $`\delta \mathrm{}`$, where the optimal feedback phase is again $`\varphi =\pi /2`$, and get for the occupation number:
$$n_{ss}=\frac{N\frac{1}{2}\eta \stackrel{~}{\gamma }G+\frac{1}{8}\stackrel{~}{\gamma }G^2}{1+\eta \stackrel{~}{\gamma }G}.$$
(68)
The minimal occupation number for the same limit we took in deriving Eq. (62) we get for $`N1`$:
$$n_{\mathrm{min}}\sqrt{\frac{1+\alpha }{2\epsilon }}\frac{1}{2}\frac{2(1+\alpha )}{8\epsilon N}.$$
(69)
This expression does not include the large term $`N/2`$ any more and thus the obtainable energy for large $`N`$ has an upper bound which is independent of $`N`$, thus feedback cooling alone can give a thermal (symmetric) state with a temperature below the Doppler temperature.
## V Conclusion
In this paper we have studied quantum feedback cooling of a trapped ion in front of a mirror. This work is motivated by recent experiments BushevThesis , and – as shown in FeedbackCoolingExp – provides a quantitative understanding of the experimental results.
In the setup discussed in this paper the final temperatures are limited by the collection efficiency, $`\epsilon `$, and the constant scattering of photons for the position measurement. This combination of heating due to the recoil, and laser cooling due to the red detuning of the laser leads to a steady state temperature (Doppler limit). The effect of quantum feedback cooling is studied as an additional cooling mechanism on top of the ongoing laser cooling. For the experimentally relevant parameters this leads to sub-Doppler cooling, but it seems difficult to achieve ground state cooling in the trap along these lines. As shown in a parallel publication Rabl , we can devise a purely *dispersive* and thus non-invasive readout of the velocity of the trapped ion based on the variation of the index of refraction with velocity, i.e. based on electromagnetically induced transparency. Such a scheme allows, under idealized conditions, ground state cooling of the ion purely by quantum feedback.
###### Acknowledgements.
The authors thank R. Blatt, F. Dubin, J. Eschner, and D. Rotter for discussions which motivated the present work. Research at the University of Innsbruck is supported by the Austrian Science Foundation, EU projects and the Institute of Quantum Information.
## Appendix A Derivation of the Quantum Stochastic Master Equation (13)
Starting from the stochastic master equation (10) we define the reduced density matrix $`\widehat{W}(t)\mathrm{Tr}_b\left\{|\mathrm{\Psi }(t)\mathrm{\Psi }(t)|\right\}`$. Note that $`\widehat{W}(t)`$ is now a trace-class operator for the internal electronic, the motional and the mirror mode bath degrees of freedom. We calculate
$$d\widehat{W}(t)=\mathrm{Tr}_b\left\{|\mathrm{\Psi }(t+dt)\mathrm{\Psi }(t+dt)||\mathrm{\Psi }(t)\mathrm{\Psi }(t)|\right\}$$
(70)
by inserting $`|\mathrm{\Psi }(t+dt)=|\mathrm{\Psi }(t)+d|\mathrm{\Psi }(t)`$ from Eq. (10). Using the Ito rules $`dB_u(t)dB_u^{}^{}(t)=\delta (uu^{})dt`$, and cyclic property of the trace for background bath operators, all terms of the form
$`\mathrm{Tr}_b\left\{dB_u^{}(t)|\mathrm{\Psi }(t)\mathrm{\Psi }(t)|\right\}=`$ (71)
$`=\mathrm{Tr}_b\left\{|\mathrm{\Psi }(t)\mathrm{\Psi }(t)|dB_u(t)\right\}=0`$
vanish because the initial bath state is the vacuum state. With these rules we obtain Eq. (13).
## Appendix B Adiabatic Elimination, Lamb-Dicke Limit and Laser Cooling
This appendix fills in the details of deriving the QSME (17) from the QSSE (10) under the assumption of weak driving and small Lamb-Dicke parameter. Note that we will need to consider two different Lamb-Dicke parameters due to the exciting laser which is not collinear with the $`z`$-axis. As in section II.1 we denote $`\stackrel{~}{\eta }\eta \mathrm{sin}\chi `$. Inserting the ansatz (15) into the QSSE (10) and transforming to an interaction picture with respect to $`H_T`$ we get
$`|\psi _e\left(t\right)=`$ $`{\displaystyle \frac{\mathrm{i}\mathrm{\Omega }}{2}}({\displaystyle \frac{1\frac{1}{2}\stackrel{~}{\eta }^2a^{}a}{\mathrm{i}\mathrm{\Delta }_L+\frac{\mathrm{\Gamma }}{2}}}+`$ (72)
$`+`$ $`{\displaystyle \frac{\mathrm{i}\stackrel{~}{\eta }a\mathrm{e}^{\mathrm{i}\nu _Tt}}{\mathrm{i}\left(\mathrm{\Delta }_L\nu _T\right)+\frac{\mathrm{\Gamma }}{2}}}+{\displaystyle \frac{\mathrm{i}\stackrel{~}{\eta }a^{}\mathrm{e}^{\mathrm{i}\nu _Tt}}{\mathrm{i}\left(\mathrm{\Delta }_L+\nu _T\right)+\frac{\mathrm{\Gamma }}{2}}})|\psi _g\left(t\right).`$
We insert this expression back into (10). We obtain
$`d|\psi _g\left(t\right)=`$ $`\{{\displaystyle \frac{\mathrm{\Omega }^2}{4}}[{\displaystyle \frac{1}{\mathrm{i}\mathrm{\Delta }_L+\frac{\mathrm{\Gamma }}{2}}}+𝒪\left(\mathrm{e}^{\pm 2\mathrm{i}\nu _Tt}\right)+`$ (73)
$`+{\displaystyle \frac{\stackrel{~}{\eta }^2a^{}a}{\mathrm{i}\left(\mathrm{\Delta }_L\nu _T\right)+\frac{\mathrm{\Gamma }}{2}}}+{\displaystyle \frac{\stackrel{~}{\eta }^2aa^{}}{\mathrm{i}\left(\mathrm{\Delta }_L+\nu _T\right)+\frac{\mathrm{\Gamma }}{2}}}]dt+`$
$`+`$ $`[{\displaystyle \frac{\mathrm{i}\mathrm{\Omega }}{2}}{\displaystyle \frac{1}{\mathrm{i}\mathrm{\Delta }_L+\frac{\mathrm{\Gamma }}{2}}}dC_1^{}`$
$``$ $`\left({\displaystyle \frac{\mathrm{i}\stackrel{~}{\eta }a\mathrm{e}^{\mathrm{i}\nu _Tt}}{\mathrm{i}\left(\mathrm{\Delta }_L\nu _T\right)\frac{\mathrm{\Gamma }}{2}}}+{\displaystyle \frac{\mathrm{i}\stackrel{~}{\eta }a^{}\mathrm{e}^{\mathrm{i}\nu _Tt}}{\mathrm{i}\left(\mathrm{\Delta }_L+\nu _T\right)\frac{\mathrm{\Gamma }}{2}}}\right)dC_1^{}+`$
$`+`$ $`{\displaystyle \frac{\eta \left(a\mathrm{e}^{\mathrm{i}\nu _Tt}+a^{}\mathrm{e}^{\mathrm{i}\nu _Tt}\right)}{\mathrm{i}\mathrm{\Delta }_L+\frac{\mathrm{\Gamma }}{2}}}dC_2^{}]\}|\psi _g\left(t\right)(\mathrm{I})`$
with
$`dC_1^{}`$ $`\sqrt{\mathrm{\Gamma }_b}{\displaystyle 𝑑u\sqrt{N\left(u\right)}𝑑B_u^{}}+\sqrt{\mathrm{\Gamma }_m}\mathrm{sin}(k_{eg}L)dB_m^{},`$ (74)
$`dC_2^{}`$ $`\mathrm{i}\sqrt{\mathrm{\Gamma }_b}{\displaystyle 𝑑u\sqrt{N\left(u\right)}u𝑑B_u^{}}+\mathrm{cos}(k_{eg}L)dB_m^{}.`$ (75)
Consistent with the above approximations we neglect here and in the following terms oscillating at twice the trap frequency $`\nu _T`$. Physically speaking, the fourth line of Eq. (73) will correspond together with third line to a heating and cooling term, and the last line describes a diffusive term (cf. Fig. 2).
Taking the trace over the background modes to define a reduced density operator $`\widehat{w}(t)`$ according to (16) we use the Ito rules, e.g.
$$\mathrm{Tr}_b\left\{dB_u^{}\left(t\right)|\psi _g\left(t\right)\psi _g\left(t\right)|dB_u^{}\left(t\right)\right\}=\delta \left(uu^{}\right)\rho \left(t\right)dt$$
to derive Eq. (17).
## Appendix C Homodyne photodetection and the diffusion approximation
As we have seen in Sec. II.6, the statistics of the detected photons in the mirror mode are determined by the Poissonian stochastic variable $`dN_c(t)`$. Like in homodyne detection, where a strong local oscillator beats with the photodetection signal from a quantum system, an elastic scattering term beats with the signal given by the coupling of the light to the ion’s motion (cf. Eq. (31)). The parameter which gives the difference in the magnitudes of these terms is the Lamb-Dicke parameter $`\eta `$. We split the stochastic variable $`dN_c`$ into a constant (deterministic) part and a remaining stochastic part:
$$dN_c(t)\frac{1}{2}\gamma dt+\eta dY_c\left(t\right).$$
(76)
The stochastic expectation value of this equation is already known from Eq. (31):
$$dY_c(t)=\gamma \stackrel{~}{z}_c(t)dt.$$
(77)
We check the distribution properties by calculating
$`dY_c^2\left(t\right)`$ $`=\left({\displaystyle \frac{dN_c\left(t\right)\frac{1}{2}\gamma dt}{\eta }}\right)^2={\displaystyle \frac{dN_c\left(t\right)}{\eta ^2}}=`$ (78)
$`={\displaystyle \frac{\frac{1}{2}\gamma dt+\eta dY_c\left(t\right)}{\eta ^2}}\stackrel{\eta 1}{}{\displaystyle \frac{\gamma }{2\eta ^2}}dt,`$
which tells us that the stochastic variable has Gaussian properties, and thus is associated with a white noise probability distribution. Thus $`dY_c(t)=\sqrt{\gamma /2}/\eta dW(t)`$ where $`dW(t)`$ is a Wiener increment.
The evolution of the system conditioned on measuring the photocurrent can be seen by expanding the first bracket in the stochastic master equation (30) to first order in the Lamb-Dicke parameter $`\eta `$, and noting that the second bracket in (30) is
$$dN_c(t)dN_c(t)=\sqrt{\gamma /2}dW(t).$$
(79)
Thus, using the formal derivative $`\xi (t)=dW(t)/dt`$ we obtain the conditioned equation for the reduced density matrix, Eq. (34).
## Appendix D Equations of motion for the moments of the Wigner function
In Sec. IV we use a Wigner function representation for the density matrix and get an Fokker Planck equation (53) equivalent to the master equation Eq. (51) with the drift matrix (54) and the diffusion term (55). The equations of motion for the first and second moments of the Wigner function in terms of the normalized position and momentum variables $`\overline{z}=x_1`$ and $`\overline{p}=x_2`$, respectively, are:
$`{\displaystyle \frac{}{t}}x_i(t)_W=`$ $`{\displaystyle \underset{j}{}}\kappa _{ij}x_j(t)_W,`$ (80)
$`{\displaystyle \frac{}{t}}x_kx_l_W=`$ $`2D_{kl}+2D_{lk}`$ (81)
$`{\displaystyle \underset{j}{}}\left[\kappa _{kj}x_lx_j_W+\kappa _{lj}x_kx_j_W\right].`$
For a constant drift matrix, the equations for the first moments are trivial, and if the eigenvalues of $`\kappa `$ are positive, the steady state value is zero for both moments. We will therefore not concentrate on the first moments. We will give the equations for the second moments which are relevant for the number expectation value, and for this purpose we define a vector of second moments
$$𝐲(t)=(\overline{z}^2(t)_W,\overline{p}^2(t)_W,\overline{z}(t)\overline{p}(t)_W)^T.$$
(82)
We can write the equation of motion in a compact form as
$$\dot{𝐲}(t)=M𝐲(t)+𝐮$$
(83)
where the evolution matrix is
$$\frac{M}{\mathrm{\Gamma }_{\mathrm{eff}}}=\left(\begin{array}{ccc}1& 0& \stackrel{~}{\delta }\\ 0& 12\stackrel{~}{G}\mathrm{sin}\varphi & \stackrel{~}{G}\mathrm{cos}\varphi +\stackrel{~}{\delta }\\ 2\stackrel{~}{G}\mathrm{cos}\varphi +2\stackrel{~}{\delta }& 2\stackrel{~}{\delta }& 1\stackrel{~}{G}\mathrm{sin}\varphi \end{array}\right).$$
(84)
Here $`\stackrel{~}{G}G\eta \stackrel{~}{\gamma }`$ and
$$𝐮=\frac{\mathrm{\Gamma }_{\mathrm{eff}}}{4}(2N+1,2N+1+\frac{\stackrel{~}{\gamma }}{2}G^2,0)^T.$$
(85)
The steady state results are obtained by setting $`\dot{𝐲}(t)=0`$, which yields
$$𝐲_{ss}=M^1𝐮$$
(86)
and we can calculate Eqs. (58) and (65) with $`n=y_1+y_21/2`$. |
warning/0506/math-ph0506062.html | ar5iv | text | # Coalgebras and quantization
## 1 The Hopf algebra of normal products
Taking the example of a scalar field, we start from the co-algebra $`C`$ generated as a vector space by the Wick powers $`\phi ^n(x)`$, where $`n`$ is a nonnegative integer ($`x`$ is a point of $`^d`$). The coproduct of $`C`$ is given by
$`\mathrm{\Delta }_C\phi ^n(x)`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n}{k}}\right)\phi ^k(x)\phi ^{nk}(x),`$
its co-unit is $`\epsilon _C(\phi ^n(x))=\delta _{n,0}`$.
The coalgebra structure of $`C`$ enables us to define a commutative and cocommutative bialgebra $`B`$ which is the symmetric algebra $`S(C)`$ as an algebra, equipped with the coproduct $`\mathrm{\Delta }`$ defined by $`\mathrm{\Delta }u=\mathrm{\Delta }_Cu`$ if $`uS^1(C)=C`$ and extended to $`B`$ by algebra morphism (i.e. $`\mathrm{\Delta }(uv)=\mathrm{\Delta }u\mathrm{\Delta }v`$ for any $`u`$ and $`v`$ in $`B`$). The co-unit $`\epsilon `$ of $`B`$ is defined similarly by $`\epsilon (u)=\epsilon _C(u)`$ if $`uS^1(C)=C`$ and extended to $`B`$ by algebra morphism (i.e. $`\epsilon (uv)=\epsilon (u)\epsilon (v)`$ for any $`u`$ and $`v`$ in $`B`$). This product is called the normal product or Wick product in quantum field theory. The coproduct of basis elements of $`B`$ is
$`\mathrm{\Delta }\phi ^{n_1}(x_1)\mathrm{}\phi ^{n_p}(x_p)`$ $`=`$ $`{\displaystyle \underset{k_1=0}{\overset{n_1}{}}}\mathrm{}{\displaystyle \underset{k_p=0}{\overset{n_p}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n_1}{k_1}}\right)\mathrm{}\left({\displaystyle \genfrac{}{}{0pt}{}{n_p}{k_p}}\right)\phi ^{k_1}(x_1)\mathrm{}\phi ^{k_p}(x_p)\phi ^{n_1k_1}(x_1)\mathrm{}\phi ^{n_pk_p}(x_p).`$
Note that the co-unit is the vacuum expectation value : $`\epsilon (u)=0|u|0`$.
The bialgebra $`B`$ is transformed into a connected Hopf algebra $`H`$ through a quotient by the ideal (and co-ideal) generated by elements of the form $`a\epsilon _C(a)1`$, where $`aC`$ and 1 is the unit of $`S(C)`$.
## 2 Quantization
Quantization is now achieved by defining a co-quasi-triangular structure (co-QTS) on $`H`$. In $`H`$, a co-QTS $``$ is entirely determined by the value of $`(a,b)`$ for $`a`$ and $`b`$ in $`C`$. Two co-QTS are used: the operator co-QTS determined by $`(\phi ^m(x),\phi ^n(y))=\delta _{m,n}n!D_+(xy)^n,`$ where $`D_+(x)`$ is the free Wightman function and the chronological co-QTS determined by $`(\phi ^m(x),\phi ^n(y))=\delta _{m,n}n!D(xy)^n,`$ where $`D(x)`$ is a regularized version of the Feynman propagator, or Green function of the free field (see quantum field theory).
A co-quasi-triangular structure generates a twisted product defined by $`uv=(u{}_{\left(1\right)}{}^{},v{}_{\left(1\right)}{}^{})u{}_{\left(2\right)}{}^{}v{}_{\left(2\right)}{}^{}.`$ This expression is called Wick’s theorem in quantum field theory. When the co-QTS is determined by the Wightman function, the twisted product is equivalent to the star product of deformation quantization . When the co-QTS is determined by the Feynman propagator, the twisted product is commutative. It is called the chronological or time-ordered product or $`T`$-product and it plays a basic role in the perturbation theory of quantum fields. More generally, for $`a_1,\mathrm{},a_p`$ in $`C`$, the chronological product is $`T(a_1\mathrm{}a_p)=a_1\mathrm{}a_p`$. For any $`uH`$, $`T(u)=t(u{}_{\left(1\right)}{}^{})u_{\left(2\right)}`$, where $`t(u)=\epsilon (T(u))`$ is given by
$`t(\varphi ^{n_1}(x_1)\mathrm{}\varphi ^{n_p}(x_p))`$ $`=`$ $`n_1!\mathrm{}n_p!{\displaystyle \underset{M}{}}{\displaystyle \underset{i=1}{\overset{p1}{}}}{\displaystyle \underset{j=i+1}{\overset{p}{}}}{\displaystyle \frac{D(x_i,x_j)^{m_{ij}}}{m_{ij}!}},`$ (1)
where the sum is over all symmetric $`p\times p`$ matrices $`M`$ of nonnegative integers $`m_{ij}`$ such that $`_{j=1}^pm_{ij}=n_j`$ and $`m_{ii}=0`$ for all $`i`$. Each matrix $`M`$ is the adjacency matrix of a graph which is called a Feynman graph.
## 3 A second co-algebraic structure
There is a second co-product on $`H`$ defined by $`\mathrm{\Delta }^{}u=u1+1u`$ if $`uC`$, extended to $`H`$ by algebra morphism. For example
$`\mathrm{\Delta }^{}\phi ^{n_1}(x_1)\phi ^{n_2}(x_2)`$ $`=`$ $`\phi ^{n_1}(x_1)\phi ^{n_2}(x_2)1+1\phi ^{n_1}(x_1)\phi ^{n_2}(x_2)+\phi ^{n_1}(x_1)\phi ^{n_2}(x_2)`$
$`+\phi ^{n_2}(x_2)\phi ^{n_1}(x_1).`$
More generally $`\mathrm{\Delta }^{}\phi ^{n_1}(x_1)\mathrm{}\phi ^{n_p}(x_p)`$ is determined by taking all subsets $`I`$ of $`\{1,\mathrm{},p\}`$ and defining
$`\mathrm{\Delta }^{}\phi ^{n_1}(x_1)\mathrm{}\phi ^{n_p}(x_p)`$ $`=`$ $`{\displaystyle \underset{I}{}}\left({\displaystyle \underset{iI}{}}\phi ^{n_i}(x_i)\right)\left({\displaystyle \underset{jI}{}}\phi ^{n_j}(x_j)\right),`$
with the convention that $`_{iI}\phi ^{n_i}(x_i)=1`$ if $`I=\mathrm{}`$. This second coproduct is very natural and it has been implicitly used for a long time in physics .
We denote by $`(H,\mathrm{\Delta }^{})`$ the co-algebra which is equal to $`H`$ as a vector space, with co-product $`\mathrm{\Delta }^{}`$ and co-unit $`\epsilon `$ (the co-unit of the Hopf algebra $`H`$). The co-algebra $`(H,\mathrm{\Delta }^{})`$ and the Hopf algebra $`H`$ have an important relation: $`(H,\mathrm{\Delta }^{})`$ is a co-module co-algebra over $`H`$ . In other words, $`(H,\mathrm{\Delta }^{})`$ is a co-module over $`H`$ (with the co-action $`\psi =\mathrm{\Delta }`$) satisfying the compatibility property
$`(\mathrm{\Delta }^{}\mathrm{Id})\psi `$ $`=`$ $`(\mathrm{Id}\mathrm{Id}\mu )(\mathrm{Id}\tau \mathrm{Id})(\psi \psi )\mathrm{\Delta }^{},`$
where $`\tau (uv)=vu`$ and $`\mu (uv)=uv`$.
The co-product $`\mathrm{\Delta }^{}`$ enables us to define the connected and the renormalized chronological products. The *reduced co-product* is $`\underset{¯}{\mathrm{\Delta }}^{}u=\mathrm{\Delta }^{}uu11u`$, its iteration is $`\underset{¯}{\mathrm{\Delta }}_{}^{}{}_{}{}^{(0)}=\mathrm{Id}`$, $`\underset{¯}{\mathrm{\Delta }}_{}^{}{}_{}{}^{(1)}=\underset{¯}{\mathrm{\Delta }}^{}`$, $`\underset{¯}{\mathrm{\Delta }}_{}^{}{}_{}{}^{(n+1)}=(\underset{¯}{\mathrm{\Delta }}^{}\mathrm{Id}^n)\underset{¯}{\mathrm{\Delta }}_{}^{}{}_{}{}^{(n)}`$, and its action on $`uH`$ is denoted by
$`\underset{¯}{\mathrm{\Delta }}_{}^{}{}_{}{}^{(n1)}u`$ $`=`$ $`{\displaystyle }u{}_{\left(\underset{¯}{1}^{}\right)}{}^{}\mathrm{}u{}_{\left(\underset{¯}{n}^{}\right)}{}^{}.`$
The connected chronological product is now defined by
$`T_c(u)`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^n}{n}}T(u{}_{\left(\underset{¯}{1}^{}\right)}{}^{})\mathrm{}T(u{}_{(\underset{¯}{n}^{})}{}^{}),`$
for $`u\mathrm{ker}\epsilon `$. Because $`(H,\mathrm{\Delta }^{})`$ is a co-module co-algebra over $`H`$, $`T_c(u)=t_c(u{}_{\left(1\right)}{}^{})u_{\left(2\right)}`$, where $`t_c(u)`$ is given by eq.(1) with a sum over adjacency matrices $`M`$ corresponding to connected graphs. The renormalized chronological product is defined by
$`T_R(u)`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n!}}T\left(𝒪(u{}_{\left(\underset{¯}{1}^{}\right)}{}^{})\mathrm{}𝒪(u{}_{(\underset{¯}{n}^{})}{}^{})\right),`$
for $`u\mathrm{ker}\epsilon `$, where $`𝒪`$ is a linear map from $`H`$ to $`C`$ called a *generalized vertex* . The renormalized chronological product implements the renormalization of quantum field theory.
Note that a similar construction is possible in a noncommutative context where the symmetric algebra $`S(C)`$ is replaced by the tensor algebra $`T(C)`$. The second co-product $`\mathrm{\Delta }^{}`$ is then the deconcatenation co-product. |
warning/0506/nucl-th0506070.html | ar5iv | text | # An exactly solvable limit of low energy QCD
## I Introduction
As significant advances have been made over the past few years in lattice gauge QCD lattice , progress in developing other, approximate solutions to low energy QCD has not been as impressive. Many models, such as the quark model, the bag model, the flux tube model, and many others have been utilized to capture selective feature of the theory. Traditional approaches to develop models with a more rigorous relation to QCD are based on the covariant representation of the theory. Recently, however, fixed gauge approaches have also been intensively pursued as they offer a bridge between QCD and the more traditional (non-relativistic) many-body problems in nuclear and condensed-matter physics. The disadvantage of a fixed gauge approach is, however, that it is not manifestly Lorentz invariant and thus there are fewer restrictions on the dynamical operators resulting in general in complicated Hamiltonians. Fortunately experience gained form studies of other many-body systems can be helpful in identifying approximation schemes relevant for studies of particular aspects of the dynamics. Attempts to enlighten the non-perturbative structure of QCD using many-body techniques have recently been undertaken by several authors adam-gluon ; adam-conf and even the confinement scenario has been realized within the Coulomb gaugeadam-conf .
Though much progress has been made, schematic models are still very useful to shed some light onto the non-perturbative structure of QCD. For example, in gluon99 the gluon sector was investigated restricting the quark-gluon dynamics to an effective Hamiltonian for a fixed number of modes. The gluon spectrum was adjusted to reproduce lattice gauge calculations bali ; peardon and several other states have been predicted and confirmed by lattice gauge calculations. In qcd-model1 a Lipkin type model was introduced where the fermion sector consisted of two levels, one at positive and the other at negative energy, and a coupling to a boson level, occupied by color-spin zero gluon pairs, was considered. Only meson states were described. In qcd-model2 the model was extended to include baryons. The nucleon resonances (especially the Roper resonance) and $`\mathrm{\Delta }`$ resonances were well described. One drawback of these model, however is that they are purely phenomenological and contain several parameters. Our long term goal is to investigate classes of schematic models which are derived from QCD. For example, one expects the high energy quark and gluon modes be largely irrelevant in determining the structure of the vacuum, and lowest excitations, e.g. the pion or the $`\rho `$ meson states, color confinement, etc. One advantage of such schematic models is that they are QCD and will depend on only one or few well determined coupling constant(s). Preferably, such models should allow analytic or alternatively nearly analytic solutions, the latter requiring at most a numerical diagonalization.
Since it is a good practice to start from the simple cases first, here we present exact solutions to a schematic model, which has all the most drastic approximations. The goal is to identify which of these play what role in the low energy dynamics. Hopefully by systematically relaxing these approximations an intuitive picture of QCD will emerge. The model presented here is derived from QCD, under the restriction to $`SU(2)`$ in color and flavor. Only quarks and antiquarks will be considered. The interaction via gluons will be simulated via an effective, interaction reduced to the lowest momentum modes. These modes will correspond to the quarks and antiquarks restricted to be in a single spacial orbital level ($`S`$-state). Apart from this, the model will depend only on one parameter related to the energy of the color excitations thus ultimately irrelevant.
The paper is organized as follows. In section II the derivation of the schematic model, starting from QCD, is given and the appropriate particle content is identified. In section III the eigenvalues and the basis states will be constructed and the physical states are investigated in section IV. In section V conclusions will be drawn and future developments discussed.
## II Derivation of the Model Hamiltonian
As discussed above, by choosing an appropriate gauge a set of degrees of freedom can be selected which appears most natural for description of certain features of QCD. In this case physical (gauge independent) quantities may be simpler to calculate when the ”correct” gauge is chosen. For example, to compute various deep inelastic amplitudes it is advantageous to formulate QCD in the light-cone gauge, while to compute low energy spectra Coulomb gauge seems to be the natural choice. The Coulomb gauge has been extensively studied in adam-conf ; TDLee . The Gauss’s law can be used to eliminate the longitudinal component of the electric field which leaves only the transverse gluons representing generalized coordinates and their conjugated momenta (given by the transverse electric fields). Schematically the Coulomb gauge Hamiltonian has the following structure adam-conf ; TDLee ,
$$H=K_q+K_g+V_{qqg}+V_{g^3}+V_{g^4}+V_C.$$
(1)
Here $`K_q`$ and $`K_g`$ are the kinetic energies of the quarks -antiquarks and gluons, respectively, and are given by the Dirac and Yang-Mills Hamiltonians. The next three terms have polynomial dependence on the canonical degrees of freedom and represent the local (anti) quark -gluon interaction, triple- and quartic- gluon coupling, respectively. Finally $`V_C`$ is the non-abelian generalization of the Coulomb potential. In an abelian case, $`V_C=\alpha 𝑑𝐱𝑑𝐲\rho (𝐱)|𝐱𝐲|^1\rho (𝐲)`$ represents the Coulomb energy between matter charges, which are described by the charge density $`\rho (𝐱)`$. For simplicity we have already dropped the Faddeev-Popov determinant, which as shown in Szczepaniak:2003ve can be accounted for by redefining the gluon wave functional.
In a non-abelian theory like QCD the Coulomb potential depends not only on the relative separation between charges but also on the distribution of the gauge fields around them,
$$V_C=g^2𝑑𝐱𝑑𝐲\rho ^a(𝐱)\left[\frac{1}{1\lambda ^{}}\frac{1}{^2}\frac{1}{1\lambda }\right]_{a𝐱;b𝐲}\rho ^b(𝐲).$$
(2)
The matrix elements of $`1\lambda `$ are given by
$$[1\lambda ]_{a𝐱;b𝐲}=\delta _{ab}\delta (𝐱𝐲)gf_{acb}𝐀^c(𝐱)\delta (𝐱𝐲),$$
(3)
and the color-charge density is given by $`\rho ^a(x)=\psi (𝐱)T^a\psi (𝐱)f_{abc}𝐀^b(𝐱)𝐄^c(𝐱)`$. $`𝐀^a`$, $`𝐄^a`$ represent $`a=1,\mathrm{}N_c^21`$ transverse gluon coordinates and conjugate momenta, respectively and $`T^a=T_{ij}^a,i,j=1\mathrm{}N_C`$, and $`f_{abc}`$ are the generators of the fundamental and adjoint representations of the color $`SU(N_c)`$ group. Thus, unlike QED, in QCD to define the potential between a state containing matter (quark, antiquark) sources it is necessary to know the gluon wave functional of the state. It was shown in var-conf using a variational ansatz for the gluon wave functional of the vacuum that $`V_C`$ leads to a confining interaction between matter sources. Such an attractive interaction destabilizes the vacuum and leads to formation of quark-antiquark condensates and chiral symmetry breaking. The underlying mechanism is analogous to BCS superconductivity.
In the following we want to investigate the minimal requirements, e.g. the minimal number of degrees of freedom in a schematic model which yields the pattern of chiral symmetry breaking consistent with that expected in QCD. Since the necessary condition for the condensate is existence of an attractive interaction we remove the gluon degrees of freedom (for example by fixing the gluon wave functional) and approximating the Coulomb kernel by a contact potential between quark charge densities. Under such approximation the Coulomb gauge Hamiltonian of Eq. (1) reduces to,
$`H`$ $`=`$ $`{\displaystyle 𝑑𝐱\psi ^{}(𝐱)\left[i\stackrel{}{\alpha }\stackrel{}{}+\beta m_0\right]\psi (𝐱)}`$ (4)
$`+`$ $`g{\displaystyle 𝑑𝐱\rho ^a(𝐱)\rho ^a(𝐱)},`$
with the color charge density originating from quarks only, $`\rho ^a(𝐱)=\psi ^{}(𝐱)T^a\psi (𝐱)`$, and the coupling $`g`$ which has mass dimension $`2`$ will be determined later. The quark fields $`\psi (𝐱)`$ represent $`N_C\times N_f`$ degrees of freedom. The generators of the flavor axial rotations are
$$Q_5^\alpha =𝑑𝐱\psi ^{}(𝐱)\gamma _5T^\alpha \psi (𝐱),$$
(5)
with $`T^\alpha `$ being the generators of flavor, $`SU(N_f)`$. In the limit of vanishing quark mass, $`m_0=0`$, the Hamiltonian is invariant under flavor-axial rotations,
$$\underset{m_0=0}{lim}[Q_5^\alpha ,H]=0,$$
(6)
while for a finite bare mass
$$[Q_5^\alpha ,H]=2m_0P_5^\alpha ,$$
(7)
with
$$P_5^\alpha =𝑑𝐱\psi ^{}(𝐱)\gamma ^0\gamma _5T^\alpha \psi (𝐱).$$
(8)
To obtain the particle content of the spectrum of this Hamiltonian we first rewrite it in a basis of massive quarks and anti-quarks defined by the operators $`b(cf\lambda 𝐤)`$ and $`d(cf\lambda 𝐤)`$, respectively with $`c,f,\lambda ,𝐤`$ referring to color, flavor, spin component and momentum, and related to the fields in the standard way
$`\psi (𝐱)`$ $`=`$ $`{\displaystyle \underset{cf\lambda =\pm 1/2}{}}{\displaystyle }{\displaystyle \frac{d𝐤}{(2\pi )^3}}e^{i𝐱𝐤}[u(\lambda ,𝐤)b(cf\lambda 𝐤)`$ (9)
$`+`$ $`v(\lambda ,𝐤)d^{}(cf\lambda 𝐤)].`$
Here $`u`$ and $`v`$ are the eigenstates of the free Dirac Hamiltonian describing a fermion of mass $`m`$, which is not yet specified but is anticipated to be the constituent quark mass. In terms of these quark operators the Hamiltonian is given by,
$$H=H_q+H_{q\overline{q}}+V.$$
(10)
Here $`H_q`$ contains operators proportional to $`b^{}b`$ and $`d^{}d`$, $`H_{q\overline{q}}`$ contains pair creation and annihilation operators proportional to $`b^{}d^{}`$ and $`db`$, and $`V`$ contains normal-ordered four-fermion operators. Since we are interested in studying the low energy phenomena we make the following simplification. First we confine quarks to a finite box of volume $`𝒱`$. The momentum states become discrete, with $`𝐤𝐧`$ and $`𝐤=2\pi 𝐧/𝒱^{1/3}`$, so that
$$\frac{d𝐤}{(2\pi )^3}\frac{1}{𝒱}\underset{𝐧}{}.$$
(11)
In the finite volume it is also useful to rescale the particle operators,
$$b(cf\lambda 𝐤)\stackrel{~}{b}(cf\lambda 𝐧),b(cf\lambda 𝐤)=𝒱^{1/2}\stackrel{~}{b}(cf\lambda 𝐤),$$
(12)
and the same of the anti-quark operator $`d`$. The new operator are dimensionless and satisfy
$$\{\stackrel{~}{b}(cf\lambda 𝐧),\stackrel{~}{b}^{}(c^{}f^{}\lambda ^{}𝐧^{})\}=\delta _{cc^{}}\delta _{ff^{}}\delta _{\lambda \lambda ^{}}\delta _{\mathrm{𝐧𝐧}^{}}.$$
(13)
In the following we will rename $`\stackrel{~}{b},\stackrel{~}{d}`$ back as $`b`$ and $`d`$, respectively. The final approximation is to retain only the lowest momentum states, e.g. $`𝐧=0`$. Thus from now on we drop the momentum index on the quark operators. The next level of approximations would include the $`P`$\- and higher waves. Within this approximation the Hamiltonian becomes,
$`H`$ $`=`$ $`{\displaystyle \underset{1}{}}(+m_0)b_1^{}b_1+{\displaystyle \underset{1}{}}(+m_0)d_1^{}d_1`$ (14)
$``$ $`{\displaystyle \underset{1234}{}}V_{qq}(1234)b_1^{}b_2^{}b_3b_4{\displaystyle \underset{1234}{}}V_{\overline{q}\overline{q}}(1234)d_1^{}d_2^{}d_3d_4`$
$``$ $`2{\displaystyle \underset{1234}{}}V_{q\overline{q}}(1234)b_1^{}d_2^{}d_3b_4.`$
with
$$=\frac{gC_F}{𝒱}\underset{𝐧}{\overset{𝐧_{max}}{}}\delta _{𝐧0}=\frac{gC_F}{𝒱},$$
(15)
and
$`V_{qq}(1234)`$ $`=`$ $`{\displaystyle \frac{g}{𝒱}}T_{c_1c_3}^aT_{c_2c_4}^a\left[\delta _{f_1f_3}\delta _{\lambda _1\lambda _3}\right]\left[\delta _{f_2f_4}\delta _{\lambda _2\lambda _4}\right]`$
$`V_{\overline{q}\overline{q}}(1234)`$ $`=`$ $`{\displaystyle \frac{g}{𝒱}}T_{c_3c_1}^aT_{c_4c_2}^a\left[\delta _{f_1f_3}\delta _{\lambda _1\lambda _3}\right]\left[\delta _{f_2f_4}\delta _{\lambda _2\lambda _4}\right]`$
$`W_{q\overline{q}}(1234)`$ $`=`$ $`{\displaystyle \frac{g}{𝒱}}T_{c_1c_4}^aT_{c_3c_2}^a\left[\delta _{f_1f_4}\delta _{\lambda _1\lambda _4}\right]\left[\delta _{f_2f_3}\delta _{\lambda _2\lambda _3}\right].`$
Here $`1=(f_1,c_1,\lambda _1)`$ etc. denote all remaining (discrete) quantum numbers of the particle labeled by $`1`$; $`c_1`$ denotes color, $`f_1`$ flavor and $`\lambda _1`$ spin projection. It is worth noting at this point that with $`S`$-orbitals only the pair creation part of the Hamiltonian vanishes. Scalar quark-antiquark pairs have quarks in relative spin-one coupled to one unit of orbital angular momentum which vanishes for $`S`$-waves. Within these approximations the flavor axial charge generators become
$$Q_5^\alpha =\underset{12}{}\left(b_1^{}Q_{12}^\alpha d_2^{}+d_1Q_{12}^\alpha b_2\right),$$
(17)
with
$$Q_{12}^\alpha =T_{f_1f_2}^\alpha \delta _{c_1c_2}\delta _{\lambda _1\lambda _2},$$
(18)
and the pseudo-scalar charges $`P_5^\alpha `$ become
$$P_5^\alpha =\underset{12}{}\left(b_1^{}Q_{12}^\alpha d_2^{}d_1Q_{12}^\alpha b_2\right).$$
(19)
We also note that
$$[P_5^a,H]=2m_0Q_5^a,$$
(20)
is still satisfied. For completeness, the vector flavor charges $`V^\alpha `$,
$$V^\alpha =𝑑𝐱\psi ^{}(𝐱)T^\alpha \psi (𝐱),$$
(21)
become
$$V^\alpha =\underset{12}{}\left(b_1^{}V_{12}^\alpha b_2d_1^{}V_{12}^\alpha d_2\right),$$
(22)
with
$$V_{12}^\alpha =T_{f_1f_2}^\alpha \delta _{c_1c_2}\delta _{\lambda _1\lambda _2}.$$
(23)
The Hamiltonian contains four parts. A non-interaction part, quark-quark, and antiquark-antiquark potentials and a quark-antiquark potential. We recall some basic properties of the particle operators. In the following we concentrate on the case of two colors and two flavors. Generalization to arbitrary $`N_C`$ and $`N_f`$ is straightforward. The creation and annihilation operators carry color, $`c`$, flavor, $`f`$ and spin, $`\lambda `$ indices and these all range from $`\frac{1}{2}`$ to $`+\frac{1}{2}`$. We distinguish now between co- and contravariant indices in order to denote the different transformation properties of the fermion creation and annihilation operators. We denote the creation and annihilation operators for quarks by $`b_\alpha ^{}`$ and $`b^\alpha `$, respectively, where $`\alpha `$ is a shorthand notation for $`(cf\lambda )`$. Subsequently representation of $`SU(2)`$-color, flavor, spin will be similarly denoted by three numbers $`(S_cS_fS)`$, where $`S_c`$ is the color angular momentum and similar for flavor, $`S_f`$, and spin, $`S`$. Similarly for the antiquark operators we have $`d^\alpha `$ for the creation and $`d_\alpha `$ for the annihilation operators. The anticommutation relations are now given by $`\{b^\beta ,b_\alpha ^{}\}=\{d_\beta ,d^\alpha \}=\delta _\beta ^\alpha `$. The indices are lowered according to the following convention. If $`a^\alpha `$ denotes any of the four operators ($`b^{}`$, $`b`$, $`d`$, or $`d^{}`$) with an upper index, lowering this index corresponds to,
$$a^{cf\lambda }=(1)^{\frac{1}{2}c}(1)^{\frac{1}{2}f}(1)^{\frac{1}{2}\lambda }a_{cf\lambda }.$$
(24)
We can now rewrite the Hamiltonian. The non-interacting part is trivial and given by,
$`H_q`$ $`=`$ $`(+m_0)(\widehat{n}_q+\widehat{n}_{\overline{q}}).`$ (25)
with $`\widehat{n}_q=b_\alpha ^{}b^\alpha `$ and $`\widehat{n}_{\overline{q}}=d^\alpha d_\alpha `$ being the quark and antiquark number operators, respectively. The quark-quark interaction is given by
$`V_{qq}`$ $`=`$ $`{\displaystyle \underset{c^{}sf^{}s\lambda ^{}s}{}}{\displaystyle \frac{g}{𝒱}}T_{c_1c_3}^aT_{c_2c_4}^a\left[\delta _{f_1f_3}\right]\left[\delta _{f_2f_4}\right]`$ (26)
$`\times `$ $`b_{c_1f_1\lambda _1}^{}b_{c_2f_2\lambda _2}^{}b^{c_3f_3\lambda _3}b^{c_4f_4\lambda _4}.`$
Using
$`T_{c_1c_3}^aT_{c_2c_4}^a={\displaystyle \frac{1}{2}}\left(\delta _{c_1c_4}\delta _{c_3c_2}{\displaystyle \frac{1}{2}}\delta _{c_1c_3}\delta _{c_2c_4}\right).`$ (27)
and joining operators with common indices through the anticommutation relations, we obtain in an intermediate step
$`V_{qq}={\displaystyle \frac{3g}{4𝒱}}\widehat{n}_q{\displaystyle \frac{g}{4𝒱}}\widehat{n}_q^2`$
$`+`$ $`{\displaystyle \frac{g}{2𝒱}}{\displaystyle \underset{c_1c_2}{}}\left({\displaystyle \underset{f_1\lambda _1}{}}b_{c_1f_1\lambda _1}^{}b^{c_2f_1\lambda _1}\right)\left({\displaystyle \underset{f_2\lambda _2}{}}b_{c_2f_2\lambda _2}^{}b^{c_1f_2\lambda _2}\right).`$
Finally using Eq. (24) and coupling to definite color, flavor and spin, we arrive at
$`V_{qq}`$ $`=`$ $`{\displaystyle \frac{3}{4}}{\displaystyle \frac{g}{𝒱}}\widehat{n}_q{\displaystyle \frac{2g}{𝒱}}\sqrt{3}\left[\left[b^{}b\right]^{100}\left[b^{}b\right]^{100}\right]_{000}^{000},`$
where $`\left[A^{\mathrm{\Gamma }_1}B^{\mathrm{\Gamma }_2}\right]_\mu ^\mathrm{\Gamma }`$ with $`\mathrm{\Gamma }=S_cS_fS`$ and $`\mu =cf\lambda `$ denotes the coupling of $`A`$ and $`B`$ in color, flavor and spin,
$$\left[A^{\mathrm{\Gamma }_1}B^{\mathrm{\Gamma }_2}\right]_\mu ^\mathrm{\Gamma }=\underset{\mu _1\mu _2}{}\mathrm{\Gamma }_1\mu _1,\mathrm{\Gamma }_2\mu _2|\mathrm{\Gamma }\mu A_{\mu _1}^{\mathrm{\Gamma }_1}B_{\mu _2}^{\mathrm{\Gamma }_2},$$
(30)
and $`\mathrm{\Gamma }_1\mu _1,\mathrm{\Gamma }_2\mu _2|\mathrm{\Gamma }\mu `$ is the product of three Clebsch-Gordan coefficients in color, flavor and spin.
Note, that the quadratic dependence on the quark number operator is canceled and only the linear dependence remains. The last term in Eq. (LABEL:hqq3) represents the color angular momentum squared, whose components in spherical basis are given by,
$`S_{q,m}^c`$ $`=`$ $`\sqrt{2}\left[b^{}b\right]_{m00}^{100}`$
$`S_{\overline{q},m}^c`$ $`=`$ $`\sqrt{2}\left[d^{}d\right]_{m00}^{100},`$ (31)
for the quark and antiquark part, respectively. With this, the final form of $`V_{qq}`$ is
$`V_{qq}`$ $`=`$ $`{\displaystyle \frac{3}{4}}{\displaystyle \frac{g}{𝒱}}\widehat{n}_q+{\displaystyle \frac{g}{𝒱}}\left(S_q^cS_q^c\right),`$ (32)
and we used $`\left[S_q^cS_q^c\right]=\left(𝐒_q^c𝐒_q^c\right)`$. In a complete analogy one can show that the antiquark-antiquark part is found to be,
$`V_{\overline{q}\overline{q}}`$ $`=`$ $`{\displaystyle \frac{3}{4}}{\displaystyle \frac{g}{𝒱}}\widehat{n}_{\overline{q}}+{\displaystyle \frac{g}{𝒱}}\left(𝐒_{\overline{q}}^c𝐒_{\overline{q}}^c\right).`$ (33)
And finally for the quark-antiquark interaction is given by,
$`V_{q\overline{q}}`$ $`=`$ $`{\displaystyle \frac{2g}{𝒱}}\left(𝐒_q^c𝐒_{\overline{q}}^c\right),`$ (34)
Summing all terms leaves us with a surprisingly simple Hamiltonian whose interactions are easily identified,
$`H`$ $`=`$ $`\left(+m_0{\displaystyle \frac{3}{4}}{\displaystyle \frac{g}{𝒱}}\right)\left(\widehat{n}_q+\widehat{n}_{\overline{q}}\right)+{\displaystyle \frac{g}{𝒱}}𝐒_c^2,`$ (35)
where $`𝐒_c^2=\left(𝐒_q^c+𝐒_{\overline{q}}^c\right)^2`$ is the total color angular momentum squared.
## III The Spectrum
The basis used to diagonalize $`H`$ is determined by the number of degrees of freedom each quark (antiquark) carries. There are eight degrees of freedom: two spin times two flavor and times two color components. The Fock space is thus finite and contains maximally eight quarks and eight antiquarks. The group structure for each sector is given by hamermesh
$`U(8)`$ $`U_c(2)`$ $`U_{fS}(4)`$
$`\left[1^{n_q}\right]`$ $`\left[h_1h_2\right]`$ $`\left[2^{h_2}1^{h_1h_2}\right]`$
$`U_{fS}(4)`$ $`U_f(2)`$ $`U_S(2)`$
$`\left[2^{h_2}1^{h_1h_2}\right]`$ $`S_f`$ $`S,`$ (36)
The notation $`\left[p_1p_2\mathrm{}p_n\right]`$ refers to the Young diagrams hamermesh , which describes the symmetry under permutation of a given irreducible representation (irrep) of a unitary group. In Eq. (36) we have $`h_1+h_2=n_q`$ and the reduction of the flavor-spin group $`U_{fS}(4)`$ is given in hamermesh . In Table 1 we give a list of the color-flavor-spin content as a result of Eq. (36).
We now consider meson-like excitations, i.e. the Fock sector with equal number of quarks and antiquarks. For this case, the basis can be labeled by the following set of quantum numbers,
$`n_{\overline{q}}=n_q;(S_{\overline{q}}^c,S_q^c)S^cm^c;(S_{\overline{q}}^f,S_q^f)S^fm^f;(S_{\overline{q}},S_q)Sm,`$
(37)
where $`m^c`$, $`m^f`$ and $`m`$ refer to the magnetic color, flavor and spin projection. The eigenvalue of the Hamiltonian with respect to such states is given by
$`E`$ $`=`$ $`\left(+m_0{\displaystyle \frac{3}{4}}{\displaystyle \frac{g}{𝒱}}\right)(n_q+n_{\overline{q}})+{\displaystyle \frac{g}{𝒱}}S_c(S_c+1).`$
For physical states with no net color only the first term contributes, and using Eq. (15) we find $`E=m_0`$ and the spectrum is degenerate with respect to flavor and spin. Color excitations are separated by a finite energy gap which is an artifact of the contact approximation for the quark interactions. In full QCD the splitting is expected to be infinite as the potential between the quarks grows with the relative separation. Nevertheless, one can investigate the structure of colored excitations in the model, which might play a role in models like the quark-gluon glass condensate glass-qg , important at high densities.
The energy solutions are simple and degenerate for all states with the same color. At a first glance one might think that the physical states should be a certain sum over all basis states (Eq. (37)) with the same color. As a consequence, in our schematic model one would look for adequate superposition of the degenerate states in order to construct, e.g., the physical vacuum state. One criterion used can be to reproduce the quark condensate (see next section). However, as we will show further below, arguments of coninuity require that the lowest state has to be the vacuum state $`|0`$. To get more physical solutions, it will be necessary to introduce an interaction which lifts the large degeneracy of the Hamiltonian.
## IV Physical states and the chiral limit
In the chiral limit, $`m_0=0`$ all color singlet states have zero energy. The vacuum state should be identified as a state with all scalar quantum numbers. Since the single quark-antiquark pair in the $`S`$-wave has pseudoscalar quantum numbers, the vacuum will be given by a superposition of states with an even number of quark-antiquark pairs with total color, flavor and spin zero. The most general (unnormalized) vacuum state can be schematically written as
$$|𝐳=|0+\underset{n=1}{\overset{4}{}}z_n\left(b^{}b^{}d^{}d^{}\right)^n|0.$$
(39)
Since in the chiral limit all $`J^{PC}=0^{++}`$ states are degenerate in this model we cannot distinguish between the true vacuum and, for example the $`\sigma `$ meson. Thus we take for the vacuum a state given by the sum of the perturbative vacuum $`|0`$ and the state with the lowest number (two) of the quark-antiquark pairs coupled to definite color, flavor and spin. Each pair can be written in the following equivalent form
$`|S_cS_fS`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left[\left[b^{}b^{}\right]^{S_cS_fS}\left[d^{}d^{}\right]^{S_cS_fS}\right]_{000}^{000}|0.`$
Because the two-quark state has to be antisymmetric (the same for the antiquarks) the only allowed color, flavor and spin values are $`(S_cS_fS)=(000),(110),(101)`$ and $`(011)`$. The coupling of first two quarks and then two antiquarks to a total color, flavor and spin zero can be re-expressed easily in terms of the coupling of two quark-antiquark pairs as follows,
$`\left[\left[b^{}b^{}\right]^{S_cS_fS}\left[d^{}d^{}\right]^{S_cS_fS}\right]_{000}^{000}={\displaystyle \underset{S_c^{}S_f^{}S^{}}{}}`$ (50)
$`\left\{\begin{array}{ccc}\frac{1}{2}& \frac{1}{2}& S_c\\ \frac{1}{2}& \frac{1}{2}& S_c\\ S_c^{}& S_c^{}& 0\end{array}\right\}\left\{\begin{array}{ccc}\frac{1}{2}& \frac{1}{2}& S_f\\ \frac{1}{2}& \frac{1}{2}& S_f\\ S_f^{}& S_f^{}& 0\end{array}\right\}\left\{\begin{array}{ccc}\frac{1}{2}& \frac{1}{2}& S\\ \frac{1}{2}& \frac{1}{2}& S\\ S^{}& S^{}& 0\end{array}\right\}`$
$`\times \left[\left[b^{}d^{}\right]^{S_c^{}S_f^{}S^{}}\left[b^{}d^{}\right]^{S_c^{}S_f^{}S^{}}\right]_{000}^{000}`$ , (51)
where the symbols $`\left\{\mathrm{}\right\}`$ refer to the usual 9-j symbols edmonds . For the vacuum we thus take the normalized state in the form,
$`|z_0z_1={\displaystyle \frac{1}{\sqrt{1+2\rho ^2}}}\left(|0+{\displaystyle \underset{S_c}{}}z_{S_c}{\displaystyle \underset{S_fS}{}}|S_cS_fS\right),`$ (52)
with $`|S_cS_fS`$ given in Eq. (LABEL:two-qqbar). Here we assumed that due to the degeneracy of the states with the same color, there is no dependence of the trial state parameters $`z_{S_c},S_c=0,1`$ on flavor and spin. In general the $`z`$-values are complex and can be written as $`z_{S_c}=\rho _{S_c}e^{i\varphi _{S_c}}`$, with $`\rho _0=\rho cos(\varphi )`$ and $`\rho _1=\rho sin(\varphi )`$, and $`\rho =|z_0|^2+|z_1|^2`$ being the total radius. In such a vacuum expectation values of $`\widehat{n}_q`$ and $`\widehat{n}_{\overline{q}}`$, which determine the quark condensate, are given by
$`z_0z_1|\widehat{n}_q|z_0z_1=z_0z_1|\widehat{n}_{\overline{q}}|z_0z_1`$ $`=`$ $`{\displaystyle \frac{4\rho ^2}{1+2\rho ^2}}.`$ (53)
In the limit $`z_{S_c}=0`$, $`\widehat{n}_q=\widehat{n}_{\overline{q}}=0`$ as expected for the perturbative vacuum. For large values of $`\rho `$, each expectation value approaches 2, as it has to be, because then the main contribution comes from the two quark-antiquark pairs. Using Eq. (53) it is possible to define the collective potential as the expectation value of the Hamiltonian, the result is
$$V(z_0,z_1)=z_0,z_1|H|z_0,z_1=\left(\frac{3}{4}\frac{g}{𝒱}\right)\frac{8\rho ^2}{1+2\rho ^2},$$
(54)
which corresponds near $`\rho =0`$ to a harmonic oscillator and the potential saturates for $`\rho \mathrm{}`$ at $`4\left(\frac{3}{4}\frac{g}{V}\right)`$. The use of such trial states played primordial role in nuclear physics to help understand the structure of a complicated many body problem iba and might be here also of great value when a more sofisticated Hamiltonian is used. Because, as we showed above, the factor which contains $``$ is zero one obtains a flat potential which reflects the complete degeneracy of color zero states. As already mentioned, the $`z`$ parameters are complex, but the expectation value above depends only on the total radius $`\rho `$. This implies that equipotential lines flow along constant $`\rho `$ with arbitrary angles $`\varphi _{S_c}`$ and $`\sqrt{\rho _0^2+\rho _1^2}=\rho `$.
To further determine parameters of the vacuum one can consider the quark condensate $`\overline{q}q=\overline{u}u=\overline{d}d=\overline{\psi }(0)\psi (0)/2(225\text{MeV})^31\text{fm}^3`$ reinders ,
$$\overline{q}q=\frac{1}{𝒱}\left[N_CN_S\frac{1}{2}\left(\widehat{n}_q+\widehat{n}_{\overline{q}}\right)\right]=\frac{4}{𝒱}\frac{1+\rho ^2}{1+2\rho ^2}.$$
(55)
One might be tempted to use this to determine $`\rho `$ for given volume (e.g. taking as a volume of sphere of radius of $`0.8\text{ fm}`$ would yield $`\rho =0.67`$). This is however not correct since there are further constraints from the spontaneous realization of chiral symmetry breaking. Away from the chiral limit, $`m_00`$ each additional quark-antiquark pair raises the energy by $`2m_0`$. Thus, for $`m_00`$ the vacuum has to be given by the single state $`|0`$, so $`\rho =0`$. If there is no phase transition at $`m_0=0`$ then for all $`m_0`$ the vacuum should be identified with the $`|0`$ state and $`\rho =0`$. The quark condensate is then entirely determined by the volume and the total number of degrees of freedom,
$$𝒱=N_CN_S\overline{q}q^1=2.7\text{fm}^3$$
(56)
As expected for spontaneous breaking the generators of chiral symmetry, Eq. (17), which can be also written as
$`Q_f^5`$ $`=`$ $`{\displaystyle \frac{\sqrt{N_CN_fN_S}}{2}}\left(\left[b^{}d^{}\right]_{0f0}^{010}+\left[db\right]_{0f0}^{010}\right),`$ (57)
do not annihilate the vacuum, instead they mix the vacuum with the single pion state,
$$\pi ,f^{}|Q_f^5|0=\delta _{f^{}f}f_\pi m_\pi 𝒱$$
(58)
with $`f_\pi =93\text{ MeV}`$ being the pion decay constant. Chiral symmetry, and relativistic normalization of single particle states,
$$\pi ,f^{}|\pi f=2m_\pi 𝒱,$$
(59)
implies that in the chiral limit $`m_00`$, $`m_\pi =O(m_0^2)`$ and $`f_\pi =O(1)`$. Since pion has $`J^{PC}=0^+`$ quantum numbers and is generated by the axial rotation from the vacuum the most general (unnormalized) pion state is given by,
$$|\pi b^{}d^{}\left[|0+\underset{n=1,4}{}w_i\left(b^{}b^{}d^{}d^{}\right)^n|0\right]$$
(60)
Mixing with the vacuum through the axial charge, as given by Eq. (58), constraints the quark-antiquark component to
$`|\pi ,f`$ $`=`$ $`f_\pi m_\pi 𝒱{\displaystyle \frac{2}{\sqrt{N_CN_fN_S}}}[b^{}d^{}]_{0f0}^{010}(|0`$ (61)
$`+`$ $`{\displaystyle \underset{n=1}{\overset{4}{}}}w_i\left(b^{}b^{}d^{}d^{}\right)^n|0),`$
However, all states in the expansion in Eq. (61) are eigenstates of the Hamiltonian with increasing eigenvalues and physical state cannot be given such a linear combination. We thus conclude that the single pion state should be identified with the valence component alone,
$$|\pi ,f=f_\pi m_\pi 𝒱\frac{2}{\sqrt{N_CN_fN_S}}[b^{}d^{}]_{0f0}^{010}|0$$
(62)
With the pion mass related to the bare quark mass by $`m_\pi =2m_0`$. The normalization condition of Eq. (59) then leads to
$`f_\pi `$ $`=`$ $`\sqrt{{\displaystyle \frac{N_CN_fN_S}{2m_\pi 𝒱}}}=\sqrt{{\displaystyle \frac{N_f\overline{q}q}{2m_\pi }}}=200\sqrt{N_f}\text{ MeV}.`$
The identification of other physical states with the spectrum given in Table 1 is now straightforward. Since the number of quarks and antiquarks are well defined and each additional (anti)quark raises energy by $`m_0`$ the spectrum of single meson and baryon states would correspond to stated with a single $`q\overline{q}`$ pair and three quarks respectively. States with other numbers of quarks or antiquark should be identified with multi-particle states e.g. $`qqqq\overline{q}`$ with a meson-baryon state. Colored states are split from the physical color singlet states by $`gS_c(S_c+1)/𝒱`$ where $`S_c`$ is a half-integer or integer total color for an odd or even number of quarks and antiquarks in the state, respectively, and $`g`$ is the effective strength of the colored interactions. We thus see that it is now possible to take the limit $`g\mathrm{}`$ which is expected for the zero-mode component of a confining interactions without affecting the physical spectrum.
## V Summary
Models play an important role in understanding complicated dynamical structures. Our goal here was not to build the most sophisticated model of low energy QCD, but on the contrary to identify the most basis starting point for such an endeavor. Starting from the underlying QCD interactions in the Coulomb gauge we have defined an approximations scheme which gave us a model for the interactions of the quark zero modes. The model is exactly solvable and physical states can be identified with help of the symmetry patterns observed in the physical spectrum. In particular spontaneous breaking of chiral symmetry enables to identify the vacuum state and the single pion state and conservation of the particle number by our model Hamiltonian then leads to mapping between the representations of the underlying $`U(N_C\times N_f\times N_S)`$ symmetry and the physical states. We worked with the $`N_C=2`$ number of colors, but extension to $`N_C=3`$ is straightforward since the coupling and recoupling methods in $`SU(2)`$ can be readily extended to $`SU(3)`$ (see, for example the appendix of Ref. jutta ). For the basis, instead of $`U(8)`$ we would start from $`U(12)`$ if flavor is still $`SU(2)`$ or $`U(18)`$ if flavor is also $`SU(3)`$. The reductions are known (see Ref. qcd-model1 ; qcd-model2 ; ramon ). In the model we find splitting between physical states to be proportional to the total bare mass of the quarks and anti-quars independently of the strength of the color or confining interaction. The color interaction is then responsible for lifting the energy of the color non-singlet states. Even though the model respects the pattern of chiral symmetry breaking the chiral behavior of the physical constants, e.g. the pion mass and the pion decay constant is not as expected. The pion mass turns out to be a linear and not quadratic function of the symmetry breaking parameter, $`m_0`$, and the decay constant depends on $`m_0`$. This is an expected behavior for large values of $`m_0`$, or the nonrelativistic quark model. It is not surprising that our schematic model away from the exact chiral limit of $`m_0=0`$ immediately follows the pattern of a heavy quark theory since the model conserves the quark number. This in turn is the consequence of reduction of the quark degrees of freedom. With the gauge degrees of freedom integrated out and quark Fock space reduced to the zero modes there are no pair production interactions in the Coulomb gauge. This suggests that by extending the Fock space to include a limited number of non-zero momentum modes and/or addinggluon degrees of freedom it may be possible to address the low energy phenomena in a model with a finite number of degrees of freedom.
## Acknowledgments
This work belongs to the DGAPA project IN119002. It was partially supported by the National Research Councils of Mexico (CONACYT) and the US Department of Energy grant DE-FG0287ER40365. P.O. Hess acknowledges the support of the Nuclear Theory Center at Indiana University where this work was initiated. |
warning/0506/hep-th0506217.html | ar5iv | text | # On Dual Formulation of Gravity. II. Metric and Affine Connection.
## Abstract
In this note we construct a dual formulation of gravity where the main dynamical object is affine connection. We start with the well known first order Palatini formulation but in (Anti) de Sitter space instead of flat Minkowski space as a background. The final result obtained by solving equations for the metric is the Lagrangian written by Eddington in his book in 1924. Also there is an interesting connection with attempts to construct gravitational analog of Born-Infeld electrodynamics.
In general, by dual formulation we mean any situation where the very same particle is described by different tensor fields. The most simple and straightforward way to obtain such dual formulation based on the use of first order ”parent” Lagrangians. As is well known in flat Minkowski space such dualization procedure leads to different results for massive and massless particles. At the same time in (Anti) de Sitter space-time gauge invariance requires introduction quadratic mass-like terms into the Lagrangians. As a result dualization for massless particles in (Anti) de Sitter spaces goes exactly in the same way as that for massive particles. As an example, we have recently shown that by using well known tetrad formalism it is possible to obtain dual formulation of gravity with the Lorentz connection being the main dynamical field while tetrad is just auxiliary fields which could be expressed in terms of Lorentz connection and its derivatives. But there exist another well known first order formalism for gravity usually called Palatini formalism, the main components being the metric and affine connection. Such formalism differs drastically from the tetrad one because affine connection is not a gauge invariant object (or, geometrically, it is not a covariant tensor) and does not have its own gauge invariance. In spite of this difference, as we are going to show in this note, it is also possible to apply the same dualization procedure to obtain a formulation of gravity where the main dynamical field is the affine connection. Rather naturally and at the same time surprisingly the final result is nothing else but the Lagrangian written by Eddington in 1924 !
Let us start with the first order Lagrangian describing free massless spin-2 particle in flat Minkowski space:
$$_0=h^{\mu \nu }(_\alpha \mathrm{\Gamma }_{\mu \nu }{}_{}{}^{\alpha }_\mu \mathrm{\Gamma }_\nu )+\eta ^{\mu \nu }(\mathrm{\Gamma }_{\mu \nu }{}_{}{}^{\alpha }\mathrm{\Gamma }_{\alpha }^{}\mathrm{\Gamma }_{\mu \alpha }{}_{}{}^{\beta }\mathrm{\Gamma }_{\nu \beta }^{}{}_{}{}^{\alpha })$$
(1)
Here $`h_{\mu \nu }`$ is symmetric second rank tensor while $`\mathrm{\Gamma }_{\mu \nu }^\alpha `$ is assumed to be symmetric on the lower pair of indices. We denote $`\mathrm{\Gamma }_\alpha =\mathrm{\Gamma }_{\alpha \beta }^\beta `$, $`\mathrm{\Gamma }^\alpha =\eta ^{\mu \nu }\mathrm{\Gamma }_{\mu \nu }^\alpha `$ (note, that $`\mathrm{\Gamma }_\alpha `$ and $`\mathrm{\Gamma }^\alpha `$ are in general different objects). This Lagrangian is invariant under the following local gauge transformations:
$$\delta h_{\mu \nu }=_\mu \xi _\nu +_\nu \xi _\mu \eta _{\mu \nu }(\xi ),\delta \mathrm{\Gamma }_{\mu \nu }{}_{}{}^{\alpha }=_\mu _\nu \xi ^\alpha $$
(2)
As is well known, if one solves the algebraic equation of motion for the $`\mathrm{\Gamma }`$ field and put the result back into the Lagrangian one obtains usual second order Lagrangian for the symmetric tensor $`h_{\mu \nu }`$. In order to have a possibility to construct dual formulation where the main dynamical object is $`\mathrm{\Gamma }`$ we move from the flat Minkowski space to (Anti) de Sitter space. Let $`\overline{g}_{\mu \nu }`$ be a metric for this space (it is not a dynamical quantity, just a background field here) and $`D_\mu `$ — derivatives covariant with respect to background connection winch is torsionless and metric compatible:
$$D_\alpha \overline{g}_{\mu \nu }=0,[D_\mu ,D_\nu ]v_\alpha =\overline{R}_{\mu \nu ,\alpha }{}_{}{}^{\beta }(\overline{g})v_\beta =\kappa (\overline{g}_{\mu \alpha }\delta _\nu {}_{}{}^{\beta }\delta _\mu {}_{}{}^{\beta }\overline{g}_{\nu \alpha }^{})v_\beta $$
(3)
where $`\kappa =2\mathrm{\Lambda }/(d1)(d2)`$. First of all we have to replace in the Lagrangian as well as in the gauge transformations the flat metric $`\eta _{\mu \nu }`$ by $`\overline{g}_{\mu \nu }`$ and partial derivatives $`_\mu `$ by covariant ones $`D_\mu `$:
$`_0`$ $`=`$ $`h^{\mu \nu }(D_\alpha \mathrm{\Gamma }_{\mu \nu }{}_{}{}^{\alpha }D_\mu \mathrm{\Gamma }_\nu )+\overline{g}^{\mu \nu }(\mathrm{\Gamma }_{\mu \nu }{}_{}{}^{\alpha }\mathrm{\Gamma }_{\alpha }^{}\mathrm{\Gamma }_{\mu \alpha }{}_{}{}^{\beta }\mathrm{\Gamma }_{\nu \beta }^{}{}_{}{}^{\alpha })`$
$`\delta h_{\mu \nu }`$ $`=`$ $`D_\mu \xi _\nu +D_\nu \xi _\mu \overline{g}_{\mu \nu }(D\xi ),\delta \mathrm{\Gamma }_{\mu \nu }{}_{}{}^{\alpha }={\displaystyle \frac{1}{2}}(D_\mu D_\nu +D_\nu D_\mu )\xi ^\alpha `$ (4)
Now (just because covariant derivatives do not commute) our Lagrangian is not invariant under the gauge transformations. Indeed, simple calculations give:
$$\delta _0=\kappa [(d2)\mathrm{\Gamma }_\mu \xi ^\mu \frac{d3}{2}\mathrm{\Gamma }^\mu \xi _\mu \frac{3d1}{2}h^{\mu \nu }D_\mu \xi _\nu +h(D\xi )]$$
But gauge invariance could be easily restored by adding terms quadratic in $`h_{\mu \nu }`$ field to the Lagrangian as well as appropriate corrections for the gauge transformations:
$`\mathrm{\Delta }_0`$ $`=`$ $`{\displaystyle \frac{\kappa (d1)}{2}}[h^{\mu \nu }h_{\mu \nu }{\displaystyle \frac{1}{d2}}h^2]`$
$`\delta ^{}\mathrm{\Gamma }_{\mu \nu }^\alpha `$ $`=`$ $`{\displaystyle \frac{\kappa }{2}}(\delta _\mu {}_{}{}^{\alpha }\xi _{\nu }^{}+\delta _\nu {}_{}{}^{\alpha }\xi _{\mu }^{})\kappa \overline{g}_{\mu \nu }\xi ^\alpha `$ (5)
Now one can easily solve the equations for the $`h_{\mu \nu }`$ field, which are also algebraic now, to obtain:
$$h_{\mu \nu }=\frac{1}{\kappa (d1)}[R_{\mu \nu }\frac{1}{2}\overline{g}_{\mu \nu }R]$$
(6)
where we introduced a symmetric second rank tensor (it is not a full Ricci tensor yet, only the first part of it):
$$R_{(\mu \nu )}=\frac{1}{2}(D_\mu \mathrm{\Gamma }_\nu +D_\nu \mathrm{\Gamma }_\mu )D_\alpha \mathrm{\Gamma }_{\mu \nu }^\alpha $$
(7)
Then if we put this expression back into the initial first order Lagrangian we obtain dual second order formulation for massless spin-2 particle in terms of $`\mathrm{\Gamma }`$ field:
$$_{II}=\frac{1}{\kappa (d1)}[R^{\mu \nu }R_{\mu \nu }\frac{1}{2}R^2]+\overline{g}^{\mu \nu }(\mathrm{\Gamma }_{\mu \nu }{}_{}{}^{\alpha }\mathrm{\Gamma }_{\alpha }^{}\mathrm{\Gamma }_{\mu \alpha }{}_{}{}^{\beta }\mathrm{\Gamma }_{\nu \beta }^{}{}_{}{}^{\alpha })$$
(8)
A natural question arises: our field $`\mathrm{\Gamma }_{\mu \nu }^\alpha `$ has a lot of independent components (40 in $`d=4`$ instead of two helicities for massless spin-2 particle), so there should exist a large gauge symmetry in such a model. And indeed, it is easy to check that the kinetic terms in our second order Lagrangian are invariant under the local ”affine” transformations:
$`\delta \mathrm{\Gamma }_{\mu \nu }^\alpha `$ $`=`$ $`_\mu z_\nu {}_{}{}^{\alpha }+_\nu z_\mu {}_{}{}^{\alpha }+{\displaystyle \frac{1}{d1}}[\delta _\mu {}_{}{}^{\alpha }(z)_{\nu }^{}+\delta _\nu {}_{}{}^{\alpha }(z)_{\mu }^{}]`$ (9)
$`{\displaystyle \frac{1}{d1}}[\delta _\mu {}_{}{}^{\alpha }_{\nu }^{}z+\delta _\nu {}_{}{}^{\alpha }_{\mu }^{}z]`$
where $`z_\mu ^\nu `$ is arbitrary second rank tensor and $`z=z_\mu ^\mu `$.
Now, having in our disposal an alternative description for massless spin-2 particle, it is natural to see how an interaction in such dual theory looks like. Nice feature of Palatini formulation is that switching on an interaction is a simple one step procedure . But as we have seen, it is very important for the possibility to construct dual formulations to work not in a flat Minkowski space but in (Anti) de Sitter space. So we start with the usual Lagrangian with the cosmological term:
$$=\sqrt{g}g^{\mu \nu }R_{\mu \nu }+\mathrm{\Lambda }\sqrt{g}$$
(10)
where now
$$R_{\mu \nu }=\frac{1}{2}(D_\mu \mathrm{\Gamma }_\nu +D_\nu \mathrm{\Gamma }_\mu )D_\alpha \mathrm{\Gamma }_{\mu \nu }{}_{}{}^{\alpha }+\mathrm{\Gamma }_{\mu \nu }{}_{}{}^{\alpha }\mathrm{\Gamma }_{\alpha }^{}\mathrm{\Gamma }_{\mu \alpha }{}_{}{}^{\beta }\mathrm{\Gamma }_{\nu \beta }^{}^\alpha $$
(11)
Then we introduce a convenient combination $`\widehat{g}^{\mu \nu }=\sqrt{g}g^{\mu \nu }`$ and rewrite a Lagrangian as:
$$=\widehat{g}^{\nu \nu }R_{\mu \nu }+\mathrm{\Lambda }det(\widehat{g}^{\mu \nu })^{\frac{1}{d2}}$$
(12)
The crucial point here is that first term contains $`\widehat{g}`$ only linearly. As a result it is possible to get complete nonlinear solution of the $`\widehat{g}`$ equations. We obtain (up to some numerical coefficients):
$$\widehat{g}^{\mu \nu }\sqrt{det(R_{\mu \nu })}(R^{\mu \nu })^1$$
(13)
At last, if we put this expression back into the first order Lagrangian we obtain (again up to normalization) a very simple and elegant Lagrangian:
$$=\sqrt{det(R_{\mu \nu })}$$
(14)
And it is just a Lagrangian written by Eddington eighty years ago in his book ! This result is very natural because this Lagrangian is the only invariant that could be constructed out of the affine connection alone, without any use of metric or any other objects, but it is exiting that this Lagrangian turns out to be dual formulation of usual gravity theory. Let us stress once again that working in a flat Minkowski space it is very hard if at all possible to give any reasonable physical interpretation to such model. But let us consider this model on a (Anti) de Sitter background. For that purpose we represent a total connection as $`\mathrm{\Gamma }_{\mu \nu }{}_{}{}^{\alpha }=\overline{\mathrm{\Gamma }}_{\mu \nu }{}_{}{}^{\alpha }+\stackrel{~}{\mathrm{\Gamma }}_{\mu \nu }^\alpha `$ where $`\overline{\mathrm{\Gamma }}_{\mu \nu }^\alpha `$ is a background connection while $`\stackrel{~}{\mathrm{\Gamma }}_{\mu \nu }^\alpha `$ — small perturbation around it (see e.g. ). Then for the curvature tensor we will have:
$$R_{\mu \nu ,\alpha }{}_{}{}^{\beta }=\overline{R}_{\mu \nu ,\alpha }{}_{}{}^{\beta }+[D_\mu \stackrel{~}{\mathrm{\Gamma }}_{\nu \alpha }{}_{}{}^{\beta }+\stackrel{~}{\mathrm{\Gamma }}_{\mu \alpha }{}_{}{}^{\rho }\stackrel{~}{\mathrm{\Gamma }}_{\rho \nu }^{}{}_{}{}^{\beta }(\mu \nu )]$$
(15)
where $`\overline{R}_{\mu \nu ,\alpha }^\beta `$ is a curvature tensor for the background connection and $`D_\mu `$ is a derivative covariant with respect to $`\overline{\mathrm{\Gamma }}`$. Then for the constant curvature space we have $`\overline{R}_{\mu \nu }=\mathrm{\Lambda }\overline{g}_{\mu \nu }`$ so the Lagrangian takes the form:
$$=\sqrt{det(\mathrm{\Lambda }\overline{g}_{\mu \nu }+\stackrel{~}{R}_{\mu \nu })}$$
(16)
It is interesting that Lagrangians of such kind have already been investigated e.g. in attempts to construct gravitational analog of the Born-Infeld electrodynamics. But now the interpretation of the Lagrangian is drastically different. Indeed, let us use the well known decomposition for the determinant
$$\sqrt{det(I+A)}=1+\frac{1}{2}Sp(A)+\frac{1}{8}(Sp(A))^2\frac{1}{4}Sp(A^2)+\mathrm{}.$$
where $`A`$ is any matrix. Then if we consider the curvature $`R_{\mu \nu }`$ as being expressed in terms of metric and its second derivatives the first linear terms gives scalar curvature while quadratic terms give higher derivative terms leading to the appearance of ghosts. But here the main dynamical quantity is affine connection $`\mathrm{\Gamma }`$ and the curvature $`R_{\mu \nu }`$ contains only first derivatives. As a result a term linear in $`R`$ is just a total derivative and could be dropped out of the action, while the quadratic terms give exactly the kinetic terms we obtained above.
Finally, let us add some comments on possible interaction with matter in such formulation of gravity. The most clear and straightforward way to obtain these interactions is to start with usual interactions in first order form and then try to go to the dual formulation. For example, for the scalar field we get:
$``$ $`=`$ $`\sqrt{g}[g^{\mu \nu }R_{\mu \nu }+\mathrm{\Lambda }+{\displaystyle \frac{1}{2}}g^{\mu \nu }_\mu \phi _\nu \phi {\displaystyle \frac{m^2}{2}}\phi ^2]=`$ (17)
$`=`$ $`\widehat{g}^{\mu \nu }(R_{\mu \nu }+{\displaystyle \frac{1}{2}}_\mu \phi _\nu \phi )+(\mathrm{\Lambda }{\displaystyle \frac{m^2}{2}}\phi ^2)det(\widehat{g}^{\mu \nu })^{\frac{1}{d2}}`$
and the second line shows that the main effect is the replacement of $`R_{\mu \nu }`$ by $`R_{\mu \nu }+\frac{1}{2}_\mu \phi _\nu \phi `$ (compare ). Also if scalar field has nonzero mass the cosmological constant $`\mathrm{\Lambda }`$ is replaced by field dependent combination $`\mathrm{\Lambda }\frac{m^2}{2}\phi ^2`$. But for the vector field (even massless) the situation turns out to be much more complicated because even for the minimal interaction:
$``$ $`=`$ $`\sqrt{g}[g^{\mu \nu }R_{\mu \nu }+\mathrm{\Lambda }{\displaystyle \frac{1}{4}}g^{\mu \alpha }g^{\nu \beta }F_{\mu \nu }F_{\alpha \beta }]=`$ (18)
$`=`$ $`\widehat{g}^{\nu \nu }R_{\mu \nu }+\mathrm{\Lambda }det(\widehat{g}^{\mu \nu })^{\frac{1}{d2}}{\displaystyle \frac{1}{4}}det(\widehat{g}^{\mu \nu })^{\frac{1}{d2}}\widehat{g}^{\mu \alpha }\widehat{g}^{\nu \alpha }F_{\mu \nu }F_{\alpha \beta }`$
equations for the $`\widehat{g}`$ become highly nonlinear. But in a weak field approximation such model could reproduce a correct kinetic term for the vector field. Note also that the corrections to the $`R_{\mu \nu }`$ tensor here start with the terms quadratic in $`F_{\mu \nu }`$ and there is no term linear in it in contrast with .
Thus we have shown that the dualization procedure based on the use of (Anti) de Sitter background space could be applied to the gravity theory in a Palatini formalism and leads to the formulation in terms of affine connection. In this, the final Lagrangian coincides with that of Eddington . A number of interesting question arises, for example, whose related with the gauge symmetries of such formulation, which deserve further study. |
warning/0506/nucl-th0506062.html | ar5iv | text | # Relativistic Model of Triquark Structure
## I INTRODUCTION
There has been a great deal of interest in the study of pentaquarks and a large number of experiments have been carried out \[1-11\]. The existence of pentaquarks is uncertain since they have been seen in some experiments and not in others. Various groups hope to clarify this situation by performing more precise experimental searches. The $`\mathrm{\Theta }^+(1540)`$ which decays to a kaon and a nucleon has been seen in several experiments. It has been interpreted as a pentaquark with a $`udud\overline{s}`$ structure . A pentaquark $`\mathrm{\Theta }_c^0`$ with the assumed structure $`udud\overline{c}`$ has also been observed recently. A recent review may be found in Ref. .
We were particularly interested in the diquark-triquark model of Karliner and Lipkin which has been applied in the study of pentaquarks , and we have made use of a variant of that model in our work . One problem for the theorist has been the very small widths of the observed pentaquarks. We have studied that question in a relativistic diquark-triquark model and found we could explain the small widths seen in experiment . In our model, as in that of Refs. , the pentaquark is described as scalar diquark coupled to a triquark. \[See Fig. 1.\] Using the insight gained in our analysis of the nucleon, which made use of a quark-diquark model , we took the scalar diquark mass to be 400 MeV in our study of the pentaquark.
There have been many studies of diquark structure making use of the Nambu-Jona-Lasinio (NJL) model. Some of that work is reviewed in Ref. . One may suggest that, in addition to studies of diquark structure, a study of triquark structure may be of interest. In our earlier work the triquark mass was taken to be 800 MeV on phenomenological grounds. We note that a calculation of triquark properties, using the operator product expansion and including direct instanton contributions, obtained a triquark of $`ud\overline{s}`$ structure of mass 800 MeV . An additional triquark state was found at 900 MeV in Ref. . (Another work making use of QCD sum rules yields quite small values for the width of the $`\mathrm{\Theta }^+(1540)`$ pentaquark .)
Once we decide to study triquark structure, we face the following problem. The triquark of mass 800 MeV is very close to the threshold for decay to a 400 MeV diquark and a 450 MeV strange quark. Similarly, the triquark mass is close to the mass of a $`u`$ (or $`d`$) quark of 350 MeV and a kaon of mass 495 MeV. This difficultly cannot be overcome by including a confinement model since there is no confining interaction between a kaon and a quark. Therefore, in the present work we have made use of the quark propagator obtained in Ref. . In that work we considered quark propagation in the presence of a gluon condensate and found that the quark propagator had no on-mass-shell poles. That is, the quark was a non-propagating mode in the presence of the condensate. As we will see, the use of the quark propagator of Ref. enables us to proceed in our analysis of triquark structure. (In an early work we used the form of the propagator discussed in Ref. in a study of nontopological solitons .)
The organization of our work is as follows. In Section II we review our model of quark propagation in the presence of a gluon condensate. In that model the quark propagator has no on-mass-shell poles. In Section III we present the equation for the vertex describing triquark decay to the channels: i) a $`u`$ quark plus a $`\text{K}^0`$ meson, ii) a scalar diquark plus a $`\overline{s}`$ quark and iii) a $`d`$ quark and a $`\text{K}^+`$ meson. \[See Figs. 2-4.\] In Section IV we describe the results of our analysis and in Section V we present some further comments and conclusions. The Appendix contains a discussion of the normalization of the wave functions of the scalar diquark and the kaon.
## II Quark Propagator in The Presence of a Gluon Condensate
In an earlier work we discussed quark propagator in the presence of a condensate of the form $`<g^2A_\mu ^aA_a^\mu >`$ which has recently been shown to be the Landau gauge version of a more general gauge invariant expression. We have discussed quark propagation in the presence of such a condensate treating the vacuum as a random medium of gluon fields . It is found that the quark propagator has no on-mass-shell poles indicating that the quark cannot propagate over extended distances. As an example, we show one of the momentum-dependent mass functions obtained in our model in Fig. 5. It may be seen that the equation $`p^2M^2(p^2)=0`$ has no solution. In our work we have modified the solution shown in Fig. 5 to have a constituent mass value of 350 MeV for the up (or down) quark for spacelike $`k^2`$ and for a small region of timelike $`k^2`$ near $`k^2=0`$. We use a simplified form for the momentum-dependent mass function.
$`M(k^2)=[k^2+c^2]^{1/2}\text{for}k^2>m_q^2c^2,`$ (2.1)
and
$`M(k^2)=m_q\text{for}k^2<m_q^2c^2,`$ (2.2)
with $`c=0.3`$ GeV and with $`m_q`$ being the quark mass which we take to be 350 MeV for the up (or down) quark and 450 MeV for the strange quark. In contrast to the result shown in the Fig. 5, we use the constituent quark mass for $`M(k^2)`$ when $`k^2<m_q^2c^2`$. That is more in keeping with standard phenomenology, since the result shown in Fig. 5 does not capture the behavior expected for the constituent quark mass for the spacelike values of $`k^2`$.
## III Dynamical Equations For the Triquark vertex function
In order to construct a bound-state triquark wave function we consider the diagrams of Figs. 2-4. In Fig. 2 we show the vertex (open circle) for the virtual triquark decay to a $`u`$ quark of momentum $`Pk`$ and a $`K^0`$ meson of momentum $`k`$. The $`K^0`$ and $`K^+`$ mesons, and the diquark are taken to be on mass shell in our formalism. (Such restrictions arise when we complete the $`k_0^{}`$ integral in the complex $`k_0^{}`$ plane.) On the right-hand side of Fig. 2 we see the triquark component consisting of a $`\overline{s}`$ quark and a scalar diquark. The final state is reached by the exchange of a $`d`$ quark of momentum $`kk^{}`$. In the second figure on the right we have a $`d`$ quark and a $`K^+`$ meson in the intermediate state, with exchange of a $`\overline{s}`$ quark yielding the final $`K^0`$ and $`u`$ quark. Similar comments pertain to the processes shown in Figs. 3 and 4.
The triquark vertex depends upon $`P`$ and $`k`$ and has a Dirac index $`\alpha `$: $`\mathrm{\Gamma }_\alpha (P,k)`$. We introduce coupling constants $`g_1`$ and $`g_2`$ which correspond to the coupling of either the kaon or scalar diquark to their quark components. \[See the Appendix.\] By completing the integral over the $`k_0^{}`$ variable, we find we may place the kaon and the diquark on mass shell, leaving a three-dimensional integral over $`\stackrel{}{k}^{}`$. It is also useful to solve for $`\stackrel{}{k}\mathrm{\Gamma }_\alpha (P,k)`$ rather than $`\mathrm{\Gamma }_\alpha (P,k)`$. We take $`\stackrel{}{P}=0`$ and find a bound state at a specific value of $`P^0`$. The equation we solve may be written with the Dirac indices explicit:
$`|\stackrel{}{k}|\mathrm{\Gamma }_\alpha (\stackrel{}{k})`$ $`=`$ $`{\displaystyle \frac{g_1g_2}{(2\pi )^2}}{\displaystyle _1^1}𝑑u{\displaystyle _0^{k_{cut}}}d|\stackrel{}{k^{}}|{\displaystyle \frac{|\stackrel{}{k}||\stackrel{}{k^{}}|}{2E_{mes}(k^{})}}`$
$`\times `$ $`F[(k^{}k/2)^2]F\left[\left(k{\displaystyle \frac{k^{}+P}{2}}\right)^2\right]`$
$`\times `$ $`{\displaystyle \frac{[\text{ / }k\text{ / }k^{}M_1((kk^{})^2)]_{\alpha \mu }}{(kk^{})^2M_1^2(kk^{})^2)}}{\displaystyle \frac{[\text{ / }k^{}+M_2(k^2)]_{\mu \beta }}{k^2M_2^2(k^2)}}[|\stackrel{}{k^{}}|\mathrm{\Gamma }_\beta (\stackrel{}{k^{}})].`$
The values of $`k_0`$ and $`k_0^{}`$ are fixed using the on-mass-shell conditions for the kaon and diquark. Here, $`M_1`$ and $`M_2`$ are either $`M_u,M_d`$ or $`M_s`$ depending upon which diagram of Figs. 2-4 is being considered. The $`F^{}`$s are form factors introduced for the diquark and kaon vertices which appear in the figures as small filled circles.
The form factor for the final-state kaon or diquark is
$`F[(k^{}k/2)^2]=\mathrm{exp}\left[{\displaystyle \frac{1}{\alpha ^2}}|(k^{}k/2)^2|\right],`$ (3.2)
with
$`(k^{}k/2)^2=(P^0E_{mi}(\stackrel{}{k^{}})E_{mf}(\stackrel{}{k})/2)^2(\stackrel{}{k^{}}^2+\stackrel{}{k}^2/4|\stackrel{}{k}||\stackrel{}{k^{}}|u),`$ (3.3)
where $`E_{mi}(\stackrel{}{k^{}})`$ refers to either the kaon or diquark of momentum $`\stackrel{}{k^{}}`$ and $`E_{mf}(\stackrel{}{k})`$ refers to the final-state kaon or diquark.
The form factor for the intermediate-state kaon or diquark is
$`F\left[\left(k{\displaystyle \frac{k^{}+P}{2}}\right)^2\right]=\mathrm{exp}\left[{\displaystyle \frac{1}{\alpha ^2}}\left|\left(k{\displaystyle \frac{k^{}+P}{2}}\right)^2\right|\right],`$ (3.4)
with
$`\left(k{\displaystyle \frac{k^{}+P}{2}}\right)^2=(E_{mf}(\stackrel{}{k^{}})P^0E_{mi}(\stackrel{}{k})/2)^2(\stackrel{}{k}^2+\stackrel{}{k^{}}^2/4|\stackrel{}{k}||\stackrel{}{k^{}}|u).`$ (3.5)
We remark that there are two terms to be considered on the right-hand side of Eq. (3.1) when we relate that expression to the diagrams of Figs. 2-4. There is also an implied sum over the Dirac indices, $`\mu `$ and $`\beta `$. Since there are three decay channels ($`K^0u`$, $`0^+\overline{s}`$ and $`K^+d`$) and four Dirac indices (0,1,2,3), there are twelve vertex functions to consider. If we take $`N`$ points for each vertex function, we need to evaluate a $`12N`$ by $`12N`$ matrix when searching for the bound-state eigenvalue. In our calculation it is useful to take $`\stackrel{}{k}`$ along the $`z`$-axis so the vector $`\stackrel{}{k^{}}`$ has components $`\stackrel{}{k^{}}=\stackrel{}{k^{}}(\mathrm{sin}\theta \mathrm{cos}\varphi ,\mathrm{sin}\theta \mathrm{sin}\varphi ,\mathrm{cos}\theta )`$. We have used $`u=\mathrm{cos}\theta `$ in writing Eq. (2.1).
Once the eigenvalue is found, we may then obtain the vertex functions or the corresponding wave function. The wave function is
$`\psi _\alpha (\stackrel{}{k})={\displaystyle \frac{1}{k_0^2\stackrel{}{k}^2M^2(k^2)}}\mathrm{\Gamma }_\alpha (\stackrel{}{k}),`$ (3.6)
where $`k^0=P^0E_{mes}(\stackrel{}{k})`$ and $`P^0`$ is the eigenvalue. \[See Figs. 2-4.\]
## IV Wave Functions of the Triquark
In Fig. 6 we show the (unnormalized) wave functions found in our analysis. The dashed line shows the diquark-(strange quark) component, while the solid line exhibits the $`K^0u`$ and $`K^+d`$ components which are equal in this model. The small components of these wave function are shown in the lower part of the figure. We may write the four-component wave function of Eq. (3.6) as
$`\psi _s(\stackrel{}{k})=\left(\begin{array}{c}R_u(k)|s\\ \stackrel{}{\sigma }\widehat{k}R_l(k)|s\end{array}\right).`$ (4.3)
The upper and lower components of the wave function are shown in Fig. 6. There are three wave functions of the form of Eq. (4.1) corresponding to the channels $`0^++s`$, $`K^0+u`$ and $`K^++d`$. As noted above, the wave functions for the $`K^0+u`$ and $`K^++d`$ components are equal in our model.
## V Discussion
Our interest in triquark structure is related to the diquark-triquark model of pentaquark structure \[14-17\]. As stated earlier, it is not clear that pentaquarks exist because of various contradictory results obtained in experimental studies. Recent work of Karliner and Lipkin appears to be quite important for the interpretation of experimental searches for the pentaquark. These authors claim that : ”Significant signal-background interference effects can occur in experiments like $`\gamma N\overline{K}\mathrm{\Theta }^+`$ as a narrow $`I=0`$ resonance in a definite final state against a non-resonant background, with an experimental resolution coarser than the expected resonance width. We show that when the signal and background have roughly the same magnitude, destructive interference can easily combine with a limited experimental resolution to completely destroy the resonance signal. Whether or not this actually occurs depends critically on the yet unknown phase of the $`I=0`$ and $`I=1`$ amplitudes …”.
In the present work we have introduced a model of triquark structure. In order to carry out our calculation we have used a quark self-energy that does not give rise to on-mass-shell poles. Similar results for gluon propagation in the presence of a condensate are presented in Ref. , where it is shown that the gluon is also a nonpropagating mode in the presence of the gluon condensate.
It would be desirable to improve the model presented in our work and to see if there are other useful applications of the quark propagator used in this work. Whether our triquark model may be used in a more detailed description of the pentaquark than that we have presented previously remains to be seen.
## Appendix A
In order to calculate the normalization factor for our kaon or diquark we consider the diagram shown in Fig. 7. We define
$`N`$ $`=`$ $`\text{Tr}{\displaystyle \frac{d^4k}{(2\pi )^4}\frac{\text{ / }p\text{ / }k+M_1[(pk)^2]}{(pk)^2M_1^2[(pk)^2]}\frac{\text{ / }k+M_2(k^2)}{k^2M_2^2(k^2)}}`$
$`\times (\gamma \widehat{n}){\displaystyle \frac{\text{ / }k+M_2(k^2)}{k^2M_2^2(k^2)}}f^2(p,k),`$
where $`\widehat{n}=(1,0,0,0)`$. (We may also identify $`g=1/\sqrt{N}`$ as the effective coupling constant at the kaon or diquark vertex.) In Eq. (A1) $`f(p,k)`$ is a form factor defined at the kaon or diquark vertex. When $`\stackrel{}{p}=0`$, we have
$`f(p,k)=\mathrm{exp}\left[{\displaystyle \frac{1}{\alpha ^2}}|(p_0/2k_0)^2\stackrel{}{k}^2|\right].`$ (A2)
In order to calculate $`N`$ of Eq. (A1) we need the value of the trace
Trace $`=`$ $`\text{Tr}[\text{ / }p\text{ / }k+M_1[(pk)^2][\text{ / }k+M_2(k^2)](\gamma \widehat{n})[\text{ / }k+M_2(k^2)]`$
$`=`$ $`8k_0^2p_04(k_0^2\stackrel{}{k}^2)p_0+4p_0M_2^2(k^2)`$
$`4k_0(k_0^2\stackrel{}{k}^2)4k_0M_1^2(k^2)+8k_0M_1[(pk)^2]M_2(k^2).`$
We obtain
$`N`$ $`=`$ $`{\displaystyle \frac{1}{4\pi ^3}}{\displaystyle _{k_{max}^0}^{k_{max}^0}}𝑑k_0{\displaystyle _0^{k_{max}}}\stackrel{}{k}^2d|\stackrel{}{k}|\left({\displaystyle \frac{1}{(p_0k_0)^2\stackrel{}{k}^2M_1[(pk)^2]}}\right)`$
$`f^2(p_0,k)\left({\displaystyle \frac{1}{k_0^2\stackrel{}{k}^2M_2(k^2)}}\right)(\text{Trace}).`$
Here $`k`$ is a four-vector $`k=(k_0,\stackrel{}{k})`$. In our analysis we put $`k_{max}=k_{max}^0=0.6`$ GeV. (For the diquark, we find $`1/\sqrt{N}=11.21`$. We use the same value for the kaon since that value is only $`1.5\%`$ greater than $`1/\sqrt{N}`$ for the diquark.) |
warning/0506/cond-mat0506657.html | ar5iv | text | # SOLID FRICTION FROM STICK-SLIP DOWN TO PINNING AND AGING
##
An established tradition, when writing about solid friction, is to date its emergence as a well identified scientific question from Leonardo da Vinci (ca 1500). This makes it an unusually longstanding problem, since both its fundamental physical aspects and its modelization for the purpose of studies of e.g. seismic fault dynamics are still under lively debate nowadays (see for example Dawson , MRSAmontons , Bo , Scholz ). A very important progress in its phenomenologic modelling, for macroscopic solids, was accomplished with the formulation of the so-called rate and state constitutive laws, which emerged in the 70’s from the work on rock friction of Dieterich Dieterich , Rice and Ruina RR .
It was later shown that, inspite of their simplicity in terms of number of parameters, they provide an excellent description of most of the salient features of the low-velocity frictional dynamics of a wide variety of materials, ranging from granite to paper. Such an amazing “universality” is naturally appealing for the physicist, since it suggests :
– The possibility of a unified, largely material-independent, description on the underlying microscopic level.
– One step further, a possible feedback in terms of limitations and/or refinement of the mechanical constitutive laws.
It is on this particular approach, deliberately different from that of tribology proper <sup>1</sup><sup>1</sup>1 This means, in particular, that we limit ourselves to systems and sliding regimes such that wear and frictional self-heating play a negligible role. , but which parallels those presently developed in the fields of plasticity of amorphous materials Langer1 and rheology of jammed systems livre Nagel , that we concentrate here.
## I Introduction
### I.1 From Amontons-Coulomb to Rate-and-State
For more than two centuries, the description of solid friction commonly used in mechanical modelling was provided by the classical Amontons-Coulomb laws. They state that, when a nominally planar solid block lying on top of a planar track is submitted to a normal force $`W`$ and a tangential one $`F`$ (Figure 1) :
\- No motion occurs as long as $`F`$ is smaller than a finite threshold $`F_s`$.
\- Sliding is dissipative, and the corresponding dynamical friction force $`F_d`$ is constant and equal to $`F_s`$.
\- Their common value $`F`$ is proportional to the normal load $`W`$ and, for a given $`W`$, independent of the macroscopic contact area $`\mathrm{\Sigma }`$. Hence, the frictional behavior of a couple of materials is characterized by a single number, the friction coefficient :
$$\mu =\frac{F}{W}$$
(1)
Note that this behavior (static threshold plus velocity-independent dynamic force) is the exact analogue, for friction, of the Hill rigid-plastic model of plasticity (threshold plus constant yield stress), and thus probably carries unphysical singularities as is the case of the Hill model.
Indeed, various departures from this description or its implications have gradually emerged, the most salient of which are the following :
(i) In general, the static friction coefficient $`\mu _s`$ is larger than the dynamic one $`\mu _d`$.
(ii) $`\mu _s`$ is not a mere number, but a slowly increasing function of the so-called waiting time $`t_w`$, i.e. the duration of static contact prior to sliding.
(iii) When measured in stationary sliding, $`\mu _d`$ is not constant. In particular, for low enough velocities (typically $`<100\mu `$m.sec<sup>-1</sup>) it is a slowly decreasing function of $`V`$.
That a stationary $`\mu _d`$ can be measured is then by itself a puzzle, since such a velocity-weakening characteristic is well known to make steady motion always unstable : a positive (negative) velocity fluctuation induces a decrease (increase) of the friction force, leading to further accelaration (deceleration), thereby getting amplified. Indeed, when a slider is pulled at a low enough velocity through a driving stage, necessarily of finite stiffness, non steady sliding is frequently observed : the motion is jerky, alternating between “stick” periods of rest during which the stage stores elastic energy, and slip events (Figure 2). However, as skilled mechanical engineers have known already for long, these stick-slip oscillations disappear upon stiffening the driving stage, and steady sliding is realized.
This paradox points towards some inadequacy in the implicit assumption underlying the above instability argument - namely that the stationary $`\mu _d(V)`$ function can be used as such to describe non steady motion. In other words, one suspects that the dynamic friction force does not only depends on the rate variable, i.e. the instantaneous sliding velocity.
(iv) Non steady friction is hysteretic, as illustrated on Figure 3, which shows the instantaneous friction force associated with a velocity cycle for which inertia is negligible.
(v) Dieterich Dieterich has studied the frictional response to a sudden jump of the driving velocity from $`V_i`$ to $`V_f`$ (Figure 4). He showed that it exhibits a transient the span of which is controlled by a characteristic length $`D_0`$, of micrometric order. That is, its duration $`\mathrm{\Delta }tD_0/V_f`$.
These observations can be translated into the following statements :
$``$ A frictional interface at rest becomes stronger as time lapses : it strengthens when aging (see (ii)).
$``$ When sliding, it becomes weaker (see (i)) : it can be said to rejuvenate upon sliding. Moreover, its dynamic age in steady motion decreases with increasing $`V`$ (see (iii)).
$``$ This, together with (iv) and (v), means that its physical state evolves with a dynamics characterized by the length $`D_0`$ and coupled to the sliding dynamics itself.
Rate and state constitutive laws model such behaviors in terms of a single (or a few) dynamical state variable(s), the physical nature of which often remains unspecified. In such a situation, the physicist’s approach aims at identifying the physical content embedded in the state variables, and at justifying on this basis their phenomenological dynamics. One may then hope, as a side benefit of importance, to be able to predict their limits of validity and, possibly, to propose some further extensions.
### I.2 Spatial scales
For this purpose, a natural first step is to identify the relevant length scales. In the case of friction between macroscopic solids, this immediately leads to distinguishing between two classes of systems.
(A) Rough hard solids for which “reasonable” loading levels (apparent pressures $`W/\mathrm{\Sigma }`$ well below elastic moduli) do not result in an intimate molecular contact along the whole interfacial area.
For them a first relevant length scale is the average size of the microcontacts – the spots which form the real area of contact $`\mathrm{\Sigma }_r`$. It lies usually in the micrometric range. Moreover, we will see that a detailed analysis leads to attributing ultimately frictional dissipation to elementary mechanical instabilities involving molecular rearrangements on the nanometric scale.
(B) Soft and/or smooth solids which are able to get into intimate contact everywhere along their interface. For them, the nanometric scale retains its importance, while the mesoscopic one becomes irrelevant.
In both cases, due to the long range of elastic interactions, a third important length scale is the global size of the system $`L`$. Its value has a crucial bearing on the spatial (in)homogeneity of the frictional motion, hence on the complexity of the sliding dynamics.
An important consequence, not to be overlooked, results straightforwardly from the identification of these spatial scales. Reducing the “state” of the interface to one or a few variables necessarily implies a statistical averaging, which can be meaningful only if performed on a large number of units. This entails that a rate and state phenomenology of friction is legitimate only on scales larger than a finite cutoff. In particular, when using discretized (block) models for e.g. numerical studies of extended interface dynamics, the basic block size must remain much larger than either the average intercontact distance (class (A)) or at least the ultimate nanometric scale (class (B)). That is, we contend that the question of the regularity of the continuum limit of discretized models — a subject of debate in the field of mode II interfacial fracture along seismic faults Rice-Cochard — though of course mathematically sound, is physically irrelevant : if it should turn out that a continuum description of the fracture head zone would imply lengths smaller than the intercontact distance, the interface could no longer be assumed to have homogeneous frictional properties, but should be explicily treated as a juxtaposition of frictional (contact) and non-frictional (non-contact) regions.
### I.3 Outline
Following Tabor Tabor , it is both useful and physically sound to express the friction force between two solids as
$$F=\sigma _s.\mathrm{\Sigma }_r$$
(2)
where $`\mathrm{\Sigma }_r`$ is the real contact area, and the stress $`\sigma _s`$ is the so-called shear strength of the interface.
Section II is devoted to friction at multicontact interfaces between rough hard materials (MCI), an interfacial configuration which is prevalent when dealing with macroscopic bodies. We first analyze their geometry, then show that the geometric factor $`\mathrm{\Sigma }_r`$ plays, for these systems, the part of a state-dependent variable, governed by what we will call the geometric age $`\varphi `$. We then show that, for MCI, the rate dependence of the rheological factor $`\sigma _s`$ can be assigned to local mechanical instabilities within the nanometer-thick molecular “junctions” forming the real contacts. On this basis, the Rice-Ruina constitutive law, which they originally formulated on a phenomenological basis, results as a good approximation for the low velocity frictional behavior of MCI. We sketch out its consequences in terms of the sliding dynamics of a driven spring-block system. We also point out its various limits of validity.
These limits turn out to be of two different types.
$``$ On the one hand, within the physical framework which identifies “state” with geometric age, the phenomenological functional expression of $`\mathrm{\Sigma }_r(\varphi )`$ is necessarily aproximate, and its limits can be evaluated.
$``$ On the other hand, the detailed analysis of experiments brings to light the limits of the above physical framework itself. That is, it prompts the idea that a second underlying slow dynamics manifests itself through the rheological factor : geometric age is not the whole story, frictional contacts also have a structural age.
This leads us to concentrate, in Section III, on the rheology of frictional contacts. We analyze it for various configurations (rough-on-flat MCI, surface force aparatus (SFA), extended soft contacts) corresponding to various confinement and compacity levels. In all cases, it appears that structural aging/rejuvenation mechanisms are indeed at work in the contact-forming nanometer-thick interfacial junctions. The associated dynamics is quite strongly system-dependent. As a guide for future investigation, we suggest a first level of classification which distinguishes between two main types of dissipative behaviors, namely :
(i) jammed junction plasticity, akin to that of soft glassy materials;
(ii) adsorption–desorption controlled dynamics.
MCI microcontacts belong to class (i), gel/glass contacts to class (ii). Both mechanisms probably contribute to friction in SFA contacts, while the case of elastomer/glass remains open to discussion.
## II Multicontact interfaces
The renewal of interest for solid friction in the past 20 years, triggered by the work of rock mechanicians, was aimed at modelling friction along seismic faults. For this reason, it naturally focussed on MCI. It has resulted in bringing to light robust features shared by a wide variety of materials. We will see that this has led to the building a complete framework of description, namely :
— a predictive constitutive law
— an underlying physical interpretation opening onto further questions concerning contact rheology (Section III).
This is why MCI friction is treated here extensively.
### II.1 Geometry of multicontact interfaces
#### II.1.1 Surface roughness
Consider two macroscopic solids (Figure 1), referred to as slider and track, with nominally planar contact surfaces. Nominally means here that they are flat on large scales comparable with the macroscopic slider lateral size. It is well known that, due to the unavoidable presence of steps, atomic flatness cannot be realized over lengths much beyond the micrometric scale. An exception is provided by lamellar solids, in particular cleaved mica for which this range may reach up to $`1`$ cm. Optically smooth surfaces, with r.m.s. roughness on the order of a few nanometers over $`1`$ cm<sup>2</sup>, are uncommon, examples being provided by float glass and some highly polished metals. In general, natural or ground surfaces exhibit a r.m.s. roughness between $`0.1`$ and a few microns on $`1`$ cm<sup>2</sup>.
Such random surfaces exhibit asperities with distributed heights, so that physical contact between them occurs only at random spots, the microcontacts between load-bearing asperities. They form what we call multicontact interfaces (MCI) (Figure 5).
It is Tabor Tabor , in his pioneering work, who emphasized the importance of distinguishing between the apparent and the real areas of contact, $`\mathrm{\Sigma }_{app}`$ and $`\mathrm{\Sigma }_r`$, and of evaluating the average microcontact size $`\overline{a}`$. It has emerged (see below) that typical values for $`\mathrm{\Sigma }_r/\mathrm{\Sigma }_{app}`$ and $`\overline{a}`$ are respectively $`10^3`$ and a few microns. This means that a typical MCI is formed of a sparse set of microcontacts with average separation $`100\mu `$m – a picture which has been confirmed by direct optical observation (Diet-Kilg and Figure 5). It is intuitively clear that increasing the normal load $`W`$ results in decreasing the distance between the average surface planes (the so-called “closure”), hence in increasing the number of microcontacts as well as the area of preexisting ones : $`\mathrm{\Sigma }_r`$ increases with $`W`$, as illustrated by Dieterich’s observations Diet-Kilg . Can one make the functional nature of this dependence explicit? Does it depend on the statistical nature of the surface roughness and if so, how? Could it be that this dependence would explain the Amontons proportionality between $`F`$ and $`W`$?
#### II.1.2 The single microcontact
As a first step, let us recall a few basic results from contact mechanics concerning a single contact. Let us focus on the a single pair of contacting asperities, modelled as elastic spherical caps with a common radius of curvature $`R`$ and Young modulus $`E`$, pressed together by a normal force $`w`$ acting along their intercenter axis. The problem of the resulting elastic contact was fully solved by Hertz and is equivalent to that between a rigid plane and an elastic sphere of radius $`R^{}=R/2`$ and Young modulus $`E^{}=E/2(1\nu ^2)`$, with $`\nu `$ the Poisson modulus Johnson .
Its solution can be evaluated as follows. Due to the spherical geometry, the contact diameter $`a`$, the “compression” $`\delta `$ (Figure 6) and $`R^{}`$ are related by $`a^2R^{}\delta `$. As can be checked from the full Hertz solution Johnson ,the elastic energy is essentially stored within a depth on the order of the contact radius, hence the relevant strain level $`ϵ\delta /a`$, so that the average normal stress $`\overline{p}w/a^2E^{}ϵ\delta /a`$. From this it immediately results that, dimensionally
$$a\left(\frac{wR^{}}{E}\right)^{1/3}\delta \left(\frac{w^2}{R^{}E^2}\right)^{1/3}\overline{p}\left(\frac{wE^2}{R^2}\right)^{1/3}$$
(3)
The exact expressions (see Johnson ) only differ by multiplicative constants of order unity.
Now, most solids are truly elastic only up to a yield stress $`Y`$, beyond which they start deforming plastically. So, the Hertz expressions above lose validity when the maximum normal stress below the contact reaches this level. Upon increasing the normal load, the size of the plastified region increases, until it occupies a volume $`a^3`$. At this stage the contact has become “fully plastic” and deforms so that the normal stress remains quasi-constant : $`\overline{p}H`$, where $`H`$ is the hardness <sup>2</sup><sup>2</sup>2On the basis of geometric arguments, one estimates that $`H3Y`$Tabor hardness . of the softer material. At room temperature, for metals, the ratio $`H/E`$ ranges in general between $`10^2`$ and $`10^3`$, while for polymer glasses, $`H/E10^1`$$`10^2`$. Imagine now, following Bowden and Tabor Tabor , that the apparent pressure $`p_{app}=W/\mathrm{\Sigma }_{app}`$ is large enough for this regime to be reached in all the microcontacts forming a MCI. Then, the real area of contact is such that $`\mathrm{\Sigma }_r/\mathrm{\Sigma }_{app}p_{app}/H`$, and $`\mathrm{\Sigma }_rW/H`$.
So, in the fully plastic regime, $`\mathrm{\Sigma }_r`$ is proportional to the normal load, and Amontons’s law simply follows. Indeed the friction coefficient reads :
$$\mu =\frac{F}{W}=\frac{\mathrm{\Sigma }_r}{W}\frac{\sigma _s}{H}$$
(4)
Since, in this approximation, $`\sigma _s`$ is the shear strength under the constant normal stress $`H`$, $`\mu `$ is effectively $`W`$-independent. However rough it is, this approximation is illuminating in several respects.
On the one hand, the a priori surprising fact that the values of dry friction coefficients depend only weakly on the mechanical properties of materials and are commonly a fraction of unity can now be translated into the statement : interfacial shear strength values are roughly comparable with bulk yield stress levels — a point which will be of qualitative importance later on.
On the other hand, it permits to get a somewhat more precise idea about the minimum loading level necessary to form a MCI between given materials. Consider as an example a $`1`$ cm thick steel plate under its own weight. Then $`p_{app}10^3`$ Pa, while $`H10^9`$ Pa, hence $`\mathrm{\Sigma }_r/\mathrm{\Sigma }_{app}10^6`$ : a $`1`$ cm<sup>2</sup> surface would form no more than about $`10`$ microcontacts of area $`10\mu `$m<sup>2</sup> (see below) – hardly enough to form a decent statistical set!
Tabor’s suggestion, which he developed for the case of metals, later gave rise to a number of discussions about the relevance of the full plasticity assumption to wider classes of materials which also obey Amontons’s law while being less ductile than metals. This opened onto the problem of modelling the elastic-plastic contact between random surfaces.
#### II.1.3 Area of contact between random surfaces
The crucial contribution to this question was made by Greenwood and Williamson GW . Their model reduces the characterization of each of the contacting random surfaces to :
– the statistical distribution of asperity summit heights above some average plane $`\varphi (z)`$. These asperities are assumed to have spherical tips.
– the asperity radius of curvature $`R`$, assumed to be unique.
– the number of summits $`N`$ on the apparent surface $`\mathrm{\Sigma }_{app}`$.
They show that the problem of contact between two such nominally flat elastic random surfaces can be cast into that of a “composite” random surface and a rigid plane Johnson .
Let $`d`$ be the separation between the average plane and the rigid one (Figure 7), $`n`$ the number of microcontacts at this separation, assumed to be dilute enough to be treated as independent Hertz contacts. Each asperity with height $`z>d`$ contributes to $`n`$ :
$$n=N_d^{\mathrm{}}𝑑z\varphi (z)$$
(5)
The compression of a contact is $`(zd)`$, its area $`\pi R(zd)`$, and it bears the load $`(4/3)\pi R^{1/2}(zd)^{3/2}`$ (eq.(3), Johnson ) so that the real area of contact
$$\mathrm{\Sigma }_r=N_d^{\mathrm{}}𝑑z\pi R(zd)\varphi (z)$$
(6)
and $`d`$ is related to the given total normal load $`W`$ by :
$$W=N_d^{\mathrm{}}𝑑z(4/3)\pi R^{1/2}(zd)^{3/2}\varphi (z)$$
(7)
The main physical content of the Greenwood -Williamson (GW) results already emerges under the schematic assumption of an exponential distribution $`\varphi _0(z)=s^1\mathrm{exp}(z/s)`$ for $`z>0`$ and zero otherwise. One then gets :
$$n=\frac{1}{\sqrt{\pi }}\frac{W}{ER^{1/2}s^{3/2}}\mathrm{\Sigma }_r=\sqrt{\pi }\left(\frac{R}{s}\right)^{1/2}\frac{W}{E}$$
(8)
Both $`\mathrm{\Sigma }_r`$ and the number of contacts are proportional to $`W`$, and individual contact sizes grow under increasing load in such a way that the average contact radius :
$$\overline{a}=\left(\frac{\mathrm{\Sigma }_r}{\pi n}\right)^{1/2}=\sqrt{Rs}$$
(9)
remains a load-independent constant.
GW have shown numerically that these properties are conserved, to a very good approximation, for a gaussian distribution of summit heights, in the load range, ranging over several decades, such that $`1nN`$. A rough evaluation of $`N`$ is obtained by representing $`\mathrm{\Sigma }_{app}`$ as densely paved by average asperities of curvature $`R^1`$ and height $`s`$, so that $`N\mathrm{\Sigma }_{app}/Rs`$. The upper limit then translates into
$$\frac{p_{app}}{E}\sqrt{\frac{s}{R}}\mathrm{or}\mathrm{\Sigma }_r\mathrm{\Sigma }_{app}$$
(10)
Note that this condition also ensures that elastic interactions between microcontacts can reasonably be neglected.
A considerable amount of work Nayak has been devoted to evaluating the two geometrical parameters involved in the GW model on the basis of surface topography measurements. Curvature determinations have been hotly debated, since they involve a second order derivative of the profile, which makes them strongly noise-sensitive. When measured with profilometers of lateral resolution of micrometric order, typical values of $`R`$ for surfaces blasted or lapped with abrasive powders commonly lie in the $`10`$$`100\mu `$m range Pauline Proc Roy Soc GW Scholz-Engel , while r.m.s. roughnesses $`s1\mu `$m. Then the corresponding average contact radii $`\overline{a}3`$$`10\mu `$m, an estimate which has been confirmed from the Dieterich-Kilgore visualizations Diet-Kilg .
Modern experimental improvements have revealed that surface profiles are in general more complex than was initially assumed by GW, exhibiting multiscale roughness distributions EBouchaud . They are often modelled as self-affine surfaces <sup>3</sup><sup>3</sup>3The statistical properties of a self affine surface $`z(x,y)`$ are invariant under the scaling transformation $`x\zeta x,y\zeta y,z\zeta ^{3D_f}z`$, where $`D_f`$ is the fractal dimension..
Persson Bofractal has recently proposed a theoretical treatment of the resulting contact problem. His calculation of the contact area between a flat compliant medium and a rigid self-affine surface follows a renormalization scheme for the stress distribution for profiles at growing stages of magnification $`\zeta `$. This function he proves to obey, for complete contact between the surfaces, a Fokker-Planck-like equation, with $`\zeta `$ playing the part of time. The real case of incomplete contact is then simulated via a boundary condition eliminating from the area of contact the regions sustaining tensile stress. This enables him to express the solution in terms of the elastic response of the compliant medium to imposed displacements specified all along its surface — while partial contact actually corresponds to a mixed situation of given displacements within contacts and given (null) normal stress elsewhere. It is therefore difficult to assess the influence of the ansatz on the main result of interest here, namely that the real contact area still obeys the Amontons proportionality $`\mathrm{\Sigma }_rW`$.
Note that Persson’s result depends crucially on the physical existence of a lower cut-off of the spatial scales, which is in practice of mesoscopic order. Indeed, even in the highly improbable case where the fresh surfaces would remain fractal down to the atomic scale, the nano-asperities would get plastified, and thus smoothed out, under the very high pressures they would experience upon contact.
This remark brings up the important question of the limit of validity of the pure elastic contact models. Clearly, when thinking in GW’s terms, plastic deformation of asperities starts coming into play when the average pressure $`\overline{p}`$ becomes comparable with the yield stress $`Y`$. From equations (8) and (9) , this condition can be expressed in terms of the dimensionless plasticity index
$$\psi =\frac{E}{Y}\sqrt{\frac{s}{R}}$$
(11)
which depends both on the topographic and the mechanical properties of the interface.
The elastic regime is limited to the region $`\psi <1`$, for which GW have shown that the fraction of plastified contacts is negligible. Persson’s more elaborate theory yields a comparable evaluation, with $`s`$ and $`R`$ in equation (11) now understood to be those on the large scale (upper space cutoff). This last result takes into account consistently the above-mentioned plastification on small lateral scales.
On the other hand, for $`\psi `$ larger than a few units, most microcontacts are in a state of full plastic deformation. Multiscale roughness loses relevance, and one may resort to the fully plastic version of the GW model. The matter which has flown being assumed to redistribute on a scale much larger than contact radii (no “piling up”), geometry imposes that the area of a contact involving an asperity with initial compression $`(zd)`$ is $`2\pi R(zd)`$. It bears the load $`\pi a^2H`$, with $`H`$ the hardness. Then, trivially
$$\mathrm{\Sigma }_r=\frac{W}{H}\overline{a}=\sqrt{2Rs}$$
(12)
One may therefore assert the robustness of three characteristics of multicontact interfaces, namely, to a good approximation :
– The real area of contact is proportional to the normal load, and independent of the apparent area $`\mathrm{\Sigma }_{app}`$.
– The average contact radius is load-independent, thus introducing a mesoscopic length scale $`\overline{a}`$ defined by the topography of the contacting surfaces.
– In a wide variety of cases, $`\overline{a}`$ is of micrometric order.
All these statements have been directly confirmed from the analysis, by Dieterich and Kilgore Diet-Kilg , of their optical images of several MCI between transparent solids.
For metals, $`E/Y10^2`$$`10^3`$. It would therefore take unrealistically small values of $`s/R`$, of order $`10^4`$$`10^6`$, for metal/metal MCI to be in the pure elastic regime. They are fully plastic, as anticipated by Bowden and Tabor.
The opposite case is that of elastomers, which are fully (visco)elastic up to strains $`1`$. Their interfaces, such as that between tire and road, when analyzed in the frame of Persson’s predictions Bofractal Bofracvisc should provide a test of his theory.
MCI involving polymeric glasses ($`E/Y10`$$`100`$) pertain to the intermediate regime, where $`\psi 1`$ : a fraction only of the microcontacts flow significantly. Although no quantitative theory is available, we believe that, in view of their robustness, the above-mentioned three main results of the GW model hold for all non-purely elastic MCI. Confirmation of this statement has been provided by experimental investigations of various interfacial properties, such as shear stiffness, which all exhibit the Amontons linear dependence on normal load (see Appendix A).
In fine, we can summarize the above results into the following simple statement for the friction coefficient of a MCI :
$$\mu =\frac{\sigma _s}{\overline{p}}$$
(13)
$`\overline{p}`$ is the average normal stress borne by the microcontacts. It is load-independent. Its expression in terms of the material properties (Young modulus and yield stress) and of the topographic surface characteristics depends upon the value of the plasticity index $`\psi `$ (eq.(11)). For $`\psi 1`$ (fully plastic regime) $`\overline{p}=H3Y`$. In the opposite, fully elastic case $`\psi 1`$, $`\overline{p}E\sqrt{s/R}H/\psi `$, and in the intermediate elasto-plastic range it is expected to extrapolate smoothly with $`\psi `$ between these two limits.
Let us insist that the interfacial shear strength $`\sigma _s`$ in expression (13), which is likely to depend upon the contact pressure, is load-independent for a given MCI, since such is $`\overline{p}`$ itself. This result is specific of multicontact interfaces, for which Amontons’s law is of purely geometric origin. It must be contrasted with the case of intimate single contacts such as those studied with the surface force apparatus (SFA) or extended gel glass ones (see Section III). For these systems, Amontons’s law, when observed, must be attributed to the pressure dependence of $`\sigma _s`$ itself Robbins-500 ans .
### II.2 Geometric age : a major state variable for MCI
#### II.2.1 Time dependence of the static threshold:
The classical description of solid friction states that there exists, for any given interface, a well defined static friction coefficient $`\mu _s=F_s/W`$ such that, as long as the applied shear force $`F<F_s`$, no sliding occurs.
This amounts to asserting that, at a depinning threshold $`F_s`$, the interface commutes from a purely elastic, reversible response to external shearing, to an irreversible, dissipatively flowing one. Most likely, and in view of unavoidable disorder and temperature effects, such a strong statement, which implies that the interface would undergo an abrupt unjamming transition, is only approximate. We come back to this point in $`\mathrm{\S }`$ II.C.4, where we discuss in detail the intrinsic difficulties and limits associated with the definition of such a threshold. Let us only state at this point that, in order to minimize ambiguities and make comparison between data meaningful, it is important to define carefully the protocol used to measure $`\mu _s`$.
A series of such stop and go experiments consists in :
(i) Preparing the initial interfacial state reproducibly, by sliding steadily at a chosen velocity $`V_{prep}`$, then stopping the pulling.
(ii) Waiting at rest for a given waiting time $`t_w`$ under a specified shear stress. This may either be self-selected by the system (natural arrest stress) or imposed at some value below this level.
(iii) Resuming loading at a prescribed velocity $`V_{load}`$, for example $`V_{load}=V_{prep}`$.
A typical shear force response is displayed on Figure 8. $`\mu _s`$ is then conventionally identified with the peak level. Experiments of this type have been perfored on a number of MCI. They reveal that $`\mu _s`$ is not a mere number characterizing a couple of solids, but a slowly increasing function of the waiting time. More specifically, $`\mu _s(t_w)`$ varies logarithmically over several decades of $`t_w`$ ranging up from about $`1`$ sec (the typical fast limit for such mechanical experiments). This behavior, displayed on Figure 9, holds for a wide variety of materials including metals Dokos , rocks Dieterich Marone , glassy polymers Pauline1 and paper Heslot .
As a MCI ages at rest, it strengthens logarithmically. Moreover, the slope
$$B=\frac{d\mu _s}{d(\mathrm{ln}t_w)}$$
(14)
is found, for all the above-mentioned materials, to be roughly (see Figure 9) on the order of $`10^2`$.
Such a “generic” behavior is striking, and suggests that it results from a robust physical mechanism.
#### II.2.2 Creep growth of the real area of contact
Tabor’s decomposition of the friction force (eq.(2)), together with the Greenwood-Williamson analysis, immediately raises the possibility that this strengthening might be attributable to slow growth of the real area of contact. Indeed, we have seen that, in general, the average normal stress $`\overline{p}`$ borne by the microcontacts is comparable with the bulk yield stresses of the contacting materials <sup>4</sup><sup>4</sup>4 More precisely, for an asymmetric interface, $`\overline{p}`$ is on the order of the yield stress of the softer material.. At this stress level, which prevails in a volume $`a^3`$ (with $`a`$ the microcontact radius), one expects materials to undergo plastic creep, resulting in the slow growth of microcontacts, hence of $`\mathrm{\Sigma }_r`$.
Confirmation of this idea was obtained by Dieterich and Kilgore Diet-Kilg who were able to measure directly the time evolution of $`\mathrm{\Sigma }_r`$ on their optical images of the microcontacts forming the MCI between two transparent glassy polymer blocks. They found that $`\mathrm{\Sigma }_r(t_w)`$ does grow logarithmically, at a rate compatible with that measured in microindentation experiments.
The question is then to ascertain whether or not this area growth is sufficient to explain that of $`\mu _s`$. This has been investigated in detail by Berthoud et al Pauline1 . They studied the temperature dependence of the static aging slope $`B`$ for symmetric MCI involving the polymeric glasses PMMA and PS (polystyrene) between room temperature and the vicinity of the bulk glass transitions. They analyzed their results in terms of a model for the growth of $`\mathrm{\Sigma }_r`$ due to Bréchet and Estrin BE .
This model schematizes the creep-induced growth of a microcontact as follows. Once a microcontact has been created, at time $`t=0`$, in a first stage of fast plastic flow, of duration negligible on the scale of the later evolution, the normal stress $`\sigma `$ sets to an ”initial” value $`\sigma _Y`$ of the order of the yield stress in the relevant geometry <sup>5</sup><sup>5</sup>5In uniaxial loading, $`\sigma _Y=Y`$, while for a sphere-sphere contact $`\sigma _Y=H3Y`$ Johnson .. Once this state is reached, plastic evolution continues via creep, the rate of which is given by an expression à la Nabarro-Herring :
$$\dot{ϵ}=\dot{ϵ}_0\mathrm{exp}\sigma /S$$
(15)
where $`\dot{ϵ}`$ is the compressive strain rate (treated as a scalar), $`S`$ the so-called strain-rate sensitivity of the flow stress, and $`\dot{ϵ}_0`$ a $`\sigma `$-independent Arrhenius factor. Since this creep law reflects a thermally activated process, both $`\dot{ϵ}_0`$ and $`S`$ are $`T`$-dependent.
Such plastic deformation occurs at constant volume, so that :
$$\dot{ϵ}=\frac{1}{a_0^2}\frac{d(a^2)}{dt}$$
(16)
with $`a`$ the contact radius, $`a_0`$ its initial value. Since $`\sigma =w/\pi a^2(1ϵ)w/\pi a_0^2+𝒪(ϵ^2)`$, with $`w`$ the normal load. Then, from eqs.(14) and (15), and approximating the MCI by a set of identical average GW contacts, the real area of contact evolves as :
$$\mathrm{\Sigma }_r(t)=\mathrm{\Sigma }_{r0}\left[1+m\mathrm{ln}\left(1+\frac{t}{\tau }\right)+𝒪(\mathrm{ln}^2)\right]$$
(17)
with
$$m=\frac{S}{\sigma _Y},\tau =\dot{ϵ}_0\frac{1}{m}e^{1/m}$$
(18)
Typically, at room temperature, for the systems of interest here, $`m`$ values are on the order of a few $`10^2`$. On the other hand, no reliable evaluation of the cross-over time $`\tau `$ can be made a priori, due not only to its exponential dependence on $`m`$ but even more to the lack of any precise data on $`\dot{ϵ}_0`$. All MCI studied up to now, except for polymer glass ones close to the glass temperature, exhibit a linear logarithmic growth of $`\mu _s`$, hence of $`\mathrm{\Sigma }_r`$, for waiting times above $`1`$ sec. This only enables us to state, at this stage, that in general $`\tau `$ is smaller than this limiting value. Then, assuming that the Tabor interfacial shear strength $`\sigma _s`$ exhibits no aging, one obtains from eq.(17), in the accessible range $`t\tau `$, for the log-slope of the static friction coefficient :
$$B=\mu _{s0}m$$
(19)
where $`\mu _{s0}=\sigma _s\mathrm{\Sigma }_{r0}/W=\sigma _s/\sigma _Y`$ is the static threshold base value at short times $`t\tau `$ <sup>6</sup><sup>6</sup>6Strictly speaking the value of $`\sigma _s`$ which appears in this expression should be that at the driving velocity (see $`\mathrm{\S }`$II.C.3 and Pauline2 . Since friction coefficients of MCI are usually a fraction of unity, one thus expects, at room temperature, $`B`$ to lie in the $`10^2`$ range.
The analysis of experimental data for $`B(T)`$ in ref.Pauline1 is quite intricate, in view of several difficulties concerned with :
– the fact that $`\tau `$ and, hence, $`\mu _{s0}`$, could be accessed in the realizable waiting time range only close to the glass transition temperatures $`T_g`$.
– the problems related with mapping the authors’ bulk data for $`\sigma _0`$ and $`S`$ onto the sphere-sphere geometry.
Berthoud et al were nevertheless able to show that expression (19) accounts semiquantitatively for the roughly tenfold increase of $`B`$ between its values, of order $`10^2`$ at room temperature as expected from expression (19), and the vicinity of the bulk $`T_g`$’s. However, it must be noted that eq.(19) is found to systematically underestimate the experimental data, the relative misfit increasing significantly on approaching $`T_g`$. This hints towards the fact that, while creep-induced growth of the real area of contact is responsible for most of the static strengthening, some aging of the interfacial strength is not completely ruled out : for example, close to $`T_g`$, such symmetric polymer glass contacts might exhibit partial “healing” due to interdiffusion of polymer chains.
Clearly, the Bréchet-Estrin model only describes plastic creep of the contacting asperities. However, quasi-logarithmic static aging was also observed by Ronsin et al Coey-Rons on rough rubber/rough glass MCIs. For such materials, contact area growth is obviously of viscoelastic origin. We show in Appendix B, following Hui et al Hui , that the Greenwood-Williamson model can be worked out explicitly for linear viscoelastic materials, with a very simple outcome. Namely, the GW results (eq.(8)) for the elastic MCI still hold formally, provided that the inverse Young modulus $`1/E`$ is replaced by the so-called creep compliance $`J(t)`$, which measures the delayed strain response to a unit instantaneous stress jump. Hence the real area of contact becomes a slowly increasing function of contact duration. A logarithmic increase is known to provide a reasonable approximation for $`J(t)`$ over time decades, for materials with a very wide spectrum of relaxation times such as rubbers.
#### II.2.3 Geometric age as a dynamical state variable
In the light of the foregoing analysis we will for the moment attribute the strengthening of MCI static thresholds to the sole time dependence of the area factor in Tabor’s expression. $`\mathrm{\Sigma }_r`$ thus becomes a function of the geometric age of the interface $`\varphi `$ by which $`t`$ should now be replaced in equation (17). More precisely, this age is defined as follows.
$``$ For a non-moving MCI, $`\varphi (t)`$ is simply the time which has been spent at rest at time $`t`$.
$``$ Consider now a MCI sliding at the constant velocity $`V`$. As motion proceeds, a given microcontact, once created (born) is gradually sheared until it slides, then disappears (dies) when the relative displacement between the partner asperities reaches a value on the order of a fraction of the contact diameter. Since, under constant normal load, the average number of microcontacts is conserved, any contact death is, on average, associated with the birth of of a new microcontact at an uncorrelated position. Since $`\mathrm{\Sigma }_r\mathrm{\Sigma }_{app}`$, we can safely consider that the newborn is formed between fresh asperities, which have not yet experienced creep <sup>7</sup><sup>7</sup>7This approximation, which neglects wear, is validated by the observed stability of the frictional characteristics of a MCI over slid lengths as large as a few $`10`$ cm unpub – a distance over which we can estimate the number of contact configuration renewals to $`10^{45}`$..
Contact renewal therefore limits the age of the MCI to the average lifetime of a given configuration of microcontacts, which can be written phenomenologically as :
$$\varphi _{ss}=\frac{D_0}{V}$$
(20)
In other words, motion interrupts aging, since interfacial configuration memory is destroyed after sliding the characteristic length $`D_0`$, which we expect to lie in the micrometer range.
Note that the larger $`V`$, the younger the steady sliding MCI, hence the smaller $`\mathrm{\Sigma }_r(\varphi )`$: geometric aging thus immediately appears as a candidate process for explaining the $`V`$-weakening behavior of the steady sliding dynamical friction coefficient $`\mu _d(V)`$ mentioned in Section I.
$``$ Let us now turn to non-steady sliding at the instantaneous velocity $`\dot{x}(t)`$. $`\varphi `$ is no longer time-independent, since it keeps track of geometric aging over the time necessary to slide the memory length $`D_0`$. In this sense, $`D_0`$ is the length over which memory of the history of motion is preserved.
The most simple phenomenological expression for $`\varphi `$ accounting for this behavior is :
$$\varphi (t)=_{t_0}^t𝑑t_1\mathrm{exp}\left[\frac{(x(t)x(t_1)}{D_0}\right]$$
(21)
where $`t_0`$ is the time at which the two solids were first brought into contact. Note that, in agreement, as needed, with the previously defined expressions. :
– for the static MCI equation (21) yields $`\varphi (t)=(tt_0)t_w`$,
– steady sliding corresponds to the limit $`t_0\mathrm{}`$, so that $`\varphi (t)`$ as defined from eq.(21) reduces to the constant $`\varphi _{ss}(V)=D_0/V`$,
Geometric age thus becomes a dynamical variable, coupled to the instantaneous velocity $`\dot{x}(t)`$ by the non-linear differential equation, equivalent to expression (21):
$$\dot{\varphi }=1\frac{\dot{x}\varphi }{D_0}$$
(22)
which was first proposed by Rice and Ruina RR .
That the memory of the interfacial state of sliding MCIs is indeed characterized by a length, a possibility first suggested by Rabinowicz Rabino in an often overlooked pioneer work, was established by Dieterich Dieterich on the basis of his systematic exploration of frictional transients following velocity jumps. In these experiments the slider is set into steady motion by driving it, through a “spring” of stiffness $`K`$ <sup>8</sup><sup>8</sup>8More precisely, $`K`$ is the equivalent stiffness of the driving stage plus slider system., at an initial velocity $`V=V_i`$. At $`t=0`$, $`V`$ is suddenly jumped to $`V_f`$, and one measures the force response. As can be seen on Figure 4, this exhibits a two-step transient.
(i) In a first, very rapid stage, for upward (resp. downward) velocity jumps, the instantaneous dynamic friction coefficient increases (resp. decreases).
(ii) This so-called “direct effect” is followed by a much slower monotonous variation in the reverse direction, ending at the level $`\mu _d(V_f)`$.
Dieterich found that the duration of the transient, which is dominated by that of stage (ii), scales as $`1/V_f`$. From this he identified a value of $`D_05\mu `$m for a granite/granite MCI. One may then conclude that the slow part of the transient corresponds to the gradual relaxation of the real contact area from its initial steady value $`\mathrm{\Sigma }_r(\varphi =D_0/V_i)`$ towards its final one $`\mathrm{\Sigma }_r(\varphi =D_0/V_f)`$.
In conclusion of this analysis, it appears reasonable to write tentatively Tabor’s expression as :
$$F=\mathrm{\Sigma }_r(\varphi )\sigma _s(\dot{x})$$
(23)
which assigns the whole state dependence to the area factor, while assuming that the shear strength $`\sigma _s`$ only depends on the instantaneous sliding velocity. The validity of this assumption will be discussed at length in the forthcoming sections.
Equations (23) and (22), together with Newton’s equation :
$$M\ddot{x}=F_{drive}(x,t)F(\varphi ,\dot{x})$$
(24)
where $`F_{drive}`$ is the externally applied driving force, provide a closed set of equations for the frictional motion of the slider once the form of the functional $`\dot{x}`$-dependence of $`F`$ has been identified.
### II.3 Junction rheology : gross features
#### II.3.1 Junctions at multicontact interfaces
Up to now, we have concentrated on the analysis of the contact geometry of MCI which, though an important prerequisite, does not yet touch upon our main question, namely that of the origin of frictional dissipation and of the detailed nature of the associated rheology described by the interfacial shear strength $`\sigma _s`$. Let us first try to identify the regions which are the seat of frictional dissipation.
We have restricted our definition of MCI to the normal load range such that microcontacts are numerous enough to form a good statistical set, but sparse enough for elastic interactions between them via the bulk materials to be negligible. This entails that we can now simply focus on the behavior of a single sheared microcontact. Such a unit (Figure 10) is constituted of the bulk of the two asperities and of an interfacial layer in which molecules from both surfaces have come into adhesive contact.
This layer, which we call the junction, has a disordered structure. This is obvious when the two solids are polymer glasses, since in this case the junction is formed by polymer tails and loops protruding from the amorphous bulks. In the case of crystalline bulks, structural and chemical surface disorder, which prevails under usual (non atomically planar and non ultra high vacuum) conditions together with plastic smoothing out of nano-scale roughness certainly lead to structural disorder.
Such highly defective structures, probably less dense than the bulks, are likely to result in lower resistance against plastic shear deformation. So, we will assume that the junction is a disordered quasi bidimensional medium with thickness $`h`$ in the nanometric range, where shear naturally localizes. We will show that this assumption is borne out by the analysis of experimental data on low velocity friction.
#### II.3.2 A threshold rheology
Stating the existence of a well-defined static friction coefficient can be formulated equivalently as the fact that, when the sliding velocity $`V0`$, $`\sigma _s`$ does not vanish, but tends towards a finite limit $`\sigma ^{}`$. That is, from a rheological viewpoint, frictional contacts behave as yield stress fluids.
The yield stress $`\sigma ^{}`$ therefore appears as the threshold beyond which the interfacial quasi–$`2D`$ solid flows plastically.
At lower stress levels, in this simple picture, the interface does not slide, it is pinned and responds to shear as a solid i.e., in principle, elastically. Due to the nanometric thickness of the junctions, the corresponding stiffness is always much larger than that of the microcontact-forming asperities themselves (see Appendix A) <sup>9</sup><sup>9</sup>9A rough estimate for the corresponding ratio is $`\overline{a}/h10^3`$, with $`\overline{a}`$ the average microcontact radius, $`h`$ the junction thickness. , it is therefore not accessible in MCI configurations.
It is common knowledge that plastic flow does not set in as an ideally sharp transition : at non-zero temperatures, thermal activation induces creep below the nominal yield stress, which can be attained, in principle, only by loading at extremely high rates. In other words, the smaller the loading rate, the more fuzzy the threshold. So, one expects that a constant shear load close below the static threshold should induce creep-like frictional sliding. This was indeed observed by Heslot et al Heslot . <sup>10</sup><sup>10</sup>10Note, however, that with MCI this creep is certainly amplified by the geometric rejuvenation (weakening) associated with sliding..
#### II.3.3 Beyond threshold : rate effects
##### Velocity jumps : the direct effect.
We already mentioned in $`\mathrm{\S }`$II.B.3 that the force response (see Figure 4) to a jump of the driving velocity from $`V_i`$ to $`V_f`$ exhibits, previous to the slow transient attributable to geometric age adaptation, a first much faster part. The associated force jump $`W\mathrm{\Delta }\mu _d`$ is positive (resp. negative) for $`V_f>V_i`$ (resp. $`V_f<V_i`$), hence the term “direct effect”. Dieterich has shown that :
$$\mathrm{\Delta }\mu _d(V_iV_f)=A\mathrm{ln}\left(\frac{V_f}{V_i}\right)$$
(25)
where $`A`$ is a constant for a given couple of materials. Its values lie in the $`10^2`$ range. This holds for all the MCI studied up to now, for velocities between about $`0.1`$ and $`100\mu `$m.sec<sup>-1</sup>.
Note that, as long as inertia is negligible, motion is quasi-static, i.e. after the jump at $`t=0`$ :
$$K\left(V_ftx\right)=W\left[\mu _d(\varphi ,\dot{x})\mu _d^{st}\left(V_i\right)\right]$$
(26)
$`x(t)`$ measures the position of some reference point on the slider along the pulling direction, $`x(0)=0`$, and $`\mu _d^{st}(V)=\mu _d(D_0/V,V)`$ is the friction coefficient for steady sliding at velocity $`V`$.
The peak value of the transient force ($`d\mu _d/dt=0`$) therefore occurs for $`\dot{x}=V_f`$. Velocity jump experiments are performed with as stiff as possible driving stages such that, as can be checked by direct measurements, the distance slept during the fast part of the transient be much smaller than the memory length $`D_0`$. Under such conditions the corresponding variation of geometric age from its initial value $`\varphi _i=D_0/V_i`$ is negligible, and the direct effect is fully attributable to the rate dependence of the shear strength $`\sigma _s`$.
From equation (23) one then expects that :
$$\mathrm{\Delta }\mu \left(V_iV_f\right)=\frac{\mathrm{\Sigma }_r(D_0/V_i)}{W}\left[\sigma _s(V_f)\sigma _s(V_i)\right]$$
(27)
That is, making use of equations (17) and (25)
$$\mathrm{\Delta }\sigma _s=\sigma _s(V_f)\sigma _s(V_i)=\frac{A\mathrm{ln}\left(V_f/V_i\right)}{\mathrm{\Sigma }_{r0}\left[1+m\mathrm{ln}\left(1+\frac{D_0}{V_i\tau }\right)\right]}$$
(28)
Since $`m10^2`$, the log term in the denominator can be neglected, and $`\mathrm{\Delta }\sigma _s\frac{AW}{\mathrm{\Sigma }_{r0}}\mathrm{ln}\left(\frac{V_f}{V_i}\right)=A\sigma _Y\mathrm{ln}\left(\frac{V_f}{V_i}\right)`$, from which we can write the following empirical expression for $`\sigma _s`$ :
$$\sigma _s(\dot{x})=\sigma _{s0}\left[1+\alpha \mathrm{ln}\frac{\dot{x}}{V}+𝒪(\mathrm{ln}^2)\right]$$
(29)
with
$$\alpha =\frac{A\sigma _Y}{\sigma _{s0}}$$
(30)
and $`\sigma _{s0}\sigma _s(V_0)`$ is the shear strength at the reference velocity $`V_0`$, which may be chosen anywhere in the range ($`0.1`$$`100\mu `$m/sec) where equation (25) holds.
The necessity of resorting to an expression involving a finite reference velocity is imposed by the formal divergence of the logarithm in the vanishing $`\dot{x}`$ limit. This divergence is of course unphysical and only points out the limits of the empirical approach. We will return to this point later, when discussing the behavior to be expected for $`\sigma _s(\dot{x})`$ in the very low and large velocity regimes.
Note finally that $`\sigma _{s0}/\sigma _Y`$ is the friction coefficient $`\mu _{s0}`$ appearing in equation (19), it is therefore roughly of order $`1`$ and, in order of magnitude :
$$\alpha A10^2$$
(31)
##### Steady sliding friction coefficient.
The velocity jump experiments initiated by Dieterich have played a pioneering role to evidence the state and rate character of MCI friction and to separate clearly the age (state) and rheologic (rate) contributions to the dynamic friction coefficient $`\mu _d`$. However, it is the studies of the velocity and temperature dependences of $`\mu _d(V)`$ under steady sliding conditions which have led to a systematic confirmation of the validity of expression (29) and opened the way to its more quantitative explicitation.
The typical experimental $`V`$-dependence of $`\mu _d`$ is illustrated on Figure 11, on the cases of granite, paper and PMMA symetric MCI at room temperature. It is seen that :
(i) $`\mu _d`$ is $`V`$-weakening ($`d\mu _d/dV<0`$) in the whole low velocity regime explored in such experiments.
(ii) Its variations are quasi-logarithmic, except in the higher $`V`$ range ($`V100\mu `$m/sec) where it exhibits a saturating trend. For all materials, the log-slope $`\beta _d=d\mu _d/d(\mathrm{ln}V)`$ lies, once more, in the $`10^2`$ range.
Such results can be confronted with the predictions obtained from equations (23), (17) and (29) which yield, in steady motion where the geometric age $`\varphi =D_0/V`$ :
$`\mu _d(V)`$ $`=`$ $`{\displaystyle \frac{\sigma _{s0}}{\sigma _Y}}\left[1+m\mathrm{ln}\left(1+{\displaystyle \frac{D_0}{V\tau }}\right)+\alpha \mathrm{ln}{\displaystyle \frac{V}{V_0}}\right]`$ (32)
$`=`$ $`{\displaystyle \frac{\sigma _{s0}}{\sigma _Y}}+B\mathrm{ln}\left(1+{\displaystyle \frac{D_0}{V\tau }}\right)+A\mathrm{ln}{\displaystyle \frac{V}{V_0}}`$
where we have taken advantage of the smallness of $`m`$ and $`\alpha `$ to neglect the $`𝒪(\mathrm{ln}^2)`$ terms, and made use of equations (19) and (30).
For $`VD_0/\tau `$, equation (32) gives for the log-slope of $`\mu _d`$ the value :
$$\beta _d=BA$$
(33)
The fact that $`\beta _d`$ is positive entails that, at room temperature, $`B>A`$. Moreover, the $`\beta _d`$ data are found to be quantitatively compatible with the available independently measured $`B`$ and $`A`$.
On the other hand, the saturating trend at the larger $`V`$’s appears as resulting from that of the real area of contact in the small age limit, which becomes relevant for $`VD_0/\tau `$. More precisely, equation (32) predicts a broad minimum of $`\mu _d`$, which should occur for $`V_{min}=(D_0/\tau )(BA)/AD_0/\tau `$.
That is, expression (32) predicts that, beyond $`V_{min}`$, dynamic friction should become velocity-strengthening. This natural physical consequence of the nature of geometric aging, which has been checked experimentally on the case of paper Heslot , has been overlooked up to now in most studies of fault dynamics. It is important in view of its bearing upon the nature of the sliding dynamics.
We have seen in $`\mathrm{\S }`$II.B.2 that parameter $`B`$ varies with temperature in a way which is explained by its physical origin – the thermally activated creep of $`\mathrm{\Sigma }_r`$. On the same polymer glass systems, Berthoud et al have shown Pauline2 that $`A`$ increases quasi-linearly <sup>11</sup><sup>11</sup>11Note, however, that the T-range in these experiments is small. when increasing $`T`$ up to the vicinity of $`T_g`$. Blanpied et al Blanpied and Nakatani Nakatani have investigated the behavior of, respectively, $`\beta _d`$ and $`A`$ for granite over wide $`T`$-ranges, from room temperature to $`800^{}C`$ (still well below the melting point). $`\beta _d`$ was found to decrease with increasing $`T`$, enough for hot granite to become velocity-strengthening – a point of importance in the context of deep seismicity. $`A`$, on the contrary, increases in a quasi-linear manner. We defer a more detailed analysis and interpretation of these results, in particular those concerned with $`A`$, to $`\mathrm{\S }`$ II.C.4 below, where we propose a physical model for the rate effect.
#### II.3.4 Threshold rheology as a signature of multistability
The primary feature of frictional rheology is the existence of a threshold. That is, the force $`F`$ needed for an interface to slide at velocity $`V`$ remains finite for vanishing $`V`$’s. This is obviously at odds with the standard picture of dissipation at low rates, namely : if the externally imposed rate of shear is much smaller than the internal relaxation ones, the system evolves quasi-adiabatically, so that dissipation vanishes linearly with $`v`$. As first formulated by M. Brillouin in 1904 Brillouin , in order for $`F`$ to exhibit a finite threshold, it is therefore necessary that, however slow the drive, the sheared medium evolves through a succession of adiabatic adaptation periods interspeded with fast instability events. Each of these events then corresponds, even at vanishing $`v`$, to a finite energy loss, thus accounting for the threshold behavior.
In modern language, a threshold rheology implies multistability, as has been developed at length in various fields such as magnetic hysteresis, wetting dynamics, charge density wave and type-II superconducting transport Fisher .
##### A toy model for junction rheology
Let us first illustrate this idea here on a toy model of a sheared junction <sup>12</sup><sup>12</sup>12This model was first developed by Caroli and Nozières CNTrieste . Although their formal results directly apply here, let us point out that the original physical interpretation – that the elementary mechanically unstable units were the interasperity microcontacts as a whole – was not correct CNinter . Let us insist again that the relevant instabilities do not occur on the micrometric scale but, within the inter-asperity junction, on the nanometric scale.. We represent it, as shown on Figure 12, as formed by a set of identical blobs of elastic matter attached to two stiff plates with vertical separation $`h`$, and area $`\mathrm{\Sigma }`$. The blobs are randomly distributed along the surfaces of these plates, and form $`n`$ contacts per unit area. They compress each other when making contact, in which case they interact via the (repulsive) pinning potential $`U(x)`$, where $`x`$ measures their relative position along the shearing direction (Figure 12). As $`x`$ is increased, some contacts are being destroyed while others are created at random positions, so that $`n`$ remains constant. We assume, for simplicity, that horizontal displacements are purely one-dimensional.
The system is sheared by imposing to the upper plate the displacement $`X`$ with respect to the lower one. Let $`k`$ be the shear stiffness of the system formed by two contacting blobs. Under the action of the pinning potential, the center of each blob contacting surface undergoes an elastic shear displacement $`u/2`$ along $`Ox`$, so that the horizontal separation $`x=X+u`$. The total energy of a contact is thus :
$$(X,u)=U(X+u)+\frac{1}{2}ku^2$$
(34)
Assume for the moment that $`T=0`$. Then, for fixed $`X`$, $`u`$ sets at the value $`u^{}(X)`$ which minimizes $``$, such that :
$$U^{}(X+u^{})=ku^{}$$
(35)
Two cases must then be distinguished, depending on whether the instantaneous equilibrium thus defined is unique (monostable contact) or not (bistable contact).
Monostable contact:
If $`k>k_0=max(U^{\prime \prime })U_0/a^2`$, – with $`a`$ the range of $`U`$, $`U_0`$ its maximum value – the solution of eq.(35) is unique whatever the value of the reference coordinate $`X`$ (see Figure 13).
Assume that we impose a constant external shear velocity $`\dot{X}=V`$. As $`X`$ increases, $`u`$ nearly follows its instantaneous equilibrium value $`u^{}(X)`$, and the contact energy evolves as $`^{}(X)(X,u^{}(X))`$. Indeed, the excess elastic energy is dissipated out of the contact region of size $`a`$ via acoustic radiation, at the characteristic rate $`\tau _{eq}^1c/a`$, with $`c`$ a sound velocity, hence much larger than the shear rate $`\dot{\gamma }V/hV/a`$.
The instantaneous pinning force exerted by the contact on the upper plate is simply :
$$f_p=\frac{d^{}}{dX}=\frac{}{X}|_{u=u^{}}\frac{}{u}|_{u=u^{}}\frac{du^{}}{dX}$$
(36)
i.e., taking advantage of eq.(35)
$$f_p(X)=U^{}\left(X+u^{}(X)\right)$$
(37)
and the net work spent to sweep through the contact :
$$w=_{\mathrm{}}^{\mathrm{}}dXf_p(X)=U(X+u^{})\left)\right|_{\mathrm{}}^{\mathrm{}}=0$$
(38)
The total instantaneous friction force $`F(V)`$ for a junction of area $`\mathrm{\Sigma }`$ is the sum of the pinning forces $`f_p`$ over the $`n\mathrm{\Sigma }`$ randomly positioned contacts of transverse range $`a`$, so that :
$$F=\frac{n\mathrm{\Sigma }}{a}_{\mathrm{}}^{\mathrm{}}𝑑X\left(f_p(X)\right)=n\mathrm{\Sigma }\frac{w}{a}=0$$
(39)
As could be expected, for an elastically monostable junction, frictional dissipation vanishes when $`V0`$.
Bistable contact :
In the opposite case where $`k<k_0`$, in a finite range of values of $`X`$ equation (35) has three solutions $`u_i^{}(X)`$ (Figure 14), the two extreme ones $`u_{1,2}^{}`$ corresponding to minima $`_{1,2}^{}(X)`$ of the energy $`(X,u)`$, while $`u_0^{}`$ is associated with a maximum. As $`X`$ is slowly increased from $`\mathrm{}`$, the instantaneous $`u^{}`$ evolves continuously with $``$, i.e. follows branch $`u_1^{}`$ (Figure 14b). The $`_1^{}`$ equilibrium gradually evolves from fully to meta-stable, while the barrier $`\mathrm{\Delta }_b(X)`$ separating the two basins of $``$ decreases until the spinodal limit $`X=X_+`$ where it vanishes. At this bifurcation point, branch $`(1)`$ terminates, $`u_1^{}`$ becomes unstable, and the only choice for the system is to relax towards the, now single, lower minimum $`_2^{}`$. The corresponding finite energy difference $`\mathrm{\Delta }_{12}`$ is, again, lost via acoustic radiation, i.e. instantaneously on the scale of the drive. The pinning force $`f_p`$ exhibits the hysteresis cycle common to such bistable systems (Figure 14c). The work for sweeping through the contact (the hatched area noted ($`1`$) on Figure 14c) $`w=\mathrm{\Delta }_{12}`$ is now finite.
The friction force F on the whole junction corresponds to a random distribution of $`X`$, i.e. to a uniform population $`P(X)dX=n\mathrm{\Sigma }dX/a`$ of branch ($`1`$) of the hysteresis cycle (Figure 15a), so that $`F_d=n\mathrm{\Sigma }w/a`$. Multistability results in a threshold rheology and, as long as $`Vc`$, the dynamic friction force is velocity independent <sup>13</sup><sup>13</sup>13The corrections resulting from imperfect adiabatic adaptation are easily shown to be of relative order $`(V/c)^{2/3}`$ CNTrieste .
Let us now assume that, starting fom this steady sliding regime, we suddenly suppress the external shear force. The interblob contacts must react so as to bring the system to global mechanical equilibrium, where the sum of the pinning forces vanishes. The only way for the junction to realize this situation is via a recoil $`\mathrm{\Delta }`$ of the upper plate : each individual reference coordinate $`X_i`$ decreases to $`(X_i\mathrm{\Delta })`$, yielding a new distribution $`P_{eq}(X)`$ such that the two shaded areas on Figure 15b be equal. $`P_{eq}`$ has a discontinuity at the Maxwell plateau of the force cycle. If, starting from this state, we now gradually increase the external shear force $`F`$, $`P_{eq}`$ rigidly shifts to the right by $`\delta X(F)`$. As long as $`\delta X<\mathrm{\Delta }`$ , no irreversible “jump” occurs, the displacement is therefore fully reversible : the junction is in an elastic regime. This terminates when the discontinuity of the shifted distribution reaches the spinodal limit $`X_+`$, i.e. when $`\delta X=\mathrm{\Delta }`$, i.e. $`F=F_d`$. In this stop and go scenario, the static threshold force is equal to the dynamic one.
However, due to multistability, global equilibrium can be realized by a huge number of configurations of the interblob contacts. Let us assume that we create the unsheared junction by bringing the two plates into contact. Depending on the details of this “mechanical quench”, each newly formed interblob contact may correspond to a state on either of the two branches of the force cycle, the only condition to be satisfied being that :
$$𝑑X\underset{i=1,2}{}P_i(X)f_{pi}(X)=0$$
(40)
where $`i`$ labels the branches of the force cycle. The elastic regime ends when the edge of the shifted $`P_1`$ reaches $`X_+`$. This point, where irreversible sliding starts, is the static threshold. Contrary to dynamic friction, it is not an intrinsic property of the junction but depends on its past history. If, for example, $`P_1(X_+)0`$, dissipation starts at vanishingly small shear. This has been illustrated in ref. jamming where it was shown that the dispersion in measured values of $`\mu _s`$ for a PMMA MCI is strongly reduced by preparing the initial state in a controlled way, i.e. by sliding before repose.
Rate effect at finite temperature :
Consider a bistable contact, swept along branch $`1`$ of the force cycle, i.e. sitting at the left minimum of the $`(u)`$ potential (Figure 14. As soon as $`T0`$, thermal noise is able to activate jumps over the barrier $`\mathrm{\Delta }_b(X)`$, i.e. to provoke premature jumps onto branch $`2`$ before the spinodal limit $`X_+`$ is reached. For fixed $`X`$, the jumping rate is the Kramers one :
$$\frac{1}{\tau (X)}\omega _a\mathrm{exp}\left[\frac{\mathrm{\Delta }_b(X)}{k_BT}\right]$$
(41)
with $`\omega _a`$ an attempt frequency, here typically $`c/a`$. As $`X`$ approaches $`X_+`$, $`\mathrm{\Delta }_b`$ decreases smoothly to zero, and $`\tau (X)`$ decreases exponentially up to a cut-off fixed by viscous (acoustic) dissipation. In practice, only close below the spinodal limit does the activated jumping rate become non negligible. The steady sliding distribution $`P`$ is now controlled by the competition between :
– the advection imposed by the external drive of $`X`$ at velocity $`V`$,
– the thermally induced jumps from branch $`1`$ to branch $`2`$.
This is expressed by the evolution equation <sup>14</sup><sup>14</sup>14This expression assumes implicitly that the drive is slow enough for the instantaneous jumping rate to be that for the non-advected system:
$$\frac{P}{t}=V\frac{P}{X}\frac{P}{\tau (X)}$$
(42)
the steady solution $`P_{st}`$ of which reads :
$$P_{st}(X)=Const\mathrm{exp}\left[_{\mathrm{}}^X\frac{dX^{}}{V\tau (X^{})}\right]$$
(43)
Due to the exponential variation of $`\tau `$, $`P_{st}`$ is quasi- completely depleted between $`X_+`$ and a cut-off $`X_c`$ below which activation plays a negligible role. The friction force is thus smaller than its $`T=0`$ limit. Moreover, the larger $`V`$, the less time the contact spends in the vicinity of a given $`X`$, the less the probability for it to jump prematurely. Hence, $`X_c`$ increases with the driving velocity, and the friction force increases accordingly : thermal noise results in a velocity-strenghtening rheology.
This rate dependence can be evaluated analytically as long as the the barrier height evolution for $`XX_c`$ can be approximated simply. For not too large $`V`$, where variations are linear, it is found (see ref. Bo , Chapter 11) that :
$$F(V)=F(V_0)\left[1+\alpha _{th}\mathrm{ln}\left(\frac{V}{V_0}\right)\right]$$
(44)
with $`V_0`$ is a reference velocity in the above-mentioned range of validity,and :
$$\alpha _{th}=\frac{k_BT}{\overline{\sigma }v_{act}}$$
(45)
and $`\overline{\sigma }nw/a`$ is on the order of the $`T=0`$ frictional stress. At larger velocities, when $`X_c`$ comes very close below $`X_c`$, $`\mathrm{\Delta }_b(X_cX)^{3/2}`$, and the rate dependence of $`F`$ commutes from log-linear to $`(\mathrm{ln}V)^{2/3}`$ CNTrieste \- a behavior which has been observed by Sills and Overney Overney in a nanoscale friction experiment on glassy polystyrene performed with an atomic force microscope.
##### From the toy model to the real junction :
When noticing that the functional form predicted by the toy model (eq.(44)) for the interfacial rheology does fit the empirical expression (eq.(29)) deduced from experiments, one is led to go one step further. That is, let us take for a moment the toy model at face value and compare experimental results for coefficient $`\alpha `$ (eq.(30)) with $`\alpha _{th}`$ (eq.(45)). As mentioned above, Nakatani Nakatani found that, for granite, $`\alpha `$ increases quasi-linearly with temperature, in agreement with equation (45) – a result confirmed for PMMA in the more restricted velocity range investigated in Pauline2 . These authors have then been able to deduce values of the volume $`v_{act}`$. In both cases, they find it to be on the order of a few nm<sup>3</sup>.
This naturally leads to conclude that, inspite of its crudeness, the toy model does capture the main physical mechanisms responsible for the frictional rheology of MCI, namely :
(i) When sheared, the highly confined junctions between asperities are the seat of mechanical instabilities, each of which primarily affects a small region containing a few atomic (or, for polymer glasses, monomer) units.
(ii) The strengthening, logarithmic, rate effect can be assigned to premature “flips” of these clusters induced by thermal noise.
Feature (i)) is fully consistent with our description of the junction as a quasi-$`2D`$ disordered medium, solidified under the high confinement conditions imposed by the bulks of the contacting asperities, but mechanically weaker than these bulks – probably due to a higher fraction of quenched free volume. Indeed, a number of numerical studies, pioneered by Argon et al Argon (see also Falk and Langer Falk and Malandro and Lacks Lacks ) of sheared glasses, both polymeric and molecular, at temperatures $`T_g`$, have shown that elastic dissipation in these systems occurs essentially via sudden collective rearrangements of clusters of typically nanometric volume. These clusters, termed by Falk and Langer shear transformation zones (STZ), are randomly located <sup>15</sup><sup>15</sup>15 In the absence of shear localization and their average density is constant at constant shear rate $`\dot{\gamma }`$ <sup>16</sup><sup>16</sup>16To which extent and how this density depends on $`\dot{\gamma }`$ is an important though still unsettled question. : on average, once a STZ has flipped into a lower energy local equilibrium state, another one appears at an uncorrelated position. That is, we may interpret the interblob bistable contacts of the toy model as a sketchy representation of STZ’s, the volume $`v_{act}`$ providing a rough evaluation of the average zone volume.
Feature (ii) opens onto a much more difficult, still essentially unsolved question. In the toy model, the bistable units are completely decoupled, and their contributions to the friction stress are simply additive. The real junction is, of course, dense everywhere, so that each STZ must be understood as embedded in a quasi-$`2D`$ elastic medium, and coupled to the adjacent deformable asperity bulks which, though stiffer, do transmit stresses. The flip of a STZ is therefore akin to the tranformation of an Eshelby inclusion Eshelby : as already emphasized in Argon , it gives rise to a multipolar force signal which deforms the surrounding medium Ajdari , hence resulting in loading steps on the other STZ’s, which were “on their way towards flipping”. The randomness of these signals in time and space is the source of the so-called dynamic noise, which acts in parallel with the thermal one. Elastic coupling being long ranged, dynamic noise is likely to trigger cascades (avalanches) of correlated flips, which are not accounted for in mean field approximations – hence the difficulty of this class of problems Fisher .
Moreover, clearly, the larger the imposed strain rate $`\dot{\gamma }`$, the more frequent the flips hence, roughly speaking, the larger the effective strength of the dynamic noise – which must vanish in the vanishing $`\dot{\gamma }`$ limit. So, if these effects can be modelled in terms of an effective temperature, this must be a growing function of $`\dot{\gamma }`$. Up to now, no theory of this effect is available, and it is therefore not taken into account by the recent phenomenological theories of plasticity of amorphous media : the STZ-based one Falk and the soft glass rheology of Sollich et al. Sollich . Qualitatively, one expects such rheologies to exhibit a crossover between two limiting regimes : a thermal noise-controlled one at low $`\dot{\gamma }`$ and/or high $`T`$, a dynamic noise-controlled one in the opposite conditions.
Our above analysis of friction leads us to conclude that the rheology of MCI junctions is controlled , at and above room temperature, by thermal noise, in the investigated range of $`\dot{\gamma }=V/h`$ i.e. typically $`\dot{\gamma }10^5`$.
Persson has studied numerically a model which represents the junction, as sketched on Figure (16), as a dense set of elastically coupled pinned blocks driven by a rigid plate. He found (see Bo , chap. 11) that dynamic noise effects were negligibly small, thus justifying a thermal noise-controlled mean field rheology (equation (44)). Note however that the smallness of dynamic noise effects in his model is likely to result from the infinite stiffness of the confining plates, which should lead to exponential screening of elastic couplings.
Our analysis of MCI rheology can therefore be summed up by the artist’s view of Figure (17). The ”threshold-plus-logarithmic” behavior fitted by the toy model results should therefore be understood as corresponding to an intermediate regime, with limits on both the high and low velocity sides.
In view of the current high activity in the field of soft glass rheology, one may be optimistic about the emergence of theoretical predictions about the position of the high-$`V`$ crossover towards the non thermal regime. However, it must be pointed out that MCI friction is certainly not a good tool for experimental investigation of this question since, at the necessary sliding velocities, in the mm/sec range or larger, one would have to cope with the tricky problems related with e.g. frictional heating.
The opposite low-$`V`$ limit is essentially academic. Let us come back to the toy model : as the stress decreases, the cut-off $`X_c`$ of the contact distribution approaches form above the position corresponding to the Maxwell plateau, and the rate of back jumps from branch $`2`$ to $`1`$ becomes non negligible and comparable with the $`12`$ one. At the same time, the interbranch barrier $`\mathrm{\Delta }_b`$ increases to $`\sigma _Yv_{act}`$, making jumps extremely rare. This results in the classical Eyring ($`sinh`$) behavior for $`\sigma _s(V)`$, ending in an extremely fast quasi-linear drop to zero, with slope $`\eta /h`$, where $`\eta `$ is the viscosity of the glassy junction. The crossover velocity $`V^{}`$ is in practice much to small to be observable.
### II.4 Sliding dynamics of a MCI :
#### II.4.1 The Rice-Ruina friction law :
We can now sum up the above phenomenological analysis into a model of MCI dynamical friction which reads, with the help of equations (17), (22),(23) and ((29) :
$$F(\varphi ,\dot{x})=\sigma _{s0}\mathrm{\Sigma }_{r0}\left[1+m\mathrm{ln}\left(1+\frac{\varphi }{\tau }\right)\right]\left[1+\alpha \mathrm{ln}\left(\frac{\dot{x}}{V_0}\right)\right]$$
(46)
$$\dot{\varphi }=1\frac{\dot{x}\varphi }{D_0}$$
(47)
We saw that both $`m`$ and $`\alpha `$ are typically of order $`10^2`$, so that it appears reasonable to neglect the $`m\alpha \mathrm{ln}^2`$ term in expression (46) for the friction force. Moreover, we saw that the short time cutoff $`\tau `$ associated with the early stage of area creep becomes relevant only for velocities $`D_0/\tau `$, on the order of a few hundred $`\mu `$m/sec at least. For the moderately accelerated instationary dynamics which we will analyze below, in order of magnitude $`\varphi D_0/\dot{x}`$ so that, in the low velocity regime considered here, $`\varphi /\tau 1`$. Equation (46) then reduces to the following expression for the dynamic friction coefficient :
$$\mu _d(\varphi ,\dot{x})=\mu _d(V_0)+B\mathrm{ln}\left(\frac{\varphi V_0}{D_0}\right)+A\mathrm{ln}\left(\frac{\dot{x}}{V_0}\right)$$
(48)
where the reference velocity $`V_0`$ can be chosen at will in the above-mentioned velocity range. The constants $`B`$ and $`A`$ are defined in equations (19) and (30).
Equations (47) and (48) constitute the Rice-Ruina (RR) model, proposed by these authors in 1983 on the basis of Dieterich’s experiments.
#### II.4.2 The RR dynamics of a driven block :
The question then immediately arises of analyzing and testing experimentally the predictions of the model in terms of sliding dynamics. First of all, does it correctly describe the stick-slip (SS) oscillations of a driven block, and their disappearance at high enough stiffness?
Consider the system depicted on Figure 1. Let us assume for the moment that the sliding velocity is uniform along the interface (homogeneous sliding), i.e. that the block has a single degree of freedom, the position $`x(t)`$ of e.g. its center of mass. Its equation of motion reads :
$$M\ddot{x}=K\left(xx_0(t)\right)W\mu _d(\varphi ,\dot{x})$$
(49)
where $`\mu _d`$ and $`\varphi `$ are specified by equations (48) and (47), and $`(xx_0(t))`$ is the instantaneous elongation of the driving spring.
At any pulling velocity $`V`$ there always exists a steady sliding solution, namely :
$$\dot{x}=V\varphi =\frac{D_0}{V}x(t)x_0(t)=\frac{W}{K}\mu _d(\frac{D_0}{V},V)$$
(50)
In view of the non-linearities of the friction law, one must wonder about its dynamic stability, i.e. perform a standard linear stability analysis : setting $`\dot{x}=V+\delta \dot{x}(t)`$, $`\varphi =D_0/V+\delta \varphi (t)`$, one linearizes the dynamical equations in $`\delta \dot{x}`$, $`\delta \varphi `$. Thanks to the time invariance of the basic state, the solutions for them are of the form $`Const\mathrm{exp}(i\mathrm{\Omega }t)`$. Steady sliding is stable (resp. unstable) when $`Im\mathrm{\Omega }>0`$ (resp. $`<0`$).
This calculation is performed in Appendix C. It shows that, for given values of $`M`$, $`W`$ and $`V`$, steady motion is stable for $`K>K_c`$, with the critical stiffness given by:
$$K_c\frac{D_0}{W}=\left(\mu _\varphi \mu _{\dot{x}}\right)\left[1+\frac{MV^2}{WD_0\mu _{\dot{x}}}\right]$$
(51)
At the bifurcation point ($`K=K_c`$) $`\mathrm{\Omega }`$ is pure real and has the value :
$$\mathrm{\Omega }_c=\frac{V}{D_0}\sqrt{\frac{\mu _\varphi \mu _{\dot{x}}}{\mu _{\dot{x}}}}$$
(52)
That is, the corresponding bifurcation is of the Hopf type, signalling that for $`KK_c`$ motion should become oscillatory.
In equations (51), (52), $`\mu _\varphi =\mu _d/(\mathrm{ln}\varphi )`$, $`\mu _{\dot{x}}=\mu _d/(\mathrm{ln}\dot{x})V`$, both derivatives being evaluated at $`\varphi =D_0/V,\dot{x}=V`$. The RR expression results in $`\mu _\varphi =B`$, $`\mu _{\dot{x}}=A`$. In the experiments aiming at characterizing the SS bifurcation, the block was sliding under its own weight : $`W=Mg`$. Then with $`D_01\mu `$m$`/sec`$, $`A10^2`$, the inertial correction in the second factor of the r.h.s. of equation (51) is of order $`10^5(V_{\mu m/s})^2`$. So, inertia is negligible for $`V100\mu `$m$`/sec`$. Then the RR model predicts that the critical stiffness
$$K_c=\left(BA\right)\frac{W}{D_0}$$
(53)
should be velocity-independent, while the critical pulsation
$$\mathrm{\Omega }_c=\frac{V}{D_0}\sqrt{\frac{BA}{A}}$$
(54)
The SS bifurcation has been studied on symmetric MCI involving paper Heslot , glassy PMMA and PS Pauline2 . In these experiments, $`K`$ was kept constant. Then, when $`M`$ was increased at constant $`V`$, the initially steady motion was observed to bifurcate at a value $`(K/M)_c`$ beyond which it develops oscillations whose amplitude grows continuously as $`K/M`$ decreases, until true sticking phases appear (Figure 18). Tracking the bifurcation at various $`V`$’s leads to the stability diagram in the plane ($`K/M,V`$) of the control parameters, a typical example of which is displayed on Figure 19. It is seen that $`(K/W)_c`$ is not strictly constant, but decreases slowly with increasing $`V`$ at a rate on the order of $`10`$ per cent per decade. On the other hand, $`V/\mathrm{\Omega }_c`$ was measured to be constant up to experimental accuracy <sup>17</sup><sup>17</sup>17 The poor accuracy on $`\mathrm{\Omega }_c`$ was due in particular to the difficulty of extrapolating measurements necessarily performed at finite amplitude to the continuous bifurcation while, in this case, non linearities develop very fast HopfNL . Heslot . So, the predictions of the RR model appear qualitatively satisfactory.
The parameters $`A,B,D_0`$ of the model are then evaluated as follows : measurements of $`d\mu _d/d(\mathrm{ln}V)`$ yield $`(BA)`$, $`D_0`$ is obtained by equation (53) from the value of $`(K/W)_c`$ at some reference velocity – e.g. $`1\mu `$m$`/sec`$ – then $`A`$ is determined with the help of equation (54). For example, for a PMMA MCI, one thus obtains Cochard : $`A=\mathrm{1.2\hspace{0.17em}10}^2\pm \mathrm{2\hspace{0.17em}10}^3,B=\mathrm{2.3\hspace{0.17em}10}^2\pm \mathrm{2\hspace{0.17em}10}^3,D_0=0.4\pm 0.04\mu `$m. Such parameter values are fully compatible with those deduced from static aging and velocity jump experiments.
In order to test more thoroughly the validity of the phenomenological model, we have analyzed whether it correctly predicts a few other salient features of MCI block dynamics, among which :
(i) Non-linear development of SS-like oscillations close to the bifurcation: it must be mentioned first that the RR model as such (constant $`\mu _\varphi =B`$ and $`\mu _{\dot{x}}=A`$) results in a highly singular behavior for $`KK_c`$ Ranjith , namely unbounded growth of the velocity at finite time. This unphysical behavior, in contradiction with the observation of a continuous bifurcation, can be traced back to the simplifying assumption that $`\mu _\varphi `$ and $`\mu _{\dot{x}}`$ are mere constants, which is also responsible for the fact that the model fails to accounts for the non zero slope of the $`K_c(V)`$ bifurcation curve. Baumberger et al HopfNL relaxed this assumption by assigning the slope of $`K_c(V)`$ to a $`(\mathrm{ln}^2\varphi )`$ correction to expression (48) for the friction coefficient. A standard perturbation expansion then results <sup>18</sup><sup>18</sup>18The expansion parameter is found to be $`ϵA/\mu _2`$ with $`\mu _2=[^2\mu _d/(\mathrm{ln}x)^2]`$, so that for the RR model ($`\mu _2=0`$) the expansion explodes. For PMMA, $`\mu _2`$ is measured to be $`10^3`$, hence $`ϵ10`$ : non linearities develop fast in the SS regime.. This extension of the RR constitutive law accounted very satisfactorily for the growth of the oscillation amplitude and of the frequency shift with $`(K_cK)`$.
(ii) Creep-like sliding motion following cessation of the drive: As can be seen on Figure 8 once the drive at velocity $`V`$ is stopped, the block continues to slide while slowing down until it stops at a finite force level. Baumberger and Gauthier TBrelax analyzed the dependence of such relaxation curves on the control parameters $`K/W`$, $`V`$. They showed that the RR model predicts that the total distance slipped before arrest should be $`V`$-independent, while experiments show a $`100\%`$ increase on one decade of $`V`$. Again, this discrepancy is cured by including into equation (48) the same $`\mathrm{ln}^2\varphi `$ correction as for non linear effects close to the bifurcation. An example of the resulting fit of relaxation data is shown on Figure 20.
These two examples indicate the limits to be assigned to the rate and state model in its simplest version. While, clearly, the RR law does capture the essential features of MCI friction – namely geometric aging and dissipation governed by thermally-assisted mechanical instabilities on the nanometer scale – it is insufficient to account for finer dynamical features involving finite departures from stationary sliding. As described above, phenomenological corrections to the expression of $`\mu _d`$ permit to extend the validity of the model. However, such extensions must be regarded with proper caution : clearly, the non-zero slope of the $`K_c(V)`$ bifurcation curve indicates that an extension of expression (48) is needed. But this leaves open the question of the respective weights of the various possible corrections which we have mentioned when deriving expression (48). Not only is a $`\mathrm{ln}^2\varphi `$ term possible, but also corrections of the form $`(\mathrm{ln}\dot{x})^2`$ and $`(\mathrm{ln}\varphi )(\mathrm{ln}\dot{x})`$ as well as higher order ones. Moreover, corrections associated with the short age time cut-off $`\tau `$ might become relevant. This would mean introducing many more parameters, in principle to be fitted from strongly non-linear dynamical features. As shown by the detailed analysis of the response to to large amplitude normal force oscillations oscnorm Cochard , such fits become of more doubtful value as the dynamical complexity increases. This points towards the fact that formal extensions, or regularizations, of such a constitutive law, though they may appear useful, are in general difficult to legitimate in detail on a physical basis. So they must be considered with a critical eye when using them for predictive purposes.
#### II.4.3 Limitations of block models for extended MCI :
Seismic faults may be represented as spatially extended MCI under both normal and tangential loading. Studying their dynamics then means solving for the elastodynamic motion of two “semi-infinite” deformable media in frictional contact. It was the need for a realistic constitutive law for fault friction which motivated the work which resulted in the formulation of the RR model.
Once such a law is available, one is naturally led to plugging it into a continuum mechanics description. On the other hand, note that, since seismic events amount to the nucleation and propagation of interfacial cracks, they involve variations of deformation fields down to the very small scales relevant close to crack tips.
This leads us to an important remark resulting from the above analysis of the physics involved in the RR law. We saw that this law emerges from two averaging processes :
(i) the rheology described by the interfacial strength $`\sigma _s`$ results from averaging over dissipative events involving nanometric regions within junctions; its expression is thus valid only on scales much above the nanometer one.
(ii) the creep growth of the real area of contact is translated into an geometric age which is an average over a large number of microcontacts. It thus only makes sense on a scale much larger than intercontact distances, i.e. typically millimetric.
As already mentioned, this makes the much discussed question of the ill-posedness Ranjith-Rice of the continuous limit a formal one. In view of the above remarks, which lead to the existence of a physically based small space scale cut-off, the most physical recipee for curing mathematical ill-posedness seems to be the one proposed by Simoes and Martins Simoes : they wash out the UV singularity by replacing the local friction relation between interfacial normal and tangential stresses $`\left(\sigma _t(x)=[\mu _d\sigma _n]_x\right)`$ by a non local one $`\left(\sigma _t(x)=𝑑x^{}w(xx^{})[\mu _d\sigma _n]_x^{}\right)`$, with $`w`$ a spreading function of finite width.
In practice, the dynamic complexity due to the non-linearities of MCI friction Carlson-Langer can only be studied numerically. This is implemented by discretizing the system into elastic blocks of size $`L\times L\times L`$ <sup>19</sup><sup>19</sup>19Choosing blocks much longer in the transverse than in the longitudinal direction (i.e., in Burridge-Knopoff-like models, compressive stiffnesses much larger than shear ones) may lead to dynamical artefacts : for the so-called small events which involve only a few blocks, the associated elastic fields only affect depths, on the order of their lateral extent, much smaller than the block height. It is then illegitimate to neglect internal block degrees of freedom.. RiceRiceBK has pointed out that the existence of the stick-slip instability imposes a natural upper limit on this block size. Indeed, such a discretization implies that the degrees of freedom of a block reduce to those of, say, its center of mass, internal ones being irrelevant. A block of size $`L`$ has a tangential stiffness $`K_LEL`$, with $`E`$ an elastic modulus of the bulk. It bears the normal load $`W_L=\sigma _nL^2`$, with $`\sigma _n`$ the average normal far field stress. It is stable against SS if $`K_LD_0/W_L>(BA)`$, i.e. $`L<L_R`$ with :
$$L_R=\frac{E}{\sigma _n}\frac{1}{BA}D_0$$
(55)
Assume that a block has a size $`L_R<L<2L_R`$. Cut it into $`8`$ subblocks of size $`L/2`$. Each interfacial subblock is stable vis-à-vis SS, while the block is not. So the whole dynamic complexity on scale $`L`$ is associated with the motion of the relative subblock position, which cannot therefore be considered irrelevant.
The maximum size of a discretization block is therefore $`L_R`$. From equation (55), with $`\sigma _n/E10^3`$, $`BA10^2`$, $`D_01`$$`10\mu `$m, $`L_R`$ lies in the meter range.
## III JUNCTION RHEOLOGY : STRUCTURAL AGING/REJUVENATION EFFECTS
It is clear, at this stage, that the rough-on-rough MCI configuration only provides quite an indirect access to the analysis of junction rheology. Indeed, information about the interfacial strength $`\sigma _s`$ can be extracted from friction data only after ”deconvoluting” it from geometric aging effects, necessarily at the expense of precision. This might lead to overlooking finer physical features of dissipation within sheared junctions.
That such might indeed be the case emerged from the experimental study by Bureau et al jamming of the response of a rough/rough PMMA MCI to an oscillating shear force :
$$F(t)=F_0+f\mathrm{cos}\omega t$$
(56)
biased about a value $`F_0`$ below the static threshold. For very small $`f`$, the MCI responds elastically. As $`f`$ is increased, the system enters a regime in which the oscillatory response is superimposed on a slow, self-decelerating, gross sliding motion which corresponds to a saturating, finite displacement. This regime prevails in a narrow amplitude range where $`F_{max}=F_0+fF_s`$, where $`F_s`$ is a static threshold. It was analyzed, on the basis of the RR model, as resulting essentially from the geometric age dynamics. During each force oscillation, the system alternates between sliding, hence rejuvenating, when $`F(t)`$ is close below its maximum, and slowing down strongly, thus aging during the rest of the period. For $`f`$ levels such that aging barely wins, gross motion decelerates and the system finally jams. Above a threshold $`f_>`$, rejuvenation wins, resulting in indefinite accelerated sliding.
The RR model accounts satisfactorily for the long time dynamics in the jamming regime. However, it exhibits, for small slid distances on the order of a few $`100`$ nm, a clear discrepancy with experimental results which strongly suggests that another rejuvenation mechanism, distinct from the geometric one, hence ignored by the RR description, might be at work within the junctions.
### III.1 Accessing junction rheology directly : suitable configurations
This confirms the interest of studying junction rheology directly, with the help of interfacial configurations which either are free of geometric aging or permit to circumvent this effect. This can be achieved in three different ways.
#### III.1.1 Rough-on-flat multicontact interfaces
So far, we have restricted the definition of MCI to interfaces of macroscopic extent between two rough surfaces. It is then tempting, in order to get direct information about junction rheology, to try and realize interfaces la Greenwood which would not undergo the microcontact birth and death process responsible, for rough-on-rough systems, for the non trivial dependence of geometric age on the sliding dynamical history.
Such a configuration has been realized by Bureau et al Lionel EPJB who studied friction between a rough PMMA slider (roughness $`\mu `$m) and a flat and smooth plate made of float glass (roughness $``$ nm) <sup>20</sup><sup>20</sup>20 Since glass is much harder than PMMA, no ploughing of the track by the slider asperities takes place..
The interface is of the GW type but, due to translational invariance of the track, the integrity of the contact population is now preserved upon sliding. The load-bearing asperities creep as in a regular MCI but, since aging is not interrupted by sliding, the geometric age is simply the time $`T`$ elapsed from the creation of the interface. The real area of contact $`\mathrm{\Sigma }_r`$ grows logarithmically with $`T`$ so that, in steady motion at velocity $`V`$, the friction force $`F=\mathrm{\Sigma }_r(\varphi =T).\sigma _s`$(V), becomes an increasing function of $`T`$, as can be seen on Figure 21. One can then take advantage of the slowing down of the logarithmic growth mode by letting the interface ”mature” by waiting up to large $`T`$’s (typically $`10`$ hours), then performing experiments during a comparatively short period (e.g. $`1`$hour) during which $`\mathrm{\Sigma }_r`$ remains quasi-constant, so that friction force variations are directly attributable to those of the interfacial strength $`\sigma _s`$.
#### III.1.2 Junctions in the Surface Force Apparatus
The surface force apparatus (SFA) was first developed Tabor-Winterton IsraelSFA to study adhesion. In a SFA, two atomically planar mica surfaces, slightly curved so as to realize a cross-cylinder geometry (Figure 22), are brought into contact in the presence of a fluid (the lubricant). The interplate distance is decreased down to nanometric values on the order of a few molecular sizes. In this boundary lubrication regime where the fluid is very highly confined, the cylinders deform elastically, leading to a Hertz planar circular contact, with radius commonly in the $`10\mu `$m range.
The SFA was later modified JacobSSR so as to allow for shear motion and measurement of friction forces. Since the area is measured optically, the normal stress borne by the contact is directly accessible, as well as the shear stress associated with frictional dissipation in the junction formed by the lubricant. In many instances (see below) such boundary lubrication layers behave as weak disordered solids, exhibiting an elastic regime at low shear levels.
#### III.1.3 Extended soft contacts
We refer here to the contacts which form between a highly compliant (soft) slider with a smooth surface and a hard flat smooth track. High slider compliance ensures that molecular adhesive contact is realized everywhere along the interface, even in the presence of submillimetric roughness.
Such contacts can be achieved in either of two configurations :
flat-on-flat: up to now, this geometry has been used with sliders made of moulded hydrogels (Young moduli in the $`1`$$`10`$ kPa range) Gong gelEPJ . The contact lateral extent is usually of order centimeters.
ball-on-flat: this is the configuration commonly used to study friction at elastomer/glass interfaces Schall Barquins Chaudhury . Pressing the ball onto the flat under controlled normal load results in the formation of a circular adhesive Hertz contact Johnson . Since elastomer Young moduli typically lie in the $`1`$$`10`$ MPa range, contact radii in such experiments are commonly of order millimeters.
### III.2 MCI junctions revisited :
The results described here were obtained by Bureau et al Lionel EPJB on MCI formed by rough PMMA sliding on smooth (float) glass. The surface state of the glass was controlled by grafting onto it a single monolayer of short silane molecules, which passivate strongly adhesive sites.
As already mentioned, such interfaces, when “old enough”, can be considered to work at constant area of contact $`\mathrm{\Sigma }_r`$, thus giving direct access to variations of the frictional stress <sup>21</sup><sup>21</sup>21 As the absolute value of $`\mathrm{\Sigma }_r`$ is not measurable accurately, the absolute value of $`\sigma _s`$ cannot be accessed in this configuration..
#### III.2.1 Structural aging :
The central result is illustrated on Figure 23. Namely, standard stop-and-go experiments (see $`\mathrm{\S }`$II.B.1) reveal the presence of a static friction peak, the amplitude of which grows logarithmically with the waiting time $`t_w`$ (Figure 24).
Contrary to what is the case for rough/rough MCI, this peak cannot be associated with geometric aging, which is not operative here. On the other hand, the very existence of a $`\mu _s(t_w)`$ larger than the dynamic $`\mu _d`$ means that the interface is the seat of an aging-when-waiting vs rejuvenating-when-sliding process, which must therefore necessarily take place within the junctions and affect their structure. The existence of this structural aging entails that $`\sigma _s`$ is not, as we assumed up to now, a mere function of the instantaneous sliding velocity $`\dot{x}`$, but does itself depend upon the dynamical history of the interface.
Moreover, one observes that the structural aging rate depends noticeably on the value of the tangential stress applied during the stop phase (see Figure 24) : the lower its level, the slower aging is (see figure 5 of reference Lionel EPJB ).
Finally, in many cases – and especially for small waiting stresses – structural aging becomes observable only after a finite, strongly sample-dependent, latency time $`\tau _L`$.
#### III.2.2 Steady sliding dynamic friction :
Figure 25 shows a typical example of the velocity dependence observed for $`\mu _d`$. It exhibits a minimum at $`V=V_{min}`$. The velocity-weakening observed for $`V<V_{min}`$ provides another evidence that the structural age decreases with $`V`$, i.e. that sliding induces rejuvenation.
For $`V>V_{min}`$, $`\mu _d`$ increases quasi-logarithmically ($`\mu _d/\mu _{min}(1+\alpha \mathrm{ln}(V/V_{min}))`$. That is, we recover here the expression of $`\sigma _s`$ (equation (29)) deduced from the analysis of experiments on rough/rough MCI, which were used to formulate the RR model. Moreover, the reduced slope $`\alpha 4.10^2`$ has a magnitude comparable with the value ($`\alpha 5.10^2`$) found for rough/rough PMMA junctions. This enables us to conclude that already quite close above $`V_{min}`$, the rejuvenation effect is quasi-saturated.
#### III.2.3 Discussion
This phenomenology, the fine details of which were masked, for rough/rough sytems, by the larger effects of geometric aging <sup>22</sup><sup>22</sup>22As can be seen on Figure 24, the log-slope of $`\mu _s`$ associated with structural aging is on the order of a few $`10^3`$, typically from $`10`$ to $`4`$ times smaller than the aging slopes for rough/rough MCI. This leads one to conclude that, in the latter case, the growth of $`\mu _s`$, though dominated by the geometric effect, does contain a small contribution due to structural aging. , appears qualitatively consistent with our description of junctions as confined weak glassy media. Indeed, it is closely akin to the behavior of sheared bulk soft glassy materials, such as colloidal glasses Derec Abou and pastes Cloitre . The rheology of these systems has also been modelized Sollich Fielding Lequeux as resulting from the interplay between :
— aging at rest, i.e. strengthening via slow relaxation down the highly multistable energy landscape characteristic of quenched disordered, jammed, media.
— rejuvenation by motion, which can be seen as reshuffling the populations of the local minima.
In the sliding (i.e. plastic) regime, the rate of the localized dissipative events invoked in $`\mathrm{\S }`$II.C.4 should be increased by the decrease of the energy barriers due to stress induced biasing. The “young”, shallow energy states which were depopulated as aging pushed the system into “older”, deeper ones, thus can get repopulated : the sheared soft glass rejuvenates.
In the same perspective, one may tentatively interpret the above-mentioned dependence of the aging rate on waiting stress level as follows. In a stop-and-go experiment, the “stop” acts as a mechanical quench of the strongly rejuvenated formerly sliding state. Aging then restarts in a landscape with barriers whose height is decreased by stress-induced biasing, hence accelerates when the stress level increases <sup>23</sup><sup>23</sup>23 Viasnoff and Lequeux Viasnoff have found that applying to a colloidal glass, aging under stress-free conditions, an oscillating shear stress of finite duration does result in a strong perturbation of the aging process. However, it is important to note that, in their experiment, the stress amplitude is larger than the yield stress of the material, which certainly results in some non-stationary flow. This contrasts with the situation of ref. Lionel EPJB , where the applied waiting stress does not provoke any sliding..
### III.3 Boundary Lubrication Junctions
Over the past twenty years a considerable body of results has come out of friction experiments performed in the SFA, in the boundary lubrication (BL) regime. It is much too vast to be extensively reviewed here. We limit ourselves to those salient features which we consider to be most relevant in the perspective of this article.
#### III.3.1 Confinement-induced solidification
When the “lubricant” – a liquid in the bulk phase at the temperature of the experiment – is compressed between the two mica surfaces, it is gradually squeezed out until its thickness is reduced to a few molecular sizes. Then, for materials composed of small, spherical or short-chain, molecules (e.g. OMCTS – an inert silicone – squalane, linear alcanes) one observes JacobSSR that the normal force-vs-thickness $`F_N(D)`$ curve develops oscillations whose amplitude grows rapidly as $`D`$ decreases (see Figure 26). This behavior can be ascribed to the layering of the confined fluid, each successive minimum being reached by squeezing out one more full layer.
Thin enough (typically 3–4 molecular layers) such junctions, when sheared, exhibit an elastic response at low shear stress levels, until they reach a static threshold where they yield and start sliding Homola Yoshi van Alsten . This demonstrates that they have solidified under the effect of the high confinement <sup>24</sup><sup>24</sup>24 For some fluids made of short-chain molecules, e.g. hexadecane Yoshi , the solid-like response only buils up after a finite amount of shear-induced sliding – suggesting that in this case initial sliding helps ordering.. This behavior has been confirmed by numerical studies Thompson .
Klein and Kumacheva Klein found that, in the case of OMCTS (and at variance with previous results reported in van Alsten ), the liquid-to-solid transition occurs abruptly when decreasing the number of layers from $`7`$ to $`6`$. Whether or not such abruptness is a general feature has not been documented on different systems.
Finally, it was shown in Yoshi that increasing the temperature makes the confined solid weaker, i.e. leads to a decrease of its yield stress.
It is important to note that the highest normal stresses which can be realized in the SFA are on the order of $`10`$ MPa. That is, they are at least one order of magnitude smaller than the levels met in MCI microcontacts, which we have seen to be comparable with the yield stress of the confining bulk solids.
#### III.3.2 Structural aging
Stop and go experiments show that the static threshold of solidified BL junctions increases with the waiting time $`t_w`$ spent at rest. This strengthening, which reveals structural aging, has been studied quantitatively on a variety of lubricants, including hexadecane Yoshi , squalane Dru1 Gour and a star-shaped polymer melt Yamada . Strengthening also manifests itself through the increase with $`t_w`$ of the layer elastic stiffness Reiter .
In all cases, strengthening is slow, namely either linear on a logarithmic $`t_w`$-scale Yoshi Reiter or somewhat faster Gour . As illustrated on figure 27, finite latency times, ranging up to about $`10`$ seconds, previous to aging, have been observed on hexadecane Yoshi and several lubricants including squalane Dru1 . This result was later contradicted, for squalane, in reference Gour – a discrepancy which is possibly ascribable to the difficulty of defining such a notion with high precision. Indeed, at short waiting times, the height of the stiction peak (the difference between the static threshold force and the dynamic one at the imposed driving velocity) becomes very small. Since accuracy is limited by the noise level, it may be difficult to ascertain whether one should invoke complete latency or a gradual increase of the aging rate.
Finally, Drummond and Israelachvili Dru1 have compared, for squalane and a star-shaped lubricant (PAO), aging under zero and a finite shear stress close below the dynamic level. Contrary to what was observed on PMMA/glass MCI, they find that, for a given $`t_w`$, the stiction peak is systematically higher when waiting under zero shear stress. This leads the authors to associate rejuvenation upon sliding, responsible for the force drop producing the stiction peak, with shear-induced molecular alignment, whose relaxation is certainly faster at zero waiting stress. We will come back to this in the discussion of $`\mathrm{\S }`$III.F.
#### III.3.3 Sliding dynamics
All BL systems with a finite static threshold exhibit the same dynamical behavior in the explored low driving velocity domain (ranging from $`10^2\mu `$m/sec to a few $`10\mu `$m/sec). Namely, as illustrated by figure 28 :
— For $`V<V_{c1}`$ periodic stick-slip (SS) motion.
— In a finite range $`V_{c1}<V<V_{c2}`$, the motion exhibits intermittency, in the form of randomly spaced stick-slip peaks Dru2 .
— For $`V>V_{c2}`$, these peaks disappear, sliding becomes steady Dru2 and, in this regime, dynamic friction is $`V`$-weakening.
In most cases studied, the amplitude of the SS force signal is roughly constant up to $`V_{c2}`$ <sup>25</sup><sup>25</sup>25 Whether or not such a discontinuous transition exhibits hysteresis upon increasing/decreasing $`V`$ is not documented, to our knowledge..
It is worth mentioning that a qualitatively similar behavior is observed with gelatin/glass extended sliding contacts gelatPRL . In this latter case it is associated with the fact that, in the SS regime, sliding is spatially inhomogeneous – it occurs via propagating slip pulses (see $`\mathrm{\S }`$III.C.4 below). In the intermittent regime above $`V_{c1}`$, the corresponding large force peaks are interspeded with smaller ones due to nucleation, propagation and death of small slip events within the contact. This remark suggests the possibility that sliding in BL junctions might also exhibit spatial complexity. Such a model has been proposed by Persson Bo , who invokes the nucleation and propagation of ”shear-melted islands”.
In an extensive study of squalane junctions, Gourdon and Israelachvili Gour have, in particular, investigated the effect of normal load on the sliding dynamics. They find that the dynamics described above prevails at high normal stresses ($`P5`$ MPa). However, at lower normal stresses ($`2`$ MPa), a different behavior is observed : for $`VV_{c1}`$, the amplitude of the erratic stick-slip decreases continuously down to zero.
Finally, using a SFA with a large sliding range ($`500\mu `$m), Drummond and Israelachvili Dru1 have studied, for squalane, the shape of friction force transients following velocity jumps. These transients are long-lived, and characterized by a distance comparable with the contact diameter.
### III.4 Extended Soft Contacts : Gelatin/Glass Friction
They are particularly easy to realize with hydrogels, such as gelatin, which exhibit both ultra-low shear moduli (typically a few kPa) and an extended elastic regime. Their high solvent content (typically $`80\%`$) endows them with specific frictional properties resulting from their poroelastic character DLJohnson . For this reason, we separate this case from that of non-swollen elastomers.
Gong and Osada Gong have performed extensive studies of :
— the influence on the dynamic friction level of the (attractive vs repulsive) nature of gel/substrate interactions. They have shown that repulsive interactions result in the formation of a thin interfacial layer of solvent, hence to hydrodynamic lubrication. For attractive couples, friction is strongly influenced by the formation of transient adhesive polymer-substrate bonds, which they modelized by adapting a description due to Schallamach (see below).
— the dependence of the friction force on normal loads in these different situations.
However, up to now, the only system on which the nature of the sliding dynamics has been fully analyzed is gelatin on glass gelEPJ , on which we now concentrate.
#### III.4.1 Static threshold
The experiments described here were performed under zero normal load, adhesion being strong enough to lead to solid friction, i.e. to a finite static threshold.
When a gelatin block is sheared (Figure 29), as the force increases, the gel/glass contact first remains pinned, and the block deforms elastically. At a threshold stress $`\sigma _0`$ sliding sets in. This occurs via nucleation at the trailing block edge of an interfacial fracture without measurable vertical opening, which propagates forward with a constant velocity $`V_{tip}`$ in the mm–cm/sec range, i.e. much smaller than that of transverse sound (typically $`1`$ m/sec) <sup>26</sup><sup>26</sup>26Note that, while the shear modulus is controlled by the elasticity of the loose polymer network, the undrained bulk modulus is that of the solvent, so hydrogels can be considered incompressible.. By changing the gelatin concentration, hence the mesh size $`\xi `$, and the solvent viscosity $`\eta _s`$, it appears that $`V_{tip}D/\xi `$, with $`D`$ the collective diffusion coefficient DLJohnson , measured from quasi-elastic light scattering experiments Tanaka : $`DG\xi ^2/\eta _sk_BT/\eta _s\xi `$, with $`Gk_BT/\xi ^3`$ the shear modulus.
This indicates that the interfacial pinning, which must be severed in order to crack, is provided by bonds with a typical lateral spacing $`d`$ on the order of the mesh size.
Indeed, the value of the self-selected $`V_{tip}`$ can be interpreted as follows. The energy released by bond snapping events is emitted at the frequency $`\omega V_{tip}/d`$, and wave vector $`qd^1`$. The lowest deformation mode of the gel block is the so-called Biot mode, associated with network diffusion within the solvent $`\omega _D=Dq^2`$. It resonates with emission, hence becomes most efficient to dissipate the fracture energy, when $`\omega _D\omega `$, from which it emerges that $`d\xi `$ <sup>27</sup><sup>27</sup>27It is worth mentioning that Rubinstein et al Fineberg heve recently observed that sliding between purely elastic solids also sets in via propagation of interfacial cracks, but in the elastic case these propagate at velocities in the sonic range. The low values of $`V_{tip}`$ for gelatin are therefore, in contrast, a signature of poroelasticity..
Now, one also observes, as shown on figure 30, that the threshold $`\sigma _0`$ for crack nucleation is not a constant, but increases logarithmically with waiting time at rest $`t_w`$ (in the explored range $`t_w1sec`$). This indicates that adhesive bonding between the substrate and the polymer tails and/or loops dangling from the disordered network proceeds on two time scales : a rather fast one, which leads to finite static pinning with a bond areal density $`\xi ^2`$, followed by a slow, thermally activated relaxation of the adsorbed polymer configurations leading to the slow increase of this density.
#### III.4.2 Sliding dynamics
When the gel block is driven at velocity $`V`$, two types of dynamics are observed.
(i) Above a critical velocity $`V_c`$ (ranging, depending on gel composition, from $`25`$ to $`250\mu `$m/sec), sliding is steady. The sliding stress $`\sigma _s(V)`$ is velocity-strengthening, and corresponds to a shear-thinning rheology, namely :
$$\sigma _s(V)V^{1\alpha }$$
with $`\alpha =0.6\pm 0.07`$. This is reminiscent of the behavior of polymer solutions Bird , and suggests that, in this sliding regime, interfacial pinning plays only a minor role, dissipation being due to the viscosity of a layer of thickness the mesh size (on the order of $`10`$ nm) formed by a solution of polymer segments attached to the network. If such is the case, the shear strain rate $`\dot{\gamma }V/\xi `$, and the effective viscosity
$$\eta _{eff}=\frac{\sigma _s}{\dot{\gamma }}\eta _s(\dot{\gamma }\tau _r)^\alpha $$
with $`\tau _R`$ a typical relaxation time for a chain length $`\xi `$ in solution. Using for $`\tau _R`$ the Rouse time one gets, for a variety of gelatin concentrations and solvent viscosities, the collapse of data shown on Figure 31, which lends excellent support to the above interpretation.
(ii) For $`V<V_c`$, periodic stick-slip prevails. Optical observation shows that during the “slip phase” sliding is not homogeneous, but proceeds by propagation of self-healing pulses (Figure 29). Pulse heads are the previously described crack tips, behind which the local slip velocity $`v`$ decreases from the standard quasi-diverging field towards the driving value $`V`$. However, this regular decrease is suddenly interrupted when $`v`$ reaches a value equal to the critical value $`V_c`$ which turns out to coincide with the upper limit of the stick-slip regime. At $`v=V_c`$, resticking (healing) occurs quasi-discontinuously.
When $`V`$ approaches $`V_c`$ from below, the pulse length increases, so that the transition to steady sliding takes place via increasing time spacing of pulses of quasi-constant amplitude, while being non hysteretic. Moreover, for $`VV_c`$, an intermittent behavior, already described above ($`\mathrm{\S }`$III.C.3) is observed – the interpulse spacing becomes erratic, and they are interspeded with smaller force signals associated to short-lived small propagating slip events. A more detailed characterization of this complex dynamical behavior will certainly be of interest.
#### III.4.3 Rate and state interpretation
Clearly, as for the previously considered cases, static strengthening means interfacial aging at rest, which we expect to be associated with rejuvenation, i.e. weakening, upon sliding, as proved by the reproducibility of $`\sigma _0(t_w)`$ in stop and go experiments.
This is consistent with the abruptness of the resticking process at the trailing edge of the self-healing pulses. Indeed, assume that the steady-sliding characteristics is $`V`$-weakening for $`V<V_{min}`$. Close behind the pulse head, the local $`v`$, larger than $`V_{min}`$, lies on the locally stable branch of $`\sigma _s(V)`$. As the distance behind the head increases, $`v`$ decreases. When it reaches $`V_{min}`$, the interface becomes unstable on a small spatial scale (smaller than the optical resolution). In other words, since the driving stiffness, provided by the very compliant gel block, is extremely low, the interface jumps down to $`V=0`$ at quasi-constant stress. This leads to concluding that the resticking velocity $`V_c`$ corresponds to a minimum of $`\sigma _s(V)`$ above which, as appeared in the previous paragraph, interfacial pinning should be negligible. This interpretation is also consistent with the coincidence of the resticking velocity and the disappearance of stick-slip as soon as the asymptotic value $`V`$ of the local velocity $`v`$, needed for steady sliding to get established, reaches the stable branch of $`\sigma _s(V)`$.
So, here again, it seems natural to try and define a structural age of the gel/glass interface. Quite clearly, the corresponding dynamic state variable must somehow measure the pinning energy associated with the formation of polymer-glass adhesive bonds.
Already $`50`$ years ago, Schallamach Schall proposed, in the context of dry rubber friction, a seminal model of such a dynamic state variable, which remains the basis of all more recent theories of gel Gong and rubber Charitat Chernyak friction. In this model, briefly summarized an discussed in Appendix D, the elementary mechanical instability (see $`\mathrm{\S }`$II.C.4) is the snapping of a bonded molecule out of its adsorption site under shear loading through an elastic spring (here the polymer segment). Once this energy barrier has been jumped over, under the combined effect of loading and thermal activation, the molecule is freely advected until it readsorbs via a thermally activated process. Of course, readsorption is thwarted by advection, which limits the time available for this process to take place. The state variable is thus the density of bonded molecules. The friction stress is the product of this $`V`$-weakening density, and of the $`V`$-strengthening average depinning force (the faster advection, the less time for escaping above high barriers). So it exhibits a bell shape with a maximum at $`V=V_{max}`$. The low velocity regime is an Eyring viscous one, with the bond density close to its equilibrium saturation level. In the $`V`$-weakening high velocity regime the depinning force saturates, $`\sigma _s(V)`$ is then fully controlled by dynamic rejuvenation, hence vanishes for $`V\mathrm{}`$.
Models of the Schallamach type call for two important remarks :
(i) The extent of the Eyring-like low $`V`$ regime ($`V<\stackrel{~}{V}`$, see Figure 35) crucially depends on the height of the desorption barrier. When, as often assumed – explicitly or not – it is on the order of a few $`k_BT`$, $`\stackrel{~}{V}`$ is non-negligible. A strong consequence is that, in such cases, no static threshold should be observable except if one could load at exceedingly large velocities. This points towards the interest of systematically associating studies of dynamic friction in such systems with stop and go experiments.
Gelatin/glass interfaces exhibit well defined static force peaks and crack tips. We must therefore conclude that, in this case, we are dealing with strong bonds – in which case $`\stackrel{~}{V}`$ becomes negligibly small (see end of $`\mathrm{\S }`$II.C.4).
(ii) If we stick to the Schallamach prediction, in the case of gelatin, $`\sigma _s(V)`$ should be uniformly $`V`$-weakening, since in the explored range $`VV_{max}`$. How can one then explain the existence of the $`V>V_c`$ strengthening behavior?
We meet here with an important shortcoming of the Schallamach model. Namely, unpinned molecules are assumed to glide freely, that is, viscous dissipation in the interfacial layer is omitted - although, since its thickness is typically nanometric, shear rate levels are high. In the case of gelatin we saw that, on the contrary, it is this contribution which fully accounts for the observed non-newtonian friction at the essentially depinned interface. This illustrates the importance of taking into account in this kind of problem, not only interfacial pinning, but, as well, junction viscosity. When both effects are of comparable importance, this open problem becomes a true nanofluidics one, as far as the size of the confined polymer chains is precisely comparable with the thickness of the “channel” , i.e. the junction.
In summary, once complemented with standard viscous dissipation in the sliding junction, the Schallamach picture, which identifies the state variable with the density of adsorbed bonds, provides a sound basis on which to modelize gel friction. Clearly, a more elaborate modelization of the dangling chains, taking into account the existence of a multiplicity of possible bonding configurations Semenov and a detailed description of chain dynamics Charitat would be needed to make such models more quantitative, as well as to touch upon the question of long term aging.
### III.5 Extended Soft Contacts : Dry Elastomer Friction
Due to its importance in a number of applied fields (e.g., to name only a few, traction of tires, windshield wiper efficiency and durability), this question has been a subject of very active investigation in the past six decades. We restrict ourselves here to a sketchy summary of the commonly accepted concepts used for interpretation of interface geometry vs molecular adhesion effects. For extensive bibliographies, see Roberts Savkoor Chaudhury .
We consider here the case of elastomers in extended contact with much stiffer solids, which can thus be assumed non deformable. Bulk viscoelastic dissipation is very important in rubbers. Its strong temperature dependence follows the so-called WLF time-temperature superposition principle Ferry . Grosch Grosch , and Ludema and Tabor Ludema , have shown that one must distinguish between contributions to the sliding friction stress governed respectively by bulk and interfacial dissipation.
#### III.5.1 Bulk dissipation
It comes into play as soon as the stiff partner surface exhibits non negligible roughness - on a scale larger than the average cross-link spacing. Assume for a moment that the roughness has a characteristic wavelength $`\lambda `$. Then, in stationary sliding, a given material point within the elastomer feels a stress field modulated at a frequency $`\omega V/\lambda `$, which penetrates down to a depth $`\lambda `$ below the surface. This results in a dissipation governed by the loss modulus of the bulk, $`G^{\prime \prime }(\omega )`$, i.e. obeying a $`VT`$ superposition principle. One is thus able to account for the most salient feature of rubber friction, namely the existence of a maximum of the $`\sigma _s(V)`$ curve, at a velocity $`V_m`$, which increases with the temperature $`T`$, above which stick-slip motion sets in Grosch .
#### III.5.2 Interfacial dissipation
Grosch Grosch identified, in his experimental results on friction of several elastomers on silicon carbide paper, an additional feature – a shoulder on the $`\sigma _s(V)`$ curve at a velocity $`V_m\omega _m\mathrm{\Lambda }`$, where $`\omega _m`$ corresponds to the maximum of the loss modulus and $`\mathrm{\Lambda }6`$nm. This shoulder turned into a broad maximum when sliding on smooth “wavy” glass. He attributed this to interfacial dissipation – i.e. to the snapping out of advected adhesive bonds as modelized by Schallamach (see $`\mathrm{\S }`$III.D.3).
It is only recently that progress in surface control enabled Vorvolakos and Chaudhury (VC) Chaudhury to reconsider this question. They used two model systems, constituted of cross-linked PDMS sliding on (i) a self-assembled silane monolayer supported by a Si wafer (ii) a thin film of glassy PS. They were thus able to get rid of surface roughness effects, hence of bulk dissipation.
Silane coverage results in very low energy surfaces, so one expects polymer substrate adhesive bonds to be rather weak in this case and, hence, Schallamach’s model to be appropriate. Indeed, VC show that their results for $`\sigma _s(V)`$ in this system (see figure 32), which extend over $`5`$ velocity decades, are well explained qualitatively in this frame <sup>28</sup><sup>28</sup>28 Interfacial viscous dissipation is neglected in the VC analysis. As already mentioned, it should come into play, especially in the $`V`$-weakening regime $`V>V_m`$, where it is certainly needed to interpret the observed stick-slip behavior.. Moerover, they deduce from a time-temperature superposition argument an energy scale $`25`$ kJ.mol<sup>-1</sup>, i.e. $`0.3`$ eV/bond. Whether this energy, much larger than the van der Waals ones expected for PDMS/silane interactions, is relevant to adhesive bond snapping in this system remains an open question.
PDMS bonding to PS is likely to be stronger than that on silane. Indeed, consideration of the VC results in the low velocity limit leads to raising the question of the possible existence of a finite static friction threshold. Systematic investigation of the transient behavior at the onset of sliding would certainly be enlightening.
Finally, it is worth mentioning that non homogeneous sliding has been observed Schallwaves Roberts at some smooth soft rubber/glass interfaces. This occurs via propagating self healing pulses. But, contrary to the gelatin case where contact is maintained in the sliding regions, these Schallamach waves are of the mode-I fracture type. That is, they consist of regions of interfacial detachment with macroscopic opening which re-adhere at the back edge. So, in this regime, dissipation is essentially controlled by adhesion hysteresis Barquins . Though high compliance and large viscoelastic losses are thought to promote this sliding mode, no prediction relating material properties to the dynamics, nor even to the occurrence of Schallamach waves is available up to now.
### III.6 A tentative classification : Jammed junction plasticity vs adsorption controlled dynamics
The various frictional behaviors briefly described above pertain to junctions which differ in several respects, namely :
$``$ Normal stress level, i.e. confinement pressure $`P`$, which ranges from zero (gelatin) to about the yield stress of the softer bulk material (rough on flat MCI).
$``$ Level of adhesive interactions, i.e. strength of the pinning sites provided by the confining potential. This ranges from very weak (fully silanized substrates) to strong (e.g. bare glass).
$``$ Density of adsorbable units, ranging from dense (polymer glasses) to (semi)dilute (hydrogels).
All these systems share two properties :
– they exhibit a static threshold <sup>29</sup><sup>29</sup>29 Except possibly for some elastomers under low normal stress and in contact with highly passivated substrates.
– this threshold ages slowly with time spent at rest.
The experimental signature of this structural aging is the slow (in general quasi-logarithmic) growth of the upper limit of the elastic regime. In experiments where the slider is remotely driven at constant velocity $`V`$, it results in general in a force peak. This unambiguously reveals rejuvenation upon sliding, since the sliding velocity $`v`$ is equal to the driving one $`V`$ both in the steady sliding state and at the peak, where the slid distance is minute, so that the age is basically given by the waiting time before pulling <sup>30</sup><sup>30</sup>30 When established sliding is not steady, but occurs via stick-slip, static aging results in the growth of the first peak with stopping time.. This means that the state of the sliding junction is specified, not only by its instantaneous velocity, i.e. by the shear rate $`\dot{\gamma }`$ it experiences, but also by the value of some dynamical state variable.
On the other hand, different classes of systems exhibit different steady state rheologies, from logarithmic strengthening for highly confined PMMA/glass to power law for gelatin.
While these behaviors call for rate and state descriptions, it is clear that various classes of systems must be distinguished, in terms of the microscopic nature of the state variable $`\phi `$, which is necessarily related with the junction structure and, hence, with the pinning mechanism giving rise to static friction. The existing body of experimental results discussed in $`\mathrm{\S }`$III.B–E is still far from sufficient to provide a firm basis for such a phenomenology. So, the classification which we propose here is only tentative, and primarily intended to help structuring further questions to be, in a first stage, investigated experimentally.
Let us start from the very high confinement regime, realized in the MCI configuration ($`\mathrm{\S }`$III.B) where, for PMMA, the confining pressure $`P100`$ MPa is comparable with the bulk yield stress. We saw that junctions between such a polymer glass and highly silanized glass behave as soft glassy nanometer-thick media.
In other words, we claim that, in such a case, the disordered material of the dense junction jams livre Nagel , i.e. solidifies under confinement. It ages as a glass, by relaxing towards deeper local minima of the energy landscape (inherent states) and, when it flows, dissipation can be attributed to irreversible flips of molecular-sized shear transformation zones (STZ) Falk . In this case, frictional sliding is nothing but plastic flow of this interfacial glass which, being weaker than the bulk, naturally localizes shear.
Now, note that, in the MCI experiments, when the glass substrate is well silanized, no direct manifestation of the polymeric nature of the junctions is observed. However, when the same experiments are performed on a poorly silanized substrate Lionel EPJB , the previously observed rather narrow aging peak persists, but is followed by a long bumpy stress transient (see Figure 33). Moreover if, after reaching the steady state, one performs stop-and-go tests with a high stopping shear stress level, the long transient is absent when restarting motion. If, on the contrary, the stop is performed under zero shear stress, the wide bump reappears. This strongly suggests that these transients, associated with the presence of stronger corrugations of the confining substrate potential, are due to gradual stretching and alignment of the constitutive molecules, which persists when the shear stress is maintained, and relaxes otherwise. Such rubbing-induced alignment has been documented already long ago by Pooley and Tabor Pooley-Tabor on teflon, a material for which it is particularly strong. It would be desirable, following these authors, to test this interpretation by devising a set-up which would permit to change the sliding direction by e.g. $`90^{}`$.
This observation naturally raises the question of the competition between jamming and adhesive pinning to the substrate. It is reasonable to expect that, the lower the confining pressure $`P`$, the weaker the jamming, i.e. the relative contribution of STZ flips to dissipation.
In boundary lubrication experiments performed in the SFA, pressure levels are smaller that those in MCI junctions by at least two orders of magnitude. As mentioned in $`\mathrm{\S }`$III.C.2, the amplitude of the transients observed on squalane by Drummond and Israelachvili Dru2 is larger when stopping under smaller shear stress. They also seem to be associated with quite broad stress bumps. This leads one, following the authors, to associate these transients with molecular alignment. Whether or not a narrow glassy-like aging peak is observable at the highest SFA pressures would be worth investigating in detail.
This remark points toward the need for trying to bridge, for given interfacial materials, between SFA and MCI pressure levels, for example with the help of the ball-on-flat configuration. Such comparisons are necessary in order to appreciate to which extent SFA data can be safely extrapolated to the conditions prevailing in macroscopic contacts between hard solids.
The other limiting situation of purely adhesive pinning. It is exemplified by gelatin gels in contact with glass under zero or very small gelEPJ normal load ($`\mathrm{\S }`$III.D). In this case the junction is constituted of a (semi) dilute solution of protein segments in water – no jamming can be invoked. The pinning responsible for the finite static threshold is due to adsorption of the polymer molecules onto the glass substrate, and one may, at least semi-quantitatively, identify the state variable $`\phi `$ with the areal density of adsorption bonds. Such a structure is able to relax by a combination of thermal desorption, chain configuration rearrangements and readsorption which leads, at rest, to the slow increase of $`\phi `$. The number of candidate sites on a segment of course depends on the chemical nature of the network-forming polymer chains. This dynamical problem, although non trivial, might be amenable to a theoretical treatment which would extend equilibrium studies such as those of Subbotin et al. Semenov alo ng the lines of the Schallamach-like steady state model built by Charitat and Joanny Charitat for the case where only the end monomer is adsorbable.
In such a completely non-jammed junction, frictional dissipation contains two contributions due respectively to :
— Thermally helped advection-induced depinning. It is this mechanism which controls the solid friction behavior proper for such adsorption-controlled interfaces.
— Viscous flow of the polymer-water mixture in the sheared interfacial layer.
While the latter is velocity-strengthening, the former decreases with decreasing $`\phi `$, so it is basically velocity-weakening.
Elastomers probably correspond to an intermediate situation. The interfacial junctions which they form certainly have a dense polymer content, while clearly not being solid (jammed), since their bulk itself is not. Junction thickness scales as the distance between entanglements. One of the open questions is that of evaluating realistically the viscous contribution to friction of such sheared layers. Whether they exhibit, as might be expected, slow aging at rest is not documented yet, as far as all experiments on elastomer friction have focussed up to now on the steady sliding behavior and on its temperature dependence.
## IV Conclusion
The route we have chosen to follow in this review has led us from macroscopic dynamics down to gradually decreasing space scales. Solid friction thus appears as a particularly favorable area of material science in which the governing mechanisms are “universal” enough to allow for a rather general phenomenology down to the nanoscopic level.
A first level of analysis, relevant to the widespread case of macroscopic interfaces between rough hard solids, identifies two main physical ingredients :
(i) A shear-induced depinning process, in which mechanical instabilities lead to structural rearrangements of nanometric shear transformation zones within the junctions where shear localizes. It is the associated multistability of these glassy junctions which is responsible for the existence of finite friction thresholds.
(ii) Geometric aging, i.e. slow creep growth of the area of the sparse microcontacts forming such interfaces under the high contact pressures associated with this geometry. Sliding limits the duration of contact life, making geometric age a dynamical variable with a memory of the sliding history.
This description provides the physical basis of the Rice-Ruina constitutive laws which, when properly extended, account for all the main features of the low velocity frictional dynamics. It also permits to point more precisely the limits of such a constitutive law. Among these, an important one is concerned with higher velocities — typically $``$ a few $`100\mu `$m/sec – at which geometric age becomes small, leading to the saturation of the destabilizing $`V`$-weakening behvior of dynamic friction, and to its inversion into $`V`$-strengthening. Behaviors in the intermediate (mm/sec to cm/sec) range are still insufficiently documented experimentally to allow for detailed models of dissipation in this regime.
For still faster sliding, self-heating becomes non negligible, turning temperature into a new relevant state variable which affects junction rheology. An example of such an effect is that of the stick-slip dynamics of the bowed violin string, which has been modelized by Smith and Woodhouse Woodhouse in terms of the triggering by slip-induced heating the of a transition from sticking solid to viscous fluid of the rosin rubbed on the bow.
Beyond the slow dynamics leading to the notion of geometric age, another one emerges, associated with structural aging of the junctions, whose effects are masked, for rough-on-rough systems, by the larger geometric aging ones. Structural aging at rest and its counterpart – rejuvenation when sliding – is at work in all types of frictional junctions, leading to slow growth of static thresholds with time and to stick-slip-like dynamical instabilities. It affects the topography of the multistable energy landscapes which govern frictional rheology. Further experimental information will be needed in order to try and specify it in terms of precise state variables, whose nature is probably not unique.
On the basis of comparisons between different classes of interfaces, we have suggested a schematic classification in terms of two limiting types.
(i) Jammed junctions, solidified under confinement into a soft glass structure. They are relevant at high contact pressures such as met with multicontact interfaces. Multistability results from their glassy structure itself, and in this case frictional dissipation enters the wider class of problems constituted of plasticity of amorphous media and soft glass rheology.
(ii) Purely adhesive junctions for which pinning is due to adsorption bonds between the junction molecules and the substrates(s). They form at contacts under low confinement pressures involving gels and, more generally, polymers.
Many real junctions probably are a compound of these two ideal types, their frictional dissipation containing contributions of both types whose respective weights depend on pressure, temperature, density and on the physico-chemical nature of the confining surfaces.
While, in our opinion, the question of geometric age effects in multicontact solid friction is now reasonably cleared up, that of juction rheology remains a largely open issue worth of further investigation. In order for modelization to make significant progress, more detailed experimental information is needed, in particular about :
– Systematic trends associated with varying pressure and temperature, for systems with controlled surfaces.
– The effect of various aging histories, paralleling analogous studies in structural and spin glasses. In the case of adsorption-controlled junctions, information about aging might also be obtained via adhesion tests adhesion , provided that these be performed on materials with negligible viscoelasticity.
Numerical simulations of sheared junctions in the presence of realistic confining potentials should also be of great help.
Finally, the reader may have been surprised that we never used terms such as “nanofriction” or “nanotribology”. This we did purposefully, since we consider them rather misleading. Indeed, “nanotribology” is often associated with investigations of the nature of the basic dissipative processes which, as we saw, occur on the nanometric scale. However, friction, as defined in terms of a stress or a friction coefficient, results from the self averaging of such events within contacts of, at least, mesoscopic (micrometric) size. This is precisely why friction laws do not make sense on the nanometric scale.
We consider that the term nanofriction should be reserved for situations where lateral contact sizes are themselves of nanometric order, in which case different effects may emerge, such as exemplified by aging due to capillary condensation Riedo . It is also likely that, in such regimes, the nature of dissipative processes Kim is different from that in mesoscopic contacts.
## Acknowledgments
We are particularly grateful to J.R. Rice, K.L. Johnson and Y. Bréchet for encouraging our first steps in the field of friction, and letting us benefit liberally from their vast knowledge of mechanics and material science. We are deeply indebted to P. Nozières, B. Perrin, B. Velicky and F. Heslot for their precious contributions at various stages of our work on this subject, as well as to L. Bureau and O. Ronsin for a long lasting collaboration, and fruitful exchanges on these and related subjects.
## Appendix A : Shear stiffness of a MCI
Already long ago, a number of studies have tried to get independent information about the structure of MCI from the load dependence of various interfacial properties besides friction. These include electrical Tabor and thermal thermal conductances, variations of closure with normal load closure , acoustic reflection and interfacial normal and transverse mechanical stiffnesses ultrasound . In order to confront experimental results with the Greenwood-Williamson model, each of these properties must be calculated within this same frame.
Let us sketch out, as an example, the evaluation of the interfacial shear stiffness. For a single Hertz contact of radius $`a`$ between two identical spheres of shear modulus $`G`$ and Poisson ratio $`\nu `$, it was calculated by Mindlin Mindlin Johnson to be :
$$k_s=\frac{4Ga}{2\nu }$$
(57)
which expresses the fact that the strain energy due to shearing is essentially localized, in each medium, within a volume of section $`a^2`$ and height $`a`$. So, the relative shear displacement involved in the definition of $`k_s`$ is that between any two points on both sides of the contact at distances $`a`$ from it.
Consider a set of $`n`$ GW microcontacts between two blocks of thickness $`D`$, shear modulus $`G`$. One may separate this elastic system into an interfacial layer of thickness $`a`$ but $`D`$ sandwiched between two bulk regions of thickness $`D`$. As far as $`\mathrm{\Sigma }_R\mathrm{\Sigma }_{app}`$, the microcontacts are elastically independent, and the total interfacial stiffness $`K_s`$ is simply the sum of the individual ones, i.e., with the notations of $`\mathrm{\S }`$ II.C :
$$K_s=\frac{4G}{2\nu }n\overline{a}$$
(58)
From equations (8), (9)
$$K_s=\frac{4G}{(2\nu )\sqrt{\pi }E}\frac{W}{s}=\frac{2}{(2\nu )(1+\nu )}\frac{W}{s}$$
(59)
So, the GW model predicts, contrary to immediate intuition, that the interfacial stiffness of an elastic MCI is independent of the bulk material shear modulus.
Of course, experiments measure the compound response of the interface and the bulk. A more important warning comes from the fact that expression (57) is valid only in the limit of very small ratios of shear to normal forces. Indeed, Mindlin has shown Johnson that increasing shear induces, at the periphery of the contact, a slipping annulus along which, would it remain non-slipping, the friction sliding threshold would be overcome. This effect results in a non-linear weakening behavior of $`k_s`$ Johnson and of the MCI response Lionel Proc .
Consider now a GW interface in the fully plastic limit. When a small shear force is applied for the first time, it induces some further plastic flow, probably leading to some small adaptation of contact radii. Under subsequent unloading-reloading cycles, these adapted contacts respond elastically, with an individual stiffness $`Ga`$. $`K_s`$ is then immediately evaluated to be $`W/\lambda `$, where the length $`\lambda `$ is of the order of $`s/\psi `$, with $`\psi `$ the plasticity index defined in section II.C Pauline Proc Roy Soc .
The experiments of Berthoud and Baumberger Pauline Proc Roy Soc on PMMA (a polymer glass)/PMMA and AU4G (an aluminum alloy)/AU4G have confirmed the Amontons-like behavior (59) of $`K_s`$. The lengths $`\lambda `$ for the two systems were different, but both on the order of a few microns. They were found to be compatible, in order of magnitude, with the values predicted from the above model. These results confirm the robustness of the three properties of MCI, listed in section II.C, which emerge from the GW model.
In view of the difficulty of determining accurately $`s`$ and, even more, $`R`$, as well as of the numerous approximations involved in the GW model and, even more, of its elasto-plastic extension, it would be fallacious to expect a more quantitative agreement.
Amontons-like behaviors have also been measured to hold, with various degrees of accuracy depending on the specific difficulties inherent to each experimental method, on the other interfacial properties listed at the beginning of this Appendix.
## Appendix B : The viscoelastic GW model
Our derivation summarizes the work of Hui et al Hui .
Let us consider a MCI between a linear viscoelastic, incompressible and isotropic medium and a rigid one – a good approximation for contact between a rough elastomer and a rough hard solid. We want to describe this MCI in the the frame of the simplest version of the GW model (exponential distribution of summit heights).
In a linear viscoelastic material, shear stress and strain are related via a retarded elastic shear modulus $`G(t)`$, by the constitutive equation :
$$\sigma (t)=_{\mathrm{}}^t𝑑t^{}G(tt^{})\frac{dϵ(t^{})}{dt^{}}$$
(60)
which can be formally inverted into
$$ϵ(t)=_{\mathrm{}}^t𝑑t^{}J(tt^{})\frac{d\sigma (t^{})}{dt^{}}$$
(61)
$`J(t)`$, which characterizes the strain response to a unit stress step applied at $`t=0`$, increases with time, as a result of the progressive relaxation of the constitutive polymer molecules.
Let us first consider the viscoelastic analogue of the single Hertz contact ($`\mathrm{\S }`$ II.A.2), created by applying the load $`w`$ at time $`t=0`$. It is intuitively clear that, since the compliance increases, the compression $`\delta (t)`$ and the contact radius $`a(t)`$ increase monotonously. Lee and Radok Lee-Radok have shown (see also Johnson ) that in such a case the solution of the contact problem is that deduced from the Hertz one by simply replacing the shear modulus by its retarded counterpart $`G(t)`$. So, the Hertz load-radius relation (eq. (3)) becomes :
$$w=_{\mathrm{}}^tG(tt^{})\frac{d}{dt^{}}\left[\frac{4}{\pi R^{}}a^3(t^{})\right]$$
(62)
with $`a^2(t^{})=R^{}\delta (t^{})`$.
We now turn to the corresponding GW problem, contact under the normal load $`W`$ being first established at $`t=0`$. As time increases, clearly, the interfacial separation $`d(t)`$ will decrease monotonously, the number of microcontacts $`n(t)`$ and their individual area increase, leading to the increase of the real area of contact $`\mathrm{\Sigma }_r(t)`$.
The three GW equations (eqs.(5) - (7)) become :
$$n(t)=N_{d(t)}^{\mathrm{}}𝑑z\frac{1}{s}e^{z/s}$$
(63)
$$\mathrm{\Sigma }_r(t)=N_{d(t)}^{\mathrm{}}𝑑z\pi R(zd(t))\frac{1}{s}e^{z/s}$$
(64)
$$W(t)=N_{d(t)}^{\mathrm{}}𝑑z\frac{1}{s}e^{z/s}_{\mathrm{}}^t𝑑t^{}\theta (t^{}t_0(z))G(tt^{})\frac{d}{dt^{}}\left(\frac{4}{3}R^{1/2}\left(zd(t^{})\right)^{3/2}\right)$$
(65)
where $`\theta `$ is the unit step function, $`W(t)=W\theta (t)`$, and the function $`t_0(z)`$ is the inverse of $`d(t)`$, i.e.
$$d\left(t_0(z)\right)=z$$
(66)
Equation (65) is based upon relation (62) and states that each contact with summit height $`z`$ existing at time $`t`$ has contributed since $`t^{}=t_0(z)`$ when it first appeared, i.e. for $`t^{}>t_0(z)`$ or equivalently, $`d(t^{})<z`$.
Setting $`y=\left(zd(t^{})\right)/s`$ and performing the $`y`$-integration in eq.(65) yields :
$$W(t)=N\left(\pi Rs^3\right)^{1/2}_{\mathrm{}}^t𝑑t^{}G(tt^{})\frac{d}{dt^{}}\left(e^{d(t^{})/s}\right)$$
(67)
which can be inverted, with the help of definitions (60), (61) into :
$$e^{d(t)/s}=\frac{1}{\left(\pi Rs^3\right)^{1/2}}_{\mathrm{}}^t𝑑t^{}J(tt^{})\frac{dW(t^{})}{dt^{}}$$
(68)
On the other hand, from equation (64)
$$\mathrm{\Sigma }_r(t)=N\pi Re^{d(t)/s}$$
(69)
so that, finally, since $`dW/dt^{}=\delta (t^{})`$, the real area of contact evolves according to :
$$\mathrm{\Sigma }_r(t)=\sqrt{\frac{\pi R}{s}}J(t)W$$
(70)
## Appendix C : The driven block : Linear stability analysis
Consider the system defined in $`\mathrm{\S }`$II.D.2, made of a block of mass $`M`$, carrying the normal load $`W`$, driven at the constant velocity $`V`$ through a spring of stiffness $`K`$ <sup>31</sup><sup>31</sup>31We assume that, as is commonly the case, $`K`$ is much smaller than the interfacial stiffness. Finite interfacial compliance effects are evaluated in Cochard .. The equations describing the motion read :
$$M\ddot{x}=K\left(xVt\alpha \right)W\mu _d(\varphi ,\dot{x})$$
(71)
$$\dot{\varphi }=1\frac{\dot{x}\varphi }{D_0}$$
(72)
In the RR model :
$$\mu _d=\mu _{d0}+B\mathrm{ln}\frac{\varphi }{\varphi _0}+A\mathrm{ln}\frac{\dot{x}}{V_0}$$
(73)
with $`V_0`$ some reference velocity, $`\varphi _0=D_0/V_0`$. However, in view of the shortcomings of expression (73) discussed in $`\mathrm{\S }`$II.D, we will leave the functional form of $`\mu _d`$ unspecified at this stage.
Equations (71), (72) have the stationary solution :
$$x_{st}(t)=Vt\alpha \frac{W}{K}\mu _d(\frac{D_0}{V},V)$$
(74)
$$\varphi _{st}=\frac{D_0}{V}$$
(75)
In order to study its (linear) stability, we linearize equations (71), (72) in $`\delta x(t)=x(t)x_{st}(t);\delta \varphi (t)=\varphi (t)\varphi _0`$. Then :
$$\frac{M}{W}\delta \ddot{x}+\frac{K}{W}\delta x+\frac{\mu _{\dot{x}}}{V}\delta \dot{x}+\frac{V\mu _\varphi }{D_0}\delta \varphi =0$$
(76)
$$\delta \dot{\varphi }+\frac{V}{D_0}\delta \varphi +\frac{1}{V}\delta \dot{x}=0$$
(77)
with$`\mu _\varphi =\mu _d/(\mathrm{ln}\varphi )`$, $`\mu _{\dot{x}}=\mu _d/(\mathrm{ln}\dot{x})V`$, both quantities being evaluated at $`\varphi =D_0/V,\dot{x}=V`$. Note that we saw that, for slowly sliding MCI, $`\mu _{\dot{x}}>0`$, $`\mu _\varphi <0`$, and $`\mu _\varphi \mu _{\dot{x}}>0`$. The solutions of this system are of the form :
$$\left(\genfrac{}{}{0pt}{}{\delta x}{\delta \varphi }\right)=\left(\genfrac{}{}{0pt}{}{\xi }{\psi }\right)e^{i\mathrm{\Omega }t}$$
(78)
and the frequencies $`\mathrm{\Omega }`$ of these eigenmodes are the roots of :
$$i\frac{M}{W}\mathrm{\Omega }^3\left(\frac{\mu _{\dot{x}}}{V}+\frac{MV}{WD_0}\right)\mathrm{\Omega }^2+\left(\frac{K}{W}+\frac{\mu _{\dot{x}}\mu _\varphi }{D_0}\right)i\mathrm{\Omega }+\frac{K}{W}\frac{V}{D_0}=0$$
(79)
If, for all roots, $`Im\mathrm{\Omega }>0`$, all (infinitesimal) fluctuations about the stationary state decay, steady sliding is (locally) stable. For very large $`K`$ the solutions of (79) read :
$$\mathrm{\Omega }_1\frac{iV}{D_0}\mathrm{\Omega }_\pm =\pm \sqrt{\frac{K}{M}}+(1+3i)\frac{W\mu _{\dot{x}}}{MV}$$
(80)
Since $`\mu _{\dot{x}}>0`$, the system is stable in the large $`K`$ limit. As the stiffness decreases, it becomes unstable if, and when, the imaginary part of one at least of the roots vanishes. One easily checks on equation (79) that this occurs when :
$$\frac{K}{W}=\left(\frac{K}{W}\right)_c=\left(\mu _\varphi \mu _{\dot{x}}\right)\left(1+\frac{MV^2}{WD_0\mu _{\dot{x}}}\right)$$
(81)
As expected, instability only occurs when $`\mu _\varphi \mu _{\dot{x}}>0`$, i.e. (see $`\mathrm{\S }`$II) when the steady sliding $`\mu _d`$ is velocity-weakening, which is indeed the case for our systems in the $`V`$-range under consideration.
Note that inertia only comes into play in eq.(81) via the quantity $`\theta =(MV^2/WD_0\mu _{\dot{x}})`$. For ablock sliding under its own weight ($`W=Mg`$), with $`D_01\mu `$m, $`\mu _{\dot{x}}10^2`$, $`\theta 10^5\left(V_{\mu m/sec}\right)^2`$ can safely be neglected in our low velocity regime, and the position of the bifurcation line is simply given by :
$$\left(\frac{K}{W}\right)_c=\mu _\varphi \mu _{\dot{x}}$$
(82)
For the RR model, $`\mu _\varphi \mu _{\dot{x}}=BA`$ is a constant, $`(K/W)_c`$ is $`V`$-independent.
At the bifurcation, the two (complex conjugate) neutral modes oscillate at the critical frequency :
$$\mathrm{\Omega }_c=\frac{V}{D_0}\sqrt{\frac{\mu _\varphi }{\mu _{\dot{x}}}1}$$
(83)
That is, the bifurcation is of the Hopf type. A third order perturbation expansion can be performed standardly provided that $`\mu _\varphi `$ and $`\mu _{\dot{x}}`$ are not mere constants (see $`\mathrm{\S }`$II). It shows that the stick-slip bifurcation is direct - i.e. that the SS amaplitude grows continuously when decreasing e.g. $`K`$ below its critical value.
Note finally that, for a velocity-strengthening $`\mu _d`$, ($`\mu _\varphi \mu _{\dot{x}})<0`$, the system is always stable against infinitesimal perturbations. This does not preclude the possibility of a finite amplitude instability. Brockley Brockley has shown that – in the case where $`\mu _\varphi =0`$ – when one takes inertia into account, the system does exhibit a strongly hysteretic Hopf bifurcation. As explained in $`\mathrm{\S }`$II, one expects this situation to prevail, for rough-on-rough MCI, at large enough driving velocities ($`V>V_{min}`$) for which geometric aging becomes inactive. A Brockley-like regime has indeed been observed TBTrieste , for a paper/paper interface, in the cm/sec range.
## Appendix D : The Schallamach model of adsorption-controlled friction
Following Schallamach’s seminal article Schall , consider an extended interface between a soft slider and a hard flat track covered with adhesive sites which can pin the slider molecules. Represent the slider as a set of $`N_0`$ identical and independent chains of stiffness $`\kappa `$, potentially forming bonds by adsorption of their end monomer onto the track. Adsorption and subsequent desorption are thermally activated, and, when sliding, desorption is aided by advection. The number $`N`$ of bonds is therefore expected to be a function of the sliding velocity $`V`$ and of the temperature $`T`$. In steady motion, the elastic force exerted on a given bond increases linearly with time until the bond snaps off and the stored energy is dissipated.
The frictional force thus reads :
$$F=N\kappa V\overline{t}$$
(84)
where $`\overline{t}`$ is the average lifetime of a bond.
Let $`\tau `$ be the average time for which a chain remains depinned before readsorbing. The stationary number of bonds is :
$$N=N_0\frac{\overline{t}}{\overline{t}+\tau }$$
(85)
and the friction force
$$F=N_0\kappa V\tau \frac{(\overline{t}/\tau )^2}{1+(\overline{t}/\tau )}$$
(86)
In order to desorb (resp. adsorb), a chain must overcome an energy barrier of height $`\mathrm{\Delta }E=W+E`$ (resp. $`E`$), as sketched on Figure 34. When the pinned chain is stretched at velocity $`V`$, the barrier is lowered and its effective height becomes $`\mathrm{\Delta }E^{}=\mathrm{\Delta }E\kappa bVt`$, where $`b`$ is a length of atomic order. Schallamach assumes that the resulting desorption rate is :
$$r(t)=\frac{1}{\tau _0}\mathrm{exp}\left[\frac{\mathrm{\Delta }E^{}(t)}{k_BT}\right]$$
(87)
where $`\tau _0^1`$ is an attempt frequency, and $`t`$ the age of the bond.
Note that this expression for $`r`$ tacitly assumes that thermal activation is efficient enough for all bonds to break before reaching the deterministic threshold such that $`\mathrm{\Delta }E^{}=0`$, i.e. that $`\mathrm{\Delta }E^{}(\overline{t})k_BT`$. We will discuss this assumption in more detail below.
The time $`\tau `$ is that for thermal activation above the barrier $`E`$, of unspecified origin, and Schallamach sets :
$$\tau =\tau _0\mathrm{exp}\left(E/k_BT\right)$$
(88)
an expression which neglects in particular advection effects.
Consider a set of $`n_0`$ bonds all formed at the same time $`t=0`$. At time $`t`$, the number of these, $`n(t)`$, which have not desorbed obeys :
$$\frac{dn}{dt}=r(t)n$$
(89)
The average bond lifetime $`\overline{t}`$ then reads, with the help of equations (87),(89) :
$$\overline{t}=_0^{\mathrm{}}\frac{n(t)}{n_0}𝑑t=\tau _{out}\frac{V_0}{V}_0^{\mathrm{}}\frac{e^y}{y+(V_0/V)}𝑑y$$
(90)
where $`\tau _{out}=\tau _0\mathrm{exp}\left(\mathrm{\Delta }E/k_BT\right)`$ is the desorption time of the unstretched chain, and $`V_0=k_BT/\kappa b\tau _{out}`$.
This, together with expression (86), yields the friction force in the steady state. More precisely, in the small and large velocity limits :
$$\overline{t}/\tau _{out}=1\frac{V}{V_0}+\mathrm{}(VV_0)$$
(91)
$$\overline{t}/\tau _{out}=\frac{V_0}{V}\left[\mathrm{ln}\frac{V}{V_0}0.577+\mathrm{}\right](VV_0)$$
(92)
$`\overline{t}`$ decreases monotonously with $`V`$, while the elastic energy $`\kappa bV\overline{t}`$ stored before debonding increases from linearly ($`VV_0`$) to logarithmically ($`VV_0`$) as shown on Figure 35. Also shown on this figure is the average number of bonds in steady state $`N(V)`$, which decreases all the more slowly to zero in the large $`V`$ limit that $`\tau /\tau _{out}`$ is small.
As a result $`F(V)`$ exhibits a maximum at $`V=V_{max}`$ (Figure 35) resulting from the interplay of two effects :
– the decrease with $`V`$ of the number of bonds.
– the increase of the average pinning force $`\overline{f}_p=\kappa V\overline{t}`$.
In the limiting regimes :
$$F(V)N_0\kappa \frac{\tau _{out}^2}{\tau +\tau _{out}}V(VV_0)$$
(93)
$$F(V)N_0\frac{\tau _{out}}{\tau }\frac{k_BT}{b}\frac{V_0}{V}\mathrm{ln}^2\frac{V}{V_0}(VV_{max})$$
(94)
Let us now discuss in more detail Schallamach’s results and assumptions. Two weak points in the theory are immediately clear.
$``$ It predicts vanishing friction at large $`V`$, since, in this limit, complete depinning is achieved. The junction is then a mere sheared liquid layer, whose viscous dissipation has been neglected. The corresponding contribution $`F_{vis}(V)`$ to the total friction force $`F_{tot}(V)=F(V)+F_{vis}(V)`$ becomes dominant for $`VV_{max}`$. Its precise expression of course depends on the nature of the slider material (e.g. hydrogel vs elastomer).
$``$ As mentioned above, Schallamach’s expression for the desorption rate $`r`$ is valid only as long as $`\mathrm{\Delta }E^{}(\overline{t})k_BT`$. Since equation (74) results in an unboundedly growing value of $`\kappa bV\overline{t}`$ (see Figure 35), this assumption fails for $`VV^{}`$ such that $`\mathrm{\Delta }E=\kappa bV^{}\overline{t}(V^{})`$, which yields $`V^{}V_0\mathrm{exp}(\mathrm{\Delta }E/k_BT)`$. In this fast regime, we are back to the scenario of $`\mathrm{\S }`$ II.C.4. Advection-controlled deterministic bond breaking becomes more efficient than “premature” thermally activated depinning, resulting in a value of $`\overline{f}_p`$ which crosses over from the Schallamach logarithmic regime to the deterministic saturation value $`f_{sat}=\mathrm{\Delta }E/\kappa b`$. That is, as shown on Figure 35, $`\overline{f}_p`$ exhibits three regimes : an initial linear one for $`V\stackrel{~}{V}`$, a logarithmic one for $`\stackrel{~}{V}VV^{}`$, saturation for $`V>V^{}`$. The lower crossover $`\stackrel{~}{V}`$ can be evaluated as $`\stackrel{~}{V}V_0(\mathrm{\Delta }E/k_BT)\mathrm{\Delta }E/\kappa b\tau _{out}`$.
Keeping track of these two corrections results in a total friction force curve whose qualitative shape depends on the order of magnitude of the parameter $`\tau /\tau _{out}`$. Note that, within Schallamach’s assumption (equation (75)) $`\tau /\tau _{out}<1`$.
If $`\tau /\tau _{out}1`$, one easily evaluates $`V_{max}V_0\tau _{out}/\tau `$, so that
$$\frac{V}{V_{max}}\frac{\mathrm{\Delta }E}{k_BT}\frac{\tau }{\tau _{out}}\frac{\mathrm{\Delta }E}{k_BT}e^{W/k_BT}1$$
and $`F(V)`$ exhibits a wide plateau followed by a slow decrease. Whether or not the total force $`F_{tot}=F+F_{vis}`$ exhibits a $`𝒩`$-shape for $`V>V_{max}`$ depends on the effective viscosity of the unpinned junction and on the density of pinning sites $`N_0`$.
As $`\tau /\tau _{out}`$ increases towards $`1`$, the width of the plateau decreases with $`V_{max}`$, $`F(V)`$ becomes steeper on the high-$`V`$ side, so that the more likely a $`𝒩`$-shaped $`F_{tot}(V)`$, hence a $`V`$-weakening friction regime leading to stick-slip.
Finally, it is worth recalling that, when the system is loaded from rest at a prescribed velocity $`V_{load}`$ larger than, typically, $`\stackrel{~}{V}`$, its response is initially quasi-elastic, the upper limit of this regime appearing as a static threshold. Since $`\stackrel{~}{V}\mathrm{\Delta }E/\kappa b\tau _{out}`$, reasonable estimates lead to conclude that in order for $`\stackrel{~}{V}`$ to lie in the sub-$`\mu `$m/sec range, adsorption should be quite strong – corresponding to binding energies $`W`$ not far below the eV level. |
warning/0506/cond-mat0506519.html | ar5iv | text | # Critical Casimir Effect in 3He -4He films
## Abstract
Universal aspects of the thermodynamic Casimir effect in wetting films of <sup>3</sup>He-<sup>4</sup>He mixtures near their bulk tricritical point are studied within suitable models serving as representatives of the corresponding universality class. The effective forces between the boundaries of such films arising from the confinement are calculated along isotherms at several fixed concentrations of <sup>3</sup>He. Nonsymmetric boundary conditions impose nontrivial concentration profiles leading to repulsive Casimir forces which exhibit a rich behavior of the crossover between the tricritical point and the line of critical points. The theoretical results agree with published experimental data and emphasize the importance of logarithmic corrections.
Finite-size contributions to the free energy of a fluid confined between two surfaces at a distance $`L`$ give rise to an effective force between them. Theory predicts that at the bulk critical point $`T_c`$ of such a system this force becomes long-ranged as a result of critical fluctuations of the corresponding ordering degrees of freedom. This is analogous to the well-known Casimir effect in electromagnetism. This so-called critical Casimir force $`f_C`$ per unit area and in units of $`k_BT_c`$ can be expressed in terms of universal scaling functions krech:99:0 .
Only recently, sophisticated wetting experiments have provided detailed quantitative data for critical Casimir forces in various systems garcia:99:0 ; law:99:0 ; garcia:02:0 ; balibar:02:0 ; pershan . In the case of <sup>4</sup>He wetting films near the superfluid transition, these experimental studies support quantitatively theoretical predictions for $`f_C`$ ($`TT_c`$) corresponding to the universality class of the $`XY`$ model krech:91:0 . For the case of <sup>3</sup>He-<sup>4</sup>He films near the bulk tricritical point some theoretical predictions are available krech:91:0 , but those do not apply for the boundary conditions relevant for recent wetting experiments performed in these systems garcia:02:0 ; balibar:02:0 . However, the shape of the scaling function of the Casimir force depends sensitively on the type of boundary conditions (BC) and thus on the surface universality classes to which the confining surfaces belong diehl:86:0 . The experiments of Ref. garcia:02:0 report a repulsive $`f_C`$ around the tricritical point which suggests nonsymmetric BC for the superfluid order parameter (SOP). This is opposite to the case of pure <sup>4</sup>He wetting films near the $`\lambda `$-point where $`f_C`$ was found to be attractive garcia:99:0 ; krech:91:0 . For the latter system the BC seem to be very well approximated by symmetric Dirichlet-Dirichlet BC $`(O,O)`$ forming the so-called ordinary (O) surface universality class because the quantum-mechanical wave function describing the superfluid state vanishes at both interfaces krech:99:0 ; garcia:99:0 . The type of BC for <sup>3</sup>He-<sup>4</sup>He wetting films is not clear from the outset because a <sup>4</sup>He-rich layer forms near the substrate-fluid interface, which may become superfluid already above the bulk $`\lambda `$-line laheurte:78:0 whereas <sup>3</sup>He has a preference for the fluid-vapor interface. Thus the two interfaces impose a nontrivial concentration profile which in turn couples to the SOP. This leads to the hypothesis that the concentration profile induces effectively nonsymmetric $`(O,+)`$ BC for the SOP, i.e., Dirichlet boundary conditions at the fluid-vapor interface and symmetry-breaking (+) BC at the substrate-fluid interface (also known as the so-called extraordinary or normal universality class diehl:86:0 ). For the present tricritical behavior the upper critical dimension $`d^{}`$ equals 3. In this case theory predicts that for three-dimensional systems the asymptotic tricritical thermodynamic functions exhibit power laws with critical exponents taking their classical values. However, logarithmic corrections to the mean-field (MF) behavior are expected under experimental conditions LawSar .
Here we consider two complementary approaches. Field-theoretical methods and renormalization-group (RG) analyses are used to derive universal properties of the Casimir force at the tricritical point and the form of logarithmic corrections. However, these methods do not lend themselves for systematic studies of $`f_C`$ along all thermodynamic paths followed in the aforementioned experiments. In order to be able to interpret the rich variation of $`f_C`$ extracted from the capacity measurements in Ref. garcia:02:0 , to understand the emergence of the actual BC and, moreover, to predict the behavior of $`f_C`$ in the crossover region between the tricritical and the critical points, we employ the vectoralized Blume-Emery-Griffiths model (VBEG) maciolek:04:0 as a representative of the same universality class as the actual physical system. This lattice model is extended to the film geometry and treated within mean field theory (MFT).
First we derive the leading asymptotic behavior of $`f_C`$ at tricriticality for $`(O,+)`$ BC. To this end we consider the standard Ginzburg-Landau (GL) Hamiltonian for an $`O(n)`$-symmetric tricritical system ($`T=T_t`$) in a film geometry:
$$[𝚽]=d^{d1}x_0^L𝑑z\left\{\frac{1}{2}(𝚽)^2+\frac{u}{6!}(𝚽^2)^3\right\}$$
(1)
where $`L`$ is the film thickness, $`𝚽`$ is the $`n`$-component OP and $`z`$ is the distance between the confining surfaces; $`u`$ is a bare coupling constant. In a film geometry $`f_C(f^{ex}/L)=𝒯_{zz}`$ is given by the stress tensor component $`𝒯_{zz}`$ krech:99:0 , where $`f^{ex}(L)(ff_b)L`$, $`f`$ is the total free energy per unit area and per $`k_BT_t`$ and $`f_b`$ is the bulk contribution. The stress tensor is given by krech:99:0 $`𝒯_{ij}=_i𝚽_j𝚽\delta _{ij}(d2)/(4(d1))(_i_j\delta _{ij}^2)𝚽^2`$, where $``$ is the integrand of (1). We take $`𝚽=(m(z),0,\mathrm{},0)`$. Determination of the tricritical Casimir force starts from the Euler-Lagrange equation for the OP profile: $`m^{\prime \prime }(z)=(u/120)m^5(z)`$ with $`(O,+)`$ BC, i.e., $`m(0)=0\text{and}m(L)=+\mathrm{}`$. In this case the spatially constant $`𝒯_{zz}`$ can be expressed as $`(1/2)(m^{}(0))^2`$. With the scaling ansatz $`m(z)=(u/360)^{1/4}L^{1/2}\phi (z/L)`$ and $`T_{zz}=(90/u)^{1/2}L^3\mathrm{\Theta }`$, and after integrating directly the first integral of the Euler-Lagrange equation one obtains $`\mathrm{\Theta }^{1/3}={\displaystyle _0^{\mathrm{}}}𝑑p/\sqrt{1+p^6}1.40218`$ by implementing the above BC. Eventually, in units of $`k_BT_t`$ the MFT result for the tricritical Casimir force $`f_C^t`$ in the case of $`(O,+)`$ BC is $`f_C^t=2.7568(4)\left(90/u\right)^{1/2}L^3`$. Note that within MFT the parameter $`u>0`$ remains undetermined. Its value follows from using standard RG arguments. In $`d=3ϵ`$ the above MFT result yields the leading contribution in an $`ϵ`$-expansion. After removing the uv singularity via renormalization the asymptotic scaling behaviour of $`f_C^t`$ follows from substituting $`u`$ by the appropriate fixed-point value $`u^{}ϵ`$. At $`d=d^{}`$, and under spatial rescaling by a factor $`\mathrm{}`$, $`u`$ flows to its RG fixed point value $`u^{}=0`$ according to $`\overline{u}(\mathrm{})=(240\pi ^2)/\left((3n+22)|\mathrm{ln}\mathrm{}|\right)`$ eisen:88 . With the rescaling factor $`\mathrm{}=l_0/L`$, where $`l_0`$ is a microscopic length scale of the order of a few Å, this yields a logarithmic correction to the power law $`L`$-dependence of the tricritical Casimir force:
$$f_C^t0.54(3n+22)^{1/2}(\mathrm{ln}L/l_0)^{1/2}L^3.$$
(2)
Gaussian fluctuations give contributions of at least $`O(u^0)`$ which are therefore subdominant. We compare Eq. (2) for $`n=2`$ with the data obtained by Garcia and Chan garcia:02:0 for their experimental value $`L/l_0`$ 520 Å/1.3 Å. This gives $`\vartheta _tf_C^tL^36.96`$ in a good agreement with $`\vartheta _t^{exp}=8.4\pm 1.7`$, which suggests that this experiment maybe the first to have verified implicitly the existence of logarithmic corrections near the tricritical point. However, in order to extract the actual value of the universal Casimir amplitude (i.e., the numerical prefactor in Eq. (2)) the experimental data require a reanalysis based on a functional form given by Eq. (2).
Now we turn to the VBEG model and consider a $`d=3`$ simple cubic lattice consisting of $`L`$ parallel lattice layers at spacing $`a`$. Each layer has $`A`$ sites, labeled $`i,j,\mathrm{}`$ and associated with an occupation variable $`t_i=0,1`$ and a phase $`\theta _i`$ $`(0\theta _i<2\pi )`$ which mimics the phase of the <sup>4</sup>He wave function. A <sup>3</sup>He (<sup>4</sup>He) atom at site $`i`$ corresponds to $`t_i=0(1)`$. The Hamiltonian is given by
$$=J\underset{<ij>}{}t_it_j\mathrm{cos}(\theta _i\theta _j)K\underset{<ij>}{}t_it_j+\underset{i}{}\mathrm{\Delta }_it_i$$
(3)
where the first two sums run over nearest-neighbor pairs and the last one is over all lattice sites. The field $`\mathrm{\Delta }_i`$ is related to the chemical potentials of the two components of the mixture. $`\mathrm{\Delta }_i=\mathrm{\Delta }_1`$ and $`\mathrm{\Delta }_i=\mathrm{\Delta }_2`$ on the left and right surface layer, respectively, and $`\mathrm{\Delta }_i=\mathrm{\Delta }`$ otherwise. The differences $`\mathrm{\Delta }_i\mathrm{\Delta },i=1,2`$, are a measure of the relative preferences of <sup>4</sup>He atoms for the two surfaces such that $`\mathrm{\Delta }_1<\mathrm{\Delta }`$ corresponds to the preference of <sup>4</sup>He atoms for the solid substrate. At the opposite surface with the vapor we choose $`\mathrm{\Delta }_2=\mathrm{\Delta }_t`$, the tricritical bulk value. Near the tricritical point this choice is consistent with the assumption made in Ref. garcia:02:0 for the concentration profile across the wetting film, whereby at the interface with the vapor the <sup>3</sup>He mole fraction takes the bulk value. Bulk properties of this model were studied within MFT and by Monte Carlo simulations in $`d=3`$ in maciolek:04:0 . The resulting bulk phase diagram resembles that observed experimentally for <sup>3</sup>He-<sup>4</sup>He mixtures (see Fig. 1). For the film geometry we have solved this model within MFT. This yields a set of self-consistent equations for the OP in the $`l`$th layer, i.e., the concentration $`Q_l=<t_l>=1X(l)`$ of <sup>4</sup>He, and the SOP $`𝐌_l=(m_l^{(1)},m_l^{(2)})`$ with $`m_l^{(1)}=t_l\mathrm{cos}\theta _l`$ and $`m_l^{(2)}=t_l\mathrm{sin}\theta _l`$ where $`t_l`$ and $`\theta _l`$ denote the occupation number and the phase in the $`l`$th layer, respectively. In the absence of helicity one has $`𝐌_l=(m_l^{(1)},0)(m_l,0)`$, $`Q_l=I_0(\stackrel{~}{J}b_l)/\left(e^{\stackrel{~}{K}a_l+\mathrm{\Delta }_l/T}+I_0(\stackrel{~}{J}b_l)\right)`$, and $`m_l=I_1(\stackrel{~}{J}b_l)/\left(e^{\stackrel{~}{K}a_l+\mathrm{\Delta }_l/T}+I_0(\stackrel{~}{J}b_l)\right)`$. $`I_0(z)`$ and $`I_1(z)`$ are modified Bessel functions and $`\mathrm{\Delta }_l=\mathrm{\Delta }`$ for $`l1,L`$; $`\stackrel{~}{K}=qK/T`$ and $`\stackrel{~}{J}=qJ/T`$, where $`q=q_{}+2q^{}`$ is the coordination number of the lattice. $`q_{}`$ is the in-layer coordination number while each site (but not in the first and last layers) is connected to $`q^{}`$ atoms in each adjacent layer. We have introduced $`b_lm_{l1}+q_{}m_l+m_{l+1}`$ for $`l1,L`$, $`b_1q_{}m_1+m_2`$, and $`b_Lm_{L1}+q_{}m_L`$ and analogously $`a_lQ_{l1}+q_{}Q_l+Q_{l+1}`$ for $`l1,L`$, $`a_1=q_{}Q_1+Q_2`$, and $`a_L=Q_{L1}+q_{}Q_L`$. The coupled equations for $`Q_l`$ and $`m_l`$ are solved numerically; the acceptable solution minimizes the free energy. First, we have analyzed the semi-infinite system. Close to the $`\lambda `$-line we observe a higher <sup>4</sup>He concentration near the left surface, which induces a local superfluid ordering (see Fig. 2). By varying $`T`$ and $`\mathrm{\Delta }`$ one obtains a whole line of continuous surface transitions corresponding to the onset of the formation of a superfluid film near the wall; it meets the $`\lambda `$-line at the special transition point whose position depends on the value of $`\mathrm{\Delta }_1`$ (see Fig. 1). These findings are in agreement with the results of a Migdal-Kadanoff analysis peliti:85:0 . In the film geometry the Casimir force is obtained by taking a finite difference after calculating $`f^{ex}`$ for $`L_0`$ and $`L_0+1`$. Figure 3 summarizes our result for a film of width $`L=20`$, $`K/J=0.5`$, $`\mathrm{\Delta }_1/J=3`$, and $`\mathrm{\Delta }_2=\mathrm{\Delta }_t/J0.61`$. $`f_C`$ is calculated along the thermodynamic paths indicated in Fig. 1. Below $`T_t`$, $`f_C`$ is calculated along the coexistence line, infinitesimally on the superfluid branch of bulk coexistence. Our results are presented in terms of the scaling function $`\vartheta L^df_C`$ as a function of the scaling variable $`ytL^{1/\nu }=(L\xi _0/\xi )^{1/\nu }`$, where $`t=(TT_t)/T_t`$. $`\xi _0`$ is the amplitude of the correlation length $`\xi `$ and $`\nu =1`$. The surface transition does not leave a visible trace in the behavior of $`f_C`$. For <sup>3</sup>He concentration $`X<X_t`$, upon crossing the $`\lambda `$-line there is a steep variation associated with a break in slope, giving rise to the formation of shoulders which are similar to those observed experimentally garcia:02:0 . For $`X>X_t`$, when $`T`$ reaches the phase separation temperature $`f_C`$ coincides with the curve common to all values of $`X`$. This occurs with a discontinuous first derivative. The aforementioned common curve exhibits a pronounced maximum below $`T_t`$ at $`y0.74`$ and gradually decreases to zero for $`y\mathrm{}`$. Below $`T_t`$, both the concentration and the SOP profiles corresponding to this common curve display an interface-like structure separating two domains of the coexisting bulk phases (see $`t=0.0625`$ in Fig. 2). Features of $`f_C`$ in this ’soft mode’ phase can be attributed to purely interfacial effects, similarly to Ising-like films with asymmetric BC parry:92:0 . Beyond MFT a positive sign of the force can be regarded as a consequence of entropic repulsion fisher:86:0 . The maximum of $`f_C`$ is expected to occur at the temperature $`T`$ for which the interfacial width $`\xi _bL`$, i.e., $`y1`$ parry:92:0 which checks with Fig. 3 footnote . For $`XX_t0.05`$ we observe a crossover to the critical superfluid behavior of pure <sup>4</sup>He and a gradual formation of a second, less pronounced local maximum located slightly below the $`\lambda `$-line. This local maximum decreases upon departure from $`X_t`$ and finally disappears above the special transition $`S`$. This is expected, because above $`S`$ the BC turn into the type $`(O,O)`$ for which $`f_C`$ vanishes within MFT. For lower $`T`$, $`f_C`$ increases steeply upon approaching bulk coexistence revealing that interfacial effects associated with the ’soft mode’ lead to a much stronger Casimir effect than near the critical $`\lambda `$-line. The qualitative features of $`\vartheta `$ extracted from the experimental data for $`XX_t`$ (see Fig. 5 in Ref. garcia:02:0 ) are very well captured by the present lattice model. Discrepancies can be attributed to fluctuation effects neglected in the present MFT VBEG approach: (i) The discontinuities of slopes as obtained within MFT upon crossing the $`\lambda `$-line are expected to be smeared out by fluctuations. (ii) The experimental scaling function $`\vartheta `$ does not vanish at low temperatures, which may be due to Goldstone modes in the superfluid phase. However, the possibility that this behavior is an artifact of an extreme change in the dielectric constant of the film cannot be excluded privat . In the crossover regime to the critical behavior only few experimental data for the thickness of the wetting film are published. However, again the variation of film thicknesses agrees with our findings. In particular, one observes a rapid thickening of the films upon approaching the coexistence line; for some values of $`X`$ a small maximum located near the $`\lambda `$-line is also visible (compare Fig. 3). A quantitative comparison is not possible because, for our choice of surface terms in the Hamiltonian, the fixed-point BC $`(O,+)`$ cannot be reached within the VBEG model. In order to be able to extract universal properties - which requires to reach fixed-point BC - it would be necessary to introduce a surface field which couples to the SOP so that the BC (+) can be realized. Also MFT is not sufficient, even in $`d=d^{}`$. A naive correction of $`\vartheta `$ by the logarithmic factor derived within the GL model will not give the proper universal behavior. Instead, renormalization group schemes or Monte Carlo simulations have to be employed. Nonetheless, our MFT results for $`X=X_t`$, if matched at the tricritical point $`y=0`$ and after adjusting the amplitude $`\xi _0`$ so that the positions of the maximum of the scaling function are the same (i.e., $`\xi _0^{th}/a0.065`$) reproduce very well the experimental curve (see Fig. 4), especially near the maximum where we expect interfacial effects to be dominant. This is consistent because the ’soft mode’ phase does not depend on the details of the surface fields. Notice, that the experimental data nominally for $`X=X_t`$ more closely match the theory for $`X=X_t0.01`$. This raises the question as to whether the <sup>3</sup>He concentration in the film is shifted relative to the bulk one.
A.M. benefited from discussions with R. Garcia, M. Krech and S. Kondrat. This work was partially funded by KBN grant No.4 T09A 066 22. |
warning/0506/cond-mat0506282.html | ar5iv | text | # Dressed Feshbach molecules in the BEC-BCS crossover
## Abstract
We present the RPA theory of the BEC-BCS crossover in an atomic Fermi gas near a Feshbach resonance that includes the relevant two-body atomic physics exactly. This allows us to determine the probability $`Z`$ for the dressed molecules in the Bose-Einstein condensate to be in the closed channel of the Feshbach resonance and to compare with the recent experiments of Partridge et al. \[cond-mat/0505353\] with <sup>6</sup>Li. We determine for this extremely broad resonance also the condensate density of the dressed molecules throughout the BEC-BCS crossover.
Introduction. — The superfluid phase in an atomic Fermi gas near a Feshbach resonance realizes a fundamentally new state of matter, which shows a macroscopic coherence between atom pairs and molecules that is controlled by the applied magnetic field. As a result, such a gas offers the exciting possibility to study in great detail the crossover between the Bose-Einstein condensation (BEC) of diatomic molecules and the Bose-Einstein condensation of atomic Cooper pairs, i.e., the Bardeen-Cooper-Schrieffer (BCS) transition stoof1996 ; timmermans2001 ; ohashi2002 ; milstein2002 . In fact, the BEC-BCS crossover region is presently already actively being explored by a number of experimental groups around the world regal2004 ; MIT ; Duke ; Innsbruck ; ENS ; Rice .
In more detail the physics of the BEC-BCS crossover occurring near a Feshbach resonance can be understood as follows: The superfluid phase of the gas is always associated with a Bose-Einstein condensate of pairs, but the wave function of the pairs or dressed molecules is given by the linear superposition
$`𝐫|\chi _{\mathrm{dressed}}=\sqrt{Z}\chi _\mathrm{m}(𝐫)|\mathrm{closed}+\sqrt{1Z}\chi _{\mathrm{aa}}(𝐫)|\mathrm{open}.`$
In the BEC limit the applied magnetic field is taken such that the bare molecular energy level lies far below the threshold of the two-atom continuum and we have $`Z1`$. In that case we are dealing with a Bose-Einstein condensate of tightly-bound diatomic molecules and the spatial part of the pair wave function is equal to the bare molecular wave function $`\chi _\mathrm{m}(𝐫)`$. The spin part of the pair wave function is then equal to $`|\mathrm{closed}`$, i.e., the spin state of the closed channel that causes the Feshbach resonance duine2003 . In the BCS limit the bare molecular energy level lies far above the threshold of the two-atom continuum and can be adiabatically eliminated. We then have that $`Z0`$ and the spatial part of the pair wave function equals the usual BCS wave function for atomic Cooper pairs $`\chi _{\mathrm{aa}}(𝐫)`$. This Cooper-pair wave function depends also on the magnetic field, because the effective attraction between the atoms after the adiabatic elimination of the bare molecular state depends on the energy of that state. The spin state of the Cooper pairs is, however, always equal to the spin state $`|\mathrm{open}`$ of the open channel of the Feshbach problem.
The probability $`Z`$ plays therefore a crucial role in the description of the BEC-BCS crossover since it quantifies the amount of coherence between the atom pairs and the bare molecules in the gas. Unfortunately, however, the various theories falco2004a ; strinati ; levin ; mackie ; ho ; griffin ; burnett ; wetterich ; stringari that are presently being used to understand the outcome of the experiments are unable to accurately determine this quantity. This comes about because, either $`Z`$ is assumed to be zero from the outset, the many-body theory does not incorporate the two-body Feshbach physics exactly, or the theory is able to determine only the total number of bare molecules in the gas and thus requires an additional assumption about the total number of dressed molecules to extract $`Z`$. This situation is particularly unsettling because of the recent <sup>6</sup>Li experiments of Partridge et al. Rice , which have used the photodissociation rate to measure the value of $`Z`$ throughout the crossover regime. In view of the above situation it is pressing to develop an ab initio many-body theory for the calculation of $`Z`$. How that may be achieved is the main topic of this Letter.
BEC-BCS crossover theory. — Introducing creation and annihilation operators for the bare molecules and atoms, the effective grand-canonical hamiltonian of the gas with chemical potential $`\mu `$ becomes drummond ; timmermans1999 ; duine2003
$`H`$ $`={\displaystyle 𝑑𝐱\psi _\mathrm{m}^{}(𝐱)\left(\frac{\mathrm{}^2\mathbf{}^2}{4m}+\delta 2\mu \right)\psi _\mathrm{m}(𝐱)}`$ (1)
$`+{\displaystyle \underset{\sigma =,}{}}{\displaystyle 𝑑x\psi _\sigma ^{}(𝐱)\left(\frac{\mathrm{}^2\mathbf{}^2}{2m}\mu \right)\psi _\sigma (𝐱)}`$
$`+g{\displaystyle 𝑑x\left(\psi _\mathrm{m}^{}(𝐱)\psi _{}(𝐱)\psi _{}(𝐱)+\psi _{}^{}(𝐱)\psi _{}^{}(𝐱)\psi _\mathrm{m}(𝐱)\right)}`$
$`+{\displaystyle \frac{4\pi a_{\mathrm{bg}}\mathrm{}^2}{m}}{\displaystyle 𝑑x\psi _{}^{}(𝐱)\psi _{}^{}(𝐱)\psi _{}(𝐱)\psi _{}(𝐱)},`$
where the two hyperfine states of the atoms are denoted by $`|`$ and $`|`$, and the magnetic-moment difference $`\mathrm{\Delta }\mu _{\mathrm{mag}}`$ between the hyperfine states $`|\mathrm{closed}`$ and $`|\mathrm{open}(||)/\sqrt{2}`$ gives the so-called detuning from resonance $`\delta =\mathrm{\Delta }\mu _{\mathrm{mag}}(BB_0)`$. Note that the atom-molecule coupling constant $`g`$ and the background scattering length $`a_{\mathrm{bg}}`$ depend on the magnetic field $`B`$ in such a manner that the total scattering length $`a=a_{\mathrm{bg}}mg^2/4\pi \mathrm{}^2\delta `$ agrees with the Feshbach resonance of interest. This is especially important for the broad Feshbach resonance near $`834`$ Gauss that is used in all <sup>6</sup>Li experiments up to date houbiers ; thomas .
From now on we consider only the most interesting region close to resonance, where $`a_{\mathrm{bg}}a`$ and the effective interaction between the atoms is dominated by the resonant part $`g^2/\delta `$. This suggests that the last term in the right-hand side of Eq. (1) can be neglected altogether. This is, however, not true because we can neglect this term only after we have included its effect on the atom-molecule coupling falco2005 . Physically, the reason for this subtlety is that the above Hamiltonian is an effective Hamiltonian that is only valid for low energies. However, to properly account for the two-body physics near a Feshbach resonance also high-energy states are needed. The main effect of these high-energy states is to renormalize the atom-molecule coupling to $`g(𝐤)g/(1+ika_{\mathrm{bg}})`$, where $`𝐤`$ is the relative momentum of the atoms involved in the coupling duine2003 ; falco2005 . Only after having performed this substitution are we allowed to neglect the background interaction between the atoms.
Without the background interaction the atomic part of the Hamiltonian is quadratic. Using standard functional methods the fermionic fields can thus be integrated out exactly. This leads to an effective action for the molecules that at sufficiently low temperatures has a minimum at a nonzero value of $`\psi _\mathrm{m}(𝐱)\sqrt{Zn_{\mathrm{mc}}}`$, where we introduced the dressed molecular condensate density $`n_{\mathrm{mc}}`$. Neglecting fluctuations at this point leads to a mean-field theory of the BEC-BCS crossover. As mentioned in the introduction, however, this mean-field theory is unable to calculate the probability $`Z`$ since it only determines the bare molecular condensate density $`|\psi _\mathrm{m}(𝐱)|^2`$. We therefore also consider the quadratic fluctuations around the minimum of the effective molecular action, i.e., we consider the Bogoliubov theory of the bare molecules.
As expected with a spontaneously broken $`U(1)`$ symmetry associated with the presence of a Bose-Einstein condensate, the gaussian fluctuations are determined by normal and anomalous self energies of the bare molecules, which at zero temperature reduce to
$`\mathrm{}\mathrm{\Sigma }_{11}(𝐤,i\omega )`$ $`=`$ $`{\displaystyle \frac{d𝐤^{}}{(2\pi )^3}|g(𝐤^{})|^2\left\{\frac{|u_a(𝐤_+^{})|^2|u_a(𝐤_{}^{})|^2}{i\mathrm{}\omega \mathrm{}\omega _a(𝐤_+^{})\mathrm{}\omega _a(𝐤_{}^{})}\frac{|v_a(𝐤_+^{})|^2|v_a(𝐤_{}^{})|^2}{i\mathrm{}\omega +\mathrm{}\omega _a(𝐤_{}^{}{}_{+}{}^{})+\mathrm{}\omega _a(𝐤_{}^{})}+\frac{1}{2ϵ(𝐤^{})}\right\}},`$
$`\mathrm{}\mathrm{\Sigma }_{12}(𝐤,i\omega )`$ $`=`$ $`2{\displaystyle \frac{d𝐤^{}}{(2\pi )^3}|g(𝐤^{})|^2\left\{u_a(𝐤_+^{})v_a(𝐤_+^{})u_a(𝐤_{}^{})v_a(𝐤_{}^{})\frac{\mathrm{}\omega _a(𝐤_+^{})+\mathrm{}\omega _a(𝐤_{}^{})}{(\mathrm{}\omega _a(𝐤_+^{})+\mathrm{}\omega _a(𝐤_{}^{}{}_{}{}^{}))^2+(\mathrm{}\omega )^2}\right\}}.`$ (2)
Here we have also introduced the BCS dispersion $`\mathrm{}\omega _\mathrm{a}(𝐤)=\sqrt{(ϵ(𝐤)\mu )^2+|g(𝐤)|^2Zn_{\mathrm{mc}}}`$, the bare atomic dispersion $`ϵ(𝐤)=\mathrm{}^2𝐤^2/2m`$, the usual BCS coherence factors $`u_\mathrm{a}(𝐤)`$ and $`v_\mathrm{a}(𝐤)`$, and the notation $`𝐤_\pm ^{}=𝐤/2\pm 𝐤^{}`$. In terms of the above self energies the minimum of the effective action is determined by the exact Hugenholtz-Pines relation $`2\mu =\delta +\mathrm{}\mathrm{\Sigma }_{11}(\mathrm{𝟎},0)\mathrm{}\mathrm{\Sigma }_{12}(\mathrm{𝟎},0)`$, which turns out to be equal to a modified BCS gap equation
$`\delta 2\mu ={\displaystyle \frac{d𝐤}{(2\pi )^3}|g(𝐤)|^2\left(\frac{1}{2\mathrm{}\omega (𝐤)}\frac{1}{2ϵ(𝐤)}\right)}.`$ (3)
Finally, we also need the equation of state, which we for consistency reasons duine2003 obtain by differentiating the thermodynamic potential with respect to the chemical potential. Including the effect of the fluctuations we obtain for the total density of atoms
$$n=\mathrm{Tr}[G_\mathrm{a}]+2Zn_{\mathrm{mc}}\mathrm{Tr}[G_\mathrm{m}]+\frac{1}{2}\mathrm{Tr}\left[G_\mathrm{m}\frac{\mathrm{}\mathrm{\Sigma }}{\mu }\right],$$
(4)
where $`G_\mathrm{a}`$ and $`G_\mathrm{m}`$ are the Nambu ($`2\times 2`$)-matrix Green’s functions of the bare atoms and molecules, respectively. For a given density and magnetic field the last two equations determine only the bare molecular condensate density and the chemical potential. Hence, we need a third equation to determine also $`Z`$.
Dressed molecules. — Before we derive this missing equation, let us first discuss in some more detail the physics behind the maybe somewhat unexpected equation of state in Eq. (4). The first two terms represent the mean-field theory without fluctuations that is most often used in the recent literature levin ; mackie ; ho ; burnett . Because of the absence of fluctuations all the molecules are Bose-Einstein condensed and there is no depletion. The third term precisely corresponds to this depletion. Finally, the fourth term physically describes the dressing of the bare atoms and molecules. This can be made more clear by reformulation the equation of state in terms of dressed atoms and molecules, instead of bare atoms and molecules. Since every dressed molecule contains two atoms, we expect the contribution $`2n_{\mathrm{mc}}`$ from the condensate of dressed molecules. Indeed, in the BEC limit it can be shown explicitly that the atomic density $`\mathrm{Tr}[G_\mathrm{a}]=2𝑑𝐤|v_\mathrm{a}(𝐤)|^2/(2\pi )^3`$ contains exactly the expected contribution $`2(1Z)n_{\mathrm{mc}}`$ of paired atoms in the Bose-Einstein condensate of dressed molecules. The atomic density does, however, not contain the paired atoms associated with the dressed molecules that are not in the Bose-Einstein condensate. This omission is repaired by the fluctuation corrections which contain both the associated changes in the atomic density and twice the total density of dressed molecules that are not Bose-Einstein condensed, i.e., twice the dressed molecular depletion.
We are now in a position to determine $`Z`$. In principle it is equal to the residue of the pole of $`G_{\mathrm{m};11}(\mathrm{𝟎},i\omega )`$ at $`\omega =0`$. To understand the physics of that result better it is useful to consider the spectral function of the bare molecules with zero momentum, i.e., $`\rho _\mathrm{m}(\mathrm{𝟎},\omega )=\mathrm{Im}[G_{\mathrm{m};11}(\mathrm{𝟎},\omega +i0)]/\pi `$. This spectral function is closely related to the density of states of the molecules at zero momentum and thus gives detailed information on the dressed molecular content of the gas. Most importantly for our purposed, the Bose-Einstein condensate of dressed molecules gives rise to a delta-function in the spectral function exactly at zero frequency. The strength of this delta function is precisely $`Z`$, because this is the probability to take a bare molecule out of the Bose-Einstein condensate of dressed molecules. Besides this bound-state contribution, the spectral function contains also a contribution from the continuum of atomic scattering states. In the BEC limit of the crossover the continuum contribution only occurs at positive frequency and starts at a frequency of $`2\mu /\mathrm{}2\mathrm{}/ma^2`$ as shown in Fig. 1a. In the BCS limit the continuum contribution occurs both at positive and negative frequencies, which start at a frequency of about $`\pm 2g\sqrt{Zn_{\mathrm{mc}}}/\mathrm{}`$, respectively, due to the gap that exists for the creation of an atomic quasiparticle-quasihole pair. This is shown in Fig. 1b. The negative frequency part of the spectral function is especially important, because it determines the depletion of the Bose-Einstein condensate of dressed molecules. Physically it represents the dressed molecules that are stabilized by the Fermi sea falco2004b . Making use of the above physical picture, we finally obtain the desired result,
$`Z={\displaystyle \frac{1\mathrm{\Sigma }_{11}^{(1)}}{\left(1\mathrm{\Sigma }_{11}^{(1)}\right)^2+\mathrm{\Sigma }_{12}\left(\mathrm{\Sigma }_{12}^{(2)}\mathrm{\Sigma }_{11}^{(2)}\right)}},`$ (5)
where $`\mathrm{\Sigma }_{ij}^{(n)}(i)^n^n\mathrm{\Sigma }_{ij}(\mathrm{𝟎},0)/\omega ^n`$.
At nonzero momenta the spectral function is similar but now contains two delta functions at the frequencies $`\pm \omega _\mathrm{m}(𝐤)`$, which have the strength $`Z|u_\mathrm{m}(𝐤)|^2`$ and $`Z|v_\mathrm{m}(𝐤)|^2`$, respectively, with $`|u_\mathrm{m}(𝐤)|^2|v_\mathrm{m}(𝐤)|^2=1`$. This shows explicitly how at long wavelengths our theory leads to a Bogoliubov-like theory for dressed molecules with a wave function renormalization factor $`Z`$. Moreover, in agreement with the Goldstone theorem, the quasiparticle dispersion $`\omega _\mathrm{m}(𝐤)`$ always turns out to be linear at long wavelengths. In the following we therefore determine also the associated speed of (second) sound throughout the BEC-BCS crossover region.
Results and discussion. — In Fig. 2 we show our results for $`Z`$ and compare with the experimental data of Partridge et al. Rice . In general the agreement is satisfactory. This is particularly true at relatively low magnetic fields where $`Z`$ is determined by two-body physics, which is exactly incorporated into our theory. At higher magnetic fields the theoretical values of $`Z`$ are somewhat higher than the experimental ones. We believe that the reasons for this are twofold. First, the experiment is performed in an optical trap. As a result the experimental data involves an appropriate average over the density profile of the gas, which lowers the observed value of $`Z`$. Second, the photodissociation laser used in the experiment has a width which is much larger than the Fermi energy of the gas. The laser, therefore, has not sufficient resolution to probe only the Bose-Einstein condensate of dressed molecules, and probes also the dressed molecules which are not Bose-Einstein condensed. This second effect should be especially important at high magnetic fields. To disentangle these different effects, however, goes beyond the scope of the present Letter and is left for further investigation.
We also show in Fig. 2 the Bose-Einstein condensate fraction of dressed molecules $`2n_{\mathrm{mc}}/n`$ throughout the BEC-BCS crossover. In qualitative agreement with the poor man’s approach of Ref. falco2004a , the latter fraction is always substantial below the Feshbach resonance and becomes negligible only sufficiently far above the Feshbach resonance when $`k_\mathrm{F}|a|<1`$. This is an important observation, because in our theoretical description of the experiment of Partridge et al. the molecular probe couples directly to the dressed molecules, which act as distinct entities in the gas since the atom-molecule coupling is much larger than the coupling of the probe laser to the bare molecules. In this manner it is most easy to understand the experimental observation that there is initially an exponential (one-body) decay of a large part of the total atomic density with a rate that is much smaller than the bare molecular photodissociation rate.
For completeness we give in Fig. 2 also the Bose-Einstein condensate fraction of bare molecules $`2Zn_{\mathrm{mc}}/n`$ and the fluctuation corrections to the total atomic density. As expected the fluctuation corrections are very important in the crossover region and become small far away from the Feshbach resonance, where mean-field theory applies. From the fluctuations we also extract the speed of sound of the gas, which is shown in Fig. 3. Note that in the BCS limit the Anderson-Bogoliubov mode is recovered. In combination with the presence of the sharp peaks in the spectral function in Fig. 1b, this shows that also the decoupling of the amplitude and phase fluctuations of the Bose-Einstein condensate of dressed molecules that occurs in the BCS limit is correctly incorporated. We therefore conclude that the RPA-like atom-molecule theory developed here gives an excellent account of the subtle interplay between two-body and many-body physics taking place at the crossover near a Feshbach resonance.
We are very grateful for many helpful remarks and stimulating discussions with Randy Hulet. This work is supported by the Stichting voor Fundamenteel Onderzoek der Materie (FOM) and the Nederlandse Organisatie voor Wetenschaplijk Onderzoek (NWO). |
warning/0506/quant-ph0506193.html | ar5iv | text | # Security of Continuous-variable quantum cryptography using coherent states: Decline of postselection advantage
## Abstract
We investigate the security of continuous-variable (CV) quantum key distribution (QKD) using coherent states in the presence of quadrature excess noise. We consider an eavesdropping attack which uses a linear amplifier and beam splitter. This attack makes a link between beam-splitting attack and intercept-resend attack (classical teleportation attack). We also show how postselection loses its efficiency in a realistic channel.
Quantum key distribution (QKD) is a technique that allows two parties, Alice (the sender) and Bob (the receiver), to share a key which is kept secret from an eavesdropper (Eve) who has advanced computational and technological power rmp74 . To achieve a signal transmission between distant parties, controlling optical quantum states is essential. Several QKD schemes based on continuous-variable (CV) which uses the quadrature amplitude of light field have been proposed squeezed ; epr ; osci ; con ; alcon ; coherent ; postsel ; coherentR ; hirano ; namiki1 . Although usage of squeezed states or EPR states are fundamentally interesting, coherent-state protocols have practical advantage of easy state preparation. CV QKD using coherent states over a 1$`k`$m-optical-fiber path has been experimentally demonstrated at 1.55$`\mu `$m-communication wavelength hirano .
The performance of QKD is limited by the presence of the transmission loss. A simple treatment of the loss effect is beam-splitting attack (BSA) where Eve replaces the transmission path with the lossless one and a beam splitter (BS). Then she obtains the signals corresponding to the loss without making any disturbance to the signal. At first sight, over the existence of 50% loss (3dB loss), it seems to be impossible to distill the secret key using coherent-state signal because Eve can get stronger signal than Bob coherent . However, since the knowledge about the signal depends on the measurement result, coherent-state protocol can provide a secure key by conditional use of measurement results (postselection, PS) even in the presence of higher loss postsel ; hirano ; namiki1 . PS plays an important role in many implementations of quantum information processing tasks as well as QKD.
In realistic condition, besides the loss, excess Gaussian noise is imposed on the quadrature distribution hirano . Since any excess noise tapers off when the state falls into vacuum at high loss, the excess noise added by Eve near Alice’s side will disappear at Bob’s side for a long transmission distance. Then, for a sufficiently long distance, eavesdropping cannot be detectable. From this observation, it is shown that CV-QKD protocols using coherent states cannot work for arbitrary transmission distance in the presence of excess noise namiki2 . This limitation is given by an intercept-resend attack called classical teleportation attack (CTA).
The question is what kind of attack links between “direct” CTA and “indirect” BSA, and how PS works in the presence of excess noise. In this Letter, we provide an intermediate attack between BSA and CTA, and show how PS loses its advantage in the presence of noise.
We consider the realistic channel which transforms coherent state into a Gaussian mixture of coherent states as
$$|\alpha \widehat{\rho }(\alpha ,\eta ,\delta )\frac{2}{\pi \delta }e^{\frac{2|\beta |^2}{\delta }}|\sqrt{\eta }\alpha +\beta \sqrt{\eta }\alpha +\beta |d^2\beta ,$$
(1)
where $`\eta `$ is the line transmission and $`\delta `$ is quadrature excess noise. Coherent state is eigen state of $`\widehat{a}`$: $`\widehat{a}|\alpha =\alpha |\alpha `$. We define quadrature amplitude $`\widehat{x}_1`$, $`\widehat{x}_2`$ by the relation $`\widehat{a}=\widehat{x}_1+i\widehat{x}_2`$. $`\widehat{a}`$ is the annihilation operator of signal pulse mode. The quadrature variance of coherent state is given by $`(\mathrm{\Delta }x)^2=\frac{1}{4}`$. As a frame work, we assume that all noise is caused by Eve in the quantum channel and Bob has an ideal detector. Then, for coherent-state input to the channel (1), Bob observes Gaussian quadrature distribution hirano and the observed quadrature variance $`(\mathrm{\Delta }x_{\text{obs}})^2`$ is related to the excess noise as
$$(\mathrm{\Delta }x_{\text{obs}})^2=(1+\delta )(\mathrm{\Delta }x)^2.$$
(2)
Bob’s mean values of quadratures can be related to the transmission and coherent-state amplitude as
$$\widehat{x}_1+i\widehat{x}_2=\sqrt{\eta }\alpha .$$
(3)
In terms of $`\delta `$ and $`\eta `$, the limitation given by CTA is $`\delta <2\eta `$. We refer to it as classical teleportation limit (CTL) namiki2 .
Some of eavesdropping attacks which cause the state change (1) can be constructed by combining BSs and phase-insensitive amplifiers (AMP). Simplest case is that Eve uses only one BS and one AMP. If Eve uses the Amplifiers after the BS, Eve’s and Bob’s quadratures are modulated independently. Thus the results of Bob’s quadrature measurement and Eve’s state are not correlated, and the effectiveness of PS is inherently different from that of BSA namiki2 . Here we consider the other case which we call amplification-beam-splitting attack (AMPBSA) where Eve inserts BS after AMP (see FIG. 1).
Let us assume that Eve operates a phase-insensitive amplifier amp with amplifier gain $`g1`$ and inserts a BS with reflectivity $`1\kappa `$ (see Fig. 1). Then $`g`$ and $`\kappa `$ are related to Bob’s mean value and variance of quadratures as
$`\widehat{x}_1+i\widehat{x}_2`$ $`=`$ $`\sqrt{g\kappa }\alpha ,`$ (4)
$`(\mathrm{\Delta }x_{\text{obs}})^2`$ $`=`$ $`\{2(g1)\kappa +1\}(\mathrm{\Delta }x)^2.`$ (5)
Using Eqs. (2), (3), (4) and (5), we obtain
$`g`$ $`=`$ $`{\displaystyle \frac{\eta }{\eta \delta /2}},`$ (6)
$`\kappa `$ $`=`$ $`\eta \delta /2.`$ (7)
From Eq. (6), we can see that BSA ($`\delta =0`$) is the unit gain case: $`g=1`$, and CTL ($`\delta 2\eta `$) is relevant to the infinite-gain limit: $`g\mathrm{}`$.
Provided Alice sent $`|\alpha `$, Eve’s operation makes the joint state of Bob and Eve:
$`|\alpha _B|0_E`$ $``$ $`\widehat{\rho }_{BE}(\alpha ,\eta ,\delta ){\displaystyle \frac{2}{\pi \delta }}{\displaystyle e^{\frac{2}{\delta }|\beta |^2}|\sqrt{\eta }\alpha +\beta _B\sqrt{\eta }\alpha +\beta ||\xi (\sqrt{\eta }\alpha +\beta )_E\xi (\sqrt{\eta }\alpha +\beta )|d^2\beta },`$ (8)
where we defined
$`\xi \sqrt{{\displaystyle \frac{1\eta +\delta /2}{\eta \delta /2}}}.`$ (9)
The subscripts B and E stand for Bob’s and Eve’s system, respectively. We can easily see that $`\text{Tr}_E\left(\widehat{\rho }_{BE}\right)=\widehat{\rho }(\alpha ,\eta ,\delta )`$ where $`\text{Tr}_E`$ is partial trace of Eve’s system.
The form of $`\widehat{\rho }_{BE}`$ shows that Eve receives $`|\xi (\sqrt{\eta }\alpha +\beta )`$ when Bob receives $`|\sqrt{\eta }\alpha +\beta `$ i.e., the coherent state Eve receives is different from that of Bob’s only by the amplitude factor $`\xi `$ for each transmission. This simple picture is useful to explain (i)CTL and (ii)3dB loss limit coherent as follows:
(i) If $`\xi `$ becomes infinity ($`\delta 2\eta `$ or $`g\mathrm{}`$), Eve can read out the coherent-state amplitude $`\sqrt{\eta }\alpha +\beta `$ with arbitrary resolution by performing simultaneous measurement of quadratures. In other words, she can determine the state Bob receives $`|\sqrt{\eta }\alpha +\beta `$ or she can produce infinite number of copies of this state. This condition is equivalent to the case of intercept-resend attack, and it demonstrates CTL.
(ii) A sufficient condition for secure key distribution against individual attack is
$$I_{AB}I_{AE},$$
(10)
where $`I_{AB(AE)}`$ is Mutual information between Alice and Bob (Eve) pa ; gpa ; postsel . For Gaussian continuous key distribution protocols alcon ; coherent , the Mutual information is directly related to the signal-to-noise ratio (SNR); the higher is the SNR, the higher is the amount of mutual information between the parties. The case $`\xi 1`$ means Bob’s SNR is higher than that of Eve because in Eq. (8) Bob’s coherent-state amplitude is larger than that of Eve’s. Thus, the security condition given in (10) is written as
$$\xi 1.$$
(11)
If $`\delta =0`$, we obtain a reduced condition $`\eta 1/2`$ which gives the 3dB loss limit. This analysis excludes reverse-reconciliation protocol of coherentR .
Now we go into the security of PS protocol against AMPBSA. In the postselection protocol postsel ; hirano ; namiki1 , a measurement result $`x`$ higher than a given threshold $`x_00`$ is selected to distil the secret key. By setting $`x_0`$ higher, bit-error rate (BER) of Bob can be arbitrary small provided the quadrature distribution is Gaussian. In contrast to this, if Eve’s signal is independent of $`x`$, Eve’s BER remains constant under PS. So PS can make information advantage for Bob compare to Eve. This is not the case for AMPBSA. In Eq. (8), the modulation $`\beta `$ is added to both of Eve’s and Bob’s systems collectively. Through $`\beta `$, Eve’s state depends on $`x`$ and then Eve’s BER depends on PS. So the correlation between Bob’s and Eve’s state given by $`\beta `$ weakens the PS advantage.
Let us consider the case that Alice sends binary-phase-shifted coherent states $`|\pm \alpha `$ with $`\alpha >0`$ and Bob performs quadrature measurement on the correct basis $`\widehat{x}_1`$ namiki1 ; namiki2 . To describe the postselection events, we use Eve’s density operator conditioned on Bob’s measurement result $`x`$. If Alice sent $`|\alpha `$ and Bob observed $`x`$, Eve’s density operator (conditioned on $`x`$) is given by
$`\widehat{\rho }_E(\alpha |x)`$ $`=`$ $`{\displaystyle \frac{\text{Tr}_B\left(\widehat{\rho }_{BE}(\alpha ,\eta ,\delta )|x_Bx|\right)}{P_B(x|\alpha )}},`$ (12)
where
$`P_B(x|\alpha )`$ $`=`$ $`\text{Tr}\left(\widehat{\rho }_{BE}(\alpha ,\eta ,\delta )|x_Bx|\right)`$ (13)
$`=`$ $`x|\widehat{\rho }(\alpha ,\eta ,\delta )|x`$
is the probability that Bob gets quadrature value $`x`$.
Since Eve does not know the sign of $`x`$, we can estimate Eve’s information from the density operator conditioned on the absolute value of $`x`$:
$`\widehat{\rho }_E(\alpha ||x|)`$ $``$ $`P(\alpha |x)\widehat{\rho }_E(\alpha |x)+P(\alpha |x)\widehat{\rho }_E(\alpha |x),`$ (14)
where we define the probability that Alice’s choice is $`|\alpha `$ when Bob gets $`x`$:
$`P(\alpha |x)`$ $``$ $`{\displaystyle \frac{P_B(x|\alpha )}{P_B(x|\alpha )+P_B(x|\alpha )}}.`$ (15)
Therefore, if Bob’s measurement result is $`\pm x`$, Eve gets either of the two mixed states $`\widehat{\rho }_E(\pm \alpha ||x|)`$ corresponding to Alice’s choice $`|\pm \alpha `$, respectively. The next problem is how Eve differentiates the given two signal $`\widehat{\rho }_E(\pm \alpha ||x|)`$. In general Eve may choose her measurement knowing the value $`|x|`$. Here we restrict our analysis for the case that Eve performs quadrature measurement and determines the bit value according to the sign of her measurement result $`x_E`$ as Bob does. In this case she does not use the information of $`|x|`$.
For the PS protocol, Mutual information is written as
$$I_{AB(E)}=\frac{1}{2}\underset{|x|x_0}{}P_B(x)i(q_{B(E)}(x)),$$
(16)
where
$$P_B(x)=\frac{1}{2}(P_B(x|\alpha )+P_B(x|\alpha ))$$
(17)
is the probability that Bob’s measurement result is $`x`$,
$$i(q)=1+q\mathrm{log}_2q+(1q)\mathrm{log}_2(1q)$$
(18)
is Mutual information of binary symmetric channel, and $`q_{B(E)}(x)`$ is BER of Bob (Eve) conditioned on $`|x|`$. Since $`i(q)`$ is a decreasing function of $`q`$ $`(0q\frac{1}{2})`$, $`q_B(x)>q_E(x)`$ for any $`x`$ implying that any conditional use of measurement result does not satisfy inequality (10).
Eve’s BER conditioned on $`|x|`$ is the probability that the signal is $`\widehat{\rho }_E(\alpha ||x|)`$ when the sign of Eve’s measurement result $`x_E`$ is positive:
$`q_E(x,\eta ,\delta )`$ $`=`$ $`{\displaystyle \frac{_0^{\mathrm{}}x_E|\widehat{\rho }_E(\alpha ||x|)|x_Edx_E}{_0^{\mathrm{}}(x_E|\widehat{\rho }_E(\alpha ||x|)|x_E+x_E|\widehat{\rho }_E(\alpha ||x|)|x_E)dx_E}}`$ (19)
$`=`$ $`{\displaystyle \frac{1}{2}}P(\alpha |x)\text{ erfc}\left(\sqrt{2}\lambda (\delta x+\sqrt{\eta }\alpha )\right)+{\displaystyle \frac{1}{2}}P(\alpha |x)\text{ erfc}\left(\sqrt{2}\lambda (\delta x\sqrt{\eta }\alpha )\right),`$ (20)
where we use the definition of $`\widehat{\rho }_{BE}`$ (8), Eqs. (12 \- 14), and quadrature distribution of coherent state $`|x|\alpha |^2=\sqrt{2/\pi }e^{2(x\alpha )^2}`$, and we define $`\text{erfc}(x)=2/\sqrt{\pi }_s^{\mathrm{}}e^{t^2}𝑑t`$ and
$`\lambda `$ $``$ $`\sqrt{{\displaystyle \frac{(1\eta )+\delta /2}{(\eta +\delta /2)(1+\delta )}}}.`$ (21)
In what follows we set $`x>0`$ for simplicity.
For sufficiently large $`x`$, $`q_E`$ is bounded above as
$`q_E(x,\eta ,\delta )`$ $`=`$ $`{\displaystyle \frac{1}{2}}P(\alpha |x)\text{ erfc}\left(\sqrt{2}\lambda (\delta x+\sqrt{\eta }\alpha )\right)+{\displaystyle \frac{1}{2}}P(\alpha |x)\text{ erfc}\left(\sqrt{2}\lambda (\delta x\sqrt{\eta }\alpha )\right)`$ (22)
$``$ $`{\displaystyle \frac{1}{2}}\left\{\text{ erfc}\left(\sqrt{2}\lambda (\delta x+\sqrt{\eta }\alpha )\right)+\text{ erfc}\left(\sqrt{2}\lambda (\delta x\sqrt{\eta }\alpha )\right)\right\}`$
$``$ $`\text{erfc}\left(\sqrt{2}\lambda (\delta x\sqrt{\eta }\alpha )\right)`$
$``$ $`{\displaystyle \frac{1}{\sqrt{\pi }}}e^{2\lambda ^2(\delta x\sqrt{\eta }\alpha )^2}={\displaystyle \frac{1}{\sqrt{\pi }}}e^{2\lambda ^2(\delta ^2x^22\sqrt{\eta }\alpha \delta x+\eta \alpha ^2)}`$
$`<`$ $`{\displaystyle \frac{e^{2\lambda ^2\eta \alpha ^2}}{\sqrt{\pi }}}e^{\lambda ^2\delta ^2x^2}.`$
The first inequality comes from the fact $`P(\pm \alpha |x)1`$. The second inequality comes from the fact $`\text{erfc}(s)`$ is a decreasing function of $`s`$ with $`\lambda (\delta x\sqrt{\eta }\alpha )\lambda (\delta x+\sqrt{\eta }\alpha )`$. Then, we use an inequality $`\sqrt{\pi }\text{erfc}(s)/2=_s^{\mathrm{}}e^{t^2}𝑑t_s^{\mathrm{}}te^{t^2}𝑑t=e^{s^2}/2`$ for $`s1`$ assuming $`x`$ is large enough so that $`\sqrt{2}\lambda (\delta x\sqrt{\eta }\alpha )>1`$, which gives the third inequality. Further assuming $`x>4\sqrt{\eta }\alpha /\delta `$, we obtain the final expression.
Bob’s BER conditioned on $`|x|`$ is the conditional probability that the sent state is $`|\alpha `$ when his measurement result $`x`$ is positive:
$`q_B(x,\eta ,\delta )`$ $``$ $`P(\alpha |x)`$ (23)
$`=`$ $`{\displaystyle \frac{1}{1+\mathrm{exp}\left[\frac{8\sqrt{\eta }\alpha x}{1+\delta }\right]}}.`$
Since $`q_B\mathrm{exp}\left(\frac{8\sqrt{\eta }\alpha }{1+\delta }x\right)`$ and $`q_E<e^{\lambda ^2\delta ^2x^2}`$ from expression (22), for any given $`\delta >0`$, there exist sufficiently large $`x`$ where $`q_E<q_B`$ holds. In such a condition, simple PS setting higher threshold is no more advantageous. It will be more efficient to discard higher quadrature value.
Figure 2 shows, the parameter region of $`\eta `$ and $`\delta `$ which satisfies $`q_Eq_B`$ for any choice of $`x`$ for several mean photon number $`n\alpha ^2`$. This means, above the line, any kind of PS cannot achieve $`I_{AB}I_{AE}`$. We can see that larger $`n`$ seems to be more tolearant to noise. This is because, for a given $`\delta `$, if $`n`$ is larger the effect of collective noise is relatively smaller. It should be noted that choice of large $`n`$ results much smaller secure key gain for higher loss owing to BSA namiki1 . The estimation of secure key gain with some optimization of both Eve’s measurement and PS strategy is left for full paper.
In this analysis we assume an ancillary system of amplifier is just traced out (See Fig. 1). It is likely that the ancilla provides some useful information for Eve. In this sense AMPBSA may be weak attack. The condition where Eve cannot access the ancilla can be realized if Alice sends thermal coherent states $`\widehat{\rho }(\alpha ,1,\eta \delta )`$ instead of coherent states and Eve performs just BSA. This case the density operator of joint system is described by Eq. (8) with the replacement $`\xi \sqrt{(1\eta )/\eta }`$.
In conclusion, we have investigated the security of CV QKD using coherent states against amplification-beam-splitting attack. This attack makes a link between “direct” classical-teleportation attack and “indirect” beam-splitting attack. It has been shown that the postselection protocol setting higher threshold need not ensure the security in the presence of excess Gaussian noise.
We thank M. Koashi for helpful discussions. |
warning/0506/math0506437.html | ar5iv | text | # Nonlinear Connections on Gerbes, Clifford–Finsler Modules, and the Index Theorems
## 1 Introduction
In this paper, we elaborate an approach to the Atiyah–Singer index formulas for manifolds and bundle spaces provided with nonlinear connection (in brief, N–connection) structure. We follow the methods developed for bundle gerbes and bundle gerbe modules and related results from . It should be noted that bundles and gerbes and their higher generalizations ($`n`$–gerbes) can be described both in two equivalent forms: in local geometry, with local functions and forms, and in non–local geometry, by using holonomies and parallel transports, , see review .
The nonholonomic structure of a so–called N–anholonomic space (see, for instance, and ) is stated by a non–integrable, i.e. nonholonomic (equivalently, anholonomic), distribution defining a N–connection structure<sup>1</sup><sup>1</sup>1the rigorous definitions and notations are given below; for our purposes, it will be enough to consider nonholonomic spaces defined by N–connections structures (in brief, called N–anholonomic manifolds). Nonholonomic geometric configurations are naturally derived in modern gravity and string theories by using generic off–diagonal metrics, generalized connections and nonholonomic frame structures . The approach can be elaborated in general form by unifying the concepts of Riemann–Cartan and Finsler–Lagrange spaces and their generalizations on gerbes. It is also related to modelling gravitational field interactions and the Lagrange and/or Hamilton mechanics and further developments to quantum deformations , noncommutative geometry and gravity with N–anholonomic structures , supermanifolds provided with N–connection structure , as well to nonholonomic Lie and Clifford algebroids and their applications in constructing new classes of exact solutions .
The general nonholonomic manifolds fail to be spin and there are substantial difficulties in definition of curvature which can be revised and solved in the theory of nonholonomic gerbes. Some of such constructions are relevant to anomalies in quantum field theory when the obstruction to existence of spin structure is regarded as an anomaly in the global definition of spinor fields. The typical solution of this problem is to introduce some additional fields which also have an anomaly in their global definition but choose a such configuration when both anomalies cancel each another.
The failure in existence of usual spin structure is differently treated for the nonholonomic manifolds. At least for the N–anholonomic spaces, it is possible to define the curvature, which is not a trivial construction for general nonholonomic manifolds, and the so–called N–adapted (nonholonomic) Clifford structures with nontrivial N–connection. This problem was firstly solved for the Finsler–Lagrange spaces and their higher order generalizations but it can be also generalized to noncommutative geometry and gravity models with nontrivial nonholonomic structures and Lie–Clifford algebroid symmetries, see reviews and recent results in Refs. .
This work is devoted to the index theorems for nonholonomic gerbes and bundle gerbe modules adapted to the N–connection structure. The key idea is to consider the so–called N–anholonomic spin gerbe defined for any N–anholonomic manifold (such a nonholonomic gerbe is a usual spin gerbe for a vanishing N–connection structures and becomes trivial if the basic manifold is spin). We shall construct ”twisted” Dirac d–operators <sup>2</sup><sup>2</sup>2In brief, we shall write ”d–operators and d–objects” for operators and objects distinguished by a N–connection structure, see next section. and investigate their properties for Clifford N–anholonomic modules. Then, we shall define the Chern character of such of such modules and show that the usual index formula holds for such a definition but being related to the N–anholonomic structure. The final aim, the proof of index theorems for various types of N–anholonomic spaces, follows from matching up the geometric formalism of Clifford modules and nonhlonomic frames with associated N–connections.
The structure of the paper is as follows:
Section 2 contains an introduction to the geometry of N–anholonomic manifolds. There are given two equivalent definitions of N–connections, considered basic geometric objects characterizing them and defined and computed in abstract form the torsions and curvatures of N–anholonomic manifolds.
Section 3 is devoted to a study of two explicit examples of N–anholonomic manifolds: the Lagrange–Finsler spaces and Riemann–Cartan manifolds provided with N–connection structures. There are proved two main results: Result 3.1: any regular Lagrange mechanics theory, or Finsler geometry, can be canonically modelled as a N–anholonomic Riemann–Cartan manifold with the basic geometric structures (the N–connection, metric and linear connection) being defined by the fundamental Lagrange, or Finsler, function; Result 3.2: There are N–anholonomic Einstein–Cartan (in particular cases, Einstein) spaces parametrized by nontrivial N–connection structure, nonholonomic frames and, in general, non–Riemann connections defined as generic off–diagonal solutions in modern gravity.
In section 4, there are considered the lifting of N–anholonomic bundle gerbes and definition of connection and curvatures on such spaces. We define the (twisted–anholonomic) Chern characters for bundle gerbes and modules induced by N–connections and distinguished metric and linear connection strucutres. We conclude with two important results/ applications of the theory of nonholonomic gerbes: Result 4.1: Any regular Lagrange (Finsler) configuration is topologically characterized by its Chern character computed by using canonical connections defined by the Lagrangian (fundamental Finsler function). Result 4.2: the geometric constructions for a N–anholonomic Riemann–Cartan manifold (including exact solutions in gravity) can be globalized to N–anholonomic gerbe configurations and characterized by the corresponding Chern character.
Section 5 presents the main result of this paper: The twisted index formula for N–anholonomic Dirac operators and related gerbe constructions are stated by Theorem 5.2). We introduce Clifford d–algebras on N–anholonomic bundles and define twisted nonholonomic Dirac operators on Clifford gerbes. We conclude that there are certain fundamental topological characteristics derived from a regular fundamental Lagrange, or Finsler, function and that such indices classify new classes of exact solutions in gravity globalized on gravitational gerbe configurations.
The Appendix contains a set of component formulas for N–connections, metric and linear connection structures and related torsions and curvatures on N–anholonomic manifolds. They may considered for some local proofs of results in the main part of the paper, as well for some applications in modern physics.
## 2 N–Anholonomic Manifolds
We formulate a coordinate free introduction into the geometry of nonholonomic manifolds. The reader may consult details in Refs. . Here we note that there is a comprehensive study of nonholonomic and (for integrable structure) of fibred structures in Ref. following the so–called Schouten – Van Kampen and Vrǎnceanu connections . Different directions in the geometry of nonholonomic manifolds were developed for different geometric structures , in the geometry of Finsler and Lagrange spaces with applications to mechanics and modern geometry . Even from formal point of view all geometric structures on nonholonomic bundle spaces were rigorously investigated by the R. Miron’s school in Romania, various purposes and applications in modern physics requested a different class of nonholonomic manifolds with supersymmetric, noncommutative, Lie algebroid, gerbe etc generalizations . In our approaches, we use such linear and nonlinear connection structure which can be derived naturally as exact solutions in modern gravity theories and from certain Lagrangians/ Hamiltonians in the case of geometric mechanics. Some important component/coordinate formulas are given in the Appendix.
### 2.1 Nonlinear connection structures
Let $`𝐕`$ be a smooth manifold of dimension $`(n+m)`$ with a local fibred structure. Two important particular cases are those of a vector bundle, when we shall write $`𝐕=𝐄`$ (with $`𝐄`$ being the total space of a vector bundle $`\pi :`$ $`𝐄M`$ with the base space $`M)`$ and of a tangent bundle when we shall consider $`𝐕=\mathrm{𝐓𝐌}.`$ The differential of a map $`\pi :𝐕M`$ defined by fiber preserving morphisms of the tangent bundles $`T𝐕`$ and $`TM`$ is denoted by $`\pi ^{}:T𝐕TM.`$ The kernel of $`\pi ^{}`$ defines the vertical subspace $`v𝐕`$ with a related inclusion mapping $`i:v𝐕T𝐕.`$
###### Definition 2.1
A nonlinear connection (N–connection) $`𝐍`$ on a manifold $`𝐕`$ is defined by the splitting on the left of an exact sequence
$$0v𝐕\stackrel{𝑖}{}T𝐕T𝐕/v𝐕0,$$
(1)
i. e. by a morphism of submanifolds $`𝐍:T𝐕v𝐕`$ such that $`𝐍𝐢`$ is the unity in $`v𝐕.`$
The exact sequence (1) states a nonintegrable (nonholonomic, equivalently, anholonomic) distribution on $`𝐕,`$ i.e. this manifold is nonholonomic. We can say that a N–connection is defined by a global splitting into conventional horizontal (h) subspace, $`\left(h𝐕\right),`$ and vertical (v) subspace, $`\left(v𝐕\right),`$ corresponding to the Whitney sum
$$T𝐕=h𝐕_Nv𝐕$$
(2)
where $`h𝐕`$ is isomorphic to $`M.`$ We put the label $`N`$ to the symbol $``$ in order to emphasize that such a splitting is associated to a N–connection structure. In this paper, we shall omit local coordinate considerations.
For convenience, in Appendix, we give some important local formulas (see, for instance, the local representation for a N–connection (A.1)) for the basic geometric objects and formulas on spaces provided with N–connection structure. Here, we note that the concept of N–connection came from E. Cartan’s works on Finsler geometry (see a detailed historical study in Refs. and alternative approaches developed by using the Ehressmann connection ). Any manifold admitting an exact sequence of type (1) admits a N–connection structure. If $`𝐕=𝐄,`$ a N–connection exists for any vector bundle $`𝐄`$ over a paracompact manifold $`M,`$ see proof in Ref. .
The geometric objects on spaces provided with N–connection structure are denoted by ”bolfaced” symbols. Such objects may be defined in ”N–adapted” form by considering h– and v–decompositions (2). Following conventions from , one call such objects to be d–objects (i. e. they are distinguished by the N–connection; one considers d–vectors, d–forms, d–tensors, d–spinors, d–connections, ….). For instance, a d–vector is an element $`𝐗`$ of the module of the vector fields $`\chi (𝐕)`$ on $`𝐕,`$ which in N–adapted form may be written
$$𝐗=h𝐗+v𝐗\text{ or }𝐗=X_N^{}X,$$
where $`h𝐗`$ (equivalently, $`X`$) is the h–component and $`v𝐗`$ (equivalently, $`{}_{}{}^{}X)`$ is the v–component of $`𝐗.`$
A N–connection is characterized by its N–connection curvature (the Nijenhuis tensor)
$$\mathrm{\Omega }(𝐗,𝐘)[^{}X,^{}Y]+^{}[𝐗,𝐘]^{}[^{}X,𝐘]^{}[𝐗,^{}Y]$$
(3)
for any $`𝐗,𝐘\chi (𝐕),`$ where $`[𝐗,𝐘]\mathrm{𝐗𝐘}`$ $`\mathrm{𝐘𝐗}`$ and $`{}_{}{}^{}[,]`$ is the v–projection of $`[,],`$ see also the coordinate formula (A.2) in Appendix. This d–object $`\mathrm{\Omega }`$ was introduced in Ref. in order to define the curvature of a nonlinear connection in the tangent bundle over a smooth manifold. But this can be extended for any nonholonomic manifold, nonholnomic Clifford structure and any noncommutative / supersymmetric versions of bundle spaces provided with N–connection structure, i. e. with nonintegrable distributions of type (2), see .
###### Proposition 2.1
A N–connection structure on $`𝐕`$ defines a nonholonomic N–adapted frame (vielbein) structure $`𝐞=(e,^{}e)`$ and its dual $`\stackrel{~}{𝐞}=\left(\stackrel{~}{e},^{}\stackrel{~}{e}\right)`$ with $`e`$ and $`{}_{}{}^{}\stackrel{~}{e}`$ linearly depending on N–connection coefficients.
Proof. It follows from explicit local constructions, see formulas (A.4), (A.3) and (A.5) in Appendix.$`\mathrm{}`$
###### Definition 2.2
A manifold $`𝐕`$ is called N–anholonomic if it is defined a local (in general, nonintegrable) distribution (2) on its tangent space $`T𝐕,`$ i.e. $`𝐕`$ is N–anholonomic if it is enabled with a N–connection structure (1).
All spinor and gerbe constructions in this paper will be performed for N–anholonomic manifolds.
### 2.2 Curvatures and torsions of N–anholonomic manifolds
One can be defined N–adapted linear connection and metric structures on $`𝐕:`$
###### Definition 2.3
A distinguished connection (d–connection) $`𝐃`$ on a N–anholonomic manifold $`𝐕`$ is a linear connection conserving under parallelism the Whitney sum (2). For any $`𝐗\chi (𝐕),`$ one have a decomposition into h– and v–covariant derivatives,
$$𝐃_𝐗𝐗𝐃=X𝐃+^{}X𝐃=D_X+^{}D_X.$$
(4)
The symbol ”$`\mathrm{"}`$ in (4) denotes the interior product. We shall write conventionally that $`𝐃=(D,^{}D).`$
For any d–connection $`𝐃`$ on a N–anholonomic manifold $`𝐕,`$ it is possible to define the curvature and torsion tensor in usual form but adapted to the Whitney sum (2):
###### Definition 2.4
The torsion
$$𝐓(𝐗,𝐘)𝐃_𝐗𝐘𝐃_𝐘𝐗[𝐗,𝐘]$$
(5)
of a d–connection $`𝐃=(D,^{}D),`$ for any $`𝐗,𝐘\chi (𝐕),`$ has a N–adapted decomposition
$$𝐓(𝐗,𝐘)=𝐓(X,Y)+𝐓(X,^{}Y)+𝐓(^{}X,Y)+𝐓(^{}X,^{}Y).$$
(6)
By further h- and v–projections of (6), denoting $`h𝐓T`$ and $`v𝐓^{}T,`$ taking in the account that $`h[^{}X,^{}Y]=0,`$ one proves
###### Theorem 2.1
The torsion of a d–connection $`𝐃=(D,^{}D)`$ is defined by five nontrivial d–torsion fields adapted to the h– and v–splitting by the N–connection structure
$`T(X,Y)`$ $``$ $`D_XYD_YXh[X,Y],`$
$`^{}T(X,Y)`$ $``$ $`^{}[Y,X],`$
$`T(X,^{}Y)`$ $``$ $`^{}D_YXh[X,^{}Y],`$
$`^{}T(X,^{}Y)`$ $``$ $`^{}D_XY^{}[X,^{}Y],`$
$`^{}T(^{}X,^{}Y)`$ $``$ $`^{}D_X^{}Y^{}D_Y^{}X^{}[^{}X,^{}Y].`$
The d–torsions $`T(X,Y),^{}T(^{}X,^{}Y)`$ are called respectively the $`h(hh)`$–torsion, $`v(vv)`$–torsion and so on. The formulas (A.18) in Appendix present a local proof of this Theorem.
###### Definition 2.5
The curvature of a d–connection $`𝐃=(D,^{}D)`$ is defined
$$𝐑(𝐗,𝐘)𝐃_𝐗𝐃_𝐘𝐃_𝐘𝐃_𝐗𝐃_{[𝐗,𝐘]}$$
(7)
for any $`𝐗,𝐘\chi (𝐕).`$
Denoting $`h𝐑=R`$ and $`v𝐑=^{}R,`$ by straightforward calculations, one check the properties
$`R(𝐗,𝐘)^{}Z`$ $`=`$ $`0,^{}R(𝐗,𝐘)Z=0,`$
$`𝐑(𝐗,𝐘)𝐙`$ $`=`$ $`R(𝐗,𝐘)Z+^{}R(𝐗,𝐘)^{}Z`$
for any for any $`𝐗,𝐘,𝐙\chi (𝐕).`$
###### Theorem 2.2
The curvature $`𝐑`$ of a d–connection $`𝐃=(D,^{}D)`$ is completely defined by six d–curvatures
$`𝐑(X,Y)Z`$ $`=`$ $`\left(D_XD_YD_YD_XD_{[X,Y]}^{}D_{[X,Y]}\right)Z,`$
$`𝐑(X,Y)^{}Z`$ $`=`$ $`\left(D_XD_YD_YD_XD_{[X,Y]}^{}D_{[X,Y]}\right)^{}Z,`$
$`𝐑(^{}X,Y)Z`$ $`=`$ $`(^{}D_XD_YD_Y^{}D_XD_{[^{}X,Y]}{}_{}{}^{}D_{[{}_{}{}^{}X,Y]}^{})Z,`$
$`𝐑(^{}X,Y)^{}Z`$ $`=`$ $`(^{}D_X^{}D_Y^{}D_Y^{}D_XD_{[^{}X,Y]}{}_{}{}^{}D_{[^{}X,Y]}^{})^{}Z,`$
$`𝐑(^{}X,^{}Y)Z`$ $`=`$ $`(^{}D_XD_YD_Y^{}D_X^{}D_{[^{}X,^{}Y]})Z,`$
$`𝐑(^{}X,^{}Y){}_{}{}^{}Z`$ $`=`$ $`(^{}D_XD_YD_Y^{}D_X^{}D_{[^{}X,^{}Y]})^{}Z.`$
The proof of Theorems 2.1 and 2.2 is given for vector bundles provided with N–connection structure in Ref. . Similar Theorems and respective proofs hold true for superbundles , for noncommutative projective modules and for N–anholonomic metric–affine spaces , where there are also give the main formulas in abstract coordinate form. The formulas (A.23) from Appendix consist a coordinate proof of Theorem 2.2.
###### Definition 2.6
A metric structure $`\stackrel{˘}{g}`$ on a N–anholonomic space $`𝐕`$ is a symmetric covariant second rank tensor field which is not degenerated and of constant signature in any point $`𝐮𝐕.`$
In general, a metric structure is not adapted to a N–connection structure.
###### Definition 2.7
A d–metric $`𝐠=g_N^{}g`$ is a usual metric tensor which contracted to a d–vector results in a dual d–vector, d–covector (the duality being defined by the inverse of this metric tensor).
The relation between arbitrary metric structures and d–metrics is established by
###### Theorem 2.3
Any metric $`\stackrel{˘}{g}`$ can be equivalently transformed into a d–metric
$$𝐠=g(X,Y)+^{}g(^{}X,^{}Y)$$
(8)
for a corresponding N–connection structure.
Proof. We introduce denotations $`h\stackrel{˘}{g}(X,Y)g(X,Y)`$ and $`v\stackrel{˘}{g}(^{}X,^{}Y)`$ $`=^{}g(^{}X,^{}Y)`$ and try to find a N–connection when
$$\stackrel{˘}{g}(X,^{}Y)=0$$
(9)
for any $`𝐗,𝐘\chi (𝐕).`$ In local form, the equation (9) is just an algebraic equation for $`𝐍=\{N_i^a\},`$ see formulas (A.6), (A.7) and (A.8) and related explanations in Appendix. $`\mathrm{}`$
###### Definition 2.8
A d–connection $`𝐃`$ on $`𝐕`$ is said to be metric, i.e. it satisfies the metric compatibility (equivalently, metricity) conditions with a metric $`\stackrel{˘}{g}`$ and its equivalent d–metric $`𝐠,`$ if there are satisfied the conditions
$$𝐃_𝐗𝐠=\mathrm{𝟎}.$$
(10)
Considering explicit h– and v–projecting of (10), one proves
###### Proposition 2.2
A d–connection $`𝐃`$ on $`𝐕`$ is metric if and only if
$$D_Xg=0,D_X^{}g=0,^{}D_Xg=0,^{}D_X^{}g=0.$$
One holds this important
###### Conclusion 2.1
Following Propositions 2.1 and 2.2 and Theorem 2.3, we can elaborate the geometric constructions on a N–anholonomic manifold $`𝐕`$ in N–adapted form by considering N–adapted frames $`𝐞=(e,^{}e)`$ and co–frames $`\stackrel{~}{𝐞}=\left(\stackrel{~}{e},^{}\stackrel{~}{e}\right),`$ d–connection $`𝐃`$ and d–metric fields $`𝐠=[g,^{}g].`$
In Riemannian geometry, there is a preferred linear Levi–Civita connection $``$ which is metric compatible and torsionless, i.e.
$${}_{}{}^{}𝐓(𝐗,𝐘)_𝐗𝐘_𝐘𝐗[𝐗,𝐘]=0,$$
and defined by the metric structure. On a general N–anholonomic manifold $`𝐕`$ provided with a d–metric structure $`𝐠=[g,^{}g],`$ the Levi–Civita connection defined by this metric is not adapted to the N–connection, i. e. to the splitting (2). The h– and v–distributions are nonintegrable ones and any d–connection adapted to a such splitting contains nontrivial d–torsion coefficients. Nevertheless, one exists a minimal extension of the Levi–Civita connection to a canonical d–connection which is defined only by a metric $`\stackrel{˘}{g}.`$
###### Theorem 2.4
For any d–metric $`𝐠=[g,^{}g]`$ on a N–anholonomic manifold $`𝐕,`$ there is a unique metric canonical d–connection $`\widehat{𝐃}`$ satisfying the conditions $`\widehat{𝐃}𝐠=0`$ and with vanishing $`h(hh)`$–torsion, $`v(vv)`$–torsion, i. e. $`\widehat{T}(X,Y)=0`$ and $`^{}\widehat{T}(^{}X,^{}Y)=0.`$
Proof. The formulas (A.19) and (A.21) and related discussions in Appendix give a proof, in component form, of this Theorem.$`\mathrm{}`$
The following Corollary gathers some basic information about N–anholonomic manifolds.
###### Corollary 2.1
A N–connection structure defines three important geometric objects:
1. a (pseudo) Euclidean N–metric structure $`{}_{}{}^{\eta }𝐠=\eta _N^{}\eta ,`$ i.e. a d–metric with (pseudo) Euclidean metric coefficients with respect to $`\stackrel{~}{𝐞}`$ defined only by $`𝐍;`$
2. a N–metric canonical d–connection $`\widehat{𝐃}^N`$ defined only by $`{}_{}{}^{\eta }𝐠`$ and $`𝐍;`$
3. a nonmetric Berwald type linear connection $`𝐃^B.`$
Proof. Fixing a signature for the metric, $`sign^\eta 𝐠=(\pm ,\pm ,\mathrm{},\pm ),`$ we introduce these values in (A.8) we get $`{}_{}{}^{\eta }𝐠=\eta _N^{}\eta `$ of type (8), i.e. we prove the point 1. The point 2 is to be proved by an explicit construction by considering the coefficients of $`{}_{}{}^{\eta }𝐠`$ into (A.21). This way, we get a canonical d–connection induced by the N–connection coefficients and satisfying the metricity conditions (10). In an approach to Finsler geometry , one emphasizes the constructions derived for the so–called Berwald type d–connection $`𝐃^B,`$ considered to be the ”most” minimal (linear on $`\mathrm{\Omega })`$ extension of the Levi–Civita connection, see formulas (A.22). Such d–connections can be defined for an arbitrary d–metric $`𝐠=[g,^{}g],`$ or for any $`{}_{}{}^{\eta }𝐠=\eta _N^{}\eta .`$ They are only ”partially” metric because, for instance, $`D^Bg=0`$ and $`^{}D_{}^{B}{}_{}{}^{}g=0`$ but, in general, $`D_{}^{B}{}_{}{}^{}g0`$ and $`^{}D^Bg0,`$ i. e. $`𝐃^B𝐠0,`$ see Proposition 2.2. It is a more sophisticate problem to define spinors and supersymmetric physically valued models for such Finsler spaces, see discussions in . $`\mathrm{}`$
###### Remark 2.1
The geometrical objects $`\widehat{𝐃}^N,𝐃^B`$ for $`{}_{}{}^{\eta }𝐠,`$ nonholonomic bases $`𝐞=(e,^{}e)`$ and $`\stackrel{~}{𝐞}=\left(\stackrel{~}{e},^{}\stackrel{~}{e}\right),`$ see Proposition 2.1 and the N–connection curvature $`𝛀`$ (3), define completely the main properties of a N–anholonomic manifold $`𝐕.`$
It is possible to extend the constructions for any additional d–metric and canonical d–connection structures. For our considerations on nonholnomic Clifford/spinor structures, the class of metric d–connections plays a preferred role. That why we emphasize the physical importance of d–connections $`\widehat{𝐃}`$ and $`\widehat{𝐃}^N`$ instead of $`𝐃^B`$ or any other nonmetric d–connections.
Finally, in this section, we note that the d–torsions and d–curvatures on N–anholonomic manifolds can be computed for any type of d–connection structure, see Theorems 2.1 and 2.2 and the component formulas (A.18) and (A.23).
## 3 Examples of N–anholonomic spaces:
For corresponding parametrizations of the N–connection, d–metric and d–connection coefficients of a N–anholonomic space, it is possible to model various classes of (generalized) Lagrange, Finsler and Riemann–Cartan spaces. We briefly analyze three such nonholonomic geometric structures.
### 3.1 Lagrange–Finsler geometry
This class of geometries is usually defined on tangent bundles but it is possible to model such structures on general N–anholonomic manifolds, in particular in (pseudo) Riemannian and Riemann–Cartan geometry if nonholonomic frames are introduced into consideration . Let us outline the first approach when the N–anholonomic manifold $`𝐕`$ is taken to be just a tangent bundle $`(TM,\pi ,M),`$ where $`M`$ is a $`n`$–dimensional base manifold, $`\pi `$ is a surjective projection and $`TM`$ is the total space. One denotes by $`\stackrel{~}{TM}=TM\backslash \{0\}`$ where $`\{0\}`$ means the null section of map $`\pi .`$
We consider a differentiable fundamental Lagrange function $`L(x,y)`$ defined by a map $`L:(x,y)TML(x,y)`$ of class $`𝒞^{\mathrm{}}`$ on $`\stackrel{~}{TM}`$ and continuous on the null section $`0:MTM`$ of $`\pi .`$ The values $`x=\{x^i\}`$ are local coordinates on $`M`$ and $`(x,y)=(x^i,y^k)`$ are local coordinates on $`TM.`$ For simplicity, we consider this Lagrangian to be regular, i.e. with nondegenerated Hessian
$${}_{}{}^{L}g_{ij}^{}(x,y)=\frac{1}{2}\frac{^2L(x,y)}{y^iy^j}$$
(11)
when $`rank\left|g_{ij}\right|=n`$ on $`\stackrel{~}{TM}`$ and the left up ”L” is an abstract label pointing that the values are defined by the Lagrangian $`L.`$
###### Definition 3.1
A Lagrange space is a pair $`L^n=[M,L(x,y)]`$ with the tensor $`{}_{}{}^{L}g_{ij}^{}(x,y)`$ being of constant signature over $`\stackrel{~}{TM}.`$
The notion of Lagrange space was introduced by J. Kern and elaborated in details in Ref. as a natural extension of Finsler geometry.
###### Theorem 3.1
There are canonical N–connection $`{}_{}{}^{L}𝐍,`$ almost complex $`{}_{}{}^{L}𝐅,`$ d–metric $`{}_{}{}^{L}𝐠`$ and d–connection $`{}_{}{}^{L}\widehat{𝐃}`$ structures defined by a regular Lagrangian $`L(x,y)`$ and its Hessian $`{}_{}{}^{L}g_{ij}^{}(x,y)`$ (11).
Proof. The canonical $`{}_{}{}^{L}𝐍`$ is defined by certain nonlinear spray configurations related to the solutions of Euler–Lagrange equations, see the local formula (A.27) in Appendix. It is given there the explicit matrix representation of $`{}_{}{}^{L}𝐅`$ (A.28) which is a usual definition of almost complex structure, after $`{}_{}{}^{L}𝐍`$ and N–adapted bases have been constructed. The d–metric (A.29) is a local formula for $`{}_{}{}^{L}𝐠.`$ Finally, the canonical d–connection $`{}_{}{}^{L}\widehat{𝐃}`$ is a usual one but for $`{}_{}{}^{L}𝐠`$ and $`{}_{}{}^{L}𝐍`$ on $`\stackrel{~}{TM}.\mathrm{}`$
A similar Theorem can be formulated and proved for the Finsler geometry:
###### Remark 3.1
A Finsler space defined by a fundamental Finsler function $`F(x,y),`$ being homogeneous of type $`F(x,\lambda y)=\lambda F(x,y),`$ for nonzero $`\lambda ,`$ may be considered as a particular case of Lagrange geometry when $`L=F^2.`$
From the Theorem 3.1 and Remark 3.1 one follows:
###### Result 3.1
Any Lagrange mechanics with regular Lagrangian $`L(x,y)`$ (any Finsler geometry with fundamental function $`F(x,y))`$ can be modelled as a nonhlonomic Riemann–Cartan geometry with canonical structures $`{}_{}{}^{L}𝐍,`$ $`{}_{}{}^{L}𝐠`$ and $`{}_{}{}^{L}\widehat{𝐃}`$ ($`{}_{}{}^{F}𝐍,`$ $`{}_{}{}^{F}𝐠`$ and $`{}_{}{}^{F}\widehat{𝐃})`$ defined on a corresponding N–anholonomic manifold $`𝐕.`$
It was concluded that any regular Lagrange mechanics/Finsler geometry can be geometrized/modelled as an almost Kähler space with canonical N–connection distribution, see and, for N–anholonomic Fedosov manifolds, . Such approaches based on almost complex structures are related with standard sympletic geometrizations of classical mechanics and field theory, for a review of results see Ref. .
For applications in optics of nonhomogeneous media and gravity (see, for instance, Refs. ), one considers metrics of type $`g_{ij}e^{\lambda (x,y)}{}_{}{}^{L}g_{ij}^{}(x,y)`$ which can not be derived from a mechanical Lagrangian but from an effective ”energy” function. In the so–called generalized Lagrange geometry, one introduced Sasaki type metrics (A.29), see the Appendix, with any general coefficients both for the metric and N–connection.
### 3.2 N–connections and gravity
Now we show how N–anholonomic configurations can defined in gravity theories. In this case, it is convenient to work on a general manifold $`𝐕,dim𝐕=n+m`$ enabled with a global N–connection structure, instead of the tangent bundle $`\stackrel{~}{TM}.`$
For N–connection splittings of (pseudo) Riemannian–Cartan spaces of dimension $`(n+m)`$ (there were also considered (pseudo) Riemannian configurations), the Lagrange and Finsler type geometries were modelled by N–anholonomic structures as exact solutions of gravitational field equations . Inversely, all approaches to (super) string gravity theories deal with nontrivial torsion and (super) vielbein fields which under corresponding parametrizations model N–anholonomic spaces . We summarize here some geometric properties of gravitational models with nontrivial N–anholonomic structure.
###### Definition 3.2
A N–anholonomic Riemann–Cartan manifold $`{}_{}{}^{RC}𝐕`$ is defined by a d–metric $`𝐠`$ and a metric d–connection $`𝐃`$ structures adapted to an exact sequence splitting (1) defined on this manifold.
The d–metric structure $`𝐠`$ on$`{}_{}{}^{RC}𝐕`$ is of type (8) and satisfies the metricity conditions (10). With respect to a local coordinate basis, the metric $`𝐠`$ is parametrized by a generic off–diagonal metric ansatz (A.7), see Appendix. In a particular case, we can take $`𝐃=\widehat{𝐃}`$ and treat the torsion $`\widehat{𝐓}`$ as a nonholonomic frame effect induced by nonintegrable N–splitting. For more general applications, we have to consider additional torsion components, for instance, by the so–called $`H`$–field in string gravity.
Let us denote by $`Ric(𝐃)`$ and $`Sc(𝐃),`$ respectively, the Ricci tensor and curvature scalar defined by any metric d–connection $`𝐃`$ and d–metric $`𝐠`$ on $`{}_{}{}^{RC}𝐕,`$ see also the component formulas (A.24), (A.25) and (A.26) in Appendix. The Einstein equations are
$$En(𝐃)Ric(𝐃)\frac{1}{2}𝐠Sc(𝐃)=𝚼$$
(12)
where the source $`𝚼`$ reflects any contributions of matter fields and corrections from, for instance, string/brane theories of gravity. In a closed physical model, the equation (12) have to be completed with equations for the matter fields, torsion contributions and so on (for instance, in the Einstein–Cartan theory one considers algebraic equations for the torsion and its source)… It should be noted here that because of nonholonomic structure of $`{}_{}{}^{RC}𝐕,`$ the tensor $`Ric(𝐃)`$ is not symmetric and that $`𝐃\left[En(𝐃)\right]0`$ which imposes a more sophisticate form of conservation laws on such spaces with generic ”local anisotropy”, see discussion in (this is similar with the case when the nonholonomic constraints in Lagrange mechanics modifies the definition of conservation laws). A very important class of models can be elaborated when $`𝚼=diag[\lambda ^h(𝐮)g,\lambda ^v(𝐮)^{}g],`$ which defines the so–called N–anholonomic Einstein spaces.
###### Result 3.2
Various classes of vacuum and nonvacuum exact solutions of (12) parametrized by generic off–diagonal metrics, nonholonomic vielbeins and Levi–Civita or non–Riemannian connections in Einstein and extra dimension gravity models define explicit examples of N–anholonomic Einstein–Cartan (in particular, Einstein) spaces.
Such exact solutions (with noncommutative, algebroid, toroidal, ellipsoid, … symmetries) have been constructed in Refs. . We note that a subclass of N–anholonomic Einstein spaces was related to generic off–diagonal solutions in general relativity by such nonholonomic constraints when $`Ric(\widehat{𝐃})=Ric()`$ even $`\widehat{𝐃},`$ where $`\widehat{𝐃}`$ is the canonical d–connection and $``$ is the Levi–Civita connection, see formulas (A.19) and (A.20) in Appendix and details in Ref. .
A direction in modern gravity is connected to analogous gravity models when certain gravitational effects and, for instance, black hole configurations are modelled by optical and acoustic media, see a recent review or results in . Following our approach on geometric unification of gravity and Lagrange regular mechanics in terms of N–anholonomic spaces, one holds
###### Theorem 3.2
A Lagrange (Finsler) space can be canonically modelled as an exact solution of the Einstein equations (12) on a N–anholonomic Riemann–Cartan space if and only if the canonical N–connection $`{}_{}{}^{L}𝐍`$ ($`{}_{}{}^{F}𝐍`$), d–metric $`{}_{}{}^{L}𝐠`$ ($`{}_{}{}^{F}𝐠)`$ and d–connection $`{}_{}{}^{L}\widehat{𝐃}`$ ($`{}_{}{}^{F}\widehat{𝐃})`$ structures defined by the corresponding fundamental Lagrange function $`L(𝐱,𝐲)`$ (Finsler function $`F(𝐱,𝐲))`$ satisfy the gravitational field equations for certain physically reasonable sources $`𝚼.`$
Proof. We sketch the idea: It can be performed in local form by considering the Einstein tensor (A.26) defined by the $`{}_{}{}^{L}𝐍`$ ($`{}_{}{}^{F}𝐍`$) in the form (A.27) and $`{}_{}{}^{L}𝐠`$ ($`{}_{}{}^{F}𝐠)`$ in the form (A.29) inducing the canonical d–connection $`{}_{}{}^{L}\widehat{𝐃}`$ ($`{}_{}{}^{F}\widehat{𝐃}).`$ For certain zero or nonzero $`𝚼`$, such N–anholonomic configurations may be defined by exact solutions of the Einstein equations for a d–connection structure. A number of explicit examples were constructed for N–anholonomic Einstein spaces .$`\mathrm{}`$
It should be noted that Theorem 3.2 states explicit conditions when the Result 3.1 holds for N–anholonomic Einstein spaces.
###### Conclusion 3.1
Generic off–diagonal metric and vielbein structures in gravity and regular Lagrange mechanics models can be geometrized in a unified form on N–anholonomic manifolds. In general, such spaces are not spin and this presents a strong motivation for elaborating the theory of nonholonomic gerbes and related Clifford/spinor structures developed in this work.
Following this Conclusion, it is not surprizing that a lot of gravitational effects (black hole configurations, collapse scenaria, cosmological anisotropies etc) can be modelled in nonlinear fluid, acoustic or optic media.
## 4 Lifts of Nonholonomic Bundle Gerbes and Connections
In this section, we present an introduction into the geometry of lifts of nonholonomic bundle gerbes and related N–anholonomic modules. We define connections and curvatures for such bundle modules. This material reproduces, in the corresponding holonomic limits, certain fundamental results from .
### 4.1 N–anholonomic bundle gerbes and their lifts
#### 4.1.1 Local constructions
On N–anholonomic manifolds, one deals with nonintegrable h– and v–splitting of geometric objects, described by the so–called d–objects (for instance, d–vectors, d–spinors, d–tensors, d–connections, … like we considered in the previous section). It is convenient to introduce the concept of Lie d–group $`𝐆=(G,^{}G)`$ which is just a couple of two usual Lie groups $`G`$ and$`{}_{}{}^{}G`$ associated to a N–connection splitting (2). We conventionally consider a central extension of a finite dimensional of Lie d–groups $`𝐆`$ to $`\stackrel{ˇ}{𝐆},`$ defined by a map $`\pi :`$ $`\stackrel{ˇ}{𝐆}𝐆`$ such that it is defined the exact sequence
$$0_k\stackrel{ˇ}{𝐆}𝐆1$$
(13)
where $`_k=/k`$ denotes the cyclic subgroup of the circle $`U(1).`$ This sequence of d–groups splits into respective horizontal component
$$0_k\stackrel{ˇ}{G}G1$$
and vertical component
$$0_k^{}\stackrel{ˇ}{G}^{}G1.$$
Let us denote by $`𝐔`$ and $`\stackrel{ˇ}{𝐔}`$ the corresponding right principal sets: the are just $`𝐆`$ and $`\stackrel{ˇ}{𝐆}`$ but conventionally re–defined in order to consider distinguished (not mixing the h– and v–subsets) actions of $`𝐆`$ on $`𝐔`$ and $`\stackrel{ˇ}{𝐆}`$ on $`\stackrel{ˇ}{𝐔}.`$ We consider an equivariant $`\stackrel{ˇ}{𝐆}`$ bundle $`𝐯_𝐔𝐔\times 𝐯,`$ where $`𝐯`$ is a d–vector space and a finite–dimensional representation $`\rho :\stackrel{ˇ}{𝐆}GL(𝐯)`$ with the $`\stackrel{ˇ}{𝐆}`$ action
$$\stackrel{ˇ}{𝐠}(𝐮,𝐯)=(𝐮\stackrel{ˇ}{𝐠}^1,\rho (\stackrel{ˇ}{𝐠})𝐯)=\left\{\begin{array}{c}(x\stackrel{ˇ}{g}^1,\rho (\stackrel{ˇ}{g})v),\\ \left(y^{}\stackrel{ˇ}{g}^1,^{}\rho (^{}\stackrel{ˇ}{g})^{}v\right)\end{array}\right\}.$$
The pull–back of the $`_k`$ distinguished bundle $`\stackrel{ˇ}{𝐆}𝐆`$ is
$$𝚿=\tau ^{}\stackrel{ˇ}{𝐆}𝐔\times 𝐔$$
defined by $`𝚿_{(𝐮_1,𝐮_2)}\{\pi (\stackrel{ˇ}{𝐠})=\tau (𝐮_1,𝐮_2)\}`$ for any $`\stackrel{ˇ}{𝐠}\stackrel{ˇ}{𝐆},`$ where $`\tau :𝐔\times 𝐔`$ $`𝐆`$ is a canonical map $`𝐮_1\tau (𝐮_1,𝐮_2)𝐮_2`$ translating $`𝐮_1`$ into $`𝐮_2.`$ So, $`𝚿_{(𝐮_1,𝐮_2)}`$ is the set of all distinguished lifts of $`\tau (𝐮_1,𝐮_2)`$ to $`\stackrel{ˇ}{𝐆}`$ and $`𝚿`$ is the $`_k`$–principal bundle provided with a trivial N–connection (in this case, with zero N–connection curvature). The bundle $`𝚿`$ has a module the $`\stackrel{ˇ}{𝐆}`$–equivariant bundle $`𝐯_𝐔𝐔.`$ This follows from the fact that for any two pull–backs $`𝐔_1`$ and $`𝐔_2`$ of $`𝐯_𝐔`$ as two respective projections $`𝐔\times 𝐔𝐔`$ one has that $`𝚿_{(𝐮_1,𝐮_2)}\stackrel{ˇ}{𝐆}`$ transforms the distinguished fiber in $`(𝐮_1,𝐮_2)`$ of $`𝐔_1`$ into the corresponding one of $`𝐔_2`$ related by the representation map $`\rho .`$ Having also the $`\stackrel{ˇ}{𝐆}`$–equivariance, of $`End(𝐯_𝐔),`$ we can write $`End(𝐯_𝐔)/𝐆=\mathrm{𝐄𝐧𝐝}(𝐯)`$ for distinguished endomorphysms.
#### 4.1.2 Global constructions
The above presented constructions can be globalized to the case of N–anholonomic manifold $`𝐕`$ instead of the d–vector space $`𝐯`$ (the d–objects with trivial splitting can be considered for any point of $`𝐕).`$ The procedure is completely similar to that given for ”holonomic” manifolds in but it should be performed in a form to preserve the N–connection splitting (2). This may be achieved by applied globalizing the bundle $`𝚿`$ and transforming it into a nonholonomic bundle.
Having in mind the distinguished extension (13), we replace the set $`𝐔`$ by a principal N–anholonomic $`𝐆`$–bundle $`\pi :𝐁𝐕`$ and consider the product $`𝐁\times 𝐁𝐕`$ instead of $`𝐔\times 𝐔.`$ Like for a trivial point of $`𝐔,`$ the globalized map $`\tau :𝐁\times 𝐁𝐆`$ allows us to introduce $`𝚿=\tau ^{}\stackrel{ˇ}{𝐆}`$ being the $`_k`$–bundle over $`𝐁\times 𝐁`$ which defines a lifting N–anholonomic bundle gerbe if to follow the terminology for holonomic constructions, .
We can consider d–tensor objects of weight $`q`$ as $`𝚿^q`$–modules being nonholonomic variants of bundle gerbe modules for the N–anholonomic bundle gerbe $`𝚿^q𝚿^q.`$ In more explicit form, we use a $`\stackrel{ˇ}{𝐆}`$–equivariant bundle $`𝐖𝐁`$ for the action of $`\stackrel{ˇ}{𝐆}`$ of $`𝐁:`$
###### Definition 4.1
The N–anholonomic $`\stackrel{ˇ}{𝐆}`$–equivariant bundle $`𝐖𝐁`$ with defined action of weight $`q`$ of the isotropy distinguished subgroups states $`𝐖`$ as a $`𝚿^q`$–module.
The space $`𝐖`$ can be also treated as a vector bundle direct sum of $`𝚿^q`$–modules, all adapted to the N–connection structure, i.e. preserving the h– and v–decomposition by (2). This allows us to concentrate the attention only to ”boldfaced” $`𝚿^q`$–modules carrying out all information about nonholnomic and non–trivial topological configurations. Such constructions run parallel to the usual theory of vector bundles provided with N–connection structure and in a more formalized form (unifying the approaches to gauge fields, gravity and geometrized mechanics) to N–anholonomic manifolds.
It should be noted that if $`𝐖`$ is a $`𝚿`$–module, then we get a trivial module but it can provided with a nontrivial N–connection (with nonvanishing N–connection curvature). In such cases, one works with constructions of type $`(𝐖)=End(𝐖)/𝐆`$ splitting into h– and v–subspaces and this allow us to reformulate in nonholnomic form, for N–anholonomic $`𝚿`$–modules, the main properties of such spaces formally formulated for trivial N–connection structure .
###### Proposition 4.1
The $`𝚿^q`$–modules satisfy the following N–adapted properties:
1. N–anholonomic $`𝚿^q`$–modules and bundles on $`𝐕`$ are bijective equivalent.
2. The bundle of N–adapted endomorphisms of a $`𝚿^q`$–module is a $`𝚿^0`$–module.
3. The direct sum of two $`𝚿^q`$–modules is a $`𝚿^q`$–module.
4. The d–tensor product of a $`𝚿^{q_1}`$–module to a $`𝚿^{q_2}`$–module results in a $`𝚿^{q_1+q_2}`$–module.
We omit the proof of these properties following from an explicit Cech description of the above structures in N–adapted from (dubbing the constructions from for h– and v–configurations). Here we note that the elements of cohomological classes, like $`[e]H^3(𝐕,_k)H^2(𝐕,U(1))`$ and $`\delta [e]H^3(𝐕,_k),`$ are defined for N–anholonomic manifolds, see Ref. for an introduction in $`K`$–theory and related cohomological calculus. This results in distinguished (by N–connection structure) $`K`$–group of the semi–group of N–anholonomic $`𝚿`$–modules. <sup>3</sup><sup>3</sup>3We emphasize that we wrote $`𝚿`$–modules instead of $`𝚪`$–modules because in this work the symbol $`𝚪`$ is used for d–connections.
### 4.2 Curvatures for N–anholonomic bundle gerbe modules
For a holonomic manifold, because $`_k`$ is finite, there is a natural $`\stackrel{ˇ}{𝐆}`$–equivariant flat connection $`{}_{}{}^{flat}_{𝐗}^{}`$ on any cart from a covering of bundle $`𝐯_𝐔.`$ For N–anholonomic manifolds the role of flat connection is played by metric canonical d<sub>N</sub>–connection $`\widehat{𝐃}^N`$ defined by a (pseudo) Euclidean N–metric structure $`{}_{}{}^{\eta }𝐠=\eta _N^{}\eta `$ and the N–connection $`𝐍,`$ see Corollary 2.1. If a d–metric structure $`𝐠=[g,^{}g]`$ is stated on such a N–anholonomic manifold, we shall work with the corresponding canonical d–connection $`\widehat{𝐃},`$ see Theorem 2.4. For simplicity, in this section we shall derive our constructions starting from $`\widehat{𝐃}^N`$ but we note that, in general, we can work with an arbitrary d–connection $`𝐃`$ lifted on $`𝐖`$ as a distinguished linear operator, N–adapted to (2), acting in the space of d–forms $`\omega =(\omega ^0,\omega ^1,\mathrm{}),`$ where, for instance, $`\omega ^1`$ denotes the space of 1–forms distinguished by the N–connection structure. Let us consider the d–operator
$$\stackrel{}{𝐃}:\omega ^0(𝐁,𝐖)\omega ^1(𝐁,𝐖).$$
For a necessary small open subset $`𝐔`$ $`𝐕,`$ we can identify the restriction of $`𝐁`$ to $`𝐔`$ with $`𝐔\times 𝐕`$ and their restriction of $`𝐖`$ with $`𝐖_𝐕.`$ In result, we may write
$$\stackrel{}{𝐃}=\widehat{𝐃}_𝐕^N+𝐃_B,$$
(14)
for $`𝐃_B`$ being a pull–back of a connection from the base $`𝐔.`$ This way, $`\stackrel{}{𝐃}`$ is defined as a $`𝚿`$–module d–connection if it is equivariant for the group $`\stackrel{ˇ}{𝐆}^\mathrm{\#}=\left(U(1)\times \stackrel{ˇ}{𝐆}\right)/_k`$ with $`_kU(1)\times \stackrel{ˇ}{𝐆}`$ parametrized as a h– and v–distinguished inclusions by anti–diagonal subroups. Such a d–connection satisfies the rule
$$\stackrel{}{𝐃}\left(f\phi \right)=𝐞_\mu (f)𝐞^\mu _N\phi +f_N\stackrel{}{𝐃}\phi $$
(15)
for any function $`f`$ on $`𝐕`$ and section $`\phi `$ of $`𝐖`$ where $`𝐞_\mu `$ and $`𝐞^\mu `$ are N–elongated operators (A.3) and (A.4).
Let us consider two $`𝚿^q`$–module d–connections $`\stackrel{}{𝐃}_1`$ and $`\stackrel{}{𝐃}_2`$ on $`𝐖.`$ The distorsion $`\stackrel{}{𝐏}=`$ $`\stackrel{}{𝐃}_1`$ $`\stackrel{}{𝐃}_2`$ is $`\stackrel{ˇ}{𝐆}`$–equivariant and belongs to the d–vector space $`\omega ^1(𝐁,\mathrm{𝐄𝐧𝐝}(𝐖)),`$ this follows from (15). For any vertical to $`𝐕`$ d–vector $`\lambda ,`$ one holds $`\lambda \stackrel{}{𝐏}=0.`$ This allows us to ”divide” on $`𝐆`$ and transform $`\stackrel{}{𝐏}`$ into an element of $`\omega ^1(𝐁,(𝐖)),`$ i.e. by such N–adapted distorsions we are able to generate all $`𝚿^q`$–module d–connections starting from (14). In result, we proved
###### Proposition 4.2
The set of $`𝚿^q`$–module d–connections on $`𝐖`$ is a N–distinguished affine space generated by N–adapted distorsions as elements of $`\omega ^1(𝐁,`$ $`(𝐖)).`$
The curvature of a d–connection $`\stackrel{}{𝐃}`$ (14) is to be constructed by globalizing the results of Theorem 2.2 (which is a very similar to the proof of the previous Proposition):
###### Theorem 4.1
The curvature of a $`𝚿^q`$–module d–connections on $`𝐖`$ descends to define an element $`\stackrel{}{𝐑}\omega ^2(𝐁,(𝐖)).`$
The d–connection (14) is defined by a N–adapted tensor product. This extends to a straightforward proof of a corresponding result for curvature:
###### Corollary 4.1
For any $`𝚿^q`$–module d–connections $`\stackrel{}{𝐃}`$ and $`\stackrel{}{𝐃}^{},`$ respectively, on N–anholonomic $`\stackrel{ˇ}{𝐆}`$–equivariant bundles $`𝐖`$ and $`𝐖^{},`$ we can compute the curvature of the d–tensor product connection $`\stackrel{}{𝐃}`$ on $`𝐖`$ $`𝐖^{},`$
$$𝐑_B=\stackrel{}{𝐑}1+1\stackrel{}{𝐑}^{}\omega ^2(𝐕,(𝐖𝐖^{}))=\omega ^2(𝐕,(𝐖)(𝐖^{})).$$
In order to define the (twisted) Chern character it is enough to have the data for a $`𝚿`$–module d–connection $`\stackrel{}{𝐃}`$ and its descendent curvature $`\stackrel{}{𝐑}`$ . For N–anholonomic configurations, the constructions depend on the fact if there N–anholonomic manifold is provided or not with a d–metric structure.
###### Definition 4.2
The (twisted and nonholonomic) Chern character of a $`𝚿^q`$–module is defined by the curvature of d–connection $`\stackrel{}{𝐃}`$ induced by the N–connection structure,
$$ch(\stackrel{}{𝐑})=tr\mathrm{exp}\frac{\stackrel{}{𝐑}}{2\pi i}.$$
(16)
###### Remark 4.1
If additionally to the N–connection structure on $`𝐕,`$ it is defined a d–metric structure $`𝐠,`$ the corresponding Chern character must be computed by using the $`\stackrel{}{𝐑}^{}`$ defined as a distorsion from the nonholonomic configuration stated by a d–metric $`{}_{}{}^{\mathrm{\eta }}\mathrm{g}=[\mathrm{\eta },^{}\mathrm{\eta }]`$ (inducing together with $`N`$ the canonical d–connection $`\widehat{𝐃}^N`$ and $`\stackrel{}{𝐑})`$to a d–connection $`𝐠=[g,^{}g]`$ (inducing the canonical d–connection $`\widehat{𝐃}`$ and curvature $`\stackrel{}{𝐑}^{}.`$
The values $`ch(\stackrel{}{𝐑})`$ and/or $`ch(\stackrel{}{𝐑}^{})`$ are closed and this mean that the corresponding de Rham cohomology classes are independent of the choice of $`𝚿^q`$–module d–connections if a N–connection structure is prescribed. This has a number of interesting applications in modern geometric mechanics, generalized Finsler geometry and gravity with nontrivial N–anholonomic structures:
###### Result 4.1
Any regular Lagrange, or Finsler, configuration is topologically characterized by the corresponding canonical (twisted) Chern character (16) computed by using the curvature $`{}_{}{}^{L}\stackrel{}{𝐑},`$ or $`{}_{}{}^{F}\stackrel{}{𝐑},`$ induced by the curvature (7) defined by the N–connection$`{}_{}{}^{L}𝐍,`$ or$`{}_{}{}^{F}𝐍,`$ in the form (A.27) and d–metric $`{}_{}{}^{L}𝐠,`$ or$`{}_{}{}^{F}𝐠,`$ in the form (A.29) defining the canonical d–connection $`{}_{}{}^{L}\widehat{𝐃}`$ ($`{}_{}{}^{F}\widehat{𝐃}).`$
The set of exact solutions with generic off–diagonal metrics, nonholnomic frames and various type of local anisotropy, noncommutative and/or Lie algebroid symmetries constructed in Refs. can be globalized for gerbe configurations with nontrivial N–connection structure, i.e. one holds
###### Result 4.2
The geometric objects for a N–anholonomic Riemann–Cartan manifold $`{}_{}{}^{RC}𝐕`$ can be globalized to N–anholonomic gerbe configurations and characterized by the corresponding (twisted–anholonomic) Chern character (16). This character is computed by using the curvature $`\stackrel{}{𝐑}`$ induced by the curvature (7) defined by the N–connection$`𝐍,`$ d–metric $`𝐠`$ (8) and the canonical d–connection $`\widehat{𝐃}.`$
Finally, in this section, we conclude that the last two Results state new types of (topological) symmetries and a new classification of regular Lagrange systems, Finsler spaces and Einstein–Cartan spaces provided with N–connection structure.
## 5 Nonholonomic Clifford Gerbes and Modules
This section presents a development of the geometry of N–anholonomic manifolds and related nonholonomic Clifford and Dirac structures . The reader may consult Refs. for local component representations of the results and related local calculus and proofs.
### 5.1 Clifford d–algebras and N–anholonomic bundles
This work states an explicit example of generalized spinor constructions by considering in sequence (13) the d–groups $`𝐆=\mathrm{𝐒𝐩𝐢𝐧}(n+m)`$ and $`\stackrel{ˇ}{𝐆}=\mathrm{𝐒𝐎}`$ $`(n+m)`$ where the boldfaced d–groups split respectively into h– and v–components $`\mathrm{𝐒𝐩𝐢𝐧}(n+m)=\{Spin(n),Spin(m)\}`$ and $`\mathrm{𝐒𝐎}(n+m)=\{SO(n),`$ $`SO(m)\}.`$ One get the central extension
$$0_k\mathrm{𝐒𝐩𝐢𝐧}(n+m)\mathrm{𝐒𝐎}(𝐧+𝐦)1$$
splitting into respective h– and v–components,
$$0_kSpin(n)SO(n)1$$
and
$$0_kSpin(m)SO(m)1.$$
Let us consider two real vector spaces $`v`$ and $`{}_{}{}^{}v`$ of dimension $`n`$ and $`m`$ each provided with positive defined scalar products and defining a d–vector space $`𝐯=v^{}v.`$ We denote by $`C(v)`$ and $`C(^{}v)`$ the corresponding $`_2`$ graded Clifford algebras defining a Clifford d–algebra
$$𝐂(𝐯)=C_+(v)C_{}(v)_NC_+(^{}v)C_{}(^{}v).$$
The splitting $`\pm `$ is related to the chirality operator $`\gamma =\pm `$ on $`C_\pm .`$ A hermitian Clifford d–module is a $`_2`$–graded d–vector space $`𝐯^E`$ provided with complex scalar products on the h– and v–components. The endomorphisms of spin representation $`S=S^+S^{}`$ and $`{}_{}{}^{}S=^{}S^+^{}S^{}`$ define respectively the hermitian Clifford modules for conventional h– and v–subspaces, $`C(v)=End(S)`$ and $`C(^{}v)=End(^{}S).`$ Any hermitian Clifford d–modules of finite dimension can be represented in the form $`𝐯^E=S𝐯^C`$ where $`𝐯^C`$ is a complex d–vector space on which $`𝐂(𝐯)`$ acts trivially in distinguished from. We can identify
$$𝐯^C=Hom(𝐒,𝐯^E)\text{ and }End(𝐯^C)=End(𝐯^E)$$
supposing that such maps are h- and v–split and commute with the action of $`𝐂(𝐯).`$
We take the bundle $`𝐁`$ from section 4.1.2 to be the bundle of N–adapted orthogonal frames (see Proposition 2.1) on $`T𝐕.`$ In the spin case, the construction of $`\stackrel{ˇ}{𝐆}^\mathrm{\#}`$ is that for the d–group $`\mathrm{𝐒𝐩𝐢𝐧}_c(n+m)`$ considered in for spin N–anholonomic manifolds, when $`𝚿`$ is a trivial $`_2`$ N–anholonomic bundle gerbe and $`spin`$–c when $`𝚿_c`$ is a trivial $`U(1)`$ bundle gerbe provided with N–connection stucture.
###### Definition 5.1
The lifting bundle gerbe $`𝚿`$ for the case $`𝐆=\mathrm{𝐒𝐩𝐢𝐧}(n+m)`$ and $`\stackrel{ˇ}{𝐆}=\mathrm{𝐒𝐎}(n+m)`$ is called the spin–bundle N–anholonomic gerbe.
We can consider half–spin representations $`S^\pm `$ of $`Spin(n)`$ and $`{}_{}{}^{}S_{}^{\pm }`$ of $`Spin`$ $`(m)`$ and introduce the d–spin representations
$$𝐒=\left(S^+S^{}\right)_N\left({}_{}{}^{}S_{}^{+}^{}S^{}\right).$$
(17)
###### Definition 5.2
The $`𝚿^1`$–modules associated to the N–adapted d–spin representation (17) define the N–anholonomic spin $`𝚿^1`$–modules generalizing the concept of d–spin bundles on $`𝐕.`$
The above mentioned spin constructions have a straightforward extension to even–dimensional oriented N–anholomic Riemann–Cartan manifolds (this holds always for oriented Lagrange–Finsler spaces), denoted $`𝐕^{2n}.`$ One introduces the N–anholonomic bundle of complex Clifford d–algebras of $`T^{}𝐕^{2n}`$ and consider the Clifford distinguished map (multiplication) $`𝐜:T^{}𝐕^{2n}`$ $`𝐂(𝐕^{2n}),`$ where formally $`𝐯𝐕^{2n}.`$
###### Definition 5.3
A N–anholonomic Clifford module (in brief, Clifford d–module) is a complex $`_2`$–graded hermitian N–anholonomic vector bundle
$$𝐄=E_+E_{}_N^{}E_+^{}E_{}$$
over $`𝐕^{2n}`$ satisfying the properties that $`𝐄_u`$ is a hermitian Clifford d–module for $`𝐂_u(𝐕^{2n})`$ in each point $`𝐮𝐕^{2n}`$ and that the sub–bundles $`E_+`$ and $`{}_{}{}^{}E_{+}^{}`$ are respectively orthogonal to $`E_{}`$ and $`{}_{}{}^{}E_{}^{}.`$
We consider the spin–bundle N–anholonomic gerbe $`𝚿`$ from Definition 5.1 and the pull–back of $`𝐄`$ to $`𝐁,`$ denoted $`𝐄_𝐁=\pi ^1(𝐄),`$ where $`\pi :𝐁𝐕^{2n}`$ is the bundle of N–adapted frames on $`𝐕^{2n}.`$ In any point $`b𝐁,`$ there is an isomorphism transforming $`\left(𝐄_𝐁\right)_p`$ into $`𝐂(^{n+m})`$ Clifford d–module. We have $`\left(𝐄_𝐁\right)_p=𝐒𝐯_{(b)}^C`$ for $`𝐯_{(b)}^C=Hom(𝐒,\left(𝐄_𝐁\right)_p)`$ with the homomorphisms defined on $`𝐂(^{n+m}).`$ The construction can be globalized, $`𝐄_𝐁=𝐒𝐯^C.`$ The action of $`\mathrm{𝐒𝐩𝐢𝐧}(n+m)`$ on $`𝐒`$ induces also an action $`𝐄_𝐁`$ and transforms it into a N–anholonomic $`𝚿^1`$–module. In result, we proved
###### Theorem 5.1
For a N–anholonmic spin bundle gerbe $`𝐒_𝐁,`$ every Clifford d–module $`𝐄`$ on N–anholonomic $`𝐕^{2n},`$ with its nonholonomic bundle gerbe $`𝚿`$ , has the form $`𝐄_𝐁=𝐒𝐯^C`$ for some N–anholonomic $`𝚿^1`$–module $`𝐯^C.`$
This theorem generalizes for N–anholonomic spaces some similar results given in . For spin N–anholonomic manifolds $`𝐕^{2n}`$ considered in , we have that every Clifford d–module is a d–tensor product of an N–anholonomic spin bundle with an arbitrary bundle.
### 5.2 N–anholonomic Dirac operators and gerbes
#### 5.2.1 The index topological formula for holonomic Dirac operators
Let us remember the Atiyah–Singer index topological formula for the Dirac operator :
$$ind(D_\mathrm{\Gamma }^+)dim\mathrm{ker}(D_\mathrm{\Gamma }^+)dimco\mathrm{ker}(D_\mathrm{\Gamma }^+)=<\widehat{A}(M)ch(W),[M]>$$
(18)
where the compact $`M`$ is an oriented even dimensional spin manifold with spin–bundles $`S^\pm `$ and $`\mathrm{\Gamma }`$ is a unitary connection on the vector bundle $`W,`$ see details on definitions and denotations in Ref. (below, we shall give details for N–anholonomic configurations). In this formula, we use the genus of the manifold
$$\widehat{A}(M)\left|det\frac{R}{2\mathrm{sinh}(R/2)}\right|^{1/2}$$
determined by the Riemannian curvature $`R`$ of the manifold $`M.`$ The operator $`D_\mathrm{\Gamma }^+`$ is the so–called coupled Dirac operator (first order differential operator) acting in the form $`D_\mathrm{\Gamma }^+:C^{\mathrm{}}(M,E^+)C^{\mathrm{}}(M,E^{}),`$ for $`E^\pm S^\pm W.`$ This Dirac operator can be introduced for non–spin manifolds even itself this object is not well defined. In a formal way, we can induce the Dirac operator as a compatible connection on $`E=E^+E^{}`$ treated as a Clifford module with multiplication extended to act as the identity on $`W.`$
For non–spin manifolds, one exists an index formula for Dirac operators defined on hermitian Clifford modules $`E,`$
$$ind(D_\mathrm{\Gamma }^+)=<\widehat{A}(M)ch(E/S),[M]>$$
(19)
which, if $`M`$ is spin and $`E^\pm S^\pm W,`$ the relative Chern character $`ch(E/S)`$ reduces to the Chern character of $`W,`$ i.e. to $`ch(W).`$ We note that one may be not possible to define a canonical trivialization of $`M`$ but it is supposed that one exist a canonical no were vanishing (volume) density $`[M]`$ which allows us to perform the integration. This always holds for the Riemannian manifolds. The aim of next section is to prove that formulas (18) and (19) can be correspondingly generalized for N–anholonomic manifolds provided with d–metric and d–connection structures.
#### 5.2.2 Twisted nonnholonomic Dirac operators on Clifford gerbes
Let us go to the Definition 5.1 of the spin–bundle N–anholonomic gerbe $`𝚿`$ derived for an N–anholonmic manifold $`𝐕.`$ The Clifford multiplication is parametrized by N–adapted maps between such $`𝚿`$–modules,
$$c:\left(_𝐁^n\right)^{}S_𝐁^+S_𝐁^{}\text{ and }^{}c:\left(_𝐁^m\right)^{}^{}S_𝐁^+^{}S_𝐁^{}$$
where $`_𝐁^{n+m}=\pi ^{}T𝐕`$ is the bull–back to the N–adapted frame bundle from the tangent bundle $`T𝐕`$ with $`^n`$ and $`^m`$ being the fundamental representations, respectively, of $`SO(n)`$ and $`SO(m)`$ defining the d–group $`\mathrm{𝐒𝐎}(n+m).`$ Any d–connection on $`𝐕`$ defines a canonical d–connection inducing a standard d–connection on the bundle of N–adapted frames $`𝐁.`$
###### Definition 5.4
The N–anholonomic (twisted) Dirac operator is defined:
$$𝔻^+:C^{\mathrm{}}(𝐁,S_𝐁^+)C^{\mathrm{}}(𝐁,S_𝐁^{}),$$
for $`𝚿`$–modules and
$$\stackrel{}{𝔻}:C^{\mathrm{}}(𝐁,S_𝐁^+𝐖)C^{\mathrm{}}(𝐁,S_𝐁^{}𝐖),$$
for $`𝐖`$ being a $`𝚿^1`$–module with induced canonical d–connection.
The introduced d–operators split into N–adapted components, $`𝔻^+=\left(h𝔻^+,^{}𝔻^+\right)`$ and $`\stackrel{}{𝔻}=\left(h\stackrel{}{𝔻},^{}\stackrel{}{𝔻}\right).`$ These operators are correspondingly $`Spin`$ $`(n)`$– and $`Spin(m)`$–invariant. The space $`S_𝐁^\pm 𝐖`$ is a N–anholonomic $`𝚿^0`$–module descending to bundles $`𝐄^\pm `$ on $`𝐕`$ which transforms the Dirac d–operator to be a twisted Dirac d–operator:
$$\stackrel{}{𝔻}^+:C^{\mathrm{}}(𝐁,𝐄^+)C^{\mathrm{}}(𝐁,𝐄^{}).$$
(20)
If $`𝐕`$ is a spin manifold, than the operators from Definition 5.4 descend to $`𝐕`$ and with a local decomposition $`𝐄_𝐕=𝐒_𝐕𝐖_𝐕.`$
Let us consider an N–anholonomic Clifford module $`𝐄`$ for $`𝐂(𝐕)`$ for which there is a d–connection $`{}_{}{}^{A}𝐃`$ induced by the canonical d–connection $`\widehat{𝐃}`$ and acting following the rule
$${}_{}{}^{A}𝐃[𝐜(\omega )f]=𝐜(\widehat{𝐃}\omega )\lambda +𝐜(\omega )^A𝐃\lambda $$
for any 1–form $`\omega `$ on $`𝐕`$ and $`\lambda C^{\mathrm{}}(𝐕,𝐄).`$ This action preserves the global decomposition (2). We may associate to $`{}_{}{}^{A}𝐃`$ a Dirac d–operator $`\stackrel{}{𝔻}`$ by using the sequence
$$C^{\mathrm{}}(𝐕,𝐄)\stackrel{{}_{}{}^{A}𝐃}{}C^{\mathrm{}}(𝐕,T^{}𝐕𝐄)\stackrel{𝐜}{}C^{\mathrm{}}(𝐕,𝐄).$$
This sequence splits also in h– and v–components. Because the Clifford multiplication by $`T^{}𝐕`$ results in two distinguished odd parts of $`C(𝐕),`$ we get an odd operator $`\stackrel{}{𝔻}`$ splitting in $`\pm `$ components,
$$\stackrel{}{𝔻}^\pm :C^{\mathrm{}}(𝐕,𝐄^\pm )C^{\mathrm{}}(𝐕,𝐄^{})$$
acting in distinguished form on h– and v–components,
$$\stackrel{}{𝔻}^\pm :C^{\mathrm{}}(V,𝐄^\pm )C^{\mathrm{}}(V,𝐄^{})\text{ and }\stackrel{}{𝔻}^\pm :C^{\mathrm{}}(^{}V,𝐄^\pm )C^{\mathrm{}}(^{}V,𝐄^{}).$$
Such operators are N–adapted and formal adjoint of each other respective h– and v–component of the standard functional $`L^2`$ (not confusing with the Lagrange fundamental function considered in the previous section) defining the inner product on $`C^{\mathrm{}}(𝐕,𝐄).`$
The curvature $`{}_{}{}^{A}𝐑`$ of the d–connection $`{}_{}{}^{A}𝐃`$ is a 2–form with values in
$$End(𝐄)=𝐂(𝐕)End_{𝐂(𝐕)}(𝐄).$$
In general, $`{}_{}{}^{A}𝐑`$ does not commute with the action on $`𝐂(𝐕).`$ For Riemannian manifolds, it was proposed to introduce the twisting curvature $`R_{E/S}=^ARc(R)`$ for any $`c(R)C(M)`$ satisfying the conditions $`[{}_{}{}^{A}R,c(\lambda )]=c(R(\lambda ))`$ and $`[c(R),c(\lambda )]=c(R(\lambda ))`$ for $`R`$ being the Riemannian curvature of $`M`$ and any tangent vector $`\lambda .`$ In a similar form, for N–anholonomic manifolds, we can define the twisting canonical curvature
$$\widehat{𝐑}_{E/S}=^A𝐑c(\widehat{𝐑})$$
induced by (7), see formulas (A.23) from Appendix, computed for the canonical d–connection $`\widehat{𝐃},`$ see Theorem 2.4 and (A.21). With this curvature, we may act as with the Riemannian one following the procedure of defining generalized Dirac operators from .
#### 5.2.3 Main result and concluding remarks
The material of previous section 5.2.2 consists the proof of
###### Theorem 5.2
(Twisted Index formula for N–anholonomic Dirac
operators). If $`𝐕`$ is a compact N–anholonomic manifold, then the Dirac operator $`\stackrel{}{𝔻}^+`$ (20) satisfies the index formula
$$ind(\stackrel{}{𝔻}^+)=<\widehat{A}(𝐕)ch(𝐖),[𝐕]>,$$
where the genus
$$\widehat{A}(𝐕)\left|det\frac{\widehat{𝐑}}{2\mathrm{sinh}(\widehat{𝐑}/2)}\right|^{1/2}$$
is determined by the curvature $`\widehat{𝐑}`$ of the canonical d–connection $`\widehat{𝐃}.`$
This theorem can be stated for certain particular cases of Lagrange, or Finsler, geometries and their spinor formulation, for instance, with the aim to locally anisotropic generalization of the so–called $`C`$–spaces which will present topological characteristics derived from a fundamental Lagrange (or Finsler) function or, in a new fashion, for non–spin $`C`$–gerbes associated to nonholonomic gravitational and spinor interactions. The Main Result of this work can be also applied for topological classification of new types of globalized exact solutions defining nonholonomic gravitational and matter field configurations .
Acknowledgement: The authors are grateful to C. Castro Perelman for useful discussions. S. V. thanks Prof. M. Anastasiei for kind support.
## Appendix A Some Local Formulas from N–Connection Geometry
In this Appendix, we present some component formulas and equations defining the local geometry of N–anholonomic spaces, see details in Refs. .
Locally, a N–connection, see Definition 2.1, is stated by its coefficients $`N_i^a(u),`$
$$𝐍=N_i^a(u)dx^i_a$$
(A.1)
where the local coordinates (in general, abstract ones both for holonomic and nonholonomic variables) are split in the form $`u=(x,y),`$ or $`u^\alpha =(x^i,y^a),`$ where $`i,j,k,\mathrm{}=1,2,\mathrm{},n`$ and $`a,b,c,\mathrm{}=n+1,n+2,\mathrm{},n+m`$ when $`_i=/x^i`$ and $`_a=/y^a.`$ The well known class of linear connections consists on a particular subclass with the coefficients being linear on $`y^a,`$ i.e., $`N_i^a(u)=\mathrm{\Gamma }_{bj}^a(x)y^b.`$
An explicit local calculus allows us to write the N–connection curvature (3) in the form
$$𝛀=\frac{1}{2}\mathrm{\Omega }_{ij}^adx^idx^j_a,$$
with the N–connection curvature coefficients
$$\mathrm{\Omega }_{ij}^a=\delta _{[j}N_{i]}^a=\delta _jN_i^a\delta _iN_j^a=_jN_i^a_iN_j^a+N_i^b_bN_j^aN_j^b_bN_i^a.$$
(A.2)
Any N–connection $`𝐍=N_i^a(u)`$ induces a N–adapted frame (vielbein) structure
$$𝐞_\nu =\left(e_i=_iN_i^a(u)_a,e_a=_a\right),$$
(A.3)
and the dual frame (coframe) structure
$$𝐞^\mu =\left(e^i=dx^i,e^a=dy^a+N_i^a(u)dx^i\right).$$
(A.4)
The vielbeins (A.4) satisfy the nonholonomy (equivalently, anholonomy) relations
$$[𝐞_\alpha ,𝐞_\beta ]=𝐞_\alpha 𝐞_\beta 𝐞_\beta 𝐞_\alpha =W_{\alpha \beta }^\gamma 𝐞_\gamma $$
(A.5)
with (antisymmetric) nontrivial anholonomy coefficients $`W_{ia}^b=_aN_i^b`$ and $`W_{ji}^a=\mathrm{\Omega }_{ij}^a.`$<sup>4</sup><sup>4</sup>4One preserves a relation to our previous denotations if we consider that $`𝐞_\nu =(e_i,e_a)`$ and $`𝐞^\mu =(e^i,e^a)`$ are, respectively, the former $`\delta _\nu =\delta /u^\nu =(\delta _i,_a)`$ and $`\delta ^\mu =\delta u^\mu =(d^i,\delta ^a);`$ we emphasize that operators (A.3) and (A.4) define, correspondingly, the “N–elongated” partial derivatives and differentials which are convenient for calculations on N–anholonomic manifolds. These formulas present a local proof of Proposition 2.1 when
$$𝐞=\{𝐞_\nu \}=(e=\{e_i\},^{}e=\{e_a\})$$
and
$$\stackrel{~}{𝐞}=\{𝐞^\mu \}=(\stackrel{~}{e}=\{e^i\},^{}\stackrel{~}{e}=\{e^a\}).$$
Let us consider metric structure
$$\stackrel{˘}{g}=\underset{¯}{g}_{\alpha \beta }\left(u\right)du^\alpha du^\beta $$
(A.6)
defined with respect to a local coordinate basis $`du^\alpha =(dx^i,dy^a)`$ by coefficients
$$\underset{¯}{g}_{\alpha \beta }=\left[\begin{array}{cc}g_{ij}+N_i^aN_j^bh_{ab}& N_j^eh_{ae}\\ N_i^eh_{be}& h_{ab}\end{array}\right].$$
(A.7)
In general, such a metric (A.7) is generic off–diagonal, i.e. it can not be diagonalized by any coordinate transforms and that $`N_i^a(u)`$ are any general functions. The condition (9), for $`Xe_i`$ and $`{}_{}{}^{}Y^{}e_a,`$ transform into
$$\stackrel{˘}{g}(e_i,^{}e_a)=0,\text{ equivalently }\underset{¯}{g}_{ia}N_i^bh_{ab}=0$$
where $`\underset{¯}{g}_{ia}`$ $`g(/x^i,/y^a),`$ which allows us to define in a unique form the coefficients $`N_i^b=h^{ab}\underset{¯}{g}_{ia}`$ where $`h^{ab}`$ is inverse to $`h_{ab}.`$ We can write the metric $`\stackrel{˘}{g}`$ with ansatz (A.7) in equivalent form, as a d–metric adapted to a N–connection structure, see Definition 2.7,
$$𝐠=𝐠_{\alpha \beta }\left(u\right)𝐞^\alpha 𝐞^\beta =g_{ij}\left(u\right)e^ie^j+h_{ab}\left(u\right)^{}e^a^{}e^b,$$
(A.8)
where $`g_{ij}𝐠(e_i,e_j)`$ and $`h_{ab}𝐠\left({}_{}{}^{}e_{a}^{},^{}e_b\right)`$ and the vielbeins $`𝐞_\alpha `$ and $`𝐞^\alpha `$ are respectively of type (A.3) and (A.4).
We can say that the metric $`\stackrel{˘}{g}`$ (A.6) is equivalently transformed into (A.8) by performing a frame (vielbein) transform
$$𝐞_\alpha =𝐞_\alpha ^{\underset{¯}{\alpha }}_{\underset{¯}{\alpha }}\text{ and }𝐞_{}^\beta =𝐞_{\underset{¯}{\beta }}^\beta du^{\underset{¯}{\beta }}.$$
with coefficients
$`𝐞_\alpha ^{\underset{¯}{\alpha }}(u)`$ $`=`$ $`\left[\begin{array}{cc}e_i^{\underset{¯}{i}}(u)& N_i^b(u)e_b^{\underset{¯}{a}}(u)\\ 0& e_a^{\underset{¯}{a}}(u)\end{array}\right],`$ (A.11)
$`𝐞_{\underset{¯}{\beta }}^\beta (u)`$ $`=`$ $`\left[\begin{array}{cc}e_{\underset{¯}{i}}^i(u)& N_k^b(u)e_{\underset{¯}{i}}^k(u)\\ 0& e_{\underset{¯}{a}}^a(u)\end{array}\right],`$ (A.14)
being linear on $`N_i^a.`$ We can consider that a N–anholonomic manifold $`𝐕`$ provided with metric structure $`\stackrel{˘}{g}`$ (A.6) (equivalently, with d–metric (A.8)) is a special type of a manifold provided with a global splitting into conventional “horizontal” and “vertical” subspaces (2) induced by the “off–diagonal” terms $`N_i^b(u)`$ and a prescribed type of nonholonomic frame structure (A.5).
A d–connection, see Definition 2.3, splits into h– and v–covariant derivatives, $`𝐃=D+^{}D,`$ where $`D_k=(L_{jk}^i,L_{bk}^a)`$ and $`{}_{}{}^{}D_{c}^{}=(C_{jk}^i,C_{bc}^a)`$ are correspondingly introduced as h- and v–parametrizations of (A.15),
$$L_{jk}^i=\left(𝐃_ke_j\right)e^i,L_{bk}^a=\left(𝐃_ke_b\right)e^a,C_{jc}^i=\left(𝐃_ce_j\right)e^i,C_{bc}^a=\left(𝐃_ce_b\right)e^a.$$
The components $`𝚪_{\alpha \beta }^\gamma =(L_{jk}^i,L_{bk}^a,C_{jc}^i,C_{bc}^a)`$ completely define a d–connection $`𝐃`$ on a N–anholonomic manifold $`𝐕.`$
The simplest way to perform a local covariant calculus by applying d–connections is to use N–adapted differential forms like $`𝚪_\beta ^\alpha =𝚪_{\beta \gamma }^\alpha 𝐞^\gamma `$ with the coefficients defined with respect to (A.4) and (A.3).One can introduce the d–connection 1–form
$$𝚪_\beta ^\alpha =𝚪_{\beta \gamma }^\alpha 𝐞^\gamma ,$$
when the N–adapted components of d-connection $`𝐃_\alpha =(𝐞_\alpha 𝐃)`$ are computed following formulas
$$𝚪_{\alpha \beta }^\gamma \left(u\right)=\left(𝐃_\alpha 𝐞_\beta \right)𝐞^\gamma ,$$
(A.15)
where ”$`\mathrm{"}`$ denotes the interior product. This allows us to define in local form the torsion (5) $`𝐓=\{𝒯^\alpha \},`$ where
$$𝒯^\alpha \mathrm{𝐃𝐞}^\alpha =d𝐞^\alpha +\mathrm{\Gamma }_\beta ^\alpha 𝐞^\beta $$
(A.16)
and curvature (7) $`𝐑=\{_\beta ^\alpha \},`$ where
$$_\beta ^\alpha 𝐃𝚪_\beta ^\alpha =d𝚪_\beta ^\alpha \mathrm{\Gamma }_\beta ^\gamma 𝚪_\gamma ^\alpha .$$
(A.17)
The d–torsions components of a d–connection $`𝐃,`$ see Theorem 2.1, are computed
$`T_{jk}^i`$ $`=`$ $`L_{jk}^iL_{kj}^i,T_{ja}^i=T_{aj}^i=C_{ja}^i,T_{ji}^a=\mathrm{\Omega }_{ji}^a,`$
$`T_{bi}^a`$ $`=`$ $`T_{ib}^a={\displaystyle \frac{N_i^a}{y^b}}L_{bi}^a,T_{bc}^a=C_{bc}^aC_{cb}^a.`$ (A.18)
For instance, $`T_{jk}^i`$ and $`T_{bc}^a`$ are respectively the coefficients of the $`h(hh)`$–torsion $`T(X,Y)`$ and $`v(vv)`$–torsion $`^{}T(^{}X,^{}Y).`$
The Levi–Civita linear connection $`=\{^{}𝚪_{\beta \gamma }^\alpha \},`$ with vanishing both torsion and nonmetricity $`\stackrel{˘}{g}=0,`$ is not adapted to the global splitting (2). There is a preferred, canonical d–connection structure,$`\widehat{𝐃},`$ on N–anholonomic manifold $`𝐕`$ constructed only from the metric and N–connection coefficients $`[g_{ij},h_{ab},N_i^a]`$ and satisfying the conditions $`\widehat{𝐃}𝐠=0`$ and $`\widehat{T}_{jk}^i=0`$ and $`\widehat{T}_{bc}^a=0,`$ see Theorem 2.4. By straightforward calculations with respect to the N–adapted bases (A.4) and (A.3), we can verify that the connection
$$\widehat{𝚪}_{\beta \gamma }^\alpha =^{}𝚪_{\beta \gamma }^\alpha +\widehat{𝐏}_{\beta \gamma }^\alpha $$
(A.19)
with the deformation d–tensor <sup>5</sup><sup>5</sup>5$`\widehat{𝐏}_{\beta \gamma }^\alpha `$ is a tensor field of type (1,2). As is well known, the sum of a linear connection and a tensor field of type (1,2) is a new linear connection.
$$\widehat{𝐏}_{\beta \gamma }^\alpha =(P_{jk}^i=0,P_{bk}^a=e_b(N_k^a),P_{jc}^i=\frac{1}{2}g^{ik}\mathrm{\Omega }_{kj}^ah_{ca},P_{bc}^a=0)$$
(A.20)
satisfies the conditions of the mentioned Theorem. It should be noted that, in general, the components $`\widehat{T}_{ja}^i,\widehat{T}_{ji}^a`$ and $`\widehat{T}_{bi}^a`$ are not zero. This is an anholonomic frame (or, equivalently, off–diagonal metric) effect. With respect to the N–adapted frames, the coefficients
$`\widehat{𝚪}_{\alpha \beta }^\gamma =(\widehat{L}_{jk}^i,\widehat{L}_{bk}^a,\widehat{C}_{jc}^i,\widehat{C}_{bc}^a)`$ are computed:
$`\widehat{L}_{jk}^i`$ $`=`$ $`{\displaystyle \frac{1}{2}}g^{ir}\left(e_kg_{jr}+e_jg_{kr}e_rg_{jk}\right),`$ (A.21)
$`\widehat{L}_{bk}^a`$ $`=`$ $`e_b(N_k^a)+{\displaystyle \frac{1}{2}}h^{ac}\left(e_kh_{bc}h_{dc}e_bN_k^dh_{db}e_cN_k^d\right),`$
$`\widehat{C}_{jc}^i`$ $`=`$ $`{\displaystyle \frac{1}{2}}g^{ik}e_cg_{jk},\widehat{C}_{bc}^a={\displaystyle \frac{1}{2}}h^{ad}\left(e_ch_{bd}+e_ch_{cd}e_dh_{bc}\right).`$
In some approaches to Finsler geometry , one uses the so–called Berwald d–connection $`𝐃^B`$ with the coefficients
$${}_{}{}^{B}𝚪_{\alpha \beta }^{\gamma }=({}_{}{}^{B}L_{jk}^{i}=\widehat{L}_{jk}^i,^BL_{bk}^a=e_b(N_k^a),^BC_{jc}^i=0,^BC_{bc}^a=\widehat{C}_{bc}^a).$$
(A.22)
This d–connection minimally extends the Levi–Civita connection (it is just the Levi–Civita connection if the integrability conditions are satisfied, i.e. $`\mathrm{\Omega }_{kj}^a=0,`$ see (A.20)). But, in general, for this d–connection the metricity conditions are not satisfied, for instance $`D_ag_{ij}0`$ and $`D_ih_{ab}0.`$
By a straightforward d–form calculus in (A.17), we can find the N–adapted components $`𝐑_{\beta \gamma \delta }^\alpha `$ of the curvature $`𝐑=\{_\beta ^\alpha \}`$ of a d–connection $`𝐃,`$ i.e. the d–curvatures from Theorem 2.2:
$`R_{hjk}^i`$ $`=`$ $`e_kL_{hj}^ie_jL_{hk}^i+L_{hj}^mL_{mk}^iL_{hk}^mL_{mj}^iC_{ha}^i\mathrm{\Omega }_{kj}^a,`$
$`R_{bjk}^a`$ $`=`$ $`e_kL_{bj}^ae_jL_{bk}^a+L_{bj}^cL_{ck}^aL_{bk}^cL_{cj}^aC_{bc}^a\mathrm{\Omega }_{kj}^c,`$
$`R_{jka}^i`$ $`=`$ $`e_aL_{jk}^iD_kC_{ja}^i+C_{jb}^iT_{ka}^b,`$ (A.23)
$`R_{bka}^c`$ $`=`$ $`e_aL_{bk}^cD_kC_{ba}^c+C_{bd}^cT_{ka}^c,`$
$`R_{jbc}^i`$ $`=`$ $`e_cC_{jb}^ie_bC_{jc}^i+C_{jb}^hC_{hc}^iC_{jc}^hC_{hb}^i,`$
$`R_{bcd}^a`$ $`=`$ $`e_dC_{bc}^ae_cC_{bd}^a+C_{bc}^eC_{ed}^aC_{bd}^eC_{ec}^a.`$
Contracting respectively the components of (A.23), one proves
The Ricci tensor $`𝐑_{\alpha \beta }𝐑_{\alpha \beta \tau }^\tau `$ is characterized by h- v–components, i.e. d–tensors,
$$R_{ij}R_{ijk}^k,R_{ia}R_{ika}^k,R_{ai}R_{aib}^b,R_{ab}R_{abc}^c.$$
(A.24)
It should be noted that this tensor is not symmetric for arbitrary d–connections $`𝐃.`$
The scalar curvature of a d–connection is
$${}_{}{}^{s}𝐑𝐠^{\alpha \beta }𝐑_{\alpha \beta }=g^{ij}R_{ij}+h^{ab}R_{ab},$$
(A.25)
defined by a sum the h– and v–components of (A.24) and d–metric (A.8).
The Einstein tensor is defined and computed in standard form
$$𝐆_{\alpha \beta }=𝐑_{\alpha \beta }\frac{1}{2}𝐠_{\alpha \beta }^s𝐑$$
(A.26)
For a Lagrange geometry, see Definition 3.1, by straightforward component calculations, one can be proved the fundamental results:
1. The Euler–Lagrange equations
$$\frac{d}{d\tau }\left(\frac{L}{y^i}\right)\frac{L}{x^i}=0$$
where $`y^i=\frac{dx^i}{d\tau }`$ for $`x^i(\tau )`$ depending on parameter $`\tau ,`$ are equivalent to the “nonlinear” geodesic equations
$$\frac{d^2x^i}{d\tau ^2}+2G^i(x^k,\frac{dx^j}{d\tau })=0$$
defining paths of the canonical semispray
$$S=y^i\frac{}{x^i}2G^i(x,y)\frac{}{y^i}$$
where
$$2G^i(x,y)=\frac{1}{2}^Lg^{ij}\left(\frac{^2L}{y^ix^k}y^k\frac{L}{x^i}\right)$$
with $`{}_{}{}^{L}g_{}^{ij}`$ being inverse to (11).
2. There exists on $`\stackrel{~}{TM}`$ a canonical N–connection
$${}_{}{}^{L}N_{j}^{i}=\frac{G^i(x,y)}{y^i}$$
(A.27)
defined by the fundamental Lagrange function $`L(x,y),`$ which prescribes nonholonomic frame structures of type (A.3) and (A.4), $`{}_{}{}^{L}𝐞_{\nu }^{}=(e_i,^{}e_k)`$ and $`{}_{}{}^{L}𝐞_{}^{\mu }=(e^i,^{}e^k).`$
3. The canonical N–connection (A.27), defining $`{}_{}{}^{}e_{i}^{},`$ induces naturally an almost complex structure $`𝐅:\chi (\stackrel{~}{TM})\chi (\stackrel{~}{TM}),`$ where $`\chi (\stackrel{~}{TM})`$ denotes the module of vector fields on $`\stackrel{~}{TM},`$
$$𝐅(e_i)=^{}e_i\text{ and }𝐅(^{}e_i)=e_i,$$
when
$$𝐅=^{}e_ie^ie_i^{}e^i$$
(A.28)
satisfies the condition $`𝐅𝐅=𝐈,`$ i. e. $`F_\beta ^\alpha F_\gamma ^\beta =\delta _\gamma ^\alpha ,`$ where $`\delta _\gamma ^\alpha `$ is the Kronecker symbol and “$``$” denotes the interior product.
4. On $`\stackrel{~}{TM},`$ there is a canonical metric structure
$${}_{}{}^{L}𝐠=^Lg_{ij}(x,y)e^ie^j+^Lg_{ij}(x,y)^{}e^i^{}e^j$$
(A.29)
constructed as a Sasaki type lift from $`M.`$
5. There is also a canonical d–connection structure $`{}_{}{}^{L}\widehat{𝚪}_{\alpha \beta }^{\gamma }`$ defined only by the components of $`{}_{}{}^{L}N_{j}^{i}`$ and $`{}_{}{}^{L}g_{ij}^{},`$ i.e. by the coefficients of metric (A.29) which in its turn is induced by a regular Lagrangian. The values $`{}_{}{}^{L}\widehat{𝚪}_{\alpha \beta }^{\gamma }=(^L\widehat{L}_{jk}^i,^L\widehat{C}_{bc}^a)`$ are computed just as similar values from (A.21). We note that on $`\stackrel{~}{TM}`$ there are couples of distinguished sets of h- and v–components. |
warning/0506/math0506299.html | ar5iv | text | # Discrete Lagrangian and Hamiltonian Mechanics on Lie groupoids
## 1. Introduction
During the last decade, much effort has been devoted to construction of geometric integrators for Lagrangian systems using a discrete variational principle (see and references therein). In particular, this effort has been concentrated for the case of discrete Lagrangian functions $`L`$ on the cartesian product $`Q\times Q`$ of a differentiable manifold. This cartesian product plays the role of a “discretized version” of the standard velocity phase space $`TQ`$. Applying a natural discrete variational principle, one obtains a second order recursion operator $`\xi :Q\times QQ\times Q`$ assigning to each input pair $`(x,y)`$ the output pair $`(y,z)`$. When the discrete Lagrangian is an approximation of a continuous Lagrangian function (more appropriately, when the discrete Lagrangian approximates the integral action for $`L`$) we obtain a numerical integrator which inherits some of the geometric properties of the continuous Lagrangian (symplecticity, momentum preservation). Although this type of geometric integrators have been mainly considered for conservative systems, the extension to geometric integrators for more involved situations is relatively easy, since, in some sense, many of the constructions mimic the corresponding ones for the continuous counterpart. In this sense, it has been recently shown how discrete variational mechanics can include forced or dissipative systems, holonomic constraints, explicitely time-dependent systems, frictional contact, nonholonomic constraints, multisymplectic fields theories… All these geometric integrators have demonstrated, in worked examples, an exceptionally good longtime behavior and obviously this research is of great interest for numerical and geometric considerations (see ).
On the other hand, Moser and Veselov consider also discrete Lagrangian systems evolving on a Lie group. All this examples leads to A. Weinstein to study discrete mechanics on Lie groupoids, which is a structure that includes as particular examples the case of cartesian products $`Q\times Q`$ as well as Lie groups.
A Lie groupoid $`G`$ is a natural generalization of the concept of a Lie group, where now not all elements are composable. The product $`g_1g_2`$ of two elements is only defined on the set of composable pairs $`G_2=\{(g,h)G\times G|\beta (g)=\alpha (h)\}`$ where $`\alpha :GM`$ and $`\beta :GM`$ are the source and target maps over a base manifold $`M`$. This concept was introduced in differential geometry by Ch. Ereshmann in the 1950’s. The infinitesimal version of a Lie groupoid $`G`$ is the Lie algebroid $`AGM`$, which is the restriction of the vertical bundle of $`\alpha `$ to the submanifold of the identities.
We may thought a Lie algebroid $`A`$ over a manifold $`M`$, with projection $`\tau :AM`$, as a generalized version of the tangent bundle to $`M`$. The geometry and dynamics on Lie algebroids have been extensively studied during the past years. In particular, one of the authors of this paper (see ) developed a geometric formalism of mechanics on Lie algebroids similar to Klein’s formalism of the ordinary Lagrangian mechanics and more recently a description of the Hamiltonian dynamics on a Lie algebroid was given in (see also ).
The key concept in this theory is the prolongation, $`𝒫^\tau A`$, of the Lie algebroid over the fibred projection $`\tau `$ (for the Lagrangian formalism) and the prolongation, $`𝒫^\tau ^{}A`$, over the dual fibred projection $`\tau ^{}:A^{}M`$ (for the Hamiltonian formalism). See for more details. Of course, when the Lie algebroid is $`A=TQ`$ we obtain that $`𝒫^\tau A=T(TQ)`$ and $`𝒫^\tau ^{}A=T(T^{}Q)`$, recovering the classical case. An alternative approach, using the linear Poisson structure on $`A^{}`$ and the canonical isomorphism between $`T^{}A`$ and $`T^{}A^{}`$ was discussed in .
Taking as starting point the results by A. Weinstein , we elucidate in this paper the geometry of Lagrangian systems on Lie groupoids and its Hamiltonian counterpart. Weinstein gave a variational derivation of the discrete Euler-Lagrange equations for a Lagrangian $`L:G`$ on a Lie groupoid $`G`$. We show that the appropriate space to develop a geometric formalism for these equations is the Lie algebroid $`𝒫^\tau GV\beta _GV\alpha G`$ (see section 3 for the definition of the Lie algebroid structure). Note that $`𝒫^\tau G`$ is the total space of the prolongation of the Lie groupoid $`G`$ over the vector bundle projection $`\tau :AGM`$, and that the Lie algebroid of $`𝒫^\tau G`$ is just the prolongation $`𝒫^\tau (AG)`$ (the space were the continuous Lagrangian Mechanics is developed). Using the Lie algebroid structure of $`𝒫^\tau G`$ we may describe discrete Mechanics on the Lie groupoid $`G`$. In particular,
* We give a variational derivation of the discrete Euler-Lagrange equations:
$$\stackrel{}{X}(g)(L)\stackrel{}{X}(h)(L)=0$$
for every section $`X`$ of $`AG`$, where the right or left arrow denotes the induced right and left-invariant vector field on $`G`$.
* We introduce two Poincaré-Cartan 1-sections $`\mathrm{\Theta }_L^+`$ and $`\mathrm{\Theta }_L^{}`$, and an unique Poincaré-Cartan 2-section, $`\mathrm{\Omega }_L`$, on the Lie algebroid $`P^\tau GG`$.
* We study the discrete Lagrangian evolution operator $`\xi :GG`$ and its preservation properties. In particular, we prove that $`(𝒫^\tau \xi ,\xi )^{}\mathrm{\Omega }_L=\mathrm{\Omega }_L`$, where $`𝒫^\tau \xi `$ is the natural prolongation of $`\xi `$ to $`𝒫^\tau G`$.
* Reduction theory is stablished in terms of morphisms of Lie groupoids.
* The associated Hamiltonian formalism is developed using the discrete Legendre transformations $`𝔽^+L:GA^{}G`$ and $`𝔽^{}L:GA^{}G`$.
* A complete characterization of the regularity of a Lagrangian on a Lie groupoid is given in terms of the symplecticity of $`\mathrm{\Omega }_L`$ or, alternatively, in terms of the regularity of the discrete Legendre transformations. In particular, Theorem 4.13 solves the question posed by Weinstein about the regularity conditions for a discrete Lagrangian function on more general Lie groupoids than the cartesian product $`Q\times Q`$. In the regular case, we define the Hamiltonian evolution operator and we prove that it defines a symplectic map.
* We prove a Noether’s theorem for discrete Mechanics on Lie groupoids.
* Finally, some illustrative examples are shown, for instance, discrete Mechanics on the cartesian product $`Q\times Q`$, on Lie groups (discrete Lie-Poisson equations), on action Lie groupoids (discrete Euler-Poincaré equations) and on gauge or Atiyah Lie groupoids (discrete Lagrange-Poincaré equations).
We expect that the results of this paper could be relevant in the construction of new geometric integrators, in particular, for the numerical integration of dynamical systems with symmetry.
The paper is structured as follows. In Section 2 we review some basic results on Lie algebroids and Lie groupoids. Section 3 is devoted to study the Lie algebroid structure of the vector bundle $`𝒫^\tau GV\beta _GV\alpha G`$. The main results of the paper appear in Section 4, where the geometric structure of discrete Mechanics on Lie groupoids is given. Finally, in Section 5, we study several examples of the theory.
## 2. Lie algebroids and Lie groupoids
### 2.1. Lie algebroids
A Lie algebroid $`A`$ over a manifold $`M`$ is a real vector bundle $`\tau :AM`$ together with a Lie bracket $`[[,]]`$ on the space $`\mathrm{\Gamma }(\tau )`$ of the global cross sections of $`\tau :AM`$ and a bundle map $`\rho :ATM`$, called the anchor map, such that if we also denote by $`\rho :\mathrm{\Gamma }(\tau )𝔛(M)`$ the homomorphism of $`C^{\mathrm{}}(M)`$-modules induced by the anchor map then
$$[[X,fY]]=f[[X,Y]]+\rho (X)(f)Y,$$
(2.1)
for $`X,Y\mathrm{\Gamma }(\tau )`$ and $`fC^{\mathrm{}}(M)`$ (see ).
If $`X,Y,Z\mathrm{\Gamma }(\tau )`$ and $`fC^{\mathrm{}}(M)`$ then, using (2.1) and the fact that $`[[,]]`$ is a Lie bracket, we obtain that
$$[[[[X,Y]],fZ]]=f([[X,[[Y,Z]]]][[Y,[[X,Z]]]])+[\rho (X),\rho (Y)](f)Z.$$
(2.2)
On the other hand, from (2.1), it follows that
$$[[[[X,Y]],fZ]]=f[[[[X,Y]],Z]]+\rho [[X,Y]](f)Z.$$
(2.3)
Thus, using (2.2), (2.3) and the fact that $`[[,]]`$ is a Lie bracket, we conclude that
$$\rho [[X,Y]]=[\rho (X),\rho (Y)],$$
that is, $`\rho :\mathrm{\Gamma }(\tau )𝔛(M)`$ is a homomorphism between the Lie algebras $`(\mathrm{\Gamma }(\tau ),[[,]])`$ and $`(𝔛(M),[,])`$.
If $`(A,[[,]],\rho )`$ is a Lie algebroid over $`M`$, one may define the differential of $`\mathrm{A}`$, $`d:\mathrm{\Gamma }(^k\tau ^{})\mathrm{\Gamma }(^{k+1}\tau ^{})`$, as follows
$`d\mu (X_0,\mathrm{},X_k)`$ $`={\displaystyle \underset{i=0}{\overset{k}{}}}(1)^i\rho (X_i)(\mu (X_0,\mathrm{},\widehat{X_i},\mathrm{},X_k))`$ (2.4)
$`+{\displaystyle \underset{i<j}{}}(1)^{i+j}\mu ([[X_i,X_j]],X_0,\mathrm{},\widehat{X_i},\mathrm{},\widehat{X_j},\mathrm{},X_k),`$
for $`\mu \mathrm{\Gamma }(^k\tau ^{})`$ and $`X_0,\mathrm{},X_k\mathrm{\Gamma }(\tau ).`$ $`d`$ is a cohomology operator, that is, $`d^2=0`$. In particular, if $`f:M`$ is a real smooth function then $`df(X)=\rho (X)f,`$ for $`X\mathrm{\Gamma }(\tau )`$. We may also define the Lie derivative with respect to a section $`X`$ of $`A`$ as the operator $`_X:\mathrm{\Gamma }(\mathrm{\Lambda }^kA^{})\mathrm{\Gamma }(\mathrm{\Lambda }^kA^{})`$ given by $`_X=i_Xd+di_X`$ (for more details, see ).
Trivial examples of Lie algebroids are a real Lie algebra $`𝔤`$ of finite dimension (in this case, the base space is a single point) and the tangent bundle $`TM`$ of a manifold $`M`$. Other examples of Lie algebroids are: i) the vertical bundle $`(\tau _P)_{|V\pi }:V\pi P`$ of a fibration $`\pi :PM`$ (and, in general, the tangent vectors to a foliation of finite dimension on a manifold $`P`$); ii) the Atiyah algebroid associated with a principal $`\mathrm{G}`$-bundle (see ); iii) the prolongation $`𝒫^\mathrm{\pi }\mathrm{A}`$ of a Lie algebroid $`\mathrm{A}`$ over a fibration $`\mathrm{\pi }:\mathrm{P}\mathrm{M}`$ (see ) and iv) the action Lie algebroid $`\mathrm{A}\mathrm{f}`$ over a map $`f:M^{}M`$ (see ).
Now, let $`(A,[[,]],\rho )`$ (resp., $`(A^{},[[,]]^{},\rho ^{})`$) be a Lie algebroid over a manifold $`M`$ (resp., $`M^{}`$) and suppose that $`\mathrm{\Psi }:AA^{}`$ is a vector bundle morphism over the map $`\mathrm{\Psi }_0:MM^{}`$. Then, the pair $`(\mathrm{\Psi },\mathrm{\Psi }_0)`$ is said to be a Lie algebroid morphism if
$$d((\mathrm{\Psi },\mathrm{\Psi }_0)^{}\varphi ^{})=(\mathrm{\Psi },\mathrm{\Psi }_0)^{}(d^{}\varphi ^{}),\text{ for all }\varphi ^{}\mathrm{\Gamma }(^k(A^{})^{})\text{ and for all }k,$$
(2.5)
where $`d`$ (resp., $`d^{}`$) is the differential of the Lie algebroid $`A`$ (resp., $`A^{}`$) (see ). In the particular case when $`M=M^{}`$ and $`\mathrm{\Psi }_0=Id`$ then (2.5) holds if and only if
$$[[\mathrm{\Psi }X,\mathrm{\Psi }Y]]^{}=\mathrm{\Psi }[[X,Y]],\rho ^{}(\mathrm{\Psi }X)=\rho (X),\text{ for }X,Y\mathrm{\Gamma }(\tau ).$$
### 2.2. Lie groupoids
In this Section, we will recall the definition of a Lie groupoid and some generalities about them are explained (for more details, see ).
A groupoid over a set $`M`$ is a set $`G`$ together with the following structural maps:
* A pair of maps $`\alpha :GM`$, the source, and $`\beta :GM`$, the target. Thus, an element $`gG`$ is thought as an arrow from $`x=\alpha (g)`$ to $`y=\beta (g)`$ in $`M`$
The maps $`\alpha `$ and $`\beta `$ define the set of composable pairs
$$G_2=\{(g,h)G\times G/\beta (g)=\alpha (h)\}.$$
* A multiplication $`m:G_2G`$, to be denoted simply by $`m(g,h)=gh`$, such that
+ $`\alpha (gh)=\alpha (g)`$ and $`\beta (gh)=\beta (h)`$.
+ $`g(hk)=(gh)k`$.
If $`g`$ is an arrow from $`x=\alpha (g)`$ to $`y=\beta (g)=\alpha (h)`$ and $`h`$ is an arrow from $`y`$ to $`z=\beta (h)`$ then $`gh`$ is the composite arrow from $`x`$ to $`z`$
* An identity section $`ϵ:MG`$ such that
+ $`ϵ(\alpha (g))g=g`$ and $`gϵ(\beta (g))=g`$.
* An inversion map $`i:GG`$, to be denoted simply by $`i(g)=g^1`$, such that
+ $`g^1g=ϵ(\beta (g))`$ and $`gg^1=ϵ(\alpha (g))`$.
A groupoid $`G`$ over a set $`M`$ will be denoted simply by the symbol $`GM`$.
The groupoid $`GM`$ is said to be a Lie groupoid if $`G`$ and $`M`$ are manifolds and all the structural maps are differentiable with $`\alpha `$ and $`\beta `$ differentiable submersions. If $`GM`$ is a Lie groupoid then $`m`$ is a submersion, $`ϵ`$ is an immersion and $`i`$ is a diffeomorphism. Moreover, if $`xM`$, $`\alpha ^1(x)`$ (resp., $`\beta ^1(x)`$) will be said the $`\alpha `$-fiber (resp., the $`\beta `$-fiber) of $`x`$.
On the other hand, if $`gG`$ then the left-translation by $`\mathrm{g}\mathrm{G}`$ and the right-translation by $`\mathrm{g}`$ are the diffeomorphisms
$$\begin{array}{ccc}l_g:\alpha ^1(\beta (g))\alpha ^1(\alpha (g))\hfill & ;\hfill & hl_g(h)=gh,\hfill \\ r_g:\beta ^1(\alpha (g))\beta ^1(\beta (g))\hfill & ;\hfill & hr_g(h)=hg.\hfill \end{array}$$
Note that $`l_g^1=l_{g^1}`$ and $`r_g^1=r_{g^1}`$.
A vector field $`\stackrel{~}{X}`$ on $`G`$ is said to be left-invariant (resp., right-invariant) if it is tangent to the fibers of $`\alpha `$ (resp., $`\beta `$) and $`\stackrel{~}{X}(gh)=(T_hl_g)(\stackrel{~}{X}_h)`$ (resp., $`\stackrel{~}{X}(gh)=(T_gr_h)(\stackrel{~}{X}(g)))`$, for $`(g,h)G_2`$.
Now, we will recall the definition of the Lie algebroid associated with $`\mathrm{G}`$.
We consider the vector bundle $`\tau :AGM`$, whose fiber at a point $`xM`$ is $`A_xG=V_{ϵ(x)}\alpha =Ker(T_{ϵ(x)}\alpha )`$. It is easy to prove that there exists a bijection between the space $`\mathrm{\Gamma }(\tau )`$ and the set of left-invariant (resp., right-invariant) vector fields on $`G`$. If $`X`$ is a section of $`\tau :AGM`$, the corresponding left-invariant (resp., right-invariant) vector field on $`G`$ will be denoted $`\stackrel{}{X}`$ (resp., $`\stackrel{}{X}`$), where
$$\stackrel{}{X}(g)=(T_{ϵ(\beta (g))}l_g)(X(\beta (g))),$$
(2.6)
$$\stackrel{}{X}(g)=(T_{ϵ(\alpha (g))}r_g)((T_{ϵ(\alpha (g))}i)(X(\alpha (g)))),$$
(2.7)
for $`gG`$. Using the above facts, we may introduce a Lie algebroid structure $`([[,]],\rho )`$ on $`AG`$, which is defined by
$$\stackrel{}{[[X,Y]]}=[\stackrel{}{X},\stackrel{}{Y}],\rho (X)(x)=(T_{ϵ(x)}\beta )(X(x)),$$
(2.8)
for $`X,Y\mathrm{\Gamma }(\tau )`$ and $`xM`$. Note that
$$\stackrel{}{[[X,Y]]}=[\stackrel{}{X},\stackrel{}{Y}],[\stackrel{}{X},\stackrel{}{Y}]=0,$$
(2.9)
$$Ti\stackrel{}{X}=\stackrel{}{X}i,Ti\stackrel{}{X}=\stackrel{}{X}i,$$
(2.10)
(for more details, see ).
Given two Lie groupoids $`GM`$ and $`G^{}M^{}`$, a morphism of Lie groupoids is a smooth map $`\mathrm{\Phi }:GG^{}`$ such that
$$(g,h)G_2(\mathrm{\Phi }(g),\mathrm{\Phi }(h))(G^{})_2$$
and
$$\mathrm{\Phi }(gh)=\mathrm{\Phi }(g)\mathrm{\Phi }(h).$$
A morphism of Lie groupoids $`\mathrm{\Phi }:GG^{}`$ induces a smooth map $`\mathrm{\Phi }_0:MM^{}`$ in such a way that
$$\alpha ^{}\mathrm{\Phi }=\mathrm{\Phi }_0\alpha ,\beta ^{}\mathrm{\Phi }=\mathrm{\Phi }_0\beta ,\mathrm{\Phi }ϵ=ϵ^{}\mathrm{\Phi }_0,$$
$`\alpha `$, $`\beta `$ and $`ϵ`$ (resp., $`\alpha ^{}`$, $`\beta ^{}`$ and $`ϵ^{}`$) being the source, the target and the identity section of $`G`$ (resp., $`G^{}`$).
Suppose that $`(\mathrm{\Phi },\mathrm{\Phi }_0)`$ is a morphism between the Lie groupoids $`GM`$ and $`G^{}M^{}`$ and that $`\tau :AGM`$ (resp., $`\tau ^{}:AG^{}M^{}`$) is the Lie algebroid of $`G`$ (resp., $`G^{}`$). Then, if $`xM`$ we may consider the linear map $`A_x(\mathrm{\Phi }):A_xGA_{\mathrm{\Phi }_0(x)}G^{}`$ defined by
$$A_x(\mathrm{\Phi })(v_{ϵ(x)})=(T_{ϵ(x)}\mathrm{\Phi })(v_{ϵ(x)}),\text{ for }v_{ϵ(x)}A_xG.$$
(2.11)
In fact, we have that the pair $`(A(\mathrm{\Phi }),\mathrm{\Phi }_0)`$ is a morphism between the Lie algebroids $`\tau :AGM`$ and $`\tau ^{}:AG^{}M^{}`$ (see ).
Next, we will present some examples of Lie groupoids.
1.- Lie groups. Any Lie group $`G`$ is a Lie groupoid over $`\{𝔢\}`$, the identity element of $`G`$. The Lie algebroid associated with $`G`$ is just the Lie algebra $`𝔤`$ of $`G`$.
2.- The pair or banal groupoid. Let $`M`$ be a manifold. The product manifold $`M\times M`$ is a Lie groupoid over $`M`$ in the following way: $`\alpha `$ is the projection onto the first factor and $`\beta `$ is the projection onto the second factor; $`ϵ(x)=(x,x)`$, for all $`xM`$, $`m((x,y),(y,z))=(x,z)`$, for $`(x,y),(y,z)M\times M`$ and $`i(x,y)=(y,x)`$. $`M\times MM`$ is called the pair or banal groupoid. If $`x`$ is a point of $`M`$, it follows that
$$V_{ϵ(x)}\alpha =\{0_x\}\times T_xMT_xM\times T_xMT_{(x,x)}(M\times M).$$
Thus, the linear maps
$$\mathrm{\Psi }_x:T_xMV_{ϵ(x)}\alpha ,v_x(0_x,v_x),$$
induce an isomorphism (over the identity of $`M`$) between the Lie algebroids $`\tau _M:TMM`$ and $`\tau :A(M\times M)M.`$
3.- The Lie groupoid associated with a fibration. Let $`\pi :PM`$ be a fibration, that is, $`\pi `$ is a surjective submersion and denote by $`G_\pi `$ the subset of $`P\times P`$ given by
$$G_\pi =\{(p,p^{})P\times P/\pi (p)=\pi (p^{})\}.$$
Then, $`G_\pi `$ is a Lie groupoid over $`P`$ and the structural maps $`\alpha _\pi `$, $`\beta _\pi `$, $`m_\pi `$, $`ϵ_\pi `$ and $`i_\pi `$ are the restrictions to $`G_\pi `$ of the structural maps of the pair groupoid $`P\times PP`$.
If $`p`$ is a point of $`P`$ it follows that
$$V_{ϵ_\pi (p)}\alpha _\pi =\{(0_p,Y_p)T_pP\times T_pP/(T_p\pi )(Y_p)=0\}.$$
Thus, if $`(\tau _P)_{|V\pi }:V\pi P`$ is the vertical bundle to $`\pi `$ then the linear maps
$$(\mathrm{\Psi }_\pi )_p:V_p\pi V_{ϵ_\pi (p)}\alpha _\pi ,Y_p(0_p,Y_p)$$
induce an isomorphism (over the identity of $`M`$) between the Lie algebroids $`(\tau _P)_{|V\pi }:V\pi P`$ and $`\tau :AG_\pi P`$.
4.- Atiyah or gauge groupoids. Let $`p:QM`$ be a principal $`G`$-bundle. Then, the free action, $`\mathrm{\Phi }:G\times QQ`$, $`(g,q)\mathrm{\Phi }(g,q)=gq`$, of $`G`$ on $`Q`$ induces, in a natural way, a free action $`\mathrm{\Phi }\times \mathrm{\Phi }:G\times (Q\times Q)Q\times Q`$ of $`G`$ on $`Q\times Q`$ given by $`(\mathrm{\Phi }\times \mathrm{\Phi })(g,(q,q^{}))=(gq,gq^{})`$, for $`gG`$ and $`(q,q^{})Q\times Q`$. Moreover, one may consider the quotient manifold $`(Q\times Q)/G`$ and it admits a Lie groupoid structure over $`M`$ with structural maps given by
$$\begin{array}{ccc}\stackrel{~}{\alpha }:(Q\times Q)/GM\hfill & ;& [(q,q^{})]p(q),\hfill \\ \stackrel{~}{\beta }:(Q\times Q)/GM\hfill & ;& [(q,q^{})]p(q^{}),\hfill \\ \stackrel{~}{ϵ}:M(Q\times Q)/G\hfill & ;& x[(q,q)],\text{ if }p(q)=x,\hfill \\ \stackrel{~}{m}:((Q\times Q)/G)_2(Q\times Q)/G\hfill & ;& ([(q,q^{})],[(gq^{},q^{\prime \prime })])[(gq,q^{\prime \prime })],\hfill \\ \stackrel{~}{i}:(Q\times Q)/G(Q\times Q)/G\hfill & ;& [(q,q^{})][(q^{},q)].\hfill \end{array}$$
This Lie groupoid is called the Atiyah (gauge) groupoid associated with the principal $`\mathrm{G}`$-bundle $`\mathrm{p}:\mathrm{Q}\mathrm{M}`$ (see ).
If $`x`$ is a point of $`M`$ such that $`p(q)=x`$, with $`qQ`$, and $`p_{Q\times Q}:Q\times Q(Q\times Q)/G`$ is the canonical projection then it is clear that
$$V_{\stackrel{~}{ϵ}(x)}\stackrel{~}{\alpha }=(T_{(q,q)}p_{Q\times Q})(\{0_q\}\times T_qQ).$$
Thus, if $`\tau _Q|G:TQ/GM`$ is the Atiyah algebroid associated with the principal $`G`$-bundle $`p:GM`$ then the linear maps
$$(TQ/G)_xV_{\stackrel{~}{ϵ}(x)}\stackrel{~}{\alpha };[v_q](T_{(q,q)}p_{Q\times Q})(0_q,v_q),\text{ with }v_qT_qQ,$$
induce an isomorphism (over the identity of $`M`$) between the Lie algebroids $`\tau :A((Q\times Q)/G)M`$ and $`\tau _Q|G:TQ/GM`$.
5.- The prolongation of a Lie groupoid over a fibration. Given a Lie groupoid $`GM`$ and a fibration $`\pi :PM`$, we consider the set
$$𝒫^\pi G=P\text{ }_\pi \times _\alpha G\text{ }_\beta \times _\pi P=\{(p,g,p^{})P\times G\times P/\pi (p)=\alpha (g),\beta (g)=\pi (p^{})\}.$$
Then, $`𝒫^\pi G`$ is a Lie groupoid over $`P`$ with structural maps given by
$$\begin{array}{ccc}\alpha ^\pi :𝒫^\pi GP\hfill & ;& (p,g,p^{})p,\hfill \\ \beta ^\pi :𝒫^\pi GP\hfill & ;& (p,g,p^{})p^{},\hfill \\ ϵ^\pi :P𝒫^\pi G\hfill & ;& p(p,ϵ(\pi (p)),p),\hfill \\ m^\pi :(𝒫^\pi G)_2𝒫^\pi G\hfill & ;& ((p,g,p^{}),(p^{},h,p^{\prime \prime }))(p,gh,p^{\prime \prime }),\hfill \\ i^\pi :𝒫^\pi G𝒫^\pi G\hfill & ;& (p,g,p^{})(p^{},g^1,p).\hfill \end{array}$$
$`𝒫^\pi G`$ is called the prolongation of $`\mathrm{G}`$ over $`\mathrm{\pi }:\mathrm{P}\mathrm{M}`$.
Now, denote by $`\tau :AGM`$ the Lie algebroid of $`G`$, by $`A(𝒫^\pi G)`$ the Lie algebroid of $`𝒫^\pi G`$ and by $`𝒫^\pi (AG)`$ the prolongation of $`\tau :AGM`$ over the fibration $`\pi `$. If $`pP`$ and $`m=\pi (p)`$, then it follows that
$$A_p(𝒫^\pi G)=\{(0_p,v_{ϵ(m)},X_p)T_pP\times A_mG\times T_pP/(T_p\pi )(X_p)=(T_{ϵ(m)}\beta )(v_{ϵ(m)})\}$$
and, thus, one may consider the linear isomorphism
$$(\mathrm{\Psi }^\pi )_p:A_p(𝒫^\pi G)𝒫_p^\pi (AG),(0_p,v_{ϵ(m)},X_p)(v_{ϵ(m)},X_p).$$
(2.12)
In addition, one may prove that the maps $`(\mathrm{\Psi }^\pi )_p`$, $`pP`$, induce an isomorphism $`\mathrm{\Psi }^\pi :A(𝒫^\pi G)𝒫^\pi (AG)`$ between the Lie algebroids $`A(𝒫^\pi G)`$ and $`𝒫^\pi (AG)`$ (for more details, see ).
6.- Action Lie groupoids. Let $`GM`$ be a Lie groupoid and $`\pi :PM`$ be a smooth map. If $`P\text{ }_\pi \times _\alpha G=\{(p,g)P\times G/\pi (p)=\alpha (g)\}`$ then a right action of $`G`$ on $`\pi `$ is a smooth map
$$P\text{ }_\pi \times _\alpha GP,(p,g)pg,$$
which satisfies the following relations
$$\begin{array}{cccc}\hfill \pi (pg)& =& \beta (g),\hfill & \text{ for }(p,g)P\text{ }_\pi \times _\alpha G,\hfill \\ \hfill (pg)h& =& p(gh),\hfill & \text{ for }(p,g)P\text{ }_\pi \times _\alpha G\text{ and }(g,h)G_2,\text{ and }\hfill \\ \hfill pϵ(\pi (p))& =& p,\hfill & \text{ for }pP.\hfill \end{array}$$
Given such an action one constructs the action Lie groupoid $`P\text{ }_\pi \times _\alpha G`$ over $`P`$ by defining
$$\begin{array}{ccc}\stackrel{~}{\alpha }_\pi :P\text{ }_\pi \times _\alpha GP\hfill & ;& (p,g)p,\hfill \\ \stackrel{~}{\beta }_\pi :P\text{ }_\pi \times _\alpha GP\hfill & ;& (p,g)pg,\hfill \\ \stackrel{~}{ϵ}_\pi :PP\text{ }_\pi \times _\alpha G\hfill & ;& p(p,ϵ(\pi (p))),\hfill \\ \stackrel{~}{m}_\pi :(P\text{ }_\pi \times _\alpha G)_2P\text{ }_\pi \times _\alpha G\hfill & ;& ((p,g),(pg,h))(p,gh),\hfill \\ \stackrel{~}{i}_\pi :P\text{ }_\pi \times _\alpha GP\text{ }_\pi \times _\alpha G\hfill & ;& (p,g)(pg,g^1).\hfill \end{array}$$
Now, if $`pP`$, we consider the map $`p:\alpha ^1(\pi (p))P`$ given by
$$p(g)=pg.$$
Then, if $`\tau :AGM`$ is the Lie algebroid of $`G`$, the $``$-linear map $`\mathrm{\Phi }:\mathrm{\Gamma }(\tau )𝔛(P)`$ defined by
$$\mathrm{\Phi }(X)(p)=(T_{ϵ(\pi (p))}p)(X(\pi (p))),\text{ for }X\mathrm{\Gamma }(\tau )\text{ and }pP,$$
induces an action of $`AG`$ on $`\pi :PM`$. In addition, the Lie algebroid associated with the Lie groupoid $`P\text{ }_\pi \times _\alpha GP`$ is the action Lie algebroid $`AG\pi `$ (for more details, see ).
## 3. Lie algebroid structure on the vector bundle $`\pi ^\tau :𝒫^\tau GG`$
Let $`GM`$ be a Lie groupoid with structural maps
$$\alpha ,\beta :GM,ϵ:MG,i:GG,m:G_2G.$$
Suppose that $`\tau :AGM`$ is the Lie algebroid of $`G`$ and that $`𝒫^\tau G`$ is the prolongation of $`G`$ over the fibration $`\tau :AGM`$ (see Example $`5`$ in Section 2.2), that is,
$$𝒫^\tau G=AG\text{ }_\tau \times _\alpha G\text{ }_\beta \times _\tau AG.$$
$`𝒫^\tau G`$ is a Lie groupoid over $`AG`$ and we may define the bijective map $`\mathrm{\Theta }:𝒫^\tau GV\beta _GV\alpha `$ as follows
$$\mathrm{\Theta }(u_{ϵ(\alpha (g))},g,v_{ϵ(\beta (g))})=((T_{ϵ(\alpha (g))}(r_gi))(u_{ϵ(\alpha (g))}),(T_{ϵ(\beta (g))}l_g)(v_{ϵ(\beta (g))})),$$
for $`(u_{ϵ(\alpha (g))},g,v_{ϵ(\beta (g))})A_{\alpha (g)}G\times G\times A_{\beta (g)}G`$. Thus, $`V\beta _GV\alpha `$ is a Lie groupoid over $`AG`$ (this Lie groupoid was considered by Saunders ). We remark that the Lie algebroid of $`𝒫^\tau GV\beta _GV\alpha AG`$ is isomorphic to the prolongation of $`AG`$ over $`\tau :AGM`$ and that the prolongation of a Lie algebroid $`A`$ over the vector bundle projection $`\tau :AM`$ plays an important role in the description of Lagrangian Mechanics on $`A`$ (see ).
On the other hand, note that $`𝒫^\tau GV\beta _GV\alpha `$ is a real vector bundle over $`G`$. In this section, we will prove that the vector bundle $`\pi ^\tau :𝒫^\tau GG`$ admits an integrable Lie algebroid structure. In other words, we will prove that there exists a Lie groupoid $`HG`$ over $`G`$ such that the Lie algebroid $`AH`$ is isomorphic to the real vector bundle $`\pi ^\tau :𝒫^\tau GG`$. In addition, we will see that the Lie groupoid $`H`$ is isomorphic to the prolongations of $`G`$ over $`\alpha `$ and $`\beta `$.
It is clear that the Lie algebroids of the Lie groupoids over $`G`$
$$G_\beta =\{(g,h)G\times G/\beta (g)=\beta (h)\},G_\alpha =\{(r,s)G\times G/\alpha (r)=\alpha (s)\},$$
are just $`V\beta G`$ and $`V\alpha G`$, respectively. This fact suggests to consider the following manifold
$$G_\beta G_\alpha =\{((g,h),(r,s))G_\beta \times G_\alpha /\beta _\beta (g,h)=\alpha _\alpha (r,s)\},$$
where $`\beta _\beta :G_\beta G`$ (respectively, $`\alpha _\alpha :G_\alpha G`$) is the target (respectively, the source) of the Lie groupoid $`G_\beta G`$ (respectively, $`G_\alpha G`$).
We will identify the space $`G_\beta G_\alpha `$ with
$$\{(g,h,s)G\times G\times G/\beta (g)=\beta (h),\alpha (h)=\alpha (s)\}.$$
This space admits a Lie groupoid structure over $`G`$ with structural maps given by
$$\begin{array}{ccc}\alpha _{\beta \alpha }:G_\beta G_\alpha G\hfill & ;\hfill & (g,h,s)g,\hfill \\ \beta _{\beta \alpha }:G_\beta G_\alpha G\hfill & ;\hfill & (g,h,s)s,\hfill \\ ϵ_{\beta \alpha }:GG_\beta G_\alpha \hfill & ;\hfill & g(g,g,g),\hfill \\ m_{\beta \alpha }:(G_\beta G_\alpha )_2G_\beta G_\alpha \hfill & ;\hfill & ((g,h,s),(s,h^{},s^{}))(g,h^{}s^1h,s^{}),\hfill \\ i_{\beta \alpha }:G_\beta G_\alpha G_\beta G_\alpha \hfill & ;\hfill & (g,h,s)(s,gh^1s,g).\hfill \end{array}$$
(3.1)
Note that
$$\begin{array}{ccc}j_\beta :G_\beta G_\beta G_\alpha \hfill & ;\hfill & (g,h)j_\beta (g,h)=(g,h,h),\hfill \\ j_\alpha :G_\alpha G_\beta G_\alpha \hfill & ;\hfill & (h,s)j_\alpha (h,s)=(h,h,s),\hfill \end{array}$$
are Lie groupoid morphisms and that the map
$$m_{\beta \alpha }(j_\beta ,j_\alpha ):G_\beta G_\alpha G_\beta G_\alpha ;(g,h,s)m_{\beta \alpha }(j_\beta (g,h),j_\alpha (h,s))$$
is just the identity map. This implies that $`(G_\beta ,G_\alpha )`$ is a matched pair of Lie groupoids in the sense of Mackenzie (see also ).
Denote by $`([[,]],\rho )`$ the Lie algebroid structure on $`\tau :AGM`$.
###### Theorem 3.1.
Let $`A(G_\beta G_\alpha )G`$ be the Lie algebroid of the Lie groupoid $`G_\beta G_\alpha G`$. Then:
1. The vector bundles $`A(G_\beta G_\alpha )G`$ and $`\pi ^\tau :𝒫^\tau GV\beta _GV\alpha G`$ are isomorphic. Thus, the vector bundle $`\pi ^\tau :𝒫^\tau GV\beta _GV\alpha G`$ admits a Lie algebroid structure.
2. The anchor map $`\rho ^{𝒫^\tau G}`$ of $`\pi ^\tau :𝒫^\tau GV\beta _GV\alpha G`$ is given by
$$\rho ^{𝒫^\tau G}(X_g,Y_g)=X_g+Y_g,\text{ for }(X_g,Y_g)V_g\beta V_g\alpha ,$$
(3.2)
and the Lie bracket $`[[,]]^{𝒫^\tau G}`$ on the space $`\mathrm{\Gamma }(\pi ^\tau )`$ is characterized by the following relation
$$[[(\stackrel{}{X},\stackrel{}{Y}),(\stackrel{}{X^{}},\stackrel{}{Y^{}})]]^{𝒫^\tau G}=(\stackrel{}{[[X,X^{}]]},\stackrel{}{[[Y,Y^{}]]}),$$
(3.3)
for $`X,Y,X^{},Y^{}\mathrm{\Gamma }(\tau )`$.
###### Proof.
(i) If $`gG`$ then, from (3.1), we deduce that the vector space $`V_{ϵ_{\beta \alpha }(g)}\alpha _{\beta \alpha }`$ may be described as follows
$`V_{ϵ_{\beta \alpha }(g)}\alpha _{\beta \alpha }`$ $`=\{(0_g,X_g,Z_g)T_gG\times T_gG\times T_gG/X_gV_g\beta ,(T_g\alpha )(X_g)=(T_g\alpha )(Z_g)\}`$
$`\{(X_g,Z_g)T_gG\times T_gG/X_gV_g\beta ,(T_g\alpha )(X_g)=(T_g\alpha )(Z_g)\}.`$
Now, we will define the linear map $`\mathrm{\Psi }_g:V_{ϵ_{\beta \alpha }(g)}\alpha _{\beta \alpha }V_g\beta V_g\alpha 𝒫_g^\tau G`$ by
$$\mathrm{\Psi }_g(X_g,Z_g)=(X_g,Z_gX_g).$$
(3.4)
It is clear that $`\mathrm{\Psi }_g`$ is a linear isomorphism and
$$\mathrm{\Psi }_g^1(X_g,Y_g)=(X_g,X_g+Y_g),\text{ for }(X_g,Y_g)V_g\beta V_g\alpha 𝒫_g^\tau G.$$
(3.5)
Therefore, the collection of the maps $`\mathrm{\Psi }_g`$, $`gG`$, induces a vector bundle isomorphism $`\mathrm{\Psi }:A(G_\beta G_\alpha )𝒫^\tau GV\beta _GV\alpha `$ over the identity of $`G`$.
(ii) A direct computation, using (3.1), proves that the linear map $`T_{ϵ_{\beta \alpha }(g)}\beta _{\beta \alpha }:V_{ϵ_{\beta \alpha }(g)}\alpha _{\beta \alpha }T_gG`$ is given by
$$(T_{ϵ_{\beta \alpha }(g)}\beta _{\beta \alpha })(X_g,Z_g)=Z_g.$$
(3.6)
Consequently, from (2.8), (3.5) and (3.6), we deduce that (3.2) holds.
Next, we will prove (3.3).
Using (3.5), it follows that
$$(\mathrm{\Psi }^1(\stackrel{}{X},\stackrel{}{Y}))(g)=(0_g,\stackrel{}{X}(g),\stackrel{}{X}(g)+\stackrel{}{Y}(g))(\stackrel{}{X}(g),\stackrel{}{X}(g)+\stackrel{}{Y}(g)),$$
for $`gG`$. Denote by $`\stackrel{}{\mathrm{\Psi }^1(\stackrel{}{X},\stackrel{}{Y})}`$ the corresponding left-invariant vector field on $`G_\beta G_\alpha `$. Then, from (2.6) and (3.1), we have that
$$\stackrel{}{\mathrm{\Psi }^1(\stackrel{}{X},\stackrel{}{Y})}(g,h,s)=(0_g,\stackrel{}{X}(h),\stackrel{}{X}(s)+\stackrel{}{Y}(s)),\text{ for }(g,h,s)G_\beta G_\alpha .$$
Thus, using (2.8) and (2.9), we conclude that
$$[\stackrel{}{\mathrm{\Psi }^1(\stackrel{}{X},\stackrel{}{Y})},\stackrel{}{\mathrm{\Psi }^1(\stackrel{}{X^{}},\stackrel{}{Y^{}})}]=\stackrel{}{\mathrm{\Psi }^1(\stackrel{}{[[X,X^{}]]},\stackrel{}{[[Y,Y^{}]]})}.$$
Therefore, we obtain that (3.3) holds. ∎
The above diagram shows the Lie groupoid and Lie algebroid structures of $`𝒫^\tau G`$:
Given a section $`X`$ of $`AGM`$, we define the sections $`X^{(1,0)}`$, $`X^{(0,1)}`$ (the $`\beta `$ and $`\alpha `$\- lifts) and $`X^{(1,1)}`$ (the complete lift) of $`X`$ to $`\pi ^\tau :𝒫^\tau GG`$ as follows:
$$X^{(1,0)}(g)=(\stackrel{}{X}(g),0_g),X^{(0,1)}(g)=(0_g,\stackrel{}{X}(g))\text{and}X^{(1,1)}(g)=(\stackrel{}{X}(g),\stackrel{}{X}(g))$$
We can easily see that
$$\begin{array}{c}[[X^{(1,0)},Y^{(1,0)}]]^{𝒫^\tau G}=[[X,Y]]^{(1,0)}\hfill \\ [[X^{(0,1)},Y^{(0,1)}]]^{𝒫^\tau G}=[[X,Y]]^{(0,1)}\hfill \end{array}\text{ and }[[X^{(0,1)},Y^{(1,0)}]]^{𝒫^\tau G}=0$$
(3.7)
and, as a consequence,
$$\begin{array}{c}[[X^{(1,1)},Y^{(1,0)}]]^{𝒫^\tau G}=[[X,Y]]^{(1,0)}\hfill \\ [[X^{(1,1)},Y^{(0,1)}]]^{𝒫^\tau G}=[[X,Y]]^{(0,1)}\hfill \end{array}\text{ and }[[X^{(1,1)},Y^{(1,1)}]]^{𝒫^\tau G}=[[X,Y]]^{(1,1)}.$$
(3.8)
###### Remark 3.2.
From Theorem 3.1, we deduce that the canonical inclusions
$$(Id,0):V\beta 𝒫^\tau GV\beta _GV\alpha ,(0,Id):V\alpha 𝒫^\tau GV\beta _GV\alpha ,$$
are Lie algebroid morphisms over the identity of $`G`$. In other words, $`(V\beta ,V\alpha )`$ is a matched pair of Lie algebroids in the sense of Mokri . This fact directly follows using the following general theorem (see ): if $`(G,H)`$ is a matched pair of Lie groupoids then $`(AG,AH)`$ is a matched pair of Lie algebroids. $``$
Next, we will consider the prolongation $`𝒫^\beta G`$ of the Lie groupoid $`G`$ over the target $`\beta :GM`$. We recall that
$$𝒫^\beta G=G\text{ }_\beta \times _\alpha G\text{ }_\beta \times _\beta G=\{(g,h,s)G\times G\times G/\beta (g)=\alpha (h),\beta (h)=\beta (s)\},$$
and that $`𝒫^\beta G`$ is a Lie groupoid over $`G`$ with structural maps
$$\begin{array}{ccc}\alpha ^\beta :𝒫^\beta GG\hfill & ;\hfill & (g,h,s)g,\hfill \\ \beta ^\beta :𝒫^\beta GG\hfill & ;\hfill & (g,h,s)s,\hfill \\ ϵ^\beta :G𝒫^\beta G\hfill & ;\hfill & g(g,ϵ(\beta (g)),g),\hfill \\ m^\beta :(𝒫^\beta G)_2𝒫^\beta G\hfill & ;\hfill & ((g,h,s),(s,t,u))(g,ht,u),\hfill \\ i^\beta :𝒫^\beta G𝒫^\beta G\hfill & ;\hfill & (g,h,s)(s,h^1,g).\hfill \end{array}$$
(3.9)
Moreover, we also have that the Lie algebroid of $`𝒫^\beta G`$ may be identified with the prolongation $`𝒫^\beta (AG)`$ of $`AG`$ over $`\beta :GM`$. We remark that
$$𝒫_g^\beta (AG)=\{(v_{ϵ(\beta (g))},X_g)A_{\beta (g)}G\times T_gG/(T_{ϵ(\beta (g))}\beta )(v_{ϵ(\beta (g))})=(T_g\beta )(X_g)\}$$
for $`gG`$.
###### Theorem 3.3.
Let $`\mathrm{\Phi }^\beta :G_\beta G_\alpha 𝒫^\beta G`$ be the map defined by
$$\mathrm{\Phi }^\beta (g,h,s)=(g,h^1s,s),$$
(3.10)
for $`(g,h,s)G_\beta G_\alpha `$. Then:
1. $`\mathrm{\Phi }^\beta `$ is a Lie groupoid isomorphism over the identity of $`G`$.
2. If $`A(\mathrm{\Phi }^\beta ):A(G_\beta G_\alpha )A(𝒫^\beta G)`$ is the corresponding Lie algebroid isomorphism then, under the identifications
$$A(G_\beta G_\alpha )𝒫^\tau GV\beta _GV\alpha ,A(𝒫^\beta G)𝒫^\beta (AG),$$
$`A(\mathrm{\Phi }^\beta )`$ is given by
$$A_g(\mathrm{\Phi }^\beta )(X_g,Y_g)=((T_gl_{g^1})(Y_g),X_g+Y_g),$$
(3.11)
for $`(X_g,Y_g)V_g\beta V_g\alpha `$, where $`l_{g^1}:\alpha ^1(\alpha (g))\alpha ^1(\beta (g))`$ is the left-translation by $`g^1`$.
###### Proof.
(i) A direct computation, using (3.1) and (3.9), proves the result.
(ii) If $`gG`$ we have that
$$\begin{array}{ccc}A_g(G_\beta G_\alpha )=V_{ϵ_{\beta \alpha }(g)}\alpha _{\beta \alpha }\hfill & =& \{(0_g,X_g,Z_g)T_gG\times T_gG\times T_gG/X_gV_g\beta ,\hfill \\ & & (T_g\alpha )(X_g)=(T_g\alpha )(Z_g)\},\hfill \\ A_g(𝒫^\beta G)=V_{ϵ^\beta (g)}\alpha ^\beta \hfill & =& \{(0_g,v_{ϵ(\beta (g))},Y_g)T_gG\times A_{\beta (g)}G\times T_gG/\hfill \\ & & (T_g\beta )(Y_g)=(T_{ϵ(\beta (g))}\beta )(v_{ϵ(\beta (g))})\}.\hfill \end{array}$$
Now, if $`(0_g,X_g,Z_g)V_{ϵ_{\beta \alpha }(g)}\alpha _{\beta \alpha }`$ then, from (3.10), we deduce that
$`(T_{ϵ_{\beta \alpha }(g)}\mathrm{\Phi }^\beta )`$ $`(0_g,X_g,Z_g)=(T_{ϵ_{\beta \alpha }(g)}\mathrm{\Phi }^\beta )(0_g,0_g,Z_gX_g)+(T_{ϵ_{\beta \alpha }(g)}\mathrm{\Phi }^\beta )(0_g,X_g,X_g)`$
$`=(0_g,(T_gl_{g^1})(Z_gX_g),Z_gX_g)+(T_{ϵ_{\beta \alpha }(g)}\mathrm{\Phi }^\beta )(0_g,X_g,X_g).`$
On the other hand, suppose that $`\beta (g)=xM`$ and that $`\gamma :(\epsilon ,\epsilon )\beta ^1(x)`$ is a curve in $`\beta ^1(x)`$ such that $`\gamma (0)=g`$ and $`\gamma ^{}(0)=X_g`$. Then, one may consider the curve $`\stackrel{~}{\gamma }:(\epsilon ,\epsilon )G_\beta G_\alpha `$ on $`G_\beta G_\alpha `$ given by
$$\stackrel{~}{\gamma }(t)=(g,\gamma (t),\gamma (t))$$
and it follows that
$$\stackrel{~}{\gamma }(0)=(g,g,g),\stackrel{~}{\gamma }^{}(0)=(0_g,X_g,X_g).$$
Moreover, we obtain that
$$\overline{\gamma }(t)=(\mathrm{\Phi }^\beta \stackrel{~}{\gamma })(t)=(g,ϵ(x),\gamma (t)),\text{ for all }t$$
and thus
$$\overline{\gamma }^{}(0)=(0_g,0_{ϵ(\beta (g))},X_g).$$
This proves that
$$(T_{ϵ_{\beta \alpha }(g)}\mathrm{\Phi }^\beta )(0_g,X_g,Z_g)=(0_g,(T_gl_{g^1})(Z_gX_g),Z_g).$$
(3.12)
Finally, using (2.11), (2.12), (3.5) and (3.12), we deduce that (3.11) holds. ∎
Next, we will consider the prolongation $`𝒫^\alpha G`$ of the Lie groupoid $`G`$ over the source $`\alpha :GM`$. We recall that
$$𝒫^\alpha G=G\text{ }_\alpha \times _\alpha G\text{ }_\beta \times _\alpha G=\{(g,h,s)G\times G\times G/\alpha (g)=\alpha (h),\beta (h)=\alpha (s)\}$$
and that $`𝒫^\alpha G`$ is a Lie groupoid over $`G`$ with structural maps
$$\begin{array}{ccc}\alpha ^\alpha :𝒫^\alpha GG\hfill & ;& (g,h,s)g,\hfill \\ \beta ^\alpha :𝒫^\alpha GG\hfill & ;& (g,h,s)s,\hfill \\ ϵ^\alpha :G𝒫^\alpha G\hfill & ;& g(g,ϵ(\alpha (g)),g),\hfill \\ m^\alpha :(𝒫^\alpha G)_2𝒫^\alpha G\hfill & ;& ((g,h,s),(s,t,u))(g,ht,u),\hfill \\ i^\alpha :𝒫^\alpha G𝒫^\alpha G\hfill & ;& (g,h,s)(s,h^1,g).\hfill \end{array}$$
Moreover, we also have that the Lie algebroid of $`𝒫^\alpha G`$ may be identified with the prolongation $`𝒫^\alpha (AG)`$ of $`AG`$ over $`\alpha :GM`$. We remark that
$$𝒫_g^\alpha (AG)=\{(v_{ϵ(\alpha (g))},X_g)A_{\alpha (g)}G\times T_gG/(T_{ϵ(\alpha (g))}\beta )(v_{ϵ(\alpha (g))})=(T_g\alpha )(X_g)\},$$
for $`gG`$.
###### Theorem 3.4.
Let $`\mathrm{\Phi }^\alpha :G_\beta G_\alpha 𝒫^\alpha G`$ be the map defined by
$$\mathrm{\Phi }^\alpha (g,h,s)=(g,gh^1,s),$$
for $`(g,h,s)G_\beta G_\alpha `$. Then:
1. $`\mathrm{\Phi }^\alpha `$ is a Lie groupoid isomorphism over the identity of $`G`$.
2. If $`A(\mathrm{\Phi }^\alpha ):A(G_\beta G_\alpha )A(𝒫^\alpha G)`$ is the corresponding Lie algebroid isomorphism then, under the canonical identifications
$$A(G_\beta G_\alpha )𝒫^\tau GV\beta _GV\alpha ,A(𝒫^\alpha G)𝒫^\alpha (AG),$$
$`A(\mathrm{\Phi }^\alpha )`$ is given by
$$A_g(\mathrm{\Phi }^\alpha )(X_g,Y_g)=(T_g(ir_{g^1})(X_g),X_g+Y_g),$$
(3.13)
for $`(X_g,Y_g)V_g\beta V_g\alpha `$, where $`r_{g^1}:\beta ^1(\beta (g))\beta ^1(\alpha (g))`$ is the right-translation by $`g^1`$.
###### Proof.
Proceeding as in the proof of Theorem 3.3, we deduce the result. ∎
## 4. Mechanics on Lie Groupoids
In this section, we introduce Lagrangian (Hamiltonian) Mechanics on an arbitrary Lie groupoid and we will also analyze its geometrical properties. This construction may be considered as a discrete version of the construction of the Lagrangian (Hamiltonian) Mechanics on Lie algebroids proposed in (see also ). We first discuss discrete Euler-Lagrange equations following a similar approach to , using a variational procedure. Secondly, we intrinsically define and discuss the discrete Poincaré-Cartan sections, Legendre transformations, regularity of the Lagrangian and Noether’s theorem.
### 4.1. Discrete Euler-Lagrange equations
Let $`G`$ be a Lie groupoid with structural maps
$$\alpha ,\beta :GM,ϵ:MG,i:GG,m:G_2G.$$
Denote by $`\tau :AGM`$ the Lie algebroid of $`G`$.
A discrete Lagrangian is a function $`L:G`$. Fixed $`gG`$, we define the set of admissible sequences with values in $`G`$:
$$𝒞_g^N=\{(g_1,\mathrm{},g_N)G^N/(g_k,g_{k+1})G_2\text{ for }k=1,\mathrm{},N1\text{ and }g_1\mathrm{}g_n=g\}.$$
Given a tangent vector at $`(g_1,\mathrm{},g_N)`$ to the manifold $`𝒞_g^N`$, we may write it as the tangent vector at $`t=0`$ of a curve in $`𝒞_g^N`$, $`t(\epsilon ,\epsilon )c(t)`$ which passes through $`(g_1,\mathrm{},g_N)`$ at $`t=0`$. This type of curves is of the form
$$c(t)=(g_1h_1(t),h_1^1(t)g_2h_2(t),\mathrm{},h_{N2}^1(t)g_{N1}h_{N1}(t),h_{N1}^1(t)g_N)$$
where $`h_k(t)\alpha ^1(\beta (g_k)),`$ for all $`t,`$ and $`h_k(0)=ϵ(\beta (g_k))`$ for $`k=1,\mathrm{},N1`$.
Therefore, we may identify the tangent space to $`𝒞_g^N`$ at $`(g_1,\mathrm{},g_N)`$ with
$$T_{(g_1,\mathrm{},g_N)}𝒞_g^N\{(v_1,\mathrm{},v_{N1})/v_kA_{x_k}G\text{ and }x_k=\beta (g_k),1kN1\}.$$
Observe that each $`v_k`$ is the tangent vector to the $`\alpha `$-vertical curve $`h_k`$ at $`t=0`$.
The curve $`c`$ is called a variation of $`(g_1,\mathrm{},g_N)`$ and $`(v_1,v_2,\mathrm{},v_{N1})`$ is called an infinitesimal variation of $`(g_1,\mathrm{},g_N)`$.
Define the discrete action sum associated to the discrete Lagrangian $`L:G`$
$$\begin{array}{ccccc}\hfill 𝒮L& :& \hfill 𝒞_g^N& & \hfill \\ & & \hfill (g_1,\mathrm{},g_N)& & \underset{k=1}{\overset{N}{}}L(g_k).\hfill \end{array}$$
We now proceed, as in the continuous case, to derive the discrete equations of motion applying Hamilton’s principle of critical action. For it, we consider variations of the discrete action sum.
###### Definition 4.1 (Discrete Hamilton’s principle ).
Given $`gG`$, an admissible sequence $`(g_1,\mathrm{},g_N)𝒞_g^N`$ is a solution of the Lagrangian system determined by $`L:G`$ if and only if $`(g_1,\mathrm{},g_N)`$ is a critical point of $`𝒮L`$.
Fist of all, in order to characterize the critical points, we need to calculate:
$`{\displaystyle \frac{d}{dt}}|_{t=0}𝒮L(c(t))`$ $`=`$ $`{\displaystyle \frac{d}{dt}}|_{t=0}\{L(g_1h_1(t))+L(h_1^1(t)g_2h_2(t))`$
$`+\mathrm{}+L(h_{N2}^1(t)g_{N1}h_{N1}(t))+L(h_{N1}^1(t)g_N)\}.`$
Therefore,
$`{\displaystyle \frac{d}{dt}}|_{t=0}𝒮L(c(t))`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{N1}{}}}\left(\mathrm{d}^{}(Ll_{g_k})(ϵ(x_k))(v_k)+\mathrm{d}^{}(Lr_{g_{k+1}}i)(ϵ(x_k))(v_k)\right)`$
where $`\mathrm{d}^{}`$ is the standard differential on $`G`$, i.e., $`\mathrm{d}^{}`$ is the differential of the Lie algebroid $`\tau _G:TGG.`$ Since the critical condition is $`{\displaystyle \frac{d}{dt}}|_{t=0}𝒮L(c(t))=0`$ then, applying (2.6) and (2.7), we may rewrite this condition as
$$0=\underset{k=1}{\overset{N1}{}}\left[\stackrel{}{X}_k(g_k)(L)\stackrel{}{X}_k(g_{k+1})(L)\right]=\underset{k=1}{\overset{N1}{}}\left[dL,X_k^{(0,1)}(g_k)dL,X_k^{(1,0)}(g_{k+1})\right]$$
where $`d`$ is the differential of the Lie algebroid $`\pi ^\tau :𝒫^\tau GV\beta _GV\alpha G`$ and $`X_k`$ is a section of $`\tau :AGM`$ such that $`X_k(x_k)=v_k.`$
For $`N=2`$ we obtain that $`(g_1,g_2)G_2`$ (with $`\beta (g_1)=\alpha (g_2)=x`$) is a solution if
$$\mathrm{d}^{}\left[Ll_{g_1}+Lr_{g_2}i\right](ϵ(x))_{|A_xG}=0$$
or, alternatively,
$$\stackrel{}{X}(g_1)(L)\stackrel{}{X}(g_2)(L)=0$$
for every section $`X`$ of $`AG`$. These equations will be called discrete Euler-Lagrange equations.
Thus, we may define the discrete Euler-Lagrange operator:
$$D_{\text{DEL}}L:G_2A^{}G,$$
where $`A^{}G`$ is the dual of $`AG`$. This operator is given by
$$D_{\text{DEL}}L(g,h)=d^0\left[Ll_g+Lr_hi\right](ϵ(x))_{|A_xG}$$
with $`\beta (g)=\alpha (h)=x`$.
In conclusion, we have characterized the solutions of the Lagrangian system determined by $`L:G`$ as the sequences $`(g_1,\mathrm{},g_N)`$, with $`(g_k,g_{k+1})G_2`$, for each $`k\{1,\mathrm{},N1\}`$, and
$$D_{\text{DEL}}L(g_k,g_{k+1})=0,1kN1.$$
### 4.2. Discrete Poincaré-Cartan sections
Given a Lagrangian function $`L:G`$, we will study the geometrical properties of the discrete Euler-Lagrange equations.
Consider the Lie algebroid $`\pi ^\tau :P^\tau GV\beta _GV\alpha G`$, and define the Poincaré-Cartan 1-sections $`\mathrm{\Theta }_L^{},\mathrm{\Theta }_L^+\mathrm{\Gamma }((\pi ^\tau )^{})`$ as follows
$$\mathrm{\Theta }_L^{}(g)(X_g,Y_g)=X_g(L),\mathrm{\Theta }_L^+(g)(X_g,Y_g)=Y_g(L),$$
(4.1)
for each $`gG`$ and $`(X_g,Y_g)V_g\beta V_g\alpha `$. From the definition, we have that
$$\mathrm{\Theta }_L^{}(g)(X^{(1,0)}(g))=\stackrel{}{X}(g)(L)\text{and}\mathrm{\Theta }_L^{}(g)(X^{(0,1)}(g))=0,$$
and similarly
$$\mathrm{\Theta }_L^+(g)(X^{(0,1)}(g))=\stackrel{}{X}(g)(L)\text{and}\mathrm{\Theta }_L^+(g)(X^{(1,0)}(g))=0,$$
for $`X\mathrm{\Gamma }(\tau ).`$
We also have that $`dL=\mathrm{\Theta }_L^+\mathrm{\Theta }_L^{}`$ and so, using $`d^2=0,`$ it follows that $`d\mathrm{\Theta }_L^+=d\mathrm{\Theta }_L^{}`$. This means that there exists a unique 2-section $`\mathrm{\Omega }_L=d\mathrm{\Theta }_L^+=d\mathrm{\Theta }_L^{}`$, that will be called the Poincaré-Cartan 2-section. This 2-section will be important for studying symplecticity of the discrete Euler-Lagrange equations.
###### Proposition 4.2.
If $`X`$ and $`Y`$ are sections of the Lie algebroid $`AG`$ then
$$\mathrm{\Omega }_L(X^{(1,0)},Y^{(1,0)})=0,\mathrm{\Omega }_L(X^{(0,1)},Y^{(0,1)})=0,$$
and
$$\mathrm{\Omega }_L(X^{(1,0)},Y^{(0,1)})=\stackrel{}{X}(\stackrel{}{Y}L)\text{ and }\mathrm{\Omega }_L(X^{(0,1)},Y^{(1,0)})=\stackrel{}{Y}(\stackrel{}{X}L).$$
###### Proof.
A direct computation proves the result. ∎
###### Remark 4.3.
Remark 4.3. Let $`g`$ be an element of $`G`$ such that $`\alpha (g)=x`$ and $`\beta (g)=y`$. Suppose that $`U`$ and $`V`$ are open subsets of $`M`$, with $`xU`$ and $`yV`$, and that $`\{X_i\}`$ and $`\{Y_j\}`$ are local bases of $`\mathrm{\Gamma }(\tau )`$ on $`U`$ and $`V`$, respectively. Then, $`\{X_i^{(1,0)},Y_j^{(0,1)}\}`$ is a local basis of $`\mathrm{\Gamma }(\pi ^\tau )`$ on the open subset $`\alpha ^1(U)\beta ^1(V)`$. Moreover, if we denote by $`\{(X^i)^{(1,0)},(Y^j)^{(0,1)}\}`$ the dual basis of $`\{X_i^{(1,0)},Y_j^{(0,1)}\}`$, we have that on the open subset $`\alpha ^1(U)\beta ^1(V)`$
$$\begin{array}{c}\mathrm{\Theta }_L^{}=\stackrel{}{X_i}(L)(X^i)^{(1,0)},\mathrm{\Theta }_L^+=\stackrel{}{Y_j}(L)(Y^j)^{(0,1)},\hfill \\ \mathrm{\Omega }_L=\stackrel{}{X_i}(\stackrel{}{Y_j}(L))(X^i)^{(1,0)}(Y^j)^{(0,1)}.\hfill \end{array}$$
$``$
Finally, we obtain some useful expressions of the Poincaré-Cartan 1-sections using the Lie algebroid isomorphisms introduced in Theorems 3.3 and 3.4.
We recall that the maps
$$\begin{array}{ccc}A(\mathrm{\Phi }^\beta ):A(G_\beta G_\alpha )V\beta _GV\alpha \hfill & & \hfill A(𝒫^\beta G)𝒫^\beta (AG)\\ A(\mathrm{\Phi }^\alpha ):A(G_\beta G_\alpha )V\beta _GV\alpha \hfill & & \hfill A(𝒫^\alpha G)𝒫^\alpha (AG)\end{array}$$
given by (3.11) and (3.13) are Lie algebroid isomorphisms over the identity of $`G`$. Moreover, if $`(v_{ϵ(\beta (g))},Z_g)𝒫_g^\beta (AG)`$ then, from (3.11), it follows that
$$A_g(\mathrm{\Phi }^\beta )^1(v_{ϵ(\beta (g))},Z_g)=(Z_g(T_{ϵ(\beta (g))}l_g)(v_{ϵ(\beta (g))}),(T_{ϵ(\beta (g))}l_g)(v_{ϵ(\beta (g))})).$$
(4.2)
On the other hand, if $`(v_{ϵ(\alpha (h))},Z_h)𝒫_h^\alpha (AG)`$ then, using (3.13), we deduce that
$$A_h(\mathrm{\Phi }^\alpha )^1(v_{ϵ(\alpha (h))},Z_h)=(T_{ϵ(\alpha (h))}(r_hi)(v_{ϵ(\alpha (h))}),Z_hT_{ϵ(\alpha (h))}(r_hi)(v_{ϵ(\alpha (h))})).$$
(4.3)
Now, we introduce the sections $`\mathrm{\Theta }_L^\alpha \mathrm{\Gamma }((\tau ^\alpha )^{})`$ and $`\mathrm{\Theta }_L^\beta \mathrm{\Gamma }((\tau ^\beta )^{})`$ given by
$$\mathrm{\Theta }_L^\alpha =(A(\mathrm{\Phi }^\alpha )^1,Id)^{}(\mathrm{\Theta }_L^{}),\mathrm{\Theta }_L^\beta =(A(\mathrm{\Phi }^\beta )^1,Id)^{}(\mathrm{\Theta }_L^+).$$
(4.4)
Using (4.2) and (4.3), we obtain that
$`\mathrm{\Theta }_L^\alpha (h)(v_{ϵ(\alpha (h))},Z_h)`$ $`=v_{ϵ(\alpha (h))}(Lr_hi),`$ (4.5)
$`\mathrm{\Theta }_L^\beta (g)(v_{ϵ(\beta (g))},Z_g)`$ $`=v_{ϵ(\beta (g))}(Ll_g),`$ (4.6)
for $`(v_{ϵ(\alpha (h))},Z_h)𝒫_h^\alpha (AG)`$ and $`(v_{ϵ(\beta (g))},Z_g)𝒫_g^\beta (AG)`$.
#### 4.2.1. Poincaré-Cartan $`1`$-sections: variational motivation
Now, we follow a variational procedure to construct the $`1`$-sections $`\mathrm{\Theta }_L^+`$ and $`\mathrm{\Theta }_L^{}`$. We begin by calculating the extremals of $`𝒮L`$ for variations that do not fix the point $`gG`$. For it, we consider the manifold
$$𝒞^N=\{(g_1,\mathrm{},g_N)G^N/(g_k,g_{k+1})G_2\text{ for each }k,\mathrm{\hspace{0.33em}\hspace{0.33em}1}kN1\}.$$
If $`c:(\epsilon ,\epsilon )𝒞^N`$ is a curve in $`𝒞^N`$ and $`c(0)=(g_1,\mathrm{},g_N)`$ then there exist $`N+1`$ curves $`h_k:(\epsilon ,\epsilon )\alpha ^1(\beta (g_k))`$, for $`0kN,`$ with $`h_k(0)=ϵ(\beta (g_k))`$ and $`g_0=g_1^1`$, such that
$$c(t)=(h_0^1(t)g_1h_1(t),h_1^1(t)g_2h_2(t),\mathrm{},h_{N2}^1(t)g_{N1}h_{N1}(t),h_{N1}^1(t)g_Nh_N(t))$$
for $`t(\epsilon ,\epsilon ).`$ Thus, the tangent space to $`𝒞^N`$ at $`(g_1,\mathrm{},g_N)`$ may be identified with the vector space $`A_{\beta (g_0)}G\times A_{\beta (g_1)}G\times \mathrm{}\times A_{\beta (g_N)}G,`$ that is,
$$T_{(g_1,g_2,\mathrm{},g_N)}𝒞^N\left\{(v_0,v_1,\mathrm{},v_N)/v_kA_{x_k}G,x_k=\beta (g_k),0kN\right\}.$$
Now, proceeding as in Section 4.1, we introduce the action sum
$$𝒮L:𝒞^N,SL(g_1,\mathrm{},g_N)=\underset{k=1}{\overset{N}{}}L(g_k).$$
Then,
$`{\displaystyle \frac{d}{dt}}|_{t=0}𝒮L(c(t))`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{N1}{}}}\left[\mathrm{d}^{}(Ll_{g_k})(ϵ(x_k))(v_k)+\mathrm{d}^{}(Lr_{g_{k+1}}i)(ϵ(x_{k+1}))(v_k)\right]`$ (4.7)
$`+\mathrm{d}^{}(Lr_{g_1}i)(ϵ(x_0))(v_0)+\mathrm{d}^{}(Ll_{g_N})(ϵ(x_N))(v_N).`$
Therefore, if $`X_0,\mathrm{},X_N`$ are sections of $`\tau :AGM`$ satisfying, $`X_k(x_k)=v_k,`$ for all $`k,`$ we have that
$`{\displaystyle \frac{d}{dt}}|_{t=0}𝒮L(c(t))={\displaystyle \underset{k=1}{\overset{N1}{}}}\left[\stackrel{}{X}_k(g_k)(L)\stackrel{}{X}_k(g_{k+1})(L)\right]\stackrel{}{X_0}(g_1)(L)+\stackrel{}{X}_N(g_N)(L)`$
$`={\displaystyle \underset{k=1}{\overset{N1}{}}}(D_{\text{DEL}}L(g_k,g_{k+1}))(v_k)+\mathrm{\Theta }_L^{}(g_1)(X_0^{(1,0)}(g_1))+\mathrm{\Theta }_L^+(g_N)(X_N^{(0,1)}(g_N)).`$
Note that it is in the last two terms (that arise from the boundary variations) where appear the Poincaré-Cartan 1-sections.
### 4.3. Discrete Lagrangian evolution operator
We say that a differentiable mapping $`\xi :GG`$ is a discrete flow or a discrete Lagrangian evolution operator for $`\mathrm{L}`$ if it verifies the following properties:
1. $`\text{graph}(\xi )G_2`$, that is, $`(g,\xi (g))G_2`$, $`gG`$ ($`\xi `$ is a second order operator).
2. $`(g,\xi (g))`$ is a solution of the discrete Euler-Lagrange equations, for all $`gG`$, that is, $`(D_{\text{DEL}}L)(g,\xi (g))=0,`$ for all $`gG.`$
In such a case
$$\mathrm{d}^{}(Ll_g+Lr_{\xi (g)}i)(ϵ(\beta (g)))_{|A_{\beta (g)}G}=0,\text{for all }gG$$
(4.8)
or, in other terms,
$$\stackrel{}{X}(g)(L)\stackrel{}{X}(\xi (g))(L)=0$$
(4.9)
for every section $`X`$ of $`AG`$ and every $`gG.`$
Now, we define the prolongation $`𝒫^\tau \xi :V\beta _GV\alpha V\beta _GV\alpha `$ of the second order operator $`\xi :GG`$ as follows:
$$𝒫^\tau \xi =A(\mathrm{\Phi }^\alpha )^1(Id,T\xi )A(\mathrm{\Phi }^\beta )$$
(4.10)
with $`A(\mathrm{\Phi }^\alpha )`$ and $`A(\mathrm{\Phi }^\beta )`$ the isomorphisms defined in Theorems 3.3 and 3.4 and $`(Id,T\xi ):𝒫^\beta (AG)𝒫^\alpha (AG)`$ the map given by
$$(Id,T\xi )(v_{ϵ(\beta (g))},X_g)=(v_{ϵ(\beta (g))},(T_g\xi )(X_g)),\text{ for }(v_{ϵ(\beta (g))},X_g)𝒫_g^\beta (AG).$$
Since the pair $`((Id,T\xi ),\xi )`$ is a Lie algebroid morphism between the Lie algebroids $`𝒫^\beta (AG)G`$ and $`𝒫^\alpha (AG)G`$ then the pair $`(𝒫^\tau \xi ,\xi )`$ is also a Lie algebroid morphism
From the definition of $`𝒫^\tau \xi `$, we deduce that
$$𝒫_g^\tau \xi (X_g,Y_g)=((T_g(r_{g\xi (g)}i))(Y_g),(T_g\xi )(X_g)+(T_g\xi )(Y_g)T_g(r_{g\xi (g)}i)(Y_g))$$
(4.11)
for all $`(X_g,Y_g)V_g\beta V_g\alpha `$. Moreover, from (2.10) and (4.11), we obtain that
$$𝒫^\tau \xi (\stackrel{}{X}(g),\stackrel{}{Y}(g))=(\stackrel{}{Y}(\xi (g)),(T_g\xi )(\stackrel{}{X}(g)+\stackrel{}{Y}(g))+\stackrel{}{Y}(\xi (g)))$$
(4.12)
for all $`X,Y`$ sections of $`AG`$.
### 4.4. Preservation of Poincaré-Cartan sections
The following result explains the sense in which the discrete Lagrange evolution operator preserves the Poincaré-Cartan 2-section.
###### Theorem 4.4.
Let $`L:G`$ be a discrete Lagrangian on a Lie groupoid $`G`$. Then:
1. The map $`\xi `$ is a discrete Lagrangian evolution operator for $`L`$ if and only if $`(𝒫^\tau \xi ,\xi )^{}\mathrm{\Theta }_L^{}=\mathrm{\Theta }_L^+`$.
2. The map $`\xi `$ is a discrete Lagrangian evolution operator for $`L`$ if and only if $`(𝒫^\tau \xi ,\xi )^{}\mathrm{\Theta }_L^{}\mathrm{\Theta }_L^{}=dL`$.
3. If $`\xi `$ is discrete Lagrangian evolution operator then $`(𝒫^\tau \xi ,\xi )^{}\mathrm{\Omega }_L=\mathrm{\Omega }_L`$.
###### Proof.
From (4.4), it follows
$$(A(\mathrm{\Phi }^\alpha ),Id)^{}(\mathrm{\Theta }_L^\alpha )=\mathrm{\Theta }_L^{},(A(\mathrm{\Phi }^\beta ),Id)^{}(\mathrm{\Theta }_L^\beta )=\mathrm{\Theta }_L^+.$$
(4.13)
On the other hand, if $`(v_{ϵ(\beta (g))},X_g)𝒫_g^\beta (AG)`$ then, using (4.5) and (4.6), we have that
$$\left\{((\text{Id},T\xi ),\xi )^{}(\mathrm{\Theta }_L^\alpha )\right\}(g)(v_{ϵ(\beta (g))},X_g)=v_{ϵ(\beta (g))}(Lr_{\xi (g)}i)$$
and
$$\mathrm{\Theta }_L^\beta (g)(v_{ϵ(\beta (g))},X_g)=v_{ϵ(\beta (g))}(Ll_g).$$
Thus, $`((\text{Id},T\xi ),\xi )^{}\mathrm{\Theta }_L^\alpha =\mathrm{\Theta }_L^\beta `$ if and only if $`\xi `$ is a discrete Lagrangian evolution operator for $`L.`$ Therefore, using this fact and (4.13), we prove (i).
The second property follows from (i) by taking into account that $`dL=\mathrm{\Theta }_L^+\mathrm{\Theta }_L^{}`$. Finally, (iii) follows using (ii) and the fact that $`(𝒫^\tau \xi ,\xi )`$ is a Lie algebroid morphism.
###### Remark 4.5.
Now, we present a proof of the preservation of the Poincaré-Cartan $`2`$-section using variational arguments. Given a discrete Lagrangian evolution operator $`\xi :GG`$ for $`L,`$ we may consider the function $`𝒮_\xi L:G`$ given by
$$(𝒮_\xi L)(g)=L(g)+L(\xi (g)),\text{ for }gG.$$
If $`d`$ is the differential on the Lie algebroid $`V\beta _GV\alpha G`$ and $`X,Y`$ are sections of $`\tau :AGM`$ then, using (4.1), (4.9) and (4.12), we obtain that
$`d(𝒮_\xi L)(g)(\stackrel{}{X}(g),\stackrel{}{Y}(g))`$ $`=`$ $`\stackrel{}{X}(g)L+\stackrel{}{Y}(g)L+(T_g\xi )(\stackrel{}{X}(g))L+(T_g\xi )(\stackrel{}{Y}(g))L`$
$`=`$ $`\stackrel{}{X}(g)L+\stackrel{}{Y}(g)L\stackrel{}{Y}(\xi (g))L+(T_g\xi )(\stackrel{}{X}(g))L`$
$`+(T_g\xi )(\stackrel{}{Y}(g))L+\stackrel{}{Y}(\xi (g))L`$
$`=`$ $`\mathrm{\Theta }_L^{}(g)(\stackrel{}{X}(g),\stackrel{}{Y}(g))+\left[(𝒫^\tau \xi ,\xi )^{}\mathrm{\Theta }_L^+\right](g)(\stackrel{}{X}(g),\stackrel{}{Y}(g)).`$
This implies that
$$(𝒫^\tau \xi ,\xi )^{}\mathrm{\Theta }_L^+\mathrm{\Theta }_L^{}=d(𝒮_\xi L).$$
Thus, we conclude that $`(𝒫^\tau \xi ,\xi )^{}\mathrm{\Omega }_L=\mathrm{\Omega }_L`$. $``$
### 4.5. Lie groupoid morphisms and reduction
Let $`(\mathrm{\Phi },\mathrm{\Phi }_0)`$ be a Lie groupoid morphism between the Lie groupoids $`GM`$ and $`G^{}M^{}`$. The prolongation $`𝒫^\tau \mathrm{\Phi }:V\beta _GV\alpha V\beta ^{}_G^{}V\alpha ^{}`$ of the morphism $`(\mathrm{\Phi },\mathrm{\Phi }_0)`$ is defined by
$$𝒫_g^\tau \mathrm{\Phi }(V,W)=(T_g\mathrm{\Phi }(V),T_g\mathrm{\Phi }(W))$$
(4.14)
for every $`(V,W)V_g\beta V_g\alpha `$. It is easy to see that $`(𝒫^\tau \mathrm{\Phi },\mathrm{\Phi })`$ is a morphism of Lie algebroids.
###### Theorem 4.6.
Let $`(\mathrm{\Phi },\mathrm{\Phi }_0)`$ be a morphism of Lie groupoids from $`GM`$ to $`G^{}M^{}`$. Let $`L`$ and $`L^{}`$ be discrete Lagrangian functions on $`G`$ and $`G^{}`$, respectively, related by $`L=L^{}\mathrm{\Phi }`$. Then:
1. for every $`(g,h)G_2`$ and every $`vA_{\beta (g)}G`$ we have that
$$D_{\mathrm{DEL}}L(g,h)(v)=D_{\mathrm{DEL}}L^{}(\mathrm{\Phi }(g),\mathrm{\Phi }(h))\left(A_{\beta (g)}\mathrm{\Phi }(v)\right).$$
(4.15)
2. $`(𝒫^\tau \mathrm{\Phi },\mathrm{\Phi })^{}\mathrm{\Theta }_L^{}^+=\mathrm{\Theta }_L^+,`$
3. $`(𝒫^\tau \mathrm{\Phi },\mathrm{\Phi })^{}\mathrm{\Theta }_L^{}^{}=\mathrm{\Theta }_L^{},`$
4. $`(𝒫^\tau \mathrm{\Phi },\mathrm{\Phi })^{}\mathrm{\Omega }_L^{}=\mathrm{\Omega }_L.`$
###### Proof.
To prove the first we notice that, if $`(\mathrm{\Phi },\mathrm{\Phi }_0)`$ is a morphism of Lie groupoids, then we have that $`\mathrm{\Phi }l_g=l_{\mathrm{\Phi }(g)}\mathrm{\Phi }`$ and $`\mathrm{\Phi }r_h=r_{\mathrm{\Phi }(h)}\mathrm{\Phi }`$, from where we get
$`D_{\mathrm{DEL}}L(g,h)(v)`$ $`=Tl_g(v)L+Tr_h(Ti(v))L`$
$`=Tl_g(v)(L^{}\mathrm{\Phi })+Tr_h(Ti(v))(L^{}\mathrm{\Phi })`$
$`=T\mathrm{\Phi }(Tl_g(v))L^{}+T\mathrm{\Phi }(Tr_h(Ti(v)))L^{}`$
$`=Tl_{\mathrm{\Phi }(g)}(T\mathrm{\Phi }(v))L^{}+Tr_{\mathrm{\Phi }(h)}(T\mathrm{\Phi }(Ti(v)))L^{}`$
$`=Tl_{\mathrm{\Phi }(g)}(T\mathrm{\Phi }(v))L^{}+Tr_{\mathrm{\Phi }(h)}(Ti^{}(T\mathrm{\Phi }(v)))L^{}`$
$`=D_{\mathrm{DEL}}L^{}(\mathrm{\Phi }(g),\mathrm{\Phi }(h))\left(A_{\beta (g)}\mathrm{\Phi }(v)\right),`$
where we have also used that $`i^{}\mathrm{\Phi }=\mathrm{\Phi }i`$ and $`A_{\beta (g)}\mathrm{\Phi }(v)=T\mathrm{\Phi }(v)`$.
For the proof of the second, we have that
$`(𝒫^\tau \mathrm{\Phi },\mathrm{\Phi })^{}\mathrm{\Theta }_L^{}^+(g),(V,W)`$ $`=\mathrm{\Theta }_L^{}^+(\mathrm{\Phi }(g)),(T_g\mathrm{\Phi }(V),T_g\mathrm{\Phi }(W))`$
$`=((T_g\mathrm{\Phi })(W))L^{}=WL=\mathrm{\Theta }_L^+(g),(V,W),`$
for every $`(V,W)𝒫_g^\tau G`$. The proof of the third is similar to the second, and finally, for the proof of (iv) we just take the differential in (ii). ∎
As an immediate consequence of the above theorem we have that
###### Corollary 4.7.
Let $`(\mathrm{\Phi },\mathrm{\Phi }_0)`$ be a morphism of Lie groupoids from $`GM`$ to $`G^{}M^{}`$ and suppose that $`(g,h)G_2`$.
1. If $`(\mathrm{\Phi }(g),\mathrm{\Phi }(h))`$ is a solution of the discrete Euler-Lagrange equations for $`L^{}=L\mathrm{\Phi }`$, then $`(g,h)`$ is a solution of the discrete Euler-Lagrange equations for $`L`$.
2. If $`\mathrm{\Phi }`$ is a submersion then $`(g,h)`$ is a solution of the discrete Euler-Lagrange equations for $`L`$ if and only if $`(\mathrm{\Phi }(g),\mathrm{\Phi }(h))`$ is a solution of the discrete Euler-Lagrange equations for $`L^{}`$.
3. If $`\mathrm{\Phi }`$ is an immersion, then $`(g,h)`$ is a solution of the discrete Euler-Lagrange equations for $`L`$ if and only if $`D_{\mathrm{DEL}}L(\mathrm{\Phi }(g),\mathrm{\Phi }(h))`$ vanishes over $`\mathrm{Im}(A_\beta (g)\mathrm{\Phi })`$.
The case when $`\mathrm{\Phi }`$ is an inmersion may be useful to modelize holonomic mechanics on Lie groupoids, which is an imprescindible tool for explicitely construct geometric integrators (see ).
The particular case when $`\mathrm{\Phi }`$ is a submersion is relevant for reduction (see Section 5.5 in this paper).
### 4.6. Discrete Legendre transformations
Given a Lagrangian $`L:G`$ we define, just as the standard case , two discrete Legendre transformations $`𝔽^{}L:GA^{}G`$ and $`𝔽^+L:GA^{}G`$ as follows
$`(𝔽^{}L)(h)(v_{ϵ(\alpha (h))})`$ $`=`$ $`v_{ϵ(\alpha (h))}(Lr_hi),\text{ for }v_{ϵ(\alpha (h))}A_{\alpha (h)}G,`$
$`(𝔽^+L)(g)(v_{ϵ(\beta (g))})`$ $`=`$ $`v_{ϵ(\beta (g))}(Ll_g),\text{ for }v_{ϵ(\beta (g))}A_{\beta (g)}G.`$
###### Remark 4.8.
Note that $`(𝔽^{}L)(h)A_{\alpha (h)}^{}G.`$ Furthermore, if $`U`$ is an open subset of $`M`$ such that $`\alpha (h)U`$ and $`\{X_i\}`$ is a local basis of $`\mathrm{\Gamma }(\tau )`$ on $`U`$ then
$$𝔽^{}L=\stackrel{}{X_i}(L)(X^i\alpha ),$$
on $`\alpha ^1(U),`$ where $`\{X^i\}`$ is the dual basis of $`\{X_i\}`$. In a similar way, if $`V`$ is an open subset of $`M`$ such that $`\beta (g)V`$ and $`\{Y_j\}`$ is a local basis of $`\mathrm{\Gamma }(\tau )`$ on $`V`$ then
$$𝔽^+L=\stackrel{}{Y_j}(L)(Y^j\beta ),$$
on $`\beta ^1(V).`$ $``$
Next, we consider the prolongation $`\tau ^\tau ^{}:𝒫^\tau ^{}(AG)A^{}G`$ of the Lie algebroid $`\tau :AGM`$ over the fibration $`\tau ^{}:A^{}GM`$, that is,
$$\begin{array}{ccc}\hfill 𝒫_v^{}^\tau ^{}(AG)& =& \{(v_{\tau ^{}(v^{})},X_v^{})A_{\tau ^{}(v^{})}G\times T_v^{}(A^{}G)/(T_{\tau ^{}(v^{})}\beta )(v_{\tau ^{}(v^{})})\hfill \\ & & =(T_v^{}\tau ^{})(X_v^{})\}\hfill \end{array}$$
for $`v^{}A^{}G`$. Then, we may introduce the canonical section $`\mathrm{\Theta }`$ of the vector bundle $`(\tau ^\tau ^{})^{}:(𝒫^\tau ^{}AG)^{}A^{}G`$ as follows:
$$\mathrm{\Theta }(v^{})(v_{\tau ^{}(v^{})},X_v^{})=v^{}(v_{\tau ^{}(v^{})}),$$
(4.16)
for $`v^{}A^{}G`$ and $`(v_{\tau ^{}(v^{})},X_v^{})𝒫_v^{}^\tau ^{}(AG)`$. $`\mathrm{\Theta }`$ is called the Liouville section. Moreover, we define the canonical symplectic section $`\mathrm{\Omega }`$ associated with $`AG`$ by $`\mathrm{\Omega }=d\mathrm{\Theta }`$, where $`d`$ is the differential on the Lie algebroid $`\tau ^\tau ^{}:𝒫^\tau ^{}(AG)A^{}G`$. It is easy to prove that $`\mathrm{\Omega }`$ is nondegenerate and closed, that is, it is a symplectic section of $`𝒫^\tau ^{}(AG)`$ (see ).
Now, let $`𝒫^\tau 𝔽^{}L`$ be the prolongation of $`𝔽^{}L`$ defined by
$$𝒫^\tau 𝔽^{}L=(\text{Id},T𝔽^{}L)A(\mathrm{\Phi }^\alpha ):𝒫^\tau GV\beta _GV\alpha 𝒫^\tau ^{}(AG),$$
(4.17)
where $`A(\mathrm{\Phi }^\alpha ):𝒫^\tau GV\beta _GV\alpha 𝒫^\alpha (AG)`$ is the Lie algebroid isomorphism (over the identity of $`G`$) defined by (3.13) and $`(Id,T𝔽^{}L):𝒫^\alpha (AG)𝒫^\tau ^{}(AG)`$ is the map given by
$$(Id,T𝔽^{}L)(v_{ϵ(\alpha (h))},X_h)=(v_{ϵ(\alpha (h))},(T_h𝔽^{}L)(X_h)),$$
for $`(v_{ϵ(\alpha (h))},X_h)𝒫_h^\alpha (AG).`$ Since the pair $`((Id,T𝔽^{}L),𝔽^{}L)`$ is a morphism between the Lie algebroids $`𝒫^\alpha (AG)G`$ and $`𝒫^\tau ^{}(AG)A^{}G`$, we deduce that $`(𝒫^\tau 𝔽^{}L,𝔽^{}L)`$ is also a morphism between the Lie algebroids $`𝒫^\tau GV\beta _GV\alpha G`$ and $`𝒫^\tau ^{}(AG)A^{}G.`$ The following diagram illustrates the above situation:
The prolongation $`𝒫^\tau 𝔽^{}L`$ can be explicitly written as
$$𝒫_h^\tau 𝔽^{}L(X_h,Y_h)=(T_h(ir_{h^1})(X_h),(T_h𝔽^{}L)(X_h)+(T_h𝔽^{}L)(Y_h)),$$
(4.18)
for $`hG`$ and $`(X_h,Y_h)V_h\beta V_h\alpha `$.
###### Proposition 4.9.
If $`\mathrm{\Theta }`$ is the Liouville section of the vector bundle $`(𝒫^\tau ^{}(AG))^{}A^{}G`$ and $`\mathrm{\Omega }=d\mathrm{\Theta }`$ is the canonical symplectic section of $`^2(𝒫^\tau ^{}(AG))^{}A^{}G`$ then
$`(𝒫^\tau (𝔽^{}L),𝔽^{}L)^{}\mathrm{\Theta }=\mathrm{\Theta }_L^{},`$ $`(𝒫^\tau (𝔽^{}L),𝔽^{}L)^{}\mathrm{\Omega }=\mathrm{\Omega }_L.`$
###### Proof.
Let $`\mathrm{\Theta }_L^\alpha `$ be the section of $`(\tau ^\alpha )^{}:𝒫^\alpha (AG)^{}G`$ defined by (4.4). Then, from (4.5) and (4.16), we deduce that
$$((Id,T𝔽^{}L),𝔽^{}L)^{}\mathrm{\Theta }=\mathrm{\Theta }_L^\alpha .$$
Thus, using (4.4), we obtain that
$$(𝒫^\tau 𝔽^{}L,𝔽^{}L)^{}\mathrm{\Theta }=\mathrm{\Theta }_L^{}.$$
Therefore, since the pair $`(𝒫^\tau 𝔽^{}L,𝔽^{}L)`$ is a Lie algebroid morphism, it follows that
$$(𝒫^\tau 𝔽^{}L,𝔽^{}L)^{}\mathrm{\Omega }=\mathrm{\Omega }_L.$$
Now, we consider the prolongation $`𝒫^\tau 𝔽^+L`$ of $`𝔽^+L`$ defined by
$$𝒫^\tau 𝔽^+L=(Id,T𝔽^+L)A(\mathrm{\Phi }^\beta ):𝒫^\tau GV\beta _GV\alpha 𝒫^\tau ^{}(AG),$$
(4.19)
where $`A(\mathrm{\Phi }^\beta ):𝒫^\tau GV\beta _GV\alpha 𝒫^\beta (AG)`$ is the Lie algebroid isomorphism (over the identity of $`G`$) defined by (3.11) and $`(Id,T𝔽^+L):𝒫^\beta (AG)𝒫^\tau ^{}(AG)`$ is the map given by
$$(Id,T𝔽^+L)(v_{ϵ(\beta (g))},X_g)=(v_{ϵ(\beta (g))},(T_g𝔽^+L)(X_g)),$$
for $`(v_{ϵ(\beta (g))},X_g)𝒫_g^\beta (AG).`$ As above, the pair $`(𝒫^\tau 𝔽^+L,𝔽^+L)`$ is a morphism between the Lie algebroids $`𝒫^\tau GV\beta _GV\alpha G`$ and $`𝒫^\tau ^{}(AG)A^{}G`$ and the following diagram illustrates the situation
We also have:
###### Proposition 4.10.
If $`\mathrm{\Theta }`$ is the Liouville section of the vector bundle $`(𝒫^\tau ^{}(AG))^{}A^{}G`$ and $`\mathrm{\Omega }=d\mathrm{\Theta }`$ is the canonical symplectic section of $`^2(𝒫^\tau ^{}(AG))^{}A^{}G`$ then
$`(𝒫^\tau (𝔽^+L),𝔽^+L)^{}\mathrm{\Theta }=\mathrm{\Theta }_L^+,`$ $`(𝒫^\tau (𝔽^+L),𝔽^+L)^{}\mathrm{\Omega }=\mathrm{\Omega }_L.`$
###### Remark 4.11.
$`(i)`$ If $`\xi :GG`$ is a smooth map then $`\xi `$ is a discrete Lagrangian evolution operator for $`L`$ if and only if $`𝔽^{}L\xi =𝔽^+L.`$
$`(ii)`$ If $`(g,h)G_2`$ we have that
$$(D_{\text{DEL}}L)(g,h)=𝔽^+L(g)𝔽^{}L(h).$$
(4.20)
$``$
### 4.7. Discrete regular Lagrangians
First of all, we will introduce the notion of a discrete regular Lagrangian.
###### Definition 4.12.
A Lagrangian $`L:G`$ on a Lie groupoid $`G`$ is said to be regular if the Poincaré-Cartan $`2`$-section $`\mathrm{\Omega }_L`$ is symplectic on the Lie algebroid $`𝒫^\tau GV\beta _GV\alpha G.`$
Next, we will obtain necessary and sufficient conditions for a discrete Lagrangian on a Lie groupoid to be regular.
###### Theorem 4.13.
Let $`L:G`$ be a Lagrangian function. Then:
a) The following conditions are equivalent:
1. $`L`$ is regular.
2. The Legendre transformation $`𝔽^{}L`$ is a local diffeomorphism.
3. The Legendre transformation $`𝔽^+L`$ is a local diffeomorphism.
b) If $`L:G`$ is regular and $`(g_0,h_0)G_2`$ is a solution of the discrete Euler-Lagrange equations for $`L`$ then there exist two open subsets $`U_0`$ and $`V_0`$ of $`G`$, with $`g_0U_0`$ and $`h_0V_0,`$ and there exists a (local) discrete Lagrangian evolution operator $`\xi _L:U_0V_0`$ such that:
1. $`\xi _L(g_0)=h_0,`$
2. $`\xi _L`$ is a diffeomorphism and
3. $`\xi _L`$ is unique, that is, if $`U_0^{}`$ is an open subset of $`G`$, with $`g_0U_0^{}`$ and $`\xi _L^{}:U_0^{}G`$ is a (local) discrete Lagrangian evolution operator then $`\xi _{L|U_0U_0^{}}^{}=\xi _{L|U_0U_0^{}}`$.
###### Proof.
$`a)`$ First we will deduce the equivalence of the three conditions
(i)$``$(ii) If $`hG`$, we need to prove that $`T_h(𝔽^{}L):T_hGT_{𝔽^{}L(h)}A^{}G`$ is a linear isomorphism. Assume that there exists $`Y_hT_hG`$ such that $`T_h(𝔽^{}L)(Y_h)=0`$. Since $`\tau ^{}𝔽^{}L=\alpha ,`$ then $`(T_h\alpha )(Y_h)=0`$, that is, $`Y_hV_h\alpha `$.
Therefore, $`(0_h,Y_h)V_h\beta V_h\alpha `$ and, from (4.18), we have that $`𝒫_h^\tau (𝔽^{}L)(0_h,Y_h)=0`$. Moreover, $`(𝒫_h^\tau (𝔽^{}L))^{}\mathrm{\Omega }(𝔽^{}L(h))=\mathrm{\Omega }_L(h)`$ and $`\mathrm{\Omega }(𝔽^{}L(h))`$ and $`\mathrm{\Omega }_L(h)`$ are nondegenerate. Therefore, we deduce that $`𝒫_h^\tau (𝔽^{}L)`$ is a linear isomorphism. This implies that $`Y_h=0`$. This proves that $`T_h(𝔽^{}L):T_hGT_{𝔽^{}L(h)}(A^{}G)`$ is a linear isomorphism. In the same way we deduce (i)$``$(iii).
(ii)$``$(i) We will assume that $`𝔽^+L`$ is a local diffeomorphism, so that
$$𝒫_g^\tau 𝔽^+L:𝒫_g^\tau GV_g\beta V_g\alpha 𝒫_{𝔽^+L(g)}^\tau ^{}(AG)$$
is a linear isomorphism, for all $`gG`$.
On the other hand, if $`\mathrm{\Omega }`$ is the canonical symplectic section of the vector bundle $`^2(𝒫^\tau ^{}(AG))^{}A^{}G`$ then, from Proposition 4.10, we deduce that
$$(𝒫_g^\tau 𝔽^+L)^{}(\mathrm{\Omega }(𝔽^+L(g)))=\mathrm{\Omega }_L(g).$$
Thus, since $`\mathrm{\Omega }(𝔽^+L(g))`$ is nondegenerate, we conclude that $`\mathrm{\Omega }_L(g)`$ is also nondegenerate, for all $`gG.`$ Using the same arguments we deduce (iii)$``$(i).
$`b)`$ Using Remark 4.11, we have that
$$(𝔽^+L)(g_0)=(𝔽^{}L)(h_0)=\mu _0A^{}G.$$
Thus, from the first part of this theorem, it follows that there exit two open subsets $`U_0`$ and $`V_0`$ of $`G`$, with $`g_0U_0`$ and $`h_0V_0`$, and an open subset $`W_0`$ of $`A^{}G`$ such that $`\mu _0W_0`$ and
$$𝔽^+L:U_0W_0,𝔽^{}L:V_0W_0$$
are diffeomorphisms. Therefore, using Remark 4.11, we deduce that
$$\xi _L=[(𝔽^{}L)^1(𝔽^+L)]_{|U_0}:U_0V_0$$
is a (local) discrete Lagrangian evolution operator. Moreover, it is clear that $`\xi _L(g_0)=h_0`$ and, from the first part of this theorem, we have that $`\xi _L`$ is a diffeomorphism.
Finally, if $`U_0^{}`$ is an open subset of $`G`$, with $`g_0U_0^{}`$, and $`\xi _L^{}:U_0^{}G`$ is another (local) discrete Lagrangian evolution operator then $`\xi _{L|U_0U_0^{}}^{}:U_0U_0^{}G`$ is also a (local) discrete Lagrangian evolution operator. Consequently, using Remark 4.11, we conclude that
$$\xi _{L|U_0U_0^{}}^{}=[(𝔽^{}L)^1(𝔽^+L)]_{|U_0U_0^{}}=\xi _{L|U_0U_0^{}}.$$
###### Remark 4.14.
Using Remark 4.3, we deduce that the Lagrangian $`L`$ is regular if and only if for every $`gG`$ and every local basis $`\{X_i\}`$ (respectively, $`\{Y_j\}`$) of $`\mathrm{\Gamma }(\tau )`$ on an open subset $`U`$ (respectively, $`V`$) of $`M`$ such that $`\alpha (g)U`$ (respectively, $`\beta (g)V`$) we have that the matrix $`\stackrel{}{X_i}(\stackrel{}{Y_j}(L))`$ is regular on $`\alpha ^1(U)\beta ^1(V)`$. $``$
Let $`L:G`$ be a regular discrete Lagrangian on $`G`$. If $`f:G`$ is a real $`C^{\mathrm{}}`$-function on $`G`$ then, using Theorem 4.13, it follows that there exists a unique $`\xi _f\mathrm{\Gamma }(\pi ^\tau )`$ such that
$$i_{\xi _f}\mathrm{\Omega }_L=df,$$
(4.21)
$`d`$ being the differential of the Lie algebroid $`\pi ^\tau :𝒫^\tau GV\beta _GV\alpha G`$. $`\xi _f`$ is called the Hamiltonian section associated to $`\mathrm{f}`$ with respect to $`\mathrm{\Omega }_\mathrm{L}`$.
Now, one may introduce a bracket of real functions on $`G`$ as follows:
$$\{,\}_L:C^{\mathrm{}}(G)\times C^{\mathrm{}}(G)C^{\mathrm{}}(G),\{f,g\}_L=\mathrm{\Omega }_L(\xi _f,\xi _g).$$
(4.22)
Note that, from (4.21) and Propositions 4.9 and 4.10, we obtain that
$$(𝒫^\tau 𝔽^\pm L)\xi _{\overline{f}𝔽^\pm L}=\xi _{\overline{f}}𝔽^\pm L,$$
(4.23)
for $`\overline{f}C^{\mathrm{}}(A^{}G),`$ where $`𝒫^\tau 𝔽^\pm L:𝒫^\tau GV\beta _GV\alpha 𝒫^\tau ^{}(AG)`$ is the prolongation of $`𝔽^\pm L`$ (see Section 4.6) and $`\xi _{\overline{f}}`$ is the Hamiltonian section associated to the real function $`\overline{f}`$ on $`A^{}G`$ with respect to the canonical symplectic section $`\mathrm{\Omega }`$ on $`^2(𝒫^\tau ^{}(AG))^{}A^{}G,`$ that is, $`i_{\xi _{\overline{f}}}\mathrm{\Omega }=d\overline{f}.`$
On the other hand, we consider the canonical linear Poisson bracket $`\{,\}:C^{\mathrm{}}(A^{}G)\times C^{\mathrm{}}(A^{}G)C^{\mathrm{}}(A^{}G)`$ on $`A^{}G`$ defined by (see )
$$\{\overline{f},\overline{g}\}=\mathrm{\Omega }(\xi _{\overline{f}},\xi _{\overline{g}}),\text{ for }\overline{f},\overline{g}C^{\mathrm{}}(A^{}G).$$
(4.24)
We have that (see )
$$[[\xi _{\overline{f}},\xi _{\overline{g}}]]^{\tau ^\tau ^{}}=\xi _{\{\overline{f},\overline{g}\}}.$$
Moreover, from (4.22), (4.23), (4.24) and Propositions 4.9 and 4.10, we deduce that
$$\{\overline{f}𝔽^\pm L,\overline{g}𝔽^\pm L\}_L=\{\overline{f},\overline{g}\}𝔽^\pm L.$$
Using the above facts, we may prove the following result.
###### Proposition 4.15.
Let $`L:G`$ be a regular discrete Lagrangian.
1. The Hamiltonian sections with respect to $`\mathrm{\Omega }_L`$ form a Lie subalgebra of the Lie algebra $`(\mathrm{\Gamma }(\pi ^\tau ),[[,]]^{𝒫^\tau G}).`$
2. The Lie groupoid $`G`$ endowed with the bracket $`\{,\}_L`$ is a Poisson manifold, that is, $`\{,\}_L`$ is skew-symmetric, it is a derivation in each argument with respect to the usual product of functions and it satisfies the Jacobi identity.
3. The Legendre transformations $`𝔽^\pm L:GA^{}G`$ are local Poisson isomorphisms.
### 4.8. Discrete Hamiltonian evolution operator
Let $`L:G`$ be a regular Lagrangian and assume, without the loss of generality, that the Legendre transformations $`𝔽^+L`$ and $`𝔽^{}L`$ are global diffeomorphisms. Then, $`\xi _L=(𝔽^{}L)^1(𝔽^+L)`$ is the discrete Euler-Lagrange evolution operator and one may define the discrete Hamiltonian evolution operator, $`\stackrel{~}{\xi }_L:A^{}GA^{}G`$, by
$$\stackrel{~}{\xi }_L=𝔽^+L\xi _L(𝔽^+L)^1.$$
(4.25)
From Remark 4.11, we have the following alternative definitions
$$\stackrel{~}{\xi }_L=𝔽^{}L\xi _L(𝔽^{}L)^1,\stackrel{~}{\xi }_L=𝔽^+L(𝔽^{}L)^1$$
of the discrete Hamiltonian evolution operator. The following commutative diagram illustrates the situation
Define the prolongation $`𝒫^\tau ^{}\stackrel{~}{\xi }_L:𝒫^\tau ^{}(AG)𝒫^\tau ^{}(AG)`$ of $`\stackrel{~}{\xi }_L`$ by
$$𝒫^\tau ^{}\stackrel{~}{\xi }_L=𝒫^\tau 𝔽^+L𝒫^\tau \xi _L(𝒫^\tau 𝔽^+L)^1,$$
or, alternatively (see (4.10), (4.17) and (4.19)),
$$𝒫^\tau ^{}\stackrel{~}{\xi }_L=𝒫^\tau 𝔽^+L(𝒫^\tau 𝔽^{}L)^1,𝒫^\tau ^{}\stackrel{~}{\xi }_L=𝒫^\tau 𝔽^{}L𝒫^\tau \xi _L(𝒫^\tau 𝔽^{}L)^1.$$
(4.26)
###### Proposition 4.16.
If $`\mathrm{\Theta }`$ is the Liouville section of the vector bundle $`(𝒫^\tau ^{}(AG))^{}A^{}G`$ and $`\mathrm{\Omega }=d\mathrm{\Theta }`$ is the canonical symplectic section of $`^2(𝒫^\tau ^{}(AG))^{}A^{}G`$ then
$$(𝒫^\tau ^{}\stackrel{~}{\xi }_L,\stackrel{~}{\xi }_L)^{}\mathrm{\Theta }=\mathrm{\Theta }+d(L(𝔽^{}L)^1),(𝒫^\tau ^{}\stackrel{~}{\xi }_L,\stackrel{~}{\xi }_L)^{}\mathrm{\Omega }=\mathrm{\Omega }.$$
Moreover, $`\stackrel{~}{\xi }_L`$ is a Poisson morphism for the canonical Poisson bracket on $`A^{}G`$.
###### Proof.
The result follows using (4.25), (4.26) and Theorem 4.4 and Propositions 4.9 and 4.15. ∎
### 4.9. Noether’s theorem
Recall that classical Noether’s theorem states that a continuous symmetry of a Lagrangian leads to constants of the motion. In this section, we prove a discrete version of Noether’s theorem, i.e., a theorem relating invariance of the discrete Lagrangian under some transformation with the existence of constants of the motion.
###### Definition 4.17.
A section $`X`$ of $`AG`$ is said to be a Noether’s symmetry of the Lagrangian $`L`$ if there exists a function $`fC^{\mathrm{}}(M)`$ such that
$$dL(X^{(1,1)})=\beta ^{}f\alpha ^{}f.$$
In this case, $`L`$ is said to be quasi-invariant under $`X.`$
When $`dL(X^{(1,1)})=\stackrel{}{X}L+\stackrel{}{X}L=0`$, we will say that $`L`$ is invariant under $`X`$ or that $`X`$ is an infinitesimal symmetry of the discrete Lagrangian $`L`$.
###### Remark 4.18.
The infinitesimal invariance of the Lagrangian corresponds to a finite invariance property as follows. Let $`\mathrm{\Phi }_s`$ the flow of $`\stackrel{}{X}`$ and $`\gamma (s)=\mathrm{\Phi }_s(ϵ(x))`$ be its integral curve with $`\gamma (0)=ϵ(x),`$ where $`x=\beta (g)`$. Then, the integral curve of $`\stackrel{}{X}`$ at $`g`$ is $`sr_{\gamma (s)}g=g\gamma (s)`$, and the integral curve of $`\stackrel{}{X}`$ through $`ϵ(x)`$ is $`s\gamma (s)^1`$. On the other hand, if $`(h,h^{})G_2`$ and $`Y_hV_h\beta ,`$ $`Z_h^{}V_h^{}\alpha `$ then
$$(T_{(h,h^{})}m)(Y_h,Z_h^{})=(T_hr_h^{})(Y_h)+(T_h^{}l_h)(Z_h^{}).$$
Using the above facts, we deduce that the integral curve $`\mu `$ of the vector field $`\stackrel{}{X}+\stackrel{}{X}`$ on $`G`$ satisfying $`\mu (0)=g`$ is
$$\mu (s)=\gamma (s)^1g\gamma (s),\text{ for all }s.$$
Thus, the invariance of the Lagrangian may be written as
$$L\left(\gamma (s)^1g\gamma (s)\right)=L(g),\text{ for all }s.$$
$``$
If $`L:G`$ is a regular discrete Lagrangian, by a constant of the motion we mean a function $`F`$ invariant under the discrete Euler-Lagrange evolution operator $`\xi _L`$, that is, $`F\xi _L=F`$.
###### Theorem 4.19 (Discrete Noether’s theorem).
If $`X`$ is a Noether symmetry of a discrete Lagrangian $`L`$, then the function $`F=\mathrm{\Theta }_L^{}(X^{(1,1)})\alpha ^{}f`$ is a constant of the motion for the discrete dynamics defined by $`L`$.
###### Proof.
We first notice that $`\mathrm{\Theta }_L^{}(X^{(1,1)})=\stackrel{}{X}L`$ so that the function $`F`$ is $`F=\stackrel{}{X}L\alpha ^{}f`$.
If the Lagrangian $`L`$ is quasi-invariant under $`X`$ and $`g`$ is a point in $`G`$, then
$$\stackrel{}{X}(g)(L)+\stackrel{}{X}(g)(L)=f(\beta (g))f(\alpha (g)),$$
so that
$$\stackrel{}{X}(g)(L)=\stackrel{}{X}(g)(L)+f(\beta (g))f(\alpha (g)).$$
We substrate $`\stackrel{}{X}(\xi _L(g))(L)`$ to both sides of the above expression, so that
$`\stackrel{}{X}(g)(L)\stackrel{}{X}(\xi _L(g))(L)`$ $`=[\stackrel{}{X}(g)(L)f(\alpha (g))][\stackrel{}{X}(\xi (g))(L)f(\alpha (\xi _L(g))]`$
$`=F(g)F(\xi _L(g)),`$
from where the result immediately follows using (4.9). ∎
###### Proposition 4.20.
If $`X`$ is a Noether symmetry of the discrete Lagrangian $`L`$ then
$$_{X^{(1,1)}}\mathrm{\Theta }_L^{}=d(\alpha ^{}f).$$
(4.27)
Thus, if $`L`$ is regular, the complete lift $`X^{(1,1)}`$ is a Hamiltonian section with Hamiltonian function $`F=\mathrm{\Theta }_L^{}(X^{(1,1)})\alpha ^{}f`$, i.e. $`i_{X^{(1,1)}}\mathrm{\Omega }_L=dF`$.
###### Proof.
Indeed, if $`dL(X^{(1,1)})=\beta ^{}f\alpha ^{}f`$ and $`Y`$ is a section of $`AG`$, we have that (see Proposition 4.2),
$`(_{X^{(1,1)}}\mathrm{\Theta }_L^{})(Y^{(1,0)})`$ $`=`$ $`\mathrm{\Omega }_L(X^{(1,1)},Y^{(1,0)})+d(i_{X^{(1,1)}}\mathrm{\Theta }_L^{})(Y^{(1,0)})`$
$`=`$ $`\stackrel{}{Y}(\stackrel{}{X}L)+\stackrel{}{Y}(\stackrel{}{X}L)=\stackrel{}{Y}(\alpha ^{}f\beta ^{}f)`$
$`=`$ $`d(\alpha ^{}f)(Y^{(1,0)}).`$
On the other hand, using (2.9) and Proposition 4.2, we deduce that
$`(_{X^{(1,1)}}\mathrm{\Theta }_L^{})(Y^{(0,1)})`$ $`=`$ $`\mathrm{\Omega }_L(X^{(1,1)},Y^{(0,1)})+d(i_{X^{(1,1)}}\mathrm{\Theta }_L^{})(Y^{(0,1)})`$
$`=`$ $`\stackrel{}{X}(\stackrel{}{Y}L)+\stackrel{}{Y}(\stackrel{}{X}L)=[\stackrel{}{Y},\stackrel{}{X}](L)=0=d(\alpha ^{}f)(Y^{(0,1)}).`$
Thus, (4.27) holds. From (4.27), it follows that
$$i_{X^{(1,1)}}\mathrm{\Omega }_L=i_{X^{(1,1)}}d\mathrm{\Theta }_L^{}=di_{X^{(1,1)}}\mathrm{\Theta }_L^{}_{X^{(1,1)}}\mathrm{\Theta }_L^{}=d[\mathrm{\Theta }_L^{}(X^{(1,1)})\alpha ^{}f]=dF,$$
which completes the proof. ∎
We also have
###### Proposition 4.21.
The vector space of Noether symmetries of the Lagrangian $`L:G`$ is a Lie subalgebra of Lie algebra $`(\mathrm{\Gamma }(\tau ),[[,]])`$.
###### Proof.
Suppose that $`X`$ and $`Y`$ are Noether symmetries of $`L`$ and that
$$dL(X^{(1,1)})=\stackrel{}{X}L+\stackrel{}{X}L=\beta ^{}f\alpha ^{}f,$$
(4.28)
$$dL(Y^{(1,1)})=\stackrel{}{Y}L+\stackrel{}{Y}L=\beta ^{}g\alpha ^{}g,$$
(4.29)
with $`f,gC^{\mathrm{}}(M)`$. Then, using (2.9), (3.2) and (3.3), we have that
$$dL([[X,Y]]^{(1,1)})=\stackrel{}{X}(\stackrel{}{Y}L)\stackrel{}{Y}(\stackrel{}{X}L)+\stackrel{}{X}(\stackrel{}{Y}L)\stackrel{}{Y}(\stackrel{}{X}L).$$
(4.30)
On the other hand, from (2.6), (2.7), (4.28) and (4.29), we deduce that
$$\begin{array}{ccc}\stackrel{}{X}(\stackrel{}{Y}L)=\stackrel{}{X}(\stackrel{}{Y}L)\alpha ^{}(\rho (X)(g)),\hfill & & \stackrel{}{Y}(\stackrel{}{X}L)=\stackrel{}{Y}(\stackrel{}{X}L)\alpha ^{}(\rho (Y)(f)),\hfill \\ \stackrel{}{X}(\stackrel{}{Y}L)=\stackrel{}{X}(\stackrel{}{Y}L)+\beta ^{}(\rho (X)(g)),\hfill & & \stackrel{}{Y}(\stackrel{}{X}L)=\stackrel{}{Y}(\stackrel{}{X}L)+\beta ^{}(\rho (Y)(f)).\hfill \end{array}$$
Thus, using (2.9) and (4.30), we obtain that
$$dL([[X,Y]]^{(1,1)})=\beta ^{}(\rho (X)(g)\rho (Y)(f))\alpha ^{}(\rho (X)(g)\rho (Y)(f)).$$
Therefore, $`[[X,Y]]`$ is a Noether symmetry of $`L.`$
###### Remark 4.22.
If $`L:G`$ is a regular discrete Lagrangian then, from Propositions 4.20 and 4.21, it follows that the complete lifts of Noether symmetries of $`L`$ are a Lie subalgebra of the Lie algebra of Hamiltonian sections with respect to $`\mathrm{\Omega }_L`$. $``$
## 5. Examples
### 5.1. Pair or Banal groupoid
We consider the pair (banal) groupoid $`G=M\times M`$, where the structural maps are
$`\alpha (x,y)=x,\beta (x,y)=y,ϵ(x)=(x,x),i(x,y)=(y,x),`$
$`m((x,y),(y,z))=(x,z).`$
We know that the Lie algebroid of $`G`$ is isomorphic to the standard Lie algebroid $`\tau _M:TMM`$ and the map
$$\mathrm{\Psi }:AG=V_{ϵ(M)}\alpha TM,(0_x,v_x)T_xM\times T_xM\mathrm{\Psi }_x(0_x,v_x)=v_x,\text{ for }xM\text{,}$$
induces an isomorphism (over the identity of $`M`$) between $`AG`$ and $`TM`$. If $`X`$ is a section of $`\tau _M:AGTMM`$, that is, $`X`$ is a vector field on $`M`$ then $`\stackrel{}{X}`$ and $`\stackrel{}{X}`$ are the vector fields on $`M\times M`$ given by
$$\stackrel{}{X}(x,y)=(X(x),0_y)T_xM\times T_yM\text{and}\stackrel{}{X}(x,y)=(0_x,X(y))T_xM\times T_yM,$$
for $`(x,y)M\times M`$. On the other hand, if $`(x,y)M\times M`$ we have that the map
$$\begin{array}{cccc}& 𝒫_{(x,y)}^{\tau _M}GV_{(x,y)}\beta V_{(x,y)}\alpha & & T_{(x,y)}(M\times M)T_xM\times T_yM,\\ & ((v_x,0_y),(0_x,v_y))& & (v_x,v_y)\end{array}$$
induces an isomorphism (over the identity of $`M\times M`$) between the Lie algebroids $`\pi ^{\tau _M}:𝒫^{\tau _M}GV\beta _GV\alpha G=M\times M`$ and $`\tau _{(M\times M)}:T(M\times M)M\times M`$.
Now, given a discrete Lagrangian $`L:M\times M`$ then the discrete Euler-lagrange equations for $`L`$ are:
$$\stackrel{}{X}(x,y)(L)\stackrel{}{X}(y,z)(L)=0,\text{ for all }X𝔛(M),$$
(5.1)
which are equivalent to the classical discrete Euler-Lagrange equations
$$D_2L(x,y)+D_1L(y,z)=0$$
(see, for instance, ). The Poincaré-Cartan $`1`$-sections $`\mathrm{\Theta }_L^{}`$ and $`\mathrm{\Theta }_L^+`$ on $`\pi ^{\tau _M}:𝒫^{\tau _M}GT(M\times M)G=M\times M`$ are the $`1`$-forms on $`M\times M`$ defined by
$$\mathrm{\Theta }_L^{}(x,y)(v_x,v_y)=v_x(L),\mathrm{\Theta }_L^+(x,y)(v_x,v_y)=v_y(L),$$
for $`(x,y)M\times M`$ and $`(v_x,v_y)T_xM\times T_yMT_{(x,y)}(M\times M)`$.
In addition, if $`\xi :G=M\times MG=M\times M`$ is a discrete Lagrangian evolution operator then the prolongation of $`\xi `$
$$𝒫^{\tau _M}\xi :𝒫^{\tau _M}GT(M\times M)𝒫^{\tau _M}GT(M\times M)$$
is just the tangent map to $`\xi `$ and, thus, we have that
$$\xi ^{}\mathrm{\Omega }_L=\mathrm{\Omega }_L,$$
$`\mathrm{\Omega }_L=d\mathrm{\Theta }_L^{}=d\mathrm{\Theta }_L^+`$ being the Poincaré-Cartan $`2`$-form on $`M\times M`$. The Legendre transformations $`𝔽^{}L:G=M\times MA^{}GT^{}M`$ and $`𝔽^+L:G=M\times MA^{}GT^{}M`$ associated with $`L`$ are the maps given by
$$𝔽^{}L(x,y)=D_1L(x,y)T_x^{}M,𝔽^+L(x,y)=D_2L(x,y)T_y^{}M$$
for $`(x,y)M\times M`$. The Lagrangian $`L`$ is regular if and only if the matrix $`\left({\displaystyle \frac{^2L}{xy}}\right)`$ is regular. Finally, a Noether symmetry is a vector field $`X`$ on $`M`$ such that
$$D_1L(x,y)(X(x))+D_2L(x,y)(X(y))=f(y)f(x),$$
for $`(x,y)M\times M`$, where $`f:M`$ is a real $`C^{\mathrm{}}`$-function on $`M`$. If $`X`$ is a Noether symmetry then
$$xF(x)=D_1L(x,y)(X(x))f(x)$$
is a constant of the motion.
In conclusion, we recover all the geometrical formulation of the classical discrete Mechanics on the discrete state space $`M\times M`$ (see, for instance, ).
### 5.2. Lie groups
We consider a Lie group $`G`$ as a groupoid over one point $`M=\{𝔢\}`$, the identity element of $`G`$. The structural maps are
$$\alpha (g)=𝔢,\beta (g)=𝔢,ϵ(𝔢)=𝔢,i(g)=g^1,m(g,h)=gh,\text{ for }g,hG.$$
The Lie algebroid associated with $`G`$ is just the Lie algebra $`𝔤=T_𝔢G`$ of $`G`$. Given $`\xi 𝔤`$ we have the left and right invariant vector fields:
$$\stackrel{}{\xi }(g)=(T_𝔢l_g)(\xi ),\stackrel{}{\xi }(g)=(T_𝔢r_g)(\xi ),\text{ for }gG.$$
Thus, given a Lagrangian $`L:G`$ its discrete Euler-Lagrange equations are:
$$(T_𝔢l_{g_k})(\xi )(L)(T_𝔢r_{g_{k+1}})(\xi )(L)=0,\text{ for all }\xi 𝔤\text{ and }g_k,g_{k+1}G,$$
or, $`(l_{g_k}^{}dL)(𝔢)=(r_{g_{k+1}}^{}dL)(𝔢)`$. Denote by $`\mu _k=(r_{g_k}^{}dL)(𝔢)`$ then the discrete Euler-Lagrange equations are written as
$$\mu _{k+1}=Ad_{g_k}^{}\mu _k,$$
(5.2)
where $`Ad:G\times 𝔤𝔤`$ is the adjoint action of $`G`$ on $`𝔤`$. These equations are known as the discrete Lie-Poisson equations (see ).
Finally, an infinitesimal symmetry of $`L`$ is an element $`\xi 𝔤`$ such that $`(T_𝔢l_g)(\xi )(L)=(T_𝔢r_g)(\xi )(L)`$, and then the associated constant of the motion is $`F(g)=(T_𝔢l_g)(\xi )(L)=(T_𝔢r_g)(\xi )(L)`$. Observe that all the Noether’s symmetries are infinitesimal symmetries of $`L`$.
### 5.3. Transformation or action Lie groupoid
Let $`H`$ be a Lie group and $`:M\times HM`$, $`(x,h)M\times Hxh,`$ a right action of $`H`$ on $`M`$. As we know, $`H`$ is a Lie groupoid over the identity element $`𝔢`$ of $`H`$ and we will denote by $`\alpha ,\beta ,ϵ,m`$ and $`i`$ the structural maps of $`H`$. If $`\pi :M\{𝔢\}`$ is the constant map then is clear that the space
$$M_\pi \times _\alpha H=\{(x,h)M\times H/\pi (x)=\alpha (h)\}$$
is the cartesian product $`G=M\times H`$ and that $`:M\times HM`$ induces an action of the Lie groupoid $`H`$ over the map $`\pi :M\{𝔢\}`$ in the sense of Section 2.2 (see Example $`6`$ in Section 2.2). Thus, we may consider the action Lie groupoid $`G=M\times H`$ over $`M`$ with structural maps given by
$$\begin{array}{c}\stackrel{~}{\alpha }_\pi (x,h)=x,\stackrel{~}{\beta }_\pi (x,h)=xh,\stackrel{~}{ϵ}_\pi (x)=(x,𝔢),\hfill \\ \stackrel{~}{m}_\pi ((x,h),(xh,h^{}))=(x,hh^{}),\stackrel{~}{i}_\pi (x,h)=(xh,h^1).\hfill \end{array}$$
(5.3)
Now, let $`𝔥=T_𝔢H`$ be the Lie algebra of $`H`$ and $`\mathrm{\Phi }:𝔥𝔛(M)`$ the map given by
$$\mathrm{\Phi }(\eta )=\eta _M,\text{for }\eta 𝔥,$$
where $`\eta _M`$ is the infinitesimal generator of the action $`:M\times HM`$ corresponding to $`\eta `$. Then, $`\mathrm{\Phi }`$ defines an action of the Lie algebroid $`𝔥\{\text{a point}\}`$ over the projection $`\pi :M\{\text{a point}\}`$ and the corresponding action Lie algebroid $`pr_1:M\times 𝔥M`$ is just the Lie algebroid of $`G=M\times H`$ (see Example $`6`$ in Section 2.2).
We have that $`\mathrm{\Gamma }(pr_1)\{\stackrel{~}{\eta }:M𝔥/\stackrel{~}{\eta }\text{ is smooth }\}`$ and that the Lie algebroid structure $`([[,]]_\mathrm{\Phi },\rho _\mathrm{\Phi })`$ on $`pr_1:M\times HM`$ is given by
$$[[\stackrel{~}{\eta },\stackrel{~}{\mu }]]_\mathrm{\Phi }(x)=[\stackrel{~}{\eta }(x),\stackrel{~}{\mu }(x)]+(\stackrel{~}{\eta }(x))_M(x)(\stackrel{~}{\mu })(\stackrel{~}{\mu }(x))_M(x)(\stackrel{~}{\eta }),\rho _\mathrm{\Phi }(\stackrel{~}{\eta })(x)=(\stackrel{~}{\eta }(x))_M(x),$$
for $`\stackrel{~}{\eta },\stackrel{~}{\mu }\mathrm{\Gamma }(pr_1)`$ and $`xM.`$ Here, $`[,]`$ denotes the Lie bracket of $`𝔥`$.
If $`(x,h)G=M\times H`$ then the left-translation $`l_{(x,h)}:\stackrel{~}{\alpha }_\pi ^1(xh)\stackrel{~}{\alpha }_\pi ^1(x)`$ and the right-translation $`r_{(x,h)}:\stackrel{~}{\beta }_\pi ^1(x)\stackrel{~}{\beta }_\pi ^1(xh)`$ are given
$$l_{(x,h)}(xh,h^{})=(x,hh^{}),r_{(x,h)}(x(h^{})^1,h^{})=(x(h^{})^1,h^{}h).$$
(5.4)
Now, if $`\eta 𝔥`$ then $`\eta `$ defines a constant section $`C_\eta :M𝔥`$ of $`pr_1:M\times 𝔥M`$ and, using (2.6), (2.7), (5.3) and (5.4), we have that the left-invariant and the right-invariant vector fields $`\stackrel{}{C}_\eta `$ and $`\stackrel{}{C}_\eta `$, respectively, on $`M\times H`$ are defined by
$$\stackrel{}{C}_\eta (x,h)=(\eta _M(x),\stackrel{}{\eta }(h)),\stackrel{}{C}_\eta (x,h)=(0_x,\stackrel{}{\eta }(h)),$$
(5.5)
for $`(x,h)G=M\times H.`$
Note that if $`\{\eta _i\}`$ is a basis of $`𝔥`$ then $`\{C_{\eta _i}\}`$ is a global basis of $`\mathrm{\Gamma }(pr_1).`$
Next, suppose that $`L:G=M\times H`$ is a Lagrangian function and for every $`hH`$ (resp., $`xM`$) we will denote by $`L_h`$ (resp., $`L_x`$) the real function on $`M`$ (resp., on $`H`$) given by $`L_h(y)=L(y,h)`$ (resp., $`L_x(h^{})=L(x,h^{}))`$. Then, a composable pair $`((x,h_k),(xh_k,h_{k+1}))G_2`$ is a solution of the discrete Euler-Lagrange equations for $`L`$ if
$$\stackrel{}{C}_\eta (x,h_k)(L)\stackrel{}{C}_\eta (xh_k,h_{k+1})(L)=0,\text{ for all }\eta 𝔥,$$
or, in other terms (see (5.5))
$$\{(T_𝔢l_{h_k})(\eta )\}(L_x)\{(T_𝔢r_{h_{k+1}})(\eta )\}(L_{xh_k})+\eta _M(xh_k)(L_{h_{k+1}})=0,\text{ for all }\eta 𝔥.$$
As in the case of Lie groups, denote by $`\mu _k(x,h_k)=d(L_xr_{h_k})(𝔢).`$ Then, the discrete Euler-Lagrange equations for $`L`$ are written as
$$\mu _{k+1}(xh_k,h_{k+1})=Ad_{h_k}^{}\mu _k(x,h_k)+d(L_{h_{k+1}}((xh_k)))(e),$$
where $`(xh_k):HM`$ is the map defined by
$$(xh_k)(h)=x(h_kh),\text{ for }hH.$$
In the particular case when $`M`$ is the orbit of $`aV`$ under a representation of $`G`$ on a real vector space $`V`$, the resultant equations were obtained by Bobenko and Suris, see , and they were called the discrete Euler-Poincaré equations.
Finally, an element $`\xi 𝔥`$ is an infinitesimal symmetry of $`L`$ if
$$\xi _M(x)(L_h)\stackrel{}{\xi }(h)(L_x)+\stackrel{}{\xi }(h)(L_x)=f(xh)f(x)$$
where $`f:M`$ is a real $`C^{\mathrm{}}`$-function on $`M`$. The associated constant of the motion is
$$F(x,h)=\xi _M(x)(L_h)+\stackrel{}{\xi }(h)(L_x)f(x),$$
for $`(x,h)M\times H`$.
#### The heavy top
As a concrete example of a system on a transformation Lie groupoid we consider a discretization of the heavy top. In the continuous theory , the configuration manifold is the transformation Lie algebroid $`\tau :S^2\times 𝔰𝔬(3)S^2`$ with Lagrangian
$$L_c(\mathrm{\Gamma },\mathrm{\Omega })=\frac{1}{2}\mathrm{\Omega }I\mathrm{\Omega }mgl\mathrm{\Gamma }\mathrm{e},$$
where $`\mathrm{\Omega }^3𝔰𝔬(3)`$ is the angular velocity, $`\mathrm{\Gamma }`$ is the direction opposite to the gravity and $`\mathrm{e}`$ is a unit vector in the direction from the fixed point to the center of mass, all them expressed in a frame fixed to the body. The constants $`m`$, $`g`$ and $`l`$ are respectively the mass of the body, the strength of the gravitational acceleration and the distance from the fixed point to the center of mass. The matrix $`I`$ is the inertia tensor of the body. In order to discretize this Lagrangian it is better to express it in terms of the matrices $`\widehat{\mathrm{\Omega }}𝔰𝔬(3)`$ such that $`\widehat{\mathrm{\Omega }}v=\mathrm{\Omega }\times v`$. Then
$$L_c(\mathrm{\Gamma },\mathrm{\Omega })=\frac{1}{2}\mathrm{Tr}(\widehat{\mathrm{\Omega }}II\widehat{\mathrm{\Omega }}^T)mgl\mathrm{\Gamma }\mathrm{e}.$$
where $`II=\frac{1}{2}\mathrm{Tr}(I)I_3I`$. We can define a discrete Lagrangian $`L:G=S^2\times SO(3)`$ for the heavy top by
$$L(\mathrm{\Gamma }_k,W_k)=\frac{1}{h}\mathrm{Tr}(IIW_k)hmgl\mathrm{\Gamma }_k\mathrm{e}.$$
which is obtained by the rule $`\widehat{\mathrm{\Omega }}=R^T\dot{R}\frac{1}{h}R_k^T(R_{k+1}R_k)=\frac{1}{h}(W_kI_3)`$, where $`W_k=R_k^TR_{k+1}`$.
The value of the action on an admissible variation is
$`\lambda (t)`$ $`=L(\mathrm{\Gamma }_k,W_ke^{tK})+L(e^{tK}\mathrm{\Gamma }_{k+1},e^{tK}W_{k+1})`$
$`={\displaystyle \frac{1}{h}}\left[\mathrm{Tr}(IIW_ke^{tK})+mglh^2\mathrm{\Gamma }_k\mathrm{e}+\mathrm{Tr}(IIe^{tK}W_{k+1})+mglh^2(e^{tK}\mathrm{\Gamma }_{k+1})\mathrm{e}\right],`$
where $`\mathrm{\Gamma }_{k+1}=W^T\mathrm{\Gamma }_k`$ (since the above pairs must be composable) and $`K𝔰𝔬(3)`$ is arbitrary. Taking the derivative at $`t=0`$ and after some straightforward manipulations we get the DEL equations
$$M_{k+1}W_k^TM_kW_k+mglh^2(\widehat{\mathrm{\Gamma }_{k+1}\times \mathrm{e}})=0$$
where $`M=WIIIIW^T`$. Finally, in terms of the axial vector $`\mathrm{\Pi }`$ in $`^3`$ defined by $`\widehat{\mathrm{\Pi }}=M`$, we can write the equations in the form
$$\mathrm{\Pi }_{k+1}=W_k^T\mathrm{\Pi }_k+mglh^2\mathrm{\Gamma }_{k+1}\times \mathrm{e}.$$
###### Remark 5.1.
The above equations are to be solved as follows. From $`\mathrm{\Gamma }_k,W_k`$ we obtain $`\mathrm{\Gamma }_{k+1}=W_k\mathrm{\Gamma }_k`$ and $`\mathrm{\Pi }_k`$ from $`\widehat{\mathrm{\Pi }}_k=W_kIIIIW_k^T`$. The DEL equation gives $`\mathrm{\Pi }_{k+1}`$ in terms of the above data. Finally we get $`W_{k+1}`$ as the solution of the equation $`\widehat{\mathrm{\Pi }}_{k+1}=W_{k+1}IIIIW_{k+1}^T`$, as in . $``$
In the continuous theory, the section $`X(\mathrm{\Gamma })=(\mathrm{\Gamma },\mathrm{\Gamma })`$ of $`S^2\times 𝔰𝔬(3)S^2`$ is a symmetry of the Lagrangian (see ). We will show next that such a section is also a symmetry of the discrete Lagrangian. Indeed, it is easy to see that the left and right vector fields associated to $`X`$ coincide $`\stackrel{}{X}=\stackrel{}{X}`$ and are both equal to
$$\stackrel{}{X}(\mathrm{\Gamma },W)=((\mathrm{\Gamma },0),(W,\widehat{\mathrm{\Gamma }}W)))TG=TS^2\times TSO(3).$$
Thus $`\rho ^{𝒫^\tau G}(X^{(1,1)})=0`$ so that $`X`$ is a symmetry of the Lagrangian. In fact it is a symmetry of any discrete Lagrangian defined on $`G=S^2\times SO(3)`$. The associated constant of motion is
$$(\stackrel{}{X}L)(W,\mathrm{\Gamma })=\mathrm{Tr}(II\widehat{\mathrm{\Gamma }}W)=\frac{1}{2}\mathrm{Tr}[(WIIIIW^T)\widehat{\mathrm{\Gamma }}]=\mathrm{\Pi }\mathrm{\Gamma },$$
i.e. (minus) the angular momentum in the direction of the vector $`\mathrm{\Gamma }`$.
### 5.4. Atiyah or gauge groupoids
Let $`p:QM`$ be a principal $`G`$-bundle. A discrete connection on $`p:QM`$ is a map $`𝒜_d:Q\times QG`$ such that
$$𝒜_d(gq,hq^{})=h𝒜_d(q,q^{})g^1\text{and}𝒜_d(q,q)=𝔢$$
(5.6)
for $`g,hG`$ and $`q,q^{}Q`$, $`𝔢`$ being the identity in the group $`G`$ (see ). We remark that a discrete principal connection may be considered as the discrete version of an standard (continuous) connection on $`p:QM`$. In fact, if $`𝒜_d:Q\times QG`$ is such a connection then it induces, in a natural way, a continuous connection $`𝒜_c:TQ𝔤`$ defined by
$$𝒜_c(v_q)=(T_{(q,q)}𝒜_d)(0_q,v_q),$$
for $`v_qT_qQ`$. Moreover, if we choose a local trivialization of the principal bundle $`p:QM`$ to be $`G\times U`$, where $`U`$ is an open subset of $`M`$ then, from (5.6), it follows that there exists a map $`A:U\times UG`$ such that
$$𝒜_d((g,x),(g^{},y))=g^{}A(x,y)g^1,\text{and}A(x,x)=𝔢,$$
for $`(g,x),(g^{},x^{})G\times U`$ (for more details, see ).
On the other hand, using the discrete connection $`𝒜_d`$, one may identify the open subset $`(p^1(U)\times p^1(U))/G((G\times U)\times (G\times U))/G`$ of the Atiyah groupoid $`(Q\times Q)/G`$ with the product manifold $`(U\times U)\times G`$. Indeed, it is easy to prove that the map
$$((G\times U)\times (G\times U))/G(U\times U)\times G,$$
$$[((g,x),(g^{},y))]((x,y),𝒜_d((e,x),(g^1g^{},y)))=((x,y),g^1g^{}A(x,y)),$$
is bijective. Thus, the restriction to $`((G\times U)\times (G\times U))/G`$ of the Lie groupoid structure on $`(Q\times Q)/G`$ induces a Lie groupoid structure in $`(U\times U)\times G`$ with source, target and identity section given by
$$\begin{array}{cc}\alpha :(U\times U)\times GU;\hfill & \hfill ((x,y),g)x,\\ \beta :(U\times U)\times GU;\hfill & \hfill ((x,y),g)y,\\ ϵ:U(U\times U)\times G;\hfill & \hfill x((x,x),𝔢),\end{array}$$
and with multiplication $`m:((U\times U)\times G)_2(U\times U)\times G`$ and inversion $`i:(U\times U)\times G(U\times U)\times G`$ defined by
$$\begin{array}{ccc}\hfill m(((x,y),g),((y,z),h))& =& ((x,z),gA(x,y)^1hA(y,z)^1A(x,z)),\hfill \\ \hfill i((x,y),g)& =& ((y,x),A(x,y)g^1A(y,x)).\hfill \end{array}$$
(5.7)
The fibre over the point $`xU`$ of the Lie algebroid $`A((U\times U)\times G)`$ may be identified with the vector space $`T_xU\times 𝔤`$. Thus, a section of $`A((U\times U)\times G)`$ is a pair $`(X,\stackrel{~}{\xi })`$, where $`X`$ is a vector field on $`U`$ and $`\stackrel{~}{\xi }`$ is a map from $`U`$ on $`𝔤`$. Note that the space $`\mathrm{\Gamma }(A((U\times U)\times G))`$ is generated by sections of the form $`(X,0)`$ and $`(0,C_\xi )`$, with $`X𝔛(U)`$, $`\xi 𝔤`$ and $`C_\xi :U𝔤`$ being the constant map $`C_\xi (x)=\xi `$, for all $`xU`$. Moreover, an straightforward computation, using (5.7), proves that the vector fields $`\stackrel{}{(X,0)}`$, $`\stackrel{}{(X,0)}`$, $`\stackrel{}{(0,C_\xi )}`$ and $`\stackrel{}{(0,C_\xi )}`$ on $`(U\times U)\times G`$ are given by
$`\stackrel{}{(X,0)}((x,y),g)`$ $`=(0_x,X(y),(T_{A(x,y)}l_{gA(x,y)^1}((T_yA_x)(X(y)))+`$ (5.8)
$`(Ad_{A(x,y)^1}(T_yA_y)(X(y)))^l(g))),`$
$`\stackrel{}{(X,0)}((x,y),g)`$ $`=(X(x),0_y,(T_{A(x,y)}l_{gA(x,y)^1}((T_xA_y)(X(x))))+`$
$`(Ad_{A(x,y)^1}(T_xA_x)(X(x)))^l(g))),`$
$`\stackrel{}{(0,C_\xi )}((x,y),g)`$ $`=(0_x,0_y,(Ad_{A(x,y)^1}\xi )^l(g)),`$
$`\stackrel{}{(0,C_\xi )}((x,y),g)`$ $`=(0_x,0_y,\xi ^r(g)),`$
for $`((x,y),g)(U\times U)\times G`$, where $`l_h:GG`$ denotes the left-translation in $`G`$ by $`hG`$, $`Ad:G\times 𝔤𝔤`$ is the adjoint action of the Lie group $`G`$ on $`𝔤`$, $`\eta ^l`$ (respectively, $`\eta ^r`$) is the left-invariant (respectively, right-invariant) vector field on $`G`$ such that $`\eta ^l(𝔢)=\eta `$ (respectively, $`\eta ^r(𝔢)=\eta `$) and $`A_x:UG`$ and $`A_y:UG`$ are the maps defined by
$$A_x(y)=A_y(x)=A(x,y).$$
Now, suppose that $`L:(Q\times Q)/G`$ is a Lagrangian function on the Atiyah groupoid $`(Q\times Q)/G`$. Then, the discrete Euler-Lagrange equations for $`L`$ are
$$\begin{array}{ccc}\stackrel{}{(X,0)}((x,y),g_k)(L)\stackrel{}{(X,0)}((y,z),g_{k+1})(L)\hfill & =& \hfill 0,\\ \stackrel{}{(0,C_\xi )}((x,y),g_k)(L)\stackrel{}{(0,C_\xi )}((y,z),g_{k+1})(L)\hfill & =& \hfill 0,\end{array}$$
with $`X𝔛(U)`$, $`\xi 𝔤`$ and $`(((x,y),g_k),((y,z),g_{k+1}))((U\times U)\times G)_2`$.
From (5.8), it follows that the above equations may be written as
$`D_2L((x,y),g_k)+D_1L((y,z),g_{k+1})+df_{AL}[x,y,g_k](y)+`$
$`+df_{AL}[y,z,g_{k+1}](y)+df_{ALI}^1[x,y,g_k](y)+df_{ALI}^2[y,z,g_k+1](y)=0,`$ (5.9)
$`d(L_{(x,y,)}l_{g_k}I_{A(x,y)^1})(𝔢)d(L_{(y,z,)}r_{g_{k+1}})(𝔢)=0,`$ (5.10)
where $`I_{\overline{g}}:GG`$ denotes the interior automorphism in $`G`$ of $`\overline{g}G`$, $`L_{(\overline{x},\overline{y},)}:G`$ is the function $`L_{(\overline{x},\overline{y},)}(g)=L(\overline{x},\overline{y},g)`$, and $`f_{AL}[\overline{x},\overline{y},\overline{g}]`$, $`f_{ALI}^1[\overline{x},\overline{y},\overline{g}]`$ and $`f_{ALI}^2[\overline{x},\overline{y},\overline{g}]`$ are the real functions on $`U`$ given by
$`f_{AL}[\overline{x},\overline{y},\overline{g}](y)`$ $`=L(\overline{x},\overline{y},\overline{g}A(\overline{x},\overline{y})^1A(\overline{x},y)),`$
$`f_{ALI}^1[\overline{x},\overline{y},\overline{g}](y)`$ $`=L(\overline{x},\overline{y},\overline{g}A(\overline{x},\overline{y})^1A(\overline{y},y)A(\overline{x},\overline{y})),`$
$`f_{ALI}^2[\overline{x},\overline{y},\overline{g}](y)`$ $`=L(\overline{x},\overline{y},\overline{g}A(\overline{x},\overline{y})^1A(y,\overline{y})A(\overline{x},\overline{y})),`$
for $`\overline{x},\overline{y},yU`$ and $`gG`$. These equations may be considered as the discrete version of the Lagrange-Poincaré equations for a $`G`$-invariant continuous Lagrangian (see for the local expression of the Lagrange-Poincaré equations).
Note that if $`A:U\times UG`$ is the constant map $`A(x,y)=𝔢`$, for all $`(x,y)U\times U`$, or, in other words, $`𝒜_d`$ is the trivial connection then equations (5.4) and (5.10) may be written as
$$\begin{array}{c}D_2L((x,y),g_k)+D_1L((y,z),g_{k+1})=0,\hfill \\ \mu _{k+1}(y,z)=Ad_{g_k}^{}\mu _k(x,y),\hfill \end{array}$$
(5.11)
where
$$\mu _k(\overline{x},\overline{y})=d(r_{g_k}^{}L_{(\overline{x},\overline{y},)})(𝔢)$$
for $`(\overline{x},\overline{y})U\times U`$ (compare equations (5.11) with equations (5.1) and (5.2)).
#### Discrete Elroy’s beanie
As an example of a lagrangian system on an Atiyah groupoid, we consider a discretization of the Elroy’s beanie, which is, probably, the most simple example of a dynamical system with a non-Abelian Lie group of symmetries. The continuous system consists in two planar rigid bodies attached at their centers of mass, moving freely in the plane. The configuration space is $`Q=SE(2)\times S^1`$ with coordinates $`(x,y,\theta ,\psi )`$, where the three first coordinates describe the position and orientation of the center of mass of the first body and the last one the relative orientation between both bodies. The continuous system is described by a Lagrangian $`L_c(x,y,\theta ,\psi ,\dot{x},\dot{y},\dot{\theta },\dot{\psi })=\frac{1}{2}m(\dot{x}^2+\dot{y}^2)+\frac{1}{2}I_1\dot{\theta }^2+\frac{1}{2}I_2(\dot{\theta }+\dot{\psi })^2V(\psi )`$ where $`m`$ denotes the mass of the system, $`I_1`$ and $`I_2`$ are the inertias of the first and the second body, respectively, and $`V`$ is the potential energy. The system admits reduction by $`SE(2)`$ symmetry. In fact, the reduced lagrangian $`l_c:TQ/SE(2)S^1\times \times 𝔰𝔢(2)`$ is
$$l_c(\psi ,\dot{\psi },\mathrm{\Omega }_1,\mathrm{\Omega }_2,\mathrm{\Omega }_3)=\frac{1}{2}m(\mathrm{\Omega }_1^2+\mathrm{\Omega }_2^2)+\frac{1}{2}(I_1+I_2)\mathrm{\Omega }_3^2+\frac{1}{2}\frac{I_1I_2}{I_1+I_2}\dot{\psi }^2V(\psi )$$
where $`𝔰𝔢(2)`$ is the Lie algebra of $`SE(2)`$, $`\mathrm{\Omega }_1=\xi _1`$, $`\mathrm{\Omega }_2=\xi _2`$, $`\mathrm{\Omega }_3=\xi _3\frac{I_2}{I_1+I_2}\dot{\psi }`$ and $`(\xi _1,\xi _2,\xi _3)`$ are the coordinates of an element of $`𝔰𝔢(2)`$ with respect to the basis $`e_1=\left(\begin{array}{ccc}0& 0& 1\\ 0& 0& 0\\ 0& 0& 0\end{array}\right)`$, $`e_2=\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& 1\\ 0& 0& 0\end{array}\right)`$ and $`e_3=\left(\begin{array}{ccc}0& 1& 0\\ 1& 0& 0\\ 0& 0& 0\end{array}\right)`$. Note that $`\xi _1=\dot{x}\mathrm{cos}\theta +\dot{y}\mathrm{sin}\theta `$, $`\xi _2=\dot{x}\mathrm{sin}\theta +\dot{y}\mathrm{cos}\theta `$ and $`\xi _3=\dot{\theta }{\displaystyle \frac{I_2}{I_1+I_2}}\dot{\psi }`$ (for more details, see ).
In order to discretize this system, consider $`g_k=\left(\begin{array}{ccc}\mathrm{cos}\theta _k& \mathrm{sin}\theta _k& x_k\\ \mathrm{sin}\theta _k& \mathrm{cos}\theta _k& y_k\\ 0& 0& 1\end{array}\right)SE(2)`$. We construct the discrete connection $`𝒜_d:(SE(2)\times S^1)\times (SE(2)\times S^1)SE(2)`$ defined by $`𝒜_d((g_k,\psi _k),(g_{k+1},\psi _{k+1}))=g_{k+1}A(\psi _k,\psi _{k+1})g_k^1`$, where
$$A(\psi _k,\psi _{k+1})=\left(\begin{array}{ccc}\mathrm{cos}\left(\frac{I_2}{I_1+I_2}\mathrm{\Delta }\psi _k\right)& \mathrm{sin}\left(\frac{I_2}{I_1+I_2}\mathrm{\Delta }\psi _k\right)& 0\\ \mathrm{sin}\left(\frac{I_2}{I_1+I_2}\mathrm{\Delta }\psi _k\right)& \mathrm{cos}\left(\frac{I_2}{I_1+I_2}\mathrm{\Delta }\psi _k\right)& 0\\ 0& 0& 1\end{array}\right)$$
Here, $`\mathrm{\Delta }\psi _k=\psi _{k+1}\psi _k`$. The discrete connection $`𝒜_d`$ precisely induces the mechanical connection associated with the $`SE(2)`$-invariant metric $`𝒢`$ on $`Q`$:
$$𝒢=mdxdx+mdydy+(I_1+I_2)d\theta d\theta +I_2d\theta d\psi +I_2d\psi d\theta +I_2d\psi d\psi $$
We remark that the continuous Lagrangian $`L_c`$ is the kinetic energy associated with $`𝒢`$ minus the potential energy $`V`$.
Next, we consider the Atiyah groupoid $`(Q\times Q)/SE(2)`$. As we know, using the discrete connection $`𝒜_d`$, one may define a local isomorphism between the Atiyah groupoid $`(Q\times Q)/SE(2)`$ and the product manifold $`U\times U\times SE(2)`$, $`U`$ being an open subset of $``$. Then, as a local discretization of the reduced Lagrangian $`l_c`$, we introduce the discrete Lagrangian $`l_d`$ on $`U\times U\times SE(2)`$ given by
$`l_d(\psi _k,\psi _{k+1},\mathrm{\Omega }_{(1)k},\mathrm{\Omega }_{(2)k},\mathrm{\Omega }_{(3)k})={\displaystyle \frac{1}{2h^2}}m\left[\mathrm{\Omega }_{(1)k}^2+\mathrm{\Omega }_{(2)k}^2\right]`$
$`+{\displaystyle \frac{(I_1+I_2)}{h^2}}\left[1\mathrm{cos}(\mathrm{\Omega }_{(3)k})\right]+{\displaystyle \frac{1}{2}}{\displaystyle \frac{I_1I_2}{I_1+I_2}}\left({\displaystyle \frac{\mathrm{\Delta }\psi _k}{h}}\right)^2V({\displaystyle \frac{\psi _k+\psi _{k+1}}{2}})`$
where $`\mathrm{\Omega }_{(1)k}=\mathrm{\Delta }x_k\mathrm{cos}\theta _k+\mathrm{\Delta }y_k\mathrm{sin}\theta _k,`$ $`\mathrm{\Omega }_{(2)k}=\mathrm{\Delta }x_k\mathrm{sin}\theta _k+\mathrm{\Delta }y_k\mathrm{cos}\theta _k`$ and $`\mathrm{\Omega }_{(3)k}=\mathrm{\Delta }\theta _k\frac{I_2}{I_1+I_2}\mathrm{\Delta }\psi _k`$.
Now, if we denote by $`\overline{q}_k=(\psi _k,\psi _{k+1},\mathrm{\Omega }_{(1)k},\mathrm{\Omega }_{(2)k},\mathrm{\Omega }_{(3)k})`$ then
$`\stackrel{}{(0,C_{e_1})}|_{\overline{q}_k}=\mathrm{cos}\left(\mathrm{\Omega }_{\left(3\right)k}+{\displaystyle \frac{I_2}{I_1+I_2}}\mathrm{\Delta }\psi _k\right){\displaystyle \frac{}{\mathrm{\Omega }_{\left(1\right)k}}}\mathrm{sin}\left(\mathrm{\Omega }_{\left(3\right)k}+{\displaystyle \frac{I_2}{I_1+I_2}}\mathrm{\Delta }\psi _k\right){\displaystyle \frac{}{\mathrm{\Omega }_{\left(2\right)k}}}`$
$`\stackrel{}{(0,C_{e_2})}|_{\overline{q}_k}=\mathrm{sin}\left(\mathrm{\Omega }_{\left(3\right)k}+{\displaystyle \frac{I_2}{I_1+I_2}}\mathrm{\Delta }\psi _k\right){\displaystyle \frac{}{\mathrm{\Omega }_{\left(1\right)k}}}+\mathrm{cos}\left(\mathrm{\Omega }_{\left(3\right)k}+{\displaystyle \frac{I_2}{I_1+I_2}}\mathrm{\Delta }\psi _k\right){\displaystyle \frac{}{\mathrm{\Omega }_{\left(2\right)k}}}`$
$`\stackrel{}{(0,C_{e_3})}|_{\overline{q}_k}={\displaystyle \frac{}{\mathrm{\Omega }_{(3)k}}},\stackrel{}{(0,C_{e_1})}|_{\overline{q}_k}={\displaystyle \frac{}{\mathrm{\Omega }_{(1)k}}},\stackrel{}{(0,C_{e_2})}|_{\overline{q}_k}={\displaystyle \frac{}{\mathrm{\Omega }_{(2)k}}}`$
$`\stackrel{}{(0,C_{e_3})}|_{\overline{q}_k}={\displaystyle \frac{}{\mathrm{\Omega }_{\left(3\right)k}}}+\mathrm{\Omega }_{\left(2\right)k}{\displaystyle \frac{}{\mathrm{\Omega }_{\left(1\right)k}}}\mathrm{\Omega }_{\left(1\right)k}{\displaystyle \frac{}{\mathrm{\Omega }_{\left(2\right)k}}},\stackrel{}{({\displaystyle \frac{}{\psi }},0)}|_{\overline{q}_k}={\displaystyle \frac{}{\psi _{k+1}}},`$
$`\stackrel{}{({\displaystyle \frac{}{\psi }},0)}|_{\overline{q}_k}={\displaystyle \frac{}{\psi _k}}+{\displaystyle \frac{I_2}{I_1+I_2}}\mathrm{\Omega }_{\left(2\right)k}{\displaystyle \frac{}{\mathrm{\Omega }_{\left(1\right)k}}}{\displaystyle \frac{I_2}{I_1+I_2}}\mathrm{\Omega }_{\left(1\right)k}{\displaystyle \frac{}{\mathrm{\Omega }_{\left(2\right)k}}}`$
Thus, the reduced Discrete Euler-Lagrange equations
$$\stackrel{}{(0,C_{e_i})}|_{\overline{q}_k}l_d\stackrel{}{(0,C_{e_i})}|_{\overline{q}_{k+1}}l_d=0,\stackrel{}{(\frac{}{\psi },0)}|_{\overline{q}_k}l_d\stackrel{}{(\frac{}{\psi },0)}|_{\overline{q}_{k+1}}l_d=0$$
are
$$\{\begin{array}{ccc}& & \mathrm{\Omega }_{(1)k+1}=\mathrm{\Omega }_{(1)k}\mathrm{cos}(\mathrm{\Omega }_{(3)k}+\frac{I_2}{I_1+I_2}\mathrm{\Delta }\psi _k)\mathrm{\Omega }_{(2)k}\mathrm{sin}(\mathrm{\Omega }_{(3)k}+\frac{I_2}{I_1+I_2}\mathrm{\Delta }\psi _k)\hfill \\ & & \mathrm{\Omega }_{(2)k+1}=\mathrm{\Omega }_{(1)k}\mathrm{sin}(\mathrm{\Omega }_{(3)k}+\frac{I_2}{I_1+I_2}\mathrm{\Delta }\psi _k)+\mathrm{\Omega }_{(2)k}\mathrm{cos}(\mathrm{\Omega }_{(3)k}+\frac{I_2}{I_1+I_2}\mathrm{\Delta }\psi _k)\hfill \\ & & \mathrm{\Omega }_{(3)k+1}=\mathrm{\Omega }_{(3)k}\hfill \\ & & \frac{I_1I_2}{I_1+I_2}\frac{\psi _{k+2}2\psi _{k+1}+\psi _k}{h^2}=\frac{1}{2}\left(\frac{V}{\psi }(\frac{\psi _{k+2}+\psi _{k+1}}{2})+\frac{V}{\psi }(\frac{\psi _{k+1}+\psi _k}{2})\right)\hfill \end{array}$$
These equations are a discretization of the corresponding reduced equations for the continuous system (see ). In a forthcoming paper , we will give a complete description of this example comparing with the continuous equations.
### 5.5. Reduction of discrete Lagrangian systems
Next, we will present some examples of Lie groupoid epimorphisms which allow to do reduction.
$``$ Let $`G`$ be a Lie group and consider the pair groupoid $`G\times G`$ over $`G`$. Consider also $`G`$ as a groupoid over one point. Then we have that the map
$$\begin{array}{cccc}\hfill \mathrm{\Phi }_l:& G\times G& & G\\ & (g,h)& & g^1h\end{array}$$
is a Lie groupoid morphism, which is obviously a submersion. Thus, using Corollary 4.7, it follows that the discrete Euler-Lagrange equations for a left invariant discrete Lagrangian on $`G\times G`$ reduce to the discrete Lie-Poisson equations on $`G`$ for the reduced Lagrangian. This case appears in as was first noticed by , and also appear later in .
Alternatively, one can do reduction of a right-invariant Lagrangian by using the morphism
$$\begin{array}{cccc}\hfill \mathrm{\Phi }_r:& G\times G& & G\\ & (g,h)& & gh^1\end{array}$$
$``$ Let $`G`$ be a Lie group acting on a manifold $`M`$ by the left. We consider a discrete Lagrangian on $`G\times G`$ which depends on the variables of $`M`$ as parameters $`L_m(g,h)`$. In general, the Lagrangian will not be invariant under the action of $`G`$, that is $`L_m(g,h)L_m(rg,rh)`$. Nevertheless, it can happen that $`L_m(rg,rh)=L_{r^1m}(g,h)`$. In such cases we can consider the Lie groupoid $`G\times G\times M`$ over $`G\times M`$ where accordingly one consider the elements in $`M`$ as parameters. Then the Lagrangian can be considered as a function on the groupoid $`G\times G\times M`$ given by $`L(g,h,m)L_m(g,h)`$ so that the above property reads $`L(rg,rh,rm)=L(g,h,m)`$. Thus we define the reduction map
$$\begin{array}{cccc}\hfill \mathrm{\Phi }:& G\times G\times M& & G\times M\\ & (g,h,m)& & (g^1h,g^1m)\end{array}$$
where on $`G\times M`$ we consider the transformation Lie groupoid defined by the right action $`mg=g^1m`$. Since this map is a submersion, the Euler-Lagrange equations on $`G\times G\times M`$ reduces to the Euler-Lagrange equations on $`G\times M`$. This case occurs in the Lagrange top that was considered as an example in Section 5.3 (see also ).
$``$ Another interesting case is that of a $`G`$-invariant Lagrangian $`L`$ defined on the pair groupoid $`L:Q\times Q`$, where $`p:QM`$ is a $`G`$-principal bundle. In this case we can reduce to the Atiyah gauge groupoid by means of the map
$$\begin{array}{cccc}\hfill \mathrm{\Phi }:& Q\times Q& & (Q\times Q)/G\\ & (q,q^{})& & [(q,q^{})]\end{array}$$
Thus the discrete Euler-Lagrange equations reduce to the so called discrete Lagrange-Poincaré equations.
## 6. Conclusions and outlook
In this paper we have elucidated the geometrical framework for discrete Mechanics on Lie groupoids. Using as a main tool the natural Lie algebroid structure on the vector bundle $`\pi ^\tau :𝒫^\tau GG`$ we have found intrinsic expressions for the discrete Euler-Lagrange equations. We introduce the Poincaré-Cartan sections, the discrete Legendre transformations and the discrete evolution operator in the Lagrangian and in the Hamiltonian formalism. The notion of regularity has been completely characterized and we prove the symplecticity of the discrete evolution operators. Moreover, we have studied the symmetries of discrete Lagrangians on Lie groupoids relating them with constants of the motion via Noether’s Theorems. The applicablity of these developments has been stated in several interesting examples, in particular for the case of discrete Lagrange-Poincaré equations. In fact, the general theory of discrete symmetry reduction naturally follows from our results.
In this paper we have confined ourselves to the geometrical aspects of mechanics on Lie groupoids. In a forthcoming paper (see ) we will study the construction of geometric integrators for mechanical systems on Lie algebroids. We will introduce the exact discrete Lagrangian and we will discuss different discretizations of a continuous Lagrangian and its numerical implementation.
Another different aspect we will work on it in the future is to develop natural extensions of the above theories for forced systems and systems with holonomic and nonholonomic constraints. |
warning/0506/quant-ph0506016.html | ar5iv | text | # Measuring the quality factor of a microwave cavity using superconduting qubit devices
## I Introduction
Superconducting (SC) Josephson junctions are considered promising qubits for quantum information processing. This “artificial atom”, with well-defined discrete energy levels, provides a platform to test fundamental quantum effects, e.g., cavity quantum electrodynamics (QED). The study of the cavity QED of a SC qubit, e.g., in Ref. you , can also open new directions for studying the interaction between light and solid state quantum devices. These can result in novel controllable electro-optical quantum devices in the microwave regime, such as microwave single-photon generators and detectors. Cavity QED can allow the transfer of information among SC qubits via photons, used as information bus.
Recently, different information buses using bosonic systems, which play a role analogous to a single-mode light field, have been proposed to mediate the interaction between the SC qubits. These bosonic “information bus” systems can be modelled by: nanomechanical resonators (e.g., in Refs. nano ); large junctions (e.g., Ref. wang ); current-biased large junctions (e.g., Refs. large ), and LC oscillators (e.g., Refs. lc ). However, the enormous versatility provided by photons should stimulate physicists to pay more attention to SC qubits interacting via photons,while embedded inside a QED cavity.
Several theoretical proposals have analyzed the interaction between SC qubits and quantized saidi ; you ; you1 ; liu ; gao ; vourdas ; zagoskin or classical fields zhou ; paspalakis ; liu1 . The strong coupling of a single photon to a SC charge qubit has been experimentally demonstrated wallraff by using a one-dimensional transmission line resonator blais . But, the QED effect of the SC qubit inside higher-dimensional cavities has not been experimentally observed. The main roadblocks seem to be: i) whether the cavity quality factor $`Q`$ can still be maintained high enough when the SC qubit is placed inside the cavity. Different from atoms, the effect of the SC qubit on the $`Q`$ value of the cavity is not negligible due to its complex structure and larger size. ii) The higher-dimensional QED cavity has relatively large mode volume, making the interaction between the cavity field and the qubit not be strong enough for the required quantum operations within the decoherence time. iii) The transfer of information among different SC qubits requires the qubit-photon interaction to be switched on/off by the external classical flux on time scales of the inverse Josephson energy. A higher cavity $`Q`$ value, a stronger qubit-photon interaction, and a faster switching interaction for the SC qubit QED experiments, seem difficult to achieve anytime soon.
In view of the above problems, it would be desirable to explore the possibility to demonstrate a variety of relatively simple cavity QED phenomena with a SC qubit. The determination of the cavity $`Q`$ value is a very important first step for the experiments on cavity QED with SC qubits. However, theoretical calculations of the $`Q`$ value are not always easy to perform because of the complexity of the circuit. Recent experiments pkd on broadband SC detectors showed that the $`Q`$ value of the SC device can reach $`2\times 10^6`$, which indicates that relatively simple experiments using cavity QED with a SC qubit are possible.
In this paper, we propose an experimentally feasible method which can be used to demonstrate a simple cavity QED effect of the SC qubit. For instance, superpositions of two macroscopic quantum states of a single-mode microwave cavity field can be created by the interaction between a SC charge qubit and the cavity field. At this stage, the injected light field is initially a coherent state, which can be easily prepared. The decoherence of the created superposition states can be further determined by measuring either the Wigner function of the cavity field or the charge qubit states. Then the cavity $`Q`$ value can be inferred from this decoherence measurement. Our proposal only needs few operations with a relatively low Q value. Also, we do not need to assume a very fast sweep rate of the external magnetic field for switching on/off the qubit-field interaction. Furthermore, the qubit-field interaction is not necessarily resonant.
We begin in Sec. II with a brief overview of the qubit-field interaction. In Sec. III, we discuss how to prepare superpositions of two different cavity field states under the condition of large detuning. In Sec. IV, the cavity $`Q`$ value is determined by the tomographic reconstruction of the cavity field Wigner function. In Sec. V, we show an alternative method to determine the $`Q`$ value according to the states of the qubit. Finally, we list our conclusions.
## II Theoretical model
We briefly review a model of a SC charge qubit inside a cavity. The Hamiltonian can be written as you ; you1 ; liu ; gao
$`H=\mathrm{}\omega a^{}a+E_z\sigma _z`$ (1)
$`E_J\sigma _x\mathrm{cos}\left[{\displaystyle \frac{\pi }{\mathrm{\Phi }_0}}\left(\mathrm{\Phi }_cI+\eta a+\eta ^{}a^{}\right)\right],`$
where the first two terms respectively represent the free Hamiltonians of the cavity field with frequency $`\omega `$ for the photon creation (annihilation) operator $`a^{}(a)`$, and the qubit charging energy
$$E_z=2E_{\mathrm{ch}}(12n_g),$$
(2)
which depends on the gate charge $`n_g`$. The single-electron charging energy is $`E_{\mathrm{ch}}=e^2/2(C_g+2C_J)`$ with the capacitors $`C_g`$ and $`C_J`$ of the gate and the Josephson junction, respectively. The dimensionless gate charge, $`n_g=C_gV_g/2e`$, is controlled by the gate voltage $`V_g`$. Here, $`\sigma _z`$, $`\sigma _x`$ are the Pauli operators, and the charge excited state $`|e`$ and ground state $`|g`$ correspond to the eigenstates $`|=\left(\begin{array}{c}0\hfill \\ 1\hfill \end{array}\right)`$ and $`|=\left(\begin{array}{c}1\hfill \\ 0\hfill \end{array}\right)`$ of the spin operator $`\sigma _z`$, respectively. $`I`$ is an identity operator. The third term is the nonlinear qubit-photon interaction. $`E_J`$ is the Josephson energy for a single junction. The parameter $`\eta `$ is defined as $`\eta =_S𝐮(𝐫)𝑑𝐬`$ with the mode function of the cavity field $`𝐮(𝐫)`$, $`S`$ is the surface defined by the contour of the SQUID. We can decompose the cosine in Eq. (1) into classical and quantized parts. The quantized parts $`\mathrm{sin}[\pi (\eta a+H.c.)/\mathrm{\Phi }_0]`$ and $`\mathrm{cos}[\pi (\eta a+H.c.)/\mathrm{\Phi }_0]`$ can be further expanded as a power series in $`a(a^{})`$. To estimate the qubit-photon coupling constant, the qubit is assumed to be inside a full-wave cavity with the standing-wave form for a single-mode magnetic field scully
$$B_x=i\sqrt{\frac{\mathrm{}\omega }{\epsilon _0Vc^2}}(aa^{})\mathrm{cos}(kz).$$
(3)
The polarization of the magnetic field is along the normal direction of the surface area of the SQUID, located at an antinode of the standing-wave mode. The mode function $`\sqrt{\mathrm{}\omega /\epsilon _0Vc^2}\mathrm{cos}(kz)`$ can be assumed to be independent of the integral area because the maximum linear dimension of the SQUID, e.g., even for $`50\mu `$m, is much less than $`0.1`$ cm, the shortest microwave wavelength of the cavity field. Then, in the microwave regime, the estimated range of values for $`\pi \eta /\mathrm{\Phi }_0`$ is: $`8.55\times 10^6\pi \eta /\mathrm{\Phi }_01.9\times 10^3`$, for a fixed area of the SQUID, e.g., $`50\mu `$m $`\times 50\mu `$m. If the light field is not so strong (e.g., the average number of photons inside the cavity $`N=a^{}a100`$), then we can only keep the first order of $`\pi \eta /\mathrm{\Phi }_0`$ and safely neglect all higher orders. Thus, the Hamiltonian (1) becomes
$`H=\mathrm{}\omega a^{}a+E_z\sigma _zE_J\sigma _x\mathrm{cos}({\displaystyle \frac{\pi \mathrm{\Phi }_c}{\mathrm{\Phi }_0}})`$
$`+{\displaystyle \frac{\pi E_J}{\mathrm{\Phi }_0}}\mathrm{sin}({\displaystyle \frac{\pi \mathrm{\Phi }_c}{\mathrm{\Phi }_0}})\left(\eta a\sigma _++\eta ^{}a^{}\sigma _{}\right).`$ (4)
It is clear that the qubit-photon interaction can be controlled by the classical flux $`\mathrm{\Phi }_c`$, after neglecting higher-orders in $`\pi \eta /\mathrm{\Phi }_0`$.
## III Preparation of macroscopic superposition states
The qubit-photon system can be initialized by adjusting the gate voltage $`V_g`$ and the external flux $`\mathrm{\Phi }_c`$ such that $`n_g=1/2`$ and $`\mathrm{\Phi }_c=0`$, then the dynamics of the qubit-field is governed by the Hamiltonian
$$H_1=\mathrm{}\omega a^{}aE_J\sigma _x.$$
(5)
Now there is no interaction between the cavity field and the qubit; thus, the cavity field and the qubit evolve according to Eq. (5). We assume that the qubit-photon system works at low temperatures $`T`$ ( e.g., $`T=30`$ mK in Ref. nakamura ), then the mean number of thermal photons $`n_{th}`$ in the cavity can be negligible in the microwave regime liu , and the cavity is approximately considered in the zero temperature environment. The initial state of the cavity field is prepared by injecting a single-mode coherent light
$$|\alpha =\mathrm{exp}\left\{\frac{|\alpha |^2}{2}\right\}\underset{n=0}{}\frac{\alpha ^n}{\sqrt{n!}}|n,$$
(6)
into the cavity. Here, without loss of generality, $`\alpha `$ is assumed to be a real number, and $`a|\alpha =\alpha |\alpha `$. The qubit is assumed to be initially in the ground state $`|g`$. After a time interval $`\tau _1=\mathrm{}\pi /4E_J`$, the qubit ground state $`|g`$ is transformed as $`|g\left(|g+i|e\right)/\sqrt{2}`$; then, the qubit-photon state evolves into
$$|\psi (\tau _1)=\frac{1}{\sqrt{2}}(|g+i|e)|\alpha .$$
(7)
Here we have neglected the free evolution phase factor $`e^{i\omega \tau _1}`$ in $`\alpha `$.
Now, we assume that the gate voltage and the magnetic flux are switched to $`n_g1/2`$ (this value of $`n_g`$ will be specified later) and $`\mathrm{\Phi }_c=\mathrm{\Phi }_0/2`$, respectively. Then the qubit-photon interaction appears and the effective Hamiltonian governing the dynamic evolution of the qubit-photon can be written as (see Appendix A)
$$H_2=\mathrm{}\omega _{}a^{}a+\frac{1}{2}\mathrm{}\mathrm{\Omega }\sigma _z+\mathrm{}\frac{|g|^2}{\mathrm{\Delta }}\left(1+2a^{}a\right)|ee|,$$
(8)
with $`\omega _{}=\omega |g|^2/\mathrm{\Delta }`$ and $`g=(\pi \eta E_J)/(\mathrm{}\mathrm{\Phi }_0)`$. The detuning $`\mathrm{\Delta }=\mathrm{\Omega }\omega >0`$ between the qubit transition frequency $`\mathrm{\Omega }=4E_{ch}(12n_g)/\mathrm{}`$ and the cavity field frequency $`\omega `$ is assumed to satisfy the large detuning condition
$$\frac{\pi E_J|\eta |}{\mathrm{}\mathrm{\Phi }_0\mathrm{\Delta }}1.$$
(9)
The unitary evolution operator corresponding to Eq. (8) can be written as
$`U(t)`$ $`=`$ $`\mathrm{exp}\left[i\left(\omega _{}a^{}a+{\displaystyle \frac{\mathrm{\Omega }}{2}}\sigma _z\right)t\right]`$ (10)
$`\times `$ $`\mathrm{exp}\left[itF(a^{}a)|ee|\right],`$
here, the operator $`F(a^{}a)`$ is expressed as
$$F(a^{}a)=\frac{|g|^2}{\mathrm{\Delta }}(1+2a^{}a).$$
(11)
With an evolution time $`\tau _2`$, the state (7) evolves into
$$|\psi (\tau _2)=\frac{1}{\sqrt{2}}\left[|g|\beta +i\mathrm{exp}(i\theta )|e|\beta ^{}\right]$$
(12)
where a global phase $`\mathrm{exp}(i\mathrm{\Omega }\tau _2/2)`$ has been neglected, $`\theta =(\mathrm{\Omega }|g|^2/\mathrm{\Delta })\tau _2`$, $`\beta =\alpha \mathrm{exp}[i\omega _{}\tau _2]`$, and $`\beta ^{}=\beta \mathrm{exp}(i\varphi )`$, with $`\varphi =2|g|^2\tau _2/\mathrm{\Delta }`$. Equation (12) shows that a phase shift $`\varphi `$ is generated for the coherent state $`|\beta `$ of the cavity field when the qubit is in the excited state $`|e`$, but the qubit ground state $`|g`$ does not induce an extra phase for the coherent state $`|\beta `$.
The gate voltage and the magnetic field are now adjusted such that the conditions $`n_g=1/2`$ and $`\mathrm{\Phi }_c=0`$ are satisfied; then the qubit-photon interaction is switched off. Now let the system evolve a time $`\tau ^{}=\tau _1=\mathrm{}\pi /4E_J`$, then Eq. (12) becomes
$`|\psi (\tau _2)`$ $`=`$ $`{\displaystyle \frac{1}{2}}|g[|\beta \mathrm{exp}(i\theta )|\beta ^{}]`$ (13)
$`+`$ $`i{\displaystyle \frac{1}{2}}|e[|\beta +\mathrm{exp}(i\theta )|\beta ^{}],`$
where a free phase factor $`e^{i\omega \tau _1}`$ in the cavity field states $`|\beta `$ and $`|\beta e^{i\varphi }`$ has been neglected.
The superpositions of two distinguished coherent states can be conditionally generated by measuring the charge states of the qubit as,
$$|\beta _\pm =N_\pm ^1[|\beta \pm \mathrm{exp}(i\theta )|\beta ^{}],$$
(14)
where the $`+`$ ($``$) correspond to the measurement results $`|e`$ ($`|g`$), and the normalized constants $`N_\pm `$ are determined by
$$N_\pm ^2=2\pm 2\mathrm{cos}\theta ^{}\mathrm{exp}\left[2|\alpha |^2\mathrm{sin}^2\left(\frac{\varphi }{2}\right)\right],$$
(15)
where $`\theta ^{}=|\alpha |^2\mathrm{sin}\varphi \theta `$, and the relation $`|\beta |^2=|\alpha |^2`$ is used.
Due to $`\mathrm{\Phi }_c=0`$, after the superpositions in Eq.(14) are created, the dynamic evolution of the cavity field is only affected by its dissipation, characterized by the decay rate $`\gamma `$, which can be expressed by virtue of the cavity quality factor $`Q`$ as $`Q=\omega /\gamma `$. Now let the cavity field described by the states $`|\beta _+`$ or $`|\beta _{}`$ evolve a time $`\tau _3`$; then the reduced density matrices of the superpositions can be described by
$`\rho _\pm (\tau _3)={\displaystyle \frac{1}{N_\pm ^2}}\{|\beta u\beta u|+|\beta ^{}u\beta ^{}u|`$
$`\pm C|\beta ^{}u\beta u|\pm C^{}|\beta u\beta ^{}u|\},`$ (16)
where
$$C=\mathrm{exp}(i\theta )\mathrm{exp}\left\{|\alpha |^2(1e^{i\varphi })(u^21)\right\}$$
(17)
and
$$uu(\tau _3)=\mathrm{exp}\left(\frac{\gamma }{2}\tau _3\right)=\mathrm{exp}\left(\frac{\omega \tau _3}{2Q}\right).$$
(18)
It is clearly shown that the mixed state in Eq. (III) is strongly affected by the $`Q`$ value. Equation (III) is derived for zero temperature since thermal photons are negligible at low-temperature. Equations (III-18) show that the information of the cavity quality factor $`Q`$ can be encoded in a reduced density matrix of the cavity field. The $`Q`$ value can be determined using two different methods, after encoding its information in Eqs.(III-18). Below, we will discuss these two approaches.
## IV Measuring $`Q`$ by photon state tomography
The state $`\rho `$ of the optical field is generally visualized when it is represented by a Wigner function hans in the position $`x`$ and momentum $`p`$ space, which is written as
$$W(x,p)=\frac{1}{\pi }_{\mathrm{}}^{\mathrm{}}xx^{}|\rho |x+x^{}e^{i2px^{}}dx^{}.$$
(19)
The Wigner function $`W(x,p)`$ can be experimentally measured by state tomographic techniques hans . For any two coherent states, $`|\alpha `$ and $`|\beta `$, the Wigner function $`W(x,p)`$ can be represented as
$`W_{\alpha ,\beta }(x,p)`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}xx^{}|\alpha \beta |x+x^{}e^{i2px^{}}dx^{}`$ (20)
$`=`$ $`{\displaystyle \frac{1}{\pi }}\mathrm{exp}\left\{{\displaystyle \frac{1}{2}}(|\alpha |^2+|\beta |^22\alpha \beta ^{})\right\}`$
$`\times `$ $`\mathrm{exp}\left\{(xq_1)^2(p+iq_2)^2\right\}`$
with $`q_1=(\alpha +\beta ^{})/\sqrt{2}`$ and $`q_2=(\alpha \beta ^{})/\sqrt{2}`$. The Wigner functions $`W_\pm (x,p)`$ and $`W_{\pm D}(x,p)`$ for the states (14) and (III) were calculated (see Appendix B) by using Eq. (20). Comparing the tomographically measured results for the states (14) and (III), the $`Q`$ factor of the cavity can be finally determined, as explained below by using an example.
We further numerically calculate the Wigner functions $`W_\pm (x,p)`$ and $`W_{\pm D}(x,p)`$ of the states (14) and (III) from the SC qubit parameters and given operation durations. Using current values for experimental data, the basic physical parameters can be specified as follows. We assume that the SC Cooper-pair box is made from aluminum, with a BCS energy gap of $`2.4`$K (about 50 GHz) lehnert , the charge energy $`E_{\mathrm{ch}}`$ and the Josephson energy $`E_\mathrm{J}`$ are $`4E_{\mathrm{ch}}/h=149`$ GHz and $`2E_\mathrm{J}/h=13.0`$ GHz, respectively lehnert . The frequency of the cavity field is taken as $`40`$ GHz, corresponding to a wavelength $`0.75`$ cm. The above numbers show that the SC energy gap is the largest energy, so the quasi-particle excitation on the island can be well suppressed at low temperatures, e.g., $`20`$ mK . The SQUID area is assumed to be about $`50\mu `$m $`\times 50\mu `$m, then the absolute value $`|g|`$ of the qubit-photon coupling constant is about $`|g|=4\times 10^6`$ rad s<sup>-1</sup>.
Let us now prepare entangled qubit-photon states as in Eq. (13). Any gate charge value $`n_g`$ in Eq. (2), in which the large detuning condition in Eq. (9) is satisfied, can be chosen to realize our proposal. For concreteness, we give an example. The gate voltage is adjusted such that the gate charge is $`n_g0.634233`$, which can be experimentally achieved lehnert , then the detuning $`\mathrm{\Delta }=\mathrm{\Omega }\omega 9.0\times 10^6`$ rad s<sup>-1</sup>. Thus, $`\mathrm{\Omega }`$ is about $`40`$ GHz plus $`1.4`$ MHz, and $`|g|^2/\mathrm{\Delta }\mathrm{\hspace{0.17em}0.27}`$ MHz. We can also find that $`\mathrm{\Delta }/|g|2.3`$, so a large-detuning condition can be used xm ; pt . For a given Josephson energy, $`2E_J/h=13.0`$ GHz, the operation times $`\tau _1=4.8\times 10^{12}`$ s, required to prepare a superposition of $`|e`$ and $`|g`$ with equal probabilities, is much less than the qubit relaxation time $`T_1=1.3\mu `$s and dephasing time $`T_2=5`$ ns. We can choose the duration $`\tau _2`$ for a given input coherent state $`|\alpha `$ with the condition, that the distance $`|\beta \beta ^{}|`$ between two coherent states $`|\beta `$ and $`|\beta ^{}`$ satisfies
$$|\beta \beta ^{}|=2|\alpha |\mathrm{sin}\left(\frac{\varphi }{2}\right)>1.$$
(21)
So the lower bound of the duration $`\tau _2`$ can be given as
$$\tau _2=\frac{\mathrm{\Delta }}{|g|^2}\mathrm{arcsin}\left(\frac{1}{2|\alpha |}\right),$$
(22)
when $`0\varphi \pi `$. Equation (22) shows that a shorter $`\tau _2`$ can be obtained for a higher intensity $`|\alpha |`$ with fixed detuning $`\mathrm{\Delta }`$ and coupling constant $`g`$.
As an example, we plot the Wigner function of the superposition $`|\beta _+`$ in Fig. 1(a) for an input coherent light $`|\alpha `$ with the mean photon number $`\overline{n}=|\alpha |^2=16`$. We choose a simple case $`\varphi =\pi `$, corresponding to the operation time $`\tau _20.93\mu `$s, which is less than the qubit lifetime $`T_1`$ and the cavity field lifetime $`T_{\mathrm{ph}}2\mu `$s for a bad cavity with $`Q=5\times 10^5`$. In such a case, $`\beta ^{}=\beta `$ and the phase $`\theta `$ is about $`0.996[\mathrm{mod}\mathrm{\hspace{0.17em}2}\pi ]`$ rad. Other parameters used in Fig. 1 are given above. If we set the evolution time $`\tau _3=0.1\mu `$s, then the Wigner function of Eq. (III) for the above cavity quality factor is shown in Fig. 1(b). The central structure in Fig. 1(a) represents the coherence of the quantum state. In Fig. 1(b), we find that the height of the Wigner function $`W_{+D}(x,p)`$, especially for the central structure, is reduced by the environment. Comparing Fig. 1(a) and Fig. 1(b), it is found that the coherence of the superposed states is suppressed by the environment, and the decoherence of superpositions is tied to the energy dissipation of the cavity field. Then, the $`Q`$ value can in principle be estimated by measuring the Wigner functions of Eqs. (III) and (14), and comparing these two kinds of results.
## V Determining $`Q`$ by readout of charge states
The determination of the $`Q`$ value by measuring the Wigner function needs optical instruments. In solid state experiments, the charge states are typically measured. Instead of using optical instruments, it would be desirable to obtain the $`Q`$ value by measuring charge states. This will be our goal here. The process to achieve this can be described as follows.
i) According to the measurements on the charge qubit states in Eq.(13), the qubit-photon states are projected to $`|g|\beta _{}`$ or $`|e|\beta _+`$, respectively. After the evolution time $`\tau _3^{}`$, a $`\pi /2`$ quantum operation is performed on the qubit with the duration $`T=\mathrm{}\pi /4E_J`$. Then, the qubit ground state $`|g`$, or excited state $`|e`$, is transformed into the superposition $`(|g+i|e)/\sqrt{2}`$, or $`(i|g+|e)/\sqrt{2}`$, and the photon states $`|\beta _\pm `$ evolve into mixed states after the evolution time $`\tau =\tau _3^{}+T`$, and the photon-qubit states can be expressed as
$$\rho _{Q+F}=\frac{1}{2}(|g\pm i|e)(g|ie|)\rho _\pm (\tau ),$$
(23)
with subscripts $`Q`$ and $`F`$ denoting the qubit and cavity field, respectively. The reduced density matrices $`\rho _\pm (\tau )`$ take the same form as in Eq. (III) with $`\tau `$ replacing $`\tau _3`$.
ii) After the above procedure, the qubit-photon interaction is switched on by applying the external magnetic flux $`\mathrm{\Phi }_e=\mathrm{\Phi }_0/2`$. By using Eq. (10), Eq. (23) evolves into
$`2\rho _{A+F}^{(1)}`$ $`=`$ $`|gg|U_1(\tau _4)\rho _\pm (\tau )U_1^{}(\tau _4)`$
$`+`$ $`|ee|U_2(\tau _4)\rho _\pm (\tau )U_2^{}(\tau _4)`$
$``$ $`i\mathrm{exp}(i\mathrm{\Omega }_{}\tau _4)|ge|U_1(\tau _4)\rho _\pm (\tau )U_2^{}(\tau _4)`$
$`\pm `$ $`i\mathrm{exp}(+i\mathrm{\Omega }_{}\tau _4)|eg|U_2(\tau _4)\rho _\pm (\tau )U_1^{}(\tau _4)`$
with $`\mathrm{\Omega }_{}=\mathrm{\Omega }|g|^2/\mathrm{\Delta }`$, and a shorter evolution time $`\tau _4`$. For example, $`\tau _4`$ is less than the lifetime $`T_1`$ of the qubit at least. The time evolution operators $`U_1(\tau _4)`$ and $`U_2(\tau _4)`$ in Eq. (V) are
$`U_1(\tau _4)`$ $`=`$ $`\mathrm{exp}\left[i\omega _{}a^{}a\tau _4\right],`$ (25a)
$`U_2(\tau _4)`$ $`=`$ $`\mathrm{exp}\left[i\omega _+a^{}a\tau _4\right].`$ (25b)
with $`\omega _\pm =\omega \pm |g|^2/\mathrm{\Delta }`$. After this qubit-photon interaction, the information of the $`Q`$ value is encoded.
iii) The qubit-photon coupling is switched off and a $`\pi /2`$ rotation is made on the qubit. If the state of the cavity field is prepared to $`|\beta _{}`$ of Eq. (14) in the first step, then the qubit is in the ground state $`|g`$. After measuring the qubit states, the photon states are projected to
$$\rho _{e/g}=\frac{1}{4}(A\pm B)$$
(26a)
where the sign $`\mathrm{`}\mathrm{`}+\mathrm{"}`$ corresponds to the excited state $`|e`$ measurement, but $`\mathrm{`}\mathrm{`}\mathrm{"}`$ corresponds to the ground state $`|g`$ measurement. The operators $`A`$ and $`B`$ are
$`A`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{2}{}}}U_i(\tau _4)\rho _{}(\tau )U_i^{}(\tau _4),`$ (26b)
$`B`$ $`=`$ $`2\mathrm{Re}[\mathrm{exp}(i\mathrm{\Omega }_{}\tau _4)U_1(\tau _4)\rho _{}(\tau )U_2^{}(\tau _4)].`$ (26c)
After tracing out the cavity field state, the probabilities corresponding to measuring charge states $`|e`$ and $`|g`$ are
$`P_{e/g}(\tau )`$ $`=`$ $`\mathrm{Tr}_F\left(\rho _{e/g}\right)`$ (27)
$`=`$ $`{\displaystyle \frac{1}{2}}\left\{1\pm \mathrm{Re}(\mathrm{Tr}_F[\mathrm{exp}(i\phi )\rho _{}(\tau )])\right\}`$
with $`\phi =(\mathrm{\Omega }_{}2|g|^2a^{}a/\mathrm{\Delta })\tau _4`$. Then the measurement probabilities are related to the $`Q`$ values. Substituting $`\rho _{}(\tau )`$ into Eq.(27), we can obtain
$`\mathrm{Re}\left\{\mathrm{Tr}[\mathrm{exp}(i\phi )\rho _{}(\tau )]\right\}`$ (28a)
$`={\displaystyle \frac{2}{N_{}^2}}\mathrm{exp}\left[2\alpha (\tau )\mathrm{sin}^2\varphi ^{}\right]\mathrm{cos}\left[\mathrm{\Omega }_{}\tau _4\alpha (\tau )\mathrm{sin}(2\varphi ^{})\right]`$
$`{\displaystyle \frac{1}{N_{}^2}}\mathrm{cos}\left[\theta _{}|\alpha |^2\mathrm{sin}\varphi +\theta \mathrm{\Omega }_{}\tau _4\right]\mathrm{exp}\left(+G_{}\mathrm{\Gamma }\right)`$
$`{\displaystyle \frac{1}{N_{}^2}}\mathrm{cos}\left[\theta _++|\alpha |^2\mathrm{sin}\varphi \theta \mathrm{\Omega }_{}\tau _4\right]\mathrm{exp}\left(G_+\mathrm{\Gamma }\right)`$
with the parameters
$`\varphi ^{}`$ $`=`$ $`{\displaystyle \frac{|g|^2}{\mathrm{\Delta }}}\tau _4,`$ (28b)
$`\mathrm{\Gamma }`$ $`=`$ $`2|\alpha |^2\mathrm{sin}^2({\displaystyle \frac{\varphi }{2}}),`$ (28c)
$`\alpha (\tau )`$ $`=`$ $`|\alpha u(\tau )|^2,`$ (28d)
$`G_\pm `$ $`=`$ $`2\alpha (\tau )\mathrm{sin}\varphi ^{}\mathrm{sin}(\varphi \pm \varphi ^{}),`$ (28e)
$`\theta _\pm `$ $`=`$ $`2\alpha (\tau )\mathrm{cos}(\varphi \pm \varphi ^{})\mathrm{sin}\varphi ^{}.`$ (28f)
From Eq. (28), we find that $`\varphi ^{}`$ should satisfy the condition $`\varphi ^{}n\pi `$ for $`\varphi =\pi `$, in order to describe the dissipation effect; here $`n`$ is an integer. Generally speaking, if one of the functions $`G_+`$, $`G_{}`$, $`\theta _+`$, $`\theta _{}`$, $`\mathrm{sin}\varphi ^{}`$, or $`\mathrm{sin}(2\varphi ^{})`$ is nonzero, then this is enough to encode the $`Q`$ value, which can be obtained from Eq. (28), together with Eq. (18), using $`\tau `$ instead of $`\tau _3`$.
However, if the superposition of the cavity fields is prepared to the state $`|\beta _+`$ in the first step, then the ground and excited state measurements make the cavity field collapse to state
$$\rho _{g/e}^{}=\frac{1}{4}(A^{}\pm B^{}),$$
(29)
where $`A^{}`$ and $`B^{}`$ have the same forms as Eqs. (26b) and (26c), just with the replacement of $`\rho _{}(\tau )`$ by $`\rho _+(\tau )`$.
The probabilities $`P_g^{}(\tau )`$ and $`P_e^{}(\tau )`$ to measure the qubit states $`|g`$ and $`|e`$ corresponding to the prepared state $`|\beta _+`$ of Eq. (14) after a dissipation interval $`\tau `$, can also be obtained as
$$P_{g/e}^{}(\tau )=\frac{1}{2}\left\{1\pm \mathrm{Re}(\mathrm{Tr}_F[\mathrm{exp}(i\phi )\rho _+(\tau )])\right\}$$
(30)
where $`\mathrm{Re}\{\mathrm{Tr}_F[\mathrm{exp}(i\phi )\rho _+(\tau )]\}`$ can be obtained by replacing $`N_{}`$ with $`N_+`$, and replacing the sign $`\mathrm{`}\mathrm{`}\mathrm{"}`$ before the second and third terms with the sign $`\mathrm{`}\mathrm{`}+\mathrm{"}`$ in Eq. (28a)
To determine the Q values by probing the charge states, the measurement should be made for two times, the first measurement is for the preparation of the superpositions of the cavity field. After the first measurement, we make a suitable qubit rotation, and then make the qubit interact with the cavity field for a duration $`\tau _4`$. Finally, the second measurement is made and the $`Q`$ information is encoded in the measured probabilities. The different evolution times $`\tau `$ correspond to the different measuring probabilities for given $`\tau _4`$ and other parameters $`|\alpha |`$, $`\mathrm{\Delta }`$, and so on. For example, the probabilities $`P_{e/g}(\tau )`$ for several special cases are discussed as follows when the prepared state is $`|\beta _{}`$. If we assume that the qubit rotations and qubit-photon dispersive interaction are made without energy dissipation of the cavity field, e.g., $`\tau =0`$, then the measuring probabilities $`P_{e/g}(\tau =0)`$ only encode the information of the cavity field but do not include the quality factor $`Q`$. If the coherence of the states $`|\beta _\pm `$ nearly disappears after time $`\tau `$, then the state $`|\beta _{}`$ becomes a classical statistical mixture
$$\rho _{}(\tau )=\frac{1}{N_{}^2}\left[|\beta u(\tau )\beta u(\tau )|+|\beta ^{}u(\tau )\beta ^{}u(\tau )|\right].$$
(31)
The probabilities $`P_{e/g}(\tau )`$ are then reduced to
$`P_{e/g}(\tau )`$ $`=`$ $`{\displaystyle \frac{1}{N_{}^2}}\pm {\displaystyle \frac{\mathrm{exp}\left[2\alpha (\tau )\mathrm{sin}^2\varphi ^{}\right]}{N_{}^2}}`$
$`\times `$ $`\mathrm{cos}\left[\mathrm{\Omega }_{}\tau _4\alpha (\tau )\mathrm{sin}(2\varphi ^{})\right],`$
which tends to $`1/2`$ for $`|\alpha |^21`$. If $`\tau >1/\gamma =t_{\mathrm{ph}}`$ of the single-photon state lifetime, then the photons of the states $`|\beta _\pm `$ are completely dissipated into the environment. In this case, the cavity quality factor $`Q`$ cannot be encoded in the probabilities $`P_{e/g}(\tau )`$ even with some qubit and qubit-photon states operations.
As an example, let us consider how $`P_g(\tau )`$ varies with the evolution time $`\tau `$ with the cavity field dissipation. We assume that the evolution time $`\tau _4=(\pi /2)(\mathrm{\Delta }/|g|^2)`$, that is, $`\varphi ^{}=\pi /2`$. Then, the $`\tau `$-dependent probabilities $`P_g(\tau )`$ for the initially prepared state $`|\beta _{}`$ are given in Fig. 2(a) with the same parameters as in Fig. 1, except with different cavity quality factors $`Q`$. In order to see how the probability $`P_g(\tau )`$ changes with the intensity $`|\alpha |^2`$ of the input coherent state $`|\alpha `$, we plot $`P_g(\tau )`$ in Fig. 2(b) with the same parameters as Fig. 2(a) except changing the intensity to $`|\alpha |^2=4`$ from $`|\alpha |^2=16`$. Figure 2 shows that both the higher quality factor $`Q`$ and weaker intensity $`|\alpha |^2`$ of the input cavity field correspond to a larger probability $`P_g`$ of the ground state for the fixed evolution time $`\tau `$. For fixed $`Q`$ and $`\tau `$, the weaker intensity $`|\alpha |^2`$ corresponds to a higher measuring probability. We plot $`P_g(\tau )`$ in Fig. 2 considering the simple case $`\varphi =\pi `$. However, if we consider another $`\varphi `$, then $`|\alpha |^2`$ should be chosen such that it satisfies the condition $`2|\alpha |\mathrm{sin}(\varphi /2)>1`$. In conclusion, the quality factor $`Q`$ can be determined from the probabilities $`P_{e/g}(\tau )`$ of measuring the qubit states with a finite cavity field evolution time $`\tau `$.
## VI Discussions and conclusions
We discussed how to measure the cavity quality factor $`Q`$ by using the interaction between a single-mode microwave cavity field and a controllable superconducting charge qubit. Two methods are proposed. One measures the Wigner function of the state (III) by using a standard optical method hans . Another approach measures the qubit states. Using this last method, the information of the $`Q`$ value can be encoded into the reduced density matrix of the cavity field, and at the same time the qubit makes a $`\pi /2`$ rotation. Thus, with a suitable qubit-photon interaction time, information on the $`Q`$ value is then transferred to the qubit-photon states. Finally, after another $`\pi /2`$ rotation, the charge qubit states are measured, and the $`Q`$ value can be obtained, as shown in Fig. 2, Eqs. (28) and Eq. (18). However, it should be noticed that it is easy to measure charge states than to measure photon states in superconducting circuits.
Our proposal shows that a cavity QED experiment with a SC qubit can be performed even for a relatively low $`Q`$ values, e.g., $`Q10^6`$. Initially, a coherent state is injected into the cavity, which is relatively easy to do experimentally. Although all rotations of the qubit are chosen as $`\pi /2`$ to demonstrate our proposal, other rotations can also be used to achieve our goal.
To simplify these studies and without loss of generality, we have assumed two components $`|\beta `$ and $`|\beta \mathrm{exp}(i\varphi )`$ for superpositions with $`\varphi =\pi `$ phase difference in our numerical demonstrations. Of course, other superpositions can also be used to realize our purpose. The only condition to satisfy is that the distance between the two states $`|\beta `$ and $`|\beta \mathrm{exp}(i\varphi )`$ should be larger than one. In order to obtain a numerical estimate for the detuning, we specify a value of the gate charge number $`n_g`$. However, any gate charge that satisfies the large-detuning condition can be chosen to realize our proposal.
Although we did not give a detailed description of another resonance-based approach, it should be pointed out that the $`Q`$ values can also be determined by virtue of the resonant qubit-photon interaction. For example, if the superpositions liu of the vacuum and the single photon state are experimentally prepared, then we can follow the same steps as in Sec. III and IV to obtain the $`Q`$ value. This method rinner has been applied to micromasers, where the qubits are two-level atoms. However, the coherent states and non-resonant qubit-photon interaction should be easier to do experimentally than the approach using single-photon states and resonant qubit-photon interaction.
Our proposal can also be generalized to the models used in Refs. wallraff ; blais , which are experimentally accessible. We hope that our proposal can open new doors to experimentally test the $`Q`$ value and motivate further experiments on cavity quantum electrodynamics with SC qubits.
## VII acknowledgments
This work was supported in part by the National Security Agency (NSA) and Advanced Research and Development Activity (ARDA) under Air Force Office of Research (AFOSR) contract number F49620-02-1-0334, and by the National Science Foundation grant No. EIA-0130383.
## Appendix A Effective hamiltonian with larger detuning
The Hamiltonian $`H=H_0+H_1`$ of the two-level atom interacting with a single-mode cavity field can be written as
$`H_0`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{}\mathrm{\Omega }\sigma _z+\mathrm{}\omega a^{}a,`$ (33a)
$`H_1`$ $`=`$ $`\mathrm{}(ga^{}\sigma _{}+g^{}a\sigma _+)`$ (33b)
with a complex number $`g`$. Let us assume $`\mathrm{\Delta }=\mathrm{\Omega }\omega >0`$ and $`g/\mathrm{\Delta }1`$. The eigenstates and corresponding eigenvalues of the free Hamiltonian $`H_0`$ are
$`|e|n`$ $``$ $`n\mathrm{}\omega +{\displaystyle \frac{1}{2}}\mathrm{}\mathrm{\Omega },`$ (34a)
$`|g|m`$ $``$ $`m\mathrm{}\omega {\displaystyle \frac{1}{2}}\mathrm{}\mathrm{\Omega }`$ (34b)
In the interaction picture, any state can be written as
$$|\psi (t)=U(t,t_0)|\psi (t_0)$$
(35)
with
$`U(t,t_0)`$ $`=`$ $`1+{\displaystyle \frac{1}{i\mathrm{}}}{\displaystyle _{t_0}^t}H_{\mathrm{int}}(t_1)dt_1`$
$`+`$ $`\left({\displaystyle \frac{1}{i\mathrm{}}}\right)^2{\displaystyle _{t_0}^t}{\displaystyle _{t_0}^{t_1}}H_{\mathrm{int}}(t_1)H_{\mathrm{int}}(t_2)dt_1dt_2+\mathrm{}`$
here $`H_{\mathrm{int}}=U_0^{}(t)H_1U_0(t)`$ with $`U_0(t)=\mathrm{exp}\{iH_0t/\mathrm{}\}`$. In the basis $`\{|E_l\}=\{|e|n,|g|m\}`$, Eq. (A) can be expressed as
$`U(t,t_0)`$ $`=`$ $`1+`$
$`+`$ $`{\displaystyle \frac{1}{i\mathrm{}}}{\displaystyle _{t_0}^t}{\displaystyle \underset{l,m}{}}|E_lE_l|H_{\mathrm{int}}(t_1)|E_mE_m|\mathrm{d}t_1+\mathrm{}.`$
After neglecting the fast-oscillating factor and keeping the first order term in $`g/\mathrm{\Delta }`$, $`U(t,t_0)`$
$`U(t,t_0)=U(t,0)=U(t)`$ (37)
$`1i{\displaystyle \frac{|g|^2}{\mathrm{\Delta }}}{\displaystyle _0^t}dt_1[(n+1)|e,ne,n|n|g,ng,n|],`$
where we assume $`t_0=0`$. Finally, we obtain the effective Hamiltonian in the interaction picture as
$$H_{\mathrm{eff}}=\mathrm{}\frac{|g|^2}{\mathrm{\Delta }}(|ee|aa^{}|gg|a^{}a).$$
(38)
Returning Eq. (38) to the Schrödinger picture, Eq. (10) is obtained. This method can be generalized to obtain the effective Hamiltonian of the model with many two-level system interacting with a common single-mode field. Equation (10) can also be obtained by using the Fröhlich-Nakajima transformation Fro ; Nakajima ; wu ; sun .
## Appendix B Wigner functions of superposition and mixed states
For completeness, we explicitly write the Wigner functions $`W_\pm (x,p)`$ of the superposition states in Eq. (14) as follows:
$`W_\pm (x,p)`$
$`={\displaystyle \frac{1}{\pi N_\pm ^2}}\{\mathrm{exp}[(x\sqrt{2}\mathrm{Re}\beta )^2(p\sqrt{2}\mathrm{Im}\beta )^2]`$
$`+\mathrm{exp}\left[\left(x\sqrt{2}\mathrm{Re}\beta ^{}\right)^2\left(p\sqrt{2}\mathrm{Im}\beta ^{}\right)^2\right]`$
$`\pm 2\mathrm{R}\mathrm{e}\left[P\mathrm{exp}((x\mathrm{}_1)^2(p+i\mathrm{}_2)^2)\right]\},`$ (39)
with
$`P`$ $`=`$ $`\mathrm{exp}(i\theta )\mathrm{exp}[|\alpha |^2(1e^{i\varphi })],`$ (40)
$`\mathrm{}_1`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(\beta +\beta ^{}),`$ (41)
$`\mathrm{}_2`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(\beta \beta ^{}).`$ (42)
The Wigner functions $`W_{\pm D}(x,p)`$ of the mixed states in Eq. (III) with dissipation can be written as
$`W_{\pm D}(x,p)`$
$`={\displaystyle \frac{1}{\pi N_\pm ^2}}\{\mathrm{exp}[(xu\sqrt{2}\mathrm{Re}\beta )^2(pu\sqrt{2}\mathrm{Im}\beta )^2]`$
$`+\mathrm{exp}\left[\left(xu\sqrt{2}\mathrm{Re}\beta ^{}\right)^2\left(pu\sqrt{2}\mathrm{Im}\beta ^{}\right)^2\right]`$
$`\pm 2\mathrm{R}\mathrm{e}\left[P\mathrm{exp}((xu\mathrm{}_1)^2(p+iu\mathrm{}_2)^2)\right]\}.`$ (43) |
warning/0506/cs0506051.html | ar5iv | text | # Comparison of two different implementations of a finite-difference-method for first-order pde in Mathematica^"®" and Matlab^"®"
## 1 Physics and analytical solution
The growth rate of microcracks in a brittle material can be discribed by a mesoscopic equation. Here the specialized version for uniaxial loading is presented.
$`{\displaystyle \frac{f(l,t)}{t}}`$ $`=`$ $`{\displaystyle \frac{1}{l^2}}{\displaystyle \frac{l^2v_l(l,t)f(l,t)}{l}},`$ (1)
$`f(l,t)`$ is the distribution function for the crack length $`l`$ at time $`t`$, $`v_l=\dot{l}`$ is the growth velocity of the cracks. A Rice-Griffith-like dynamic is assumed for crack growth, which gives
$`\dot{l}`$ $`=`$ $`\{\begin{array}{cc}\alpha ^{}+\beta ^{}l\sigma (t)^2\hfill & \text{, if }\alpha ^{}\beta ^{}l\sigma ^2\hfill \\ 0\hfill & \text{, otherwise}\hfill \end{array}`$ (2)
The theory is given in detail in . For an exponential (or a step-wise) initial condition and constant loading speed it is possible to give an exact analytical solution, which is also presented in and looks like:
$`f(l,t)`$ $`=`$ $`\{\begin{array}{cc}l^2e^{\frac{\beta ^{}v_\sigma ^2}{3}t^3}F\left(le^{\frac{\beta ^{}v_\sigma ^2}{3}t^3}+\frac{\alpha ^{}}{(9\beta ^{}v_\sigma ^2)^{1/3}}\mathrm{\Gamma }(1/3,0,\frac{\beta ^{}v_\sigma ^2}{3}t^3)\right)\hfill & \text{, if}\alpha ^{}\beta ^{}v_\sigma ^2lt^2\hfill \\ & \\ f(l,0)\hfill & \text{, otherwise}.\hfill \end{array}`$ (3)
## 2 The numerical algorithm
The partial differential equations are first order in time and crack length. Two different algorithms, upwind and fcts (forward time centered space), have been tested. The symmetric algorithm (ftcs-method) is in this case a bit more stable than the upwind, which is somewhat astonishing. Both algorithms can be found in .
In the symmetric algorithm $`f(l,t+1)`$ is calculated from $`f(l1,t)`$, $`f(l,t)`$ and $`f(l+1,t)`$ by
$`f(l,t+1)`$ $`=`$ $`f(l,t)(1(3\beta \sigma ^2(t){\displaystyle \frac{2\alpha }{dll}})dt)\left(dll\beta \sigma ^2(t)\right)\alpha )\times `$ (4)
$`\times {\displaystyle \frac{dt}{2dl}}\left(f(l+1,t)f(l1,t)\right)`$
This is what the main part of the implementation in Mathematica$`^\text{®}`$ looks like:
```
For[t = 1, t < TMAX, t++, For[l = 1, l < LMAX, l++,
If[l*dl*(t*dt)^2*be - al > 0,
If[l == 0, f[l, t + 1] = f[l, t],
f[l, t + 1] =
f[l, t]*(1 - (3*be*(t*dt)^2 - (2*al)/(dl*l))*dt) -
(dl*l*be*(t*dt)^2 - al)*dt/dl*(f[l + 1, t] - f[l - 1, t])],
f[l, t + 1] = f[l, t]
]
]
]
```
This is what the main part of the implementation in Matlab$`^\text{®}`$ looks like:
```
for t=1:1:TMAX-1
for l=1:1:LMAX-1
if (l.*dl.*sigma(t).*beta-alpha > 0)
f(l,t+1)=f(l,t).*(1-(3.*beta.*(sigma(t)).^2 - ...
(2.*alpha)./(dl.*l)).*dt)-(dl.*l.*beta.*(sigma(t)).^2 - ...
alpha).*dt./(2*dl).*(f(l+1,t)-f(l-1,t));
else
f(l,t+1)=f(l,t);
end
end
end
```
## 3 Results obtained by Mathematica$`^\text{®}`$
Mathematica$`^\text{®}`$ is a general purpose computer algebra system (cas) by Wolfram Research Inc..
As one can see the solution shows some wave-like efects, which can be interpreted as a sign for numerical instability. This instability is a result of the discontinous initial condition.
## 4 Results obtained by Matlab$`^\text{®}`$
Matlab$`^\text{®}`$ is a tool for numerical mathematics, especially designed for matrix manipulation, by The MathWorks Inc..
Here in both cases huge instabilities occur. As soon as some cracks start growing the numerical error goes to infinity. The following pictures show an extract of the above pictures.
## 5 Comparison of the results and conclusion
Obviously the results, obtained with both programs, show huge errors. But in contrast to the results calculated with Matlab$`^\text{®}`$ , one can use the results obtained with Mathematica$`^\text{®}`$ at least for some rough predictions.
The difference in the two software packages, used for these simulations, is that Matlab$`^\text{®}`$ uses floating-point variables of precision “double” (16 byte), whereas Mathematica$`^\text{®}`$ is capable of both numerical and symolic computation. Therefore it is possible that Mathematica$`^\text{®}`$ uses a much higher precision to perform some of the operations than Matlab$`^\text{®}`$ .
## Acknowledgement
We thank Dr. Christina Papenfuß and Dr. Peter Ván for discussions. Financial support by the DAAD, OTKA and by the VISHAY Company, 95100 Selb, Germany, is greatfully acknowledged. |
warning/0506/gr-qc0506090.html | ar5iv | text | # Form invariant transformations between n– and m– dimensional flat Friedmann–Robertson–Walker cosmologies
## I Introduction
Over the past several years there has been much interest in examining cosmology in higher dimensions to see if the standard four–dimensional Friedmann–Robertson–Walker (FRW) cosmology can be recovered. The idea that our physical four–dimensional Universe is embedded in a higher–dimensional spacetime has also also attracted the attention of particle physicists and astrophysicists. Theoretical motivation for such attempts can be found within the framework of many theories of unification, among them string, superstring and M theory, require extra spatial dimensions to be consistent. Until today a number of important solutions of Einstein equations in higher dimensions have been obtained and studied, and they have led to important generalizations and wider understanding of gravitational fields. In this respect, of interest are the works on n-dimensional black holesMyers , Kaluza–Klein inflationary cosmologies Abbott , circularly symmetric perfect fluids Garcia , black holes on branes Emparan , and recently, contributions on braneworld scenarios Lukas ; Randall ; Himemoto .
It is interesting to note that some authors have also considered phenomenological analysis in higher dimensional cosmology. For example, the phenomenological analysis of five-dimensional cosmology was stimulated by the work of Binetruy, Deffayet, and Langlois Binetruy , and subsequently by the Randall-Sundrum model Randall .
On the other hand, scalar fields play a crucial role in describing cosmological models. In the standard big-bang theory such fields are included for solving most of the problems found at very early times in the evolution of the universe, and are called the “inflaton” scalar field Guth ; Linde ; Albrecht . This scalar field is characterized by its scalar potential.
At the same time, measurements of the luminosity–redshift relations observed for the discovered type Ia supernovae with redshift $`z>0.35`$ Perlmutter ; Garnavich , indicate that at present the universe is expanding with an accelerated fashion suggesting a net negative pressure for the universe. One plausible explanation of this astronomical observation is based on the introduction of a scalar field, which is called the “quintessence” or “dark energy” scalar field.
Although these scalar fields are quite different in nature, there are authors who think that the “inflaton” and the “quintessence” fields might be of the same nature, in which a very specific scalar potential form is used Peebles .
In Ref. lmath it was shown that in several physical problems the Einstein field equations for flat FRW cosmological models and Bianchi I-type metric containing a scalar field can be linearized and solved by writing them in invariant form. In all these cases explicit use has been made of the non-local transformation group. The symmetry transformations that preserve the form of the Einstein equations introduce an alternative concept of equivalence between different physical problems Chimento1 . Cosmological models are equivalent when the corresponding dynamical equations are form invariant under the action of that group. Hence, it will be interesting to investigate the consequences of this group when the dimension of the space time is taken to be a free parameter of the theory. Notice that the multidimensional point of view has been used in general relativity to extract information or to endow with properties fields and/or physical systems belonging to spaces of different dimensions.
In this sense, the purpose of the present work is to illustrate how a group of symmetry transformations acts on n– and m– dimensional flat FRW cosmologies which satisfy Einstein equations. This group relates the energy density and the isotropic pressure of the cosmic fluid (source variables) to the expansion rate (geometrical variable) linking two different cosmologies, one of which could be accelerated. Hence, even when the energy density is a dimensional invariant we can get assisted inflation Chimento1 -Chimento1' driven by the freedom associated with the dimension of the space time. In the case of requiring the condition (10) of Ref. Cataldo the linked cosmologies become identical, they share the same scale factor, or there is a duality between contracting and superaccelerated expanding scenarios associated with phantom cosmologies Chimento2 -todos , i.e. the scale factor of one of them is the inverse of that of the other. The above formulation also can be applied to a self–interacting scalar field using its conventional perfect fluid description.
The outline of the present paper is as follows: In Sec. II we review the well known Einstein equations for the FRW metrics in n– and m– dimensional gravities coupled to a perfect fluid. The case for constant barotropic indices is discussed in detail. In Sec. III we briefly review the field equations for the FRW metrics coupled to a scalar field. In Sec IV some conclusions are given.
## II Dimensional form invariance symmetry in flat FRW spacetimes
We shall assume the spherically symmetric flat FRW metric of an $`n`$–dimensional spacetime given by
$`ds^2=dt^2+a__n(t)^2\left(dr^2+r^2d\mathrm{\Omega }_{n2}^2\right),`$ (1)
where the spherical sector, related to $`n2`$ angular variables $`\theta _i`$, with $`i`$ running from $`1`$ to $`(n2)`$, is determined to be $`d\mathrm{\Omega }_{_{n2}}^{}{}_{}{}^{2}=d\theta _{_1}^{}{}_{}{}^{2}+\mathrm{sin}^2\theta __1\mathrm{d}\theta _{_2}^{}{}_{}{}^{2}+\mathrm{}+\mathrm{sin}^2\theta __1\mathrm{}\mathrm{sin}^2\theta _{_{\mathrm{n}3}}\mathrm{d}\theta _{_{\mathrm{n}2}}^{}{}_{}{}^{2}`$, for $`n3`$. The Einstein equations for an $`n`$–dimensional spacetime are given by
$`G_{_{\alpha \beta }}=R_{_{\alpha \beta }}{\displaystyle \frac{R}{2}}g_{_{\alpha \beta }}=\kappa __nT_{_{\alpha \beta }},`$
where Greek indices run from $`1`$ to $`n`$, and $`\kappa __n`$ stands for the multidimensional gravitational constant.
The independent Einstein equations for the n–dimensional FRW metric (1) filled with a perfect fluid are:
$`G_{t}^{}{}_{}{}^{t}={\displaystyle \frac{(n1)(n2)}{2}}{\displaystyle \frac{\dot{a}__n^2}{a__n^2}}=\kappa __n\rho __n,`$ (2)
$`G_{r}^{}{}_{}{}^{r}=(n2){\displaystyle \frac{\ddot{a}__n}{a__n}}+{\displaystyle \frac{(n2)(n3)}{2}}{\displaystyle \frac{\dot{a}__n^2}{a__n^2}}=\kappa __np__n,`$ (3)
where $`\kappa __n`$, $`a__n`$, $`\rho __n`$ and $`p__n`$ are the gravitational constant, scale factor, energy density and the pressure in an n–dimensional spacetime respectively. Dots denote differentiation with respect to $`t`$. The dependent Einstein equations are related as $`G_{\theta _{_{n2}}}^{}{}_{}{}^{\theta _{_{n2}}}=\mathrm{}=G_{\theta __1}^{}{}_{}{}^{\theta __1}=G_{r}^{}{}_{}{}^{r}`$.
We can replace Eq. (3) by the conservation equation:
$`\dot{\rho }__n+(n1){\displaystyle \frac{\dot{a}__n}{a__n}}(\rho __n+p__n)=0,`$ (4)
which, as is well known, is derivable from the equation $`T_{}^{\alpha \beta }{}_{_{;\beta }}{}^{}=0`$. Thus the Einstein equations for an n–dimensional flat FRW cosmology are given by Eq (2) and Eq (4), which we shall rewrite in the form:
$`\alpha H^{\mathrm{\hspace{0.17em}2}}=\kappa \rho ,\dot{\rho }+\beta H(\rho +p)=0,`$ (5)
where $`\alpha =(n1)(n2)/2`$, $`H=\dot{a}/a`$, $`\beta =(n1)`$ and we have omitted the subindex $`n`$.
For a different m–dimensional flat FRW cosmology the Einstein equations are given by
$`\overline{\alpha }\overline{H}^{\mathrm{\hspace{0.17em}2}}=\overline{\kappa }\overline{\rho },\dot{\overline{\rho }}+\overline{\beta }\overline{H}(\overline{\rho }+\overline{p})=0,`$ (6)
where $`\overline{\alpha }=(m1)(m2)/2`$, $`\overline{H}=\dot{\overline{a}}/\overline{a}`$, $`\overline{\beta }=(m1)`$, and $`\overline{a}`$, $`\overline{\kappa }`$, $`\overline{\rho }`$ and $`\overline{p}`$ are the scale factor, gravitational constant, energy density and the pressure in an m–dimensional spacetime respectively.
By “invariant form” we shall mean that the system of equations (5) transform into Eqs. (6) under the symmetry transformations:
$`\overline{\rho }=\overline{\rho }(\rho ),`$ (7)
$`\overline{H}=\pm \theta \sqrt{{\displaystyle \frac{\overline{\rho }}{\rho }}}H,`$ (8)
$`\overline{\rho }+\overline{p}=\pm {\displaystyle \frac{\beta }{\overline{\beta }\theta }}\sqrt{{\displaystyle \frac{\rho }{\overline{\rho }}}}{\displaystyle \frac{d\overline{\rho }}{d\rho }}(\rho +p),`$ (9)
where $`\theta =(\alpha \overline{\kappa }/\overline{\alpha }\kappa )^{1/2}`$ and $`\overline{\rho }=\overline{\rho }(\rho )`$ is an invertible function. Notice that always $`\theta ^2=\alpha \overline{\kappa }/\overline{\alpha }\kappa 0`$ since $`n3`$ and $`m3`$, and that these form invariant transformations are defined without imposing any restriction on the cosmic fluid. When the dimension of both cosmologies coincides, then we have $`\alpha =\overline{\alpha }`$, $`\beta =\overline{\beta }`$, $`\theta =1`$, and these transformations reduce to that of Ref. Chimento1 and are independent of the dimension where the cosmic fluid “lives”.
The invariant quantities associated with the set of transformations (7)-(9) are
$`{\displaystyle \frac{\overline{\alpha }\overline{H}^2}{\overline{\kappa }\overline{\rho }}}={\displaystyle \frac{\alpha H^2}{\kappa \rho }},`$ (10)
$`{\displaystyle \frac{d\overline{\rho }}{\overline{\beta }\overline{H}(\overline{\rho }+\overline{p})}}={\displaystyle \frac{d\rho }{\beta H(\rho +p)}}.`$ (11)
The first invariant expresses that the expansion of the universe is proportional to the multidimensional gravitational constant and to the energy density contained in the universe. However, the expansion dims with the dimension of the space time because it is proportional to the factor $`1/\alpha `$. The second invariant expresses the fact that the transformations do not modify the cosmic time.
In the case of considering perfect fluids with equations of state $`p=(\gamma 1)\rho `$ and $`\overline{p}=(\overline{\gamma }1)\overline{\rho }`$ in n and m–dimensional spacetimes respectively, we conclude that the barotropic indices $`\gamma `$ and $`\overline{\gamma }`$ transform as
$`\overline{\gamma }={\displaystyle \frac{\overline{\rho }+\overline{p}}{\overline{\rho }}}=\pm {\displaystyle \frac{\beta }{\overline{\beta }\theta }}\left({\displaystyle \frac{\rho }{\overline{\rho }}}\right)^{3/2}{\displaystyle \frac{d\overline{\rho }}{d\rho }}\gamma `$ (12)
under the symmetry transformations (7)–(9). In what follows, the upper and the lower signs will be referred to as the $`(+)`$ and $`()`$ branches respectively.
These general form–invariant transformations relate cosmologies in two different dimensions. For instance, they can be used for generating a new m–dimensional FRW cosmology from a given cosmology in (3+1)–dimensions (with $`n=4`$), or in (2+1)–dimensions (with $`n=3`$), where a lot of them are known.
In this direction we investigate the consequences of a simple example generated by the following transformation between energy densities:
$`\overline{\rho }=b^2\rho ,`$ (13)
with $`b`$ a positive constant. Inserting the latter in (8) and (12) we find that $`a`$ and $`\overline{a}`$ are related to each other by
$`\overline{a}=a^{\pm b\theta },b\theta =\pm {\displaystyle \frac{\beta \gamma }{\overline{\beta }\overline{\gamma }}},`$ (14)
where without loss of generality the constant of proportionality has been set equal to unity. Hence, the deceleration parameter $`q(t)=H^2\ddot{a}/a`$ transforms as
$$\overline{q}=1\pm \frac{1}{b\theta }(q+1).$$
(15)
When the energy density is a dimensional invariant, i.e. for the condition $`\overline{\rho }=\rho `$ or $`b=1`$, we get the relation $`\theta =\pm \beta \gamma /(\overline{\beta }\overline{\gamma })`$ which may be interpreted as a constraint for the barotropic indices $`\overline{\gamma }`$ and $`\gamma `$, since the pressures are not the same in both dimensions Nota . In this case an expanding universe with a positive deceleration parameter, (+) branch, transforms into an accelerated one if $`\theta `$ is taken to be large enough. This means that by adequately selecting the dimension of the space time we can get assisted inflation. For instance, for constant $`\gamma `$ and $`\overline{\gamma }`$, the Einstein equations lead to power law solutions:
$`\overline{a}=t^{2/\overline{\beta }\overline{\gamma }}=t^{\pm 2\theta /\beta \gamma }=a^{\pm \theta },`$
after using Eq. (14) for $`b=1`$. Then, for the (+) branch and $`2\theta >\beta \gamma `$ we obtain an accelerated expansion.
It is interesting to investigate the choice $`b\theta =1`$, because from (14) we have $`\overline{a}=a^{\pm 1}`$. Now we pay attention to the $`()`$ branch since in this case the symmetry transformation $`aa^1`$ (duality) maps the initial singularity at $`t=0`$, $`a(0)=0`$, into other kind of singularity $`\overline{a}(0)=\mathrm{}`$, i.e, the scale factor $`\overline{a}`$ and the scalar curvature $`\overline{R}`$ diverge at a finite time. In particular, for $`\gamma >0`$ the $`()`$ branch of the power law solution $`\overline{a}=(t)^{2/\beta \gamma }`$ defined for $`t<0`$ diverges in the future at $`t=0`$. This kind of singularity dubbed “big rip” is a characteristic of phantom or ghost cosmologies. Hence, using the condition $`\beta \gamma =\overline{\beta }\overline{\gamma }`$, we obtain the relation between the n– and m– dimensional flat FRW cosmologies:
$$\frac{\overline{\gamma }\overline{\kappa }}{m2}=\frac{\gamma \kappa }{n2},$$
(16)
defining the phantom sector of our model.
The case of considering structural invariance of the scale factors $`a(t)`$ in n– and m–dimensional FRW cosmologies, i.e. dimensional invariance of the scale factor as it was assumed in Ref. Cataldo , corresponds to selecting the (+) branch of the transformations (7)–(9) generated by Eq. (13) with $`b\theta =1`$. Thus, the energy densities corresponding to n– and m–dimensional FRW cosmologies are related by $`\overline{\rho }=\rho /\theta ^2`$ and barotropic indexes transform as
$`(n1)\gamma `$ $`(m1)\overline{\gamma },`$ (17)
where we have used the notation of Ref. Cataldo . Note that the results we have obtained by applying the transformations (7)–(9) enlarge those of Ref. Cataldo and add the duality between contracting and expanding cosmologies through the $`()`$ branch of the transformations which was not considered in the previous paper.
Finally, when the dimension of both cosmologies coincides we have $`\alpha =\overline{\alpha }`$, $`\beta =\overline{\beta }`$, $`\theta =1`$, and Eqs. (7)–(9) (or Eqs. (6)–(8) of Ref. Chimento1 ) are independent of the dimension where the cosmic fluid “lives”.
## III The scalar field case
Let us consider a self–interacting scalar field $`\varphi `$ driven by a potential $`V(\varphi )`$ having an associated perfect fluid energy tensor with energy density and pressure given by
$`\rho ={\displaystyle \frac{1}{2}}\dot{\varphi }^2+V(\varphi ),p={\displaystyle \frac{1}{2}}\dot{\varphi }^2V(\varphi ),`$ (18)
$`\overline{\rho }={\displaystyle \frac{1}{2}}\dot{\overline{\varphi }}^2+\overline{V}(\overline{\varphi }),\overline{p}={\displaystyle \frac{1}{2}}\dot{\overline{\varphi }}^2\overline{V}(\overline{\varphi }),`$ (19)
in n–and m–dimensional FRW space respectively. Now, Eqs. (18)-(19) along with Eqs. (7)-(9) give the rules of transformation for $`\varphi `$ and $`V`$
$`\dot{\overline{\varphi }}^2=\overline{\rho }+\overline{p}=\pm {\displaystyle \frac{\beta }{\overline{\beta }\theta }}\sqrt{{\displaystyle \frac{\rho }{\overline{\rho }}}}{\displaystyle \frac{d\overline{\rho }}{d\rho }}\dot{\varphi }^2,`$ (20)
$`\overline{V}=\overline{\rho }{\displaystyle \frac{1}{2}}{\displaystyle \frac{\beta }{\overline{\beta }\theta }}\sqrt{{\displaystyle \frac{\rho }{\overline{\rho }}}}{\displaystyle \frac{d\overline{\rho }}{d\rho }}\dot{\varphi }^2.`$ (21)
To illustrate an application of the latter we assume the transformation (13)-(14), then we obtain
$`\dot{\overline{\varphi }}^2=\pm {\displaystyle \frac{\beta b}{\overline{\beta }\theta }}\dot{\varphi }^2,`$ (22)
$`\overline{V}={\displaystyle \frac{b}{2}}\left[b{\displaystyle \frac{\beta }{\overline{\beta }\theta }}\right]\dot{\varphi }^2+b^2V(\varphi ).`$ (23)
In addition, the scale factor transforms according to Eq. (14). Notice that the $`(+)`$ branch gives the m–dimensional analog of the n–dimensional original cosmological model, while the $`()`$ branch leads us, as in the previous section, to phantom cosmologies. The transformed scalar fields are related to the original one by a dimensional generalization of the transformation considered in Ref. Chimento2 . It represents a generalization of the Wick rotation.
There is an interesting case to be investigated, for instance, let us consider the restricted group of transformations defined by the condition $`\overline{V}V`$. Then, from Eq. (23) we get $`V\dot{\varphi }^2`$ and $`\rho \dot{\varphi }^2`$. In this case Eq. (5) can be solved by assuming a power law scale factor with a scalar field of the form $`\varphi \mathrm{ln}t`$. The final solution is
$`a=t^{2/\beta \gamma },`$ (24)
$`V={\displaystyle \frac{(2\gamma )\varphi _0^2}{2\gamma }}e^{2\varphi /\varphi _0},\varphi =\varphi _0\mathrm{ln}t,`$ (25)
where $`\varphi _0=(2/\beta )(\alpha /\kappa \gamma )^{1/2}`$. These equations represent the dimensional generalization of the ordinary exponential potential and its associated power law solutions. The respective solution in the m-dimensional flat FRW space time is obtained by inserting the above n-dimensional solution (25) into the transformations (22) and (23); so a straightforward calculation gives
$$\overline{V}=\frac{(2\overline{\gamma })\overline{\varphi _0}^2}{2\overline{\gamma }}e^{2\overline{\varphi }/\overline{\varphi }_0},\varphi =\overline{\varphi }_0\mathrm{ln}t,$$
(26)
where we have used Eq. (14) and $`\overline{\varphi }_0=b\varphi _0(\overline{\gamma }/\gamma )^{1/2}`$. The scale factor is given by Eqs. (14) and (24):
$$\overline{a}=t^{2/\overline{\beta }\overline{\gamma }}=t^{\pm 2b\theta /\beta \gamma }=a^{\pm b\theta }.$$
(27)
This example shows that the form invariant transformations can be used to generate new cosmological solutions in an m-dimensional gravity from a seed one in an n-dimensional gravity.
## IV Conclusions
The main goal of the present work is to illustrate how a group of symmetry transformations acts on n– and m– dimensional flat FRW cosmologies which satisfy Einstein equations. Cosmological models are equivalent since the corresponding dynamical equations become form invariant under the action of this group. For two different cosmologies, i.e. n– and m– dimensional flat FRW metrics, this group relates their energy densities, isotropic pressures and the scale factors to generic dimensional parameters $`\alpha `$, $`\beta `$, $`\overline{\alpha }`$ and $`\overline{\beta }`$. If the dimension of both cosmologies coincides, i.e. $`n=m`$, then the group of symmetry transformations relates their energy densities, isotropic pressures and the scale factors only. In addition, a form invariant symmetry transformation which violates the dominant energy condition induces a duality between contracting and superaccelerated expanding scenarios generating phantom cosmologies. All these multidimensional considerations can also be formulated for the scalar field associating a perfect fluid description with the stress energy tensor.
Finally, these general form–invariance transformations can be considered as an algorithm for generating a new m–dimensional FRW cosmology from a known n–dimensional cosmology. For instance, we can use as a seed solution one given known cosmology in (3+1)–dimensional gravity, where there exist a lot of solutions.
## V acknowledgements
One of the authors (MC) thanks A.A. García for valuable comments. The authors thank Paul Minning for carefully reading this manuscript. This work was partially supported by CONICYT through grants FONDECYT N<sup>0</sup> 1051086, 1030469 and 1040624 (MC), and Project X224 by the University of Buenos Aires (LPC). It was also supported by the Dirección de Investigación de la Universidad del Bío–Bío (MC) and Consejo Nacional de Investigaciones Científicas y Técnicas (LPC). |
warning/0506/astro-ph0506142.html | ar5iv | text | # Upper limits on gravitational waves emission in association with the Dec 27 2004 giant flare of SGR1806-20
($`^{\colorbox[rgb]{1,1,1}{$$}}`$Corresponding author: cerdonio@pd.infn.it)
## Abstract
At the time when the giant flare of SGR1806-20 occurred, the AURIGA “bar” gw detector was on the air with a noise performance close to stationary gaussian. This allows to set relevant upper limits, at a number of frequencies in the vicinities of 900 $`\colorbox[rgb]{1,1,1}{$H$}\colorbox[rgb]{1,1,1}{$z$}`$, on the amplitude of the damped gw wave trains, which, according to current models, could have been emitted, due to the excitation of normal modes of the star associated with the peak in X-rays luminosity.
PACS : 04.80.Nn, 95.55.Ym
On 27 December 2004 the Soft Gamma-ray Repeater SGR1806-20 gave a giant flare, which was observed by a number of instruments Borkowsky .
The fluence, if the emission is assumed isotropic, at the distance of $`\colorbox[rgb]{1,1,1}{$d$}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$15$}\colorbox[rgb]{1,1,1}{$k$}\colorbox[rgb]{1,1,1}{$p$}\colorbox[rgb]{1,1,1}{$c$}`$ would imply an energy some hundred times larger than any other known giant flare Hurley ; Terasawa . Soft gamma-ray repeaters are thought to be magnetars (see Hurley and refs. therin). It has been suggested Hurley ; Schwartz that the extreme energy event of 27 December 2004 is due to a catastrophic instability involving global crustal failure and magnetic reconnection Thomson1 . Observations by CLUSTER and TC-2, in combination of data from GEOTAIL, gave evidence that the steep initial rise contains two exponential phases, of e-folding times 4.9 $`\colorbox[rgb]{1,1,1}{$m$}\colorbox[rgb]{1,1,1}{$s$}`$ and 67 $`\colorbox[rgb]{1,1,1}{$m$}\colorbox[rgb]{1,1,1}{$s$}`$ respectively, which covered the 24 $`\colorbox[rgb]{1,1,1}{$m$}\colorbox[rgb]{1,1,1}{$s$}`$ before the time of the peak intensity $`\colorbox[rgb]{1,1,1}{$t$}_\colorbox[rgb]{1,1,1}{$p$}`$; all the timescales support the notion of a sudden reconfiguration of the stars magnetic field, producing large fractures in the crust Schwartz . In particular these authors remark that the intermediate $`\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$5$}\colorbox[rgb]{1,1,1}{$m$}\colorbox[rgb]{1,1,1}{$s$}`$ time is naturally explained if the rising time is limited by ithe propagation of a triggering fracture of size $`\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$5$}\colorbox[rgb]{1,1,1}{$k$}\colorbox[rgb]{1,1,1}{$m$}`$, as it would be predicted by the theory of reference Thomson2 .
According to a few somewhat different models, as a consequence of crustal cracking de Freitas or reconfiguration of the moment of inertia tensor Ioka , non-radial kHz oscillation modes of the neutron star would be excited, giving emission of gravitational waves (gw), possibly at frequencies where the gw bar detector AURIGA AURIGA is sensitive (see insert of Fig. 1). Both the above quoted models predict gw emission, starting very close to $`\colorbox[rgb]{1,1,1}{$t$}_\colorbox[rgb]{1,1,1}{$p$}`$, which involves $`\colorbox[rgb]{1,1,1}{$k$}\colorbox[rgb]{1,1,1}{$H$}\colorbox[rgb]{1,1,1}{$z$}`$ non radial modes of oscillation of a neutron star with few hundred $`\colorbox[rgb]{1,1,1}{$m$}\colorbox[rgb]{1,1,1}{$s$}`$ damping time. The expected waveforms can be approximately parametrized as $`\colorbox[rgb]{1,1,1}{$h$}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$t$}\colorbox[rgb]{1,1,1}{$)$}\colorbox[rgb]{1,1,1}{$=$}\colorbox[rgb]{1,1,1}{$h$}_\colorbox[rgb]{1,1,1}{$0$}\colorbox[rgb]{1,1,1}{$\mathrm{exp}$}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$t$}\colorbox[rgb]{1,1,1}{$/$}\colorbox[rgb]{1,1,1}{$\tau $}_\colorbox[rgb]{1,1,1}{$s$}\colorbox[rgb]{1,1,1}{$)$}\colorbox[rgb]{1,1,1}{$\mathrm{sin}$}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$2$}\colorbox[rgb]{1,1,1}{$\pi $}\colorbox[rgb]{1,1,1}{$f$}_\colorbox[rgb]{1,1,1}{$s$}\colorbox[rgb]{1,1,1}{$t$}\colorbox[rgb]{1,1,1}{$)$}`$, where $`\colorbox[rgb]{1,1,1}{$h$}_\colorbox[rgb]{1,1,1}{$0$}`$ is the maximum gw amplitude, $`\colorbox[rgb]{1,1,1}{$f$}_\colorbox[rgb]{1,1,1}{$s$}`$ and $`\colorbox[rgb]{1,1,1}{$\tau $}_\colorbox[rgb]{1,1,1}{$s$}`$ are the frequencies and damping times of normal modes; the polarization of the wave is not known. The frequencies of the various modes are still under study and depend on a variety of factors as EOS, temperature, density, age, rotational state of the star, etc. miniutti so that we are unable to anticipate with any confidence what specific set of gw emission frequencies could be the one expected for a magnetar ready to undergo a supergiant flare. Still the lowest lying modes, $`\colorbox[rgb]{1,1,1}{$g$}\colorbox[rgb]{1,1,1}{$$}`$, $`\colorbox[rgb]{1,1,1}{$f$}\colorbox[rgb]{1,1,1}{$$}`$ and, marginally, $`\colorbox[rgb]{1,1,1}{$p$}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$m$}\colorbox[rgb]{1,1,1}{$o$}\colorbox[rgb]{1,1,1}{$d$}\colorbox[rgb]{1,1,1}{$e$}\colorbox[rgb]{1,1,1}{$s$}`$, could well be in the frequency range $`\colorbox[rgb]{1,1,1}{$500$}\colorbox[rgb]{1,1,1}{$÷$}\colorbox[rgb]{1,1,1}{$1500$}\colorbox[rgb]{1,1,1}{$H$}\colorbox[rgb]{1,1,1}{$z$}`$, depending on the status of the star.
Within a factor 10 in gw amplitude, AURIGA is sensitive to gws from $`\colorbox[rgb]{1,1,1}{$$}`$ 800 to 1050 $`\colorbox[rgb]{1,1,1}{$H$}\colorbox[rgb]{1,1,1}{$z$}`$. Here we limit the analysis to the most sensitive part of the band, namely between $`\colorbox[rgb]{1,1,1}{$850$}`$ and $`\colorbox[rgb]{1,1,1}{$950$}\colorbox[rgb]{1,1,1}{$H$}\colorbox[rgb]{1,1,1}{$z$}`$ (see the insert in Fig. 1), where the detector sensitivity varies no more than a factor 4 in amplitude. Since the last upgrading of the suspensions on Dec $`\colorbox[rgb]{1,1,1}{$2$}^{\colorbox[rgb]{1,1,1}{$n$}\colorbox[rgb]{1,1,1}{$d$}}`$ 2004 the detector is well behaved in the sense that performs stationary gaussian, after epochs of environmental disturbances are vetoed by means of auxiliary channels (i.e. signals at frequencies where the detector is gw insensitive). During nights and week-ends the vetoed epochs becomes less frequent and shorter, so that the detector achieve close to 90% stationary gaussian operation; in particular on time spans of minutes we can use the data, without even applying vetoes. This is the case for the time span of about $`\colorbox[rgb]{1,1,1}{$\pm $}\colorbox[rgb]{1,1,1}{$100$}\colorbox[rgb]{1,1,1}{$s$}`$ around the epoch of the 27 December 2004 giant flare of SGR1806-20, which we use in this analysis.
We show in the following that the noise is driven by a zero mean stochastic gaussian process with a stationary correlation function. For what concern the directional sensitivity, the orientation of AURIGA in respect to the direction of SGR1806-20 was such that the antenna pattern, averaged over polarizations, gave maximal sensitivity at the time of the giant flare. Then we have a unique opportunity to search in our data for gravitational waves emitted at the peak time of the giant flare. We take the peak time $`\colorbox[rgb]{1,1,1}{$t$}_\colorbox[rgb]{1,1,1}{$p$}`$ to be 21:30:26.68 UT of 27 December 2004 after taking into account the time difference between the arrival time at INTEGRAL int and at AURIGA site of $`\colorbox[rgb]{1,1,1}{$133.427$}\colorbox[rgb]{1,1,1}{$m$}\colorbox[rgb]{1,1,1}{$s$}`$.Lichti This time corresponds also to the peak position of the CLUSTER data which show, after the last exponential rise, the evident start of a phase in which damping occurs until the signal gets below $`\colorbox[rgb]{1,1,1}{$1$}\colorbox[rgb]{1,1,1}{$/$}\colorbox[rgb]{1,1,1}{$10$}`$ of the peak value, $`\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$300$}\colorbox[rgb]{1,1,1}{$m$}\colorbox[rgb]{1,1,1}{$s$}`$ after the peak.Schwartz Following the models quoted above, in both cases we can assume the peak time $`\colorbox[rgb]{1,1,1}{$t$}_\colorbox[rgb]{1,1,1}{$p$}`$ as the start of the gw excitation and $`\colorbox[rgb]{1,1,1}{$\tau $}_\colorbox[rgb]{1,1,1}{$s$}`$ = 100 $`\colorbox[rgb]{1,1,1}{$m$}\colorbox[rgb]{1,1,1}{$s$}`$, that is 1/3 of 300 $`\colorbox[rgb]{1,1,1}{$m$}\colorbox[rgb]{1,1,1}{$s$}`$, as the corresponding damping time. In order to extract the signal power first we reconstruct the gw amplitude $`\colorbox[rgb]{1,1,1}{$h$}_\colorbox[rgb]{1,1,1}{$r$}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$t$}\colorbox[rgb]{1,1,1}{$)$}`$ at input through the detector transfer function. Then we slice the gw sensitive frequency band of AURIGA in contiguous and non-overlapping sub-bands $`\colorbox[rgb]{1,1,1}{$f$}_\colorbox[rgb]{1,1,1}{$j$}`$ of constant width $`\colorbox[rgb]{1,1,1}{$\mathrm{\Delta }$}\colorbox[rgb]{1,1,1}{$f$}`$, and centered in $`\colorbox[rgb]{1,1,1}{$f$}_\colorbox[rgb]{1,1,1}{$j$}^\colorbox[rgb]{1,1,1}{$c$}\colorbox[rgb]{1,1,1}{$=$}\colorbox[rgb]{1,1,1}{$f$}_\colorbox[rgb]{1,1,1}{$j$}\colorbox[rgb]{1,1,1}{$+$}\colorbox[rgb]{1,1,1}{$\mathrm{\Delta }$}\colorbox[rgb]{1,1,1}{$f$}\colorbox[rgb]{1,1,1}{$/$}\colorbox[rgb]{1,1,1}{$2$}`$, by means of digital top-hat filters in the frequency domain $`\colorbox[rgb]{1,1,1}{$T$}_\colorbox[rgb]{1,1,1}{$j$}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$f$}\colorbox[rgb]{1,1,1}{$)$}\colorbox[rgb]{1,1,1}{$=$}\colorbox[rgb]{1,1,1}{$\vartheta $}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$|$}\colorbox[rgb]{1,1,1}{$f$}\colorbox[rgb]{1,1,1}{$|$}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$f$}_\colorbox[rgb]{1,1,1}{$j$}\colorbox[rgb]{1,1,1}{$)$}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$\vartheta $}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$|$}\colorbox[rgb]{1,1,1}{$f$}\colorbox[rgb]{1,1,1}{$|$}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$f$}_\colorbox[rgb]{1,1,1}{$j$}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$\mathrm{\Delta }$}\colorbox[rgb]{1,1,1}{$f$}\colorbox[rgb]{1,1,1}{$)$}`$. Within each sub-band, we compute the the equivalent input signal power over a time span $`\colorbox[rgb]{1,1,1}{$\mathrm{\Delta }$}\colorbox[rgb]{1,1,1}{$t$}`$
$$\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$j$}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$$}_{\colorbox[rgb]{1,1,1}{$\mathrm{\Delta }$}\colorbox[rgb]{1,1,1}{$t$}}\colorbox[rgb]{1,1,1}{$T$}_\colorbox[rgb]{1,1,1}{$j$}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$h$}_\colorbox[rgb]{1,1,1}{$r$}^\colorbox[rgb]{1,1,1}{$2$}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$t$}\colorbox[rgb]{1,1,1}{$+$}\colorbox[rgb]{1,1,1}{$k$}\colorbox[rgb]{1,1,1}{$\mathrm{\Delta }$}\colorbox[rgb]{1,1,1}{$t$}\colorbox[rgb]{1,1,1}{$)$}\colorbox[rgb]{1,1,1}{$𝑑$}\colorbox[rgb]{1,1,1}{$t$}\colorbox[rgb]{1,1,1}{$,$}$$
(1)
where $`\colorbox[rgb]{1,1,1}{$$}`$ stands for time convolution. The $`\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$j$}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$t$}\colorbox[rgb]{1,1,1}{$)$}`$ is sampled every $`\colorbox[rgb]{1,1,1}{$\mathrm{\Delta }$}\colorbox[rgb]{1,1,1}{$t$}`$ to construct the time series $`\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$j$}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$k$}\colorbox[rgb]{1,1,1}{$)$}`$ with $`\colorbox[rgb]{1,1,1}{$k$}`$ integer. We decided a priori a fixed partition of the time frequency plane: $`\colorbox[rgb]{1,1,1}{$\mathrm{\Delta }$}\colorbox[rgb]{1,1,1}{$f$}\colorbox[rgb]{1,1,1}{$=$}\colorbox[rgb]{1,1,1}{$1$}\colorbox[rgb]{1,1,1}{$/$}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$2$}\colorbox[rgb]{1,1,1}{$\tau $}_\colorbox[rgb]{1,1,1}{$s$}\colorbox[rgb]{1,1,1}{$)$}\colorbox[rgb]{1,1,1}{$=$}\colorbox[rgb]{1,1,1}{$5$}\colorbox[rgb]{1,1,1}{$H$}\colorbox[rgb]{1,1,1}{$z$}`$ and $`\colorbox[rgb]{1,1,1}{$\mathrm{\Delta }$}\colorbox[rgb]{1,1,1}{$t$}\colorbox[rgb]{1,1,1}{$=$}\colorbox[rgb]{1,1,1}{$201.5$}\colorbox[rgb]{1,1,1}{$m$}\colorbox[rgb]{1,1,1}{$s$}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$2$}\colorbox[rgb]{1,1,1}{$\tau $}_\colorbox[rgb]{1,1,1}{$s$}`$. For each sub-band $`\colorbox[rgb]{1,1,1}{$f$}_\colorbox[rgb]{1,1,1}{$j$}`$, we analyzed the resulting time series of $`\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$j$}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$k$}\colorbox[rgb]{1,1,1}{$)$}`$ over a time span of $`\colorbox[rgb]{1,1,1}{$\pm $}\colorbox[rgb]{1,1,1}{$100$}\colorbox[rgb]{1,1,1}{$s$}`$ around the peak time $`\colorbox[rgb]{1,1,1}{$t$}_\colorbox[rgb]{1,1,1}{$p$}`$ to check the “off source” noise statistics. The $`\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$j$}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$k$}\colorbox[rgb]{1,1,1}{$)$}`$ sample including the peak time $`\colorbox[rgb]{1,1,1}{$t$}_\colorbox[rgb]{1,1,1}{$p$}`$ is then compared to the measured noise statistics, looking for any evidence of excess power. To be more precise, the a priori choice of our sampling time made $`\colorbox[rgb]{1,1,1}{$t$}_\colorbox[rgb]{1,1,1}{$p$}`$ to fall $`\colorbox[rgb]{1,1,1}{$120$}\colorbox[rgb]{1,1,1}{$m$}\colorbox[rgb]{1,1,1}{$s$}`$ after the beginning of the integration time $`\colorbox[rgb]{1,1,1}{$\mathrm{\Delta }$}\colorbox[rgb]{1,1,1}{$t$}`$ of the “on source” sample. Fig. 1 shows how $`\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$j$}`$ fluctuates on the time spans of $`\colorbox[rgb]{1,1,1}{$\pm $}\colorbox[rgb]{1,1,1}{$5$}\colorbox[rgb]{1,1,1}{$s$}`$ around the time of the flare $`\colorbox[rgb]{1,1,1}{$t$}_\colorbox[rgb]{1,1,1}{$p$}`$ for the sub-band $`\colorbox[rgb]{1,1,1}{$f$}_\colorbox[rgb]{1,1,1}{$j$}\colorbox[rgb]{1,1,1}{$=$}\colorbox[rgb]{1,1,1}{$930$}\colorbox[rgb]{1,1,1}{$H$}\colorbox[rgb]{1,1,1}{$z$}`$. A gw emission at frequency $`\colorbox[rgb]{1,1,1}{$f$}_\colorbox[rgb]{1,1,1}{$s$}`$ would give an excess power in the band $`\colorbox[rgb]{1,1,1}{$\mathrm{\Delta }$}\colorbox[rgb]{1,1,1}{$f$}`$ centered at the $`\colorbox[rgb]{1,1,1}{$f$}_\colorbox[rgb]{1,1,1}{$j$}^\colorbox[rgb]{1,1,1}{$c$}`$ such that $`\colorbox[rgb]{1,1,1}{$|$}\colorbox[rgb]{1,1,1}{$f$}_\colorbox[rgb]{1,1,1}{$s$}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$f$}_\colorbox[rgb]{1,1,1}{$j$}^\colorbox[rgb]{1,1,1}{$c$}\colorbox[rgb]{1,1,1}{$|$}\colorbox[rgb]{1,1,1}{$<$}\colorbox[rgb]{1,1,1}{$\mathrm{\Delta }$}\colorbox[rgb]{1,1,1}{$f$}\colorbox[rgb]{1,1,1}{$/$}\colorbox[rgb]{1,1,1}{$2$}`$. The released energy would be maximum in the “on source” sample. The excess signal power in each sub-band $`\colorbox[rgb]{1,1,1}{$\mathrm{\Delta }$}\colorbox[rgb]{1,1,1}{$f$}`$ can be easily calculated from the expected waveform and reads
$`\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$s$}`$ $`\colorbox[rgb]{1,1,1}{$$}`$ $`\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$h$}_\colorbox[rgb]{1,1,1}{$0$}^\colorbox[rgb]{1,1,1}{$2$}\colorbox[rgb]{1,1,1}{$\tau $}_\colorbox[rgb]{1,1,1}{$s$}\colorbox[rgb]{1,1,1}{$/$}\colorbox[rgb]{1,1,1}{$4$}\colorbox[rgb]{1,1,1}{$)$}\colorbox[rgb]{1,1,1}{$\{$}\colorbox[rgb]{1,1,1}{$[$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$1$}}{\colorbox[rgb]{1,1,1}{$2$}}}\colorbox[rgb]{1,1,1}{$+$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$1$}}{\colorbox[rgb]{1,1,1}{$\pi $}}}\colorbox[rgb]{1,1,1}{$\mathrm{tan}$}^\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$1$}\colorbox[rgb]{1,1,1}{$\left($}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$2$}\colorbox[rgb]{1,1,1}{$x$}}{\colorbox[rgb]{1,1,1}{$\delta $}^\colorbox[rgb]{1,1,1}{$2$}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$x$}^\colorbox[rgb]{1,1,1}{$2$}}}\colorbox[rgb]{1,1,1}{$\right)$}\colorbox[rgb]{1,1,1}{$]$}`$ (2)
$`\colorbox[rgb]{1,1,1}{$+$}`$ $`\colorbox[rgb]{1,1,1}{$O$}\colorbox[rgb]{1,1,1}{$\left($}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$1$}}{\colorbox[rgb]{1,1,1}{$f$}_\colorbox[rgb]{1,1,1}{$s$}\colorbox[rgb]{1,1,1}{$\tau $}_\colorbox[rgb]{1,1,1}{$s$}}}\colorbox[rgb]{1,1,1}{$)$}^\colorbox[rgb]{1,1,1}{$2$}\colorbox[rgb]{1,1,1}{$\right\}$}\colorbox[rgb]{1,1,1}{$,$}`$
where $`\colorbox[rgb]{1,1,1}{$x$}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$2$}\colorbox[rgb]{1,1,1}{$\pi $}\colorbox[rgb]{1,1,1}{$\tau $}_\colorbox[rgb]{1,1,1}{$s$}\colorbox[rgb]{1,1,1}{$)$}^\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$1$}\colorbox[rgb]{1,1,1}{$/$}\colorbox[rgb]{1,1,1}{$\mathrm{\Delta }$}\colorbox[rgb]{1,1,1}{$f$}`$ and $`\colorbox[rgb]{1,1,1}{$\delta $}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$|$}\colorbox[rgb]{1,1,1}{$f$}_\colorbox[rgb]{1,1,1}{$s$}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$f$}_\colorbox[rgb]{1,1,1}{$j$}^\colorbox[rgb]{1,1,1}{$c$}\colorbox[rgb]{1,1,1}{$|$}\colorbox[rgb]{1,1,1}{$/$}\colorbox[rgb]{1,1,1}{$\mathrm{\Delta }$}\colorbox[rgb]{1,1,1}{$f$}`$ are the ratios between the signal bandwidth ($`\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$1$}\colorbox[rgb]{1,1,1}{$/$}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$2$}\colorbox[rgb]{1,1,1}{$\pi $}\colorbox[rgb]{1,1,1}{$\tau $}_\colorbox[rgb]{1,1,1}{$s$}\colorbox[rgb]{1,1,1}{$)$}`$) and the detuning of the signal frequency and the bandwidth $`\colorbox[rgb]{1,1,1}{$\mathrm{\Delta }$}\colorbox[rgb]{1,1,1}{$f$}`$, respectively. With our choice of parameters for the analysis, the excess signal power is approximately (within a few % error)
$$\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$s$}\colorbox[rgb]{1,1,1}{$$}\frac{\colorbox[rgb]{1,1,1}{$h$}_\colorbox[rgb]{1,1,1}{$0$}^\colorbox[rgb]{1,1,1}{$2$}\colorbox[rgb]{1,1,1}{$\tau $}_\colorbox[rgb]{1,1,1}{$s$}}{\colorbox[rgb]{1,1,1}{$6$}}\colorbox[rgb]{1,1,1}{$\left[$}\colorbox[rgb]{1,1,1}{$1$}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$\left($}\frac{\colorbox[rgb]{1,1,1}{$f$}_\colorbox[rgb]{1,1,1}{$s$}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$f$}_\colorbox[rgb]{1,1,1}{$j$}^\colorbox[rgb]{1,1,1}{$c$}}{\colorbox[rgb]{1,1,1}{$\mathrm{\Delta }$}\colorbox[rgb]{1,1,1}{$f$}_{\colorbox[rgb]{1,1,1}{$e$}\colorbox[rgb]{1,1,1}{$f$}\colorbox[rgb]{1,1,1}{$f$}}}\colorbox[rgb]{1,1,1}{$\right)$}^\colorbox[rgb]{1,1,1}{$2$}\colorbox[rgb]{1,1,1}{$\right]$}\colorbox[rgb]{1,1,1}{$,$}$$
(3)
where $`\colorbox[rgb]{1,1,1}{$\mathrm{\Delta }$}\colorbox[rgb]{1,1,1}{$f$}_{\colorbox[rgb]{1,1,1}{$e$}\colorbox[rgb]{1,1,1}{$f$}\colorbox[rgb]{1,1,1}{$f$}}\colorbox[rgb]{1,1,1}{$=$}\colorbox[rgb]{1,1,1}{$4$}\colorbox[rgb]{1,1,1}{$H$}\colorbox[rgb]{1,1,1}{$z$}`$. To check the statistics of the “off source” samples, we histogram each time series $`\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$j$}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$k$}\colorbox[rgb]{1,1,1}{$)$}`$ and compare them with the predicted probability density functions assuming gaussian noise, by fitting for the variance separately in each sub-band. The fitting probability density function is a $`\colorbox[rgb]{1,1,1}{$\chi $}^\colorbox[rgb]{1,1,1}{$2$}`$ distribution with $`\colorbox[rgb]{1,1,1}{$\alpha $}`$ effective degrees of freedom $`\colorbox[rgb]{1,1,1}{$p$}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$;$}\colorbox[rgb]{1,1,1}{$\sigma $}^\colorbox[rgb]{1,1,1}{$2$}\colorbox[rgb]{1,1,1}{$)$}\colorbox[rgb]{1,1,1}{$=$}\colorbox[rgb]{1,1,1}{$2$}^{\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$\alpha $}\colorbox[rgb]{1,1,1}{$/$}\colorbox[rgb]{1,1,1}{$2$}}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$/$}\colorbox[rgb]{1,1,1}{$\sigma $}^\colorbox[rgb]{1,1,1}{$2$}\colorbox[rgb]{1,1,1}{$)$}^{\colorbox[rgb]{1,1,1}{$\alpha $}\colorbox[rgb]{1,1,1}{$/$}\colorbox[rgb]{1,1,1}{$2$}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$1$}}\colorbox[rgb]{1,1,1}{$\mathrm{exp}$}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$/$}\colorbox[rgb]{1,1,1}{$2$}\colorbox[rgb]{1,1,1}{$\sigma $}^\colorbox[rgb]{1,1,1}{$2$}\colorbox[rgb]{1,1,1}{$)$}\colorbox[rgb]{1,1,1}{$/$}\colorbox[rgb]{1,1,1}{$\mathrm{\Gamma }$}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$\alpha $}\colorbox[rgb]{1,1,1}{$/$}\colorbox[rgb]{1,1,1}{$2$}\colorbox[rgb]{1,1,1}{$)$}\colorbox[rgb]{1,1,1}{$/$}\colorbox[rgb]{1,1,1}{$\sigma $}^\colorbox[rgb]{1,1,1}{$2$}`$, where $`\colorbox[rgb]{1,1,1}{$\sigma $}^\colorbox[rgb]{1,1,1}{$2$}`$ is the variance of the underlying gaussian stochastic process. We show in Fig. 2 the close agreement with prediction of the data for the frequency bin $`\colorbox[rgb]{1,1,1}{$f$}_\colorbox[rgb]{1,1,1}{$j$}\colorbox[rgb]{1,1,1}{$=$}\colorbox[rgb]{1,1,1}{$930$}\colorbox[rgb]{1,1,1}{$H$}\colorbox[rgb]{1,1,1}{$z$}`$, over the extended time span of $`\colorbox[rgb]{1,1,1}{$\pm $}\colorbox[rgb]{1,1,1}{$100$}\colorbox[rgb]{1,1,1}{$s$}`$.
The results for all other frequency bins are similar. The goodness of the fit has been checked by a $`\colorbox[rgb]{1,1,1}{$\chi $}^\colorbox[rgb]{1,1,1}{$2$}`$ test, and the resulting p-values for all the sub-bands are consistent with a uniform distribution in the unit interval, as expected (see the insert of Fig. 2 and Tab.1). In Table 1 we report the parameter $`\colorbox[rgb]{1,1,1}{$\sigma $}^\colorbox[rgb]{1,1,1}{$2$}`$ of the fit of $`\colorbox[rgb]{1,1,1}{$p$}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$)$}`$ to the experimental data. We find that the dependence on the effective degrees of freedom $`\colorbox[rgb]{1,1,1}{$\alpha $}`$ of the $`\colorbox[rgb]{1,1,1}{$p$}`$-levels is very weak and, within the statistical errors, we can fix $`\colorbox[rgb]{1,1,1}{$\alpha $}\colorbox[rgb]{1,1,1}{$=$}\colorbox[rgb]{1,1,1}{$3.6$}`$ for all the frequency bins. The $`\colorbox[rgb]{1,1,1}{$p$}`$-level distribution is uniform in the unit interval (see the insert of Fig. 2). The stationary behavior, at least for timescales of few minutes, is shown by the constancy in time of the parameters needed to fit the noise model.
We take advantage of the classical theory of hypothesis testing to establish if the samples $`\colorbox[rgb]{1,1,1}{$$}_{\colorbox[rgb]{1,1,1}{$t$}_\colorbox[rgb]{1,1,1}{$p$}}`$ corresponding to the arrival time $`\colorbox[rgb]{1,1,1}{$t$}_\colorbox[rgb]{1,1,1}{$p$}`$ are affected by the presence of a gw signal. To test the null hypothesis $`\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$0$}`$, i.e. that the sample is drawn from the estimated noise probability distribution in absence of signals, we set a threshold $`\colorbox[rgb]{1,1,1}{$$}_{\colorbox[rgb]{1,1,1}{$c$}\colorbox[rgb]{1,1,1}{$r$}}`$ corresponding to a confidence level (C.L.) $`\colorbox[rgb]{1,1,1}{$p$}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$<$}\colorbox[rgb]{1,1,1}{$$}_{\colorbox[rgb]{1,1,1}{$c$}\colorbox[rgb]{1,1,1}{$r$}}\colorbox[rgb]{1,1,1}{$)$}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$1$}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$p$}_{\colorbox[rgb]{1,1,1}{$c$}\colorbox[rgb]{1,1,1}{$r$}}`$. The threshold for $`\colorbox[rgb]{1,1,1}{$1$}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$p$}_{\colorbox[rgb]{1,1,1}{$c$}\colorbox[rgb]{1,1,1}{$r$}}\colorbox[rgb]{1,1,1}{$=$}\colorbox[rgb]{1,1,1}{$95$}`$ % C.L. corresponds to $`\colorbox[rgb]{1,1,1}{$$}_{\colorbox[rgb]{1,1,1}{$c$}\colorbox[rgb]{1,1,1}{$r$}}\colorbox[rgb]{1,1,1}{$=$}\colorbox[rgb]{1,1,1}{$8.8$}\colorbox[rgb]{1,1,1}{$\times $}\colorbox[rgb]{1,1,1}{$\sigma $}_\colorbox[rgb]{1,1,1}{$j$}^\colorbox[rgb]{1,1,1}{$2$}`$. Thus one sees from Table 1 that no excess of gw power is found at $`\colorbox[rgb]{1,1,1}{$t$}_\colorbox[rgb]{1,1,1}{$p$}`$ and therefore we have to set up upper limits. We set conservative confidence intervals for $`\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$s$}`$ using a confidence belt construction hagi which ensures non-uniform coverage greater or equal to 90%. The confidence belt construction proceeds as follows. Assume that the signal magnitude is $`\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$s$}`$. The measured $`\colorbox[rgb]{1,1,1}{$$}`$ in each sub-band (Eq. 1) obeys a non-central $`\colorbox[rgb]{1,1,1}{$\chi $}^\colorbox[rgb]{1,1,1}{$2$}`$ distribution with central parameter equal to $`\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$s$}\colorbox[rgb]{1,1,1}{$/$}\colorbox[rgb]{1,1,1}{$\sigma $}^\colorbox[rgb]{1,1,1}{$2$}`$ (here we drop for simplicity the index of the sub-band). Its corresponding probability density function can be written as:
$`\colorbox[rgb]{1,1,1}{$p$}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$;$}\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$s$}\colorbox[rgb]{1,1,1}{$,$}\colorbox[rgb]{1,1,1}{$\sigma $}\colorbox[rgb]{1,1,1}{$)$}`$ $`\colorbox[rgb]{1,1,1}{$=$}`$ $`{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$1$}}{\colorbox[rgb]{1,1,1}{$2$}\colorbox[rgb]{1,1,1}{$\sigma $}^\colorbox[rgb]{1,1,1}{$2$}}}\colorbox[rgb]{1,1,1}{$\mathrm{exp}$}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$+$}\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$s$}}{\colorbox[rgb]{1,1,1}{$2$}\colorbox[rgb]{1,1,1}{$\sigma $}^\colorbox[rgb]{1,1,1}{$2$}}}\colorbox[rgb]{1,1,1}{$)$}\colorbox[rgb]{1,1,1}{$\left($}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$$}}{\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$s$}}}\colorbox[rgb]{1,1,1}{$\right)$}^{\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$\alpha $}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$2$}\colorbox[rgb]{1,1,1}{$)$}\colorbox[rgb]{1,1,1}{$/$}\colorbox[rgb]{1,1,1}{$4$}}\colorbox[rgb]{1,1,1}{$\times $}`$ (4)
$`\colorbox[rgb]{1,1,1}{$\times $}`$ $`\colorbox[rgb]{1,1,1}{$I$}_{\colorbox[rgb]{1,1,1}{$\alpha $}\colorbox[rgb]{1,1,1}{$/$}\colorbox[rgb]{1,1,1}{$2$}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$1$}}\colorbox[rgb]{1,1,1}{$\left($}\sqrt{\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$s$}\colorbox[rgb]{1,1,1}{$/$}\colorbox[rgb]{1,1,1}{$\sigma $}^\colorbox[rgb]{1,1,1}{$2$}}\colorbox[rgb]{1,1,1}{$\right)$}\colorbox[rgb]{1,1,1}{$,$}`$
where $`\colorbox[rgb]{1,1,1}{$I$}_\colorbox[rgb]{1,1,1}{$k$}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$x$}\colorbox[rgb]{1,1,1}{$)$}`$ are the modified Bessel functions of the first kind of order $`\colorbox[rgb]{1,1,1}{$k$}`$. The $`\colorbox[rgb]{1,1,1}{$q$}`$-quantile of this distribution, $`\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$q$}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$s$}\colorbox[rgb]{1,1,1}{$,$}\colorbox[rgb]{1,1,1}{$\sigma $}\colorbox[rgb]{1,1,1}{$)$}`$, is implicitly defined by $`\colorbox[rgb]{1,1,1}{$q$}\colorbox[rgb]{1,1,1}{$=$}\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$0$}^\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$q$}\colorbox[rgb]{1,1,1}{$p$}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$;$}\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$s$}\colorbox[rgb]{1,1,1}{$,$}\colorbox[rgb]{1,1,1}{$\sigma $}\colorbox[rgb]{1,1,1}{$)$}\colorbox[rgb]{1,1,1}{$𝑑$}\colorbox[rgb]{1,1,1}{$$}`$. For each value of the unknown $`\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$s$}`$ we define the 95% confidence belt boundaries $`\colorbox[rgb]{1,1,1}{$$}_{\colorbox[rgb]{1,1,1}{$h$}\colorbox[rgb]{1,1,1}{$i$}}`$ and $`\colorbox[rgb]{1,1,1}{$$}_{\colorbox[rgb]{1,1,1}{$l$}\colorbox[rgb]{1,1,1}{$o$}\colorbox[rgb]{1,1,1}{$w$}}`$ as
$$\begin{array}{c}\colorbox[rgb]{1,1,1}{$$}_{\colorbox[rgb]{1,1,1}{$h$}\colorbox[rgb]{1,1,1}{$i$}}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$s$}\colorbox[rgb]{1,1,1}{$,$}\colorbox[rgb]{1,1,1}{$\sigma $}\colorbox[rgb]{1,1,1}{$)$}\colorbox[rgb]{1,1,1}{$=$}\colorbox[rgb]{1,1,1}{$\{$}\begin{array}{cccccccccccccccccccc}\colorbox[rgb]{1,1,1}{$0$}& \colorbox[rgb]{1,1,1}{$\mathrm{if}$}\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$s$}\colorbox[rgb]{1,1,1}{$<$}\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$s$}^{\colorbox[rgb]{1,1,1}{$c$}\colorbox[rgb]{1,1,1}{$r$}}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$\sigma $}\colorbox[rgb]{1,1,1}{$)$}& & & & & & & & & & & & & & & & & & \\ \colorbox[rgb]{1,1,1}{$$}_{\colorbox[rgb]{1,1,1}{$5$}\colorbox[rgb]{1,1,1}{$\%$}}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$s$}\colorbox[rgb]{1,1,1}{$,$}\colorbox[rgb]{1,1,1}{$\sigma $}\colorbox[rgb]{1,1,1}{$)$}& \colorbox[rgb]{1,1,1}{$\mathrm{otherwise}$}& & & & & & & & & & & & & & & & & & \end{array}\hfill \\ \colorbox[rgb]{1,1,1}{$$}_{\colorbox[rgb]{1,1,1}{$l$}\colorbox[rgb]{1,1,1}{$o$}\colorbox[rgb]{1,1,1}{$w$}}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$s$}\colorbox[rgb]{1,1,1}{$,$}\colorbox[rgb]{1,1,1}{$\sigma $}\colorbox[rgb]{1,1,1}{$)$}\colorbox[rgb]{1,1,1}{$=$}\colorbox[rgb]{1,1,1}{$$}_{\colorbox[rgb]{1,1,1}{$95$}\colorbox[rgb]{1,1,1}{$\%$}}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$s$}\colorbox[rgb]{1,1,1}{$,$}\colorbox[rgb]{1,1,1}{$\sigma $}\colorbox[rgb]{1,1,1}{$)$}\hfill \end{array}$$
(5)
where $`\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$s$}^{\colorbox[rgb]{1,1,1}{$c$}\colorbox[rgb]{1,1,1}{$r$}}`$ is implicitely defined by $`\colorbox[rgb]{1,1,1}{$$}_{\colorbox[rgb]{1,1,1}{$5$}\colorbox[rgb]{1,1,1}{$\%$}}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$s$}^{\colorbox[rgb]{1,1,1}{$c$}\colorbox[rgb]{1,1,1}{$r$}}\colorbox[rgb]{1,1,1}{$,$}\colorbox[rgb]{1,1,1}{$\sigma $}\colorbox[rgb]{1,1,1}{$)$}\colorbox[rgb]{1,1,1}{$=$}\colorbox[rgb]{1,1,1}{$$}_{\colorbox[rgb]{1,1,1}{$95$}\colorbox[rgb]{1,1,1}{$\%$}}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$0$}\colorbox[rgb]{1,1,1}{$,$}\colorbox[rgb]{1,1,1}{$\sigma $}\colorbox[rgb]{1,1,1}{$)$}`$. This confidence belt defines a set of confidence intervals on $`\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$s$}`$, whose frequentist coverage is – by construction – 90% for $`\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$s$}\colorbox[rgb]{1,1,1}{$>$}\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$s$}^{\colorbox[rgb]{1,1,1}{$c$}\colorbox[rgb]{1,1,1}{$r$}}`$, and 95% for $`\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$s$}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$s$}^{\colorbox[rgb]{1,1,1}{$c$}\colorbox[rgb]{1,1,1}{$r$}}`$. In other words, for every value of $`\colorbox[rgb]{1,1,1}{$$}_{\colorbox[rgb]{1,1,1}{$t$}_\colorbox[rgb]{1,1,1}{$p$}}`$ from each sub-band, if $`\colorbox[rgb]{1,1,1}{$$}_{\colorbox[rgb]{1,1,1}{$t$}_\colorbox[rgb]{1,1,1}{$p$}}\colorbox[rgb]{1,1,1}{$<$}\colorbox[rgb]{1,1,1}{$$}_{\colorbox[rgb]{1,1,1}{$95$}\colorbox[rgb]{1,1,1}{$\%$}}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$0$}\colorbox[rgb]{1,1,1}{$,$}\colorbox[rgb]{1,1,1}{$\sigma $}_\colorbox[rgb]{1,1,1}{$j$}\colorbox[rgb]{1,1,1}{$)$}`$ we set an upper limit equal to $`\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$s$}^{\colorbox[rgb]{1,1,1}{$c$}\colorbox[rgb]{1,1,1}{$r$}}`$, otherwise our procedure gives a two-sided confidence interval. In all sub-bands we obtain upper limits, which can be written as $`\colorbox[rgb]{1,1,1}{$$}_\colorbox[rgb]{1,1,1}{$s$}^{\colorbox[rgb]{1,1,1}{$c$}\colorbox[rgb]{1,1,1}{$r$}}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$18$}\colorbox[rgb]{1,1,1}{$\times $}\colorbox[rgb]{1,1,1}{$\sigma $}_\colorbox[rgb]{1,1,1}{$j$}^\colorbox[rgb]{1,1,1}{$2$}`$. These limits ranges from $`\colorbox[rgb]{1,1,1}{$$}^{\colorbox[rgb]{1,1,1}{$1$}\colorbox[rgb]{1,1,1}{$/$}\colorbox[rgb]{1,1,1}{$2$}}\colorbox[rgb]{1,1,1}{$=$}\colorbox[rgb]{1,1,1}{$3.5$}\colorbox[rgb]{1,1,1}{$\times $}\colorbox[rgb]{1,1,1}{$10$}^{\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$21$}}\colorbox[rgb]{1,1,1}{$H$}\colorbox[rgb]{1,1,1}{$z$}^{\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$1$}\colorbox[rgb]{1,1,1}{$/$}\colorbox[rgb]{1,1,1}{$2$}}`$ to $`\colorbox[rgb]{1,1,1}{$$}^{\colorbox[rgb]{1,1,1}{$1$}\colorbox[rgb]{1,1,1}{$/$}\colorbox[rgb]{1,1,1}{$2$}}\colorbox[rgb]{1,1,1}{$=$}\colorbox[rgb]{1,1,1}{$1.4$}\colorbox[rgb]{1,1,1}{$\times $}\colorbox[rgb]{1,1,1}{$10$}^{\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$21$}}\colorbox[rgb]{1,1,1}{$H$}\colorbox[rgb]{1,1,1}{$z$}^{\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$1$}\colorbox[rgb]{1,1,1}{$/$}\colorbox[rgb]{1,1,1}{$2$}}`$, according to AURIGA sensitivity.
The initial amplitude of the neutron star normal modes $`\colorbox[rgb]{1,1,1}{$h$}_\colorbox[rgb]{1,1,1}{$0$}`$ is related to $`\colorbox[rgb]{1,1,1}{$$}`$ by Eq. (3) that gives, for the best upper limit, $`\colorbox[rgb]{1,1,1}{$h$}_\colorbox[rgb]{1,1,1}{$0$}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$2.7$}\colorbox[rgb]{1,1,1}{$\times $}\colorbox[rgb]{1,1,1}{$10$}^{\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$20$}}`$. We discuss now the upper limit in terms of the total gw energy $`\colorbox[rgb]{1,1,1}{$ϵ$}_{\colorbox[rgb]{1,1,1}{$g$}\colorbox[rgb]{1,1,1}{$w$}}\colorbox[rgb]{1,1,1}{$=$}\colorbox[rgb]{1,1,1}{$E$}_{\colorbox[rgb]{1,1,1}{$g$}\colorbox[rgb]{1,1,1}{$w$}}\colorbox[rgb]{1,1,1}{$/$}\colorbox[rgb]{1,1,1}{$M$}_{\colorbox[rgb]{1,1,1}{$$}}\colorbox[rgb]{1,1,1}{$c$}^\colorbox[rgb]{1,1,1}{$2$}`$ emitted by the normal modes excitation during the peak of the giant flare of SGR1806-20. The well known formula of the quadrupolar radiation, for the expected gw signal de Freitas , can be written as $`\colorbox[rgb]{1,1,1}{$h$}_\colorbox[rgb]{1,1,1}{$0$}\colorbox[rgb]{1,1,1}{$=$}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$ϵ$}_{\colorbox[rgb]{1,1,1}{$g$}\colorbox[rgb]{1,1,1}{$w$}}\colorbox[rgb]{1,1,1}{$c$}\colorbox[rgb]{1,1,1}{$R$}_\colorbox[rgb]{1,1,1}{$S$}\colorbox[rgb]{1,1,1}{$/$}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$4$}\colorbox[rgb]{1,1,1}{$\pi $}^\colorbox[rgb]{1,1,1}{$2$}\colorbox[rgb]{1,1,1}{$\tau $}_\colorbox[rgb]{1,1,1}{$s$}\colorbox[rgb]{1,1,1}{$)$}^{\colorbox[rgb]{1,1,1}{$1$}\colorbox[rgb]{1,1,1}{$/$}\colorbox[rgb]{1,1,1}{$2$}}\colorbox[rgb]{1,1,1}{$/$}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$f$}_\colorbox[rgb]{1,1,1}{$s$}\colorbox[rgb]{1,1,1}{$d$}\colorbox[rgb]{1,1,1}{$)$}`$, where $`\colorbox[rgb]{1,1,1}{$R$}_\colorbox[rgb]{1,1,1}{$S$}`$ is the Swartzchild radius of one solar mass black hole. Thus the resultant upper limit on $`\colorbox[rgb]{1,1,1}{$ϵ$}_{\colorbox[rgb]{1,1,1}{$g$}\colorbox[rgb]{1,1,1}{$w$}}`$ reads
$`\colorbox[rgb]{1,1,1}{$ϵ$}_{\colorbox[rgb]{1,1,1}{$g$}\colorbox[rgb]{1,1,1}{$w$}}`$ $`\colorbox[rgb]{1,1,1}{$$}`$ $`\colorbox[rgb]{1,1,1}{$3$}\colorbox[rgb]{1,1,1}{$\times $}\colorbox[rgb]{1,1,1}{$10$}^\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$6$}\colorbox[rgb]{1,1,1}{$\left($}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$$}}{\colorbox[rgb]{1,1,1}{$1.3$}\colorbox[rgb]{1,1,1}{$\times $}\colorbox[rgb]{1,1,1}{$10$}^{\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$41$}}\colorbox[rgb]{1,1,1}{$H$}\colorbox[rgb]{1,1,1}{$z$}^\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$1$}}}\colorbox[rgb]{1,1,1}{$\right)$}\colorbox[rgb]{1,1,1}{$\times $}`$ (6)
$`\colorbox[rgb]{1,1,1}{$\times $}`$ $`\colorbox[rgb]{1,1,1}{$\left($}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$15$}\colorbox[rgb]{1,1,1}{$k$}\colorbox[rgb]{1,1,1}{$p$}\colorbox[rgb]{1,1,1}{$c$}}{\colorbox[rgb]{1,1,1}{$d$}}}\colorbox[rgb]{1,1,1}{$\right)$}\colorbox[rgb]{1,1,1}{$\left($}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$930$}\colorbox[rgb]{1,1,1}{$H$}\colorbox[rgb]{1,1,1}{$z$}}{\colorbox[rgb]{1,1,1}{$f$}_\colorbox[rgb]{1,1,1}{$s$}}}\colorbox[rgb]{1,1,1}{$\right)$}^\colorbox[rgb]{1,1,1}{$2$}\colorbox[rgb]{1,1,1}{$\left($}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\tau $}_\colorbox[rgb]{1,1,1}{$s$}}{\colorbox[rgb]{1,1,1}{$0.1$}\colorbox[rgb]{1,1,1}{$s$}}}\colorbox[rgb]{1,1,1}{$\right)$}\colorbox[rgb]{1,1,1}{$.$}`$
We should notice that a gw bar detector has a polarization dependent sensitivity; hence, for an unpolarized or linearly polarized gw, the result in Eq. (6) should be multiplied by a factor $`\colorbox[rgb]{1,1,1}{$2$}`$ or $`\colorbox[rgb]{1,1,1}{$\mathrm{cos}$}^\colorbox[rgb]{1,1,1}{$2$}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$2$}\colorbox[rgb]{1,1,1}{$\psi $}\colorbox[rgb]{1,1,1}{$)$}`$ respectively, where $`\colorbox[rgb]{1,1,1}{$\psi $}`$ is the angle between the bar axis and the polarization of the wave. We conclude that, if the star ever emitted gws from excitation of its normal modes at any of the frequencies studied here, in the time span $`\colorbox[rgb]{1,1,1}{$\mathrm{\Delta }$}\colorbox[rgb]{1,1,1}{$t$}`$ containing the flare time $`\colorbox[rgb]{1,1,1}{$t$}_\colorbox[rgb]{1,1,1}{$p$}`$, the gw amplitudes and energetics are limited as above.
If the giant flare of SGR1806-20 on 27 December 2004 is indeed some 100 times more energetic (however see ref. Yamazaki ) and if the gw luminosity scales with the em luminosity, then, for the frequencies considered, our upper limits come close to the predictions of the models of refs. de Freitas ; Ioka which give an energetics of the order of $`\colorbox[rgb]{1,1,1}{$ϵ$}_{\colorbox[rgb]{1,1,1}{$g$}\colorbox[rgb]{1,1,1}{$w$}}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$5$}\colorbox[rgb]{1,1,1}{$\times $}\colorbox[rgb]{1,1,1}{$10$}^\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$6$}`$. The method used here is of course sub-optimal and the upper limits are somewhat weaker than the “optimal” matched filter. In any case this work shows that, as there is the specific peak time $`\colorbox[rgb]{1,1,1}{$t$}_\colorbox[rgb]{1,1,1}{$p$}`$ to be used as external trigger, it is worth to make searches even with a single detector if its noise is well behaved. An extension of such searches involving the gw detectors on the air in a coincidence search, would also allow to use the information of the gw travel delays between the detectors to select against spuria, and would give the most exhaustive and efficient search, in terms of frequency coverage and confidence in improving the limits, if not to get a candidate detection.
###### Acknowledgements.
We are gratefull to Roberto Turolla for a critical reading of the manuscript. |
warning/0506/hep-ex0506061.html | ar5iv | text | # RECENT RESULTS ON THE 𝐵_𝑐 MESON
## 1 Introduction
The $`B_c`$ is the lowest mass bound state of $`c\overline{b}`$ or $`b\overline{c}`$. It is expected to be a pseudoscalar meson, the last such meson predicted by the Standard Model. The $`B_c`$ is unique in that either one of its quarks can decay, leaving the other as a spectator. One decay path is $`\overline{b}\overline{c}W^+`$, often leading to final states containing a $`J/\psi `$. Although these modes do not have the largest branching ratios, they have proven to be the easiest to observe due to the clean $`J/\psi `$ tag. If the $`c`$ quark decays, leaving the $`b`$ as a spectator, the final states often contain $`B_s`$ mesons. These decay modes have yet to be observed.
There are many predictions for the mass and lifetime of the $`B_c`$. Heavy quark effective potentials predict the mass to be between 6.2 and 6.3 GeV/c<sup>2</sup>. $`^\mathrm{?}`$ Since the $`B_c`$ lies between the well-measured $`b\overline{b}`$ and the $`c\overline{c}`$ states, it would be surprising if it’s mass were very different from the predictions. Expectations for the $`B_c`$ lifetime are more variable, with predictions ranging from 0.4 to 1.4 ps, so that a measurement of the $`B_c`$ lifetime will discriminate between the different models. The $`B_c`$ is expected to have a rich spectroscopy of narrow excited states, which, if observable, would also constrain the heavy quark potential.
The first observation of $`B_c`$ came from CDF in Run I. They observed $`B_c`$ in the semileptonic decay modes $`B_cJ/\psi \mu \nu `$ and $`B_cJ/\psi e\nu `$. This result has been published for some time.$`^\mathrm{?}`$ There are two new results from Run II, which are reported here. D0 has made an observation of the $`B_c`$ in the semileptonic mode $`B_cJ/\psi \mu X`$ $`^\mathrm{?}`$, and CDF has evidence for the exclusive final state $`B_cJ/\psi \pi `$ $`^\mathrm{?}`$. Both of these results are preliminary.
## 2 D0 Observation of $`B_cJ/\psi \mu X`$
The D0 result is based on 260 pb<sup>-1</sup> of data. The analysis proceeds by reconstructing a good $`J/\psi `$ from a dimuon data sample. A good 3-D vertex is required, and a $`J/\psi `$ mass constraint is applied. Then events are selected in which a third muon track can be associated with the $`J/\psi `$ vertex. A background control sample is defined in the same way, except that the third track is required to not be tagged as a muon. Since there is a missing neutrino in the final state, it is not possible to calculate the true proper time. A “pseudo-proper time” is calculated event-by-event, and this distribution is corrected on average based on Monte Carlo simulations.
Backgrounds arise from prompt $`J/\psi `$s which accidently vertex with a third muon, and also from real heavy flavor decays which pick up a third muon. These background samples are treated separately in the analysis, since they will not have the same distributions in mass or proper time. The prompt background sample is obtained by taking events with negative pseuo-proper time (which arise from resolution effects) and reflecting the distribution about $`t=0`$. This sample is then subtracted from the full background sample to obtain the heavy flavor background sample.
Figure 1 shows the result of a combined likelihood fit for the $`B_c`$ mass and lifetime. The left plot shows the distributions in $`M_{\psi \mu }`$, the invariant mass of the $`J/\psi \mu `$ system, and the pseudo-proper time. Three contributions to the full data sample are shown: the prompt background, heavy flavor background, and $`B_c`$ signal. The right plot shows an example of the log likelihood function vs. the $`B_c`$ mass hypothesis, for an assumed lifetime value of 0.45 ps. The fitting procedure gives 95 $`\pm 12(stat)\pm 11(sys)`$ signal events. This result is currently the most statistically significant observation of the $`B_c`$. The best fit value for the mass is $`5.95_{0.13}^{+0.14}(stat)\pm 0.34(sys)`$. The best fit value for the lifetime is $`0.45_{0.10}^{+0.12}(stat)\pm 0.12(sys)`$. The mass and lifetime are uncorrelated in the fit.
## 3 CDF Evidence for $`B_cJ/\psi \pi `$
CDF has recently observed evidence for the $`B_c`$ in the exclusive final state $`B_cJ/\psi \pi `$. The analysis uses a 360 pb<sup>-1</sup> sample of data with $`J/\psi \mu \mu `$ identified at the L3 trigger. A mass constraint is applied to the muons forming the $`J/\psi `$, and a third track, assumed to be a $`\pi `$, is required to form a good 3D vertex with the $`J/\psi `$. The prominent decay $`B^+J/\psi K^+`$ is used as the reference mode to validate the Monte Carlo simulations and check the various cuts. The left plot in figure 2 shows the $`B^+J/\psi K^+`$ mass distribution, demonstrating the low level of background in this reference mode. The right plot of figure 2 shows the $`M_{J/\psi \pi }`$ mass distribution in the mass region used for the $`B_c`$ search.
A blind analysis was performed in the following way. A significance function was defined as $`\mathrm{\Sigma }=\frac{S}{1.5+\sqrt{B}}`$ where $`S`$ represents the number of signal events within $`\pm 2\sigma `$ of the assumed mass, with $`\sigma `$ determined from Monte Carlo. $`B`$ is the number of background events in the same mass region, with the background taken to be linear. In searching for the $`B_c`$ signal, a likelihood fit was done in bins of $`M_{J/\psi \pi }`$ 10 MeV/c<sup>2</sup> wide.
Before the $`M_{J/\psi \pi }`$ distribution in the data was revealed, a Monte Carlo study was performed using 1000 Monte Carlo samples which consisted of only background. In these samples, the background consisted of two components: a combinatorical part, which was taken to be linear, and a contribution from various partically reconstructed $`B_c`$ decays, taken from Monte Carlo. These samples were then fit for the $`B_c`$ signal plus background, and the significance function $`\mathrm{\Sigma }`$ was calculated in all mass bins from 5.7 to 7.2 GeV/c<sup>2</sup>. Occasionally, due to statistical fluctuations, the value of $`\mathrm{\Sigma }`$ could be large. In these 1000 samples, over all mass values, the maximum value of $`\mathrm{\Sigma }`$ that occurred was 3.5. Therefore, this value of $`\mathrm{\Sigma }`$ was taken as the minimum required to claim statistically signficant evidence for $`B_cJ/\psi \pi `$.
This entire procedure was carried out before the $`M_{J/\psi \pi }`$ distribution in the data was revealed. After the requirement on $`\mathrm{\Sigma }`$ had been fixed, the same procedure was carried out on the actual data. Figure 3 shows the value of the significance function $`\mathrm{\Sigma }`$ in each mass bin for the data(left plot). The maximum value of $`\mathrm{\Sigma }`$ observed is 3.6, just above the predetermined cutoff of 3.5, at a mass value around 6.3 GeV/c<sup>2</sup>. Figure 3 also shows an expanded view of the $`M_{J/\psi \pi }`$ distribution, showing the signal region and the fit to the $`B_c`$ mass. The fit returns 18.9 $`\pm 5.7`$ signal events and a mass of 6.2879 $`\pm 0.0048(stat)\pm 0.0011(sys)`$ GeV/c<sup>2</sup>. There is as yet no lifetime determination from this analysis.
## 4 Summary
Table I summarizes the observations of and evidence for the $`B_c`$ meson. All results are in agreement with expectations from heavy quark potential models. Analysis is proceeding for both experiments. Much more data is already in hand and we can expect improved determinations of the $`B_c`$ mass and lifetime soon.
I am grateful to Sherry Towers for providing me with plots and information for the D0 analyisis, and also to Vaia Papadimitriou who was my contact for the CDF analysis.
## References |
warning/0506/hep-th0506172.html | ar5iv | text | # Large D-terms, hierarchical soft spectra and moduli stabilisation
## I Introduction
In most models of supersymmetry (SUSY) breaking, supersymmetry is broken spontaneously in a hidden sector and is then transmitted to the visible sector through some interactions, mostly gravitational. In supergravity, the hidden sector typically contains a set of chiral fields whose auxiliary components attain a vev at the minimum breaking supersymmetry spontaneously. It is generally preferred to have a dynamical explanation to this phenomenon. This breaking is communicated to the visible sector through tree level (and higher order) gravitational interactions. After integrating out the heavy fields, including the hidden sector fields, the resulting effective lagrangian contains renormalisable supersymmetry breaking soft terms nilles ; carlos . At the full supergravity (SUGRA) level, the soft terms are typically given in terms of the gauge kinetic function $`f`$, the Kähler potential $`K`$ and the superpotential $`W`$. Thus in the global limit, the structure of the soft terms crucially depends on the forms these functions take in SUGRA. For example, if the Kähler and the gauge kinetic functions are canonical, this will lead to a universal soft spectrum with mSUGRA boundary conditions.
While analysis of the above type are suitable for simplest classes of supersymmetry breaking models, for more complex situations it is useful to have general expressions for soft terms weldonsoni ; kaplunovskylouis . Such situations can typically arise when supersymmetry breaking has its origins in string theory. Given that we do not yet have a concrete model of supersymmetry breaking in string theory, it is much more advantageous to parameterise this breaking in terms of a few parameters. In terms of effective supergravity lagrangians derived from string theory, the breaking can be parameterised as the vevs of the auxiliary fields of the chiral superfields associated with the higher dimensional gravitational multiplet, namely the dilaton field $`S`$ and the moduli fields $`T_i`$, which effectively act as hidden sector fields. The main advantage of such parameterisations is that they could capture the generic features of soft spectrum emanating from a class of models without completely resorting to explicit model building. These features could then be contrasted with the phenomenological requirements. Detailed analysis parametrising the resultant soft terms for the heterotic case have been presented in ibanez . Recently, they have been further extended to the case of Type-I strings ibaneztypeI .
The above analysis which has been very useful can however, be considered as incomplete. This is because, they have implicitly assumed only $`F`$-type breaking of supersymmetry (only auxiliary fields of the chiral multiplets get a vev). In a more generic scenario, it is well known that there could be D-type susy breaking contributions too dps . These can arise for example in models based on anomalous $`U(1)`$ symmetries bd . Furthermore, in effective lagrangians from the Type II orientifolds with intersecting D-branes, one can expect such D-term contributions to be naturally present. Given these motivations, it is natural to extend the previous analysis by considering D-type SUSY breaking sources. In section two, we present these general expressions of soft terms, initially for generic fields, then for the specific case of the matter fields.
That D-type could have strong impact on the pattern of soft masses has been known for some time, particularly for the limit when the D-terms are small (less than the corresponding F-type contributions), $`\stackrel{<}{}𝒪(m_{3/2}^2)`$ dps . A more dramatic impact could be expected if the D-terms are large and as allowed by the cosmological constant limit, within the range, $`𝒪(m_{3/2}^2)\stackrel{<}{}D\stackrel{<}{}𝒪(m_{3/2}M_{Pl})`$. For example, considering only pure F-type breaking, leads to a typical spectrum of the soft masses, where the gaugino and the Higgsino masses are roughly proportional to the gravitino mass, $`m_{3/2}`$, whereas the scalar mass squared and the $`B`$-terms are proportional to $`m_{3/2}^2`$. Adding large D-type sources could significantly alter this simple pattern by generating a splitting between the fermionic and scalar superpartners, by an amount proportional to $`D`$. In the extreme limit, this would mean that the scalars can have masses close to the intermediate scale. Following the works of Ref.ibanez ; ibaneztypeI , we parameterise the soft terms in three particular cases (section III) : (i) mixed D and anomaly mediation (ii) mixed D and S mediation and (iii) mixed D and T mediation. In the mixed D and anomaly mediated scenario, scalar masses can be everywhere between the weak scale and an intermediate scale whereas the gaugino masses, B-term and the $`\mu `$ term are proportional to the gravitino mass, $`m_{3/2}`$, which can be taken close to the weak scale. In the mixed D and S(T) mediated scenarios, the hierarchy between scalar and fermionic superpartners is parameterised by an angle $`\gamma _{S(T)}`$, which could be constrained by phenomenology.
The splitting due to the D-terms could well have another important application in understanding the origins of recently proposed “split supersymmetry” models. Influenced by multivacua structure in string theory as a possible new view on the cosmological constant problem bp , these models question the solution of the gauge hierarchy problem through low energy SUSY split . In this proposal, not all superpartners are required to be at a scale close to TeV. Instead, it is sufficient if the fermionic superpartners stay close to the weak scale, whereas the scalar superpartners can be present at scales as high as $`10^9`$ GeV. This way, one keeps the nice features of gauge coupling unification and the viable dark matter candidate of low energy supersymmetry, while getting rid of unwanted features associated with large flavour changing neutral current effects and CP violation problemssplitpheno . In section IV, we address the question of attaining split supersymmetry by including D-mediation. As we will see, though it is not automatic to have exact split spectrum in these models, specifically due to the $`B`$ term, we can nevertheless envisage models where it is possible to generate hierarchical spectrum and we will present explicit models of this type.
So far we have not addressed the issue of the origin of such large D-terms. We address this issue in sections V and VI. Unlike in the heterotic case, in Type I/II string theories, Fayet-Iliopoulos terms can appear at the tree level and thus it is possible to generate SUSY breaking with large D-terms. We will present an explicit example in the context of intersecting D-branes Type I orientifold models with four stacks of D9 branes, each stack containing four coincident branes. However, a related question concerns the stabilisation of the moduli as these FI terms are field-dependent. We find that standard mechanisms like gaugino condensation, suitably combined with other mechanisms of moduli stabilisation as, e.g. three-form fluxes in IIB orientifolds, are still applicable even in the limit of large D-terms. We present an example detailing this point. We close with a summary. A preliminary version of our results was reported in dv .
## II General Expressions Including D-breaking
In the following, we will present general expressions for the soft terms including D-type supersymmetry breaking terms. As is the case with any general analysis, we will not address the question of the origins of these SUSY breaking vevs for either F-terms or D-terms. We will assume SUSY to be broken with both these types of breaking and proceed to derive the soft terms. As a starting point, let us recall the form the scalar potential in supergravity<sup>4</sup><sup>4</sup>4Most of the expressions are presented in Planck units, namely, we set $`M_P=1`$. However, at many instances, we keep $`M_P`$ explicitly to make the discussions clearer.:
$$V=e^G(G^MG_M3)+\frac{1}{2}\underset{A}{}g_A^2D_A^2.$$
(1)
Here $`G=K+\mathrm{ln}|W|^2`$, with $`K`$ being the Kähler potential and W, the superpotential and $`1/g_A^2=Ref_A`$, where $`f_A`$ is the gauge kinetic function. The $`F`$ terms in the scalar potential are given by $`G_M=G/z^M`$, where $`z`$ represents the scalar part of a chiral superfield. The index $`M`$ runs over all the chiral superfields present, matter as well as hidden sector and/or moduli fields. The D-terms, $`D_A`$ carry the obvious notation with the index $`A`$ running over all the $`U(1)`$ factors present<sup>5</sup><sup>5</sup>5Note that the D-terms can be explicitly given in terms of the fields, derivatives of the Kähler potential and a FI term. We will make use of this form in a later subsection. For the present, we just note that we consider FI terms to be moduli dependent..
While deriving the soft terms, a couple of constraints need to be satisfied. First, at the minimum, both $`D`$ and $`F`$ terms contribute to supersymmetry breaking and thus to the vacuum energy. This can be canceled by the superpotential (W) vev which gives mass to the gravitino. We will impose this fine-tuning condition on the potential. This means:
$$<V>=<e^G(G^MG_M3)+\frac{1}{2}\underset{A}{}g_A^2D_A^2>=0.$$
(2)
Second is the necessary condition for the existence of the minima: $`<_KV>=<_KV>=0.`$ Here $``$ denotes the covariant derivate on the Kähler manifold defined by $`_KV_M=_KV_M\mathrm{\Gamma }_{KM}^LV_L`$. Using the definition of the potential, eq.(1) and eq.(2), this implies<sup>6</sup><sup>6</sup>6Strictly speaking, there is a contribution proportional to the derivative of the gauge kinetic function in the minimisation condition, eq.(3). As we are concerned with the general expressions for the matter field soft terms, these contributions will be proportional to matter field vevs which are much smaller than the moduli vevs and therefore we will neglect them here.:
$$<e^G(G^M_KG_M+G_K)+\underset{A}{}g_A^2D_A(_KD_A\frac{1}{2}G_KD_A)>=0$$
(3)
We will use the eqs.(2,3) while deriving general expressions for the soft terms. In the present subsection we will not distinguish between the matter and hidden/moduli fields, but present generic expressions for the various scalar couplings in the theory. To start with, we will consider the case of the scalar mass squared matrix, which is defined as
$$M_0^2=\begin{array}{cc}m_{I\overline{J}}^2& m_{IJ}^2\\ m_{\overline{I}\overline{J}}^2& m_{\overline{I}J}^2\end{array},$$
(4)
where the various entries are defined by:
$`m_{I\overline{J}}^2=<_I_{\overline{J}}V>=<_I_{\overline{J}}V>`$ (5)
$`m_{IJ}^2=<_I_JV>=<_I_JV>.`$ (6)
Using the definition of the potential in eq.(1) and the conditions, eqs.(2, 3), we find the most general expressions for the bilinear couplings to be of the form :
$`m_{I\overline{J}}^2`$ $`=`$ $`e^G(G_{I\overline{J}}+_IG^{\overline{K}}_{\overline{J}}G_{\overline{K}}R_{I\overline{J}K\overline{L}}G^KG^{\overline{L}})+{\displaystyle \frac{1}{2}}{\displaystyle \underset{A}{}}g_A^2D_A^2(G_{\overline{J}}G_IG_{I\overline{J}})`$ (7)
$``$ $`{\displaystyle \underset{A}{}}g_A^2D_A(G_{\overline{J}}_ID_A+G_I_{\overline{J}}D_A_I_{\overline{J}}D_A)+{\displaystyle \underset{A}{}}g_A^2_ID_A_{\overline{J}}D_A,`$
$`m_{IJ}^2`$ $`=`$ $`e^G(2_JG_I+G^K_I_JG_K){\displaystyle \underset{A}{}}g_A^2D_A(G_J_ID_A+G_I_JD_A_I_JD_A)`$ (8)
$``$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{A}{}}g_A^2D_A^2(G_IG_J+_IG_J+{\displaystyle \frac{1}{2}}g_A^2_I_Jf_A)+{\displaystyle \underset{A}{}}g_A^2_ID_A_JD_A,`$
where we have neglected the vacuum brackets for simplicity<sup>7</sup><sup>7</sup>7From now on we will neglect vacuum brackets in the rest of the paper, unless and otherwise specified.. Here, $`f_A`$ represents the gauge kinetic function. The term containing the second derivative $`_I_Jf_A`$ gives contributions to the $`B_\mu `$ term from operators of the form $`d^2\theta W^\alpha W_\alpha H_1H_2`$ in superfields. A similar term of the form, $`_If_A_{\overline{J}}\overline{f}_A`$ could contribute to $`m_{I\overline{J}}^2`$. However since this contribution is proportional to the vevs of the matter fields, as we will discuss in the next section, we neglect it here<sup>8</sup><sup>8</sup>8For the same reason, we do not write down the contributions from gauge kinetic function in the $`A_{ijk}`$ term discussed below.. The function, $`R_{I\overline{J}K\overline{L}}`$ represents the Riemann (curvature) tensor of the Kähler manifold whose definition can be found in any of the standard texts wessbagger . The next step would be to derive the expression for the trilinear couplings, which we define as<sup>9</sup><sup>9</sup>9For MSSM fields this definition is equivalent to the naive one of using ordinary derivatives giving the A-term. :
$$A_{IJK}=<_I_J_KV>$$
(9)
which takes the form :
$`A_{IJK}`$ $`=`$ $`e^G(G_K(2_JG_I+G^M_I_JG_M)+G_J(2_KG_I+G^M_I_KG_M)`$ (10)
$`+`$ $`G_I(2_KG_J+G^M_J_KG_M)+2_I_KG_J+_J_KG_I+G^M_I_J_KG_M)`$
$``$ $`g_A^2D_A(_ID_A{\displaystyle \frac{1}{2}}G_ID_A)(G_JG_K+_JG_K)g_A^2D_A(_JD_A{\displaystyle \frac{1}{2}}G_JD_A)(G_IG_K+_IG_K)`$
$``$ $`g_A^2D_A(_KD_A{\displaystyle \frac{1}{2}}G_KD_A)(G_IG_J+_IG_J){\displaystyle \frac{1}{2}}g_A^2D_A^2(G_IG_JG_K+G_I_JG_K`$
$`+`$ $`G_K_IG_J+G_J_IG_K+_I_JG_K)+g_A^2(_I_JD_A_KD_A+_JD^A_I_KD_A`$
$`+`$ $`_ID_A_J_KD_A+D_A_I_J_KD_A).`$
Note that as for the $`B`$ term, there can be contributions to the A-term also from gauge kinetic function, which can be represented by operators of type $`d^2\theta W^\alpha W_\alpha h_{ijk}Q_iQ_jQ_k`$ with $`Q_i`$ representing the matter fields. These are typically of the order $`m_{3/2}^2/M_{Pl}`$ and thus they give negligibly small contributions unless $`m_{3/2}`$ has intermediate scale values. In the next subsection, we will use these expressions to get the expressions of soft masses for the matter fields.
### II.1 General Expressions of soft terms for Matter Fields
We will define matter fields by setting their vevs to zero. This would mean that both the $`F`$ and $`D`$ contributions proportional matter field vevs to be zero at the leading order. Thus, we have:
$$\mathrm{\Phi }^i=0,G^i=0,_iD_A=0,$$
with $`\mathrm{\Phi }`$ representing the scalar part of a matter field. From now on, to distinguish matter and hidden/moduli fields, we denote matter (moduli/hidden sector) fields by using latin(greek) indices. To derive the soft terms for matter fields from the general scalar couplings presented in the previous section, along with using the definitions above, we have to remove the supersymmetric contributions from them. Further, we identify the gravitino mass to be $`m_{3/2}=<e^{G/2}>`$. Taking all these modifications in consideration, the final set of equations are of the form:
$`m_{i\overline{j}}^2`$ $`=`$ $`m_{3/2}^2(G_{i\overline{j}}R_{i\overline{j}\alpha \overline{\beta }}G^\alpha G^{\overline{\beta }}){\displaystyle \frac{1}{2}}{\displaystyle \underset{A}{}}g_A^2D_A^2G_{i\overline{j}}+{\displaystyle \underset{A}{}}g_A^2D_A_i_{\overline{j}}D_A,`$ (11)
$`m_{ij}^2`$ $`=`$ $`m_{3/2}^2(2_iG_j+G^\alpha _i_jG_\alpha ){\displaystyle \frac{1}{2}}{\displaystyle \underset{A}{}}g_A^2D_A^2(_iG_j+{\displaystyle \frac{g_A^2}{2}}_i_jf_A)+{\displaystyle \underset{A}{}}g_A^2D_A_i_jD_A,`$ (12)
$`A_{ijk}`$ $`=`$ $`m_{3/2}^2\left(3_i_jG_k+G^\alpha _i_j_kG_\alpha \right){\displaystyle \frac{1}{2}}{\displaystyle \underset{A}{}}g_A^2D_A^2_i_jG_k+{\displaystyle \underset{A}{}}g_A^2D_A_i_j_kD_A,`$ (13)
$`\mu _{ij}`$ $`=`$ $`m_{3/2}_iG_j,M_{1/2}^A={\displaystyle \frac{1}{2}}(Ref_A)^1m_{3/2}f_{A\alpha }G^\alpha ,`$ (14)
where we have now also supplemented the scalar equations with those for the $`\mu `$ and the gaugino masses. In the above $`f_{A\alpha }=f_A/z^\alpha `$. Note that these expression reduces to the standard formwessbagger ; fabio ; ibanez in the limit where $`D_A`$ goes to zero.
Note that the above soft terms are not in a canonically normalised basis for the kinetic terms. This can be seen from their action which has the form $`g_{i\overline{j}}z^iz^{\overline{j}}m_{i\overline{j}}^2z^iz^{\overline{j}}`$, where $`z^i`$ represents a matter scalar field. To go to the normalised basis, one can define vielbeins such as : $`z^i=e_a^iz^a,z^{\overline{j}}=e_{\overline{b}}^{\overline{j}}z^{\overline{b}}`$ such that $`e_a^ie_{\overline{b}}^{\overline{j}}=g^{i\overline{j}}`$. Using this transformations, we have the soft mass in the normalised basis to be given by
$$\overline{m}_{a\overline{b}}^2=e_a^im_{i\overline{j}}^2e_{\overline{b}}^{\overline{j}},$$
where $`\overline{m}_{a\overline{b}}^2`$ represents the normalised masses. Similar analysis can be extended for other soft terms. In order to keep a compact notation, we however do not present the general expressions in the normalised form.
While these expressions are given for the tree level potential, higher order corrections can play a significant role, depending on the specifics of the model of supersymmetry breaking. In models with small tree-level contributions, the dominant set of corrections are of anomaly mediated typeamsb1 which are proportional to the gravitino mass $`m_{3/2}`$. These contributions are typically not modified in the presence of D-terms and have to be anyway included. For the gauginos, the most general form of these expressions have been presented in amsb2 and are given by
$$M_{1/2}^{}_{}{}^{}A=\frac{g_A^2}{16\pi ^2}\left(3T_G^AT_R^A(T_G^AT_R^A)K_\alpha G^\alpha \frac{2T_R^A}{d_R^A}(\mathrm{log}detK|_R^\mathrm{"})_{,\alpha }G^\alpha \right)m_{3/2}.$$
(15)
Here, $`T_G`$ is the Dynkin index of the adjoint representation, normalised to $`N`$ for $`SU(N)`$, $`T_R`$ is the Dynkin index associated with the representation $`R`$ of dimension $`d_R`$, normalised to $`1/2`$ for the fundamental of $`SU(N)`$ and $`K|_R^\mathrm{"}`$ is the Kähler metric restricted to the representation $`R`$. This expression reduces to the following when all the vevs are much less than $`M_P`$:
$$M_{1/2}^{}_{}{}^{}A=\frac{g_A^2b_0^A}{16\pi ^2}m_{3/2},$$
(16)
where the beta function $`b_0^A`$ was given as $`3T_G^AT_R^A`$ in the previous expression. In addition to the gauginos, the scalar mass terms as well as the B-term and the A-terms receive corrections. In the case of gauginos, as long as the tree-level F-term contributions are present, the anomaly mediated contributions remain sub-dominant, whereas in the case of scalar soft terms, both the D-term as well as the F-term contributions have to be suppressed for the anomaly mediated contributions to dominate.
### II.2 Implications of large D-terms on the soft parameters
Eqs.(11-14) give the modified expressions for the soft terms after including non-zero D-type SUSY breaking contributions in supergravity. Whereas the scalar couplings receive corrections from the D-type terms, the gaugino masses are unaffected by D-mediated effects. The $`\mu `$ term, could be visualised as a soft mass in supergravity by using the Giudice-Masiero mechanismgm . The expression presented in the previous sub-section takes care of this situation and it is seen that $`D`$ terms do not effect the $`\mu `$ term either. However, the exact implications on the soft terms by the inclusion of the D-terms depend on (a) the structure of the D-terms and (b) the magnitude of them. We will address these two issues below.
In the presence of anomalous non-linearly realised abelian gauge symmetries
$`\delta V_A=\mathrm{\Lambda }_A+\overline{\mathrm{\Lambda }}_A,\delta z^i=\mathrm{\Lambda }_AX_A^iz^i𝒱_A^i\mathrm{\Lambda }_A,`$
$`\delta T^\alpha =\eta _A^\alpha \mathrm{\Lambda }_A𝒱_A^\alpha \mathrm{\Lambda }_A,`$ (17)
where $`𝒱_A^I`$ are the Killing potentials, the auxiliary D-terms, defined by
$$_{\overline{J}}D_A=𝒱_A^IK_{\overline{J},I}=𝒱_A^iK_{\overline{J},i}+𝒱_A^\alpha K_{\overline{J},\alpha }$$
(18)
are explicitly given by
$$D_A=z^IX_I^A\frac{K}{z^I}+\xi _A=\overline{z}^{\overline{I}}X_I^A\frac{K}{\overline{z}^{\overline{I}}}+\xi _A,\xi _A\eta _A^\alpha _\alpha K,$$
(19)
where $`X_I^A`$ represents the $`U(1)_A`$ charges of the fields $`z^I`$ and $`\xi _A`$ denotes the Fayet-Iliopoulos term for the $`U(1)_A`$ factors. Note that the equality between the two last terms is a straightforward consequence of the gauge invariance of the Kähler potential. We consider the Fayet-Iliopoulos terms to be moduli dependent and we will not explicitly discuss here the various possible mechanisms of moduli stabilisation<sup>10</sup><sup>10</sup>10After moduli stabilisation, the anomalous $`U(1)`$’s become gauged R-symmetries freedman .. We have in the vacuum, after setting the matter fields vevs to zero
$`_jD_A=\overline{v}_{\overline{\beta }}X_{\overline{\beta }}^AK_{\overline{j}\beta }+\eta _A^{\overline{\alpha }}K_{j\overline{\alpha }}=0,_i_jD_A=0,`$
$`_i_{\overline{j}}D_A`$ $`=`$ $`K_{i\overline{j}}X_i^A+(\overline{v}^{\overline{l}}X_A^{\overline{l}}_{\overline{l}}+\eta _A^{\overline{\alpha }}_{\overline{\alpha }})K_{i\overline{j}},_i_j_lD_A=0`$ (20)
By using (20), the soft terms for the matter fields reduce to
$`m_{i\overline{j}}^2`$ $`=`$ $`m_{3/2}^2\left(G_{i\overline{j}}R_{i\overline{j}\alpha \overline{\beta }}G^\alpha G^{\overline{\beta }}\right)+{\displaystyle \underset{A}{}}g_A^2D_A\left(X_i^A+\overline{v}_{\overline{l}}X_{\overline{l}}^A_{\overline{l}}+\eta _A^{\overline{\alpha }}_{\overline{\alpha }}{\displaystyle \frac{1}{2}}D_A\right)G_{i\overline{j}},`$ (21)
$`m_{ij}^2`$ $`=`$ $`m_{3/2}^2\left(2_iG_j+G^\alpha _i_jG_\alpha \right){\displaystyle \frac{1}{2}}{\displaystyle \underset{A}{}}g_A^2D_A^2(_iG_j+{\displaystyle \frac{g_A^2}{2}}_i_jf_A),`$ (22)
$`A_{ijk}`$ $`=`$ $`m_{3/2}^2\left(3_i_jG_k+G^\alpha _i_j_kG_\alpha \right){\displaystyle \frac{1}{2}}_i_jG_k{\displaystyle \underset{A}{}}g_A^2D_A^2.`$ (23)
Let us now try to quantify how large the D-terms can be. To do this, let us consider the following generic forms for the Kähler and the superpotential :
$`K`$ $`=`$ $`\stackrel{~}{K}(T_\alpha ,T_{\overline{\beta }})+H_{i\overline{j}}(T_\alpha ,T_\beta )Q_iQ_{\overline{j}}^{}+(Z_{ij}(T_\alpha ,T_{\overline{\beta }})Q_iQ_j+h.c)+\mathrm{}`$ (24)
$`W`$ $`=`$ $`Y_{ijk}(T_\alpha )Q_iQ_jQ_k+\stackrel{~}{W}(T_\alpha )+\mathrm{},`$ (25)
where $`T_\alpha `$ represent moduli/hidden sector fields and $`Q_i`$ represent the matter fields. Using these equations let us now revisit the condition (2)
$$m_{3/2}^2\left(K^{\alpha \overline{\beta }}(K_\alpha K_{\overline{\beta }}+\frac{M_P^2}{W}(K_\alpha W_{\overline{\beta }}+W_\alpha K_{\overline{\beta }})+\frac{M_P^4}{|W|^2}W_\alpha W_{\overline{\beta }})3M_P^2\right)+\frac{1}{2}g_A^2D_A^2=0.$$
(26)
From the above we see that, as long as the D-terms are in the limit, $`D𝒪(m_{3/2}^2)`$, they would not contribute significantly to the vacuum energy. However, when they lie within the limit
$$m_{3/2}^2\stackrel{<}{}D_A\stackrel{<}{}m_{3/2}M_P,$$
(27)
they could be contributing significantly. The upper limit is obtained when one assumes D-term contributions to dominate over the F-term contributions or are of the same order as them. This particular limit is what we are interested in the present work as this has not been exploited in a general manner as presented here. From the generic set of soft parameters presented above, it is obvious that splitting between fermionic and scalar superpartners can be ‘naturally’ achieved once the D-terms lie within the above range. Quantitatively, if in a given model the gravitino mass is of $`𝒪`$(1 TeV), the upper limit on the D-term would be of the order of intermediate scale $``$ $`(10^{10})`$ GeV. It is obvious that as one increases the gravitino mass closer to the intermediate scale $`(10^910^{12})`$ GeV, the upper bound on the D-terms become close to the GUT scale. These upper bounds are essentially the magnitude required to cancel the cosmological constant in the limit where the F-terms tend to zero.
Given this limit, let us now try to understand in more detail how large D-terms would generate large splittings between superpartners. The equations for the gaugino and $`\mu `$-term remain unchanged as we have mentioned. The following features of the spectra are easy to extract without actually being specific about the model:
* (i). Scalar Mass Terms: The most dominant contribution to the scalar masses from the $`D`$-terms are the ones which are linear in $`D`$ which for $`m_{3/2}`$ TeV push the scalar masses to intermediate energy scale. Note that these terms depend on the charges of the fields under the additional $`U(1)`$ gauge group, thus putting a constraint that these charges to be of definite sign. If all the three generations of the sfermions have the same charges under the $`U(1)`$ groups, this term would also be universal. Otherwise, there are off-diagonal entries which are generated in the mass matrices, which could of suppressed by some powers in the expansion parameter $`ϵ_\beta =v_\beta /M_P`$, with $`v_\beta `$ representing the vev of some flavon field.
* (ii). Higgs mass terms and the $`B_\mu `$: The Higgs masses follow almost the same requirements as the soft masses. Usually, their charges are linked with the Giudice-Masiero mechanismgm . The $`B_\mu `$ term is however special. Unlike the Higgs mass terms, it does not receive large contributions from D-terms, whose contributions can be utmost of $`𝒪(m_{3/2}^2)`$. If the splitting between the Higgs masses and the $`B_\mu `$ is too large, it could lead to unphysical regions in $`\mathrm{tan}\beta `$. This could be easily seen by noting that
$$\mathrm{sin}2\beta =\frac{2B_\mu }{m_{H_1}^2+m_{H_2}^2+2\mu ^2}.$$
(28)
In the limit of large Higgs mass parameters $`m_{H_1}^2,m_{H_2}^2`$, one has to think of ways to enhance the $`B_\mu `$ term. We will present one such example in the next section.
* (iii). A-terms: Even if the D-terms are large, the A-terms are typically proportional to $`𝒪(m_{3/2})`$. No large enhancement is present. This is expected as A-terms break R-symmetries. They get related to the D-terms due to the constraints of cosmological constant cancellation, but as the scale of R-symmetry breaking is set by the gravitino mass, this naturally sets the A-terms to be of same order.
* (iv). Gaugino Masses: All along we have been commenting that the presence of SUSY breaking D-terms would not change the results for the gaugino masses presented there. This is only true as long as there are no additional fermions in the model. In the presence of additional fermions and non-zero D-terms, gauginos can get Dirac masses through operators of the form nelson
$$h_ad^2\theta \frac{\chi ^aW_\alpha ^aW_X}{M_Pl}=h_a\frac{D_X}{M_P}\psi ^a\lambda ^a+\mathrm{}=m_D^a\psi ^a\lambda ^a+\mathrm{},$$
(29)
where $`\chi ^a`$ represent here fields in the adjoint representation of the Standard Model gauge group with (mirror fermions) which mix with the gauginos and $`m_D^a`$ represent the Dirac mass for the gauginos. These mixing terms could lead to the Majorana masses for the gauginos by a seesaw mechanism $`(m_D^a)^2/M_a`$ if the mirror fermions obtain large R-symmetry breaking Majorana masses, $`M_a`$. In the present work, we do not concentrate on building models of this type.
## III Parametrization of soft terms in Type I/II string models with large D-terms
The soft terms in effective string supergravities from Type-I/II string theories have been parameterized in ibaneztypeI where pure $`F`$-type breaking has been assumed. In the present section we will extend this analysis by considering D-type SUSY breaking terms too. In each of this case, we present parameterizations of the soft terms which could be readily be useful for phenomenological studies.
### III.1 D-dominated supersymmetry breaking
The first case we consider is that of a scenario where F-terms are absent or negligible. We assume that supersymmetry breaking is achieved by pure D-terms. However, we will still require that the gravitino get a mass. This would enable us to cancel the cosmological constant even in the pure D-breaking limit<sup>11</sup><sup>11</sup>11An earlier proposal in this direction has been presented in dvalipomarol .. The scale of the gravitino mass is assumed to be not very far from the weak scale. With these conditions, the potential, eq.(1) takes the form:
$$V=\frac{1}{2}\underset{A}{}g_A^2D_A^23m_{3/2}^2M_P^2.$$
(30)
It is obvious from the above equation that requiring that the potential should vanish at the minimum (for the cosmological constant cancellation), implies that the D-terms should be
$$<D>=\frac{\sqrt{6}}{g}m_{3/2}M_P.$$
(31)
A more subtler constraint comes from the existence of a minimum, eq.(3). In this limit, it takes the form $`g_A^2D_A(_\beta D_A)=0`$. It is clear that for a single $`U(1)`$ gauge group, this would mean at the minimum either the vev to vanish or the D-term to vanish. Both these conditions are not acceptable to us. The situation would not change even if one adds more flavon fields. Thus we rule out the case of single $`U(1)`$ with pure D-breaking. The minimum case we can think of is that case with two $`U(1)`$ gauge groups with two charged fields.
We parameterise the SUSY breaking D-terms, consistently with the vanishing of the cosmological constant, as
$$<D_A>=\frac{\sqrt{6}}{g_A}\theta _Am_{3/2}M_P,$$
(32)
where $`\theta _A`$ are defined such that $`_A\theta _A^2=1`$. Then the soft terms reduce to the following form :
$`m_{i\overline{j}}^2`$ $`=`$ $`2m_{3/2}^2G_{i\overline{j}}+\sqrt{6}m_{3/2}M_P{\displaystyle \underset{A}{}}g_A\theta _A(X_i^A+v_{\overline{l}}X_{\overline{l}}^A_{\overline{l}}+\eta _A^{\overline{\alpha }}_{\overline{\alpha }})G_{i\overline{j}},`$
$`m_{ij}^2`$ $`=`$ $`m_{3/2}^2\left(_iG_j+{\displaystyle \frac{3}{2}}{\displaystyle \underset{A}{}}g_A^2\theta _A^2_i_jf_A\right),`$
$`\mu _{ij}`$ $`=`$ $`m_{3/2}_iG_j,`$
$`A_{ijk}^{}`$ $`=`$ $`m_{3/2}\lambda _{ijk}(\gamma _i+\gamma _j+\gamma _k),`$
$`M_{1/2}^{}_{}{}^{}A`$ $`=`$ $`{\displaystyle \frac{g_A^2b_0^A}{16\pi ^2}}m_{3/2}.`$ (33)
The gaugino masses vanish at the tree level in this limit. They are generated by anomaly mediated contributions as listed above. Similar thing happens for the A-parameters, which are determined by their anomalous dimensions ($`\gamma _i`$) as given above. Note that the above mass formulae are given at the high scale. One has to evolve these masses at the weak scale to make contact with weak scale phenomenology. The present scenario describes a new situation where the non-holomorphic scalar soft masses are given by dominant D-type supersymmetry breaking terms, whereas the gauginos, described by the beta-functions, the supersymmetric fermion masses (in particular the $`\mu `$ term of MSSM) are proportional to the gravitino mass and have therefore much lower values. If all the $`U(1)`$ groups are in the visible sector with large D-terms and positive charges, such a situation is not phenomenologically viable, since there is no possibility of tuning one Higgs doublet to be very light. However, if some of the $`U(1)`$ lie in the hidden sector with some others in the visible sector and the angles $`\theta _A`$ in the visible sector are all small, then the scenario with pure D-breaking becomes viable. In this last case, all the soft terms can be at the TeV scale, thus making contact with a low energy physics of the MSSM type. It would be interesting to see how this new structure of soft terms would feature with respect to low-energy constraints like electroweak symmetry breaking, dark matter, LEP Higgs bounds and other constraints. Note that a situation like split SUSY could be difficult to incorporate here.
### III.2 D-breaking with dilaton and moduli supersymmetry breaking
The above discussion presents an extreme situation i.e. completely absent F-type breaking. However such an extreme limit is not required to realise split supersymmetry breaking. The general analysis presented in the previous section shows that it is enough to have $`g_A^2D_A>>m_{3/2}^2`$. We present here soft terms for a case where, for simplicity, there is only one $`U(1)`$ large D-term and we assume that the auxiliary field of the dilaton or the overall modulus superfields also contribute to supersymmetry breaking.
Note that such a situation can arise naturally when one considers effective lagrangians of Type I string theory for an orientifold with only $`D9`$ branes. We provide expressions for the case of orbifold theories (Calabi-Yau spaces are also particular cases of the expressions below) in the large volume limit. In this limit, the gauge kinetic function and the Kähler potential $`K`$ will have the general form quevedoreview :
$`f_A^B`$ $`=`$ $`S\delta _A^B+\mathrm{},`$
$`K`$ $`=`$ $`\mathrm{log}(S+S^{})3\mathrm{log}(T+T^{}\delta _{GS}V)+{\displaystyle \underset{i}{}}(T+T^{}\delta _{GS}V)^{n_i}|\varphi _i|^2`$ (34)
$`+`$ $`{\displaystyle \underset{ijk}{}}(Z_{ijk}(T+T^{}\delta _{GS}V)^{n_i}\overline{\varphi }_i\varphi _j\varphi _k+\mathrm{h}.\mathrm{c}.)+\mathrm{},`$
where we have used the by now standard notation with $`S`$ representing the dilaton field, $`T`$ representing the overall volume modulus, $`\varphi _i`$ represent matter fields and $`n_i`$ modular weights of matter fields. We are assuming from now on that the modulus $`T`$ is the one mixing with the anomalous $`U(1)`$ gauge field, such that the gauge invariant combination $`T+T^{}\delta _{GS}V`$ should consistently appear in the Kähler potential and in the couplings to the matter fields. This can be explicitly realized in intersecting brane models, as we will illustrate later on. The last term in the Kähler potential in (34) accommodate the possibility of $`\mu `$ terms and simultaneously, that of the $`B_\mu `$ term. We parameterize the supersymmetry breaking contributions from the two sets of auxiliary fields as :
$$<G_S>=\sqrt{3}(\frac{M_P}{S+S^{}})\mathrm{cos}\gamma _S,<D>=\frac{\sqrt{6}}{g}m_{3/2}M_P\mathrm{sin}\gamma _S.$$
(35)
We then obtain the soft terms
$`m_{i\overline{j}}^2`$ $`=`$ $`(13\mathrm{sin}^2\gamma _S)m_{3/2}^2G_{i\overline{j}}+\sqrt{6}gm_{3/2}M_P\mathrm{sin}\gamma _S(X_i+\overline{v}_{\overline{l}}X_{\overline{l}}_{\overline{l}}+\delta _{GS}_{\overline{T}})G_{i\overline{j}}`$
$`m_{ij}^2`$ $`=`$ $`(23\mathrm{sin}^2\gamma _S)m_{3/2}^2_iG_j{\displaystyle \frac{3}{2}}m_{3/2}^2g^2\mathrm{sin}^2\gamma _S_i_jf,`$
$`A_{ijk}`$ $`=`$ $`3m_{3/2}^2\mathrm{cos}^2\gamma _S_i_jG_k,`$
$`M_{1/2}^A`$ $`=`$ $`{\displaystyle \frac{\sqrt{3}}{2}}m_{3/2}\mathrm{cos}\gamma _S,`$ (36)
whereas the $`\mu `$ term is unchanged (33). In the complementary case where the only F-type source of supersymmetry breaking comes from the $`T`$ field, the appropriate parametrization is
$$<G_T>=3(\frac{M_P}{T+T^{}})\mathrm{cos}\gamma _T,<D>=\frac{\sqrt{6}}{g}m_{3/2}M_P\mathrm{sin}\gamma _T.$$
(37)
The soft terms in this case are given by
$`m_{i\overline{j}}^2`$ $`=`$ $`(1+n_i\mathrm{cos}^2\gamma _T3\mathrm{sin}^2\gamma _T)m_{3/2}^2G_{i\overline{j}}+\sqrt{6}gm_{3/2}M_P\mathrm{sin}\gamma _T(X_i+\overline{v}_{\overline{l}}X_{\overline{l}}_{\overline{l}}+\delta _{GS}_{\overline{T}})G_{i\overline{j}},`$
$`m_{ij}^2`$ $`=`$ $`[2+(n_i+n_j)\mathrm{cos}\gamma _T3\mathrm{sin}^2\gamma _T]m_{3/2}^2_iG_j{\displaystyle \frac{3}{2}}m_{3/2}^2g^2\mathrm{sin}^2\gamma _T_i_jf,`$
$`A_{ijk}`$ $`=`$ $`m_{3/2}^2[3\mathrm{cos}^2\gamma _T+(n_i+n_j+n_k)\mathrm{cos}\gamma _T]_i_jG_k,`$
$`M_{1/2}^{}_{}{}^{}A`$ $`=`$ $`{\displaystyle \frac{g_A^2}{8\pi ^2}}\left(3T_G^A\mathrm{sin}^2{\displaystyle \frac{\gamma _T}{2}}T_R^A(\mathrm{sin}^2{\displaystyle \frac{\gamma _T}{2}}(1+n_i)\mathrm{cos}\gamma _T)\right)m_{3/2}.`$ (38)
Several simplifying assumptions were used in deriving (38). For reasons already explained, the analytic scalar masses come from a Giudice-Masiero term in the Kähler potential of the type $`\varphi ^{}Q_iQ_j+\mathrm{h}.\mathrm{c}.`$, where $`\varphi `$ is a flavon type field with a large vev. The Yukawa couplings were assumed, in the large volume limit, to become T-modulus independent, otherwise new contributions appear in the trilinear A-terms. The natural values of modular weights for charged D9 branes charged fields are $`n_i=1`$. Finally for phenomenological studies, the angles $`\gamma _{S,T}`$ can be used as independent parameters to be constrained by low energy physics.
## IV Split Supersymmetry
The requirement of split supersymmetry type soft spectra are as follows :
(i) Scalar soft terms : $`m_{\stackrel{~}{f}}^2𝒪(10^610^{15})`$ GeV, $`(\stackrel{~}{f}=Q,u^c,d^c,L,e^c)`$
(ii). Higgs mass parameters $`m_{H_1}^2m_{H_2}^2B_\mu 𝒪(10^610^{15})`$ GeV, with one of the Higgs mass eigenvalues fine tuned to be around the electroweak scale.
(iii). The gaugino masses and the $`\mu `$ term are around the weak scale.
As a starting point, let us consider for a moment that all D-term contributions are negligible or zero. In such a case, we see that most likely the mass squared terms are proportional to $`m_{3/2}^2`$ whereas the gaugino masses are proportional to $`m_{3/2}`$. Thus, it is difficult to expect a large splitting within the masses of the superpartners in supergravity theories with pure or dominant F-type SUSY breaking. In principle, such a splitting can be arranged by choosing suitable parameter space within the goldstino directions in certain classes of effective lagrangians coming from heterotic strings. However, it is not clear how much these parameter spaces would remain stable under radiative corrections. Another approach for creating a split would be to assume some R-symmetries<sup>12</sup><sup>12</sup>12Or even a charge symmetry accompanied by F-breaking of charged chiral superfield, wells . protecting the fermion superpartners. In this case, the gravitino mass needs to be pushed to very high values, whereas the gauginos need another mechanism to achieve masses close to the weak scalesplit ; ad . However in this case, one has to invent a mechanism to suppress the anomaly mediated contributions, which could involve for example no-scale type models.
In the presence of D-terms, it is generically difficult to realise split supersymmetry like models <sup>13</sup><sup>13</sup>13See also kn .. From the discussion in the previous section, it was obvious that it is just not sufficient to choose the $`U(1)`$ charges of the scalars to be positive to realise the split spectrum since $`B_\mu `$ term does not have large D-term contributions, we need to disentangle the $`\mu `$ and the $`B_\mu `$ term by introducing a new field $`X`$ and allowing a term of the type $`XH_1H_2`$ in the superpotential. In a simple example, the field content is as follows. The model contains an additional $`U(1)`$ group, with two additional fields $`X`$ and $`\varphi `$ with charges $`+2`$ and $`1`$. The $`\varphi `$ field can act as a flavon field attaining a large vev close to the fundamental scale. The superpotential and the relevant term in the Kähler potential are specified as
$`W=W_{SSM}+\lambda _1XH_1H_2+\lambda _2X\varphi ^2+\mathrm{},`$
$`K{\displaystyle \underset{i}{}}|\varphi _i|^2+(\varphi ^{})^2H_1H_2+\mathrm{}.`$ (39)
The scalar potential at the global SUSY level is given by
$$V=\lambda _2^2(|\varphi |^4+4|X|^2|\varphi |^2)+\frac{1}{2}g^2(2|X|^2|\varphi |^2+\xi )^2+\mathrm{}.$$
(40)
For $`\xi >0`$, the stable extremum of the above potential and the auxiliary fields are given by:
$`\varphi ={\displaystyle \frac{g^2}{2\lambda _2^2+g^2}}\xi ,X=0,`$
$`F_\varphi =0,F_X={\displaystyle \frac{\lambda _2g^2}{2\lambda _2^2+g^2}}\xi ,D={\displaystyle \frac{2\lambda _2^2}{2\lambda _2^2+g^2}}\xi .`$ (41)
From the above it is clear that $`F_Xg^2D`$ and moreover of the order of the FI term $`\xi `$. This is sufficient to enable the $`B`$ term to receive large contributions through the term $`G^X_{H_1}_{H_2}G_X`$ in the eq.(22). As long as $`\xi `$ is close to an intermediate scale value, this model seems to replicate the split spectrum, if one fixes the gravitino mass around 1 TeV. However, in typical string models, the FI term is of the $`𝒪(M_{Pl}^2/16\pi ^2)`$ which would give a too large contribution to the vacuum energy. One way to get the correct order of magnitude is by incorporating the above model into a higher dimensional theory. For illustration lets us consider a 5D theory compactified over $`S^1/Z_2`$. The Standard Model and the $`X,\varphi `$ fields live on a 3D brane, whereas the gauge fields of the $`U(1)`$ are allowed to propagate in the bulk. We will use Scherk-Schwarz mechanism to break supersymmetry. The R-symmetry is also broken by this mechanism giving rise to the gravitino mass.
The various scales in the problem are $`R=tM_5^1,RM_5^3=M_P^2`$, where $`tReT`$, the modulus field. After canonically normalizing the various fields by $`\widehat{\varphi }_i=\sqrt{t/3}\varphi _i`$ and at the global supersymmetry level, the potential retains the form (40) with $`\xi M_5^2=M_P^2/t`$. The four dimensional $`U(1)`$ gauge coupling is given by $`g^2=1/t=1/(RM_5)`$, whereas the gravitino mass is given by $`m_{3/2}=\omega /R`$, where $`\omega `$ is a number of order one. The D-term contribution to the vacuum energy is then of the form
$$V_Dg^2M_5^4m_{3/2}^2M_P^2,$$
(42)
in the right order as required by the cancellation of the vacuum energy in supergravity and realisation of the split spectrum. If the no-scale structure is broken by the dynamics, the gauginos attain their masses through anomaly mediation and thus we choose the gravitino mass to be of the order of 100 TeV. The $`\mu `$ is generated by the Giudice-Masiero mechanism and is $`\mu (<\varphi >/M_5)^2m_{3/2}`$. So, this model replicates the spectrum of the split supersymmetry at the weak scale using large D-terms of the intermediate scale and a 100 TeV massive gravitino.
In the light of above discussion, an important question is in which sense the light Higgs mass tuning is preferred over the tuning of another scalar mass. Tuning of squarks or slepton masses is best described in terms of alignment in the $`3\times 3`$ flavor space. If sfermion mass matrices are very close to the diagonal, i.e. off-diagonal terms are very small compared to the diagonal ones, the tuning of a small mass eigenvalue is impossible, whereas the tuning becomes more and more likely for off-diagonal terms of the same order as the diagonal ones. In flavor models with a low energy supersymmetric spectrum, the alignment of the quark-squark and lepton-slepton mass matrices was necessary to avoid too large FCNC effects, but a serious tension between alignment and hierarchy of fermion masses was present, at least for models with only one $`U(1)`$ factor. It is ironical that, in the limit of evading FCNC effects by decoupling the undesirable scalar particles, the alignment has still to be invoked in order to minimize the likelihood of the fine-tuning of squark and slepton masses compared with the tuning of the light Higgs mass.
## V Nonperturbative moduli stabilisation and large D-terms
In string theory, the FI terms are field (moduli) dependent. If no additional dynamics is present, the moduli fields will always exhibit a runaway behaviour and the FI terms disappear. We revisit here the issue of moduli stabilisation with realisation of large D-term contributions in a context similar to, but having some new features compared to the one discussed some time ago in bd . As will become transparent, our analysis is also relevant for the issue of the uplift of the energy density in the context of KKLT type moduli stabilisation kklt ; bkq . The gauge group consists of the Standard Model supplemented by a confining hidden sector group and an anomalous $`U(1)_X`$. We consider the case of a supersymmetric $`SU(N_c)`$ gauge group with $`N_f`$ quark flavors $`Q_i^a`$ and anti-quark $`\stackrel{~}{Q}_{\overline{i}}^a`$ where $`a=1\mathrm{}N_c`$ is an index in the fundamental representation of the $`SU(N_c)`$ gauge group and $`i,\overline{i}=1\mathrm{}N_f`$ are flavor indices. In the intersecting string realisation, discussed in some detail in the next section, the hidden sector consists of a stack of $`N_c`$ magnetised D9 branes in the type I string with kinetic function $`f=S+kT`$, where S is the dilaton (super)field, T a volume (Kähler) modulus and $`k`$ is a positive or negative integer determined by the magnetic fluxes in two compact torii. The low energy dynamics is described by $`M_{\overline{j}}^i=Q^{a,i}\stackrel{~}{Q}_{\overline{j}}^a`$, the composite ”mesons” fields. In the following we denote by $`q`$ $`(\overline{q})`$ the $`U(1)_X`$ charges of the hidden sector quarks (antiquarks). Since the FI terms are T-modulus dependent, T will shift under gauge transformations
$`V_XV_X+\mathrm{\Lambda }_X+\overline{\mathrm{\Lambda }}_X,M_{\overline{j}}^ie^{2(q+\overline{q})\mathrm{\Lambda }_X}M_{\overline{j}}^i,`$
$`TT+\delta _{GS}\mathrm{\Lambda }_X,`$ (43)
where
$$\delta _{GS}=\frac{C_{N_c}}{k},C_{N_c}=\frac{1}{4\pi ^2}N_f(q+\overline{q}),$$
(44)
is uniquely fixed by the requirement that the mixed $`U(1)_XSU(N_c)^2`$ anomaly, denoted $`C_{N_c}`$ in (44), to be exactly canceled by the nonlinear transformation of $`ImT`$. Notice that the nonlinear transformation of $`T`$ forces a chiral nature of the hidden sector with respect to the anomalous abelian gauge group, which in turn triggers supersymmetry breaking bd . In order to be able to write gauge invariant mass terms for the mesons, a field with charge opposite in sign to the ones of the mesons has to be introduced, called $`\varphi `$ in what follows, of charge $`1`$ in our conventions. The dynamical scale of the hidden sector gauge group, the effective superpotential ads and the Kähler potential are
$`\mathrm{\Lambda }=M_Pe^{8\pi ^2(S+kT)/(3N_cN_f)},`$
$`W=W_0(S)+(N_cN_f)\left({\displaystyle \frac{\mathrm{\Lambda }^{3N_cN_f}}{detM}}\right)^{\frac{1}{N_cN_f}}+m_i^{\overline{j}}({\displaystyle \frac{\varphi }{M_P}})^{(q+\overline{q})}M_{\overline{j}}^i,`$
$`K=\mathrm{ln}(S+\overline{S})3\mathrm{ln}[T+\overline{T}2Tr(M^{}M)^{1/2}|\varphi |^2\delta _{GS}V],`$ (45)
where $`m_i^{\overline{j}}`$ are mass parameters. Notice first of all that the dynamical superpotential
$$W_{np}=(N_cN_f)\left(\frac{e^{8\pi ^2(S+kT)}}{detM}\right)^{\frac{1}{N_cN_f}},$$
(46)
is precisely gauge invariant when the anomaly cancellation conditions (43)-(44) are satisfied. In order to stabilise the modulus S we invoke the three-form NS-NS and RR fluxes. $`W_0`$ depends on the modulus $`S`$, $`S=S_0`$ and eventually other (complex structure) moduli of the theory and stabilises them by giving them a very large mass. If the other relevant mass scales, the FI term and the dynamical scale $`\mathrm{\Lambda }`$ have much lower values, we can safely integrate out these fields, by keeping the T modulus in the low energy dynamics. The resulting lagrangian is similar to the one invoked in the KKLT moduli stabilisation kklt with a D-term uplifting of the vacuum energy bkq . Notice however that the simple nonperturbative superpotential $`e^{aT}`$ considered in bkq cannot be gauge invariant due to the gauge transformation of $`T`$ and therefore, precisely as in the heterotic case discussed in bd , charged hidden sector matter with appropriate charges is crucial to define a consistent gauge invariant model.
Minimisation with respect to $`T`$ in (45) stabilises also the Kähler modulus. For notational simplicity we discuss in some detail the case of an supersymmetric hidden sector $`SU(2)`$ gauge group with one quark flavor $`Q^a`$ and anti-quark $`\stackrel{~}{Q}^a`$ where $`a=1,2`$ is an index in the fundamental representation of the gauge group. Due to the anomalous nature of the $`U(1)_X`$, the sum of the quark and antiquark charges, equal to the $`M`$ meson charge, is different from zero and, in our example, equal to $`+1`$. $`\varphi `$ is a field of charge $`1`$ which participate in the Yukawa coupling $`\lambda \varphi Q^a\stackrel{~}{Q}^a`$, which plays the role of meson mass after the spontaneous symmetry breaking of the $`U(1)_X`$. The fact that the meson masses come from a perturbative trilinear Yukawa coupling in this case is instrumental in producing a large D-term contribution to supersymmetry breaking. In order to provide explicitly the scalar potential, we define the canonical field $`M\chi ^2/2`$. Then the supergravity scalar potential can be found to be
$`V_F`$ $`=`$ $`{\displaystyle \frac{1}{r^3}}\left\{{\displaystyle \frac{r^2}{3}}|_TW{\displaystyle \frac{3}{r}}W|^2+{\displaystyle \frac{r}{3}}{\displaystyle \underset{i=1}{\overset{2}{}}}|_iW+\overline{\varphi _i}_TW|^23|W|^2\right\},`$
$`V_D`$ $`=`$ $`{\displaystyle \frac{1}{S+\overline{S}+k^{}(T+\overline{T})}}\left({\displaystyle \frac{3}{r}}X_i|\varphi _i|^2+3{\displaystyle \frac{\delta _{GS}M_P^2}{T+\overline{T}}}\right)^2,`$ (47)
where $`\varphi _i=\chi ,\varphi `$, we have introduced $`r(T+\overline{T}_i|\varphi _i|^2)`$, $`\delta _{GS}`$ represent the Green-Schwarz coefficient of the $`U(1)_X`$, and $`k^{}`$ is the magnetic flux on the brane providing the anomalous $`U(1)_X`$.
By inserting (45) into (47) we find a model with all moduli stabilised. If we would ignore the $`U(1)_X`$ dynamics, for example, $`T`$ would be stabilised as in kklt by solving $`D_TW=0`$. In our case, the minimum $`T_0=T`$ will be shifted due to the D and new F contributions. A full supergravity analysis of the vacuum of (47) is possible but cumbersome. Due to this after stabilizing $`S=S_0`$ and $`T=T_0`$ by solving their equations of motion, we analyse the stabilisation of the other fields, for simplicity at the global supersymmetry level, as in bd , by a suitable rescaling of the fields and Yukawa coupling $`\lambda `$. For general $`N_c`$ , $`N_f`$ at the global level, the auxiliary fields and the scalar potential are
$`(F^{\overline{M}})_i^{\overline{i}}=2[(M^{}M)^{1/2}]_{\overline{j}}^{\overline{i}}\left[(M^1)_i^{\overline{j}}\left({\displaystyle \frac{\mathrm{\Lambda }^{3N_cN_f}}{detM}}\right)^{\frac{1}{N_cN_f}}+m_i^{\overline{j}}({\displaystyle \frac{\varphi }{M_P}})^{(q+\overline{q})}\right],`$
$`\overline{F}_{\overline{\varphi }}={\displaystyle \frac{q+\overline{q}}{M_P}}({\displaystyle \frac{\varphi }{M_P}})^{q+\overline{q}1}Tr(mM),`$
$`D_X=(q+\overline{q})Tr(M^{}M)^{1/2}|\varphi |^2+k\mu ^2,`$
$`V=|F_\varphi |^2+{\displaystyle \frac{1}{2}}[(M^{}M)^{1/2}]_{\overline{i}}^{\overline{j}}(F^{\overline{M}})_i^{\overline{i}}(F^M)_{\overline{j}}^i+{\displaystyle \frac{g_X^2}{2}}D_X^2,`$ (48)
where $`\mu ^2=3C_{N_c}/k^2(T+\overline{T})`$ is a mass scale determined by the T-modulus vev. The new feature of (48) is that $`k`$ and consequently the FI term can have both signs, whereas in the effective heterotic string framework worked out in bd , the FI term had only one possible sign.
In the limit $`\mathrm{\Lambda }<<\mu `$, the vacuum structure and the pattern of supersymmetry breaking in the two cases of $`k`$ positive and negative are vastly different.
i) $`k>0`$. In this case the vacuum can be determined as in bd , where it was analysed for arbitrary $`N_f<N_c`$ and arbitrary $`q+\overline{q}>0`$ charges. Keeping one mass parameter $`m_i^{\overline{j}}=m\delta _i^{\overline{j}}`$, we find, to the lowest orders in the parameter $`ϵ`$ defined by
$$ϵ\frac{M_0}{k\mu ^2}=(\frac{\mathrm{\Lambda }}{\sqrt{k}\mu })^{\frac{3N_cN_f}{N_c}}\left[\frac{m}{M_P}(\frac{\sqrt{k}\mu }{M_P})^{q+\overline{q}1}\right]^{\frac{N_fN_c}{N_c}},$$
(49)
a hierarchically small scale of supersymmetry breaking
$`|\varphi |^2=k\mu ^2\left[1+ϵN_f(q+\overline{q})\right],M=M_0\left[1ϵ(q+\overline{q})^2{\displaystyle \frac{N_f(N_cN_f)(2N_cN_f)}{2N_c^2}}\right],`$
$`g_X^2D_X=ϵ^2\widehat{m}^2N_f^2(q+\overline{q})^2\left[1{\displaystyle \frac{N_f}{N_c}}(q+\overline{q})\right],`$
$`F_\varphi =ϵ\widehat{m}\sqrt{k}\mu N_f(q+\overline{q}),F^{\overline{M}}=K^{M\overline{M}}_MW=ϵ^2\widehat{m}k\mu ^2{\displaystyle \frac{N_f(N_cN_f)}{N_c}}(q+\overline{q})^2,`$ (50)
where $`\widehat{m}m(\sqrt{k}\mu /M_P)^{q+\overline{q}}`$.
ii) $`k<0`$. Here we specifically consider the case $`N_c=2`$, $`N_f=1`$ and $`q+\overline{q}=1`$. In this case we find, to the lowest order in the parameter $`ϵ^{}=[(g^2+2\lambda ^2)^4/8\lambda ^2g^{10}](\mathrm{\Lambda }^2/|k|\mu ^2)^5`$ , a large scale of supersymmetry breaking
$`\varphi ={\displaystyle \frac{(g^2+2\lambda ^2)^2}{2\lambda g^4}}{\displaystyle \frac{\mathrm{\Lambda }^5}{k^2\mu ^4}}\left[1+3(g^2+14\lambda ^2)ϵ^{}\right],M={\displaystyle \frac{g^2}{g^2+2\lambda ^2}}k\mu ^2\left[12(g^2+14\lambda ^2)ϵ^{}\right],`$
$`D_X{\displaystyle \frac{2\lambda ^2}{g^2+2\lambda ^2}}k\mu ^2,F_\varphi {\displaystyle \frac{\lambda g^2}{g^2+2\lambda ^2}}k\mu ^2,F^{\overline{M}}{\displaystyle \frac{g^2+2\lambda ^2}{g^2}}{\displaystyle \frac{\mathrm{\Lambda }^5}{k\mu ^2M_P^2}}.`$ (51)
Interestingly enough, this second case generate a large scale for supersymmetry breaking with large $`F_\varphi `$ and D-term contributions. At first sight, a breaking of supersymmetry at a scale larger than the dynamical scale $`\mathrm{\Lambda }`$ destroys the supersymmetric confining dynamics underlying the nonperturbative superpotential in (46). However, the breaking of supersymmetry in the hidden sector is described by the mass splitting in the “mesonic” sector, measured by the auxiliary field $`F_M`$. Its value in case ii) is very small and actually the same as in case i), suggesting that the confining dynamics is still essentially supersymmetric. In the case $`q+\overline{q}>1`$ we expect the D-term contribution to have further suppressions since the mesons masses come now from a higher dimensional operator. Within this context, we expect our general analysis of D-term contributions to supersymmetry breaking to be of relevance for further studies of phenomenological models incorporating moduli stabilisation fnop . In the following section we describe string theory realisations based on intersecting brane models leading precisely to the case $`q+\overline{q}=1`$.
## VI Intersecting brane string realisation of large D-term supersymmetry breaking
Even if reasonable from a supergravity point of view, it is not obvious that a large D-term supersymmetry breaking in string theory is possible. Indeed, it is well known from the heterotic string constructions that the presence of Fayet-Iliopoulos terms triggers vev’s for charged fields which break the gauge symmetry rather than supersymmetry dsw . This can presumably be understood by noticing that the FI terms in the heterotic string arise at one-loop and therefore, if they would break supersymmetry, they would be a radiative breaking of supersymmetry which is known to be very hard to obtain witten . It was suggested in bd that at the nonperturbative level, gaugino condensation in the hidden sector in the presence of an anomalous $`U(1)`$ symmetry can break supersymmetry. However as in the previous section for case i) introduced there, the induced D-terms are of the order (or slightly larger) than the $`F^2/M_P`$ type terms and cannot provide the large contributions we are advocating in this paper.
In the Type I or Type II strings, on the other hand, the FI terms appear generically at tree-level and we can expect the tree-level supersymmetry breaking to be possible with large D-terms. As we will see, this will realise case ii) discussed in the previous section. We present here an explicit example suggesting this is indeed possible, in the context of intersecting branes Type I orientifold models or, T-dual equivalently, with internal magnetic fields earlier ; intersecting . We discuss also various ingredients such that supersymmetry breaking to be really possible. We consider an explicit example, even if it is clear that a large class of similar models can be constructed. The model is based on the $`Z_2\times Z_2`$ Type I orbifold without discrete torsion with internal magnetic fields $`H_i^{(a)}=(m_i^{(a)}/v_in_i^{(a)})`$ in the torus $`T^i`$, where $`v_i`$ are the volumes of the three torii. The model contains four stacks of D9 branes, each stack containing four coincident branes. Three of the stacks are magnetised and the fourth one is non-magnetised, with wrapping numbers $`(m_i^{(a)},n_i^{(a)})`$ equal to
$`M_1:(m_i^{(3)},n_i^{(3)})=(0,1),(2,1),(2,1),`$
$`M_2:(m_i^{(1)},n_i^{(1)})=(2,1),(0,1),(2,1),`$
$`M_3:(m_i^{(2)},n_i^{(2)})=(2,1),(2,1),(0,1),`$
$`M_4:(m_i^{(4)},n_i^{(4)})=(0,1),(0,1),(0,1).`$ (52)
The fluxes on $`M_1`$ and $`M_2`$ generate lower dimensional anti-brane like charges whereas the fluxes on $`M_3`$ generate lower dimensional brane like charges. The RR tadpole conditions for the $`Z_2\times Z_2`$ orbifold without discrete torsion with $`(ϵ_1,ϵ_2,ϵ_3)=(1,1,1)`$ are given by
$`{\displaystyle \underset{a}{}}M_an_1^{(a)}n_2^{(a)}n_3^{(a)}=16,{\displaystyle \underset{a}{}}M_an_1^{(a)}m_2^{(a)}m_3^{(a)}=16,`$
$`{\displaystyle \underset{a}{}}M_am_1^{(a)}n_2^{(a)}m_3^{(a)}=16,{\displaystyle \underset{a}{}}M_am_1^{(a)}m_2^{(a)}n_3^{(a)}=16.`$ (53)
The massless spectrum in this class of models is determined by the intersection numbers
$`I_{ab}={\displaystyle \underset{ab}{}}(n_i^{(a)}m_i^{(b)}m_i^{(a)}n_i^{(b)}),I_{ab}={\displaystyle \underset{ab^{}}{}}(n_i^{(a)}m_i^{(b)}+m_i^{(a)}n_i^{(b)}),`$
$`I_{aO}=8(m_i^{(a)}m_i^{(a)}m_i^{(a)}+m_i^{(a)}n_i^{(a)}n_i^{(a)}+n_i^{(a)}m_i^{(a)}n_i^{(a)}n_i^{(a)}n_i^{(a)}m_i^{(a)}).`$ (54)
The contribution of the four stacks of branes to the RR tadpole conditions with wrapping numbers (52) precisely satisfy (53) when $`M_1=M_2=M_3=M_4=4`$. The gauge group of this model is $`U(2)^3SO(4)`$. The model was chosen such that the chiral massless spectrum, determined by the intersection numbers, to contain only strings stretched between different stacks of branes. More precisely, by defining $`M_i=2p_i`$, it is given by
$`\varphi _{1,\overline{a}b}^i:16\times (\overline{𝐩}_\mathrm{𝟐},𝐩_\mathrm{𝟑}),\varphi _{2,ac}^j:16\times (𝐩_\mathrm{𝟏},𝐩_\mathrm{𝟐}),`$
$`\varphi _{3,\overline{b}\overline{c}}^k:16\times (\overline{𝐩}_\mathrm{𝟏},\overline{𝐩}_\mathrm{𝟑}),`$ (55)
where the multiplicity of $`16`$ in each sector comes from the intersection numbers of various branes, whereas all other charged states are non-chiral and will get a mass. There are mixed $`U(1)_aU(1)_b^2`$ and $`U(1)_aSU(p_b)^2`$ gauge anomalies in the model, easily computable from the massless spectrum.
$$C_{ab}=\frac{1}{4\pi ^2}Tr(X_aX_b^2)=\frac{2^4}{4\pi ^2}Tr(X_aT_b^2)=\frac{2^8}{4\pi ^2}\left(\begin{array}{ccc}0& 1& 1\\ 1& 0& 1\\ 1& 1& 0\end{array}\right),$$
(56)
where $`X_a`$, $`a=1,2,3`$ are the $`U(1)_a`$ gauge factors, whereas $`T_a`$ are the nonabelian generators. They are taken care by axionic couplings of the type $`\mathrm{\Theta }_aF^aF^a`$, where $`\mathrm{\Theta }_a=Imf_a`$, where the gauge kinetic functions are given by
$$f_a=\underset{i}{}n_i^{(a)}Sn_1^{(a)}m_2^{(a)}m_3^{(a)}T_1m_1^{(a)}n_2^{(a)}m_3^{(a)}T_2m_1^{(a)}m_2^{(a)}n_3^{(a)}T_3$$
(57)
and where $`ImS,ImT_i`$ are the axion-dilaton and the three axions associated to the three internal tori. In our concrete example above, the gauge kinetic functions are explicitly $`f_1=S4T_1`$ , $`f_2=S4T_2`$, $`f_3=S+4T_3`$. The mixed gauge anomalies are taken care by the nonlinear gauge transformations
$`\delta ImT_1={\displaystyle \frac{16}{\pi ^2}}(\alpha _2+\alpha _3),\delta ImT_2={\displaystyle \frac{16}{\pi ^2}}(\alpha _1+\alpha _3)`$
$`\delta ImT_3={\displaystyle \frac{16}{\pi ^2}}(\alpha _1+\alpha _2),`$ (58)
where $`\alpha _a`$ are the gauge transformation parameters for the $`U(1)_a`$ factors. The Kähler potential contains the terms
$$\mathrm{ln}[T_1+\overline{T}_1+\frac{16}{\pi ^2}(V_2V_3)]\mathrm{ln}[T_2+\overline{T}_2+\frac{16}{\pi ^2}(V_1+V_3)]\mathrm{ln}(T_3+\overline{T}_3+\frac{16}{\pi ^2}(V_1+V_2)]$$
(59)
which generate FI terms in the effective field theory
$$\xi _1=\frac{16}{\pi ^2}(\frac{1}{t_2}+\frac{1}{t_3}),\xi _2=\frac{16}{\pi ^2}(\frac{1}{t_1}+\frac{1}{t_3}),\xi _3=\frac{16}{\pi ^2}(\frac{1}{t_1}\frac{1}{t_2}).$$
(60)
The Fayet-Iliopoulos terms can be more generally be written in terms of magnetic fluxes as
$$\xi _aH_1^{(a)}+H_2^{(a)}+H_3^{(a)}H_1^{(a)}H_2^{(a)}H_3^{(a)}$$
(61)
and in this model they satisfy the sum rule $`\xi _1\xi _2\xi _3=0`$. The D-terms on the U(1) factors of each $`U(2)=U(1)SU(2)`$ stack are given by
$`D_1={\displaystyle \underset{i,\overline{a},b}{}}|\varphi _{2,\overline{a}b}^i|^2{\displaystyle \underset{j,a,c}{}}|\varphi _{3,ac}^j|^2+\xi _1,`$
$`D_2={\displaystyle \underset{i,\overline{a},b}{}}|\varphi _{1,\overline{a}b}^i|^2+{\displaystyle \underset{k,\overline{b},\overline{c}}{}}|\varphi _{2,\overline{b}\overline{c}}^k|^2+\xi _2,`$
$`D_3={\displaystyle \underset{j,a,c}{}}|\varphi _{1,ac}^j|^2{\displaystyle \underset{k,\overline{b},\overline{c}}{}}|\varphi _{3,\overline{b}\overline{c}}^k|^2+\xi _3,`$ (62)
and satisfy also the same rule
$$D_1D_2D_3=\xi _1\xi _2\xi _3=0.$$
(63)
We believe that the interpretation of the sum rule (63) is that the D-branes tend to recombine by condensing the by-fundamental fields (55) and to provide a supersymmetric vacuum $`D_a=0`$. Notice however that the model has renormalizable superpotential terms
$$W=\lambda _{ijk}Tr(\varphi _1^i\varphi _2^j\varphi _3^k),$$
(64)
(where the trace in the gauge group space), which has a geometrical interpretation of Yukawa couplings connecting 3-fields forming a triangle in each compact torus, analogously to models of Yukawa couplings studied in yukawas . The number of these Yukawas are related, as usual, to the number of D-flat directions. The Yukawa couplings are also field dependent and depend on complex structure moduli. In the presence of the Yukawa couplings (64), the tachyonic instabilities typically related to the brane recombination process can be removed, there is generically a geometrical obstruction and the D-brane recombination is not generically the most favorable process. The FI terms are Kähler moduli dependent and they can (perturbatively ) vanish for particular points in the Kähler moduli space, which will always be dynamically preferred. In order to avoid this phenomenon, nonperturbative effects have to be invoked, for example gaugino condensation on the D9 branes, according to the discussion in the previous section. The perturbative nature of the superpotential terms (64), to be interpreted as meson masses in case ii) of the previous section, generate large D-term contributions . Cosmological constant cancellation for large D-terms in the TeV range gravitino mass ask for intermediate values of FI terms and/or very small $`U(1)`$ gauge couplings. So for large scale of supersymmetry breaking, cosmological constant is generically hard to cancel unless FI terms are much smaller than the Planck scale. Whereas this was not the case in the nonperturbative model of the previous section, the more general formula (61) suggests that it is possible that $`\xi _a<<H_i^{(a)}`$ by tuning the magnetic fluxes, in the spirit of landscape models bp ; split ; ad ; kklt ; ddg .
In order to make connection with the field theory model of the previous section, notice that if dynamics picks up an overall Kahler modulus $`T_1=T_2=T_3=T`$, then $`\xi _3=0`$ and there are two remaining anomalous abelian factors $`U(1)_1`$ , $`U(1)_2`$. In the following we discuss in some more detail a simplified model along these lines.
### VI.1 From intersecting brane models to nonperturbative moduli stabilisation
In order to stabilise all moduli, in section 4 we used nonpeturbative effects on an asymptotically-free gauge group. The explicit string example discussed previously has no asymptotically gauge factor, but we do not expect this result to be generic. In the following, we consider a model, similar to the explicit string example presented in the previous section but containing an asymptotically-free gauge factor. It is also further simplified in order to allow a simple analysis and suited to generate large D-terms at the minimum, comparable and in some regions of the parameter space dominant with respect to the F-terms.
The field content and gauge structure are summarized as follows. The model has a gauge group $`SU(N)U(1)^2`$, with the chiral superfield content
$$\varphi _1^i:(N,1,0),\varphi _{2,\overline{j}}:(\overline{N},0,1),\varphi _3:(1,1,+1),$$
(65)
where $`i,\overline{j}=1\mathrm{}N_f`$ and the notation for charges and representations are transparent in (65). The model is therefore similar to the explicit intersecting brane model of the previous section, but slightly adapted for our purposes. There is a magnetic flux pattern which does lead to the spectrum above, which by itself does not saturate the RR tadpole conditions. This can be cured by adding additional branes or by considering other orbifolds and/or additional antisymmetric field backgrounds. The $`SU(N)`$ plays the role of a hidden sector SYM gauge group with $`N_f`$ flavors, which condenses in the IR. The composite objects
$$M_{\overline{j}}^i=\varphi _1^{ia}\varphi _{2,\overline{j}}^{\overline{a}},$$
(66)
where $`a`$ is an index in the fundamental of $`SU(N)`$, are the mesons used in constructing the effective action of the theory. For simplicity we consider in the following only the overall Kähler modulus $`T`$, whereas keeping all of them would ask for stabilisation a more complicated dynamics, for example several gaugino condensates. Consistently with the cancellation of the mixed gauge anomalies, the gauge kinetic function on the condensing gauge group is $`f_{SU(N)}=S\pm N_fT`$, where the $`+`$ ($``$) signs correspond to a hidden sector with positive (negative) product of magnetic fluxes in two torii. Similar to the explicit intersecting brane model, $`T`$ transforms under gauge transformations as
$$\delta T=\pm \frac{1}{4\pi ^2}(\mathrm{\Lambda }_2\mathrm{\Lambda }_3).$$
(67)
This can also be directly checked by computing the mixed gauge anomalies
$$U(1)_2SU(N)^2:\frac{N_f}{4\pi ^2},U(1)_3SU(N)^2:\frac{N_f}{4\pi ^2},$$
(68)
which are precisely canceled by the nonlinear gauge transformation of the axion $`ImT`$ (67). In order to write the Kähler potential, first of all we place ourselves on the $`SU(N)`$ flat direction $`\varphi _1^{ia}=\varphi _{2,\overline{j}}^{\overline{a}}`$. Similarly to the KKLT proposal, we could first integrate out the dilaton and the complex structure moduli. In doing this, for $`N_f<N`$, we find the effective superpotential and Kähler potential
$`W=W_0+(NN_f)A\left[{\displaystyle \frac{e^{8\pi ^2N_fT}}{detM_{\overline{j}}^i}}\right]^{\frac{1}{NN_f}}+\lambda _{i\overline{j}}\varphi _3M_{\overline{j}}^i,`$
$`K=3\mathrm{ln}\left[T+\overline{T}\pm {\displaystyle \frac{V_3V_2}{4\pi ^2}}2Tr(\overline{M}M)^{1/2}\overline{\varphi }_3\varphi _3\right],`$ (69)
where the constant $`W_0`$ depend on the details of the three-form fluxes, the Kähler potential was computed in the weakly coupled regime of the $`SU(N)`$ flat direction and where $`A=\mathrm{exp}\{8\pi ^2S_0/(NN_f)\}`$.
The reader will notice that this model reassembles closely the model worked out in Section 4. The mass term for the mesons in (69) is actually provided, in the intersecting brane realisation of the previous section, by the Yukawa coupling (64). In analogy with the explicit intersecting brane model, there is a constraint equation
$$D_2D_3=\xi _2\xi _3=0,$$
(70)
where
$$D_2=Tr(\overline{M}M)^{1/2}|\varphi _3|^2\pm \frac{3}{4\pi ^2(T+\overline{T})},$$
(71)
which signals the presence of a flat direction, allowing the presence of the perturbative superpotential providing the meson mass term. Indeed, since there are now two D-flatness conditions and two charged fields, the existence of the flat direction (70) is needed in order to write the meson mass term, the last term in the superpotential (69). The two signs in the expressions above correspond to the two cases $`k>0`$ and $`k<0`$ in Section 4, the second case realising the high scale supersymmetry breaking with large D-terms.
## VII Summary and Outlook
In the present work, we have initiated a program to study in a general manner the implications of the large D-terms in a supergravity on soft supersymmetry breaking parameters. These terms can come from an anomalous $`U(1)`$ flavour model as has been noted in the past or tree level FI terms in an intersecting D-brane model. We have shown that explicit models based on intersecting D-branes can be constructed giving rise to large D-terms. Irrespective of the source, we have studied the implications of these terms on the soft parameters. The mass squared terms are the most affected with contributions linear in D. However the charges of the matter fields under the anomalous $`U(1)`$ can crucially determine the actual impact. The $`B_\mu `$ and $`A`$ terms also receive corrections though they are not significantly modified in terms of magnitude. As an application we have shown that split supersymmetry can be realised with specific choices of the superpotential and Kähler potential.
Particular examples with the large $`D`$ contributions are string models of supersymmetry breaking, in particular in Type I string orientifolds. We have not addressed in detail the phenomenological signatures and constraints on the parameter space within this class of the models. Such a study could be taken in conjunction with a proper flavour model à la Froggatt-Nielsen. This could be then confronted with low energy data from accelerators, dark matter physics constraints and flavour physics.
The issue of moduli stabilisation has been receiving increasing attention in the recent years. Applications to the soft masses have also been recently addressedfnop ; quevedo . Here we have revisited the issue of moduli stabilisation using non-perturbative gaugino condensation in type I orientifolds with internal magnetic fluxes which generate large D-terms. By a suitable choice of the fluxes which fix the sign of the FI term, we have shown that it is possible to stabilize moduli and generate large D-terms.
## Acknowledgments
We wish to thank S. Lavignac, Y. Mambrini, S. Pokorski and Carlos Savoy for useful discussions.
This work is supported in part by the CNRS PICS no. 2530 and 3059, INTAS grant 03-51-6346, the RTN grants MRTN-CT-2004-503369, MRTN-CT-2004-005104 and by a European Union Excellence Grant, MEXT-CT-2003-509661. E.D. thanks the Univ. of Tucson, Arizona for warm hospitality while completing this work. SKV is also supported by Indo-French Centre for Promotion of Advanced Research (CEFIPRA) project No: 2904-2 ‘Brane World Phenomenology’. |
warning/0506/cond-mat0506187.html | ar5iv | text | # Persistent and radiation-induced currents in distorted quantum rings
## I Introduction
It is well known that a static magnetic flux through a mesoscopic ring induces a dissipationless non-decaying (persistent) current at low temperature. During the last 20 years this persistent current has been heavily investigated both from the theoretical and experimental side. ABEandPC ; PCwithSO ; VRWZ98 ; Serega ; Ulloa ; r1 ; r2 ; r3 ; r4 ; r5 ; r6 ; r7 ; r8 ; r9 ; Mohanty99 ; Mohanty01 ; LossMartin ; Szopa ; magarill96 ; bulaev04 ; manolescu ; per05 In particular, theoretical investigations were focused on the effects of electron-electron interaction, disorder, spin-orbit, PCwithSO polarized nuclear spins, VRWZ98 ; Serega and shape. magarill96 ; bulaev04 More recently, the effect of electromagnetic radiation on mesoscopic rings has been investigated. It has been pointed out that the persistent current can be strongly affected by radiation manolescu and also that radiation can induce a current at zero magnetic flux. per05 ; matos05 In order to break the clockwise-anticlockwise symmetry and obtain a current at zero flux a radiation with some degree of circular polarization is needed. This can be achieved by a superposition of pairs of time-asymmetric, linearly cross-polarized picosecond pulses, matos05 or more simply by using circularly polarized radiation. per05 All these effects were mainly studied in systems with perfect ring geometry, i.e. rings where the confinement potential does not depend on the azimuthal coordinate. The only exceptions, to our knowledge, are Refs. magarill96, and bulaev04, , where the persistent current in an elliptical quantum ring magarill96 and the persistent current in a quantum ring on a surface of constant negative curvature bulaev04 were considered.
In this paper we investigate persistent and radiation-induced currents in distorted quantum rings. We consider a narrow distorted ring of uniform cross section lying on a plane (Fig. 1). The ring can be composed by several segments of different curvature or can have any smooth curvature profile. We show that the curvature of the ring enters into the Scrödinger equation via a geometrical potential term of the form $`V_{geom}=\mathrm{}^2/(8m^{}R^2)`$, where $`R`$ is the radius of curvature. Our model differs from the one used in Ref. magarill96, , where a quantum ring of a non-uniform cross section was considered. By choosing a uniform cross section the effect of the distorsion can be described in a simpler and more transparent way. Our paper is organized as follows. In Sect. II we derive the Schrödinger equation for one electron in a distorted ring in the presence of a magnetic flux. We show that, as in the ideal case, the wavefunction in the distorted ring is a periodic function of the magnetic flux $`\mathrm{\Phi }`$ with period $`\mathrm{\Phi }_0`$. Next, in Sect. III, we consider a model where a distorted quantum ring consists of four constant-curvature segments. We find the energy spectrum of such a ring and we demonstrate that the geometrical potential $`V_{geom}`$ in this case opens gaps in the electron energy spectrum. Moreover, we show that the geometrical potential can lead to bound states. The oscillation of the persistent current, and the frequency and intensity dependence of the radiation-induced currents in the distorted ring are studied in Sec. IV. The results of our investigations are summarized in Sec. V.
## II Effective Hamiltonian
We consider one electron with effective mass $`m^{}`$ confined by a potential $`V_\gamma `$ ($`\gamma `$ denotes the characteristic width of $`V_\gamma `$) to a closed curve $`C`$ on a plane. A uniform magnetic field $`H`$ perpendicular to the plane is applied. The Schrödinger equation has the form
$$\frac{1}{2m^{}}\left(\widehat{𝐩}\frac{e}{c}𝐀\right)^2\psi +V_\gamma \psi =E\psi ,$$
(1)
where $`\widehat{𝐩}`$ is the electron momentum operator and $`𝐀(𝐫)=\frac{1}{2}[𝐇,𝐫]`$ is the vector potential. Using the property $`\text{div}(𝐀)=0`$ of this gauge, Eq. (1) can be rewritten as
$$\frac{1}{2m^{}}\left(\mathrm{}^2\mathrm{\Delta }2\frac{e}{c}𝐀\widehat{𝐩}+\frac{e^2}{c^2}𝐀^2\right)\psi +V_\gamma \psi =E\psi .$$
(2)
Our goal is to obtain an effective one-dimensional Schrödinger equation in the limit of a strong transverse confinement, i.e. in the limit $`\gamma 0`$. We follow the approach proposed in Ref. costa81, and subsequently used in Ref. shevchenko01, .
Let us introduce the orthonormal coordinate system ($`s`$, $`q`$), where $`s`$ is the arc length parameter and $`q`$ is the coordinate along the normal $`𝐧(s)`$. The curve $`C`$ is then described by a vector $`𝐫(s)`$ as a function of the arc length $`s`$. In a vicinity of $`C`$ the position is therefore is described by
$$𝐑(s,q)=𝐫(s)+q𝐧(s).$$
(3)
For the sake of simplicity shevchenko01 we assume that $`V_\gamma `$ depends only on the $`q`$ coordinate describing the displacement from the reference curve $`C`$ only.
The Laplacian $`\mathrm{\Delta }`$ in the curvilinear coordinates $`s`$ and $`q`$ is given by
$$\mathrm{\Delta }_{s,q}=\frac{1}{h}\frac{}{s}\frac{1}{h}\frac{}{s}+\frac{1}{h}\frac{}{q}h\frac{}{q},$$
(4)
with
$$h=1k(s)q,$$
(5)
where $`k(s)=R^1(s)`$ is the curvature. Using the transformation to the new wave function $`\chi (s,q)`$ via $`\psi (s,q)=\chi (s,q)/\sqrt{h}`$ (note, that $`\chi (s,q)`$ is properly normalized), we can rewrite Eq. (2) as
$`{\displaystyle \frac{1}{2m^{}}}[\mathrm{}^2({\displaystyle \frac{}{s}}{\displaystyle \frac{1}{h^2}}{\displaystyle \frac{}{s}}{\displaystyle \frac{h_{ss}}{2h^3}}+{\displaystyle \frac{5h_s^2}{4h^4}}+{\displaystyle \frac{^2}{q^2}}+{\displaystyle \frac{k^2}{4h^2}})+`$
$`2i{\displaystyle \frac{e\mathrm{}}{c}}(A_s(s,q)({\displaystyle \frac{}{s}}+{\displaystyle \frac{k_sq}{2h}})+A_q(s,q)({\displaystyle \frac{}{q}}+{\displaystyle \frac{k}{2h}}))+{\displaystyle \frac{e^2}{c^2}}𝐀^2]\chi +V_\gamma \chi =E\chi ,`$ (6)
where $`h_s=h/s`$, $`h_{ss}=^2h/s^2`$, $`k_s=k/s`$, $`A_q(s,q)`$ and $`A_s(s,q)`$ are the components of $`𝐀`$ along the $`q`$ and $`s`$ directions. Next, we make the substitution $`\chi =\text{exp}\left[i\frac{e}{\mathrm{}c}\underset{0}{\overset{q}{}}A_q(s,q^{})𝑑q^{}\right]\stackrel{~}{\chi }`$, and expand $`h`$, $`A_q(s,q)`$, $`A_s(s,q)`$ in series in $`q`$ keeping only the zero-order terms in $`q`$, as in Refs. costa81, ; shevchenko01, . The Schrödinger equation (6) can then be easily separated by setting $`\stackrel{~}{\chi }(s,q)=\nu (q)\varphi (s)`$. The usual procedure yields
$$\frac{\mathrm{}^2}{2m^{}}\frac{^2\nu }{q^2}+V_\gamma \nu =E_t\nu ,$$
(7)
$`[{\displaystyle \frac{\mathrm{}^2}{2m^{}}}{\displaystyle \frac{^2}{s^2}}+i{\displaystyle \frac{e\mathrm{}}{m^{}c}}A_s(s,0){\displaystyle \frac{}{s}}{\displaystyle \frac{\mathrm{}^2k^2}{8m^{}}}{\displaystyle \frac{ie\mathrm{}}{2m^{}c}}\times `$
$`{\displaystyle \frac{A_q(s,0)}{q}}+i{\displaystyle \frac{e\mathrm{}}{m^{}c}}A_q(s,0){\displaystyle \frac{k}{2}}+{\displaystyle \frac{e^2}{2m^{}c^2}}A_s^2(s,0)]\varphi =E_l\varphi .`$ (8)
In order to further simplify Eq. (8), we perform the transformation $`\varphi (s)=\text{exp}\left[i\frac{e}{\mathrm{}c}\underset{0}{\overset{s}{}}A_s(s^{},0)𝑑s^{}\right]\stackrel{~}{\varphi }(s)`$, which gives
$`[{\displaystyle \frac{\mathrm{}^2}{2m^{}}}{\displaystyle \frac{^2}{s^2}}{\displaystyle \frac{\mathrm{}^2k^2}{8m^{}}}{\displaystyle \frac{ie\mathrm{}}{2m^{}c}}{\displaystyle \frac{A_s(s,0)}{s}}`$
$`{\displaystyle \frac{ie\mathrm{}}{2m^{}c}}a_1(s)+i{\displaystyle \frac{e\mathrm{}}{m^{}c}}a_0{\displaystyle \frac{k}{2}}]\stackrel{~}{\varphi }(s)=E_t\stackrel{~}{\varphi }(s).`$ (9)
Notice that in the curvilinear coordinates ($`s`$, $`q`$) the divergence of $`𝐀`$ is given by
$`\text{div}𝐀={\displaystyle \frac{1}{h}}{\displaystyle \frac{}{_s}}\left({\displaystyle \frac{1}{h}}A_s\right)+{\displaystyle \frac{1}{h}}{\displaystyle \frac{}{q}}\left(hA_q\right)`$
$`{\displaystyle \frac{}{_s}}A_s+{\displaystyle \frac{}{q}}A_qkA_q=0.`$ (10)
Consequently, Eq. (9) reduces to
$$\frac{\mathrm{}^2}{2m^{}}\frac{^2\stackrel{~}{\varphi }}{s^2}\frac{\mathrm{}^2k^2(s)}{8m^{}}\stackrel{~}{\varphi }=E_l\stackrel{~}{\varphi }.$$
(11)
Therefore, we have derived two decoupled equations: one describing the transverse confinement of electrons in the ring (Eq. (7)), and the second describing the longitudinal motion of the electron in the ring (Eq. (11)). The vector potential $`𝐀`$ does not explicitly appear in these two equations. However it will appear in the solution of Eq. (11) because of the boundary conditions on $`\stackrel{~}{\varphi }`$ specified below. The spectrum of Eq. (7) depends on the particular shape of the confinement potential $`V_\gamma `$. In this paper we assume that the electrons occupy only the lowest subband of the transversal confinement. Therefore, the position of this energy subband is not important. The curvature of C enters into Eq. (11) through the geometrical potential term $`\mathrm{}^2k(s)^2/8m^{}`$.
The boundary conditions for $`\stackrel{~}{\varphi }`$ are obtained from the requirements of continuity of the wave function $`\varphi (s)`$ and its derivative, i.e. $`\varphi (0)=\varphi (L)`$, $`\varphi (0)/s=\varphi (L)/s`$ ($`L`$ is the ring circumference). Using Stokes’ theorem we finally obtain
$$\stackrel{~}{\varphi }(0)=e^{i2\pi \frac{\mathrm{\Phi }}{\mathrm{\Phi }_0}}\stackrel{~}{\varphi }(L),\frac{\stackrel{~}{\varphi }(0)}{s}=\frac{\stackrel{~}{\varphi }(L)}{s}.$$
(12)
Here $`\mathrm{\Phi }`$ is the magnetic flux through the area confined by C and $`\mathrm{\Phi }_0`$ is the magnetic flux quantum. Eqs. (12) imply that all equilibrium physical properties of a narrow closed loop are periodic in $`\mathrm{\Phi }`$ with period $`\mathrm{\Phi }_0`$, as in the case of a perfect ring. r3
## III Electron energy spectrum
In this section we consider the electron spectrum in a ring with a dent-type distortion as shown in Fig. 1. Geometrically, the curve C of such a ring consists of four smoothly connected circular segments. The radius of the long segment is $`r_1`$, the three short segments have the same radius $`r_2`$ (we assume here that $`r_2<r_1`$). The angles $`\alpha `$ and $`\beta `$ are related to $`\phi `$ (for the definition of these angles see Fig. 1) as $`\alpha =2\text{Arcsin}\left[\mathrm{sin}(\phi /2)(r_1r_2)/(2r_2)\right]`$, $`\beta =\phi /2+\alpha /2`$. The advantage of a distortion with constant curvature segments is a constant geometrical potential in each segment. The picture is even more simple since the geometrical potential does not depend on the direction of bending, i.e. $`V_{geom}`$ is the same for the three short segments in Fig. 1. Therefore, we write Eq. (11) for the long segment and three short segments as
$`{\displaystyle \frac{\mathrm{}^2}{2m^{}}}{\displaystyle \frac{^2\varphi _1}{s^2}}U_0\varphi _1=E_l\varphi _1\text{for }0<s<l\text{,}`$ (13)
$`{\displaystyle \frac{\mathrm{}^2}{2m^{}}}{\displaystyle \frac{^2\varphi _2}{s^2}}=E_l\varphi _2\text{for }l<s<L,`$ (14)
where $`l=\left(2\beta +\alpha \right)r_2`$ is the total length of the short segments, $`L=(2\pi \phi )r_1+l`$, $`U_0=\frac{\mathrm{}^2}{8m^{}}(1/r_2^21/r_1^2)`$. The general solution of Eqs. (13, 14) reads
$`\varphi _1=b_1e^{ik_1s}+b_2e^{ik_1s},`$ (15)
$`\varphi _2=c_1e^{ik_2s}+c_2e^{ik_2s},`$ (16)
where $`k_1=\sqrt{\frac{2m^{}}{\mathrm{}}(E_l+U_0)}`$ and $`k_2=\sqrt{\frac{2m^{}}{\mathrm{}}E_l}`$. The wave functions $`\varphi _{1(2)}`$ are connected at $`s=l`$ via $`\varphi _1(l)=\varphi _2(l)`$, $`\varphi _1(l)/s=\varphi _2(l)/s`$ and at $`l=0,L`$ via Eqs. (12). From these boundary conditions we obtain a transcendental equation defining the energy spectrum for unbound states $`E_l>0`$
$`2\mathrm{cos}(2\pi {\displaystyle \frac{\mathrm{\Phi }}{\mathrm{\Phi }_0}})+\left[{\displaystyle \frac{k_1}{k_2}}+{\displaystyle \frac{k_2}{k_1}}\right]\mathrm{sin}(k_1l)\mathrm{sin}(k_2(Ll))`$
$`2\mathrm{cos}(k_1l)\mathrm{cos}(k_2(Ll))=0`$ (17)
and for bound states $`U_0<E_l<0`$
$`2\mathrm{cos}\left(2\pi {\displaystyle \frac{\mathrm{\Phi }}{\mathrm{\Phi }_0}}\right)+\left[{\displaystyle \frac{k_1}{\stackrel{~}{k}_2}}{\displaystyle \frac{\stackrel{~}{k}_2}{k_1}}\right]\mathrm{sin}(k_1l)\mathrm{sinh}(\stackrel{~}{k}_2(Ll))`$
$`2\mathrm{cos}(k_1l)\mathrm{cosh}(\stackrel{~}{k}_2(Ll))=0.`$ (18)
Here, $`\stackrel{~}{k}_2=\sqrt{\frac{2m^{}}{\mathrm{}}(E_l)}`$.
The calculated energy levels for weakly and strongly distorted rings as a function of the flux are given in Fig. 2. It is well known that in perfect rings the energy levels are intersecting parabolas. In distorted rings, gaps are opened at the points of intersection of the parabolas. This effect is qualitatively similar to the effect of disorder. r3 We emphasize that for a fixed radius of distortion, the gap decreases for larger values of the intersection point energy. At a fixed point of intersection, the gap increases by decreasing the radius of distortion (i.e. the gap is larger in more distorted rings). Notice finally that due to the distortion the effective circumference of the ring increases. This produces negative shifts of the energy levels which are larger at higher energy. The qualitatively new feature of the spectrum is the presence of bound states with $`E_l<0`$. Similar bound states were already discussed in elliptical quantum rings. magarill96 It is interesting that the transition from unbound to bound states in the ring is smooth: the shallow bound states are still sensitive to the magnetic flux (Fig. 2, $`r_2/r_1=0.25`$). In contrast, deep bound states have a weak sensitivity to the magnetic flux (Fig. 2, $`r_2/r_1=0.17`$). Correspondingly, there is a finite contribution to the persistent current from the shallow bound states, while the contribution from the deep bound states is small.
## IV Effect of the distortion on the current
### IV.1 Persistent currents induced by a magnetic flux
At non-zero temperature $`T`$, the current in the ring is given by
$$I=\frac{F}{\mathrm{\Phi }}.$$
(19)
Here, the free energy $`F=k_BT\underset{n}{}ln\left(1+\mathrm{exp}\frac{\mu E_n}{k_BT}\right)`$, $`k_B`$ is the Boltzman constant, and $`T`$ is the temperature. We consider a system with a fixed number of spinless electrons $`N`$. The chemical potential $`\mu `$ that enters into $`F`$ is determined by the equation
$$\underset{n}{}\frac{1}{1+e^{\frac{E_n\mu }{k_BT}}}=N.$$
(20)
In Fig. (3), we show the effect of distortion on the persistent current in quantum rings with three electrons. These results were obtained numerically using Eqs. (20) and (19) with the energy spectrum determined from Eqs. (17) and (18). At zero temperature the persistent current oscillations in a perfect ring have a saw-tooth form. The distortion of the ring produces a smoothing of the oscillations due to the opening of energy gaps at the intersection points. Fig. (3) shows that the smoothing increases for larger distortions. Notice that the persistent current as a function of the magnetic field flux in the distorted ring at $`T=0`$ looks similar to the persistent current in a perfect ring at $`T>0`$.
The distortion of the ring changes the temperature dependence of the current amplitude. The temperature dependence of the persistent current and its amplitude in a perfect ring with a fixed number of electrons at low temperatures was derived in PV using a two-level approach. Vagner83 It was found that at low temperatures the persistent current can be written as PV
$$I=\frac{2N\epsilon }{\mathrm{\Phi }_0}\left[\frac{\mathrm{sinh}\gamma }{1+\mathrm{cosh}\gamma }2\left(\frac{\mathrm{\Phi }}{\mathrm{\Phi }_0}\frac{\delta }{2}\right)\right],$$
(21)
where $`\gamma =N\epsilon (\mathrm{\Phi }/\mathrm{\Phi }_0\delta /2)/(k_BT)`$, $`\epsilon =\mathrm{}^2/(2m^{}r^2)`$, $`\delta =0`$ if $`N`$ is even and $`\delta =1`$ if $`N`$ is odd, and the overall factor two takes into account the spin degeneracy. The temperature dependence of the amplitude of the persistent current oscillations is given by
$`I_{max}={\displaystyle \frac{2N\epsilon }{\mathrm{\Phi }_0}}(\sqrt{14{\displaystyle \frac{k_BT}{N\epsilon }}}`$
$`2{\displaystyle \frac{k_BT}{N\epsilon }}\text{arccosh}[{\displaystyle \frac{M\epsilon }{2k_BT}}1]).`$ (22)
The comparison of the temperature dependence of the persistent current amplitude in the perfect and distorted rings is shown in the inset of Fig. 3. While the persistent current amplitude in the perfect ring starts to decrease at $`T=0`$, the temperature dependence in the distorted ring shows an activation energy behavior in the vicinity of $`T=0`$ due to the gaps in the energy spectrum.
In quantum rings with many electrons the main contribution to the persistent current is due to the electrons near the Fermi level. Since the energy gap disappears for large values of the intersection point energy, the persistent current at $`T=0`$ will not be smoothed. Consequently, the persistent current in a distorted ring with a large number of electrons will be as in a perfect ring with a radius $`r^{}=L/(2\pi )`$ in a weaker (for the case of a dent) magnetic field $`B^{}=BS/S^{}`$, where $`S`$ is the area of the distorted ring and $`S^{}=2\pi r^2`$.
### IV.2 Current induced by circularly polarized radiation
In a recent paper, per05 a novel mechanism for current generation in quantum rings was proposed. It was suggested that in the presence of a circularly-polarized continuous-wave (cw) radiation the light-dressed ground state of the ring is characterized by a non-zero current. The purpose of this Section is to study the influence of the ring distortion on this radiation-induced current.
Let us consider an electron confined in a distorted ring in the presence of circularly-polarized cw radiation. The single electron Hamiltonian in the dipole approximation reads
$$H=H_0+V(t)=\frac{\mathrm{}^2}{2m^{}}\frac{^2}{s^2}+U_{geom}(s)+e𝐄(t)𝐫(s),$$
(23)
where $`𝐄(t)=E_0\mathrm{cos}(\omega t)\widehat{x}\pm E_0\mathrm{cos}(\omega t)\widehat{x}`$ is the circularly-polarized electric field, $`E_0`$ is its amplitude, and $`\pm `$ corresponds to $`\sigma _\pm `$ radiation. The distorted quantum ring is considered again as made of four constant curvature segments, which allows us to use the energy spectrum and wave functions of $`H_0`$ obtained in Sec. III at $`\mathrm{\Phi }=0`$. Assuming that the radiation frequency is close to the transition between the ground and two first excited levels, we restrict our attention only to these three levels, with energies given by $`E_0`$, $`E_1`$ and $`E_2`$.
The external radiation causes transitions between these levels. The electron dynamics in the ring can be conveniently described using the a density matrix approach similar to the one used in Ref. takagahara, for quantum dots. The evolution of density matrix $`\rho `$ is given by
$$i\mathrm{}\dot{\rho }=[H,\rho ]\mathrm{\Gamma }\{\rho \},$$
(24)
where $`\mathrm{\Gamma }\{\rho \}`$ represents a relaxation terms. In the rotating wave approximation the corresponding equations for the density matrix elements are
$`\dot{\rho }_{00}=v_{01}\stackrel{~}{\rho }_{10}+v_{02}\stackrel{~}{\rho }_{20}\stackrel{~}{\rho }_{01}v_{10+}\stackrel{~}{\rho }_{02}v_{20+}+\kappa _{20}\rho _{22}+\kappa _{10}\rho _{11},`$ (25)
$`\dot{\rho }_{11}=v_{10+}\stackrel{~}{\rho }_{01}\stackrel{~}{\rho }_{10}v_{01}\kappa _{10}\rho _{11}+\kappa _{21}\rho _{22},`$ (26)
$`\dot{\rho }_{22}=v_{20+}\stackrel{~}{\rho }_{02}\stackrel{~}{\rho }_{20}v_{02}\kappa _{20}\rho _{22}\kappa _{21}\rho _{22},`$ (27)
$`\dot{\stackrel{~}{\rho }}_{01}={\displaystyle \frac{E_0E_1+\mathrm{}\omega }{i\mathrm{}}}\stackrel{~}{\rho }_{01}+v_{01}\rho _{11}+v_{02}\rho _{21}\rho _{00}v_{01}\gamma _{01}\stackrel{~}{\rho }_{01,}`$ (28)
$`\dot{\stackrel{~}{\rho }}_{02}={\displaystyle \frac{E_0E_2+\mathrm{}\omega }{i\mathrm{}}}\stackrel{~}{\rho }_{02}+v_{01}\rho _{12}+v_{02}\rho _{22}\rho _{00}v_{02}\gamma _{02}\stackrel{~}{\rho }_{02},`$ (29)
$`\dot{\rho }_{12}={\displaystyle \frac{E_1E_2}{i\mathrm{}}}\rho _{12}+v_{10+}\stackrel{~}{\rho }_{02}\stackrel{~}{\rho }_{10}v_{02}\gamma _{12}\rho _{12},`$ (30)
$`\dot{\stackrel{~}{\rho }}_{10}={\displaystyle \frac{E_1E_0\mathrm{}\omega }{i\mathrm{}}}\stackrel{~}{\rho }_{10}+v_{10+}\rho _{00}\rho _{11}v_{10+}\rho _{12}v_{20+}\gamma _{10}\stackrel{~}{\rho }_{10},`$ (31)
$`\dot{\stackrel{~}{\rho }}_{20}={\displaystyle \frac{E_2E_0\mathrm{}\omega }{i\mathrm{}}}\stackrel{~}{\rho }_{20}+v_{20+}\rho _{00}\rho _{21}v_{10+}\rho _{22}v_{20+}\gamma _{20}\stackrel{~}{\rho }_{20},`$ (32)
$`\dot{\rho }_{21}={\displaystyle \frac{E_2E_1}{i\mathrm{}}}\rho _{21}+v_{20+}\stackrel{~}{\rho }_{01}\stackrel{~}{\rho }_{20}v_{01}\gamma _{21}\rho _{21}.`$ (33)
Here, the transformations $`\rho _{01}=e^{i\omega t}\stackrel{~}{\rho }_{01}`$, $`\rho _{02}=e^{i\omega t}\stackrel{~}{\rho }_{02}`$, $`\rho _{10}=e^{i\omega t}\stackrel{~}{\rho }_{10}`$, $`\rho _{20}=e^{i\omega t}\stackrel{~}{\rho }_{20}`$ were used, $`\kappa _{ij}`$ is the relaxation rate of diagonal density matrix elements, $`\gamma _{i,j}`$ is the dephasing rate of the off-diagonal coherences $`\rho _{ij}`$. $`v_{ij\pm }=\overline{i|V(t)|je^{\pm i\omega t}/(i\mathrm{})}`$, where $`\overline{(\mathrm{})}`$ denotes an averaging over a period of $`V(t)`$. For example, in the case of $`\sigma _+`$ radiation we obtain
$`v_{01}={\displaystyle \frac{1}{i\mathrm{}}}{\displaystyle \frac{eE_0}{2}}0|x+{\displaystyle \frac{y}{i}}|1,`$ (34)
$`v_{10+}={\displaystyle \frac{1}{i\mathrm{}}}{\displaystyle \frac{eE_0}{2}}1|x{\displaystyle \frac{y}{i}}|0.`$ (35)
The persistent current is calculated using $`I=\text{Tr}\left[\rho \widehat{j}\right]`$, where $`\widehat{j}`$ is the standard quantum mechanical current operator, and $`\rho `$ is the steady-state solution of Eqs. (25)-(33). We find that in the distorted rings the current operator matrix has a form
$$\widehat{j}=\left(\begin{array}{ccc}0& 0& j_{02}\\ 0& 0& j_{12}\\ j_{20}& j_{12}& 0\end{array}\right),$$
(36)
with $`j_{ij}=j_{ji}^{}`$. Correspondingly, the persistent current is given by
$$I=2\text{Re}\left(\rho _{21}j_{12}+\rho _{20}j_{02}\right)=2\text{Re}\left(\rho _{21}j_{12}+e^{i\omega t}\stackrel{~}{\rho }_{20}j_{02}\right).$$
(37)
The first term in the right hand side of Eq. (37) is time-independent and will be referred to as the DC component of the current, the second term in the right hand side of Eq. (37) is the AC component of the current. In a perfect ring $`j_{20}=0`$, thus an AC component in the current is a signature of ring distortion.
Figs. 4 and 5 show the DC component $`2\text{Re}\left(\rho _{21}j_{12}\right)`$ and the amplitude of the AC component $`2|\text{Re}\left(\stackrel{~}{\rho }_{20}j_{02}\right)|`$ of the radiation induced current for two values of the radiation intensity and different distortion radii. The exact steady state solutions of Eqs. (25)-(33) for these plots were found numerically. In the case of a high excitation power (Fig. 4) broad current peaks are observed. In the almost perfect ring ($`r_2/r_1=0.99`$ curve in Fig. 4(a)) the DC component has a single resonance peak. As the distortion degree increases, this peak shifts to a higher energy and its amplitude decreases (see $`r_2/r_1=0.75`$ curve). For a stronger distortion ($`r_2/r_1=0.5`$) a second peak appears at a lower energy, which is related to the lower energy splitted level $`E_1`$. At $`r_2/r_1=0.25`$ and $`r_2/r_1=0.17`$ an additional negative current peak is observed. The amplitude of the AC component of the persistent current is zero in the perfect ring. This amplitude becomes different than zero in distorted quantum rings with a maximum located in the region of the $`E_2E_0`$ resonance (Fig. 4). The complex dependence of the DC and AC persistent current components on the radiation frequency indicates that significant quantum-interference effects are occurring.
The current peaks are narrower in the case of low radiation power (see Fig. 5). Fig. 5(a) shows that the DC component is suppressed in quantum rings with strong distortion. In contrast, the amplitude of the AC component becomes non-zero in the distorted rings and increases with the distortion (Fig. 5(b)). The maximum of the AC component is located in the region of the $`E_2E_0`$ resonance, as in the case of the high radiation power. We note that the sign change of $`\text{Re}\left(\stackrel{~}{\rho }_{20}j_{02}\right)`$ is responsible for the vertical lines in the peak centers in Fig. 5(b).
In the regime of low radiation power we can find an approximate solution of Eqs. (25)-(33). Eq. (32) gives the following expression for $`\stackrel{~}{\rho }_{20}`$ in the first order in $`E_0`$:
$$\stackrel{~}{\rho }_{20}=\frac{i\mathrm{}}{E_2E_0\mathrm{}\omega i\mathrm{}\gamma _{20}}v_{20+}\rho _{00}.$$
(38)
Similarly, from Eq. (28) $`\stackrel{~}{\rho }_{01}`$ in the first order in $`E_0`$ can be found. This expression for $`\stackrel{~}{\rho }_{01}`$ together with Eq. (38) and Eq. (33) yields in the second order in $`E_0`$
$`\rho _{21}={\displaystyle \frac{i\mathrm{}}{E_2E_1i\mathrm{}\gamma _{21}}}[{\displaystyle \frac{i\mathrm{}}{E_2E_0\mathrm{}\omega i\mathrm{}\gamma _{20}}}+`$
$`{\displaystyle \frac{i\mathrm{}}{E_1E_0\mathrm{}\omega i\mathrm{}\gamma _{01}}}]v_{01}v_{20+}\rho _{00}.`$ (39)
We have found that the persistent current components calculated from Eqs. (37), (38), (39) with $`\rho _{0,0}=1`$ perfectly coincides with the persistent current components calculated numerically in Fig. 5. In the case of a perfect ring the current can be seen as a $`𝒳_2`$ effect. The presence of a distorsion induces a $`𝒳_1`$ term which corresponds to the AC component.
## V Conclusions
In conclusion, we have investigated persistent and radiation-induced currents in quantum rings with distortions. We have derived an effective Schrödinger equation describing electrons in a narrow distorted quantum ring (closed loop) in the presence of an external magnetic field flux. We have shown that the electron energy spectrum is a periodic function of the magnetic flux. The ring curvature enters into the effective equations through a geometrical potential term. We have solved the equations in the case of a distorted ring consisting of four constant-curvature segments. We have considered the effect of the ring distortion on the magnetic flux-induced and radiation-induced currents. It was found that the effect on the flux-induced persistent current is more pronounced in quantum rings with a small number of electrons and lower chemical potential. The gaps at the points of intersection of the energy levels lead to a smoothing of the persistent current oscillations and to a different temperature dependence. The persistent current in a distorted ring with a large number of electrons behaves like in a perfect ring with a different radius in a renormalized magnetic field.
We have also found that the ring distortion affects radiation-induced currents. Using a density matrix approach and the rotating wave approximation, it was found that the current in distorted quantum rings acquires an AC component, in addition to the DC component characteristic of perfect rings. The frequency dependence of the DC component is modified by the distortion and shows several peaks of different sign. The frequency of the AC component is equal to the radiation frequency while its amplitude increases with the distortion. Finally, we remark that the non-trivial dependence of the DC persistent current component on the radiation frequency can be useful for quantum control schemes involving localized spins, as suggested in Ref. per05, .
We thank Prof. M. Dykman for many fruitful discussions. This research was supported by the National Science Foundation, Grant NSF DMR-0312491. |
warning/0506/cond-mat0506596.html | ar5iv | text | # Magnetic phase diagram copper metaborate CuB2O4 in magnetic field parallel c-axis: resonant, magnetic and magnetoelastic investigations
## 1 Introduction
The tetragonal crystal of copper metaborate CuB<sub>2</sub>O<sub>4</sub> has the complex magnetic structure which examination was carried out by various experimental methods including a neutron scattering , $`\mu `$SR , a magnetic resonance and magnetic measurements . Neutron investigations showed that at $`T<T_{spon}=9.5`$ K the magnetic state of a crystal is incommensurate with a magnetic propagation vector directional along a tetragonal axis. Resonant and magnetic measurements allowed to assume that in the temperature interval from $`9.5`$ K up to Néel temperature $`T_N=20`$ K the ground state is also modulated and long periodical. In a magnetic field perpendicular to the tetragonal axis both modulated states transform into the field-induced weak ferromagnetic state with the magnetic moment laying in a basal plane of a crystal.
The magnetic phase diagram of copper metaborate in the perpendicular field is given in . The purpose of the work is the experimental investigation of the magnetic phase diagram of CuB<sub>2</sub>O<sub>4</sub> in the magnetic field along a tetragonal axis. Phase transitions are explored by experimental methods sensitive to magnetic state of the crystal: electron-spin resonance (ESR), magnetic and magnetostrictive measurements.
## 2 Experimental results
Resonant and magnetic measurements are performed on single crystals of CuB<sub>2</sub>O<sub>4</sub> growthing by spontaneous crystallization method. The samples having a form of a plate with sizes of up to $`2\times 7\times 7`$ mm<sup>3</sup> which were cut out in $`(100)`$ and $`(110)`$ crystallographic planes were used in magnetostriction measurements. Resonant investigations are carried out on computer-controlled magnetic resonance spectrometer with pulsed magnetic field . Magnetostriction measurements were done using the capacitance technique. Magnetic measurements are performed by a SQUID magnetometer MPMS-5.
Temperature dependencies of resonance field and line width measured for $`Hc`$ at different frequencies have sharp anomalies (Fig.1). The maximal line broadening decreases with the increasing of a frequency and, correspondingly, a magnetic field.
To ensure that anomalies are due to the transition from incommensurate to commensurate weak ferromagnetic state the dependencies of longitudinal and transversal magnetization have been measured at $`Hc`$. The longitudinal magnetization (Fig. 2a) increases smoothly with a field increase at any temperatures below Néel temperature, however, the dependences become more nonlinear at lower temperatures. The base saturation level of a longitudinal magnetization is reached at $`T=2`$ K in a field about $`30`$ kOe after which, as show measurements up to $`350`$ kOe, the weak linear rise of a magnetization is observed. At the same time the transversal magnetization (Fig. 2b) is near zero below some critical field and have a jump about $`0.6`$ emu/g at the critical field which is increasing with temperature lowering. The magnetization in magnetic field higher than critical one increases first reaching some maximum and then decreases gradually with further increase in a field.
The field dependences of transversal magnetization are also measured at $`T>T_{spon}`$ and differ considerably from the low temperature ones. The magnetization at $`T>T_{spon}`$ rises continuously with magnetic field showing a kink point at the critical field. In this temperature interval the critical values of fields are much lower than at $`T<9.5`$ K.
Magnetostrictive examinations also allow to register transition between phases. We measured field dependences of a longitudinal and transversal magnetostriction for $`Hc`$ at various temperatures. The typical field dependences of a longitudinal magnetostriction at temperatures below $`9.5`$ K are shown in Fig. 4. All dependences have common character: weak rise of a magnetostriction above and below a critical field and the jump at the transition. The field dependences of a transversal magnetostriction have a similar view.
## 3 Discussion
Complexity of the magnetic phase diagram of copper metaborate is due to coexistence of two subsystems of copper ions with a various degree of the magnetic order and different magnetic dimensionality. It is known , that in unit cell of CuB<sub>2</sub>O<sub>4</sub> 12 copper ions $`Cu^{2+}`$ occupy two non-equivalent positions — $`4b`$ and $`8d`$. Four ions in a position $`4b`$ form the three-dimensional subsystem that is magnetically ordered below Néel temperature $`T_N=20`$ K (the strong subsystem A). Other eight ions form a weak ordered subsystem B which is the one-dimensional and is partially polarized due to an exchange interaction with ions of the strong subsystem.
The analysis of resonant properties of CuB<sub>2</sub>O<sub>4</sub> allow to assume that the resonant absorption in this crystal at orientation of a magnetic field along a tetragonal axis is connected with the weak subsystem B. The following arguments confirm such explanation. First, below Néel temperature a strong subsystem at $`Hc`$ in both commensurate and incommensurate states can be considered as easy-plane antiferromagnet. The spectrum of antiferromagnetic resonance (AFMR) of this subsystem at such orientation of a field contains a branch with nonlinear frequency-field dependence and an energy gap $`\omega _c\gamma \sqrt{2H_EH_A}`$, where $`H_E`$ and $`H_A`$ — effective fields, accordingly, of exchange and anisotropy with respect to the tetragonal axis (the second branch of AFMR is Goldstone with $`\omega =0`$). As the field of anisotropy $`H_A`$ for CuB<sub>2</sub>O<sub>4</sub> is unknown it is impossible to estimate the value of the gap but the usual its value in uniaxial antiferromagnets is about several hundreds in GHz. At the same time the magnetic resonance data at this field orientation (Fig. 5) show that the frequency-field dependences are near-linear with neglible values of the gaps in both commensurate and incommensurate states.
The data of inelastic neutron scattering also show two spin wave branches, one of which, high-energy, has a gap $`840`$ GHz at $`T=1.5`$ K and is attributed to the strong subsystem A, and another branch with small initial splitting is referred by authors to a weak subsystem.
It is visible from the inset of Fig.1 that the temperature dependence of line intensity is well described as $`IC/(T\theta )`$ with $`\theta 2`$ K what is usual for the disordered systems. At last, the full absence of anomalies of intensity and line width at the $`T_N`$ of the strong subsystem also allows to connect the observed resonant absorption to a weak subsystem B.
In our opinion the broadening of the resonance line near the temperature of phase transition is caused by the fluctuations that are increasing with the approach to the transition. This result correlates with anomalous increase of diffuse neutron scattering at phase transition . With increase of frequency and, accordingly, a resonant field the external magnetic field suppresses fluctuations more strongly, and line spreading decreases.
Field dependences of longitudinal magnetization copper metaborate, measured along a tetragonal axis, confirm the assumption about a various degree of the magnetic order in subsystems A and B. Calculations show that saturation magnetizations of subsystems along the $`c`$-axis are: $`M_S^A=13.54`$ emu/g and $`M_S^B=27.07`$ emu/g. The field dependence of longitudinal magnetization measured at $`T=2`$ K have two areas: sharp nonlinear increasing with saturation in fields approximately up to 30 kOe and practically linear subsequent increase in fields up to 350 kOe. Saturation occurs at the level corresponding to a weak subsystem, hence, this subsystem is saturated in magnetic fields up to $`30`$ kOe at $`T=2`$ K. The subsequent linear rise of magnetization is caused, mainly, by the antiferromagnetic susceptibility of the strong subsystem A.
It is obvious that the jumps of magnetization in a basal plane at $`T<9.5`$ K (Fig. 2b) are caused by the phase transition into the field-induced weak ferromagnetic state. The inset in Fig. 2 shows the magnetic hysteresis, hence at $`T<9.5`$ K the phase transition at $`Hc`$, as well as in basal plane is of the first order. There is no magnetic hysteresis of transversal magnetization at the critical field at $`T>T_{spon}`$ (Fig. 3). Thus the phase transition has the second order at the temperature range from $`9.5`$ to $`20`$ K.
The initial part of the field dependences of transversal magnetization also are partially caused by magnetization due to the projection of a field in the basic plane, and also nonideal orthogonality of measuring coils and the magnetic field. The measured transversal magnetization is the sum of contributions of both magnetic subsystems A and B, therefore the field dependence of the magnetization in the basal plane above the critical field is defined by two processes. On the one hand, the saturation of the weak subsystem along the $`c`$-axis in strong fields results in reduction of its contribution to total magnetization in the basal plane. On the other hand, there is an increase in the total magnetization caused by the increasing of a magnetic field component in the basal plane due to the nonideal its orientation along the $`c`$-axis. Clearly, that the main contribution to the last process is caused by the strong subsystem. Due to the competition of these two processes the total magnetization increases first with the increasing of a magnetic field above its critical value, then starts to fall.
The analysis of the results on magnetostriction shows that its smooth increase with a field above and below the critical value is caused by the contribution of magnetized weak-ordered subsystem B. Taking into account that both longitudinal and transversal magnetostrictions have similar field dependences we can assume that the magnetostriction is a volume dilatation at $`Hc`$. It is obvious that the jumps of magnetostriction are caused by the change of the magnetic state of CuB<sub>2</sub>O<sub>4</sub> at the phase transition.
Thus, the boundaries of magnetic phases of CuB<sub>2</sub>O<sub>4</sub> magnetized along the tetragonal axis are established by resonant, magnetic and magnetostrictive measurements. We have measured also an angular dependence of the critical field at temperatures below $`9.5`$ K, this dependence has clearly defined maximum at $`Hc`$ (Fig. 6). Because of this sharp dependence we show on the phase diagram (Fig. 7) only the most authentic phase boundaries obtained by measurements of longitudinal magnetostriction and transversal magnetization at which the magnetic field was closest to the $`c`$-axis. An estimate shows that at such angular dependence the inaccuracy of installation of samples at measurements in 2–3 degrees is capable to explain apparent spread of values of the critical field. In addition to our data, the temperature of spontaneous phase transition measured by heat capacity is marked, and the values of critical fields measured with the second optical harmonic generation well agreed with our data are also presented.
The resulted diagram is similar to the phase diagram CuB<sub>2</sub>O<sub>4</sub> magnetized in a basal plane , but the critical fields for $`Hc`$ are much higher. The state $`1`$ on the diagram corresponds to an incommensurate spiral phase, and the state $`2`$ is the field-induced weak ferromagnetic one. The nature of a state $`3`$ is unknown, we assume that it is also modulated, but the wave vector of modulation is much less than the resolution of neutron diffraction. The phase boundaries correspond to the first order transition and to the second one at temperatures, respectively, below and above $`T_{spon}=9.5`$ K.
It is necessary to tell, that the existence of phase transition from incommensurate to commensurate state in a magnetic field laying in a basal plane is not surprised. In this case the field laying in a plane of a spiral deforms this structure, transforming it to fan, and then — to commensurate state. At the same time the phase transition in a magnetic field oriented along a wave vector of spiral structure seems at first sight surprising as energies of spiral and commensurate structures in a magnetic field at such magnetization are identical. We assume that the physical reason of incommensurate - commensurate phase transition at $`T<9.5`$ K and in $`Hc`$ implies that the weak subsystem together with the strong subsystem plays the important role in a formation of spiral structure. And when the weak subsystem is saturated by a field along the tetragonal axis and its contribution to the formation of spiral structure changes.
The authors are greatly indebted to M.A Popov and S.N. Martynov for helpful discussions.
This work was supported by the Russian Foundation for Basic Research (grant RFBR 03-02-16701). |
warning/0506/nlin0506057.html | ar5iv | text | # Abstract
## Abstract
A direct method for the computation of polynomial conservation laws of polynomial systems of nonlinear partial differential equations (PDEs) in multi-dimensions is presented. The method avoids advanced differential-geometric tools. Instead, it is solely based on calculus, variational calculus, and linear algebra.
Densities are constructed as linear combinations of scaling homogeneous terms with undetermined coefficients. The variational derivative (Euler operator) is used to compute the undetermined coefficients. The homotopy operator is used to compute the fluxes.
The method is illustrated with nonlinear PDEs describing wave phenomena in fluid dynamics, plasma physics, and quantum physics. For PDEs with parameters, the method determines the conditions on the parameters so that a sequence of conserved densities might exist. The existence of a large number of conservation laws is a predictor for complete integrability. The method is algorithmic, applicable to a variety of PDEs, and can be implemented in computer algebra systems such as Mathematica, Maple, and REDUCE.
## 1 Introduction
Nonlinear partial differential equations (PDEs) that admit conservation laws arise in many disciplines of the applied sciences including physical chemistry, fluid mechanics, particle and quantum physics, plasma physics, elasticity, gas dynamics, electromagnetism, magneto-hydro-dynamics, nonlinear optics, and the bio-sciences. Conservation laws are fundamental laws of physics. They maintain that a certain quantity, e.g. momentum, mass (matter), electric charge, or energy, will not change with time during physical processes. Often the PDE itself is a conservation law, e.g. the continuity equation relating charge to current.
As shown in and the articles in this issue, computer algebra systems (CAS) like Mathematica, Maple, and REDUCE, are useful to tackle computational problems in chemistry. Finding closed-form conservation laws of nonlinear PDEs is a nice example. Using CAS interactively, we could make a judicious guess (ansatz) and find a few simple densities and fluxes. Yet, that approach is fruitless for complicated systems with nontrivial conservation laws. Furthermore, completely integrable PDEs admit infinitely many independent conservation laws. Computing them is a challenging task. It involves tedious computations which are prone to error if done with pen and paper. The most famous example is the Korteweg-de Vries (KdV) equation from soliton theory which describes water waves in shallow water, ion-acoustic waves in plasmas, etc. Our earlier work dealt primarily with the symbolic computation of conservation laws of completely integrable PDEs in $`(1+1)`$ dimensions (with independent variables $`x`$ and $`t).`$ In this paper we present a symbolic method that covers PDEs in multi-dimensions (e.g. $`x,y,z,`$ and $`t),`$ irrespective of their complete integrability. As before, our approach relies on the concept of dilation (scaling) invariance which limits it to polynomial conserved densities of polynomial PDEs in evolution form.
There are many reasons to compute conserved densities and fluxes of PDEs explicitly. Invariants often lead to new discoveries as was the case in soliton theory. We may want to verify if conserved quantities of physical importance (e.g. momentum, energy, Hamiltonians, entropy, density, charge) are intact after constitutive relations have been added to close a system. For PDEs with arbitrary parameters we may wish to compute conditions on the parameters so that the model admits conserved quantities. Conserved densities also facilitate the study of qualitative properties of PDEs, such as bi- or tri-Hamiltonian structures. They often guide the choice of solution methods or reveal the nature of special solutions. For example, an infinite sequence of conserved densities assures complete integrability of the PDE, i.e. solvability by the Inverse Scattering Transform and the existence of solitons .
Conserved densities aid in the design of numerical solvers for PDEs . Indeed, semi-discretizations that conserve discrete conserved quantities lead to stable numerical schemes (i.e. free of nonlinear instabilities and blowup). While solving differential-difference equations (DDEs), which arise in nonlinear lattices and as semi-discretizations of PDEs, one should check that their conserved quantities indeed remain unchanged. Capitalizing on the analogy between PDEs and DDEs, the techniques presented in this paper have been adapted to DDEs and fully discretized lattices .
There are various methods (see ) to compute conservation laws of nonlinear PDEs. A common approach relies on the link between conservation laws and symmetries as stated in Noether’s theorem . However, the computation of generalized (variational) symmetries, though algorithmic, is as daunting a task as the direct computation of conservation laws. Nonetheless, we draw the reader’s attention to `DE_APPLS`, a package for constructing conservation laws from symmetries available within Vessiot , a general purpose suite of Maple packages for computations on jet spaces. Other methods circumvent the existence of a variational principle but still rely on the symbolic solution of a determining system of PDEs. Despite their power, only a few of the above methods have been fully implemented in CAS (see ).
We purposely avoid Noether’s theorem, pre-knowledge of symmetries, and a Lagrangian formulation. Neither do we use differential forms or advanced differential-geometric tools. Instead, we present and implement our tools in the language of calculus, linear algebra, and variational calculus. Our down-to-earth calculus formulas are transparent, easy to use by scientists and engineers, and are readily adaptable to nonlinear DDEs (not covered in Vessiot).
To design a reliable algorithm, we must compute both densities and fluxes. For the latter, we need to invert the divergence operator which requires the integration (by parts) of an expression involving arbitrary functions. That is where the homotopy operator comes into play. Indeed, the homotopy operator reduces that problem to a standard one-dimensional integration with respect to a single auxiliary parameter. One of the first<sup>1</sup><sup>1</sup>1We refer the reader to \[26, p. 374\] for a brief history of the homotopy operator. uses of the homotopy operator in the context of conservation laws can be found in . The homotopy operator is a universal, yet little known, tool that can be applied to many problems in which integration by parts in multi-variables is needed. A literature search revealed that homotopy operators are used in integrability testing and inversion problems involving PDEs, DDEs, lattices, and beyond . Assuming the reader is unfamiliar with homotopy operators, we “demystify” the homotopy formulas and make them ready for use in nonlinear sciences.
The particular application in this paper is computation of conservation laws for which our algorithm proceeds as follows: build a candidate density as a linear combination (with undetermined coefficients) of “building blocks” that are homogeneous under the scaling symmetry of the PDE. If no such symmetry exists, construct one by introducing parameters with scaling. Subsequently, use the conservation equation and the variational derivative to derive a linear algebraic system for the undetermined coefficients. After the system is solved, use the homotopy operator to compute the flux. Implementations for $`(1+1)`$-dimensional PDEs in Mathematica and Maple can be downloaded from . Mathematica code that automates the computations for PDEs in multi-dimensions is being designed.
This paper is organized as follows. We establish some connections with vector calculus in Section 2. In Section 3 we list the nonlinear PDEs that will be used throughout the paper: the Korteweg-de Vries and Boussinesq equations from soliton theory , the Landau-Lifshitz equation for Heisenberg’s ferromagnet and a system of shallow water wave equations for ocean waves . Sections 4 and 5 cover the dilation invariance and conservation laws of those four examples. In Section 6, we introduce and apply the tools from the calculus of variations. Formulas for the variational derivative in multi-dimensions are given in Section 6.1, where we apply the Euler operator for testing the exactness of various expressions. Removal of divergence-equivalent terms with the Euler operator is discussed in Section 6.2. The higher Euler operators and homotopy operators in 1D and 2D are in Sections 6.3 and 6.4. In the latter section, we apply the homotopy operator to integrate by parts and to invert of the divergence operator. In Section 7 we present our three-step algorithm to compute conservation laws and apply it to the KdV equation (Section 7.1), the Boussinesq equation (Section 7.2), the Landau-Lifshitz equation (Section 7.3), and the shallow water wave equations (Section 7.4). Conclusions are drawn in Section 8.
## 2 Connections with Vector Calculus
We address a few issues in multivariate calculus which complement the material in : (i) To determine whether or not a vector field $`𝐅`$ is conservative, i.e. $`𝐅=\mathbf{}f`$ for some scalar field $`f,`$ one must verify that $`𝐅`$ is irrotational or curl free, that is $`\mathbf{}\times 𝐅=\mathrm{𝟎}.`$ The cross $`(\times )`$ denotes the Euclidean cross product. The field $`f`$ can be computed via standard integrations \[24, p. 518, 522\]. (ii) To test if $`𝐅`$ is the curl of some vector field $`𝐆,`$ one must check that $`𝐅`$ is incompressible or divergence free, i.e. $`\mathbf{}𝐅=0.`$ The dot $`()`$ denotes the Euclidean inner product. The components of $`𝐆`$ result from solving a coupled system of first-order PDEs \[24, p. 526\]. (iii) To verify whether or not a scalar field $`f`$ is the divergence of some vector function $`𝐅,`$ no theorem from vector calculus comes to the rescue. Furthermore, the computation of $`𝐅`$ such that $`f=\mathbf{}𝐅`$ is a nontrivial matter which requires the use of the homotopy operator or Hodge decomposition . In single variable calculus, it amounts to computing the primitive $`F=f𝑑x.`$
In multivariate calculus, all scalar fields $`f,`$ including the components $`F_i`$ of vector fields $`𝐅=(F_1,F_2,F_3),`$ are functions of the independent variables $`(x,y,z).`$ In differential geometry the above issues are addressed in much greater generality. The functions $`f`$ and $`F_i`$ can now depend on arbitrary functions $`u(x,y,z),v(x,y,z),`$ etc. and their mixed derivatives (up to a fixed order) with respect to $`(x,y,z).`$ Such functions are called differential functions . As one might expect, carrying out the gradient-, curl-, or divergence-test requires advanced algebraic machinery. For example, to test whether or not $`f=\mathbf{}𝐅`$ requires the use of the variational derivative (Euler operator) in 2D or 3D as we will show in Section 6.1. The actual computation of $`𝐅`$ requires integration by parts. That is where the higher Euler and homotopy operators enter the picture (see Sections 6.3 and 6.4).
At the moment, no major CAS have reliable routines for integrating expressions involving arbitrary functions and their derivatives. As far as we know, CAS offer no functions to test if a differential function is a divergence. Routines to symbolically invert the total divergence are certainly lacking. In Section 6 we will present these tools and apply them in Section 7.
## 3 Examples of Nonlinear PDEs
Definition: We consider a nonlinear system of evolution equations in $`(3+1)`$ dimensions,
$$𝐮_t=𝐆(𝐮,𝐮_x,𝐮_y,𝐮_z,𝐮_{2x},𝐮_{2y},𝐮_{2z},𝐮_{xy},𝐮_{xz},𝐮_{yz},\mathrm{}),$$
(1)
where $`𝐱=(x,y,z)`$ and $`t`$ are space and time variables. The vector $`𝐮(x,y,z,t)`$ has $`N`$ components $`u_i.`$ In the examples we denote the components of $`𝐮`$ by $`u,v,w,`$ etc.. Throughout the paper we use the subscript notation for partial derivatives,
$$𝐮_t=\frac{𝐮}{t},𝐮_{2x}=𝐮_{xx}=\frac{^2𝐮}{x^2},𝐮_{2x3y}=𝐮_{xxyyy}=\frac{^5𝐮}{x^2y^3},𝐮_{xy}=\frac{^2𝐮}{xy},etc..$$
(2)
We assume that $`𝐆`$ is smooth and does not explicitly depend on $`𝐱`$ and $`t.`$ There are no restrictions on the number of components, order, and degree of nonlinearity of $`𝐆.`$
We will predominantly work with polynomial systems, although systems involving one transcendental nonlinearity can also be handled . If parameters are present in (1), they will be denoted by lower-case Greek letters. Throughout this paper we work with the following four prototypical PDEs: Example 1: The ubiquitous Korteweg-de Vries (KdV) equation
$$u_t+uu_x+u_{3x}=0,$$
(3)
for $`u(x,t)`$ describes unidirectional shallow water waves and ion-acoustic waves in plasmas. Example 2: The wave equation,
$$u_{tt}u_{2x}+3uu_{2x}+3u_x^2+\alpha u_{4x}=0,$$
(4)
for $`u(x,t)`$ with real parameter $`\alpha ,`$ was proposed by Boussinesq to describe surface waves in shallow water . For what follows, we rewrite (4) as a system of evolution equations,
$$u_t+v_x=0,v_t+u_x3uu_x\alpha u_{3x}=0,$$
(5)
where $`v(x,t)`$ is an auxiliary dependent variable. Example 3: The classical Landau-Lifshitz (LL) equation,
$$𝐒_t=𝐒\times \mathrm{\Delta }𝐒+𝐒\times D𝐒,$$
(6)
without external magnetic field, models nonlinear spin waves in a continuous Heisenberg ferromagnet . The spin vector $`𝐒(x,t)`$ has real components $`(u(x,t),v(x,t),w(x,t));`$ $`\mathrm{\Delta }=\mathbf{}^2`$ is the Laplacian, and $`D=\mathrm{diag}(\alpha ,\beta ,\gamma )`$ is a diagonal matrix with real coupling constants $`\alpha ,\beta ,`$ and $`\gamma `$ along the downward diagonal. Splitting (6) into components we get
$`u_t`$ $`=`$ $`vw_{2x}wv_{2x}+(\gamma \beta )vw,`$
$`v_t`$ $`=`$ $`wu_{2x}uw_{2x}+(\alpha \gamma )uw,`$ (7)
$`w_t`$ $`=`$ $`uv_{2x}vu_{2x}+(\beta \alpha )uv.`$
The second and third equations can be obtained from the first by cyclic permutations
$$(uv,vw,wu)\mathrm{and}(\alpha \beta ,\beta \gamma ,\gamma \alpha ).$$
(8)
We take advantage of the cyclic nature of (3) in the computations in Section 7.3. Example 4: The (2+1)-dimensional shallow-water wave (SWW) equations,
$`𝐮_t+(𝐮\mathbf{})𝐮+2𝛀\times 𝐮=\mathbf{}(h\theta )+\frac{1}{2}h\mathbf{}\theta ,`$
$`h_t+\mathbf{}(h𝐮)=0,\theta _t+𝐮(\mathbf{}\theta )=0,`$ (9)
describe waves in the ocean using layered models . Vectors $`𝐮=u(x,y,t)𝐢+v(x,y,t)𝐣`$ and $`𝛀=\mathrm{\Omega }𝐤`$ are the fluid and angular velocities. $`\mathbf{}=\frac{}{x}𝐢+\frac{}{y}𝐣`$ is the gradient operator and $`𝐢`$, $`𝐣,`$ and $`𝐤`$ are unit vectors along the $`x`$, $`y,`$ and $`z`$-axes. $`\theta (x,y,t)`$ is the horizontally varying potential temperature field and $`h(x,y,t)`$ is the layer depth. System (3) can be written as
$`u_t+uu_x+vu_y2\mathrm{\Omega }v+\frac{1}{2}h\theta _x+\theta h_x=0,v_t+uv_x+vv_y+2\mathrm{\Omega }u+\frac{1}{2}h\theta _y+\theta h_y=0,`$
$`h_t+hu_x+uh_x+hv_y+vh_y=0,\theta _t+u\theta _x+v\theta _y=0.`$ (10)
## 4 Key Concept: Dilation or Scaling Invariance
Definition: System (1) is dilation invariant if it does not change under a scaling symmetry. Example: The KdV equation (3) is dilation invariant under the scaling symmetry
$$(x,t,u)(\lambda ^1x,\lambda ^3t,\lambda ^2u),$$
(11)
where $`\lambda `$ is an arbitrary parameter. Indeed, apply the chain rule to verify that $`\lambda ^5`$ factors out. Definition: The weight<sup>2</sup><sup>2</sup>2The weights are the remnants of physical units after non-dimensionalization of the PDE., $`W,`$ of a variable is the exponent in $`\lambda ^p`$ which multiplies the variable. Example: We will always replace $`x`$ by $`\lambda ^1x.`$ Thus, $`W(x)=1`$ or $`W(/x)=1.`$ In view of (11), $`W(/t)=3`$ and $`W(u)=2`$ for the KdV equation (3). Definition: The rank of a monomial is defined as the total weight of the monomial. An expression is uniform in rank if its monomial terms have equal rank. Example: All monomials in (3) have rank $`5.`$ Thus, (3) is uniform in rank with respect to (11). Weights of dependent variables and weights of $`/x,/y,`$ etc. are assumed to be non-negative and rational. Ranks must be positive integer or rational numbers. The ranks of the equations in (1) may differ from each other.
Conversely, requiring uniformity in rank for each equation in (1) we can compute the weights, and thus the scaling symmetry, with linear algebra. Example: Indeed, for the KdV equation (3) we have
$$W(u)+W(/t)=2W(u)+1=W(u)+3,$$
(12)
where we used $`W(/x)=1.`$ Solving (12) yields
$$W(u)=2,W(/t)=3.$$
(13)
This means that $`x`$ scales with $`\lambda ^1,t`$ with $`\lambda ^3,`$ and $`u`$ with $`\lambda ^2,`$ as claimed in (11). Dilation symmetries, which are special Lie-point symmetries , are common to many nonlinear PDEs. However, non-uniform PDEs can be made uniform by giving appropriate weights to parameters that appear in the PDEs, or, if necessary, by extending the set of dependent variables with auxiliary parameters with appropriate weights. Upon completion of the computations we can set the auxiliary parameters equal to $`1.`$ Example: The Boussinesq system (5) is not uniform in rank because the terms $`u_x`$ and $`\alpha u_{3x}`$ lead to a contradiction in the weight equations. To circumvent the problem we introduce an auxiliary parameter $`\beta `$ with (unknown) weight, and replace (5) by
$`u_t+v_x`$ $`=`$ $`0,`$
$`v_t+\beta u_x3uu_x\alpha u_{3x}`$ $`=`$ $`0.`$ (14)
Requiring uniformity in rank, we obtain (after some algebra)
$$W(u)=2,W(v)=3,W(\beta )=2,W(\frac{}{t})=2.$$
(15)
Therefore, (4) is invariant under the scaling symmetry
$$(x,t,u,v,\beta )(\lambda ^1x,\lambda ^2t,\lambda ^2u,\lambda ^3v,\lambda ^2\beta ).$$
(16)
Definition: System (1) is called multi-uniform in rank if it admits more than one scaling symmetry (none of which result from introducing unnecessary parameters with weights). Example: The LL system (3) is not uniform in rank unless we allow the parameters $`\alpha ,\beta ,`$ and $`\gamma `$ to have weights. Doing so, uniformity in rank requires
$`W(u)+W({\displaystyle \frac{}{t}})`$ $`=`$ $`W(v)+W(w)+2=W(\gamma )+W(v)+W(w)=W(\beta )+W(v)+W(w),`$
$`W(v)+W({\displaystyle \frac{}{t}})`$ $`=`$ $`W(u)+W(w)+2=W(\alpha )+W(u)+W(w)=W(\gamma )+W(u)+W(w),`$ (17)
$`W(w)+W({\displaystyle \frac{}{t}})`$ $`=`$ $`W(u)+W(v)+2=W(\beta )+W(u)+W(v)=W(\alpha )+W(u)+W(v).`$
Solving these relations gives $`W(u)=W(v)=W(w),W(\alpha )=W(\beta )=W(\gamma )=2,W(\frac{}{t})=W(u)+2,`$ where $`W(u)`$ is arbitrary. The LL equation is thus multi-uniform and invariant under the following class of scaling symmetries,
$$(x,t,u,v,w,\alpha ,\beta ,\gamma )(\lambda ^1x,\lambda ^{(a+2)}t,\lambda ^au,\lambda ^av,\lambda ^aw,\lambda ^2\alpha ,\lambda ^2\beta ,\lambda ^2\gamma ),$$
(18)
where $`W(u)=a`$ is an arbitrary non-negative integer or rational number. For example, if we take $`a=1`$ then $`W(u)=W(v)=W(w)=1,W(\alpha )=W(\beta )=W(\gamma )=2,W(\frac{}{t})=3`$ and (3) is invariant under $`(x,t,u,v,w,\alpha ,\beta ,\gamma )(\lambda ^1x,\lambda ^3t,\lambda u,\lambda v,\lambda w,\lambda ^2\alpha ,\lambda ^2\beta ,\lambda ^2\gamma ).`$ Taking different choices for $`a<W(\alpha )=2`$ will prove advantageous for the computations in Section 7.3. Example: The SWW equations (3) are not uniform in rank unless we give a weight to $`\mathrm{\Omega }.`$ Indeed, uniformity in rank for (3) requires, after some algebra, that
$`W(/t)=W(\mathrm{\Omega }),W(/y)=W(/x)=1,W(u)=W(v)=W(\mathrm{\Omega })1,`$ (19)
$`W(\theta )=2W(\mathrm{\Omega })W(h)2,`$
where $`W(h)`$ and $`W(\mathrm{\Omega })`$ are arbitrary. Hence, the SWW system is multi-uniform and invariant under the class of scaling symmetries,
$$(x,y,t,u,v,h,\theta ,\mathrm{\Omega })(\lambda ^1x,\lambda ^1y,\lambda ^bt,\lambda ^{b1}u,\lambda ^{b1}v,\lambda ^ah,\lambda ^{2ba2}\theta ,\lambda ^b\mathrm{\Omega }).$$
(20)
where $`W(h)=a`$ and $`W(\mathrm{\Omega })=b.`$ Various choices for $`a`$ and $`b`$ will aid the computations in Section 7.4.
## 5 Conservation Laws
Definition: A conservation law in differential form is a PDE of the form,
$$\mathrm{D}_t\rho +\mathrm{Div}𝐉=0,$$
(21)
which is satisfied on solutions of (1). The total time derivative $`\mathrm{D}_t`$ and total divergence $`\mathrm{Div}`$ are defined below. The scalar differential function $`\rho `$ is called the conserved density; the vector differential function $`𝐉`$ is the associated flux. In electromagnetism, (21) is the continuity equation relating charge density $`\rho `$ to current $`𝐉.`$
In general, $`\rho `$ and $`𝐉`$ are functions of $`𝐱,t,𝐮,`$ and the derivatives of $`𝐮`$ with respect to the components of $`𝐱.`$ In this paper, we will assume that densities and fluxes (i) are local, which means free of integral terms; (ii) polynomial in $`𝐮`$ and its derivatives; and (iii) do not explicitly depend on $`𝐱`$ and $`t.`$ So, $`\rho (𝐮,𝐮_x,𝐮_y,𝐮_z,𝐮_{2x},\mathrm{})`$ and $`𝐉(𝐮,𝐮_x,𝐮_y,𝐮_z,𝐮_{2x},\mathrm{}).`$ Definition: There is a close relationship between conservation laws and constants of motion. Indeed, integration of (21) over all space yields the integral form of a conservation law,
$$P=_{\mathrm{}}^{\mathrm{}}_{\mathrm{}}^{\mathrm{}}_{\mathrm{}}^{\mathrm{}}\rho 𝑑x𝑑y𝑑z=\mathrm{constant},$$
(22)
provided that $`𝐉`$ vanishes at infinity. As in mechanics, the $`P`$’s are called constants of motion.
The total divergence Div is computed as $`\mathrm{Div}𝐉=(\mathrm{D}_x,\mathrm{D}_y,\mathrm{D}_z)(J_1,J_2,J_3)=\mathrm{D}_xJ_1+\mathrm{D}_yJ_2+\mathrm{D}_zJ_3.`$ In the 1D case, with one spatial variable $`(x),`$ conservation equation (21) reduces to
$$\mathrm{D}_t\rho +\mathrm{D}_xJ=0,$$
(23)
where both density $`\rho `$ and flux $`J`$ are scalar differential functions. The conservation laws (21) and (23) involve total derivatives $`\mathrm{D}_t`$ and $`\mathrm{D}_x.`$ In the 1D case,
$$\mathrm{D}_t\rho =\frac{\rho }{t}+\underset{k=0}{\overset{n}{}}\frac{\rho }{u_{kx}}\mathrm{D}_x^k(u_t).$$
(24)
where $`n`$ is the order of $`u`$ in $`\rho .`$ Upon replacement of $`u_t,u_{tx},`$ etc. from $`u_t=G,`$ we get
$$\mathrm{D}_t\rho =\frac{\rho }{t}+\rho (u)^{}[G],$$
(25)
where $`\rho (u)^{}[G]`$ is the Fréchet derivative of $`\rho `$ in the direction of $`G.`$ Similarly,
$$\mathrm{D}_xJ=\frac{J}{x}+\underset{k=0}{\overset{m}{}}\frac{J}{u_{kx}}u_{(k+1)x},$$
(26)
where $`m`$ is the order of $`u`$ of $`J.`$ Example: Assume that the KdV equation (3) has a density of the form
$$\rho =c_1u^3+c_2u_x^2,$$
(27)
where $`c_1`$ and $`c_2`$ are constants. Since $`\rho `$ does not explicitly depend on time and is of order $`n=1`$ in $`u`$, application of (24) gives
$`\mathrm{D}_t\rho `$ $`=`$ $`{\displaystyle \frac{\rho }{u}}\mathrm{I}(u_t)+{\displaystyle \frac{\rho }{u_x}}\mathrm{D}_x(u_t)`$ (28)
$`=`$ $`3c_1u^2u_t+2c_2u_xu_{tx}=3c_1u^2(uu_x+u_{3x})2c_2u_x(uu_x+u_{3x})_x`$
$`=`$ $`3c_1u^2(uu_x+u_{3x})2c_2u_x(u_x^2+uu_{2x}+u_{4x})`$
$`=`$ $`(3c_1u^3u_x+3c_1u^2u_{3x}+2c_2u_x^3+2c_2uu_xu_{2x}+2c_2u_xu_{4x}),`$
where $`\mathrm{D}_x^0=\mathrm{I}`$ is the identity operator and where we used (3) to replace $`u_t`$ and $`u_{tx}.`$ The generalizations of (24) and (26) to multiple dependent variables is straightforward . Example: Taking $`𝐮(x,t)=(u(x,t),v(x,t),w(x,t)),`$
$`\mathrm{D}_t\rho `$ $`=`$ $`{\displaystyle \frac{\rho }{t}}+{\displaystyle \underset{k=0}{\overset{n_1}{}}}{\displaystyle \frac{\rho }{u_{kx}}}\mathrm{D}_x^k(u_t)+{\displaystyle \underset{k=0}{\overset{n_2}{}}}{\displaystyle \frac{\rho }{v_{kx}}}\mathrm{D}_x^k(v_t)+{\displaystyle \underset{k=0}{\overset{n_3}{}}}{\displaystyle \frac{\rho }{w_{kx}}}\mathrm{D}_x^k(w_t),`$ (29)
$`\mathrm{D}_xJ`$ $`=`$ $`{\displaystyle \frac{J}{x}}+{\displaystyle \underset{k=0}{\overset{m_1}{}}}{\displaystyle \frac{J}{u_{kx}}}u_{(k+1)x}+{\displaystyle \underset{k=0}{\overset{m_2}{}}}{\displaystyle \frac{J}{v_{kx}}}v_{(k+1)x}+{\displaystyle \underset{k=0}{\overset{m_3}{}}}{\displaystyle \frac{J}{w_{kx}}}w_{(k+1)x}`$ (30)
where $`n_1,n_2`$ and $`n_3`$ are the highest orders of $`u,v`$ and $`w`$ in $`\rho ,`$ and $`m_1,m_2`$ and $`m_3`$ are the highest orders of $`u,v,w`$ in $`J.`$ Example: Assume that the Boussinesq system (4) has a density of the form
$$\rho =c_1\beta ^2u+c_2\beta u^2+c_3u^3+c_4v^2+c_5u_xv+c_6u_x^2.$$
(31)
where $`c_1`$ through $`c_6`$ are constants. Since $`n_1=1`$ and $`n_2=0`$, application of (29) gives
$`\mathrm{D}_t\rho `$ $`=`$ $`{\displaystyle \frac{\rho }{u}}\mathrm{I}(u_t)+{\displaystyle \frac{\rho }{u_x}}\mathrm{D}_x(u_t)+{\displaystyle \frac{\rho }{v}}\mathrm{I}(v_t)`$ (32)
$`=`$ $`(c_1\beta ^2+2c_2\beta u+3c_3u^2)u_t+(c_5v+2c_6u_x)u_{tx}+(2c_4v+c_5u_x)v_t`$
$`=`$ $`((c_1\beta ^2+2c_2\beta u+3c_3u^2)v_x+(c_5v+2c_6u_x)v_{2x}+(2c_4v+c_5u_x)(\beta u_x3uu_x\alpha u_{3x})),`$
where we replaced $`u_t,u_{tx},`$ and $`v_t`$ from (4). Remark: The flux $`𝐉`$ in (21) is not uniquely defined. In 1D we use (23) and the flux is only determined up to an arbitrary constant. In 2D, the flux is only determined up to a divergence-free vector $`𝐊=(K_1,K_2)=(\mathrm{D}_y\varphi ,\mathrm{D}_x\varphi ),`$ where $`\varphi `$ is an arbitrary scalar differential function. In 3D, the flux can only be determined up to a curl term. Indeed, if $`(\rho ,𝐉)`$ is a valid density-flux pair, so is $`(\rho ,𝐉+\mathbf{}\times 𝐊)`$ for any arbitrary vector differential function $`𝐊=(K_1,K_2,K_3).`$ Recall that $`\mathbf{}\times 𝐊=(\mathrm{D}_yK_3\mathrm{D}_zK_2,\mathrm{D}_zK_1\mathrm{D}_xK_3,\mathrm{D}_xK_2\mathrm{D}_yK_1).`$ Example: The Korteweg-de Vries equation (3) is known to have infinitely many polynomial conservation laws. The first three density-flux pairs are
$`\rho ^{(1)}`$ $`=`$ $`u,J^{(1)}={\displaystyle \frac{1}{2}}u^2+u_{2x},`$
$`\rho ^{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}u^2,J^{(2)}={\displaystyle \frac{1}{3}}u^3{\displaystyle \frac{1}{2}}u_x^2+uu_{2x},`$ (33)
$`\rho ^{(3)}`$ $`=`$ $`{\displaystyle \frac{1}{3}}u^3u_x^2,J^{(3)}={\displaystyle \frac{1}{4}}u^42uu_x^2+u^2u_{2x}+u_{2x}^22u_xu_{3x}.`$
The first two express conservation of momentum and energy, respectively. They are easy to compute by hand. The third one, less obvious and requiring more work, corresponds to Boussinesq’s moment of instability . Observe that the above densities are uniform of ranks $`2,4,`$ and $`6.`$ The fluxes are also uniform of ranks $`4,6,`$ and $`8.`$
The computations get quickly out of hand for higher ranks. For example, for rank $`12`$
$$\rho ^{(6)}=\frac{1}{6}u^610u^3u_x^25u_x^4+18u^2u_{2x}^{}{}_{}{}^{2}+\frac{120}{7}u_{2x}^{}{}_{}{}^{3}\frac{108}{7}uu_{3x}^{}{}_{}{}^{2}+\frac{36}{7}u_{4x}^{}{}_{}{}^{2}$$
(34)
and $`J^{(6)}`$ is a scaling homogeneous polynomial with $`20`$ terms of rank 14 (not shown). In general, if in (23) $`\mathrm{rank}(\rho )=R`$ then $`\mathrm{rank}(J)=R+W(/t)1.`$ All the pieces in (21) are also uniform in rank. This comes as no surprise since the conservation law (21) holds on solutions of (1), hence it ‘inherits’ the dilation symmetry of (1). Example: The Boussinesq equation (4) also admits infinitely many conservation laws and is completely integrable . The first four density-flux pairs for (4) are
$`\rho ^{(1)}`$ $`=`$ $`u,J^{(1)}=v,`$
$`\rho ^{(2)}`$ $`=`$ $`v,J^{(2)}=\beta u{\displaystyle \frac{3}{2}}u^2\alpha u_{2x},`$ (35)
$`\rho ^{(3)}`$ $`=`$ $`uv,J^{(3)}={\displaystyle \frac{1}{2}}\beta u^2u^3+{\displaystyle \frac{1}{2}}v^2+{\displaystyle \frac{1}{2}}\alpha u_x^2\alpha uu_{2x},`$
$`\rho ^{(4)}`$ $`=`$ $`\beta u^2u^3+v^2+\alpha u_x^2,J^{(4)}=2\beta uv3u^2v2\alpha u_{2x}v+2\alpha u_xv_x.`$
The densities are of ranks $`2,3,5`$ and $`6,`$ respectively. The corresponding fluxes are of one rank higher. After setting $`\beta =1`$ we obtain the conserved quantities of (5) even though initially this system was not uniform in rank. Example: We assume that $`\alpha ,\beta ,`$ and $`\gamma `$ in (6) are nonzero. Cases with vanishing parameters must be investigated separately. The first six density-flux pairs are
$`\rho ^{(1)}`$ $`=`$ $`u,J^{(1)}=v_xwvw_x,(\beta =\gamma \alpha ),`$
$`\rho ^{(2)}`$ $`=`$ $`v,J^{(2)}=uw_xu_xw,(\alpha =\gamma \beta ),`$
$`\rho ^{(3)}`$ $`=`$ $`w,J^{(3)}=u_xvuv_x,(\alpha =\beta \gamma ),`$
$`\rho ^{(4)}`$ $`=`$ $`u^2+v^2+w^2,J^{(4)}=0,`$ (36)
$`\rho ^{(5)}`$ $`=`$ $`(u^2+v^2+w^2)^2,J^{(5)}=0,`$
$`\rho ^{(6)}`$ $`=`$ $`u_x^2+v_x^2+w_x^2+(\gamma \alpha )u^2+(\gamma \beta )v^2,`$
$`J^{(6)}`$ $`=`$ $`2((vw_xv_xw)u_{2x}+(u_xwuw_x)v_{2x}+(uv_xu_xv)w_{2x}.`$
$`.+(\beta \gamma )u_xvw+(\gamma \alpha )uv_xw+(\alpha \beta )uvw_x).`$
Note that the second and third density-flux pairs follow from the first pair via the cyclic permutations in (8). If we select $`W(u)=W(v)=W(w)=a=\frac{1}{4},`$ then the first three densities have rank $`\frac{1}{4}`$ and $`\rho ^{(4)},\rho ^{(5)},`$ and $`\rho ^{(6)}`$ have ranks $`\frac{1}{2},1,\frac{5}{2},`$ respectively.
The physics of (6) demand that the magnitude $`S=𝐒`$ of the spin field is constant in time. Indeed, using (3) we readily verify that $`uu_t+vv_t+ww_t=0.`$ Hence, any differentiable function of $`S`$ is also conserved in time. This is apparent in $`\rho ^{(4)}=𝐒^2`$ and $`\rho ^{(5)}=𝐒^4,`$ which are conserved (in time) even for the fully anisotropic case where $`\alpha \beta \gamma \alpha .`$
Obviously, linear combinations of conserved densities are also conserved. Hence, $`\rho ^{(6)}=𝐒_x^2+(\gamma \alpha )u^2+(\gamma \beta )v^2`$ can be replaced by
$$\stackrel{~}{\rho }^{(6)}=\frac{1}{2}(\rho ^{(6)}\gamma \rho ^{(4)})=\frac{1}{2}(u_x^2+v_x^2+w_x^2\alpha u^2\beta v^2\gamma w^2).$$
(37)
Now, $`u_x^2=\mathrm{D}_x(uu_x)uu_{2x},`$ etc., and densities are divergence-equivalent (see Section 6.2) if they differ by a total derivative with respect to $`x.`$ So, (37) is equivalent with
$$\stackrel{~}{\rho }^{(6)}=\frac{1}{2}(uu_{2x}+vv_{2x}+ww_{2x}+\alpha u^2+\beta v^2+\gamma w^2),$$
(38)
which can be compactly written as
$$\stackrel{~}{\rho }^{(6)}=\frac{1}{2}(𝐒\mathrm{\Delta }𝐒+𝐒D𝐒),$$
(39)
where $`\mathrm{\Delta }`$ is the Laplacian and $`D=\mathrm{diag}(\alpha ,\beta ,\gamma ).`$ So, $`\stackrel{~}{\rho }^{(6)}`$ is the Hamiltonian density in
$`H`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \left(𝐒\mathrm{\Delta }𝐒+𝐒D𝐒\right)𝑑x}`$ (40)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \left(𝐒_x^2+(\gamma \alpha )u^2+(\gamma \beta )v^2\right)𝑑x}.`$
The Hamiltonian $`H`$ expresses that the total energy is constant in time since $`dH/dt=0.`$ Example: The SWW equations (3) admit a Hamiltonian formulation and infinitely many conservation laws. Yet, that is still insufficient to guarantee complete integrability of PDEs in $`(2+1)`$-dimensions. The first few conserved densities and fluxes for (3) are
$$\begin{array}{cccc}\rho ^{(1)}=h,\hfill & 𝐉^{(1)}=\left(\begin{array}{c}uh\\ vh\end{array}\right),\hfill & \rho ^{(2)}=h\theta ,\hfill & 𝐉^{(2)}=(\begin{array}{c}uh\theta \\ vh\theta \end{array}),\hfill \\ \rho ^{(3)}=h\theta ^2,\hfill & 𝐉^{(3)}=(\begin{array}{c}uh\theta ^2\\ vh\theta ^2\end{array}),\hfill & & \\ \multicolumn{2}{c}{\rho ^{(4)}=(u^2+v^2)h+h^2\theta ,}& \multicolumn{2}{c}{𝐉^{(4)}=(\begin{array}{c}u^3h+uv^2h+2uh^2\theta \\ v^3h+u^2vh+2vh^2\theta \end{array}),}\\ \multicolumn{4}{c}{\rho ^{(5)}=(2\mathrm{\Omega }u_y+v_x)\theta ,}\\ \multicolumn{4}{c}{𝐉^{(5)}=\frac{1}{6}\left(\begin{array}{c}12\mathrm{\Omega }u\theta 4uu_y\theta +6uv_x\theta +2vv_y\theta +u^2\theta _y+v^2\theta _yh\theta \theta _y+h_y\theta ^2\\ 12\mathrm{\Omega }v\theta +4vv_x\theta 6vu_y\theta 2uu_x\theta u^2\theta _xv^2\theta _x+h\theta \theta _xh_x\theta ^2\end{array}\right).}\end{array}$$
(41)
As shown in , system (3) has conserved densities $`\rho =hf(\theta )`$ and $`\rho =(2\mathrm{\Omega }u_y+v_x)g(\theta ),`$ where $`f(\theta )`$ and $`g(\theta )`$ are arbitrary functions. Such densities are the integrands of the Casimirs of the Poisson bracket associated with (3). The algorithm presented in Section 7.4 only computes densities of the form $`\rho =h\theta ^k`$ and $`\rho =(2\mathrm{\Omega }u_y+v_x)\theta ^l,`$ where $`k`$ and $`l`$ are positive integers.
## 6 Tools from the Calculus of Variations
In this section we introduce the variational derivative (Euler operator), the higher Euler operators (also called Lie-Euler operators) from the calculus of variations, and the homotopy operator from homological algebra and variational bi-complexes . These tools will be applied in Section 7.
### 6.1 Variational Derivative (Euler Operator)
Definition: A scalar differential function $`f`$ is a divergence if and only if there exists a vector differential function $`𝐅`$ such that $`f=\mathrm{Div}𝐅.`$ In 1D, we say that a differential function $`f`$ is exact<sup>3</sup><sup>3</sup>3We do not use integrable to avoid confusion with complete integrability from soliton theory . if and only if there exists a scalar differential function $`F`$ such that $`f=\mathrm{D}_xF.`$ Obviously, $`F=\mathrm{D}_x^1(f)=f𝑑x`$ is then the primitive (or integral) of $`f.`$ Example: Taking $`f=\mathrm{D}_t\rho `$ from (28), that is
$$f=3c_1u^3u_x+2c_2u_x^3+2c_2uu_xu_{2x}+3c_1u^2u_{3x}+2c_2u_xu_{4x},$$
(42)
where<sup>4</sup><sup>4</sup>4Variable $`t`$ is parameter in subsequent computations. For brevity, we write $`u(x)`$ instead of $`u(x;t).`$ $`u(x;t)`$, we will compute $`c_1`$ and $`c_2`$ so that $`f`$ is exact.
Upon repeated integration by parts (by hand), we get
$$F=f𝑑x=\frac{3}{4}c_1u^4+(c_23c_1)uu_x^2+3c_1u^2u_{2x}c_2u_{2x}^2+2c_2u_xu_{3x}+(3c_1+c_2)u_x^3𝑑x.$$
(43)
Removing the “obstruction” (the last term) requires $`c_2=3c_1.`$ Substituting a solution, $`c_1=\frac{1}{3},c_2=1,`$ into (27) gives $`\rho ^{(3)}`$ in (5). Substituting that solution into (42) and (43) yields
$$f=u^3u_x2u_x^32uu_xu_{2x}+u^2u_{3x}2u_xu_{4x}$$
(44)
which is exact, together with its integral
$$F=\frac{1}{4}u^42uu_x^2+u^2u_{2x}+u_{2x}^22u_xu_{3x},$$
(45)
which is $`J^{(3)}`$ in (5). Currently, CAS like Mathematica, Maple<sup>5</sup><sup>5</sup>5Future versions of Maple will be able to “partially” integrate such expressions ., and REDUCE cannot compute (43) as a sum of terms that can be integrated out and the obstruction. Example: Consider the following example in 2D
$$f=u_xv_yu_{2x}v_yu_yv_x+u_{xy}v_x,$$
(46)
where $`u(x,y)`$ and $`v(x,y).`$ It is easy to verify by hand (via integration by parts) that $`f`$ is exact. Indeed, $`f=\mathrm{Div}𝐅`$ with
$$𝐅=(uv_yu_xv_y,uv_x+u_xv_x).$$
(47)
As far as we know, the leading CAS currently lack tools to compute $`𝐅.`$ Two questions arise: (i) How can we tell whether $`f`$ is the divergence of some differential function $`𝐅\mathrm{?}`$ (ii) How do we compute $`𝐅`$ automatically and without multivariate integration by parts of expressions involving arbitrary functions? To answer these questions we use the following tools from the calculus of variations: the variational derivative (Euler operator), its generalizations (higher Euler operators or Lie-Euler operators), and the homotopy operator. Definition: The variational derivative (zeroth Euler operator), $`_{𝐮(𝐱)}^{(\mathrm{𝟎})}`$, is defined \[26, p. 246\] by
$$_{𝐮(𝐱)}^{(\mathrm{𝟎})}=\underset{J}{}(\mathrm{D})_J\frac{}{𝐮_J},$$
(48)
where the sum is over all the unordered multi-indices $`J`$ \[26, p. 95\]. For example, in the 2D case the multi-indices corresponding to second-order derivatives can be identified with $`\{2x,2y,2z,xy,xz,yz\}.`$ Obviously, $`(D)_{2x}=(\mathrm{D}_x)^2`$$`=\mathrm{D}_x^2`$, $`(D)_{xy}`$ $`=(\mathrm{D}_x)(\mathrm{D}_y)`$ $`=\mathrm{D}_x\mathrm{D}_y,`$ etc.. For notational details see \[26, p. 95, p. 108, p. 246\].
With applications in mind, we give calculus-based formulas for the variational derivatives in 1D and 2D. The formula for 3D is analogous and can be found in . Example: In 1D, where $`𝐱=x,`$ for scalar component $`u(x)`$ we have<sup>6</sup><sup>6</sup>6Variable $`t`$ in $`𝐮(𝐱;t)`$ is suppressed because it is a parameter in the Euler operators.
$$_{u(x)}^{(0)}=\underset{k=0}{\overset{\mathrm{}}{}}(\mathrm{D}_x)^k\frac{}{u_{kx}}=\frac{}{u}\mathrm{D}_x\frac{}{u_x}+\mathrm{D}_x^2\frac{}{u_{2x}}\mathrm{D}_x^3\frac{}{u_{3x}}+\mathrm{},$$
(49)
In 2D where $`𝐱=(x,y),`$ we have for component $`u(x,y)`$
$`_{u(x,y)}^{(0,0)}`$ $`=`$ $`{\displaystyle \underset{k_x=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{k_y=0}{\overset{\mathrm{}}{}}}(\mathrm{D}_x)^{k_x}(\mathrm{D}_y)^{k_y}{\displaystyle \frac{}{u_{k_xxk_yy}}}={\displaystyle \frac{}{u}}\mathrm{D}_x{\displaystyle \frac{}{u_x}}\mathrm{D}_y{\displaystyle \frac{}{u_y}}`$ (50)
$`+\mathrm{D}_x^2{\displaystyle \frac{}{u_{2x}}}+\mathrm{D}_x\mathrm{D}_y{\displaystyle \frac{}{u_{xy}}}+\mathrm{D}_y^2{\displaystyle \frac{}{u_{2y}}}\mathrm{D}_x^3{\displaystyle \frac{}{u_{3x}}}\mathrm{},`$
Note that $`u_{k_xxk_yy}`$ stands for $`u_{xx\mathrm{}xyy\mathrm{}y}`$ where $`x`$ is repeated $`k_x`$ times and $`y`$ is repeated $`k_y`$ times. Similar formulas hold for components $`v,w,`$ etc.. The first question is then answered by the following theorem \[26, p. 248\]. Theorem: A necessary and sufficient condition for a function $`f`$ to be a divergence, i.e. there exists a differential function $`𝐅`$ so that $`f=\mathrm{Div}𝐅,`$ is that $`_{𝐮(𝐱)}^{(\mathrm{𝟎})}(f)0.`$ In words: the Euler operator annihilates divergences. If, for example, $`𝐮(𝐱)=(u(𝐱),v(𝐱))`$ then both $`_{u(𝐱)}^{(\mathrm{𝟎})}(f)`$ and $`_{v(𝐱)}^{(\mathrm{𝟎})}(f)`$ must vanish identically. For the 1D case, the theorem says that a differential function $`f`$ is exact, i.e. there exists a differential function $`F`$ so that $`f=D_xF,`$ if and only if $`_{𝐮(x)}^{(0)}(f)0.`$ Example: Avoiding integration by parts, we determine $`c_1`$ and $`c_2`$ so that $`f`$ in (42) will be exact. Since $`f`$ is of order $`4`$ in the (one) dependent variable $`u(x),`$ the zeroth Euler operator (49) terminates at $`k=4.`$ Using nothing but differentiations, we readily compute
$`_{u(x)}^{(0)}(f)`$ $`=`$ $`{\displaystyle \frac{f}{u}}\mathrm{D}_x\left({\displaystyle \frac{f}{u_x}}\right)+\mathrm{D}_x^2\left({\displaystyle \frac{f}{u_{2x}}}\right)\mathrm{D}_x^3\left({\displaystyle \frac{f}{u_{3x}}}\right)+\mathrm{D}_x^4\left({\displaystyle \frac{f}{u_{4x}}}\right)`$ (51)
$`=`$ $`9c_1u^2u_x+2c_2u_xu_{2x}+6c_1uu_{3x}\mathrm{D}_x(3c_1u^3+6c_2u_x^2+2c_2uu_{2x}+2c_2u_{4x})`$
$`+\mathrm{D}_x^2(2c_2uu_x)\mathrm{D}_x^3(3c_1u^2)+\mathrm{D}_x^4(2c_2u_x)`$
$`=`$ $`9c_1u^2u_x+2c_2u_xu_{2x}+6c_1uu_{3x}(9c_1u^2u_x+14c_2u_xu_{2x}+2c_2uu_{3x}+2c_2u_{5x})`$
$`+(6c_2u_xu_{2x}+2c_2uu_{3x})(18c_1u_xu_{2x}+6c_1uu_{3x})`$
$`=`$ $`6(3c_1+c_2)u_xu_{2x}.`$
Note that the terms in $`u^2u_x,uu_{3x},`$ and $`u_{5x}`$ dropped out. Hence, $`_{u(x)}^{(0)}(f)0`$ leads to $`3c_1+c_2=0.`$ Substituting $`c_1=\frac{1}{3},c_2=1`$ into (42) gives (44) as claimed. Example: As an example in 2D, we readily verify that $`f=u_xv_yu_{2x}v_yu_yv_x+u_{xy}v_x`$ from (46) is a divergence. Applying (50) to $`f`$ for each component of $`𝐮(x,y)=(u(x,y),v(x,y))`$ separately, we straightforwardly verify that $`_{u(x,y)}^{(0,0)}(f)0`$ and $`_{v(x,y)}^{(0,0)}(f)0.`$
### 6.2 Removing Divergences and Divergence-equivalent Terms
It is of paramount importance that densities are free of divergences for they could be moved into the flux $`𝐉`$ of the conservation law $`\mathrm{D}_t\rho +\mathrm{Div}𝐉=0.`$ We show how the Euler operator can be used to remove divergences and divergence-equivalent terms in densities. An algorithm to do so is given in . The following examples illustrate the concept behind the algorithm. Definition: Two scalar differential functions, $`f^{(1)}`$ and $`f^{(2)},`$ are divergence-equivalent if and only if they differ by the divergence of some vector $`𝐕,`$ i.e. $`f^{(1)}f^{(2)}`$ if and only if $`f^{(1)}f^{(2)}=\mathrm{Div}𝐕.`$ If a scalar expression is divergence-equivalent to zero then it is a divergence. Example: Any conserved density of the KdV equation (3) is uniform in rank with respect to $`W(u)=2`$ and $`W(/x)=1.`$ The list of all terms of, say, rank 6 is $`=\{u^3,u_x^2,uu_{2x},u_{4x}\}.`$ Now, terms $`f^{(1)}=uu_{2x}`$ and $`f^{(2)}=u_x^2`$ are divergence-equivalent because $`f^{(1)}f^{(2)}=u_x^2+uu_{2x}=\mathrm{D}_x(uu_x).`$ Using (49), note that $`_{u(x)}^{(0)}(u_x^2)=2u_{2x}`$ and $`_{u(x)}^{(0)}(uu_{2x})=2u_{2x}`$ are equal (therefore, linearly dependent). So, we discard $`uu_{2x}.`$ Moreover, $`u_{4x}=\mathrm{D}_x(u_{3x})`$ is a divergence and, as expected, $`_{u(x)}^{(0)}(u_{4x})=0.`$ So, we can discard $`u_{4x}.`$ Hence, $``$ can be replaced by $`𝒮=\{u^3,u_x^2\}`$ which is free of divergences and divergence-equivalent terms. Example: The list of non-constant terms of rank 6 for the Boussinesq equation (4) is
$$=\{\beta ^2u,\beta u^2,u^3,v^2,u_xv,u_x^2,\beta v_x,uv_x,\beta u_{2x},uu_{2x},v_{3x},u_{4x}\}.$$
(52)
Using (49), for every term $`t_i`$ in $``$ we compute $`𝐯_i=_{𝐮(x)}^{(0)}(t_i)=(_{u(x)}^{(0)}(t_i),_{v(x)}^{(0)}(t_i)).`$ If $`𝐯_i=(0,0)`$ then $`t_i`$ is discarded and so is $`𝐯_i.`$ If $`𝐯_i(0,0)`$ we verify whether or not $`𝐯_i`$ is linearly independent of the non-zero vectors $`𝐯_j,`$ $`j=1,2,\mathrm{},i1.`$ If independent, the term $`t_i`$ is kept, otherwise, $`t_i`$ is discarded and so is $`𝐯_i.`$
Upon application of (49), the first six terms in $``$ lead to linearly independent vectors $`𝐯_1`$ through $`𝐯_6.`$ Therefore, $`t_1`$ through $`t_6`$ are kept (and so are the corresponding vectors). For $`t_7=\beta v_x`$ we compute $`𝐯_7=_{𝐮(x)}^{(0)}(\beta v_x)=(0,0).`$ So, $`t_7`$ is discarded and so is $`𝐯_7.`$ For $`t_8=uv_x`$ we get $`𝐯_8=_{𝐮(x)}^{(0)}(uv_x)=(v_x,u_x)=𝐯_5.`$ So, $`t_8`$ is discarded and so is $`𝐯_8.`$
Proceeding in a similar fashion, $`t_9,t_{10},t_{11}`$ and $`t_{12}`$ are discarded. Thus, $``$ is replaced by
$$𝒮=\{\beta ^2u,\beta u^2,u^3,v^2,u_xv,u_x^2\},$$
(53)
which is free of divergences and divergence-equivalent terms. Example: In the 2D case, $`f^{(1)}=(u_xu_{2x})v_y`$ and $`f^{(2)}=(u_yu_{xy})v_x`$ are divergence-equivalent since $`f^{(1)}f^{(2)}=u_xv_yu_{2x}v_yu_yv_x+u_{xy}v_x=\mathrm{Div}(uv_yu_xv_y,uv_x+u_xv_x).`$ Using (50), note that $`_{𝐮(x,y)}^{(0)}(f^{(1)})=_{𝐮(x,y)}^{(0)}(f^{(2)})=(v_{xy}v_{xxy},u_{xy}+u_{xxy}).`$
### 6.3 Higher Euler Operators
To compute $`𝐅=\mathrm{Div}^1(f)`$ or, in the 1D case, $`F=\mathrm{D}_x^1(f)=f𝑑x,`$ we need higher-order versions of the variational derivative, called higher Euler operators or Lie-Euler operators . The general formulas in terms of differential forms are given in \[26, p. 367\]. We give calculus-based formulas for the 1D and 2D cases (see for the 3D case). Definition: The higher Euler operators in 1D (with variable $`x`$)<sup>7</sup><sup>7</sup>7Variable $`t`$ in $`𝐮(𝐱;t)`$ is suppressed because it is a parameter in the higher Euler operators. are
$$_{𝐮(x)}^{(i)}=\underset{k=i}{\overset{\mathrm{}}{}}\left(\genfrac{}{}{0pt}{}{k}{i}\right)(\mathrm{D}_x)^{ki}\frac{}{𝐮_{kx}},$$
(54)
where $`\left(\genfrac{}{}{0pt}{}{k}{i}\right)`$ is the binomial coefficient. Note that the higher Euler operator for $`i=0`$ matches the variational derivative in (49). Example: The first three higher Euler operators in 1D for component $`u(x)`$ are
$`_{u(x)}^{(1)}`$ $`=`$ $`{\displaystyle \frac{}{u_x}}2\mathrm{D}_x{\displaystyle \frac{}{u_{2x}}}+3\mathrm{D}_x^2{\displaystyle \frac{}{u_{3x}}}4\mathrm{D}_x^3{\displaystyle \frac{}{u_{4x}}}+\mathrm{},`$
$`_{u(x)}^{(2)}`$ $`=`$ $`{\displaystyle \frac{}{u_{2x}}}3\mathrm{D}_x{\displaystyle \frac{}{u_{3x}}}+6\mathrm{D}_x^2{\displaystyle \frac{}{u_{4x}}}10\mathrm{D}_x^3{\displaystyle \frac{}{u_{5x}}}+\mathrm{},`$ (55)
$`_{u(x)}^{(3)}`$ $`=`$ $`{\displaystyle \frac{}{u_{3x}}}4\mathrm{D}_x{\displaystyle \frac{}{u_{4x}}}+10\mathrm{D}_x^2{\displaystyle \frac{}{u_{5x}}}20\mathrm{D}_x^3{\displaystyle \frac{}{u_{6x}}}+\mathrm{}.`$
Definition: The higher Euler operators in 2D (with variables $`x,y`$) are
$$_{𝐮(x,y)}^{(i_x,i_y)}=\underset{k_x=i_x}{\overset{\mathrm{}}{}}\underset{k_y=i_y}{\overset{\mathrm{}}{}}\left(\genfrac{}{}{0pt}{}{k_x}{i_x}\right)\left(\genfrac{}{}{0pt}{}{k_y}{i_y}\right)(\mathrm{D}_x)^{k_xi_x}(\mathrm{D}_y)^{k_yi_y}\frac{}{𝐮_{k_xxk_yy}}.$$
(56)
Note that the higher Euler operator for $`i_x=i_y=0`$ matches the variational derivative in (50). Example: The first higher Euler operators for component $`u(x,y)`$ are
$`_{u(x,y)}^{(1,0)}`$ $`=`$ $`{\displaystyle \frac{}{u_x}}2\mathrm{D}_x{\displaystyle \frac{}{u_{2x}}}\mathrm{D}_y{\displaystyle \frac{}{u_{xy}}}+3\mathrm{D}_x^2{\displaystyle \frac{}{u_{3x}}}+2\mathrm{D}_x\mathrm{D}_y{\displaystyle \frac{}{u_{2xy}}}\mathrm{},`$
$`_{u(x,y)}^{(0,1)}`$ $`=`$ $`{\displaystyle \frac{}{u_y}}2\mathrm{D}_y{\displaystyle \frac{}{u_{2y}}}\mathrm{D}_x{\displaystyle \frac{}{u_{yx}}}+3\mathrm{D}_y^2{\displaystyle \frac{}{u_{3y}}}+2\mathrm{D}_x\mathrm{D}_y{\displaystyle \frac{}{u_{x2y}}}\mathrm{},`$ (57)
$`_{u(x,y)}^{(1,1)}`$ $`=`$ $`{\displaystyle \frac{}{u_{xy}}}2\mathrm{D}_x{\displaystyle \frac{}{u_{2xy}}}2\mathrm{D}_y{\displaystyle \frac{}{u_{x2y}}}+3\mathrm{D}_x^2{\displaystyle \frac{}{u_{3xy}}}+4\mathrm{D}_x\mathrm{D}_y{\displaystyle \frac{}{u_{2x2y}}}+\mathrm{},`$
$`_{u(x,y)}^{(2,1)}`$ $`=`$ $`{\displaystyle \frac{}{u_{2xy}}}3\mathrm{D}_x{\displaystyle \frac{}{u_{3xy}}}2\mathrm{D}_y{\displaystyle \frac{}{u_{2x2y}}}+6\mathrm{D}_x^2{\displaystyle \frac{}{u_{4xy}}}+\mathrm{D}_y^2{\displaystyle \frac{}{u_{2x3y}}}\mathrm{}.`$
The higher Euler operators are useful in their own right as the following theorem in 1D shows. Theorem: A necessary and sufficient condition for a function $`f`$ to be an $`r^{\mathrm{th}}`$ order derivative, i.e. there exists a scalar differential function $`F`$ so that $`f=D_x^rF,`$ is that $`_{𝐮(x)}^{(i)}(f)0`$ for $`i=0,1,\mathrm{},r1.`$
### 6.4 Homotopy Operators
Armed with the higher Euler operators we now turn to the homotopy operator which will allow us to reduce the computation of $`𝐅=\mathrm{Div}^1(f),`$ or in the 1D case $`F=\mathrm{D}_x^1(f)=f𝑑x,`$ to a single integral with respect to an auxiliary variable denoted by $`\lambda `$ (not to be confused with $`\lambda `$ in Section 4). Amazingly, the homotopy operator circumvents integration by parts (in multi-dimensions involving arbitrary functions) and reduces the inversion of the total divergence operator, $`\mathrm{Div},`$ to a problem of single-variable calculus.
As mentioned in Section 6.1, $`\mathrm{Div}^1`$ is defined up to a divergence-free (curl) term. In 3D, $`\mathrm{Div}^1`$ is an equivalence class $`\mathrm{Div}^1(f)=𝐅+\mathbf{}\times 𝐊`$ where $`𝐊`$ is an arbitrary vector differential function. The homotopy operator computes a particular $`𝐊.`$
The homotopy operator in terms of differential forms is given in \[26, p. 372\]. Below we give calculus-based formulas for the homotopy operators in 1D and 2D which are easy to implement in CAS. The explicit formulas in 3D are analogous (see ). Definition: The homotopy operator in 1D (with variable $`x`$)<sup>8</sup><sup>8</sup>8Variable $`t`$ in $`𝐮(𝐱;t)`$ is suppressed because it is a parameter in the homotopy operators. is
$$_{𝐮(x)}(f)=_0^1\underset{j=1}{\overset{N}{}}I_{u_j}(f)[\lambda 𝐮]\frac{d\lambda }{\lambda },$$
(58)
where $`u_j`$ is the $`j`$th component of $`𝐮`$ and the integrand $`I_{u_j}(f)`$ is given by
$$I_{u_j}(f)=\underset{i=0}{\overset{\mathrm{}}{}}\mathrm{D}_x^i\left(u_j_{u_j(x)}^{(i+1)}(f)\right).$$
(59)
The integrand involves the 1D higher Euler operators in (54). In (58), $`N`$ is the number of dependent variables and $`I_{u_j}(f)[\lambda 𝐮]`$ means that in $`I_{u_j}(f)`$ we replace $`𝐮(x)\lambda 𝐮(x),𝐮_x(x)\lambda 𝐮_x(x),`$ etc.. In practice, we first add the $`I_{u_j}(f)`$ and then scale the variables. The sum in $`I_{u_j}(f)`$ terminates at $`i=p1`$ where $`p`$ is the order of $`u_j`$ in $`f.`$
Given an exact function $`f,`$ the second question, that is how to compute $`F=\mathrm{D}_x^1(f)=f𝑑x,`$ is then answered by the following theorem \[26, p. 372\]. Theorem: For an exact function $`f,`$ one has $`F=_{𝐮(x)}(f).`$ Thus, in the 1D case, applying the homotopy operator (58) allows one to bypass integration by parts. As an experiment, one can start from some function $`\stackrel{~}{F},`$ compute $`f=\mathrm{D}_x\stackrel{~}{F},`$ then compute $`F=_{𝐮(x)}(f),`$ and finally verify that $`F\stackrel{~}{F}`$ is a constant. Example: Using the homotopy operator (58) with (59), we recompute (45). Since $`f`$ in (42) is of order $`p=4`$ in $`u_1(x)=u(x),`$ the sum in (59) terminates at $`i=p1=3.`$ Hence,
$`I_u(f)`$ $`=`$ $`u_{u(x)}^{(1)}(f)+\mathrm{D}_x\left(u_{u(x)}^{(2)}(f)\right)+\mathrm{D}_x^2\left(u_{u(x)}^{(3)}(f)\right)+\mathrm{D}_x^3\left(u_{u(x)}^{(4)}(f)\right)`$ (60)
$`=`$ $`u{\displaystyle \frac{f}{u_x}}2u\mathrm{D}_x({\displaystyle \frac{f}{u_{2x}}})+3u\mathrm{D}_x^2({\displaystyle \frac{f}{u_{3x}}})4u\mathrm{D}_x^3({\displaystyle \frac{f}{u_{4x}}})+\mathrm{D}_x(u{\displaystyle \frac{f}{u_{2x}}}3u\mathrm{D}_x({\displaystyle \frac{f}{u_{3x}}})`$
$`+6u\mathrm{D}_x^2({\displaystyle \frac{f}{u_{4x}}}))+\mathrm{D}_x^2(u{\displaystyle \frac{f}{u_{3x}}}4u\mathrm{D}_x({\displaystyle \frac{f}{u_{4x}}}))+\mathrm{D}_x^3(u{\displaystyle \frac{f}{u_{4x}}})`$
$`=`$ $`u^4+8u^2u_{2x}+6uu_{4x}+4uu_x^24\mathrm{D}_x\left(2u^2u_x+3uu_{3x}\right)+\mathrm{D}_x^2\left(u^3+8uu_{2x}\right)2\mathrm{D}_x^3\left(uu_x\right)`$
$`=`$ $`u^46uu_x^2+3u^2u_{2x}+2u_{2x}^24u_xu_{3x}.`$
Formula (58), with $`𝐮=u_1=u,`$ requires an integration with respect to $`\lambda :`$
$`F`$ $`=`$ $`_{u(x)}(f)={\displaystyle _0^1}I_u(f)[\lambda u]{\displaystyle \frac{d\lambda }{\lambda }}`$ (61)
$`=`$ $`{\displaystyle _0^1}\left(\lambda ^3u^46\lambda ^2uu_x^2+3\lambda ^2u^2u_{2x}+2\lambda u_{2x}^24\lambda u_xu_{3x}\right)𝑑\lambda `$
$`=`$ $`{\displaystyle \frac{1}{4}}u^42uu_x^2+u^2u_{2x}+u_{2x}^22u_xu_{3x}.`$
The crux of the homotopy operator method is that the integration by parts of a differential expression like (44), which involves an arbitrary function $`u(x)`$ and its derivatives, can be reduced to a standard integration of a polynomial in $`\lambda .`$ Example: For a system with $`N=2`$ components, $`𝐮(x)=(u_1(x),u_2(x))=(u(x),v(x)),`$ the homotopy operator formulas are
$$_{𝐮(x)}(f)=_0^1\left(I_u(f)+I_v(f)\right)[\lambda 𝐮]\frac{d\lambda }{\lambda },$$
(62)
with
$$I_u(f)=\underset{i=0}{\overset{\mathrm{}}{}}\mathrm{D}_x^i\left(u_{u(x)}^{(i+1)}(f)\right)\mathrm{and}I_v(f)=\underset{i=0}{\overset{\mathrm{}}{}}\mathrm{D}_x^i\left(v_{v(x)}^{(i+1)}(f)\right).$$
(63)
Example: Consider $`f=3u_xv^2\mathrm{sin}uu_x^3\mathrm{sin}u6vv_x\mathrm{cos}u+2u_xu_{2x}\mathrm{cos}u+8v_xv_{2x},`$ which is no longer polynomial in the two $`(N=2)`$ dependent variables $`u(x)`$ and $`v(x).`$ Applying Euler operator (48) for each component of $`𝐮(x)=(u(x),v(x))`$ separately, we quickly verify that $`_{u(x)}^{(0)}(f)0`$ and $`_{v(x)}^{(0)}(f)0.`$ Hence, $`f`$ is exact. Integration by parts (by hand) gives
$$F=f𝑑x=4v_x^2+u_x^2\mathrm{cos}u3v^2\mathrm{cos}u.$$
(64)
Currently, CAS like Mathematica, Maple<sup>9</sup><sup>9</sup>9Version 9.5 of Maple can integrate such expressions as a result of our interactions with the developers. and Reduce fail this integration due to the presence of trigonometric functions.
We now recompute (64) with the homotopy operator formulas (62) and (63). First,
$`I_u(f)`$ $`=`$ $`u_{u(x)}^{(1)}(f)+\mathrm{D}_x\left(u_{u(x)}^{(2)}(f)\right)=u{\displaystyle \frac{f}{u_x}}2u\mathrm{D}_x({\displaystyle \frac{f}{u_{2x}}})+\mathrm{D}_x(u{\displaystyle \frac{f}{u_{2x}}})`$ (65)
$`=`$ $`3uv^2\mathrm{sin}uuu_x^2\mathrm{sin}u+2u_x^2\mathrm{cos}u.`$
Second,
$`I_v(f)`$ $`=`$ $`v_{v(x)}^{(1)}(f)+\mathrm{D}_x(v_{v(x)}^{(2)}(f))=v{\displaystyle \frac{f}{v_x}}2v\mathrm{D}_x({\displaystyle \frac{f}{v_{2x}}})+\mathrm{D}_x(v{\displaystyle \frac{f}{v_{2x}}})`$ (66)
$`=`$ $`6v^2\mathrm{cos}u+8v_x^2.`$
Finally, using (62),
$`F`$ $`=`$ $`_{𝐮(x)}(f)={\displaystyle _0^1}\left(I_u(f)+I_v(f)\right)[\lambda 𝐮]{\displaystyle \frac{d\lambda }{\lambda }}`$ (67)
$`=`$ $`{\displaystyle _0^1}\left(3\lambda ^2uv^2\mathrm{sin}(\lambda u)\lambda ^2uu_x^2\mathrm{sin}(\lambda u)+2\lambda u_x^2\mathrm{cos}(\lambda u)6\lambda v^2\mathrm{cos}(\lambda u)+8\lambda v_x^2\right)𝑑\lambda `$
$`=`$ $`u_x^2\mathrm{cos}u3v^2\mathrm{cos}u+4v_x^2.`$
In contrast to the integration in (64) which involved arbitrary functions $`u(x)`$ and $`v(x),`$ in (67) we integrated by parts with respect to the auxiliary variable $`\lambda .`$ Any CAS can do that! We now turn to inverting the $`\mathrm{Div}`$ operator using the homotopy operator. Definition: We define the homotopy operator in 2D (with variables $`x,y`$) through its two components $`(_{𝐮(x,y)}^{(x)}(f),`$$`_{𝐮(x,y)}^{(y)}(f)).`$ The $`x`$-component of the operator is given by
$$_{𝐮(x,y)}^{(x)}(f)=_0^1\underset{j=1}{\overset{N}{}}I_{u_j}^{(x)}(f)[\lambda 𝐮]\frac{d\lambda }{\lambda },$$
(68)
with
$$I_{u_j}^{(x)}(f)=\underset{i_x=0}{\overset{\mathrm{}}{}}\underset{i_y=0}{\overset{\mathrm{}}{}}\left(\frac{1+i_x}{1+i_x+i_y}\right)\mathrm{D}_x^{i_x}\mathrm{D}_y^{i_y}\left(u_j_{u_j(x,y)}^{(1+i_x,i_y)}(f)\right).$$
(69)
Likewise, the $`y`$-component is given by
$$_{𝐮(x,y)}^{(y)}(f)=_0^1\underset{j=1}{\overset{N}{}}I_{u_j}^{(y)}(f)[\lambda 𝐮]\frac{d\lambda }{\lambda },$$
(70)
with
$$I_{u_j}^{(y)}(f)=\underset{i_x=0}{\overset{\mathrm{}}{}}\underset{i_y=0}{\overset{\mathrm{}}{}}\left(\frac{1+i_y}{1+i_x+i_y}\right)\mathrm{D}_x^{i_x}\mathrm{D}_y^{i_y}\left(u_j_{u_j(x,y)}^{(i_x,1+i_y)}(f)\right).$$
(71)
Integrands (69) and (71) involve the 2D higher Euler operators in (56). The question how to compute $`𝐅=(F_1,F_2)=\mathrm{Div}^1(f)`$ is then answered by the following theorem. Theorem: If $`f`$ is a divergence, then $`𝐅=(F_1,F_2)=\mathrm{Div}^1(f)=(_{𝐮(x,y)}^{(x)}(f),_{𝐮(x,y)}^{(y)}(f)).`$ The superscripts $`(x)`$ and $`(y)`$ remind us which components of $`𝐅`$ we are computing. As a matter of testing, we can start from some vector $`\stackrel{~}{𝐅}`$ and compute $`f=\mathrm{Div}\stackrel{~}{𝐅}.`$ Next, we compute $`𝐅=(F_1,F_2)=(_{𝐮(x,y)}^{(x)}(f),_{𝐮(x,y)}^{(y)}(f))`$ and, finally, verify that $`𝐊=\stackrel{~}{𝐅}𝐅`$ is divergence free. Example: Using (46), we show how application of the 2D homotopy operator leads to (47), up to a divergence free vector. Consider $`f=u_xv_yu_{2x}v_yu_yv_x+u_{xy}v_x`$, which is the divergence of $`𝐅`$ in (47). In order to compute $`\mathrm{Div}^1(f)`$, we use (69) for the $`u`$ component in $`𝐮(x,y)=(u(x,y),v(x,y)):`$
$`I_u^{(x)}(f)`$ $`=`$ $`u_{u(x,y)}^{(1,0)}(f)+\mathrm{D}_x\left(u_{u(x,y)}^{(2,0)}(f)\right)+{\displaystyle \frac{1}{2}}\mathrm{D}_y\left(u_{u(x,y)}^{(1,1)}(f)\right)`$ (72)
$`=`$ $`u\left({\displaystyle \frac{f}{u_x}}2\mathrm{D}_x{\displaystyle \frac{f}{u_{2x}}}\mathrm{D}_y{\displaystyle \frac{f}{u_{xy}}}\right)+\mathrm{D}_x\left(u{\displaystyle \frac{f}{u_{2x}}}\right)+{\displaystyle \frac{1}{2}}\mathrm{D}_y\left(u{\displaystyle \frac{f}{u_{xy}}}\right)`$
$`=`$ $`uv_y+{\displaystyle \frac{1}{2}}u_yv_xu_xv_y+{\displaystyle \frac{1}{2}}uv_{xy}.`$
Similarly, for the $`v`$ component of $`𝐮(x,y)=(u(x,y),v(x,y))`$ we get
$$I_v^{(x)}(f)=v_{v(x,y)}^{(1,0)}(f)=v\frac{f}{v_x}=u_yv+u_{xy}v.$$
(73)
Hence, using (68),
$`F_1`$ $`=`$ $`_{𝐮(x,y)}^{(x)}(f)={\displaystyle _0^1}\left(I_u^{(x)}(f)+I_v^{(x)}(f)\right)[\lambda 𝐮]{\displaystyle \frac{d\lambda }{\lambda }}`$ (74)
$`=`$ $`{\displaystyle _0^1}\lambda \left(uv_y+{\displaystyle \frac{1}{2}}u_yv_xu_xv_y+{\displaystyle \frac{1}{2}}uv_{xy}u_yv+u_{xy}v\right)𝑑\lambda `$
$`=`$ $`{\displaystyle \frac{1}{2}}uv_y+{\displaystyle \frac{1}{4}}u_yv_x{\displaystyle \frac{1}{2}}u_xv_y+{\displaystyle \frac{1}{4}}uv_{xy}{\displaystyle \frac{1}{2}}u_yv+{\displaystyle \frac{1}{2}}u_{xy}v.`$
Without showing the details, using (70) and (71) we compute analogously
$`F_2`$ $`=`$ $`_{𝐮(x,y)}^{(y)}(f)={\displaystyle _0^1}\left(I_u^{(y)}(f)+I_v^{(y)}(f)\right)[\lambda 𝐮]{\displaystyle \frac{d\lambda }{\lambda }}`$ (75)
$`=`$ $`{\displaystyle _0^1}\lambda \left(uv_x{\displaystyle \frac{1}{2}}uv_{2x}+{\displaystyle \frac{1}{2}}u_xv_x+u_xvu_{2x}v\right)𝑑\lambda `$
$`=`$ $`{\displaystyle \frac{1}{2}}uv_x{\displaystyle \frac{1}{4}}uv_{2x}+{\displaystyle \frac{1}{4}}u_xv_x+{\displaystyle \frac{1}{2}}u_xv{\displaystyle \frac{1}{2}}u_{2x}v.`$
Now we can readily verify that the resulting vector
$$𝐅=\left(\begin{array}{c}F_1\\ F_2\end{array}\right)=(\begin{array}{c}\frac{1}{2}uv_y+\frac{1}{4}u_yv_x\frac{1}{2}u_xv_y+\frac{1}{4}uv_{xy}\frac{1}{2}u_yv+\frac{1}{2}u_{xy}v\\ \frac{1}{2}uv_x\frac{1}{4}uv_{2x}+\frac{1}{4}u_xv_x+\frac{1}{2}u_xv\frac{1}{2}u_{2x}v\end{array})$$
(76)
differs from $`\stackrel{~}{𝐅}=(uv_yu_xv_y,uv_x+u_xv_x)`$ by a divergence-free vector
$$𝐊=\stackrel{~}{𝐅}𝐅=\left(\begin{array}{c}K_1\\ K_2\end{array}\right)=(\begin{array}{c}\frac{1}{2}uv_y\frac{1}{4}u_yv_x\frac{1}{2}u_xv_y\frac{1}{4}uv_{xy}=\frac{1}{2}u_yv\frac{1}{2}u_{xy}v\\ \frac{1}{2}uv_x+\frac{1}{4}uv_{2x}+\frac{3}{4}u_xv_x\frac{1}{2}u_xv+\frac{1}{2}u_{2x}v\end{array}).$$
(77)
As mentioned in Section 6.1, $`𝐊`$ can be written as $`(\mathrm{D}_y\varphi ,\mathrm{D}_x\varphi )`$ with $`\varphi =\frac{1}{2}uv\frac{1}{4}uv_x\frac{1}{2}u_xv.`$
## 7 Computation of Conservation Laws
In this section we apply the Euler and homotopy operators to compute density-flux pairs for the four PDEs in Section 3. Using the KdV equation (3) we illustrate the steps for a 1D case (space variable $`x`$) and dependent variable $`u(x,t).`$ Part of the computations for this example were done in the previous sections. The Boussinesq equation (4), also in 1D, is used to illustrate the case with two dependent variables $`u(x,t)`$ and $`v(x,t)`$ and an auxiliary parameter with weight.
The LL equation (3), still in 1D, has three dependent variables $`u(x,t),v(x,t)`$ and $`w(x,t).`$ The three coupling constants have weights which makes the computations more cumbersome. However, the multi-uniformity of (3) helps to reduce the complexity. The shallow-water wave equations (3) illustrate the computations in 2D (two space variables) and four dependent variables. Again, we can take advantage of the multi-uniformity of (3) to control “expression swell” of the computations.
We use a direct approach to compute conservation laws, $`\mathrm{D}_t\rho +\mathrm{Div}𝐉=0,`$ of polynomial systems of nonlinear PDEs. First, we build the candidate density $`\rho `$ as a linear combination (with constant coefficients $`c_i)`$ of terms that are uniform in rank with respect to the scaling symmetry of the PDE. In doing so, we dynamically remove divergences and divergence-equivalent terms to get the shortest possible candidate density.
Second, we evaluate $`\mathrm{D}_t\rho `$ on solutions of the PDE, thus removing all time derivatives from the problem. The resulting expression, $`E=\mathrm{D}_t\rho ,`$ must be a divergence of the as yet unknown flux. Thus, we compute $`_{𝐮(𝐱)}^{(\mathrm{𝟎})}(E)`$ and set the coefficients of like terms to zero. This leads to a linear system for the undetermined coefficients $`c_i.`$ If the given PDE has arbitrary constant parameters, then the linear system is parametrized. Careful analysis of the eliminant is needed to find all solution branches, and, when applicable, the conditions on the parameters. For each branch, the solution of the linear system is substituted into $`\rho `$ and $`E.`$
Third, we use the homotopy operator $`_{𝐮(𝐱)}`$ to compute $`𝐉=\mathrm{Div}^1(E).`$ The computations are done with our Mathematica packages . Recall that $`𝐉`$ is only defined up a curl term. Removing the curl term in $`𝐉`$ may lead to a shorter flux. Inversion of $`\mathrm{Div}`$ via the homotopy operator therefore does not guarantee the shortest flux.
### 7.1 Conservation Laws for the KdV Equation
In (5) we gave the first three density-flux pairs of (3). As an example, we will compute $`\rho ^{(3)}`$ and $`J^{(3)}.`$ The weights are $`W(/x)=1`$ and $`W(u)=2.`$ Hence, density $`\rho ^{(3)}`$ has rank 6 and flux $`J^{(3)}`$ has rank 8. The algorithm has three steps: Step 1: Construct the form of the density Start from $`𝒱=\{u\},`$ i.e. the list of dependent variables (and parameters with weight, if applicable). Construct the set $``$ of all non-constant monomials of (selected) rank $`6`$ or less (without derivatives). Thus, $`=\{u^3,u^2,u\}`$. Next, for each monomial in $`,`$ introduce the needed $`x`$-derivatives so that each term has rank 6. Since $`W(/x)=1,`$ use
$$\frac{^2u^2}{x^2}=(2uu_x)_x=2u_x^2+2uu_{2x},\frac{^4u}{x^4}=u_{4x}.$$
(78)
Ignore the highest-order terms (typically the last terms) in each of the right hand sides of (78). Augment $``$ with the remaining terms (deleting numerical factors) to get $`=\{u^3,u_x^2\}.`$
Here $`𝒮=`$ since $``$ is free of divergences and divergent-equivalent terms. Linearly combine the terms in $`𝒮`$ with constant coefficients to get the candidate density:
$$\rho =c_1u^3+c_2u_x^2,$$
(79)
which is (27). Step 2: Compute the undetermined constants $`c_i`$ Compute $`\mathrm{D}_t\rho `$ and use (3) to eliminate $`u_t`$ and $`u_{tx}.`$ As shown in (28), this gives
$$E=\mathrm{D}_t\rho =3c_1u^3u_x+3c_1u^2u_{3x}+2c_2u_x^3+2c_2uu_xu_{2x}+2c_2u_xu_{4x}.$$
(80)
Since $`E=\mathrm{D}_t\rho =\mathrm{D}_xJ,`$ the expression $`E`$ must be exact. As shown in (51), $`_{u(x)}^{(0)}(E)=6(3c_1+c_2)u_xu_{2x}0`$ leads to $`c_1=\frac{1}{3},c_2=1`$ and upon substitution into (79) to
$$\rho =\frac{1}{3}u^3u_x^2$$
(81)
which is $`\rho ^{(3)}`$ in (5). Step 3: Compute the flux $`J`$ Substitute $`c_1=\frac{1}{3},c_2=1`$ into (80) to get
$$E=u^3u_x2u_x^32uu_xu_{2x}+u^2u_{3x}2u_xu_{4x}.$$
(82)
As shown in (61), application of (58) with (59) to (82) gives
$$J=\frac{1}{4}u^42uu_x^2+u^2u_{2x}+u_{2x}^22u_xu_{3x},$$
(83)
which is $`J^{(3)}`$ in (5).
### 7.2 Conservation Laws for the Boussinesq Equation
The first few density-flux pairs of (4) were given in (5). Equation (4) has weights $`W(\frac{}{x})=1`$, $`W(u)=W(\beta )=2,`$ and $`W(v)=3.`$ We show the computation of $`\rho ^{(4)}`$ and $`J^{(4)}`$ of ranks 6 and 7, respectively. The presence of the auxiliary parameter $`\beta `$ with weight complicates matters. Step 1: Construct the form of the density Augment the list of dependent variables with the parameter $`\beta `$ (with non-zero weight). Hence, $`𝒱=\{u,v,\beta \}.`$ Construct $`=\{\beta ^2u,\beta u^2,\beta u,\beta v,u^3,u^2,u,v^2,v,uv\},`$ which contains all non-constant monomials of (chosen) rank $`6`$ or less (without derivatives). Next, for each term in $`,`$ introduce the right number of $`x`$-derivatives so that each term has rank 6. For example,
$$\frac{^2\beta u}{x^2}=\beta u_{2x},\frac{^2u^2}{x^2}=2u_x^2+2uu_{2x},\frac{^4u}{x^4}=u_{4x},\frac{(uv)}{x}=u_xv+uv_x,etc..$$
(84)
Ignore the highest-order terms (typically the last terms) in each of the right hand sides of (84). Augment $``$ with the remaining terms, after deleting numerical factors, to get $`=\{\beta ^2u,\beta u^2,u^3,v^2,u_xv,u_x^2,\beta v_x,uv_x,\beta u_{2x},uu_{2x},v_{3x},u_{4x}\}`$ as in (52). As shown in (53), removal of divergences and divergence-free terms in $``$ leads to $`𝒮=\{\beta ^2u,\beta u^2,u^3,v^2,u_xv,u_x^2\}.`$ Linearly combine the terms in $`𝒮`$ to get
$$\rho =c_1\beta ^2u+c_2\beta u^2+c_3u^3+c_4v^2+c_5u_xv+c_6u_x^2.$$
(85)
Step 2: Compute the undetermined constants $`c_i`$ Compute $`E=\mathrm{D}_t\rho `$ and use (4) to eliminate $`u_t,u_{tx},`$ and $`v_t.`$ As shown in (32), this gives
$$E=(c_1\beta ^2+2c_2\beta u+3c_3u^2)v_x+(c_5v+2c_6u_x)v_{2x}+(2c_4v+c_5u_x)(\beta u_x3uu_x\alpha u_{3x}).$$
(86)
$`E=\mathrm{D}_t\rho =\mathrm{D}_xJ`$ must be exact. Thus, apply (48) and require that $`_{u(x)}^{(0)}(E)=_{v(x)}^{(0)}(E)0.`$ Group like terms. Set their coefficients equal to zero to obtain the parametrized system
$$\beta (c_2c_4)=0,c_3+c_4=0,c_5=0,\alpha c_5=0,\beta c_5=0,\alpha c_4c_6=0.$$
(87)
Investigate the eliminant of the system. Set $`c_1=1,`$ to obtain the solution
$$c_1=1,c_2=c_4,c_3=c_4,c_5=0,c_6=\alpha c_4,$$
(88)
which is valid irrespective of the values of $`\alpha `$ and $`\beta .`$ Substitute (88) into (85) to get
$$\rho =\beta ^2u+c_4(\beta u^2u^3+v^2+\alpha u_x^2).$$
(89)
The density must be split into independent pieces. Indeed, since $`c_4`$ is arbitrary, set $`c_4=0`$ or $`c_4=1,`$ thus splitting (89) into two independent densities, $`\rho =\beta ^2u`$ and
$$\rho =\beta u^2u^3+v^2+\alpha u_x^2,$$
(90)
which are $`\rho ^{(1)}`$ and $`\rho ^{(4)}`$ in (5). Step 3: Compute the flux $`J`$ Compute the flux corresponding to $`\rho `$ in (90). Substitute (88) into (86). Take the terms in $`c_4`$ and set $`c_4=1.`$ Thus,
$$E=2\beta uv_x+2\beta u_xv3u^2v_x6uu_xv+2\alpha u_xv_{2x}2\alpha u_{3x}v.$$
(91)
Apply (62) and (63) to (91). For $`E`$ of order $`3`$ in $`u(x),`$ compute
$`I_u(E)`$ $`=`$ $`u_{u(x)}^{(1)}(E)+\mathrm{D}_x\left(u_{u(x)}^{(2)}(E)\right)+\mathrm{D}_x^2\left(u_{u(x)}^{(3)}(E)\right)`$ (92)
$`=`$ $`u{\displaystyle \frac{E}{u_x}}2u\mathrm{D}_x({\displaystyle \frac{E}{u_{2x}}})+3u\mathrm{D}_x^2({\displaystyle \frac{E}{u_{3x}}})+\mathrm{D}_x(u{\displaystyle \frac{E}{u_{2x}}}3u\mathrm{D}_x({\displaystyle \frac{E}{u_{3x}}}))+\mathrm{D}_x^2(u{\displaystyle \frac{E}{u_{3x}}})`$
$`=`$ $`2\beta uv6u^2v4\alpha uv_{2x}+6\alpha \mathrm{D}_x\left(uv_x\right)2\alpha \mathrm{D}_x^2\left(uv\right)`$
$`=`$ $`2\beta uv6u^2v2\alpha u_{2x}v+2\alpha u_xv_x.`$
With $`E`$ is of order $`2`$ in $`v(x),`$ subsequently compute
$`I_v(E)`$ $`=`$ $`v_{v(x)}^{(1)}(E)+\mathrm{D}_x\left(v_{v(x)}^{(2)}(E)\right)=v{\displaystyle \frac{E}{v_x}}2v\mathrm{D}_x({\displaystyle \frac{E}{v_{2x}}})+\mathrm{D}_x(v{\displaystyle \frac{E}{v_{2x}}})`$ (93)
$`=`$ $`2\beta uv3u^2v4\alpha uv_{2x}+2\alpha \mathrm{D}_x\left(u_xv\right)=2\beta uv3u^2v2\alpha u_{2x}v+2\alpha u_xv_x.`$
Formula (62) with $`𝐮(x)=(u(x),v(x))`$ requires an integration with respect to $`\lambda .`$ Hence,
$`J`$ $`=`$ $`_{u(x)}(E)={\displaystyle _0^1}(I_u(E)+I_v(E))[\lambda 𝐮]{\displaystyle \frac{d\lambda }{\lambda }}`$ (94)
$`=`$ $`{\displaystyle _0^1}\left(4\beta \lambda uv9\lambda ^2u^2v4\alpha \lambda u_{2x}v+4\alpha \lambda u_xv_x\right)𝑑\lambda `$
$`=`$ $`2\beta uv3u^2v2\alpha u_{2x}v+2\alpha u_xv_x,`$
which is $`J^{(4)}`$ in (5). Set $`\beta =1`$ at the end of the computations.
### 7.3 Conservation Laws for the Landau-Lifshitz Equation
Without showing full details, we will compute the six density-flux pairs (5) for (3). The weights for (3) are $`W(u)=W(v)=W(w)=a,W(\alpha )=W(\beta )=W(\gamma )=2,W(\frac{}{t})=a+2,`$ where $`a`$ is an arbitrary non-negative integer or rational number. This example involves three constants $`(\alpha ,\beta ,`$ and $`\gamma )`$ with weights which makes the computations quite cumbersome unless we take advantage of the multi-uniformity and the cyclic nature of (3). Step 1: Construct the form of the density Select a small value for $`a`$, for example $`a=\frac{1}{4},`$ and compute a density of, say, rank $`\frac{1}{4}.`$ Start from $`𝒱=\{u,v,w,\alpha ,\beta ,\gamma \},`$ i.e. the list of dependent variables and the parameters with weight. Construct the set $``$ of all non-constant monomials of the (selected) rank $`\frac{1}{4}`$ or less (without derivatives). Thus, $`=\{u,v,w\}.`$ Next, for each of the monomials in $`,`$ introduce the appropriate number of $`x`$-derivatives so that each term has rank $`\frac{1}{4}.`$ No $`x`$-derivatives are needed and the removal of divergences or divergence-equivalent terms is irrelevant. Linearly combine the terms in $`=𝒮=\{u,v,w\}`$ to get a candidate density:
$$\rho =c_1u+c_2v+c_3w.$$
(95)
It suffices to continue with $`\rho =c_1u.`$ The remaining terms follow by cyclic permutation. Step 2: Compute the undetermined constants $`c_i`$ Compute $`E=\mathrm{D}_t\rho =\mathrm{D}_t(c_1u)=c_1u_t.`$ Use (3) to eliminate $`u_t.`$ Hence,
$$E=c_1\left(vw_{2x}wv_{2x}+(\gamma \beta )vw\right),$$
(96)
which is the opposite of the right hand side of the first equation in (3). Since $`E=\mathrm{D}_t\rho =\mathrm{D}_xJ,`$ the expression $`E`$ must be exact. Obviously $`_{u(x)}^{(0)}(E)0.`$ Compute
$`_{v(x)}^{(0)}(E)`$ $`=`$ $`{\displaystyle \frac{f}{v}}(E)\mathrm{D}_x{\displaystyle \frac{f}{v_x}}(E)+\mathrm{D}_x^2{\displaystyle \frac{f}{v_{2x}}}(E)=c_1(\gamma \beta )w,`$
$`_{w(x)}^{(0)}(E)`$ $`=`$ $`{\displaystyle \frac{f}{w}}(E)\mathrm{D}_x{\displaystyle \frac{f}{w_x}}(E)+\mathrm{D}_x^2{\displaystyle \frac{f}{w_{2x}}}(E)=c_1(\gamma \beta )v.`$ (97)
Set the latter expressions identically equal to zero. Solve $`c_1(\beta \gamma )=0`$ for $`c_10.`$ Set $`c_1=1`$ to get $`\rho =u,`$ subject to the condition $`\beta =\gamma ,`$ which confirms the result in (5). Step 3: Compute the flux $`J`$ Substitute $`c_1=1`$ and $`\beta =\gamma `$ into (96) to get $`E=wv_{2x}vw_{2x}.`$ Apply the homotopy operator (58) with (59) to $`E.`$ In this example $`𝐮(x)=(u(x),v(x),w(x)).`$ Obviously, $`I_u(E)=0,`$ since there is no explicit occurrence of $`u.`$ Compute
$$I_v(E)=v_{v(x)}^{(1)}(E)+\mathrm{D}_x\left(v_{v(x)}^{(2)}(E)\right)=v\frac{E}{v_x}2v\mathrm{D}_x\left(\frac{E}{v_{2x}}\right)+\mathrm{D}_x\left(v\frac{E}{v_{2x}}\right)=v_xwvw_x,$$
(98)
and, analogously, $`I_w(E)=v_xwvw_x.`$ Finally, compute
$`J`$ $`=`$ $`_{u(x)}(E)={\displaystyle _0^1}(I_u(E)+I_v(E)+I_w(E))[\lambda 𝐮]{\displaystyle \frac{d\lambda }{\lambda }}`$ (99)
$`=`$ $`{\displaystyle _0^1}2\lambda \left(v_xwvw_x\right)𝑑\lambda =v_xwvw_x,`$
which is $`J^{(1)}`$ in (5). To compute density $`\rho ^{(4)}`$ which is quadratic in $`u,v,`$ and $`w,`$ we would start from rank $`\frac{1}{2}`$ and repeat the three steps. Step 1 would lead to $`\rho =c_1u^2+c_2uv+c_3v^2+c_4uw+c_5vw+c_6w^2.`$ Steps 2 and 3 would result in $`c_1=c_3=c_6=1.`$ So, $`\rho =u^2+v^2+w^2`$ and $`J=0,`$ which agrees with $`\rho ^{(4)}`$ and $`J^{(4)}=0`$ in (5). Starting with rank $`1,`$ Step 1 would generate a polynomial density with the 15 terms that are homogeneous of degree 4. Steps 2 and 3 would result in $`\rho ^{(5)}`$ and $`J^{(5)}=0`$ in (5). The computation of $`\rho ^{(6)}`$ and $`J^{(6)}`$ is cumbersome and long. We only indicate the strategy and give partial results. Notice that $`\rho ^{(6)}`$ would have rank $`\frac{5}{2}`$ if $`a=\frac{1}{4}`$ and the candidate density would have, amongst others, all homogeneous terms of degree nine, which is undesirable. It turns out that $`a=1`$ is a better choice to compute $`\rho ^{(6)},`$ which then has rank 4.
Start from $`𝒱`$ as above. Construct the set $`=\{u^4,u^3,u^2,u,v^4,v^3,\mathrm{},`$ $`u^3v,u^3w,\mathrm{},`$ $`u^2v^2,u^2w^2,\mathrm{},u^2vw,v^2uw,\mathrm{},`$ $`\alpha u^2,\beta u^2,`$ $`\mathrm{},\alpha uv,\mathrm{},\gamma vw,\mathrm{},uvw,uv,\mathrm{},\gamma w\}`$ of all non-constant monomials of rank $`4`$ or less (without derivatives). Next, for each of the 61 monomials in $`,`$ introduce the needed $`x`$-derivatives to make each term rank 4. Use, for example,
$$\frac{^2u^2}{x^2}=2u_x^2+2uu_{2x},\frac{u^3}{x}=3u^2u_x,\frac{^4u}{x^4}=u_{4x},\frac{u^2v}{x}=2uu_xv+u^2v_x,\frac{\alpha u}{x}=\alpha u_x.$$
(100)
Ignore the highest-order terms (typically the last terms). Augment $``$ with the remaining terms (deleting numerical factors). Next, remove all divergences and divergent-equivalent terms to get $`𝒮`$ with 47 terms (not shown). Linearly combine the terms in $`𝒮`$ to get the candidate density:
$`\rho `$ $`=`$ $`c_1\alpha u^2+c_2\beta u^2+c_3\gamma u^2+c_4u^4+c_5\alpha uv+c_6\beta uv+c_7\gamma uv+c_8u^3v+\mathrm{}+c_{12}u^2v^2`$ (101)
$`+c_{13}uv^3+\mathrm{}+c_{22}u^2vw+c_{23}uv^2w+\mathrm{}+c_{25}\alpha w^3+c_{26}\beta w^3+c_{27}\gamma w^3+\mathrm{}+c_{29}uvw^2`$
$`+c_{30}v^2w^2+\mathrm{}+c_{33}w^4+c_{34}uu_xv+c_{35}u_xv^2+c_{36}uu_xw+c_{37}u_xvw+c_{38}u_xw^2+c_{39}u_x^2`$
$`+c_{40}uv_xw+c_{41}vv_xw+c_{42}v_xw^2+c_{43}u_xv_x+c_{44}v_x^2+c_{45}u_xw_x+c_{46}v_xw_x+c_{47}w_x^2,`$
where only the most indicative terms are explicitly shown.
Compute $`\mathrm{D}_t\rho ,`$ use (3) to eliminate $`u_t,v_t,w_t,u_{tx},v_{tx},`$ and $`w_{tx},`$ and require that $`_{u(x)}^{(0)}(E)=_{v(x)}^{(0)}(E)=_{w(x)}^{(0)}(E)0`$ which leads to a linear system of 121 equations (not shown) for $`c_1`$ through $`c_{47}.`$ Substitute the solution (not shown) into (101) to obtain
$`\rho `$ $`=`$ $`(\alpha c_{25}+\beta c_{26}+\gamma c_{27})(u^2+v^2+w^2)+{\displaystyle \frac{1}{2}}c_{30}(u^2+v^2+w^2)^2`$ (102)
$`+c_{47}\left(u_x^2+v_x^2+w_x^2+(\gamma \alpha )u^2+(\gamma \beta )v^2\right),`$
where $`c_{25},c_{26},c_{27},c_{30},`$ and $`c_{47}`$ are arbitrary. Split the density into independent pieces:
$`\rho `$ $`=`$ $`u^2+v^2+w^2,\rho =(u^2+v^2+w^2)^2,`$
$`\rho `$ $`=`$ $`u_x^2+v_x^2+w_x^2+(\gamma \alpha )u^2+(\gamma \beta )v^2,`$ (103)
which are $`\rho ^{(4)},\rho ^{(5)},`$ and $`\rho ^{(6)}`$ in (5). Use the homotopy operator to compute the flux corresponding to (102):
$`J`$ $`=`$ $`2c_{47}((vw_xwv_x)u_{2x}+(wu_xuw_x)v_{2x}+(uv_xvu_x)w_{2x}.`$ (104)
$`.+(\beta \gamma )vwu_x+(\gamma \alpha )uwv_x+(\alpha \beta )uvw_x).`$
$`J^{(4)}=J^{(5)}=0`$ since the terms in $`c_{25},c_{26},c_{27},`$ and $`c_{30}`$ all dropped out. Set $`c_{47}=1`$ to get $`J^{(6)}`$ in (5).
### 7.4 Conservation Laws for the Shallow Water Wave Equations
The SWW equations (3) admit weights (19) in which $`W(h)`$ and $`W(\mathrm{\Omega })`$ are free. The fact that (3) is multi-uniform is advantageous. Indeed, we can construct a candidate $`\rho `$ which is uniform in rank for one set of weights and, subsequently, use other choices of weights to split $`\rho `$ into smaller pieces. This “divide and conquer” strategy drastically reduces the complexity of the computations as was shown in .
The first few densities and fluxes were given in (41). We compute $`\rho ^{(5)}`$ and $`J^{(5)}.`$ Note that $`\rho ^{(5)}`$ has rank 3 if we select $`W(h)=a=1`$ and $`W(\mathrm{\Omega })=b=2.`$ However, $`\rho ^{(5)}`$ has rank 4 if we take $`W(h)=a=0`$ and $`W(\mathrm{\Omega })=b=2.`$ If we set $`W(h)=0,`$ and $`W(\mathrm{\Omega })=3`$ then $`\rho ^{(5)}`$ has rank 7. So, the trick is to construct a candidate density which is scaling homogeneous for a particular (fixed) choice of $`a`$ and $`b`$ in (20) and split the density based on other choices of $`a`$ and $`b.`$ Step 1: Construct the form of the density Start from $`𝒱=\{u,v,\theta ,h,\mathrm{\Omega }\},`$ i.e. the list of variables and parameters with weights. Use (20) with $`a=1,b=2,`$ to get $`=\{\mathrm{\Omega }u,\mathrm{\Omega }v,\mathrm{},u^3,v^3,\mathrm{},u^2v,uv^2,\mathrm{},u^2,v^2,\mathrm{},u,v,\theta ,h\}`$ which has 38 monomials of (chosen) rank 3 or less (without derivatives).
All terms of rank 3 in $``$ remain untouched. To adjust the rank, differentiate each monomial of rank 2 in $``$ with respect to $`x`$ ignoring the highest-order term. For example, in $`\frac{du^2}{dx}=2uu_x,`$ the term can be ignored since it is a total derivative. The terms $`u_xv`$ and $`uv_x`$ are divergence-equivalent since $`\frac{d(uv)}{dx}=u_xv+uv_x.`$ Keep $`u_xv.`$ Likewise, differentiate each monomial of rank 2 in $``$ with respect to $`y`$ and ignore the highest-order term.
Produce the remaining terms for rank 3 by differentiating the monomials of rank 1 in $``$ with respect to $`x`$ twice, or $`y`$ twice, or once with respect to $`x`$ and $`y.`$ Again ignore the highest-order terms. Augment the set $``$ with the derivative terms of rank 3 to get $`=\{\mathrm{\Omega }u,\mathrm{\Omega }v,\mathrm{},uv^2,u_xv,u_x\theta ,u_xh,\mathrm{},u_yv,u_y\theta ,\mathrm{},\theta _yh\}`$ which has 36 terms.
Use the “divide and conquer” strategy to select from $``$ the terms which are of ranks 4 and 7 using the choices $`a=0,b=2`$ and $`a=0,b=3,`$ respectively. This gives $`𝒮=\{\mathrm{\Omega }\theta ,u_x\theta ,u_y\theta ,v_x\theta ,v_y\theta \},`$ which happens to be free of divergences and divergence-equivalent terms. So, no further reduction is needed. Linearly combine the terms in $`𝒮`$ to get
$$\rho =c_1\mathrm{\Omega }\theta +c_2u_x\theta +c_3u_y\theta +c_4v_x\theta +c_5v_y\theta .$$
(105)
Step 2: Compute the undetermined constants $`c_i`$ Compute $`E=\mathrm{D}_t\rho `$ and use (3) to remove all time derivatives:
$`E`$ $`=`$ $`\left({\displaystyle \frac{\rho }{u_x}}u_{tx}+{\displaystyle \frac{\rho }{u_y}}u_{ty}+{\displaystyle \frac{\rho }{v_x}}v_{tx}+{\displaystyle \frac{\rho }{v_y}}v_{ty}+{\displaystyle \frac{\rho }{\theta }}\theta _t\right)`$ (106)
$`=`$ $`c_2\theta (uu_x+vu_y2\mathrm{\Omega }v+\frac{1}{2}h\theta _x+\theta h_x)_x+c_3\theta (uu_x+vu_y2\mathrm{\Omega }v+\frac{1}{2}h\theta _x+\theta h_x)_y`$
$`+c_4\theta (uv_x+vv_y+2\mathrm{\Omega }u+\frac{1}{2}h\theta _y+\theta h_y)_x+c_5\theta (uv_x+vv_y+2\mathrm{\Omega }u+\frac{1}{2}h\theta _y+\theta h_y)_y`$
$`+(c_1\mathrm{\Omega }+c_2u_x+c_3u_y+c_4v_x+c_5v_y)(u\theta _x+v\theta _y).`$
Require that $`_{u(x,y)}^{(0,0)}(E)=_{v(x,y)}^{(0,0)}(E)=_{\theta (x,y)}^{(0,0)}(E)=_{h(x,y)}^{(0,0)}(E)0,`$ where, for example, $`_{u(x,y)}^{(0,0)}`$ is given in (50). Gather like terms. Equate their coefficients to zero to obtain
$$c_1+2c_3=0,c_2=c_5=0,c_12c_4=0,c_3+c_4=0.$$
(107)
Set $`c_1=2.`$ Substitute the solution
$$c_1=2,c_2=0,c_3=1,c_4=1,c_5=0.$$
(108)
into $`\rho `$ to obtain
$$\rho =2\mathrm{\Omega }\theta u_y\theta +v_x\theta ,$$
(109)
which is $`\rho ^{(5)}`$ in (41). Step 3: Compute the flux $`𝐉`$ Compute the flux corresponding to (109). To do so, substitute (108) into (106) to obtain
$`E`$ $`=`$ $`\theta (u_xv_xu_yv_y+uv_{2x}u_{2y}vu_xu_y+v_xv_yuu_{xy}+vv_{xy}`$ (110)
$`+2\mathrm{\Omega }u_x+2\mathrm{\Omega }v_y+\frac{1}{2}\theta _xh_y\frac{1}{2}\theta _yh_x)+2\mathrm{\Omega }u\theta _x+2\mathrm{\Omega }v\theta _yuu_y\theta _x`$
$`u_yv\theta _y+uv_x\theta _x+vv_x\theta _y.`$
Apply the 2D homotopy formulas in (68)-(71) to $`E=\mathrm{Div}𝐉=\mathrm{D}_xJ_1+\mathrm{D}_yJ_2.`$ So, compute
$`I_u^{(x)}(E)`$ $`=`$ $`u_{u(x,y)}^{(1,0)}(E)+\mathrm{D}_x\left(u_{u(x,y)}^{(2,0)}(E)\right)+{\displaystyle \frac{1}{2}}\mathrm{D}_y\left(u_{u(x,y)}^{(1,1)}(E)\right)`$ (111)
$`=`$ $`u\left({\displaystyle \frac{E}{u_x}}2\mathrm{D}_x\left({\displaystyle \frac{E}{u_{2x}}}\right)\mathrm{D}_y\left({\displaystyle \frac{E}{u_{xy}}}\right)\right)+\mathrm{D}_x(u{\displaystyle \frac{E}{u_{2x}}})+{\displaystyle \frac{1}{2}}\mathrm{D}_y(u{\displaystyle \frac{E}{u_{xy}}})`$
$`=`$ $`uv_x\theta +2\mathrm{\Omega }u\theta +{\displaystyle \frac{1}{2}}u^2\theta _yuu_y\theta .`$
Similarly, compute
$`I_v^{(x)}(E)`$ $`=`$ $`vv_y\theta +{\displaystyle \frac{1}{2}}v^2\theta _y+uv_x\theta ,`$ (112)
$`I_\theta ^{(x)}(E)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\theta ^2h_y+2\mathrm{\Omega }u\theta uu_y\theta +uv_x\theta ,`$ (113)
$`I_h^{(x)}(E)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\theta \theta _yh.`$ (114)
Next, compute
$`J_1(𝐮)`$ $`=`$ $`_{𝐮(x,y)}^{(x)}(E)`$ (115)
$`=`$ $`{\displaystyle _0^1}\left(I_u^{(x)}(E)+I_v^{(x)}(E)+I_\theta ^{(x)}(E)+I_h^{(x)}(E)\right)[\lambda 𝐮]{\displaystyle \frac{d\lambda }{\lambda }}`$
$`=`$ $`{\displaystyle _0^1}(4\lambda \mathrm{\Omega }u\theta +\lambda ^2(3uv_x\theta +{\displaystyle \frac{1}{2}}u^2\theta _y2uu_y\theta +vv_y\theta +{\displaystyle \frac{1}{2}}v^2\theta _y`$
$`+{\displaystyle \frac{1}{2}}\theta ^2h_y{\displaystyle \frac{1}{2}}\theta \theta _yh))d\lambda `$
$`=`$ $`2\mathrm{\Omega }u\theta {\displaystyle \frac{2}{3}}uu_y\theta +uv_x\theta +{\displaystyle \frac{1}{3}}vv_y\theta +{\displaystyle \frac{1}{6}}u^2\theta _y+{\displaystyle \frac{1}{6}}v^2\theta _y{\displaystyle \frac{1}{6}}h\theta \theta _y+{\displaystyle \frac{1}{6}}h_y\theta ^2.`$
Likewise, compute $`I_u^{(y)}(E),I_v^{(y)}(E),I_\theta ^{(y)}(E),`$ and $`I_h^{(y)}(E).`$ Finally, compute
$`J_2(𝐮)`$ $`=`$ $`_{𝐮(x,y)}^{(y)}(E)`$ (116)
$`=`$ $`{\displaystyle _0^1}\left(I_u^{(y)}(E)+I_v^{(y)}(E)+I_\theta ^{(y)}(E)+I_h^{(y)}(E)\right)[\lambda 𝐮]{\displaystyle \frac{d\lambda }{\lambda }}`$
$`=`$ $`2\mathrm{\Omega }v\theta +{\displaystyle \frac{2}{3}}vv_x\theta vu_y\theta {\displaystyle \frac{1}{3}}uu_x\theta {\displaystyle \frac{1}{6}}u^2\theta _x{\displaystyle \frac{1}{6}}v^2\theta _x+{\displaystyle \frac{1}{6}}h\theta \theta _x{\displaystyle \frac{1}{6}}h_x\theta ^2.`$
Hence,
$$𝐉=(\begin{array}{c}J_1\\ J_2\end{array})=\frac{1}{6}\left(\begin{array}{c}12\mathrm{\Omega }u\theta 4uu_y\theta +6uv_x\theta +2vv_y\theta +u^2\theta _y+v^2\theta _yh\theta \theta _y+h_y\theta ^2\\ 12\mathrm{\Omega }v\theta +4vv_x\theta 6vu_y\theta 2uu_x\theta u^2\theta _xv^2\theta _x+h\theta \theta _xh_x\theta ^2\end{array}\right),$$
(117)
which matches $`𝐉^{(5)}`$ in (41).
## 8 Conclusions
The variational derivative (zeroth Euler operator) provides a straightforward way to test exactness which is key in the computation of densities. The continuous homotopy operator is a powerful, algorithmic tool to compute fluxes explicitly. Indeed, the homotopy operator allows one to invert the total divergence operator by computing higher variational derivatives followed by a one-dimensional integration with respect to a single auxiliary parameter. The homotopy operator is a universal tool that can be applied to problems in which integration by parts (of arbitrary functions) in multi-variables is crucial.
To reach a wider audience, we intentionally did not use differential forms and the abstract framework of variational bi-complexes and homological algebra. Instead, we extracted the Euler and homotopy operators from their abstract setting and presented them in the language of standard calculus, thereby making them widely applicable to computational problems in the sciences. Our calculus-based formulas can be readily implemented in CAS.
Based on the concept of scaling invariance and using tools of the calculus of variations, we presented a three-step algorithm to symbolically compute polynomial conserved densities and fluxes of nonlinear polynomial systems of PDEs in multi-spatial dimensions. The steps are straightforward: build candidate densities as linear combinations (with undetermined constant coefficients) of terms that are homogeneous with respect to the scaling symmetry of the PDE. Subsequently, use the variational derivative to compute the coefficients, and, finally, use the homotopy operator to compute the flux.
With our method one can search for conservation laws in chemistry, physics, and engineering. Symbolic packages are well suited to assist in the search which covers many fields of application including fluid mechanics, plasma physics, electro-dynamics, gas dynamics, elasticity, nonlinear optics and acoustics, electrical networks, chemical reactions, etc..
## Acknowledgements and Dedication
This material is based upon work supported by the National Science Foundation (NSF) under Grants Nos. DMS-9732069, DMS-9912293, and CCR-9901929. Any opinions, findings, and conclusions or recommendations expressed in this material are those of the author and do not necessarily reflect the views of NSF.
The author is grateful to Bernard Deconinck, Michael Colagrosso, Mark Hickman and Jan Sanders for valuable discussions. Undergraduate students Lindsay Auble, Robert “Scott” Danford, Ingo Kabirschke, Forrest Lundstrom, Frances Martin, Kara Namanny, Adam Ringler, and Maxine von Eye are thanked for designing Mathematica code for the project.
The research was supported in part by an Undergraduate Research Fund Award from the Center for Engineering Education awarded to Ryan Sayers. On June 16, 2003, while rock climbing in Wyoming, Ryan was killed by a lightning strike. He was 20 years old. The author expresses his gratitude for the insight Ryan brought to the project. His creativity and passion for mathematics were inspiring. This paper is dedicated to him. |
warning/0506/gr-qc0506060.html | ar5iv | text | # Stellar Oscillations in Scalar-Tensor Theory of Gravity
## I Introduction
Scalar-tensor theories of gravity are an alternative or generalization of Einstein’s theory of gravity, where in addition to the tensor field a scalar field is present. The theory has been proposed in its earlier form about half century ago Fierz1956 ; Jordan1959 ; Brans1961 , and it is a viable theory of gravity for a specific range of the coupling strength of the scalar field to gravity Damour1992 ; Will1993 ; Will2001 . Actually, the existence of scalar fields is crucial (e.g. in inflationary and quintessence scenarios) to explain the accelerated expansion phases of the universe. In addition, scalar-tensor theories of gravity can be obtained from the low-energy limit of string theory and/or other gauge theories. Experimentally, the existence of a scalar field has not yet been probed, but a number of experiments in the weak field limit of general relativity set severe limits on the existence and coupling strengths of scalar fields Will2001 ; Esposito2004 .
A basic assumption is that the scalar and gravitational fields $`\phi `$ and $`g_{\mu \nu }`$ are coupled to matter via an “effective metric” $`\stackrel{~}{g}_{\mu \nu }=A^2(\phi )g_{\mu \nu }`$. The Fierz-Jordan-Brans-Dicke Fierz1956 ; Jordan1959 ; Brans1961 theory assumes that the “coupling function” has the form $`A(\phi )=\alpha _0\phi `$, i.e., it is characterized by a unique free parameter $`\alpha _0^2=(2\omega _{\mathrm{BD}}+3)^1`$, and all its predictions differ from those of general relativity by quantities of order $`\alpha _0^2`$ Damour1993 . Solar system experiments set strict limits in the value of the Brans-Dicke parameter $`\omega _{\mathrm{BD}}`$, i.e., $`\omega _{\mathrm{BD}}40000`$, which suggests a very small $`\alpha _0^2<10^5`$ (see Bertotti2003 ; Esposito2004 ).
In the early 1990s, based on a simplified version of scalar tensor theory where $`A(\phi )=\alpha _0\phi +\beta \phi ^2/2`$, Damour and Esposito-Farese Damour1993 ; Damour1996 found that for certain values of the coupling parameter $`\beta `$ the stellar models develop some strong field effects which induce significant deviations from general relativity. This sudden deviation from general relativity for specific values of the coupling constants has been named “spontaneous scalarization”. Harada Harada1998 studied in more detail models of non-rotating neutron stars in the framework of the scalar-tensor theory and he reported that “spontaneous scalarization” is possible for $`\beta 4.35`$. DeDeo and Psaltis suggested that the effects of scalar fields might be apparent in the observed redshifted lines of X-rays and $`\gamma `$-rays observed by Chandra and XMM-Newton DeDeo2003 and in quasi-periodic oscillations (QPOs) DeDeo2004 .
Recently, we have examined the possibility to obtain the information for the presence of the scalar field via gravitational wave observations of oscillating neutron stars (SotaniKokkotas2004 , hereafter Paper I). This previous study has been done using the so-called Cowling approximation. In this approximation one studies the fluid oscillations freezing the perturbations of the spacetime and of the scalar field. Even under these assumptions the effect of the scalar field on the fluid perturbations can be significant. We showed that for values of $`\beta 4.35`$ the oscillation frequencies of the fluid change drastically, and the observation of such oscillations can verify or rule out the spontaneous scalarization phenomenon.
It has been suggested that stellar oscillations can provide a unique tool for estimating the parameters of the star, i.e., mass, radius, rotation rate, magnetic field and equation of state. These ideas have been developed in the last decade in a series of papers Andersson1996 ; Andersson1998 ; Benhar1999 ; Andersson2001a ; Sotani2003 ; Sotani2004 ; Benhar2004 , where the properties of the various families of oscillation modes have been used to probe the stellar parameters. The modes which are mainly excited during the formation of a neutron star or during starquakes and emit detectable gravitational waves are the fluid $`f`$ and $`p`$ modes and the $`w`$ modes, which are associated to oscillations of the spacetime Kokkotas2001 .
The effect of the scalar field on the $`f`$ and $`p`$ modes has been examined in Paper I (in the Cowling approximation). In this article we derive the full set of equations describing the oscillations of a relativistic star, i.e., the perturbations of the fluid, the spacetime, and the scalar field. Since the stellar models are spherically symmetric, the oscillations can be classified as axial or polar depending on their parity, and we can derive perturbation equations for each class of perturbations. In the polar case we show that the wave equations describing the perturbations of the fluid and spacetime couple to the wave equation describing the perturbations of the scalar field. In other words the polar perturbations are affected not only by the presence of the scalar field in the background but also by the coupling with the wave equation describing the perturbations of the scalar field. In the axial case we find a single equation describing the perturbation of the spacetime. This equation is not coupled to perturbations of the fluid and of the scalar field: the scalar field only affects the dynamics through its influence on the background.
The paper is organized as follows. In the next section we establish our notation and briefly introduce the theoretical framework for the scalar-tensor theories of gravity we consider. In Section III we derive the perturbation equations which will be used for the numerical estimation of the oscillation frequencies. In Section IV we describe the methods that have been used to derive the axial $`w`$ modes in scalar-tensor theory of gravity and show the results. In the final Section V we discuss the results and their implications. We have included in three Appendices the details of the various analytic and numerical calculations. In Appendix A we describe the perturbations of the energy momentum tensors for the fluid and the scalar field, while in the next Appendix B we provide the analytic forms of the perturbed Einstein equations. Finally, in the last Appendix C we describe the numerical techniques that have been used to calculated the quasinormal modes. In this paper we adopt the unit of $`c=G=1`$, where $`c`$ and $`G`$ denote the speed of light and the gravitational constant, respectively, and the metric signature of $`(,+,+,+)`$.
## II Stellar models in scalar-tensor theories of gravity
In this section we will study neutron star models in scalar-tensor theory of gravity with one scalar field. This is a natural extensions of Einstein’s theory, in which gravity is mediated not only by a second rank tensor (the metric tensor $`g_{\mu \nu }`$) but also by a massless long-range scalar field $`\phi `$. The action is given by Damour1992
$$S=\frac{1}{16\pi G_{}}\sqrt{g_{}}\left(R_{}2g_{}^{\mu \nu }\phi _{,\mu }\phi _{,\nu }\right)d^4x+S_m[\mathrm{\Psi }_m,A^2(\phi )g_{\mu \nu }],$$
(1)
where all quantities with asterisks are related to the “Einstein metric” $`g_{\mu \nu }`$, then $`R_{}`$ is the curvature scalar for this metric and $`G_{}`$ is the bare gravitational coupling constant. $`\mathrm{\Psi }_m`$ represents collectively all matter fields, and $`S_m`$ denotes the action of the matter represented by $`\mathrm{\Psi }_m`$, which is coupled to the “Jordan-Fierz metric tensor” $`\stackrel{~}{g}_{\mu \nu }`$. The field equations are usually written in terms of the “Einstein metric”, but all non-gravitational experiments measure the “Jordan-Fierz” or “physical metric”. The “Jordan-Fierz metric” is related to the “Einstein metric” via the conformal transformation,
$$\stackrel{~}{g}_{\mu \nu }=A^2(\phi )g_{\mu \nu }.$$
(2)
Hereafter, we denote by a tilde quantities in the “physical frame” and by an asterisk those in the “Einstein frame”. From the variation of the action $`S`$ we get the field equations in the Einstein frame
$`G_{\mu \nu }`$ $`=8\pi G_{}T_{\mu \nu }+T_{\mu \nu }^{(\phi )},`$ (3)
$`\text{ }\text{ }\text{ }\text{ }\text{ }_{}\phi `$ $`=4\pi G_{}\alpha (\phi )T_{},`$ (4)
where $`T_{\mu \nu }^{(\phi )}`$ is the energy-momentum of the massless scalar field, i.e.,
$$T_{\mu \nu }^{(\phi )}2\left(\phi _{,\mu }\phi _{,\nu }\frac{1}{2}g_{\mu \nu }g_{}^{\alpha \beta }\phi _{,\alpha }\phi _{,\beta }\right)$$
(5)
and $`T_{}^{\mu \nu }`$ is the energy-momentum tensor in the Einstein frame which is related to the physical energy-momentum tensor $`\stackrel{~}{T}_{\mu \nu }`$ as follows,
$$T_{}^{\mu \nu }\frac{2}{\sqrt{g_{}}}\frac{\delta S_m}{\delta g_{\mu \nu }}=A^6(\phi )\stackrel{~}{T}^{\mu \nu }.$$
(6)
The scalar quantities $`T_{}`$ and $`\alpha (\phi )`$ are defined as $`T_{}T_\mu ^\mu =T_{}^{\mu \nu }g_{\mu \nu }`$ and $`\alpha (\phi )d\mathrm{ln}A(\phi )/d\phi `$. It is apparent that $`\alpha (\phi )`$ is the only field-dependent function which couples the scalar field with matter, for $`\alpha (\phi )=0`$ the theory reduces to general relativity. Finally, the law of energy-momentum conservation $`\stackrel{~}{}_\nu \stackrel{~}{T}_\mu ^\nu =0`$ is transformed into
$$_\nu T_\mu ^\nu =\alpha (\phi )T_{}_\mu \phi ,$$
(7)
and we set $`\phi _0`$ as the cosmological value of the scalar field at infinity. In this paper, we adopt the same form of conformal factor $`A(\phi )`$ as in Damour and Esposito-Farese Damour1993 , which is
$$A(\phi )=e^{\frac{1}{2}\beta \phi ^2},$$
(8)
i.e., $`\alpha (\phi )=\beta \phi `$ where $`\beta `$ is a real number. In the case $`\beta =0`$ this scalar-tensor theory reduces to general relativity, while “spontaneous scalarization” occurs for $`\beta 4.35`$ Harada1998 .
We will model the neutron stars as self-gravitating perfect fluid of cold degenerate matter in equilibrium. Then the metric describing an unperturbed, non-rotating, spherically symmetric neutron star can be written as
$$ds_{}^2=g_{\mu \nu }dx^\mu dx^\nu =e^{2\mathrm{\Phi }}dt^2+e^{2\mathrm{\Lambda }}dr^2+r^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2)$$
(9)
where $`e^{2\mathrm{\Lambda }}=12\mu (r)/r`$ while for the calculation of the mass function $`\mu (r)`$ and the “potential” function $`\mathrm{\Phi }(r)`$ the reader should refer to Paper I. Finally, the stellar matter is assumed to be a perfect fluid
$$\stackrel{~}{T}_{\mu \nu }=\left(\stackrel{~}{\rho }+\stackrel{~}{P}\right)\stackrel{~}{U}_\mu \stackrel{~}{U}_\nu +\stackrel{~}{P}\stackrel{~}{g}_{\mu \nu }.$$
(10)
where $`\stackrel{~}{U}_\mu `$ is the four-velocity of the fluid, $`\stackrel{~}{\rho }`$ is the total energy density, and $`\stackrel{~}{P}`$ is the pressure in the physical frame.
## III Basic Perturbation Equations
In this section we present the equations describing perturbations of the spacetime, scalar field, and fluid in a spherically symmetric background. The equations we provide describe the non-radial oscillations of spherically symmetric neutron stars in scalar-tensor theories. We assume, in the physical frame, using the Regge-Wheeler gauge RW1957 , the following form of the perturbed metric tensor
$$\stackrel{~}{h}_{\mu \nu }=\stackrel{~}{h}_{\mu \nu }^{()}+\stackrel{~}{h}_{\mu \nu }^{(+)},$$
(11)
where $`\stackrel{~}{h}_{\mu \nu }^{()}`$ denotes the axial (or odd parity) part of metric perturbations
$`\stackrel{~}{h}_{\mu \nu }^{()}`$ $`=`$ $`{\displaystyle \underset{l=2}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=l}{\overset{l}{}}}\left(\begin{array}{cccc}0& 0& h_{0,lm}\mathrm{sin}^1\theta _\varphi & h_{0,lm}\mathrm{sin}\theta _\theta \\ 0& 0& h_{1,lm}\mathrm{sin}^1\theta _\varphi & h_{1,lm}\mathrm{sin}\theta _\theta \\ & & 0& 0\\ & & 0& 0\end{array}\right)Y_{lm},`$ (16)
and $`\stackrel{~}{h}_{\mu \nu }^{(+)}`$ denotes the polar (or even parity) part of metric perturbations
$`\stackrel{~}{h}_{\mu \nu }^{(+)}`$ $`=`$ $`{\displaystyle \underset{l=2}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=l}{\overset{l}{}}}\left(\begin{array}{cccc}H_{0,lm}e^{2\mathrm{\Phi }}& H_{1,lm}& 0& 0\\ & H_{2,lm}e^{2\mathrm{\Lambda }}& 0& 0\\ 0& 0& r^2K_{lm}& 0\\ 0& 0& 0& r^2K_{lm}\mathrm{sin}^2\theta \end{array}\right)Y_{lm}.`$ (21)
The functions $`h_{0,lm}`$, $`h_{1,lm}`$, $`H_{0,lm}`$, $`H_{1,lm}`$, $`H_{2,lm}`$, and $`K_{lm}`$ describing the spacetime perturbations have only radial and temporal dependence while $`Y_{lm}=Y_{lm}(\theta ,\varphi )`$ is the spherical harmonic function.
Following the previous definitions the perturbed metric tensor $`h_{\mu \nu }`$ in the Einstein frame has the form:
$`h_{\mu \nu }=`$ $`{\displaystyle \frac{1}{A^2}}\stackrel{~}{h}_{\mu \nu }{\displaystyle \frac{2}{A}}g_{\mu \nu }\delta A`$ (22)
$`=`$ $`{\displaystyle \underset{l=2}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=l}{\overset{l}{}}}{\displaystyle \frac{1}{A^2}}\left(\begin{array}{cccc}\left(H_{0,lm}+2A\delta A\right)e^{2\mathrm{\Phi }}& H_{1,lm}& h_{0,lm}\mathrm{sin}^1\theta _\varphi & h_{0,lm}\mathrm{sin}\theta _\theta \\ & \left(H_{2,lm}2A\delta A\right)e^{2\mathrm{\Lambda }}& h_{1,lm}\mathrm{sin}^1\theta _\varphi & h_{1,lm}\mathrm{sin}\theta _\theta \\ & & \left(K_{lm}2A\delta A\right)r^2& 0\\ & & 0& \left(K_{lm}2A\delta A\right)r^2\mathrm{sin}^2\theta \end{array}\right)Y_{lm},`$ (27)
where $`\delta AA\beta \phi \delta \phi `$ is the perturbation of the conformal factor $`A`$ and $`\delta \phi `$ is the perturbation of the scalar field, where $`\delta \phi `$ is a function of $`t`$ and $`r`$ only. The above definition of $`h_{\mu \nu }`$ will be used to derive the perturbation equations: in other words, we will work out the perturbations in the Einstein frame and we will transform back to the physical frame whenever we need it.
By defining the new set of perturbation functions $`\widehat{H}_0`$, $`\widehat{H}_1`$, $`\widehat{H}_2`$, $`\widehat{K}`$, $`\widehat{h}_0`$, and $`\widehat{h}_1`$ as follows
$`\widehat{H}_{0,lm}`$ $`={\displaystyle \frac{1}{A^2}}\left(H_{0,lm}+2A\delta A\right),`$ $`\widehat{H}_{1,lm}`$ $`={\displaystyle \frac{1}{A^2}}H_{1,lm},`$ (28)
$`\widehat{H}_{2,lm}`$ $`={\displaystyle \frac{1}{A^2}}\left(H_{2,lm}2A\delta A\right),`$ $`\widehat{K}_{lm}`$ $`={\displaystyle \frac{1}{A^2}}\left(K_{lm}2A\delta A\right),`$ (29)
$`\widehat{h}_{0,lm}`$ $`={\displaystyle \frac{1}{A^2}}h_{0,lm},`$ $`\widehat{h}_{1,lm}`$ $`={\displaystyle \frac{1}{A^2}}h_{1,lm}.`$ (30)
The perturbed metric $`h_{\mu \nu }`$, in the Einstein frame, is simplified considerably and reduced to the “standard” Regge-Wheeler form of a perturbed spherical metric. We should notice that the scalar perturbations $`\delta A`$ are linked only with the polar perturbation functions $`H_{0,lm}`$, $`H_{1,lm}`$, $`H_{2,lm}`$, and $`K_{lm}`$. The axial perturbation functions $`h_{0,lm}`$ and $`h_{1,lm}`$ are only affected by the contribution of the scalar field to the background.
The perturbation equations will be derived by taking the variation of Equations (3) and (4)
$`\delta G_{\mu \nu }`$ $`=8\pi G_{}\delta T_{\mu \nu }+\delta T_{\mu \nu }^{(\phi )},`$ (31)
$`\text{ }\text{ }\text{ }\text{ }\text{ }_{}\delta \phi `$ $`=4\pi G_{}\delta \left[\alpha (\phi )T_{}\right],`$ (32)
The various components of $`\delta T_{\mu \nu }^{(\phi )}`$ are expressed as linear combinations of $`\delta \phi `$ and $`\stackrel{~}{h}_{\mu \nu }`$. In the Einstein frame, the perturbed energy-momentum tensor describing the matter fields $`\delta T_{\mu \nu }`$, is some linear combination of the velocity variation of the fluid $`\delta \stackrel{~}{U}^i(WY_{lm},V_\theta Y_{lm}u\mathrm{sin}^1\theta _\varphi Y_{lm},V_\varphi Y_{lm}+u\mathrm{sin}\theta _\theta Y_{lm})`$, the density and pressure variations ($`\delta \stackrel{~}{\rho }`$ and $`\delta \stackrel{~}{P}`$) together with the variation of the scalar field ($`\delta \phi `$ or $`\delta A`$) and the metric perturbation $`h_{\mu \nu }`$. The explicit form of the energy-momentum tensor is given in Appendix A.
The linearized Einstein equations (31) will be written as follows. From the $`tt`$, $`tr`$, $`rr`$ components and the sum of the components $`\theta \theta `$ and $`\varphi \varphi `$, we get
$$\underset{l,m}{}A_{lm}^{(I)}Y_{lm}=0(I=0\text{to}3),$$
(33)
where the four expressions $`A_{lm}^{(I)}=0`$ are given in Appendix B. They contain combinations of $`\widehat{H}_0`$, $`\widehat{H}_1`$, $`\widehat{H}_2`$, $`\widehat{K}`$, $`W`$, $`\delta \stackrel{~}{P}`$, $`\delta \stackrel{~}{\rho }`$, $`\delta \phi `$ and their temporal and spatial derivatives. It is worth noticing that all four equations above are descibing only polar perturbations. In a similar way, from the $`t\theta `$, $`t\varphi `$, $`r\theta `$, and $`r\varphi `$ components, we get four more equations
$$\underset{l,m}{}\left\{\alpha _{lm}^{(J)}_\theta Y_{lm}+\beta _{lm}^{(J)}\frac{1}{\mathrm{sin}\theta }_\varphi Y_{lm}\right\}=0(J=0,1),$$
(34)
$$\underset{l,m}{}\left\{\beta _{lm}^{(J)}_\theta Y_{lm}\alpha _{lm}^{(J)}\frac{1}{\mathrm{sin}\theta }_\varphi Y_{lm}\right\}=0(J=0,1).$$
(35)
Here the expressions $`\alpha _{lm}^{(J)}`$ are some linear combinations of polar perturbation functions while on the other hand, the expression $`\beta _{lm}^{(J)}`$ is a combination of only axial perturbation functions (see Appendix B).
Furthermore, from the $`\theta \varphi `$ component and the subtraction of $`\theta \theta `$ and $`\varphi \varphi `$ components, we get two more equations
$$\underset{l,m}{}\left\{s_{lm}X_{lm}t_{lm}\mathrm{sin}\theta W_{lm}\right\}=0,$$
(36)
$$\underset{l,m}{}\left\{t_{lm}X_{lm}+s_{lm}\mathrm{sin}\theta W_{lm}\right\}=0,$$
(37)
where $`s_{lm}`$ and $`t_{lm}`$ describe polar and axial type perturbations respectively while $`X_{lm}`$ and $`W_{lm}`$ are defined as
$$X_{lm}=2_\varphi \left(_\theta \frac{\mathrm{cos}\theta }{\mathrm{sin}\theta }\right)Y_{lm}\text{and}W_{lm}=\left(_\theta ^2\frac{\mathrm{cos}\theta }{\mathrm{sin}\theta }_\theta \frac{1}{\mathrm{sin}^2\theta }_\varphi ^2\right)Y_{lm}.$$
(38)
Taking the product of Equations (33) – (37) with $`\overline{Y}_{lm}`$, integrating over the solid angle and looking at components with fixed values of $`l`$ and $`m`$, we get ten partial differential equations in the variables $`t`$ and $`r`$:
$$A_{lm}^{(I)}=0,\alpha _{lm}^{(J)}=0,s_{lm}=0,(I=0\mathrm{\hspace{0.17em}\hspace{0.17em}3}\text{and}J=0,1)$$
(39)
$$\beta _{lm}^{(J)}=0,t_{lm}=0(J=0,1).$$
(40)
Equations (39) describe the polar perturbations, and Equations (40) describe the axial perturbations. The analytic expressions for Equations (40), i.e., Eqs. (82), (83) and (85), do not couple to the perturbations of the scalar field $`\delta A`$ or $`\delta \phi `$. Therefore, the perturbed scalar field is coupled only to the polar perturbations.
### III.1 Axial perturbations
It is quite easy to derive a wave equations for the axial perturbations by combining equations (83) and (85)
$`\ddot{X}`$ $`e^{\mathrm{\Phi }\mathrm{\Lambda }}\left(e^{\mathrm{\Phi }\mathrm{\Lambda }}X^{}\right)^{}+e^{2\mathrm{\Phi }}\left({\displaystyle \frac{l(l+1)}{r^2}}{\displaystyle \frac{6\mu }{r^3}}+4\pi G_{}\left(\stackrel{~}{\rho }\stackrel{~}{P}\right)A^4\right)X=0,`$ (41)
where we introduce the new function $`X(t,r)`$ defined as $`\widehat{h}_1=e^{\mathrm{\Lambda }\mathrm{\Phi }}Xr`$. This equation does not couple to the scalar field perturbations, as we have mentioned earlier, and the effects of the scaler field will enter only via the background terms. Thus for $`\beta =0`$, i.e., $`A=1`$, it reduces to the standard wave equation describing axial perturbations CF91 . Finally, from equation (83) we get the following relation,
$$\dot{u}=e^{2\mathrm{\Phi }}\dot{\widehat{h}}_0,$$
(42)
which suggests that there are no axial oscillatory fluid motions, i.e., they have zero frequency and represent stationary currents. Thus axial perturbations are described only by a single wave equation (41) which does not couple neither to polar fluid or spacetime perturbations nor to the perturbed scalar field and it can be studied independently.
### III.2 Polar perturbations
The equations describing polar perturbations can be simplified introducing a new set of perturbation functions. Introducing these functions we can reformulate the seven equations describing polar perturbations as a pair of wave equations and a constraint equation, using a procedure similar to Ref. Allen98 . The new metric perturbation functions will be
$$F(t,r)=r\widehat{K},\text{and}S(t,r)=\frac{e^{2\mathrm{\Phi }}}{r}\left(\widehat{H}_0\widehat{K}\right),$$
(43)
while the fluid perturbations can be described by variation of the enthalpy function, i.e.,
$`H(t,r)`$ $`={\displaystyle \frac{\delta \stackrel{~}{P}}{\stackrel{~}{\rho }+\stackrel{~}{P}}}.`$ (44)
The system of equations describing the polar perturbations can be reduced to the following pair of wave equations;
$`\ddot{F}`$ $`e^{2(\mathrm{\Phi }\mathrm{\Lambda })}F^{\prime \prime }e^{2\mathrm{\Phi }}\left[4\pi G_{}r\left(\stackrel{~}{P}\stackrel{~}{\rho }\right)A^4+{\displaystyle \frac{2\mu }{r^2}}\right]F^{}+2r\left(\stackrel{~}{\rho }+\stackrel{~}{P}\right)\left(1{\displaystyle \frac{1}{C_s^2}}\right)A^4e^{2\mathrm{\Phi }}H`$
$`+e^{2\mathrm{\Phi }}\left[{\displaystyle \frac{l(l+1)}{r^2}}{\displaystyle \frac{2\mu }{r^3}}4\pi G_{}\left(3\stackrel{~}{\rho }+\stackrel{~}{P}\right)A^42e^{2\mathrm{\Lambda }}\mathrm{\Psi }^2\right]F`$
$`+2e^{2\mathrm{\Lambda }}\left[1r^2\mathrm{\Psi }^24\pi G_{}r^2e^{2\mathrm{\Lambda }}\left(\stackrel{~}{\rho }+\stackrel{~}{P}\right)A^4\right]S+8e^{2\mathrm{\Phi }}\left[\mathrm{\Psi }e^{2\mathrm{\Lambda }}+4\pi G_{}\alpha r\left(\stackrel{~}{P}\stackrel{~}{\rho }\right)A^4\right]\delta \phi =0,`$ (45)
and
$`\ddot{S}`$ $`e^{2(\mathrm{\Phi }\mathrm{\Lambda })}S^{\prime \prime }e^{2\mathrm{\Phi }}\left[4\pi G_{}r\left(\stackrel{~}{P}\stackrel{~}{\rho }\right)A^4+{\displaystyle \frac{2\mu }{r^2}}\right]S^{}+e^{2\mathrm{\Phi }}\left[{\displaystyle \frac{l(l+1)}{r^2}}{\displaystyle \frac{2\mu }{r^3}}4\pi G_{}\left(3\stackrel{~}{P}+\stackrel{~}{\rho }\right)A^4+2e^{2\mathrm{\Lambda }}\mathrm{\Psi }^2\right]S`$
$`+{\displaystyle \frac{4e^{4\mathrm{\Phi }}}{r^5}}\left\{\mathrm{\Phi }_{}^{}{}_{}{}^{2}r^3e^{2\mathrm{\Lambda }}+4\pi G_{}\stackrel{~}{\rho }Ar^33\mu +{\displaystyle \frac{4\pi G_{}}{r}}\left(\stackrel{~}{\rho }3\stackrel{~}{P}\right)\alpha A^4\mathrm{\Psi }+\left[4\pi G_{}\left(\stackrel{~}{\rho }\stackrel{~}{P}\right)r^5A^42r^3\right]\mathrm{\Psi }^2\right\}F`$
$`+4e^{4\mathrm{\Phi }}\left\{\mathrm{\Psi }^3e^{2\mathrm{\Lambda }}+\left[8\pi G_{}\left(2\stackrel{~}{P}\stackrel{~}{\rho }\right)A^4+{\displaystyle \frac{10\mu }{r^3}}{\displaystyle \frac{2}{r^2}}\right]\mathrm{\Psi }\right\}\delta \phi =0.`$ (46)
From the $`tt`$ component of the perturbed Einstein equations we get the Hamiltonian constraint:
$`F^{\prime \prime }`$ $`+\left[{\displaystyle \frac{e^{2\mathrm{\Lambda }}}{r^2}}\left(\mu 4\pi G_{}r^3\stackrel{~}{\rho }A^4\right){\displaystyle \frac{1}{2}}r\mathrm{\Psi }^2\right]F^{}+{\displaystyle \frac{e^{2\mathrm{\Lambda }}}{r^3}}\left[12\pi G_{}\stackrel{~}{\rho }r^3A^4\mu rl(l+1)+{\displaystyle \frac{1}{2}}r^3\mathrm{\Psi }^2e^{2\mathrm{\Lambda }}\right]F`$
$`re^{2\mathrm{\Phi }}S^{}+e^{2\mathrm{\Phi }+2\mathrm{\Lambda }}\left[\mathrm{\Psi }^2r^2e^{2\mathrm{\Lambda }}+8\pi G_{}r^2\left(\stackrel{~}{\rho }+\stackrel{~}{P}\right)A^42e^{2\mathrm{\Lambda }}{\displaystyle \frac{l(l+1)}{2}}\right]S`$
$`+{\displaystyle \frac{8\pi G_{}r}{C_s^2}}\left(\stackrel{~}{\rho }+\stackrel{~}{P}\right)e^{2\mathrm{\Lambda }}A^4H+2r\mathrm{\Psi }\delta \phi ^{}+32\pi G_{}rA^4e^{2\mathrm{\Lambda }}\stackrel{~}{\rho }\alpha \delta \phi =0.`$ (47)
Furthermore, for the special form of the conformal factor, i.e., $`A=\mathrm{exp}(\beta \phi ^2/2)`$ or $`\alpha =\beta \phi `$, from Equation (32) we obtain a wave equation for the perturbed scalar field $`\delta \phi `$:
$`\delta \ddot{\phi }`$ $`e^{\mathrm{\Phi }\mathrm{\Lambda }}\left(e^{\mathrm{\Phi }\mathrm{\Lambda }}\delta \phi ^{}\right)^{}{\displaystyle \frac{2}{r}}e^{2\mathrm{\Phi }2\mathrm{\Lambda }}\delta \phi ^{}+e^{2\mathrm{\Phi }}\left[{\displaystyle \frac{l(l+1)}{r^2}}+4e^{2\mathrm{\Lambda }}\mathrm{\Psi }^24\pi G_{}A^4\left(\stackrel{~}{\rho }3\stackrel{~}{P}\right)\left(4\alpha ^2+\beta \right)\right]\delta \phi `$
$`=`$ $`\left[r^2e^{2\mathrm{\Lambda }}\mathrm{\Psi }^3+{\displaystyle \frac{2\mu \mathrm{\Psi }}{r}}+4\pi G_{}rA^4\left\{2r\mathrm{\Psi }\stackrel{~}{P}\alpha \left(\stackrel{~}{\rho }3\stackrel{~}{P}\right)\right\}\right]S`$
$`+e^{2\mathrm{\Phi }}\left[e^{2\mathrm{\Lambda }}\mathrm{\Psi }^3+{\displaystyle \frac{2\mu \mathrm{\Psi }}{r^3}}+4\pi G_{}A^4\left\{2\mathrm{\Psi }\stackrel{~}{P}{\displaystyle \frac{\alpha }{r}}\left(\stackrel{~}{\rho }3\stackrel{~}{P}\right)\right\}\right]F4\pi G_{}A^4\alpha e^{2\mathrm{\Phi }}\left(\stackrel{~}{\rho }+\stackrel{~}{P}\right)\left({\displaystyle \frac{1}{C_s^2}}3\right)H,`$ (48)
where $`C_s^2=d\stackrel{~}{P}/d\stackrel{~}{\rho }`$.
Finally, by linearizing equation (7), i.e., the energy-momentum conservation equations, we get a system of evolution equations for the components of the perturbed velocity and the density perturbation
$`W^{}+{\displaystyle \frac{1}{C_s^2}}e^{2\mathrm{\Lambda }2\mathrm{\Phi }}\dot{H}+e^{2\mathrm{\Lambda }2\mathrm{\Phi }}\left(3\alpha \delta \dot{\phi }+{\displaystyle \frac{r}{2}}e^{2\mathrm{\Phi }}\dot{S}+{\displaystyle \frac{3}{2r}}\dot{F}\right)`$
$`{\displaystyle \frac{l(l+1)}{r^2}}e^{2\mathrm{\Lambda }}V+\left[2\mathrm{\Phi }^{}\mathrm{\Lambda }^{}+{\displaystyle \frac{2}{r}}+3\alpha \mathrm{\Psi }{\displaystyle \frac{1}{C_s^2}}\left(\mathrm{\Phi }^{}+\alpha \mathrm{\Psi }\right)\right]W=0,`$ (49)
$`\dot{W}+H^{}+\alpha \delta \phi ^{}+{\displaystyle \frac{r}{2}}e^{2\mathrm{\Phi }}S^{}{\displaystyle \frac{1}{2r}}F^{}+e^{2\mathrm{\Phi }}\left({\displaystyle \frac{1}{2}}r\mathrm{\Phi }^{}+r^2\mathrm{\Psi }^2+{\displaystyle \frac{2\mu }{r}}e^{2\mathrm{\Lambda }}+8\pi G_{}r^2A^4e^{2\mathrm{\Lambda }}\stackrel{~}{P}\right)S`$
$`+\left(\mathrm{\Psi }^2+{\displaystyle \frac{2\mu }{r^3}}e^{2\mathrm{\Lambda }}+{\displaystyle \frac{1}{2r^2}}+8\pi G_{}A^4e^{2\mathrm{\Lambda }}\stackrel{~}{P}\right)F+(\beta 4)\mathrm{\Psi }\delta \phi =0,`$ (50)
$`\dot{V}{\displaystyle \frac{r}{2}}e^{2\mathrm{\Phi }}S{\displaystyle \frac{1}{2r}}F+H+\alpha \delta \phi =0.`$ (51)
Concluding, the dynamics of the polar perturbation is described by the system of evolution equations (45), (46) and (48) together with the constraint equation (47). The rest of the functions describing the spacetime and fluid perturbations will be computed by taking proper combinations of $`F`$, $`S`$, $`H`$ and $`\delta \phi `$.
## IV Spacetime Perturbations in Scalar-Tensor Gravity
In Paper I we have studied stellar perturbations in scalar-tensor theories of gravity freezing the spacetime and scalar field perturbations. This is the so-called Cowling approximation, in which we only consider perturbations of the fluid. In practice we worked with a system of equations similar to (49) - (51), but setting $`H_1=H_0=H_2=K=\delta \phi =0`$. Even under this approximation, spontaneous scalarization has a remarkable effect on the oscillation spectra of the $`f`$ and $`p`$ modes. Based on this observation we proposed in Paper I that a successful detection of gravitational waves from oscillating stars will provide us with a tool to constrain the phenomenon of ’spontaneous scalarization’.
The quasinormal modes of the fluid perturbations described in Paper I will be affected by the coupling to the spacetime and scalar field perturbations. This is an interesting problem on its own. However, since we have already shown that the effect of spontaneous scalarization is quite strong when we limit consideration to perturbations of the fluid, in this paper we will examine the effect of the scalar field on the quasinormal modes describing the pure spacetime oscillations, i.e., the so-called $`w`$ modes Kokkotas1992 . The $`w`$ modes are similar to quasinormal modes of black holes. They have higher frequencies and shorter damping times that the $`f`$ modes, typical frequencies being of order $`712`$ kHz and damping times of order of 0.1 ms. These properties of the $`w`$ modes are common to axial and polar perturbations Benhar1999 . In the case of polar perturbations the $`w`$ modes are associated to small fluid motions, while in the axial case there is no coupling with the fluid at all. This is the reason why here we choose to study the effect of the scalar field considering only the axial perturbations, described by the single wave equation (41). We expect the effect of the scalar field on the axial and polar $`w`$ modes to be of the same order of magnitude. It should be pointed out here that, according to recent collapse calculations Baiotti05 , the $`w`$ modes are significantly excited. This adds further motivation to our study of the effects of scalar fields on $`w`$ mode oscillations.
The equations needed to construct the equilibrium stellar configurations as well as the equations of state (EOS) used are described in Paper I. In Paper I we also considered cases where the asymptotic value of the scalar field $`\phi _00`$, here, for simplicity, we only deal with scalar fields with $`\phi _0=0`$.
To compute the quasinormal frequencies of the axial $`w`$ modes we will use two different techniques. In the first approach we carry out time evolutions of Equation (41) and Fourier transform the signal at infinity; in the second approach we assume a harmonic time dependence of the perturbations, and the corresponding boundary value problem.
### IV.1 Evolving the axial perturbation equation
The time evolution of the 1+1 equation (41) is rather simple to obtain. We set some arbitrary initial data (for example a Gaussian pulse) in equation (41) and evolve these data for some time. Then we compute the oscillation frequencies by taking the Fourier transform of the signal emitted at infinity.
In Figure 1 we show the waveforms observed at a distance of about 300km from a neutron star with Arnowit-Deser-Misner (ADM) mass $`1.4M_{}`$. It is noticable that the arrival time of wave for different values of $`\beta `$ is not the same, because the effective potential due to central neutron star changes as a function of $`\beta `$. In this figure, we can see that the waveforms for $`\beta =0`$ and for $`\beta =4`$ are identical. This result can be understood as follows. For the axial perturbation, the gravitational wave is not coupled with the perturbation of the matter and the scalar field. So the presence of the scalar field is realized only due to the modified background. On the other hand, with $`\phi _0=0`$, the central value of scalar field $`\phi _c`$ is zero for any $`\beta >4.35`$ (see SotaniKokkotas2004 ). Thus for values of $`\beta >4.35`$ the effect of the scalar field is insignificant and it will affect the results only for $`\beta <4.35`$.
### IV.2 Boundary value method
Our second approach to calculate the quasinormal frequencies of the axial $`w`$ modes is more involved than simple time evolutions. However, using time evolutions we can only identify those modes that are significantly excited by a certain set of initial data. For example, using time evolutions it is not easy to identify quasinormal modes that damp out very fast. The approach we present here allows us to calculate both slowly and highly damped quasinormal modes. We Fourier-expand the wave equation (41) as $`X(t,r)=X(r)\mathrm{exp}(i\omega t)`$ and get
$`X^{\prime \prime }+{\displaystyle \frac{2\mu }{r(r2\mu )}}X^{}+\left(1{\displaystyle \frac{2\mu }{r}}\right)^1\left[\omega ^2e^{2\mathrm{\Phi }}{\displaystyle \frac{l(l+1)}{r^2}}+{\displaystyle \frac{6\mu }{r^3}}4\pi G_{}\left(\stackrel{~}{\rho }\stackrel{~}{P}\right)A^4\right]X=0.`$ (52)
In this way we obtain an eigenvalue problem: the complex quasinormal modes $`\omega `$ can be obtained imposing appropriate boundary conditions. In our case, the boundary conditions are that $`X(r)`$ should be regular at the stellar center and that there are no incoming waves at infinity. Inside the star we can just integrate the above differential equation; outside the star we use appropriate asymptotic expansions to ensure that there is no incoming radiation. Here we adopt a variant of Leaver’s continued fraction method Leaver1985 , that has been originally used for the calculation of quasinormal modes of black holes. The procedure is described in detail in Appendix C.
### IV.3 Results
The results of the two methods we described agree very well, providing a good consistency check on our calculations. In Figures 2, 3 and 4, we present the eigenvalues. Our results suggest that the presence of a spontaneous scalarization can be inferred from the $`w`$ modes emitted by a newly born, oscillating neutron star.
In Figure 2 we show the eigenfrequencies of the $`w`$ modes for neutron star models with $`M_{ADM}=1.4M_{}`$. The plot is reminiscent of earlier calculation of these modes (see e.g. Figure 3 in Kokkotas1999 ). The modes that might be relevant for gravitatonal wave detectors are the lowest $`w`$ modes Kokkotas1992 . The $`w_{\mathrm{II}}`$ modes Nollert1993 damp out roughly twice as fast as the $`w`$ modes, but having lower frequencies they could also be relevant for detection by Earth-based interferometers. The higher-frequency $`w`$ modes ($`w_2`$, $`w_3`$, $`w_4`$, …)are difficult, if not impossible to detect.
In the study of $`w`$ modes as a tool for asteroseismology Andersson1996 ; Andersson1998 ; Benhar1999 ; Kokkotas2001 it has been suggested that a proper normalization for $`Re(\omega )`$ is to multiply it with the radius $`R`$ of the star and to scale it as a function of the compactness $`M/R`$. This phenomenological argument has been recently verified analytically by Tsui and Leung Tsui2005 . Introducing $`f=Re(\omega )/2\pi `$, it is clear that $`Rf`$ scales linearly as function of the compactness $`M/R`$. This applies both to $`w_{\mathrm{II}}`$ and $`w_1`$ modes (and even to the higher overtones). The linear relations that can be derived from Figure 3 are
$$f_{\omega _1\mathrm{mode}}(\text{kHz})=\frac{1}{\overline{R}}\left(\alpha _1\beta _1\frac{\overline{M}}{\overline{R}}\right)\text{and}f_{\omega _{\mathrm{II}}\mathrm{mode}}(\text{kHz})=\frac{1}{\overline{R}}\left(\alpha _{\mathrm{II}}+\beta _{\mathrm{II}}\frac{\overline{M}}{\overline{R}}\right),$$
(53)
where the constants $`\alpha _1`$, $`\beta _1`$, $`\alpha _{\mathrm{II}}`$ and $`\beta _{\mathrm{II}}`$ are listed in Table 1.
Another reason why it is harder to detect high-damped quasinormal modes such as the $`w_{\mathrm{II}}`$ modes for compact stars is that the effective amplitude scales as the square root of the number of oscillations Kokkotas2001 . Typically we can hardly observe more than $`23`$ cycles for highly damped quasinormal modes of black-holes and for the $`w_{\mathrm{II}}`$ modes of compact stars. Spontaneous scalarization might help in this direction. Figure 4 shows that the damping time of the $`w_{\mathrm{II}}`$ mode for stars with $`\beta 4.35`$ is significantly longer than for typical stars in general relativity. The reason is that the presence of a scalar field increases the maximum mass of the stars and their compactness. Since the damping scales with compactness, the $`w_{\mathrm{II}}`$ modes live considerably longer. On the contrary, the damping times of the $`w_1`$ modes become shorter as the compactness increases (left panel in Figure 4).
## V Conclusion
We derived the equations describing stellar perturbations in scalar-tensor theories of gravity. The presence of a scalar field affects the equilibrium model, and consequently the oscillation spectrum. The scalar field perturbations couple with the polar perturbations of the spacetime and fluid, but they don’t couple with the axial perturbations. Since the spacetime modes of polar and axial perturbations have the same qualitative behavior, we have chosen to study the effect of the scalar field on the axial perturbations.
The results show that in the presence of spontaneous scalarization, a scalar field reduces the oscillation frequency of the $`w_1`$ modes by about 10% (i.e. by about 1kHz). The decrease in frequency for the $`w_{\mathrm{II}}`$ modes is about 25% the frequency of (i.e., about 1.5 kHz). The effect on the damping time is even more pronounced. The damping of $`w_{\mathrm{II}}`$ modes can decrease by as much as 30%, while it can increase by as much as 50% for the $`w_1`$ modes. Detectors operating at these high frequencies are under development. Through a detection of the $`w`$ mode spectrum, they could provide a unique proof for the existence of scalar fields with $`\beta 4.35`$.
A more detailed model of the effect of the scalar field on the oscillation spectra requires the inclusion of a larger set of equations of state. Another open problem is the study of polar oscillations, which couple directly to the scalar field. Work in these directions is in progress.
###### Acknowledgements.
We acknowledge valuable comments by Emanuele Berti and Shijun Yoshida. This work was partially supported by a Grant for The 21st Century COE Program (Holistic Research and Education Center for Physics Self-Organization Systems) at Waseda University and the Pythagoras I research grant of GSRT.
## Appendix A The Perturbed Energy Momentum Tensor
In this Appendix we show the explicit form of the various components of the perturbed energy momentum tensor (of the fluid and of the scalar field) appearing in the perturbation equations. We will use primes for spatial derivatives and dots for temporal derivatives. For simplicity we will omit the subscript $`lm`$ in the various perturbed quantities.
The components of the perturbed energy momentum tensor $`T_{\mu \nu }^{(\phi )}`$ for the scalar field have the form
$`\delta T_{tt}^{(\phi )}`$ $`=e^{2\mathrm{\Phi }2\mathrm{\Lambda }}\left[2\mathrm{\Psi }\delta \phi ^{}\left(\widehat{H}_0+\widehat{H_2}\right)\mathrm{\Psi }^2\right]Y_{lm},`$ (54)
$`\delta T_{tr}^{(\phi )}`$ $`=\left[2\mathrm{\Psi }\delta \dot{\phi }e^{2\mathrm{\Lambda }}\mathrm{\Psi }^2\widehat{H}_1\right]Y_{lm},`$ (55)
$`\delta T_{t\theta }^{(\phi )}`$ $`=e^{2\mathrm{\Lambda }}\mathrm{\Psi }^2\widehat{h}_0{\displaystyle \frac{1}{\mathrm{sin}\theta }}_\varphi Y_{lm},`$ (56)
$`\delta T_{t\varphi }^{(\phi )}`$ $`=e^{2\mathrm{\Lambda }}\mathrm{\Psi }^2\widehat{h}_0\mathrm{sin}\theta _\theta Y_{lm},`$ (57)
$`\delta T_{rr}^{(\phi )}`$ $`=2\mathrm{\Psi }\delta \phi ^{}Y_{lm},`$ (58)
$`\delta T_{r\theta }^{(\phi )}`$ $`=\left[2\mathrm{\Psi }\delta \phi _\theta +e^{2\mathrm{\Lambda }}\mathrm{\Psi }^2\widehat{h}_1{\displaystyle \frac{1}{\mathrm{sin}\theta }}_\varphi \right]Y_{lm},`$ (59)
$`\delta T_{r\varphi }^{(\phi )}`$ $`=\left[2\mathrm{\Psi }\delta \phi _\varphi e^{2\mathrm{\Lambda }}\mathrm{\Psi }^2\widehat{h}_1\mathrm{sin}\theta _\theta \right]Y_{lm},`$ (60)
$`\delta T_{\theta \theta }^{(\phi )}`$ $`=r^2e^{2\mathrm{\Lambda }}\left[2\mathrm{\Psi }\delta \phi ^{}+\left(\widehat{H}_2\widehat{K}\right)\mathrm{\Psi }^2\right]Y_{lm},`$ (61)
$`\delta T_{\varphi \varphi }^{(\phi )}`$ $`=r^2e^{2\mathrm{\Lambda }}\left[2\mathrm{\Psi }\delta \phi ^{}+\left(\widehat{H}_2\widehat{K}\right)\mathrm{\Psi }^2\right]\mathrm{sin}^2\theta Y_{lm}.`$ (62)
In order to get the components of the perturbed energy momentum tensor for the fluid we define, in the physical frame, the variations of pressure $`\delta \stackrel{~}{P}=\delta \stackrel{~}{P}Y_{lm}`$, energy density $`\delta \stackrel{~}{\rho }=\delta \stackrel{~}{\rho }Y_{lm}`$ and the components of the perturbed 4-velocity (in the physical frame)
$`\delta \stackrel{~}{U}^t`$ $`={\displaystyle \frac{1}{2A^3}}e^\mathrm{\Phi }H_0Y_{lm},`$ (63)
$`\delta \stackrel{~}{U}^r`$ $`={\displaystyle \frac{1}{A}}e^{\mathrm{\Phi }2\mathrm{\Lambda }}WY_{lm},`$ (64)
$`\delta \stackrel{~}{U}^\theta `$ $`={\displaystyle \frac{1}{Ar^2}}e^\mathrm{\Phi }\left(V_\theta Y_{lm}u{\displaystyle \frac{1}{\mathrm{sin}\theta }}_\varphi Y_{lm}\right),`$ (65)
$`\delta \stackrel{~}{U}^\varphi `$ $`={\displaystyle \frac{1}{Ar^2\mathrm{sin}^2\theta }}e^\mathrm{\Phi }\left(V_\varphi Y_{lm}+u\mathrm{sin}\theta _\theta Y_{lm}\right),`$ (66)
where the perturbation functions $`\delta \stackrel{~}{P}`$, $`\delta \stackrel{~}{\rho }`$, $`W`$, $`V`$, and $`u`$ defined in the previous relations depend only on $`t`$ and $`r`$. Using the above definition the components of the perturbed energy-momentum tensor $`\delta T_{\mu \nu }`$ take the form
$`\delta T_{tt}`$ $`=A^4e^{2\mathrm{\Phi }}\left[\delta \stackrel{~}{\rho }\stackrel{~}{\rho }\widehat{H}_0+{\displaystyle \frac{4\stackrel{~}{\rho }}{A}}\delta A\right]Y_{lm},`$ (67)
$`\delta T_{tr}`$ $`=A^4e^{2\mathrm{\Phi }}\left[\left(\stackrel{~}{\rho }+\stackrel{~}{P}\right)W+e^{2\mathrm{\Phi }}\stackrel{~}{\rho }\widehat{H}_1\right]Y_{lm},`$ (68)
$`\delta T_{t\theta }`$ $`=A^4e^{2\mathrm{\Phi }}\left(\stackrel{~}{\rho }+\stackrel{~}{P}\right)V_\theta Y_{lm}+A^4e^{2\mathrm{\Phi }}\left[\left(\stackrel{~}{\rho }+\stackrel{~}{P}\right)u+e^{2\mathrm{\Phi }}\stackrel{~}{\rho }\widehat{h}_0\right]{\displaystyle \frac{1}{\mathrm{sin}\theta }}_\varphi Y_{lm},`$ (69)
$`\delta T_{t\varphi }`$ $`=A^4e^{2\mathrm{\Phi }}\left(\stackrel{~}{\rho }+\stackrel{~}{P}\right)V_\varphi Y_{lm}A^4e^{2\mathrm{\Phi }}\left[\left(\stackrel{~}{\rho }+\stackrel{~}{P}\right)u+e^{2\mathrm{\Phi }}\stackrel{~}{\rho }\widehat{h}_0\right]\mathrm{sin}\theta _\theta Y_{lm},`$ (70)
$`\delta T_{rr}`$ $`=A^4e^{2\mathrm{\Lambda }}\left[\delta \stackrel{~}{P}+\stackrel{~}{P}\widehat{H}_2+{\displaystyle \frac{4\stackrel{~}{P}}{A}}\delta A\right]Y_{lm},`$ (71)
$`\delta T_{r\theta }`$ $`=A^4\stackrel{~}{P}\widehat{h}_1{\displaystyle \frac{1}{\mathrm{sin}\theta }}_\varphi Y_{lm},`$ (72)
$`\delta T_{r\varphi }`$ $`=A^4\stackrel{~}{P}\widehat{h}_1\mathrm{sin}\theta _\theta Y_{lm},`$ (73)
$`\delta T_{\theta \theta }`$ $`=A^4r^2\left[\delta \stackrel{~}{P}+\stackrel{~}{P}\widehat{K}+{\displaystyle \frac{4\stackrel{~}{P}}{A}}\delta A\right]Y_{lm},`$ (74)
$`\delta T_{\varphi \varphi }`$ $`=A^4r^2\left[\delta \stackrel{~}{P}+\stackrel{~}{P}\widehat{K}+{\displaystyle \frac{4\stackrel{~}{P}}{A}}\delta A\right]\mathrm{sin}^2\theta Y_{lm}.`$ (75)
## Appendix B The components of the Linearized Einstein equations
Here we provide the explicit form of the various expressions used in equations equations (33) – (37) for the description of the perturbed Einstein equations (31). We have chosen to use the same notation as Kojima Kojima1992 to facilitate comparison.
$`A_{lm}^{(0)}=`$ $`\widehat{K}^{\prime \prime }+e^{2\mathrm{\Lambda }}\left({\displaystyle \frac{1}{2}}re^{2\mathrm{\Lambda }}\mathrm{\Psi }^2+{\displaystyle \frac{5\mu }{r^2}}{\displaystyle \frac{3}{r}}+4\pi G_{}rA^4\stackrel{~}{\rho }\right)\widehat{K}^{}+{\displaystyle \frac{1}{r}}\widehat{H}_2^{}+{\displaystyle \frac{(l1)(l+2)}{2r^2}}e^{2\mathrm{\Lambda }}\widehat{K}`$
$`+e^{2\mathrm{\Lambda }}\left({\displaystyle \frac{l(l+1)}{2r^2}}+{\displaystyle \frac{1}{r^2}}8\pi G_{}A^4\stackrel{~}{\rho }\right)\widehat{H}_22\mathrm{\Psi }\delta \phi ^{}8\pi G_{}A^4e^{2\mathrm{\Lambda }}(\delta \stackrel{~}{\rho }+4\stackrel{~}{\rho }\alpha \delta \phi ),`$ (76)
$`A_{lm}^{(1)}=`$ $`\dot{\widehat{K}}+e^{2\mathrm{\Lambda }}\left({\displaystyle \frac{1}{2}}re^{2\mathrm{\Lambda }}\mathrm{\Psi }^2+{\displaystyle \frac{3\mu }{r^2}}{\displaystyle \frac{1}{r}}+4\pi G_{}rA^4\stackrel{~}{P}\right)\dot{\widehat{K}}+{\displaystyle \frac{1}{r}}\dot{\widehat{H}}_2+{\displaystyle \frac{l(l+1)}{2r^2}}\widehat{H}_1`$
$`2\mathrm{\Psi }\delta \dot{\phi }+8\pi G_{}A^4e^{2\mathrm{\Phi }}\left(\stackrel{~}{\rho }+\stackrel{~}{P}\right)W,`$ (77)
$`A_{lm}^{(2)}=`$ $`e^{2\mathrm{\Phi }+2\mathrm{\Lambda }}\ddot{\widehat{K}}+e^{2\mathrm{\Lambda }}\left({\displaystyle \frac{1}{2}}re^{2\mathrm{\Lambda }}\mathrm{\Psi }^2{\displaystyle \frac{\mu }{r^2}}+{\displaystyle \frac{1}{r}}+4\pi G_{}rA^4\stackrel{~}{P}\right)\widehat{K}^{}+{\displaystyle \frac{2}{r}}e^{2\mathrm{\Phi }}\dot{\widehat{H}}_1{\displaystyle \frac{1}{r}}\widehat{H}_0^{}{\displaystyle \frac{(l1)(l+2)}{2r^2}}e^{2\mathrm{\Lambda }}\widehat{K}`$
$`+{\displaystyle \frac{l(l+1)}{2r^2}}e^{2\mathrm{\Lambda }}\widehat{H}_0e^{2\mathrm{\Lambda }}\left({\displaystyle \frac{1}{r^2}}+8\pi G_{}A^4\stackrel{~}{P}\right)\widehat{H}_22\mathrm{\Psi }\delta \phi ^{}8\pi G_{}A^4e^{2\mathrm{\Lambda }}\left(\delta \stackrel{~}{P}+4\stackrel{~}{P}\alpha \delta \phi \right),`$ (78)
$`A_{lm}^{(3)}=`$ $`\widehat{K}^{\prime \prime }\widehat{H}_0^{\prime \prime }e^{2\mathrm{\Phi }+2\mathrm{\Lambda }}\left(\ddot{\widehat{K}}+\ddot{\widehat{H}}_2\right)+2e^{2\mathrm{\Phi }}\dot{\widehat{H}}_1^{}e^{2\mathrm{\Lambda }}\left({\displaystyle \frac{1}{2}}re^{2\mathrm{\Lambda }}\mathrm{\Psi }^2+{\displaystyle \frac{r+\mu }{r^2}}+4\pi G_{}\left(2\stackrel{~}{P}\stackrel{~}{\rho }\right)rA^4\right)\widehat{H}_0^{}`$
$`e^{2\mathrm{\Lambda }}\left({\displaystyle \frac{1}{2}}re^{2\mathrm{\Lambda }}\mathrm{\Psi }^2{\displaystyle \frac{\mu }{r^2}}+{\displaystyle \frac{1}{r}}+4\pi G_{}rA^4\stackrel{~}{P}\right)\widehat{H}_2^{}+e^{2\mathrm{\Lambda }}\left(4\pi G_{}\left(\stackrel{~}{P}\stackrel{~}{\rho }\right)rA^4+{\displaystyle \frac{2(r\mu )}{r^2}}\right)\widehat{K}^{}`$
$`2e^{2\mathrm{\Phi }+2\mathrm{\Lambda }}\left({\displaystyle \frac{1}{2}}re^{2\mathrm{\Lambda }}\mathrm{\Psi }^2+{\displaystyle \frac{\mu r}{r^2}}+4\pi G_{}rA^4\stackrel{~}{\rho }\right)\dot{\widehat{H}}_1+{\displaystyle \frac{l(l+1)}{2r^2}}e^{2\mathrm{\Lambda }}\widehat{H}_0`$
$`e^{2\mathrm{\Lambda }}\left(16\pi G_{}A^4\stackrel{~}{P}+{\displaystyle \frac{l(l+1)}{2r^2}}\right)\widehat{H}_2+4\mathrm{\Psi }\delta \phi ^{}16\pi G_{}A^4e^{2\mathrm{\Lambda }}\left(\delta \stackrel{~}{P}+4\stackrel{~}{P}\alpha \delta \phi \right),`$ (79)
$`\alpha _{lm}^{(0)}=`$ $`{\displaystyle \frac{1}{2}}e^{2\mathrm{\Lambda }}\left[\widehat{H}_1^{}e^{2\mathrm{\Lambda }}\left(\dot{\widehat{H}}_2+\dot{\widehat{K}}\right)+e^{2\mathrm{\Lambda }}\left(4\pi G_{}\left(\stackrel{~}{P}\stackrel{~}{\rho }\right)rA^4+{\displaystyle \frac{2\mu }{r^2}}\right)\widehat{H}_1\right]+8\pi G_{}A^4e^{2\mathrm{\Phi }}\left(\stackrel{~}{\rho }+\stackrel{~}{P}\right)V,`$ (80)
$`\alpha _{lm}^{(1)}=`$ $`{\displaystyle \frac{1}{2}}[\widehat{H}_0^{}\widehat{K}^{}e^{2\mathrm{\Phi }}\dot{\widehat{H}}_1+e^{2\mathrm{\Lambda }}({\displaystyle \frac{1}{2}}re^{2\mathrm{\Lambda }}\mathrm{\Psi }^2+{\displaystyle \frac{3\mu }{r^2}}{\displaystyle \frac{1}{r}}+4\pi G_{}rA^4\stackrel{~}{P})\widehat{H}_0`$
$`+e^{2\mathrm{\Lambda }}({\displaystyle \frac{1}{2}}re^{2\mathrm{\Lambda }}\mathrm{\Psi }^2{\displaystyle \frac{\mu }{r^2}}+{\displaystyle \frac{1}{r}}+4\pi G_{}rA^4\stackrel{~}{P})\widehat{H}_2]2\mathrm{\Psi }\delta \phi ,`$ (81)
$`\beta _{lm}^{(0)}=`$ $`{\displaystyle \frac{1}{2}}e^{2\mathrm{\Lambda }}\left(\widehat{h}_0^{\prime \prime }\dot{\widehat{h}}_1^{}\right)\left\{2\pi G_{}\left(\stackrel{~}{\rho }+\stackrel{~}{P}\right)rA^4+\left({\displaystyle \frac{r}{2}}\mu \right)\mathrm{\Psi }^2\right\}\left(\widehat{h}_0^{}\dot{\widehat{h}}_1\right){\displaystyle \frac{1}{r}}e^{2\mathrm{\Lambda }}\dot{\widehat{h}}_1`$
$`{\displaystyle \frac{1}{2r^3}}\left\{l(l+1)r4\mu +8\pi G_{}\left(\stackrel{~}{\rho }+\stackrel{~}{P}\right)r^3A^42r^3e^{2\mathrm{\Lambda }}\mathrm{\Psi }^2\right\}\widehat{h}_08\pi G_{}A^4e^{2\mathrm{\Phi }}\left(\stackrel{~}{\rho }+\stackrel{~}{P}\right)u,`$ (82)
$`\beta _{lm}^{(1)}=`$ $`{\displaystyle \frac{1}{2}}e^{2\mathrm{\Phi }}\left(\dot{\widehat{h}}_0^{}\ddot{\widehat{h}}_1\right){\displaystyle \frac{1}{r}}e^{2\mathrm{\Phi }}\dot{\widehat{h}}_0{\displaystyle \frac{(l1)(l+2)}{2r^2}}\widehat{h}_1,`$ (83)
$`s_{lm}=`$ $`{\displaystyle \frac{1}{2}}\left(\widehat{H}_0\widehat{H}_2\right),`$ (84)
$`t_{lm}=`$ $`e^{2\mathrm{\Phi }}\dot{\widehat{h}}_0e^{2\mathrm{\Lambda }}\widehat{h}_1^{}{\displaystyle \frac{1}{r^2}}\left\{2\mu +4\pi G_{}\left(\stackrel{~}{P}\stackrel{~}{\rho }\right)r^3A^4\right\}\widehat{h}_1,`$ (85)
where the equations are simplified by virtue of the equations (13) – (17) in Paper I.
## Appendix C Numerical techniques
In this Appendix we present two numerical techniques to determine quasinormal modes. The first is the direct evolution of the time dependent, axial perturbation equation (41). In the second approach we assume a harmonic decomposition for the perturbation function $`X`$ of the form $`X(r,t)=X(r)e^{i\omega t}`$, and solve the equation (52) as the eigenvalue problem. We shall consider the equations (52) in the interior and the exterior of the star; then we will find the eigenvalues (quasinormal modes) by matching the interior and exterior solutions.
The largest numerical error in the interior solution occurs at stellar surface, where the pressure is zero. In order to avoid this difficulty, we integrate the perturbation equation (52) from both sides, i.e. from the stellar center $`r=0`$ and from the stellar surface $`r=R`$. Then we match the solutions at some intermediate point, e.g., $`r=R/2`$ (see for example LD83 ; Kokkotas1992 ). In order to deal with the boundary condition at infinity we adopt the continued fraction method, originally used for black hole perturbations by Leaver Leaver1985 . To use this method we must know the forms of the coefficient in the perturbation equation as functions of $`1/r`$. Because of the presence of a scalar field , we do not know the exact forms of these coefficients. Therefore we just use the asymptotic forms of the coefficients, and derive a five-term recurrence relation. We believe that the QNMs obtained using these asymptotic forms are accurate enough, because the difference between the value of $`\mu `$ at the stellar surface and at infinity is not so large.
### C.1 Interior region of the star
The numerical integration of Equation (52) inside the star will be split (for numerical reasons) into two parts. First, we will integrate Equation (52) from the center towards $`R/2`$ and then we will integrate from the surface towards the same point. The matching of the two solutions will provide a unique solution valid throughout the star.
Near the center it can be shown that $`X`$ has a behavior of the form
$`X=X_cr^{l+1}\left(1+O(r^2)\right),`$ (86)
where $`X_c`$ is some arbitrary constant. Using this boundary condition (86) and by integrating equation (52) from $`r=0`$ to the matching point $`r=R/2`$, one can obtain the values of $`X(r)`$ and $`X^{}(r)`$. For convenience we represent the two functions $`X`$ and $`X^{}`$ in the vector form $`𝐘=(X,X^{})`$ and we will call $`𝐘_0(r)`$ the solution in the range $`0rR/2`$. The next step will be to integrate equation (52) from the stellar surface towards $`R/2`$ with a set of boundary conditions at $`r=R`$ such as $`(X(R),X^{}(R))=(1,0)`$ and $`(X(R),X^{}(R))=(0,1)`$. In this way we get two independent solutions $`𝐘_1(r)`$ and $`𝐘_2(r)`$ corresponding to each one of the previous boundary conditions. Thus the solution of the perturbation equation (52) is
$`𝐘(r)`$ $`=𝐘_0(r),\text{for}0rR/2,`$ (87)
$`𝐘(r)`$ $`=a𝐘_1(r)+b𝐘_2(r),\text{for}R/2rR,`$ (88)
where $`a`$ and $`b`$ are some constant, which will be determined from the junction condition at $`r=R/2`$:
$$𝐘_0(R/2)=a𝐘_1(R/2)+b𝐘_2(R/2).$$
(89)
The determination of the two constants specifies uniquely the solution in the interior of the star for a given value of the frequency $`\omega `$ and the constant $`X_c`$. At the stellar surface the values of $`X(R)`$, $`X^{}(R)`$ are simply $`X(R)=a`$ and $`X^{}(R)=b`$.
### C.2 Exterior region of the star
The functions describing the stellar background simplify considerably outside the star. This leads to a corresponding simplification of the wave equation (52). In the exterior, the equations describing the background reduces to
$`\mu ^{}`$ $`={\displaystyle \frac{1}{2}}re^{2\mathrm{\Lambda }}\mathrm{\Psi }^2,`$ (90)
$`\mathrm{\Phi }^{}`$ $`={\displaystyle \frac{1}{2}}r\mathrm{\Psi }^2+{\displaystyle \frac{\mu }{r^2}}e^{2\mathrm{\Lambda }},`$ (91)
$`\phi ^{}`$ $`=\mathrm{\Psi },`$ (92)
$`\mathrm{\Psi }^{}`$ $`={\displaystyle \frac{2}{r^2}}(r\mu )e^{2\mathrm{\Lambda }}\mathrm{\Psi }.`$ (93)
Therefore the asymptotic form of the above background quantities are
$`\mu `$ $`=M_{ADM}+{\displaystyle \frac{\mu _1}{r}}+O\left({\displaystyle \frac{1}{r^2}}\right),`$ (94)
$`\mathrm{\Phi }`$ $`={\displaystyle \frac{M_{ADM}}{r}}+O\left({\displaystyle \frac{1}{r^2}}\right),`$ (95)
$`\phi `$ $`=\phi _0+{\displaystyle \frac{\omega _A}{r}}+O\left({\displaystyle \frac{1}{r^2}}\right),`$ (96)
where $`\mu _1=\omega _A^2/2`$ and $`\omega _A=M_{\text{ADM}}\mathrm{\Psi }_s/\mathrm{\Phi }_s^{}`$.
The perturbation equation (52) in view of the above relations, outside the star, get the form
$`\left(1{\displaystyle \frac{2\mu }{r}}\right)X^{\prime \prime }+{\displaystyle \frac{2\mu }{r^2}}X^{}+\left[\omega ^2e^{2\mathrm{\Phi }}{\displaystyle \frac{l(l+1)}{r^2}}+{\displaystyle \frac{6\mu }{r^3}}\right]X=0.`$ (97)
which is similar (in the absence of a scalar field, identical) to the Regge-Wheeler RW1957 equation describing the axial perturbations in the exterior of a spherically symmetric spacetime (either a black-hole or a neutron star). Using as boundary values for the integration the values of $`X`$ and $`X^{}`$ at the surface given by the two relations $`X(R)=a`$ and $`X^{}(R)=b`$ one can integrate equation (97) together with (90) – (93) from the stellar surface towards infinite. The numerical integration will obviously stop at some distance $`r=r_a`$, where we will have to match the numerical solution with the appropriate asymptotic boundary conditions (in this case, the absence of incoming radiation).
In order to find the asymptotic form of the solution of equation (97) we can assume a solution of the form
$$X(r)=\left(\frac{r}{2\widehat{M}}1\right)^{2i\widehat{M}\omega }e^{i\omega r}\underset{n=0}{\overset{\mathrm{}}{}}a_n\left(1\frac{r_a}{r}\right)^n,$$
(98)
where $`\widehat{M}=M_{ADM}`$. By substituting this form of the solution into the perturbation equation (97) and taking the leading orders for $`\mu `$ and $`\mathrm{\Phi }`$, i.e., keeping only the terms up to order $`1/r`$, from equations (94) and (95), we obtain a five-term recurrence relation for the expansion coefficients $`a_n`$ $`(n1)`$,
$$\alpha _na_{n+1}+\beta _na_n+\gamma _na_{n1}+\delta _na_{n2}+ϵ_na_{n3}=0,$$
(99)
where the coefficients of the recurrence relation are given by the following formulae
$`\alpha _n`$ $`=c_0n(n+1),`$ (100)
$`\beta _n`$ $`=d_0n+c_1n(n1),`$ (101)
$`\gamma _n`$ $`=e_0+d_1(n1)+c_2(n1)(n2),`$ (102)
$`\delta _n`$ $`=e_1+d_2(n2)+c_3(n2)(n3),`$ (103)
$`ϵ_n`$ $`=e_2+d_3(n3)+c_4(n3)(n4).`$ (104)
The coefficients $`c_i`$, $`d_i`$ and $`e_i`$ are functions of the background quantities and have the form
$`c_0`$ $`=1{\displaystyle \frac{2\widehat{M}}{r_a}}{\displaystyle \frac{2\mu _1}{r_{a}^{}{}_{}{}^{2}}},`$ (105)
$`c_1`$ $`=2+{\displaystyle \frac{6\widehat{M}}{r_a}}+{\displaystyle \frac{8\mu _1}{r_{a}^{}{}_{}{}^{2}}},`$ (106)
$`c_2`$ $`=1{\displaystyle \frac{6\widehat{M}}{r_a}}{\displaystyle \frac{12\mu _1}{r_{a}^{}{}_{}{}^{2}}},`$ (107)
$`c_3`$ $`={\displaystyle \frac{2\widehat{M}}{r_a}}+{\displaystyle \frac{8\mu _1}{r_{a}^{}{}_{}{}^{2}}},`$ (108)
$`c_4`$ $`={\displaystyle \frac{2\mu _1}{r_{a}^{}{}_{}{}^{2}}},`$ (109)
$`d_0`$ $`=2i\omega r_a2+{\displaystyle \frac{6\widehat{M}}{r_a}}+{\displaystyle \frac{4i\omega \mu _1}{r_a}}+{\displaystyle \frac{8i\omega \widehat{M}\mu _1}{r_{a}^{}{}_{}{}^{2}}}+{\displaystyle \frac{6\mu _1}{r_{a}^{}{}_{}{}^{2}}},`$ (110)
$`d_1`$ $`=2{\displaystyle \frac{12\widehat{M}}{r_a}}{\displaystyle \frac{8i\omega \mu _1}{r_a}}{\displaystyle \frac{24i\omega \widehat{M}\mu _1}{r_{a}^{}{}_{}{}^{2}}}{\displaystyle \frac{18\mu _1}{r_{a}^{}{}_{}{}^{2}}},`$ (111)
$`d_2`$ $`={\displaystyle \frac{2}{r_a}}\left(3\widehat{M}+2i\omega \mu _1+{\displaystyle \frac{12i\omega \widehat{M}\mu _1}{r_a}}+{\displaystyle \frac{9\mu _1}{r_a}}\right),`$ (112)
$`d_3`$ $`={\displaystyle \frac{2\mu _1}{r_{a}^{}{}_{}{}^{2}}}\left(3+4i\omega \widehat{M}\right),`$ (113)
$`e_0`$ $`=l(l+1)+2\omega ^2\mu _1+{\displaystyle \frac{6\widehat{M}}{r_a}}+{\displaystyle \frac{8\omega ^2\widehat{M}\mu _1}{r_a}}{\displaystyle \frac{2i\omega \mu _1}{r_a}}{\displaystyle \frac{8i\omega \widehat{M}\mu _1}{r_{a}^{}{}_{}{}^{2}}}+{\displaystyle \frac{6\mu _1}{r_{a}^{}{}_{}{}^{2}}},`$ (114)
$`e_1`$ $`={\displaystyle \frac{2}{r_a}}\left(3\widehat{M}4\omega ^2\widehat{M}\mu _1+i\omega \mu _1+{\displaystyle \frac{8i\omega \widehat{M}\mu _1}{r_a}}{\displaystyle \frac{6\mu _1}{r_a}}\right),`$ (115)
$`e_2`$ $`={\displaystyle \frac{2\mu _1}{r_{a}^{}{}_{}{}^{2}}}\left(34i\omega \widehat{M}\right).`$ (116)
The first four terms of the recurrence relation (99) $`a_2`$, $`a_1`$, $`a_0`$, and $`a_1`$ are provided by the values of $`X`$ and $`X^{}`$ at $`r=r_a`$, i.e.,
$$a_2=a_1=0,a_0=\frac{X(r_a)}{\mathrm{\Xi }(r_a)},\text{and}a_1=\frac{r_a}{\mathrm{\Xi }(r_a)}\left[X^{}(r_a)+\frac{i\omega r_a}{r_a2\widehat{M}}X(r_a)\right],$$
(117)
where
$$\mathrm{\Xi }(r)=\left(\frac{r}{2\widehat{M}}1\right)^{2i\widehat{M}\omega }e^{i\omega r}.$$
(118)
The five term recurrence relations have in principle four possible solutions. A high order recurrence relation can generally be reduced to a three term recurrence relation, in which case convergence criteria for the solution can be easily applied, and we can identify the solution describing only outgoing radiation Leaver1985 . To obtain a three-term recurrence relation we define new coefficients $`\widehat{\alpha }_n`$, $`\widehat{\beta }_n`$, $`\widehat{\gamma }_n`$, and $`\widehat{\delta }_n`$ as
$$\widehat{\alpha }_n=\alpha _n,\widehat{\beta }_n=\beta _n,\widehat{\gamma }_n=\gamma _n\text{for}n=1,2$$
(119)
and for $`n3`$
$$\widehat{\alpha }_n=\alpha _n,\widehat{\beta }_n=\beta _n\frac{\widehat{\alpha }_{n1}ϵ_n}{\widehat{\delta }_{n1}},\widehat{\gamma }_n=\gamma _n\frac{\widehat{\beta }_{n1}ϵ_n}{\widehat{\delta }_{n1}},\widehat{\delta }_n=\delta _n\frac{\widehat{\gamma }_{n1}ϵ_n}{\widehat{\delta }_{n1}}.$$
(120)
The original five-term recurrence relation becomes a four-term recurrence relation:
$$\widehat{\alpha }_na_{n+1}+\widehat{\beta }_na_n+\widehat{\gamma }_na_{n1}+\widehat{\delta }_na_{n2}=0.$$
(121)
Note that for the case of $`\mu _1=0`$, that is the case of the standard neutron star obtained by Einstein’s theory for gravity ($`\beta =0`$), the recurrence relation for $`a_n`$ has four terms Benhar1999 ; Sotani2001 ; Berti2004 .
The final step will be to define another set of coefficients $`\stackrel{~}{\alpha }_n`$, $`\stackrel{~}{\beta }_n`$, and $`\stackrel{~}{\gamma }_n`$:
$$\stackrel{~}{\alpha }_1=\widehat{\alpha }_1,\stackrel{~}{\beta }_1=\widehat{\beta }_1,\stackrel{~}{\gamma }_1=\widehat{\gamma }_1,$$
(122)
and for $`n2`$
$$\stackrel{~}{\alpha }_n=\widehat{\alpha }_n,\stackrel{~}{\beta }_n=\widehat{\beta }_n\frac{\stackrel{~}{\alpha }_{n1}\widehat{\delta }_n}{\stackrel{~}{\gamma }_{n1}}\text{and}\stackrel{~}{\gamma }_n=\widehat{\gamma }_n\frac{\stackrel{~}{\beta }_{n1}\widehat{\delta }_n}{\stackrel{~}{\gamma }_{n1}},$$
(123)
The four-term recurrence relation (121) is thus reduced to a three-term relation of the form
$`\stackrel{~}{\alpha }_na_{n+1}+\stackrel{~}{\beta }a_n+\stackrel{~}{\gamma }_na_{n1}=0.`$ (124)
Using this three-term recurrence relation, the boundary condition can be expressed as a continued fraction relation between $`\stackrel{~}{\alpha }_n`$, $`\stackrel{~}{\beta }_n`$, and $`\stackrel{~}{\gamma }_n`$:
$$\frac{a_1}{a_0}=\frac{\stackrel{~}{\gamma }_1}{\stackrel{~}{\beta }_1}\frac{\stackrel{~}{\alpha }_1\stackrel{~}{\gamma }_2}{\stackrel{~}{\beta }_2}\frac{\stackrel{~}{\alpha }_2\stackrel{~}{\gamma }_3}{\stackrel{~}{\beta }_3}\mathrm{},$$
(125)
that can be rewritten as
$$0=\stackrel{~}{\beta }_0\frac{\stackrel{~}{\alpha }_0\stackrel{~}{\gamma }_1}{\stackrel{~}{\beta }_1}\frac{\stackrel{~}{\alpha }_1\stackrel{~}{\gamma }_2}{\stackrel{~}{\beta }_2}\frac{\stackrel{~}{\alpha }_2\stackrel{~}{\gamma }_3}{\stackrel{~}{\beta }_3}\mathrm{}f(\omega ),$$
(126)
where $`\stackrel{~}{\beta }_0a_1/a_0`$, $`\stackrel{~}{\alpha }_01`$. The eigenfrequency $`\omega `$ of a quasinormal mode can be obtained solving the equation $`f(\omega )=0`$. |
warning/0506/nlin0506051.html | ar5iv | text | # Quantum Graphology
## 1 Introduction
Quantum graphs have recently attracted a lot of interest . A special volume containing a number of contributions can be found in . The attention is due to the fact that quantum graphs can be viewed as typical and simple examples for the large class of systems in which classically chaotic dynamics implies universal spectral correlations in the semiclassical limit . Up to now we have only a very limited understanding of the reasons for this universality. In a semiclassical approach to this problem the main stumbling block is the intricate interference between the contributions of (exponentially many) periodic orbits . Using quantum graphs as model systems it is possible to pinpoint and isolate this central problem. In graphs, an exact trace formula exists which is based on the periodic orbits of a mixing classical dynamical system . Moreover the orbits can be specified by a finite symbolic code with Markovian grammar. Based on these simplifications it is possible to rewrite the spectral form factor or any other two-point correlation functions in terms of a combinatorial problem . This combinatorial problem on graphs has been solved with promising results: It was shown that the form factor, ensemble averaged over graphs with a single non-trivial vertex and two attached bonds (2-Hydra) coincides exactly with the random-matrix result for the $`2\times 2`$ CUE . A simple algorithm which can evaluate the resulting combinatorial sum for any graph was presented in . In the short-time expansion of the form factor for $`v`$-Hydra graphs (i. e. one central node with $`v`$ bonds attached) was computed in the limit $`N\mathrm{}`$. In a periodic-orbit sum was used to prove Anderson localization in an infinite chain graph with randomized bond lengths. In the form factor of binary graphs was shown to approach the random-matrix prediction when the number of vertices increases. In the second order contribution $`2\tau ^2`$ to the form factor, was derived and was shown to be related to correlations within pairs of orbits differing in the orientation of one of the two loops resulting from a self-intersection of the orbit. Finally, in a field theoretical method was used to evaluate exactly the form factor of large graphs. Very recently, the spectral properties of quantum graphs were studied experimentally by the Warsaw group who constructed a microwave graph network.
The transport properties of open quantum graphs were also investigated quite thoroughly. In compact graphs were connected with leads to infinity and was shown that they display all the features which characterize quantum chaotic scattering. In the open quantum graphs were used to calculate shot-noise corrections while in the same system was employed in order to understand current relaxation phenomena from open chaotic systems.
Quite recently the interest on quantum graphs moved towards understanding statistical properties of wavefunctions. In the statistics of the nodal points was analyzed, while in quantum graphs were used in order to understand scaring of quantum eigenstates. A scar is a quantum eigenfunction with excess density near an unstable classical periodic orbit (PO). Such states are not expected within Random-Matrix Theory (RMT), which predicts that wavefunctions must be evenly distributed over phase space, up to quantum fluctuations . Experimental evidence and applications of scars come from systems as diverse as microwave resonators , quantum wells in a magnetic field , Faraday waves in confined geometries , open quantum dots and semiconductor diode lasers .
This contribution, is structured in the following way. In the following Section 2, the main definitions and properties of quantum graphs are given. We concentrate on the unitary bond-scattering matrix $`U`$ which can be interpreted as a quantum evolution operator on the graph. Section 3 deals with the corresponding classical dynamical system. In Section 4, the statistical properties of the eigenphase spectrum of the bond-scattering matrix $`U`$ are analyzed and related to the periodic orbits of the classical dynamics. Scaring phenomenon is discussed and analyzed in Section 5. Finally, our conclusions and outlook are summarized in the last Section 6.
## 2 Quantum Graphs: Basic Facts
We start with a presentation and discussion of the Schrödinger operator for graphs. Graphs consist of $`V`$ vertices connected by $`B`$ bonds. The valency $`v_i`$ of a vertex $`i`$ is the number of bonds meeting at that vertex. The graph is called $`v`$-regular if all the vertices have the same valency $`v`$. When the vertices $`i`$ and $`j`$ are connected, we denote the connecting bond by $`b=(i,j)`$. The same bond can also be referred to as $`\stackrel{}{b}(Min(i,j),Max(i,j))`$ or $`\stackrel{}{b}(Max(i,j),Min(i,j))`$ whenever we need to assign a direction to the bond. A bond with coinciding endpoints is called a loop. Finally, a graph is called bipartite if the vertices can be divided into two disjoint groups such that any vertices belonging to the same group are not connected.
Associated to every graph is its connectivity (adjacency) matrix $`C_{i,j}`$. It is a square matrix of size $`V`$ whose matrix elements $`C_{i,j}`$ are given in the following way
$$C_{i,j}=C_{j,i}=\left\{\begin{array}{c}1\mathrm{if}i,j\mathrm{are}\mathrm{connected}\hfill \\ 0\mathrm{otherwise}\hfill \end{array}\right\}.$$
For graphs without loops the diagonal elements of $`C`$ are zero. The connectivity matrix of connected graphs cannot be written as a block diagonal matrix. The valency of a vertex is given in terms of the connectivity matrix, by $`v_i=_{j=1}^VC_{i,j}`$ and the total number of undirected bonds is $`B=\frac{1}{2}_{i,j=1}^VC_{i,j}`$.
For the quantum description we assign to each bond $`b=(i,j)`$ a coordinate $`x_{i,j}`$ which indicates the position along the bond. $`x_{i,j}`$ takes the value $`0`$ at the vertex $`i`$ and the value $`L_{i,j}L_{j,i}`$ at the vertex $`j`$ while $`x_{j,i}`$ is zero at $`j`$ and $`L_{i,j}`$ at $`i`$. We have thus defined the length matrix $`L_{i,j}`$ with matrix elements different from zero, whenever $`C_{i,j}0`$ and $`L_{i,j}=L_{j,i}`$ for $`b=1,\mathrm{},B`$. The wave function $`\mathrm{\Psi }`$ contains $`B`$ components $`\mathrm{\Psi }_{b_1}(x_{b_1}),\mathrm{\Psi }_{b_2}(x_{b_2}),\mathrm{},\mathrm{\Psi }_{b_B}(x_{b_B})`$ where the set $`\{b_i\}_{i=1}^B`$ consists of $`B`$ different undirected bonds.
The Schrödinger operator (with $`\mathrm{}=2m=1`$) is defined on a graph in the following way: On each bond $`b`$, the component $`\mathrm{\Psi }_b`$ of the total wave function $`\mathrm{\Psi }`$ is a solution of the one-dimensional equation
$$\left(\mathrm{i}\frac{\mathrm{d}}{\mathrm{d}x}A_b\right)^2\mathrm{\Psi }_b(x)=k^2\mathrm{\Psi }_b(x).$$
(1)
We included a “magnetic vector potential” $`A_b`$ (with $`\mathrm{}e(A_b)0`$ and $`A_\stackrel{}{b}=A__\stackrel{}{b}`$) which breaks the time reversal symmetry. In most applications we shall assume that all the $`A_b`$’s are equal and the bond index will be dropped. On each of the bonds, the general solution of (1) is a superposition of two counter propagating waves
$$\mathrm{\Psi }_{b=(i,j)}=a_{i,j}\mathrm{e}^{\mathrm{i}(k+A_{i,j})x_{i,j}}+a_{j,i}\mathrm{e}^{\mathrm{i}(k+A_{j,i})x_{j,i}}$$
(2)
The coefficients $`a_{i,j}`$ form a vector $`𝐚(a_{\stackrel{}{b}_1}`$, $`\mathrm{}`$, $`a_{\stackrel{}{b}_B},`$ $`a_{_{\underset{1}{\overset{}{b}}}},\mathrm{},a_{_{\underset{B}{\overset{}{b}}}})^T`$ of complex numbers which uniquely determines an element in a $`2B`$dimensional Hilbert space. This space corresponds to ”free wave” solutions since we did not yet impose any conditions which the solutions of (1) have to satisfy at the vertices.
The most general boundary conditions at the vertices are given in terms of unitary $`v_j\times v_j`$ vertex-scattering matrices $`\sigma _{l,m}^{(j)}(k)`$, where $`l`$ and $`m`$ go over all the vertices which are connected to $`j`$. At each vertex $`j`$, incoming and outgoing components of the wave function are related by
$$a_{j,l}=\underset{m=1}{\overset{v_j}{}}\sigma _{l,m}^{(j)}(k)e^{\mathrm{i}kL_{jm}}a_{m,j},$$
(3)
which implies current conservation. The particular form
$$\sigma _{l,m}^{(j)}=\frac{2}{v_j}\delta _{l,m}$$
(4)
for the vertex-scattering matrices was shown in to be compatible with continuity of the wave function and current conservation at the vertices. (4) is referred to as Neumann boundary conditions. Bellow, we will concentrate on this type of graphs. Moreover we will always assume fully connected graphs i.e. the valency is $`v_j=v=V1,j=1,\mathrm{},V`$.
Stationary states of the graph satisfy (3) at each vertex. These conditions can be combined into
$$𝐚=U(k)𝐚,$$
(5)
such that the secular equation determining the eigenenergies and the corresponding eigenfunctions of the graph is of the form
$$det\left[IU(k,A)\right]=0.$$
(6)
Here, the unitary bond-scattering matrix
$$U(k,A)=D(k;A)T$$
(7)
acting in the $`2B`$-dimensional space of directed bonds has been introduced. The matrices $`D`$ and $`T`$ are given by
$`D_{ij,i^{}j^{}}(k,A)`$ $`=`$ $`\delta _{i,i^{}}\delta _{j,j^{}}\mathrm{e}^{\mathrm{i}kL_{ij}+\mathrm{i}A_{i,j}L_{ij}};`$ (8)
$`T_{ji,nm}`$ $`=`$ $`\delta _{n,i}C_{j,i}C_{i,m}\sigma _{j,m}^{(i)}.`$
$`T`$ contains the topology of the graph and is equivalent to the complete set of vertex-scattering matrices, while $`D`$ contains the metric information about the bonds. Hereafter, the bond lengths $`L_m`$ ($`m=1,\mathrm{},B`$) will be chosen to be incommensurate in order to avoid non-generic degeneracies.
It is instructive to interpret the action of $`U`$ on an arbitrary graph state as its time evolution over an interval corresponding to the mean bond length of the graph such that
$$𝐚(t)=U^t𝐚(0),t=0,1,2,\mathrm{}.$$
(9)
Clearly the solutions of (5) are stationary with respect to this time evolution. $`n`$ in (9) represents a discrete (topological) time counting the collisions of the particle with vertices of the graph. In this ”picture” the diagonal matrix $`D_{mn}(k)=\delta _{mn}\mathrm{}^{ıkl_m}`$ describes the free propagation along the bonds of the network while $`T`$ assigns a scattering amplitude for transitions between connected directed bonds. As we will see in the next section it specifies a Markovian random walk on the graph which is the classical analogue of Eq. (9).
## 3 Periodic orbits and classical dynamics on graphs
In this section we discuss the classical dynamics corresponding to the quantum evolution (9) implied by $`U`$. To introduce this dynamics we employ a Liouvillian approach, where a classical evolution operator assigns transition probabilities in a phase space of $`2B`$ directed bonds . If $`\rho _b(t)`$ denotes the probability to occupy the (directed) bond $`b`$ at the (discrete) topological time $`t`$, we can write down a Markovian Master equation of the form
$$\rho _b(t+1)=\underset{b^{}}{}M_{b,b^{}}\rho _b^{}(t).$$
(10)
The classical (Frobenius-Perron) evolution operator $`M`$ has matrix elements
$$M_{ij,nm}=\delta _{j,n}P_{im}^{(j)}$$
(11)
with $`P_{jiij^{}}^{(i)}`$ denoting the transition probability between the directed bonds $`b=(j,i)`$ and $`b^{}=(i,j^{})`$. To make the connection with the quantum description, we adopt the quantum transition probabilities, expressed as the absolute squares of matrix elements of $`M`$
$$P_{jj^{}}^{(i)}=\left|\sigma _{j,j^{}}^{(i)}(k)\right|^2.$$
(12)
Note that $`P_{jj^{}}^{(i)}`$ and $`M`$ do not involve any metric information on the graph.
The unitarity of the bond-scattering matrix $`U`$ guarantees $`_{b=1}^{2B}M_{b,b^{}}=1`$ and $`0M_{b,b^{}}1`$, so that the total probability that the particle is on any bond remains conserved during the evolution. The spectrum of $`M`$, denoted as $`\{\mu _b\}`$ with $`b=1,\mathrm{}2B`$, is restricted to the interior of the unit circle and $`\mu _1=1`$ is always an eigenvalue with the corresponding eigenvector $`|1=\frac{1}{2B}(1,1,\mathrm{},1)^T`$. In most cases, the eigenvalue $`1`$ is the only eigenvalue on the unit circle. Then, the evolution is ergodic since any initial density will evolve to the eigenvector $`|1`$ which corresponds to a uniform distribution (equilibrium). The rate at which equilibrium is approached is determined by the gap to the next largest eigenvalue. If this gap exists, the dynamics is also mixing.
It was shown recently that mixing dynamics alone does not suffice to guarantee universality of the spectral statistics of quantum graphs <sup>1</sup><sup>1</sup>1For an example of a mixing graph with non-universal spectral statistics, see . An additional condition proven recently by Gutzmann and Altland states that in the limit of $`B\mathrm{}`$, the spectral gap has to be constant or at least vanish slowly enough as $`\mathrm{\Delta }_g(1|\mu _2|)B^\alpha `$ with $`0\alpha <0.5`$ and $`\mu _2`$ being the second maximum eigenvalue of $`M`$. In Fig. 1 we report our numerical results for Neumann fully connected graphs. We see that this type of graph satisfies the condition requested by .
Graphs are one dimensional and the motion on the bonds is simple and stable. Ergodic (mixing) dynamics is generated because at each vertex a (Markovian) choice of one out of $`v`$ directions is made. Thus, chaos on graphs originates from the multiple connectivity of the (otherwise linear) system .
Despite the probabilistic nature of the classical dynamics, the concept of a classical orbit can be introduced. A classical orbit on a graph is an itinerary of successively connected directed bonds $`(i_1,i_2),(i_2,i_3),\mathrm{}`$. An orbit is periodic with period $`t_p`$ if for all $`k`$, $`(i_{t_p+k},i_{t_p+k+1})=(i_k,i_{k+1})`$. For graphs without loops or multiple bonds, the sequence of vertices $`i_1,i_2,\mathrm{}`$ with $`i_m[1,V]`$ and $`C_{i_m,i_{m+1}}=1`$ for all $`m`$ represents a unique code for the orbit. This is a finite coding which is governed by a Markovian grammar provided by the connectivity matrix. In this sense, the symbolic dynamics on the graph is Bernoulli. This analogy is strengthened by further evidence: The number of $`t_p`$PO’s on the graph is $`\frac{1}{t}_p\mathrm{tr}C^{t_p}`$, where $`C`$ is the connectivity matrix. Since its largest eigenvalue $`\mathrm{\Gamma }_c`$ is bounded between the minimum and the maximum valency i.e. $`\mathrm{min}v_i\mathrm{\Gamma }_c\mathrm{max}v_i`$, periodic orbits proliferate exponentially with topological entropy $`\mathrm{log}\mathrm{\Gamma }_c`$.
From the previous discussion it is clear that all periodic orbits on a graph are unstable. The classical probability to remain at a specific PO of period $`t_p`$ is $`M_p=_{t=1}^{t_p}(M^t)_{j,j}`$. As $`M_p<1`$, the probability to follow the PO decreases exponentially with time. Assuming regular graphs of valency $`v_j=v`$ we can evaluate the rate of instability as
$$M_p=\underset{s=1}{\overset{r_p}{}}\left(1\frac{2}{v}\right)^2\underset{f=1}{\overset{t_pr_p}{}}\left(\frac{2}{v}\right)^2\mathrm{}^{\mathrm{\Lambda }_pt_p}$$
(13)
where $`\mathrm{\Lambda }_p`$ plays the role of the Lyapunov exponent (LE) and $`r_p`$ is the number of vertices where back scattering occurs. For the graphs studied in this contribution, some PO’s $`p`$ and LE $`\mathrm{\Lambda }_p`$, are listed in Fig. 2. The shortest PO’s have period $`2`$ and bounce back and forth between two vertices. For large graphs $`v\mathrm{}`$ these are by far the least unstable ones, as their LE approaches 0 while all others become increasingly unstable $`\mathrm{\Lambda }_p\mathrm{ln}v`$.
## 4 The spectral statistics of $`U`$
We consider the matrix $`U(k,A)`$ defined in Eqs. (7),8). The spectrum consist of $`2B`$ points $`e^{\mathrm{i}ϵ_l(k)}`$ confined to the unit circle (eigenphases). Unitary matrices of this type are frequently studied since they are the quantum analogues of classical, area preserving maps. Their spectral fluctuations depend on the nature of the underlying classical dynamics . The quantum analogues of classically integrable maps display Poissonian statistics while in the opposite case of classically chaotic maps, the eigenphase statistics conform with the results of RMT for Dyson’s circular ensembles. To describe the spectral fluctuations of $`U`$ we consider the form factor
$$K(t,2B)=\frac{1}{2B}|\mathrm{tr}U^t|^2(t>0).$$
(14)
The average $`\mathrm{}`$ will be specified below. RMT predicts that $`K(t,2B)`$ depends on the scaled time $`\tau =\frac{t}{2B}`$ only , and explicit expressions for the orthogonal and the unitary circular ensembles are known .
Using (7), (8) we expand the matrix products in $`\mathrm{tr}U^t`$ and obtain a sum of the form
$$\mathrm{tr}U^t(k)=\underset{p𝒫_t}{}𝒜_p\mathrm{e}^{\mathrm{i}(kL_p+Al_p)}.$$
(15)
In this sum $`p`$ runs over all closed trajectories on the graph which are compatible with the connectivity matrix and which have the topological length $`t`$, i. e. they visit exactly $`t`$ vertices. For graphs, the concepts of closed trajectories and periodic orbits coincide, hence (15) can also be interpreted as a periodic-orbit sum. From (15) it is clear that $`K(t/2B)=0`$ as long as $`t`$ is smaller than the period of the shortest periodic orbit. The phase associated with an orbit is determined by its total (metric) length $`L_p=_{bp}L_b`$ and by the “magnetic flux” through the orbit. The latter is given in terms of its total directed length $`l_p`$ if we assume for simplicity that the magnitude of the magnetic vector potential is constant $`|A_b|A`$. The amplitude of the contribution from a periodic orbit by the product of all the elements of vertex-scattering matrices encountered
$$𝒜_p=\underset{j=1}{\overset{n_p}{}}\sigma _{i_{j1},i_{j+1}}^{(i_j)}\underset{[r,s,t]}{}\left(\sigma _{r,t}^{(s)}\right)^{n_p(r,s,t)},$$
(16)
i. e. for fixed boundary conditions at the vertices it is completely specified by the frequencies $`n_p(r,s,t)`$ of all transitions $`(r,s)(s,t)`$ . Inserting (15) into the definition of the form factor we obtain a double sum over periodic orbits
$$K(t/2B)=\frac{1}{2B}\underset{p,p^{}𝒫_n}{}𝒜_p𝒜_p^{}\mathrm{exp}\left\{\mathrm{i}k(L_pL_p^{})+\mathrm{i}A(l_pl_p)\right\}.$$
(17)
Now we have to specify our averaging procedure which has to respect the restrictions imposed by the underlying classical dynamics. To this end we will use the wavenumber $`k`$ for averaging i.e. $`\mathrm{}_k=lim_k\mathrm{}k^1_0^k𝑑k^{}(\mathrm{})`$ (and, if present, also the magnetic vector potential $`A`$). Provided that the bond lengths of the graph are rationally independent and that a sufficiently large interval is used for averaging, we have
$$\mathrm{}^{ık(L_pL_p^{})}_k=\delta _{L_p,L_p^{}}\mathrm{and}\mathrm{}^{ıA(l_pl_p^{})}_A=\delta _{l_p,l_p^{}}$$
(18)
i.e. only terms with $`L_p=L_p^{}`$ and $`l_p=l_p^{}`$ survive.
Note that $`L_p=L_p^{}`$ does not necessarily imply $`p=p^{}`$ or that $`p,p^{}`$ are related by some symmetry because there exist families $``$ of distinct but isometric orbits which can be used to write the result of (17) in the form
$$K(t/2B)=\underset{_n}{}\left|\underset{p}{}𝒜_p\right|^2.$$
(19)
The outer sum is over the set $`_n`$ of families, while the inner one is a coherent sum over the orbits belonging to a given family (= metric length). An example of such family for the tetrahedron is shown in Fig. 3. Eq. (19) is exact, and it represents a combinatorial problem since it does not depend any more on metric information about the graph (the bond lengths).
In general, the combinatorial problem (19) is very hard and cannot be solved in closed form. Nevertheless exact result for finite $`t`$ can always be obtained from (19) using a computer algebra system such as Maple . To this end, one has to represent $`\mathrm{tr}U^t`$ as a multivariate polynomial of degree $`t`$ in the variables $`\mathrm{}^{ikL_i}`$, i. e.
$$\mathrm{tr}U^t=\underset{𝒫_t}{}c_𝒫(\mathrm{}^{ikL_1})^{p_1}(\mathrm{}^{ikL_2})^{p_2}\mathrm{},$$
(20)
where $`𝒫_t`$ runs over all partitions of $`t`$ into non-negative integers $`t=p_1+p_2+\mathrm{}`$ . The form factor is then simply given as
$$K(t/2B)=\underset{𝒫_t}{}|c_𝒫|^2.$$
(21)
The task of finding the coefficients $`c_𝒫`$ can be expressed in Maple with standard functions. In Fig. 4 we compare the results of (21) with direct numerical averages for fully connected graphs with $`V=4`$ and $`V=5`$ vertices with and without magnetic field breaking the time-reversal symmetry. The results agree indeed to a high precision. Although this could be regarded merely as an additional confirmation of the numerical procedures used in , we see the main merit of (21) in being a very useful tool for trying to find the solution of (19) in closed form.
## 5 Wavefunction statistics
Following the quantization outlined in section 2 a quantum wavefunction is defined as a set of $`2B`$ complex amplitudes $`a_d`$, normalized according to $`_d|a_d|^2=1`$. Here we will care about stationary solution satisfying Eq. (5) (i.e. eigenstates of the graph with corresponding wavelength $`k`$). The standard localization measure is the Inverse Participation Ratio (IPR) which is defined as
$$=\underset{d=1}{\overset{2B}{}}|a_d|^4.$$
(22)
Ergodic states which occupy each directed bond with the same probability have $`=1/2B`$ and up to a constant factor depending on the presence of symmetries this is also the RMT prediction. In the other extreme $`=0.5`$ indicates a state which is restricted to a single bond only, i. e. the greatest possible degree of localization. Some representative eigenstates are shown in Fig. 5.
The key theoretical idea discussed and applied in several recent works is that wavefunction intensities in a complex system can often be separated into a product of short-time and long-time parts, the latter being a random variable. On the other hand the short time part can be evaluated using information about classical dynamics. Specifically we have that the probability amplitude $`A_d`$ to return to the original state $`|d`$ is
$$A_dd|U^t|d=\underset{m}{}\left|d|m\right|^2\mathrm{}^{iϵ_mt}.$$
(23)
The return probability is then
$$P_d(t)\left|A_d\right|^2=\underset{m,n}{}\left|d|m\right|^2\left|d|n\right|^2\mathrm{}^{i(ϵ_mϵ_n)t}$$
(24)
Averaging over initial states and over time (typically larger than the Heisenberg time $`t_H=2\pi /\mathrm{\Delta }=2B`$) we get
$$\overline{P_d(t)}^t_d\frac{1}{2B}\underset{t=1}{\overset{2B}{}}P_d(t)_d=\frac{1}{2B}\underset{t=1}{\overset{2B}{}}P(t)=\underset{m}{}\left|d|m\right|^4_d_d$$
(25)
where $`\overline{\mathrm{}}^t`$ indicates an average over time and $`\mathrm{}_d`$ over initial states. Above $`P(t)`$ indicates the averaged (over initial states) return probability. In the last equality we had used the fact that due to time-average the off -diagonal terms averaged out to zero. Eq. (25) expresses the mean IPR in terms of the quantum return probability (RP), averaged over time and initial states. The next step is to argue that the quantum short-time dynamics, can be described by the classical time evolution (see Fig. 6). The latter can be approximated semiclassically quite well based only on period-two PO’s which correspond to trajectories which bounce back and forth between two vertices. These type of orbits have the lowest Lyapunov exponent (LE) and it is expected to have the largest influence on eigenfunction localization because classical trajectories can cycle in their vicinity for a relatively long time and increase the RP beyond the ergodic average. The resulting survival probability is
$$P(t)=\{\begin{array}{cc}0;\hfill & t\mathrm{odd}\hfill \\ \left(1+\frac{2}{v}\right)^4;\hfill & t\mathrm{even}\hfill \end{array}$$
(26)
Indeed the period $`2`$ orbits totally dominate the classical and quantum RP at short times as can be seen in Fig. 6. Including the contribution of these orbits only, Kaplan obtained a mean IPR which is by a factor $`v`$ larger than the RMT expectation, in agreement with numerics . Moreover, following the same line of argumentation as in we get that the bulk of the IPR distribution scales as
$$\stackrel{~}{P}(/)=P()$$
(27)
indicating that the whole bulk of $`𝒫()`$ is effectively determined by the least unstable orbits. This result can be nicely verified from the numerical data presented in Fig. 7.
With all this evidence for their prominent role in wavefunction localization, one clearly expects to see strong scarring on the period 2 orbits. Such states would essentially be concentrated on two directed bonds and give rise to $`1/2`$. However, in this region $`𝒫()`$ is negligible (see Fig. 7). We conclude therefore that the shortest and least unstable orbits of our system produce no visible scars. Note that the same applies also to the value $`=1/4`$ expected from the V-shaped orbits of Fig. 2. In fact $`𝒫()`$ has an appreciable value only for $`1/6`$ (Fig. 7). The position of this cutoff precisely coincides with the IPR expected for states which are scarred by triangular orbits. They occupy six directed bonds since, due to time- reversal symmetry, scarring on a PO and its reversed must coincide. Indeed a closer inspection shows that the vast majority of states at $`1/6`$ look like the example shown in Fig. 5. Of course the step at $`=1/6`$, which is present for any graph size $`V`$, is incompatible with the scaling of $`P()`$ mentioned above and indeed this relation breaks down in the tails at the expected points (inset of Fig. 7).
These results provide clear evidence for the fact that enhanced wavefunction localization due to the presence of short unstable orbits and strong scarring can in principle rely on completely unrelated mechanisms and can also leave distinct traces in statistical measures such as the distribution of inverse participation ratios (IPR). As a matter of fact in we were able to identify a necessary and sufficient condition for the energies of perfect scars
$$(kL_d)\text{mod}\pi =0dp$$
(28)
where $`d`$ is a directed bond which belongs to the specific PO $`p`$. Eq. (28) is reminiscent of a simple Bohr-Sommerfeld quantization condition $`kL_p=2n\pi `$, as it applies, e. g., to strong scars in billiards. However, there is an important difference: not only does Eq. (28) require quantization of the total action $`kL_p`$ of the scarred orbit, it also implies action quantization on all the visited bonds $`d`$. This stronger condition can only be met if the lengths of all bonds on $`p`$ are rationally related. As in general the bond lengths are incommensurate there are no perfect scars for generic graphs. Nevertheless, for incommensurate bond lengths Eq. (28) can be approximated with any given precision and then visible scars are expected .
## 6 Conclusions and Outlook
We have reviewed some of our results on the statistical properties of eigenvalues and eigenfunctions of the unitary quantum time evolution operator derived from quantum graphs. We have consentrated on fully connected quantum graphs. For this familly of graphs, the gap $`\mathrm{\Delta }_g`$ between the two maximum eigenvalues of the classical evolution operator approaches $`1`$ as the number of directed bonds increases, thus satisfying the (sufficient) condition for a graph in order to show universal spectral statistics. One possible approach in understanding how universality emerge is the use of combinatorial methods to perform the periodic-orbit sums related to spectral two-point correlations.
At the same time, we show that the existing scar theory does not explain the appearance of visible scars (super-scars). As a matter of fact our numerical data indicated that enhanced wavefunction localization due to short unstable orbits and strong scarring are not the same thing.
Quantum graphs were proven throughout the years very useful models. They allowed us to gain a good understanding of the spectrum and eigenfunctions properties of quantum systems with underlying classical chaotic dynamics. Semiclassics on graphs is exact, and various quantum mechanical quantities can be written in terms of classical periodic orbits. These studies and their conclusions are by now well documented in the quantum chaos literature. But quantum chaology has various other challenges that wait to be addressed. Among them is a quantum mechanical theory of dynamical evolution which is still a missing chapter. Quantum dissipation, dephasing and irreversibility (also used in the framework of “fidelity” studies in quantum computation) of quantum chaotic motion are notions, which are related with specific aspects of this evolution. It is our believe that quantum graphs can play a prominent role in this ultimate challenge: to develop a general theory for the time evolution of quantum systems with underlying classical chaotic behavior.
## Acknowledgments
We would like to express our gratitude to Prof. Uzy Smilansky to whom we owe our interest in the subject, and to Dr. Holger Schanz who contributed a great deal to the results discussed in this contribution. |
warning/0506/hep-th0506204.html | ar5iv | text | # Bulk versus Brane in the Absorption and Emission : 5D Rotating Black Hole Case
## Abstract
The absorption and emission spectra for the minimally-coupled brane and bulk scalar fields are numerically computed when the spacetime is a $`5d`$ rotating black hole carrying the two different angular momentum parameters $`a`$ and $`b`$. The effect of the superradiant scattering in the spectra is carefully examined. It is shown that the low-energy limit of the total absorption cross section always equal to the area of the non-spherically symmetric horizon, i.e. $`4\pi (r_H^2+a^2)`$ for the brane scalar and $`2\pi ^2(r_H^2+a^2)(r_H^2+b^2)/r_H`$ for the bulk scalar where $`r_H`$ is an horizon radius. The energy amplification for the bulk scalar is roughly order of $`10^9\%`$ while that for the brane scalar is order of unity. This indicates that the effect of the superradiance is negligible for the case of the bulk scalar. Thus the standard claim that black holes radiate mainly on the brane is not changed although the effect of the superradiance is taken into account. The physical implication of this fact is discussed in the context of TeV-scale gravity.
The recent brane-world scenarios which assume the large or warped extra dimensions generally allow the emergence of the TeV-scale gravity, which opens the possibility to make the tiny black holes factory in the future high-energy colliders. In this reason the absorption and emission problems for the higher-dimensional non-rotating black holes were extensively explored recently. It was found numerically that the emission on the brane is dominant compared to the emission off the brane in the Schwarzschild black hole background. This fact strongly supports the main conclusion of Ref.. Adopting a different numerical technique used in Ref., it was also found that the higher-dimensional charged black holes also radiate mainly on the brane if the number of the extra dimensions is not too large.
For the higher-dimensional rotating black holes, however, the situation can be more complicated. For the non-rotating black holes the crucial factor which makes the Hawking evaporation on the brane to be dominant is a geometrical factor $`r_H/L<<1`$, where $`r_H`$ is an horizon radius and $`L`$ is a size of the extra dimensions. In the rotating black holes, besides this geometrical factor, there is another important factor called superradiance, which means that the incident wave is amplified by the extraction of the rotation energy of the black holes under the particular condition. The effect of the superradiance in the $`4d`$ black holes was extensively studied long ago. The black hole bomb, i.e. rotating black hole plus mirror system, was recently re-examined in detail from the aspect of the black hole stability.
The importance of the superradiance modes in the tiny rotating black holes produced by the high energy scattering in the future collider was discussed in Ref.. Especially, in Ref. it was shown that superradiance for the bulk scalar in the background of the $`5d`$ Myers-Perry rotating black hole exists when the wave energy $`\omega `$ satisfies $`0<\omega <m\mathrm{\Omega }_a+k\mathrm{\Omega }_b`$, where $`\mathrm{\Omega }_a`$ and $`\mathrm{\Omega }_b`$ are angular frequencies of the black hole and, $`m`$ and $`k`$ are the azimuthal quantum numbers of the incident scalar wave. The generic conditions for the existence of the superradiance modes in the presence of single or multiple angular momentum parameters were derived recently when the incident bulk scalar, bulk electromagnetic and bulk gravitational waves are scattered by the higher-dimensional rotating black hole.
Recently, the emission spectra for the brane fields were explored analytically in the low-energy regime and numerically in the entire range of the energy. The crucial difference of the brane fields from the bulk fields is the fact that the condition for the existence of the superradiance for the brane fields is $`0<\omega <m\mathrm{\Omega }`$ while same condition for the bulk field is $`0<\omega <_im_i\mathrm{\Omega }_i`$ as shown in Ref.. Thus, in the background of the higher-dimensional black holes carrying the multiple angular momentum parameters the bulk field can be scattered superradiantly in the more wide range of $`\omega `$ compared to the brane fields. This may change the standard claim that black holes radiate mainly on the brane. The purpose of this paper is to explore this issue by choosing the $`5d`$ Myers-Perry rotating black hole with two angular momentum parameters $`a`$ and $`b`$ as a prototype.
In the following we will compute the absorption and emission spectra in the full range of $`\omega `$ for the brane scalar and bulk scalar. As a computational technique we will adopt an appropriate numerical technique which will be explained in detail later. Although the superradiant scattering takes place more readily for the bulk scalar, its energy amplification arising due to the superradiance is shown to be roughly $`10^9\%`$ while that for the brane scalar is order of unity. This fact indicates that the effect of the superradiance does not change the standard claim, i.e. black holes radiate mainly on the brane. In order to compare the superradiant effects in $`4d`$ and $`5d`$, we carry out the calculation in the Appendix for the superradiant scattering in the $`4d`$ Kerr background. The comparision reveals a big difference ($`8`$ orders of magnitude) for the energy amplification.
The $`5d`$ rotating black hole derived by Myers and Perry is expressed by a metric
$`ds_5^2`$ $`=`$ $`dt^2+{\displaystyle \frac{r^2\rho ^2}{}}dr^2+\rho ^2d\theta ^2+(r^2+a^2)\mathrm{sin}^2\theta d\varphi ^2+(r^2+b^2)\mathrm{cos}^2\theta d\psi ^2`$ (2)
$`+{\displaystyle \frac{r_0^2}{\rho ^2}}(dt+a\mathrm{sin}^2\theta d\varphi +b\mathrm{cos}^2\theta d\psi )^2`$
where $`0\varphi ,\psi <2\pi `$, $`0\theta \pi /2`$,
$`\rho ^2`$ $`=`$ $`r^2+a^2\mathrm{cos}^2\theta +b^2\mathrm{sin}^2\theta `$ (3)
$``$ $`=`$ $`(r^2+a^2)(r^2+b^2)r_0^2r^2`$ (4)
and, $`a`$ and $`b`$ are two angular momentum parameters. The mass $`M`$, two angular momenta $`J_1`$ and $`J_2`$ and the Hawking temperature $`T_H`$ are given by
$$M=\frac{3\pi r_0^2}{8}J_1=\frac{2}{3}MaJ_2=\frac{2}{3}MbT_H=\frac{r_H^4a^2b^2}{2\pi r_H(r_H^2+a^2)(r_H^2+b^2)}$$
(5)
where $`r_H`$ is an horizon radius defined by $`=0`$ at $`r=r_H`$.
The induced $`4d`$ metric on the brane can be written as
$$ds_4^2=dt^2+\frac{r^2\rho ^2}{}dr^2+\rho ^2d\theta ^2+(r^2+a^2)\mathrm{sin}^2\theta d\varphi ^2+\frac{r_0^2}{\rho ^2}(dt+a\mathrm{sin}^2\theta d\varphi )^2$$
(6)
if the self-gravity on the brane is negligible, where we assume $`0\theta \pi `$ to cover the whole $`4d`$ spacetime. The scalar wave equation $`\mathrm{}\mathrm{\Phi }_{BR}=0`$ in the background (6) is not separable. If, however, $`b=0`$, this wave equation is separable into the following radial and angular equations:
$`{\displaystyle \frac{d}{dr}}\left(\stackrel{~}{}{\displaystyle \frac{dR_{BR}}{dr}}\right)+\left[{\displaystyle \frac{[\omega (r^2+a^2)am]^2}{\stackrel{~}{}}}\mathrm{\Lambda }_{\mathrm{}}^m\right]R_{BR}=0`$ (7)
$`{\displaystyle \frac{1}{\mathrm{sin}\theta }}{\displaystyle \frac{d}{d\theta }}\left(\mathrm{sin}\theta {\displaystyle \frac{d\mathrm{\Theta }_{BR}}{d\theta }}\right)+\left[{\displaystyle \frac{m^2}{\mathrm{sin}^2\theta }}+\omega ^2a^2\mathrm{cos}^2\theta +_{\mathrm{}}^m\right]\mathrm{\Theta }_{BR}=0`$ (8)
where $`\stackrel{~}{}=r^2+a^2r_0^2`$ and $`\mathrm{\Lambda }_{\mathrm{}}^m=_{\mathrm{}}^m+a^2\omega ^22am\omega `$. When deriving Eq.(7), we used a factorization condition $`\mathrm{\Phi }_{BR}=e^{i\omega t}e^{im\varphi }R_{BR}(r)\mathrm{\Theta }_{BR}(\theta )`$.
The eigenvalue of the angular equation $`_{\mathrm{}}^m`$ was computed in Ref. as an expansion of $`a\omega `$. Since, however, we need $`_{\mathrm{}}^m`$ when $`a\omega `$ is arbitrarily large, we would like to solve the angular equation numerically. This is easily solved as following. First, we note that $`\mathrm{\Theta }_{BR}`$ becomes the usual spherical harmonics $`|\mathrm{},m>`$ when $`a\omega =0`$. Of course, in this case $`_{\mathrm{}}^m=\mathrm{}(\mathrm{}+1)`$. When $`a\omega 0`$, we expand $`\mathrm{\Theta }_{BR}`$ as $`\mathrm{\Theta }_{BR}=_{\mathrm{}^{}}C_{\mathrm{}\mathrm{}^{}}|\mathrm{}^{},m>`$. Then the angular equation reduces to the following eigenvalue equation
$$\underset{\mathrm{}^{}}{}A_{\mathrm{}^{\prime \prime }\mathrm{}^{}}^mC_{\mathrm{}\mathrm{}^{}}=_{\mathrm{}}^mC_{\mathrm{}\mathrm{}^{\prime \prime }}$$
(9)
where
$$A_{\mathrm{}^{\prime \prime }\mathrm{}^{}}^m=\mathrm{}^{}(\mathrm{}^{}+1)\delta _{\mathrm{}^{}\mathrm{}^{\prime \prime }}a^2\omega ^2<\mathrm{}^{\prime \prime },m|\mathrm{cos}^2\theta |\mathrm{}^{},m>.$$
(10)
Thus, the coefficients $`C_{\mathrm{}\mathrm{}^{}}`$ and the separation constant $`_{\mathrm{}}^m`$ are simultaneously obtained by computing the eigenvectors and eigenvalues of the matrix $`A_{\mathrm{}^{\prime \prime }\mathrm{}^{}}^m`$. Solving the eigenvalue equation (9) numerically, one can easily compute the $`a\omega `$-dependence of $`_{\mathrm{}}^m`$.
Now, we consider the wave equation $`\mathrm{}\mathrm{\Phi }_{BL}=0`$ for the bulk scalar in the background of the $`5d`$ metric (2). The wave equation is always separable and the radial and angular equations are
$`{\displaystyle \frac{}{r}}{\displaystyle \frac{d}{dr}}\left({\displaystyle \frac{}{r}}{\displaystyle \frac{R_{BL}}{dr}}\right)+WR_{BL}=0`$ (11)
$`{\displaystyle \frac{d}{d\theta }}\left(\mathrm{sin}\theta \mathrm{cos}\theta {\displaystyle \frac{d\mathrm{\Theta }_{BL}}{d\theta }}\right)+\left[\lambda _{\mathrm{}}^{m_1m_2}\omega ^2(a^2\mathrm{sin}^2\theta +b^2\mathrm{cos}^2\theta ){\displaystyle \frac{m_1^2}{\mathrm{sin}^2\theta }}{\displaystyle \frac{m_2^2}{\mathrm{cos}^2\theta }}\right]\mathrm{sin}\theta \mathrm{cos}\theta \mathrm{\Theta }_{BL}=0`$ (12)
where
$`W=\left[\lambda _{\mathrm{}}^{m_1m_2}+\omega ^2(r^2+a^2+b^2)+{\displaystyle \frac{m_1^2(a^2b^2)}{r^2+a^2}}+{\displaystyle \frac{m_2^2(b^2a^2)}{r^2+b^2}}\right]`$ (13)
$`+r_0^2(r^2+a^2)(r^2+b^2)\left(\omega {\displaystyle \frac{m_1a}{r^2+a^2}}{\displaystyle \frac{m_2b}{r^2+b^2}}\right)^2.`$ (14)
When deriving Eq.(11), a factorization condition $`\mathrm{\Phi }_{BL}=e^{i\omega t}e^{i(m_1\varphi +m_2\psi })R_{BL}(r)\mathrm{\Theta }_{BL}(\theta )`$ is used.
The angular equation can be solved numerically in a similar way to the case of the brane field. When $`a=b=0`$, the eigenfunction of the angular equation is expressed in terms of the Jacobi’s polynomial as following
$`\mathrm{\Theta }_{BL}|\mathrm{},m_1,m_2>=2^{(m_1+m_21)/2}\sqrt{{\displaystyle \frac{(2\mathrm{}+m_1+m_2+1)\mathrm{\Gamma }[\mathrm{}+1]\mathrm{\Gamma }[\mathrm{}+m_1+m_2+1]}{\mathrm{\Gamma }[\mathrm{}+m_1+1]\mathrm{\Gamma }[\mathrm{}+m_2+1]}}}`$ (15)
$`\times (1\mathrm{cos}2\theta )^{m_1/2}(1+\mathrm{cos}2\theta )^{m_2/2}P_{\mathrm{}}^{(m_1,m_2)}(\mathrm{cos}2\theta )`$ (16)
with $`\lambda _{\mathrm{}}^{m_1m_2}=(2\mathrm{}+m_1+m_2)(2\mathrm{}+m_1+m_2+2)`$, where $`P_{\mathrm{}}^{(m_1,m_2)}`$ is a jacobi’s polynomial. When $`a`$ and $`b`$ are nonzero, we expand $`\mathrm{\Theta }_{BL}`$ as $`\mathrm{\Theta }_{BL}=_{\mathrm{}^{}}D_{\mathrm{}\mathrm{}^{}}|\mathrm{}^{},m_1,m_2>`$. Then, by the same way as the brane case the angular equation reduces to the eigenvalue problem:
$$\underset{\mathrm{}^{}}{}B_{\mathrm{}^{\prime \prime }\mathrm{}^{}}^{m_1m_2}D_{\mathrm{}\mathrm{}^{}}=\lambda _{\mathrm{}}^{m_1m_2}D_{\mathrm{}\mathrm{}^{\prime \prime }}$$
(17)
where
$$B_{\mathrm{}^{\prime \prime }\mathrm{}^{}}^{m_1m_2}=(2\mathrm{}^{}+m_1+m_2)(2\mathrm{}^{}+m_1+m_2+2)\delta _{\mathrm{}^{}\mathrm{}^{\prime \prime }}+<\mathrm{}^{\prime \prime },m_1,m_2|\widehat{H}_1|\mathrm{}^{},m_1,m_2>$$
(18)
with $`\widehat{H}_1=a^2\omega ^2\mathrm{sin}^2\theta +b^2\omega ^2\mathrm{cos}^2\theta `$. Solving the eigenvalue equation (17) numerically, one can compute $`\lambda _{\mathrm{}}^{m_1m_2}`$.
Now, we would like to discuss how to solve the radial equations in (7) and (11). If we define $`x=\omega r`$ and $`x_H=\omega r_H`$, the radial equations reduce to
$`(x^2x_H^2){\displaystyle \frac{d}{dx}}(x^2x_H^2){\displaystyle \frac{dR_{BR}}{dx}}+\left[(x^2+a^2\omega ^2am\omega )^2\mathrm{\Lambda }_{\mathrm{}}^m(x^2x_H^2)\right]R_{BR}=0`$ (19)
$`f(x,x_H){\displaystyle \frac{d}{dx}}f(x,x_H){\displaystyle \frac{dR_{BL}}{dx}}+\omega ^4WR_{BL}=0`$ (20)
where
$$f(x,x_H)=\frac{\omega ^4}{x}=\frac{x_H^2(x^2+a^2\omega ^2)(x^2+b^2\omega ^2)x^2(x_H^2+a^2\omega ^2)(x_H^2+b^2\omega ^2)}{xx_H^2}.$$
(21)
The radial equations (19) imply that if $`R`$ is a solution, $`R^{}`$ is a solution too. The Wronskians between them become
$`W[R_{BR}^{},R_{BR}]_xR_{BR}^{}{\displaystyle \frac{dR_{BR}}{dx}}R_{BR}{\displaystyle \frac{dR_{BR}^{}}{dx}}={\displaystyle \frac{𝒞_1}{x^2x_H^2}}`$ (22)
$`W[R_{BL}^{},R_{BL}]_xR_{BL}^{}{\displaystyle \frac{dR_{BL}}{dx}}R_{BL}{\displaystyle \frac{dR_{BL}^{}}{dx}}={\displaystyle \frac{𝒞_2x}{(x^2x_H^2)(x^2a^2b^2\omega ^4/x_H^2)}}`$ (23)
where $`𝒞_1`$ and $`𝒞_2`$ are integration constants.
From the radial equations (19) one can derive the near-horizon and asymptotic solutions analytically as a series form. The explicit expressions for the solutions of the radial equations convergent near horizon are
$`𝒢_{\mathrm{},BR}^m(x,x_H)=e^{\rho _4\mathrm{ln}|xx_H|}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}d_{\mathrm{},n}^m(xx_H)^n`$ (24)
$`𝒢_{\mathrm{},BL}^{(m_1,m_2)}(x,x_H)=e^{\rho _5\mathrm{ln}|xx_H|}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}d_{\mathrm{},n}^{(m_1,m_2)}(xx_H)^n`$ (25)
where
$$\rho _4=i\frac{\omega (r_H^2+a^2)(\omega m\mathrm{\Omega }_a)}{2x_H}\rho _5=i\frac{r_H(r_H^2+a^2)(r_H^2+b^2)(\omega m_1\mathrm{\Omega }_am_2\mathrm{\Omega }_b)}{2(r_H^4a^2b^2)}.$$
(26)
In Eq.(24) we choosed the sign in the exponents so that the solutions (24) become ingoing in the frame of reference of an observer co-rotating with a black hole. In Eq.(26) $`\mathrm{\Omega }_a`$ and $`\mathrm{\Omega }_b`$ are the angular frequency of the rotating black hole corresponding to the angular momentum parameters $`a`$ and $`b`$:
$$\mathrm{\Omega }_a=\frac{a}{r_H^2+a^2}\mathrm{\Omega }_b=\frac{b}{r_H^2+b^2}.$$
(27)
The recursion relations for the coefficients $`d_{\mathrm{},n}^m`$ and $`d_{\mathrm{},n}^{(m_1,m_2)}`$ can be easily derived by inserting Eq.(24) into the radial equation (19). Since the explicit expressions are too lengthy, we will not present them. It is important to note that when $`\omega <m\mathrm{\Omega }_a`$, the imaginary part of $`\rho _4`$ becomes positive. This implies that the near-horizon solution for the brane wave equation becomes the outgoing wave with respect to an observer at infinity. This guarantees that the superradiant scattering occurs at $`0<\omega <m\mathrm{\Omega }_a`$ for the brane field. As expected, the second equation in Eq.(26) implies that the superradiance exists for the bulk scalar at $`0<\omega <m_1\mathrm{\Omega }_a+m_2\mathrm{\Omega }_b`$. Using Eq.(22) one can show that the Wronskians between the near-horizon solutions are
$`W[𝒢_{\mathrm{},BR}^m,𝒢_{\mathrm{},BR}^m]_x=2i\omega {\displaystyle \frac{(r_H^2+a^2)(\omega m\mathrm{\Omega }_a)}{x^2x_H^2}}|g_{\mathrm{}}^m|^2`$ (28)
$`W[𝒢_{\mathrm{},BL}^{(m_1,m_2)},𝒢_{\mathrm{},BL}^{(m_1,m_2)}]_x=2i\omega ^2{\displaystyle \frac{(r_H^2+a^2)(r_H^2+b^2)(\omega m_1\mathrm{\Omega }_am_2\mathrm{\Omega }_b)x}{r_H(x^2x_H^2)(x^2a^2b^2\omega ^4/x_H^2)}}|g_{\mathrm{}}^{(m_1,m_2)}|^2`$ (29)
where $`g_{\mathrm{}}^md_{\mathrm{},0}^m`$ and $`g_{\mathrm{}}^{(m_1,m_2)}d_{\mathrm{},0}^{(m_1,m_2)}`$.
Next let us consider the solutions of the radial equations (19) convergent at the asymptotic regime:
$`_{\mathrm{}(\pm ),BR}^m(x,x_H)=(\pm i)^{\mathrm{}+1}e^{ix}(xx_H)^{\pm \rho _4}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\tau _{n(\pm )}^{BR}x^{(n+1)}`$ (30)
$`_{\mathrm{}(\pm ),BL}^{(m_1,m_2)}(x,x_H)=(\pm i)^{\mathrm{}+3/2}{\displaystyle \frac{e^{ix}(xx_H)^{\pm \rho _5}}{\sqrt{x}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\tau _{n(\pm )}^{BL}x^{(n+1)}.`$ (31)
$`_{(+)}`$ and $`_{()}`$ represent the ingoing and outgoing waves respectively. The recursion relations between the coefficients are not explicitly given here. With an aid of Eq.(22) it is easy to show that the Wronskians between the asymptotic solutions are
$`W[_{\mathrm{}(+),BR}^m,_{\mathrm{}(),BR}^m]_x={\displaystyle \frac{2i}{x^2x_H^2}}`$ (32)
$`W[_{\mathrm{}(+),BL}^{(m_1,m_2)},_{\mathrm{}(),BL}^{(m_1,m_2)}]_x={\displaystyle \frac{2ix}{(x^2x_H^2)(x^2a^2b^2\omega ^4/x_H^2)}}.`$ (33)
Next, we would like to show how the coefficients $`g_{\mathrm{}}^m`$ and $`g_{\mathrm{}}^{(m_1,m_2)}`$ are related to the partial scattering amplitude. For this we define the real scattering solutions $`R_{\mathrm{},BR}^m`$ and $`R_{\mathrm{},BL}^{(m_1,m_2)}`$, which behave as
$`R_{\mathrm{},BR}^m\stackrel{xx_H}{}g_{\mathrm{}}^m(xx_H)^{\rho _4}[1+O(xx_H)]`$ (34)
$`R_{\mathrm{},BL}^{(m_1,m_2)}\stackrel{xx_H}{}g_{\mathrm{}}^{(m_1,m_2)}(xx_H)^{\rho _5}[1+O(xx_H)]`$ (35)
$`R_{\mathrm{},BR}^m\stackrel{x\mathrm{}}{}i^{\mathrm{}+1}{\displaystyle \frac{2\mathrm{}+1}{2x}}\left[e^{ix+\rho _4\mathrm{ln}|xx_H|}(1)^{\mathrm{}}S_{\mathrm{}}^m(x_H)e^{ix\rho _4\mathrm{ln}|xx_H|}\right]+O(x^2)`$ (36)
$`R_{\mathrm{},BL}^{(m_1,m_2)}\stackrel{x\mathrm{}}{}\sqrt{{\displaystyle \frac{2}{\pi }}}{\displaystyle \frac{i^{\mathrm{}+3/2}(\mathrm{}+1)^2}{x^{3/2}}}\left[e^{ix+\rho _5\mathrm{ln}|xx_H|}(1)^{\mathrm{}+1/2}S_{\mathrm{}}^{(m_1,m_2)}(x_H)e^{ix\rho _5\mathrm{ln}|xx_H|}\right]+O(x^{5/2})`$ (37)
where $`S_{\mathrm{}}^m(x_H)`$ and $`S_{\mathrm{}}^{(m_1,m_2)}(x_H)`$ are the scattering amplitudes for the brane and bulk scalars respectively. From the near-horizon behavior we can understand the Wronskians between the real scattering solutions $`W[R_{\mathrm{},BR}^m,R_{\mathrm{},BR}^m]_x`$ and $`W[R_{\mathrm{},BL}^{(m_1,m_2)},R_{\mathrm{},BL}^{(m_1,m_2)}]_x`$ are exactly same with Eq.(28) respectively.
If we define the phase shifts $`\delta _{\mathrm{}}^m(x_H)=(1/2i)\mathrm{ln}S_{\mathrm{}}^m(x_H)`$ and $`\delta _{\mathrm{}}^{(m_1,m_2)}(x_H)=(1/2i)\mathrm{ln}S_{\mathrm{}}^{(m_1,m_2)}(x_H)`$, the asymptotic behavior of $`R_{\mathrm{},BR}^m`$ and $`R_{\mathrm{},BL}^{(m_1,m_2)}`$ can be written as
$`R_{\mathrm{},BR}^m\stackrel{x\mathrm{}}{}{\displaystyle \frac{2\mathrm{}+1}{x}}e^{i\delta _{\mathrm{}}^m}\mathrm{sin}\left[x+i\rho _4\mathrm{ln}|xx_H|{\displaystyle \frac{\pi \mathrm{}}{2}}+\delta _{\mathrm{}}^m\right]+O(x^2)`$ (38)
$`R_{\mathrm{},BL}^{(m_1,m_2)}\stackrel{x\mathrm{}}{}\sqrt{{\displaystyle \frac{8}{\pi }}}{\displaystyle \frac{(\mathrm{}+1)^2}{x^{3/2}}}e^{i\delta _{\mathrm{}}^{(m_1,m_2)}}\mathrm{sin}\left[x+i\rho _5\mathrm{ln}|xx_H|\pi {\displaystyle \frac{\mathrm{}+1/2}{2}}+\delta _{\mathrm{}}^{(m_1,m_2)}\right]+O(x^{5/2}).`$ (39)
Assuming that the phase shifts are the complex quantities, i.e. $`\delta _{\mathrm{}}^m\eta _{\mathrm{}}^m+i\beta _{\mathrm{}}^m`$ and $`\delta _{\mathrm{}}^{(m_1,m_2)}\eta _{\mathrm{}}^{(m_1,m_2)}+i\beta _{\mathrm{}}^{(m_1,m_2)}`$, the Wronskians derived from the asymptotic behavior (38) are
$`W[R_{\mathrm{},BR}^m,R_{\mathrm{},BR}^m]_x={\displaystyle \frac{i(2\mathrm{}+1)^2}{x^2x_H^2}}e^{2\beta _{\mathrm{}}^m}\mathrm{sinh}2\beta _{\mathrm{}}^m`$ (40)
$`W[R_{\mathrm{},BL}^{(m_1,m_2)},R_{\mathrm{},BL}^{(m_1,m_2)}]_x={\displaystyle \frac{8i(\mathrm{}+1)^4x}{\pi (x^2x_H^2)(x^2a^2b^2\omega ^4/x_H^2)}}e^{2\beta _{\mathrm{}}^{(m_1,m_2)}}\mathrm{sinh}2\beta _{\mathrm{}}^{(m_1,m_2)}.`$ (41)
Equating Eq.(40) with Eq.(28) yields
$`|g_{\mathrm{}}^m|^2={\displaystyle \frac{\left(\mathrm{}+\frac{1}{2}\right)^2}{\omega (r_H^2+a^2)(\omega m\mathrm{\Omega }_a)}}(1e^{4\beta _{\mathrm{}}^m})`$ (42)
$`|g_{\mathrm{}}^{(m_1,m_2)}|^2={\displaystyle \frac{2(\mathrm{}+1)^4r_H}{\pi \omega ^2(r_H^2+a^2)(r_H^2+b^2)(\omega m_1\mathrm{\Omega }_am_2\mathrm{\Omega }_b)}}(1e^{4\beta _{\mathrm{}}^{(m_1,m_2)}}).`$ (43)
In the first equation of Eq.(42) $`|g_{\mathrm{}}^m|^2>0`$ implies that the greybody factor (or the transmission coefficient) $`1e^{4\beta _{\mathrm{}}^m}1|S_{\mathrm{}}^m|^2`$ becomes negative when $`0<\omega <m\mathrm{\Omega }_a`$, which is nothing but the superradiant scattering. Similarly, the superradiance for the bulk scalar exists when $`\omega `$ satisfies $`0<\omega <m_1\mathrm{\Omega }_a+m_2\mathrm{\Omega }_b`$, which is easily deduced from the second equation of Eq.(42).
Now, we would like to discuss how to compute the physical quantities such as absorption cross section and emission rate from $`g_{\mathrm{}}^m`$ and $`g_{\mathrm{}}^{(m_1,m_2)}`$. For this discussion it is convenient to introduce new wave solutions $`\stackrel{~}{R}_{\mathrm{},BR}^m`$ and $`\stackrel{~}{R}_{\mathrm{},BL}^{(m_1,m_2)}`$, which differ from $`R_{\mathrm{},BR}^m`$ and $`R_{\mathrm{},BL}^{(m_1,m_2)}`$ in their normalization. They are normalized as
$`\stackrel{~}{R}_{\mathrm{},BR}^m(x)\stackrel{xx_H}{}(xx_H)^{\rho _4}\left[1+O(xx_H)\right]`$ (44)
$`\stackrel{~}{R}_{\mathrm{},BL}^{(m_1,m_2)}(x)\stackrel{xx_H}{}(xx_H)^{\rho _5}\left[1+O(xx_H)\right].`$ (45)
Since $`_{(+)}`$ and $`_{()}`$ derived in Eq.(30) are linearly independent solutions of the radial equations, we can gererally express these new wave solutions as a linear combination of $`_{(\pm )}`$:
$`\stackrel{~}{R}_{\mathrm{},BR}^m=f_{\mathrm{}()}^m(x_H)_{\mathrm{}(+),BR}^m(x,x_H)+f_{\mathrm{}(+)}^m(x_H)_{\mathrm{}(),BR}^m(x,x_H)`$ (46)
$`\stackrel{~}{R}_{\mathrm{},BL}^{(m_1,m_2)}=f_{\mathrm{}()}^{(m_1,m_2)}(x_H)_{\mathrm{}(+),BL}^{(m_1,m_2)}(x,x_H)+f_{\mathrm{}(+)}^{(m_1,m_2)}(x_H)_{\mathrm{}(),BL}^{(m_1,m_2)}(x,x_H)`$ (47)
where the coefficients $`f_\pm `$ are called the jost functions. Using Eq.(32) one can compute the jost functions in the following:
$`f_{\mathrm{}(\pm )}^m(x_H)=\pm {\displaystyle \frac{x^2x_H^2}{2i}}W[_{\mathrm{}(\pm ),BR}^m,\stackrel{~}{R}_{\mathrm{},BR}^m]_x`$ (48)
$`f_{\mathrm{}(\pm )}^{(m_1,m_2)}(x_H)=\pm {\displaystyle \frac{(x^2x_H^2)(x^2a^2b^2\omega ^4/x_H^2)}{2ix}}W[_{\mathrm{}(\pm ),BL}^{(m_1,m_2)},\stackrel{~}{R}_{\mathrm{},BL}^{(m_1,m_2)}]_x.`$ (49)
Inserting the explicit expressions of $`_{(\pm )}`$ into Eq.(46) and comparing those with the asymptotic behavior of the real scattering solutions in Eq.(34), one can derive the following relations:
$`g_{\mathrm{}}^m(x_H)={\displaystyle \frac{\mathrm{}+\frac{1}{2}}{f_{\mathrm{}()}^m(x_H)}},S_{\mathrm{}}^m(x_H)={\displaystyle \frac{f_{\mathrm{}(+)}^m(x_H)}{f_{\mathrm{}()}^m(x_H)}}`$ (51)
$`g_{\mathrm{}}^{(m_1,m_2)}(x_H)=\sqrt{{\displaystyle \frac{2}{\pi }}}{\displaystyle \frac{(\mathrm{}+1)^2}{f_{\mathrm{}()}^{(m_1,m_2)}(x_H)}},S_{\mathrm{}}^{(m_1,m_2)}(x_H)={\displaystyle \frac{f_{\mathrm{}(+)}^{(m_1,m_2)}(x_H)}{f_{\mathrm{}()}^{(m_1,m_2)}(x_H)}}.`$
Combining Eq.(42) and Eq.(51), we can compute the greybody factors in terms of the jost functions:
$`𝒯_{\mathrm{},BR}^m1|S_{\mathrm{}}^m|^2={\displaystyle \frac{\omega (r_H^2+a^2)(\omega m\mathrm{\Omega }_a)}{|f_{\mathrm{}()}^m|^2}}`$ (52)
$`𝒯_{\mathrm{},BL}^{(m_1,m_2)}1|S_{\mathrm{}}^{(m_1,m_2)}|^2={\displaystyle \frac{\omega ^2(r_H^2+a^2)(r_H^2+b^2)(\omega m_1\mathrm{\Omega }_am_2\mathrm{\Omega }_b)}{r_H|f_{\mathrm{}()}^{(m_1,m_2)}|^2}}.`$ (53)
Thus if $`0<\omega <m\mathrm{\Omega }_a`$, $`𝒯_{\mathrm{},BR}^m`$ becomes negative which indicates the existence of the superradiance for the brane scalar. Same is true when $`0<\omega <m_1\mathrm{\Omega }_a+m_2\mathrm{\Omega }_b`$ for the bulk scalar.
The partial absorption cross section $`\sigma _{\mathrm{}}^m`$ for the brane scalar and $`\sigma _{\mathrm{}}^{(m_1,m_2)}`$ for the bulk scalar are given by
$`\sigma _{\mathrm{}}^m={\displaystyle \frac{\pi }{\omega ^2}}𝒯_{\mathrm{},BR}^m={\displaystyle \frac{\pi (r_H^2+a^2)(\omega m\mathrm{\Omega }_a)}{\omega |f_{\mathrm{}()}^m|^2}}`$ (54)
$`\sigma _{\mathrm{}}^{(m_1,m_2)}={\displaystyle \frac{4\pi }{\omega ^3}}𝒯_{\mathrm{},BL}^{(m_1,m_2)}={\displaystyle \frac{4\pi (r_H^2+a^2)(r_H^2+b^2)(\omega m_1\mathrm{\Omega }_am_2\mathrm{\Omega }_b)}{x_H|f_{\mathrm{}()}^{(m_1,m_2)}|^2}}.`$ (55)
Of course, the total absorption cross sections $`\sigma _{BR}`$ and $`\sigma _{BL}`$ are algebric sum of their partial absorption cross sections:
$$\sigma _{BR}=\underset{\mathrm{},m}{}\sigma _{\mathrm{}}^m\sigma _{BL}=\underset{\mathrm{},m_1,m_2}{}\sigma _{\mathrm{}}^{(m_1,m_2)}.$$
(56)
The total emission rate $`\mathrm{\Gamma }_{BR}`$ for the brane scalar and $`\mathrm{\Gamma }_{BL}`$ for the bulk scalar are given by
$$\mathrm{\Gamma }_{BR}=\underset{\mathrm{},m}{}\mathrm{\Gamma }_{\mathrm{}}^md\omega \mathrm{\Gamma }_{BL}=\underset{\mathrm{},m_1,m_2}{}\mathrm{\Gamma }_{\mathrm{}}^{(m_1,m_2)}d\omega $$
(57)
where
$`\mathrm{\Gamma }_{\mathrm{}}^m={\displaystyle \frac{1}{e^{(\omega m\mathrm{\Omega }_a)/T_H}1}}{\displaystyle \frac{\omega ^3\sigma _{\mathrm{}}^m}{2\pi ^2}}`$ (58)
$`\mathrm{\Gamma }_{\mathrm{}}^{(m_1,m_2)}={\displaystyle \frac{1}{e^{(\omega m_1\mathrm{\Omega }_am_2\mathrm{\Omega }_b)/T_H}1}}{\displaystyle \frac{\omega ^4\sigma _{\mathrm{}}^{(m_1,m_2)}}{8\pi ^2}}`$ (59)
and $`T_H`$ is an Hawking temperature given in Eq.(5). Therefore, we can compute all physical quantities related to the scattering between the rotating black hole and the scalar field if we can compute the jost functions.
Now, we would like to present briefly how to compute the jost functions numerically. It is important to note that besides the near-horizon or asymptotic solution, we can derive the solutions from the radial equations (19) which is convergent in the neighborhood of $`x=b`$, where $`b`$ is an arbitrary point. Their expressions are
$`\phi _{\mathrm{},BR}^m(x)=(xx_H)^{\rho _4}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}D_{\mathrm{},n}^m(xb)^n`$ (60)
$`\phi _{\mathrm{},BL}^{(m_1,m_2)}(x)=(xx_H)^{\rho _5}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}D_{\mathrm{},n}^{(m_1,m_2)}(xb)^n.`$ (61)
The recursion relations between the coefficients can be explicitly derived by inserting Eq.(60) into the radial equation (19), which is not presented in this paper. Thus one can perform the matching procedure between the near-horizon and the asymptotic solutions by making use of this intermediate solutions as following. Matching procedure between the near-horizon solutions and $`\phi _{\mathrm{}}`$ generates a solution whose domain of convergence is larger than the near-horizon region. Repeat of this matching procedure would increase the convergence region more and more. Similar matching procedure between the asymptotic solutions and $`\phi _{\mathrm{}}`$ can be repeated to decrease the convergent region from the asymptotic region. Eventually, we can obtain two solutions which have common domain of convergence. Using these solutions we can compute the jost functions with an aid of Eq.(48).
Fig. 1 shows the $`log`$-plot of the superradiant scattering for the brane field. The $`y`$-axis is a $`\mathrm{ln}(100T_{\mathrm{},BR}^m)`$ for the most favour modes. Fig. 1(a), (b), and (c) correspond to respectively $`a_{}=1.5`$, $`a_{}=2.0`$, and $`a_{}=2.5`$ where $`a_{}a/r_H`$. As commented in Ref. the superradiant scattering for the higher modes becomes more and more significant as $`a_{}`$ becomes larger. This fact is evidently verified in Fig. 1. Table I shows the first three modes at each $`a_{}`$ whose maximum energy amplification is large. This Table shows that the lowest $`\mathrm{}=m=1`$ mode has a maximum energy amplification at $`a_{}=1.5`$. But at $`a_{}=2.0`$ (or $`2.5`$) $`\mathrm{}=m=2`$ (or $`\mathrm{}=m=3`$) mode has a maximum amplification. It also shows that the average amplification tends to increase with increasing $`a_{}`$.
Table I: Maximum Energy Amplification for the Several Modes of Brane Scalar
| $`a_{}=1.5`$ | | $`a_{}=2.0`$ | | $`a_{}=2.5`$ | |
| --- | --- | --- | --- | --- | --- |
| modes | maximum energy | modes | maximum energy | modes | maximum energy |
| | amplification (%) | | amplification (%) | | amplification (%) |
| $`(1,1)`$ | $`2.1604`$ | $`(2,2)`$ | $`2.6785`$ | $`(3,3)`$ | $`2.8399`$ |
| $`(2,2)`$ | $`2.0863`$ | $`(3,3)`$ | $`2.5934`$ | $`(4,4)`$ | $`2.8212`$ |
| $`(3,3)`$ | $`1.7392`$ | $`(1,1)`$ | $`2.4476`$ | $`(2,2)`$ | $`2.7748`$ |
($`(p,q)`$ means $`\mathrm{}=p`$ and $`m=q`$.)
Fig. 2 shows the $`log`$-plot of the superradiant scattering for the bulk scalar. The vertical axis is a $`\mathrm{ln}(100T_{\mathrm{},BL}^{m_1,m_2})`$ when $`b_{}=0`$, $`0.5`$ and $`1`$ with $`a_{}=0.5`$ where $`b_{}b/r_H`$. The modes in this figure are selected by comparing the maximum energy amplification at fixed $`\mathrm{}`$. Usually one of the mode which satisfies $`m_1+m_2=\mathrm{}`$ has the largest maximum amplification at given $`\mathrm{}`$. When, especially, $`a_{}=b_{}`$, the amplifications for the modes which have same $`m_1+m_2`$ are exactly identical. For example, when $`\mathrm{}=2`$, the amplications for $`(0,2)`$, $`(1,1)`$ and $`(2,0)`$ are exactly identical, where $`(p,q)`$ means $`m_1=p`$ and $`m_2=q`$. When $`ab`$, this kind of degeneracy is broken. However, still one of the modes which satisfies $`m_1+m_2=\mathrm{}`$ generally has the largest maximum amplification. In Table II the maximum amplification values for several modes are given. Table I and Table II show that the energy amplification for the brane scalar is order of unity while that for the bulk scalar is order of $`10^9\%`$ The fact that the energy amplification for the bulk scalar is smaller than for the brane scalar can be partly understood if we counter the power of the energy dependence. while the energy amplification for the brane scalar is proportional to $`w^2`$, that for the bulk scalar is proportional to $`w^3`$. Since the superradiant scattering usually takes place in the low-energy region, we can conjecture that the energy amplification for the bulk scalar can be small. However, it does not explain the big difference between bulk and brane. It is unclear at least for us how to explain this issue physically.. This means that the effect of the superradiance for the bulk scalar is negligible. This indicates that consideration of the effect of the superradiance does not change the standard claim that black holes radiate mainly on the brane. This will be confirmed in Fig. 5 explicitly.
Table II: Maximum Energy Amplification for the Several Modes of Bulk Scalar
| $`a_{}=0.5`$, $`b_{}=0`$ | | $`a_{}=0.5`$, $`b_{}=0.5`$ | | $`a_{}=0.5`$, $`b_{}=1`$ | |
| --- | --- | --- | --- | --- | --- |
| modes | maximum energy | modes | maximum energy | modes | maximum energy |
| | amplification (%) | | amplification (%) | | amplification (%) |
| $`(1,1,0)`$ | $`4.723\times 10^9`$ | $`(1,1,1)`$ | $`6.494\times 10^9`$ | $`(1,0,1)`$ | $`4.748\times 10^8`$ |
| $`(1,1,1)`$ | $`1.849\times 10^{12}`$ | $`(1,0,1)`$ | $`4.682\times 10^9`$ | $`(1,1,1)`$ | $`3.456\times 10^8`$ |
| $`(2,2,0)`$ | $`5.933\times 10^{15}`$ | $`(2,0,2)`$ | $`5.035\times 10^{15}`$ | $`(2,0,2)`$ | $`2.013\times 10^{13}`$ |
| $`(2,1,0)`$ | $`5.632\times 10^{16}`$ | $`(2,1,2)`$ | $`4.961\times 10^{15}`$ | $`(2,1,2)`$ | $`1.121\times 10^{13}`$ |
| $`(3,3,0)`$ | $`5.444\times 10^{21}`$ | $`(3,0,3)`$ | $`4.019\times 10^{21}`$ | $`(3,0,3)`$ | $`6.492\times 10^{19}`$ |
| $`(3,2,0)`$ | $`1.768\times 10^{21}`$ | $`(3,1,3)`$ | $`3.540\times 10^{21}`$ | $`(3,1,3)`$ | $`3.323\times 10^{19}`$ |
($`(p,q,r)`$ means $`\mathrm{}=p`$, $`m_1=q`$ and $`m_2=r`$.)
Fig. 3 shows the total and partial absorption cross sections for the brane scalar when $`a_{}=0.5`$, $`1`$ and $`1.5`$. The partial absorption cross section plotted in Fig. 3 is defined as $`\sigma _{\mathrm{}}=_m\sigma _{\mathrm{}}^m`$. The negative value of $`\sigma _{\mathrm{}}^m`$ for $`m>0`$ in the range of $`0<\omega <m\mathrm{\Omega }_a`$ arising due to the superradiant scattering is compensated by the positive value of $`\sigma _{\mathrm{}}^m`$ for $`m0`$ in the same range arising due to the normal scattering. Therefore, the partial absorption cross section $`\sigma _{\mathrm{}}`$ is positive in the full range of $`\omega `$.
In the low-energy limits the total absorption cross sections exactly equal to the non-spherically symmetric horizon area $`𝒜_{BR}`$ defined
$$𝒜_{BR}_0^\pi 𝑑\theta _0^{2\pi }𝑑\varphi \sqrt{g_{\theta \theta }g_{\varphi \varphi }(g_{\theta \varphi })^2}|_{r=r_H}=4\pi (r_H^2+a^2).$$
(62)
As proved in Ref. the low-energy limit of the total absorption cross section for the minimally coupled scalar always equals to the horizon area in the asymptotically flat and spherically symmetric black hole, which is called ‘universality’. Although the general proof is not given yet, our numerical investigation supports the evidence that this universality seems to be extended to the non-spherically symmetric background.
In the high-energy limits the total absorption cross sections approach to the nonzero values which are roughly same with the low-energy limits. In the intermediate region the total absorption cross sections do not exhibit an oscillatory pattern, which seems to be the effect of the extra dimensions.
Fig. 4 shows the total absorption cross sections for the bulk scalar when $`b_{}=0`$, $`0.5`$, and $`1`$ with $`a_{}=0.5`$. The low-energy limits equal to the area of the non-spherically symmetric horizon hypersurface $`𝒜_{BL}`$ defined
$$𝒜_{BL}_0^{\frac{\pi }{2}}𝑑\theta _0^{2\pi }𝑑\varphi 𝑑\psi \sqrt{g_{\theta \theta }\left[g_{\varphi \varphi }g_{\psi \psi }\left(g_{\varphi \psi }\right)^2\right]}=2\pi ^2\frac{(r_H^2+a^2)(r_H^2+b^2)}{r_H}.$$
(63)
This result also supports that the universality in Ref. holds in the non-spherically symmetric background. In the high-energy limits the total absorption cross sections seem to approach to the nonzero values, which is much smaller than the low-energy limits. This fact indicates that unlike the non-rotating black hole case, the contribution of the higher partial waves except S-wave to the total absorption cross section is too much negligible.
Fig. 5 shows the emission rate $`\mathrm{\Gamma }_{BL}/d\omega `$ for the bulk scalar and $`\mathrm{\Gamma }_{BR}/d\omega `$ for the brane scalar together. For the brane scalar we choosed $`a=0.5`$ and $`b=0`$ while for the bulk scalar $`b`$ is chosen as $`0`$, $`0.5`$ and $`1`$ with $`a=0.5`$. The wiggly pattern in $`\mathrm{\Gamma }_{BR}/d\omega `$ indicates that unlike the non-rotating black holes the contribution of the higher partial waves is not negligible. This means that the effect of the superradiant scattering is crucially significant in the brane emission. This wiggly pattern disappears in $`\mathrm{\Gamma }_{BL}/d\omega `$, which means that the effect of the superradiance is negligible. For the bulk scalar, therefore, the contribution of S-wave to the emission rate is dominant like the case of the non-rotating black hole background. Integrating the plots in Fig. 5, we can compute the total emission rate. For the brane scalar the total emission rate is $`0.00353832`$ and for the bulk scalar $`0.000343955`$, $`0.000123524`$ and $`7.44114\times 10^6`$ for the cases of $`b=0`$, $`b=0.5`$ and $`b=1`$ respectively. Thus the emission rate for the bulk scalar is much smaller than that for the brane scalar. Thus the effect of the superradiance in $`5d`$ rotating black hole background does not seem to change the main conclusion of Ref., i.e. black holes radiate mainly on the brane.
We computed the absorption and emission spectra for the brane and bulk scalar fields when the spacetime is an $`5d`$ rotating black hole carrying the two different angular momentum parameters. Although the effect of the superradiant scattering is taken into account, the main conclusion of Ref. does not seem to be changed. This is due to the fact that the energy amplification for the bulk scalar is order of $`10^9\%`$ while that for the brane scalar is order of unity. It seems to be straightforward to extend our calculation to the $`6d`$ rotating black hole background. It is of interest to check explicitly whether or not the effect of the superradiance is negligible in $`6d`$ case.
It is well-known that the Hawking radiation is highly dependent on the spin of the field. Thus, it seems to be greatly important to take the effect of spin into account in the higher-dimensional rotating black hole background. This is in progress and will be reported elsewhere.
Acknowledgement: This work was supported by the Korea Research Foundation under Grant (KRF-2003-015-C00109).
Appendix
The $`4d`$ Kerr metric is well-known in the form
$`ds^2=`$ $``$ $`(1{\displaystyle \frac{\mu r}{\mathrm{\Sigma }}})dt^2{\displaystyle \frac{2a\mu r\mathrm{sin}^2\theta }{\mathrm{\Sigma }}}dtd\varphi +{\displaystyle \frac{\mathrm{\Sigma }}{\mathrm{\Delta }}}dr^2+\mathrm{\Sigma }d\theta ^2`$ (A.1)
$`+`$ $`(r^2+a^2+{\displaystyle \frac{a^2\mu r\mathrm{sin}^2\theta }{\mathrm{\Sigma }}})\mathrm{sin}^2\theta d\varphi ^2`$ (A.2)
where $`\mathrm{\Delta }=r^2\mu r+a^2(rr_+)(rr_{})`$, $`\mathrm{\Sigma }=r^2+a^2\mathrm{cos}^2\theta `$ and $`r_\pm =(\mu \pm \sqrt{\mu ^24a^2})/2`$ are the inner and outer horizons. Then it is not difficult to show that the scalar wave equation $`\mathrm{}\mathrm{\Phi }=0`$ in this background is separable. The radial and angular equations of this wave equation reduce to
$`\mathrm{\Delta }{\displaystyle \frac{d}{dr}}\mathrm{\Delta }{\displaystyle \frac{dR}{dr}}+[(r^2w+a^2wam)^2\mathrm{\Delta }\mathrm{\Lambda }_{\mathrm{}}^m]R=0`$ (A.3)
$`{\displaystyle \frac{1}{\mathrm{sin}\theta }}{\displaystyle \frac{d}{d\theta }}(\mathrm{sin}\theta {\displaystyle \frac{dT}{d\theta }})+[{\displaystyle \frac{m^2}{\mathrm{sin}^2\theta }}+a^2w^2\mathrm{cos}^2\theta +E_{\mathrm{}}^m]T=0`$ (A.4)
where $`\mathrm{\Phi }=R(r)T(\theta )e^{im\varphi }e^{iwt}`$, $`\mathrm{\Lambda }_{\mathrm{}}^m=E_{\mathrm{}}^m+a^2w^22amw`$ and $`E_{\mathrm{}}^m`$ is a separation constant. Defining $`x=wr`$ and $`x_\pm =wr_\pm `$, one can show that the radial equation becomes
$`(xx_+)(x`$ $``$ $`x_{}){\displaystyle \frac{d}{dx}}(xx_+)(xx_{}){\displaystyle \frac{dR}{dx}}`$ (A.5)
$`+`$ $`[(x^2+a^2w^2amw)^2\mathrm{\Lambda }_{\mathrm{}}^m(xx_+)(xx_{})]R=0.`$ (A.6)
Solving the radial equation as a series form, one can derive the near-horizon solution
$$𝒢_{\mathrm{}}^m(x,x_+,x_{})=e^{\lambda _4ln|xx_+|}\underset{n=0}{\overset{\mathrm{}}{}}d_{\mathrm{},n}^m(xx_+)^n$$
(A.7)
and the asymptotic solution
$$_{\mathrm{}(\pm )}^m(x,x_+,x_{})=(\pm i)^{\mathrm{}+1}e^{ix}(xx_+)^{\pm \lambda _4}\underset{n=0}{\overset{\mathrm{}}{}}\tau _{n(\pm )}x^{(n+1)}$$
(A.8)
where the recursion relations for $`d_{\mathrm{},n}^m`$ and $`\tau _{n(\pm )}`$ can be explicitly derived by inserting (A.7) and (A.8) into (A.5). The factor $`\lambda _4`$ arises due to the regular singular nature of the radial equation and its explicit expression is
$$\lambda _4=i\frac{w(r_+^2+a^2)(wm\mathrm{\Omega })}{x_+x_{}}$$
(A.9)
where $`\mathrm{\Omega }=a/(r_+^2+a^2)`$ is an angular frequency of the black hole. If $`0<w<m\mathrm{\Omega }`$, $`Im\lambda _4`$ becomes negative which indicates that the near-horizon solution (A.7) becomes outgoing wave. Thus the superradiant scattering takes place under the condition $`0<w<m\mathrm{\Omega }`$.
The energy amplification arising due to the superradiant scattering was computed in Ref. when the spacetime background is a maximally rotating ($`a_{}a/r_+=1`$) Kerr black hole. The authors in Ref. solved the radial and angular equations (A.3) directly by adopting the different numerical technique. For our case, however, the near-horizon and asymptotic solutions (A.7) and (A.8) are used. Thus our numerical method cannot be applied to the case of the maximally rotating black hole because $`a_{}=1`$ implies the extremal limit, i.e. $`r_+=r_{}`$ and $`\lambda _4`$ goes to infinity in this limit. Applying the numerical method used in the present paper the energy amplification can be straightforwardly computed for $`a_{}<1`$.
Fig. 6 is a $`\mathrm{log}`$-plot of the energy amplification when $`a_{}=0.8`$. Although Fig. 6 is different from Fig. 1 of Ref. due to the different choice of $`a_{}`$, it indicates that the energy amplification of the scalar wave in the $`4d`$ Kerr black hole is order of $`10^1\%`$ like the case of $`a_{}=1`$. |
warning/0506/hep-th0506240.html | ar5iv | text | # KEK-TH-1024hep-th/0506240 Marginal Deformations and Classical Solutions in Open Superstring Field Theory
## 1 Introduction
String field theory is established as a framework for exploring nonperturbative structures in string theory. Motivated by Sen’s conjecture , many people studied classical solutions extensively and intriguing results were provided in bosonic open string field theory . In the supersymmetric case, the most promising theory is formulated in terms of the Wess-Zumino-Witten (WZW) like action proposed by Berkovits , which has no problem with contact term divergences . This open superstring field theory is a sufficient framework to elucidate nonperturbative phenomena of the Neveu-Schwarz (NS) sector. Indeed, the tachyon vacuum and kink solutions were found in the superstring field theory by using the level truncation scheme . On the analytical side, there are some attempts to construct exact solutions in terms of a half string formulation , a pregeometrical formulation , a conjecture of vacuum superstring field theory and an analogy with integrable systems .
In the present paper we construct analytic classical solutions in the open superstring field theory using techniques developed in bosonic open string field theory . The resulting solution consists of the identity string field, ghost fields and an operator associated with a current. Taking $`X(z)`$ as the current, we find that the action expanded around the solution can be transformed to the original action by a string field redefinition. In the redefined theory, however, the momentum is shifted in the string field and then the classical solution can be related to a background Wilson line. Generically, we anticipate that our solutions correspond to marginal boundary deformed backgrounds as in the bosonic case.
The analytic solutions are useful for studying gauge structure in string field theory. In bosonic string field theory, the analytic solution corresponding to Wilson lines can be represented as a “locally” pure gauge, and then we find that a “locally” pure gauge configuration in string field theory corresponds to a marginal deformation in conformal field theory . This correspondence is a natural generalization of that of low energy effective theories. Later we will see that the solution in the superstring field theory shares this feature of the bosonic theory.
Marginal deformations in string field theory were often studied using the level truncation scheme. We see that the effective potential for a marginal field becomes flatter as the truncation level is increased , and then the vacuum energy of the analytic solution must vanish. Unfortunately, we encounter a difficulty in calculating the vacuum energy in the bosonic theory. Though the vacuum energy formally vanishes, it is given as a kind of indefinite quantities if we calculate it by oscillator representation . However, we will see that the vacuum energy is to be exactly zero in the superstring field theory. This result is a characteristic feature of the supersymmetric case.
In string field theory, the gauge symmetry includes global symmetries generated by $`K_n=L_n(1)^nL_n`$ . It is a typical symmetry in string field theory because the symmetry mixes various component fields and it has a non-local structure. Based on an analytical approach, we find that the Wilson line parameter in the solution is invariant under the global transformation.
Although it is hard to include the Ramond (R) sector into the action, we have the equations of motion for both of the NS and R sectors . The equations of motion possess a fermionic symmetry which transforms the NS boson (R fermion) to the R fermion (NS boson). Then we expect that the superstring field theory has space-time supersymmetry. Actually, we find that global space-time supersymmetry is realized on-shell as a part of the fermionic symmetry. We show that the solution corresponding to a Wilson line is a supersymmetric solution, namely the solution is invariant under the global space-time supersymmetry transformation.
This paper is organized as follows. In section 2 we construct an analytic classical solution in the open superstring field theory. We find that the solution can be written by a well-defined Fock space expression. After discussing the vacuum energy and the theory expanded around the solution, we relate the solution to background Wilson lines. Michishita gives a covariant action of the R sector by imposing a constraint equation . We discuss the effect of the solution on the R sector in terms of the action proposed by Michishita. Moreover, we show how the solution is transformed under the global symmetry and space-time supersymmetry. In section 3 we extend the Wilson line solution to those which correspond to general marginal boundary deformations generated by supercurrents. We find that generalized solutions also have a favorite feature that the vacuum energy vanishes. We offer some comments related to our results and discuss open questions in section 4. In addition we include four appendices. We represent the identity string field by explicit oscillator expression in the large Hilbert space in appendix A. In appendix B we give a different derivation of the action expanded around a general solution which is originally given in refs. . We use an alternative expression of the action given in ref. to derive the expanded action. In appendix C we show that the fermionic symmetry contains global space-time supersymmetry, and give some comments on supersymmetry in the cubic superstring field theory and its modified theory . In appendix D we construct the analytic solution in bosonic string field theory which corresponds to a general marginal boundary deformations.
## 2 Classical solutions and background Wilson lines
The open superstring field theory action is given by
$`S[\mathrm{\Phi }]={\displaystyle \frac{1}{2g^2}}<<(e^\mathrm{\Phi }Q_\mathrm{B}e^\mathrm{\Phi })(e^\mathrm{\Phi }\eta _0e^\mathrm{\Phi }){\displaystyle _0^1}𝑑t(e^{t\mathrm{\Phi }}_te^{t\mathrm{\Phi }})\{(e^{t\mathrm{\Phi }}Q_\mathrm{B}e^{t\mathrm{\Phi }}),(e^{t\mathrm{\Phi }}\eta _0e^{t\mathrm{\Phi }})\}>>,`$ (2.1)
where $`\mathrm{\Phi }`$ denotes a string field of the GSO(+) NS sector which corresponds to a Grassmann even vertex operator of ghost number 0 and picture number 0 in the conformal field theory. CFT correlators $`\mathrm{}`$ are defined in the large Hilbert space and $`\{A,B\}AB+BA`$.<sup>1</sup><sup>1</sup>1For details of the definition, see for instance ref. . The action is invariant under the infinitesimal gauge transformation,
$`\delta e^\mathrm{\Phi }=(Q_\mathrm{B}\delta \mathrm{\Lambda })e^\mathrm{\Phi }+e^\mathrm{\Phi }(\eta _0\delta \mathrm{\Lambda }^{}),`$ (2.2)
where $`\delta \mathrm{\Lambda }`$ and $`\delta \mathrm{\Lambda }^{}`$ are infinitesimal parameters. Integrating this infinitesimal form, we can obtain the finite gauge transformation as<sup>2</sup><sup>2</sup>2The gauge transformation can be expressed as $`e^\mathrm{\Phi }^{}=ge^\mathrm{\Phi }g^{}`$ where $`Q_\mathrm{B}g=\eta _0g^{}=0`$, since each of the operators, $`Q_\mathrm{B}`$ and $`\eta _0`$, has trivial cohomology in the large Hilbert space. Here, $`g`$ and $`g^{}`$ are group elements in the “stringy gauge group” in superstring field theory.
$`e^\mathrm{\Phi }^{}=e^{Q_\mathrm{B}\mathrm{\Lambda }}e^\mathrm{\Phi }e^{\eta _0\mathrm{\Lambda }^{}},`$ (2.3)
where $`\mathrm{\Lambda }`$ and $`\mathrm{\Lambda }^{}`$ are finite parameters. Variating the action (2.1), we can derive the equation of motion to be
$`\eta _0(e^\mathrm{\Phi }Q_\mathrm{B}e^\mathrm{\Phi })=0.`$ (2.4)
For simplicity, we mainly consider superstring field theory describing the dynamics of a D brane without Chan-Paton degrees of freedom. We single out a direction on the world volume of the brane, writing the string coordinate as $`X(z,\overline{z})=(X(z)+X(\overline{z}))/2`$ and its supersymmetric partner as $`\psi (z)`$. Our later analysis can be easily extended to include Chan-Paton indices.
### 2.1 classical solutions in open superstring field theory
In this subsection, we will show that one of the classical solutions is given by
$`\mathrm{\Phi }_0=\stackrel{~}{V}_L(F)I,`$ (2.5)
where $`I`$ is the identity string field and the operator $`\stackrel{~}{V}_L`$ is defined as<sup>3</sup><sup>3</sup>3We note that $`e^{q\varphi }`$ ($`q`$ : odd) is a fermionic operator. More precisely, we need a cocycle factor to represent statistical property of the operator.
$`\stackrel{~}{V}_L(F)={\displaystyle _{C_{\mathrm{left}}}}{\displaystyle \frac{dz}{2\pi i}}F(z)\stackrel{~}{v}(z),\stackrel{~}{v}(z)={\displaystyle \frac{1}{\sqrt{2}}}c\xi e^\varphi \psi (z).`$ (2.6)
Here, $`C_{\mathrm{left}}`$ denotes a counter-clockwise path along a half of the unit circle, i.e., $`\pi /2<\sigma <\pi /2`$ for $`z=e^{i\sigma }`$. $`F(z)`$ is a function on the unit circle $`\left|z\right|=1`$ satisfying $`F(1/z)=z^2F(z)`$ .<sup>4</sup><sup>4</sup>4Under this condition, $`F(z)`$ cannot be a non-zero constant. We have to impose an additional constraint on $`F(z)`$ due to the reality of the string field as explained in the next subsection (see also appendix A).
First, we introduce half string operators similar to $`\stackrel{~}{V}_L(F)`$:
$`C_L(F)`$ $`=`$ $`{\displaystyle _{C_{\mathrm{left}}}}{\displaystyle \frac{dz}{2\pi i}}F(z)c(z),`$ (2.7)
$`V_L(F)`$ $`=`$ $`{\displaystyle _{C_{\mathrm{left}}}}{\displaystyle \frac{dz}{2\pi i}}F(z)v(z),`$ (2.8)
where $`c(z)`$ is the ghost field and $`v(z)`$ is defined as
$`v(z)=[Q_\mathrm{B},\stackrel{~}{v}(z)]={\displaystyle \frac{i}{2\sqrt{\alpha ^{}}}}cX(z)+{\displaystyle \frac{1}{\sqrt{2}}}\eta e^\varphi \psi (z).`$ (2.9)
By definition, the commutation relation $`[Q_\mathrm{B},\stackrel{~}{V}_L(F)]=V_L(F)`$ holds. We also define the operators, $`C_R(F)`$, $`V_R(F)`$ and $`\stackrel{~}{V}_R(F)`$ by replacing the integration path $`C_{\mathrm{left}}`$ by $`C_{\mathrm{right}}`$ which rotates counter-clockwise along $`\left|z\right|=1(\mathrm{Re}z0)`$. For these half string operators, we can derive their (anti-)commutation relations from similar procedures in refs. . The operator product expansions (OPEs) among local operators in the integrand are easily calculated as
$`v(z)\stackrel{~}{v}(z^{})`$ $``$ $`{\displaystyle \frac{1}{zz^{}}}{\displaystyle \frac{1}{2}}c(z^{}),`$ (2.10)
$`v(z)v(z^{})`$ $``$ $`{\displaystyle \frac{1}{zz^{}}}{\displaystyle \frac{1}{2}}(cc(z^{})\eta \eta e^{2\varphi }(z^{})).`$ (2.11)
Using these OPEs, we obtain (anti-)commutation relations of these local operators on the unit circle, $`\left|z\right|=1`$:
$`[v(z),\stackrel{~}{v}(z^{})]`$ $`=`$ $`{\displaystyle \frac{1}{2}}c(z^{})\delta (z,z^{}),`$ (2.12)
$`\{v(z),v(z^{})\}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(cc(z^{})\eta \eta e^{2\varphi }(z^{}))\delta (z,z^{}),`$ (2.13)
where the delta function is defined as $`\delta (z,z^{})=_{n=\mathrm{}}^{\mathrm{}}z^nz_{}^{}{}_{}{}^{n1}`$.<sup>5</sup><sup>5</sup>5The delta function satisfies $`f(z)={\displaystyle _{\left|z^{}\right|=1}}{\displaystyle \frac{dz^{}}{2\pi i}}f(z^{})\delta (z,z^{}),`$ (2.14) where the function $`f(z)`$ is square integrable on the unit circle $`\left|z\right|=1`$ $`(fL^2)`$. Moreover, for any $`f,gL^2`$, the delta function enjoys the properties, $`{\displaystyle _{C_{\mathrm{left}/\mathrm{right}}}}{\displaystyle \frac{dz}{2\pi i}}{\displaystyle _{C_{\mathrm{left}/\mathrm{right}}}}{\displaystyle \frac{dz^{}}{2\pi i}}f(z)g(z^{})\delta (z,z^{})={\displaystyle _{C_{\mathrm{left}/\mathrm{right}}}}{\displaystyle \frac{dz}{2\pi i}}f(z)g(z),`$ (2.15) $`{\displaystyle _{C_{\mathrm{left}}}}{\displaystyle \frac{dz}{2\pi i}}{\displaystyle _{C_{\mathrm{right}}}}{\displaystyle \frac{dz^{}}{2\pi i}}f(z)g(z^{})\delta (z,z^{})=0,`$ (2.16) which are necessary for derivation of (2.17) and (2.18).
We integrate (2.12) and (2.13) to derive (anti-)commutation relations of half string operators:
$`[V_L(F),\stackrel{~}{V}_L(G)]`$ $`=`$ $`{\displaystyle \frac{1}{2}}C_L(FG),`$ (2.17)
$`\{V_L(F),V_L(G)\}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\{Q_\mathrm{B},C_L(FG)\},`$ (2.18)
where we have used $`\{Q_\mathrm{B},c(z)\}=cc(z)\eta \eta e^{2\varphi }(z)`$ in the latter equation. The similar relations hold for the right-half string operators, and other (anti-)commutation relations become zero. Here, we emphasize that these equations hold for any functions $`F(z)`$ and $`G(z)`$ defined on the unit circle, because we have only to use the equal-time (anti-)commutation relations to derive the equations.<sup>6</sup><sup>6</sup>6 To derive (2.17) and (2.18), it is sufficient for $`F`$ and $`G`$ to be square integrable. The condition $`F(\pm i)=G(\pm i)=0`$ is unnecessary for these (anti-)commutation relations. Namely, the functions are not necessarily to be holomorphic, though we express them as functions of a complex variable.
Next, we consider some properties of the half string operators associated with the star product and the identity string field. Suppose that two string fields $`A`$ and $`B`$ are defined as $`|A=A(0)|0`$ and $`|B=B(0)|0`$, where $`A(z)`$ and $`B(z^{})`$ are conformal fields on the unit discs $`\left|z\right|1`$ and $`\left|z^{}\right|1`$, respectively. The star product $`AB`$ is defined in terms of the gluing Riemann surface by the identification $`zz^{}=1(\left|z\right|=1,\mathrm{Re}z0)`$ . Accordingly, it follows that
$`\left(\sigma (z)A\right)B=(1)^{\left|\sigma \right|\left|A\right|}A\left((z^2)^h\sigma (z^{})B\right)(zz^{}=1,\left|z\right|=1,\mathrm{Re}z0),`$ (2.19)
where $`\sigma (z)`$ is a primary field with dimension $`h`$, and $`\left|\alpha \right|`$ denotes the statistic index defined to be $`0(1)`$ if $`\alpha `$ is a bosonic (fermionic) operator. Multiplying a function $`F(z)`$ which satisfies $`F(1/z)=(z^2)^{1h}F(z)`$ <sup>7</sup><sup>7</sup>7$`F(z)`$ is the same as the function $`F_+^{(h+1)}(z)`$ in ref. . We note that our analysis is easily extended to the case of $`F_{}(z)`$ in ref. . to both sides of (2.19) and integrating it along the path $`\left|z\right|=1(\mathrm{Re}z0)`$, we can find the generic formula
$`\left(\mathrm{\Sigma }_R(F)A\right)B=(1)^{\left|\sigma \right|\left|A\right|}A\left(\mathrm{\Sigma }_L(F)B\right),`$ (2.20)
where the operator $`\mathrm{\Sigma }_{L(R)}`$ is defined as
$`\mathrm{\Sigma }_{L(R)}(F)`$ $`=`$ $`{\displaystyle _{C_{\mathrm{left}(\mathrm{right})}}}{\displaystyle \frac{dz}{2\pi i}}F(z)\sigma (z).`$ (2.21)
Similarly, we can obtain a generic formula associated with the identity string field $`I`$:
$`\mathrm{\Sigma }_L(F)I=\mathrm{\Sigma }_R(F)I.`$ (2.22)
If we choose $`c(z)`$, $`v(z)`$ and $`\stackrel{~}{v}(z)`$ as the primary field,<sup>8</sup><sup>8</sup>8The dimensions of $`c(z)`$, $`v(z)`$ and $`\stackrel{~}{v}(z)`$ are $`1`$, $`0`$ and $`0`$, respectively. we can derive the following equations from the generic formulae:
$`\left(C_R(F_1)A\right)B`$ $`=`$ $`(1)^{\left|A\right|}A\left(C_L(F_1)B\right),`$ (2.23)
$`\left(V_R(F_0)A\right)B`$ $`=`$ $`(1)^{\left|A\right|}A\left(V_L(F_0)B\right),`$ (2.24)
$`\left(\stackrel{~}{V}_R(F_0)A\right)B`$ $`=`$ $`A\left(\stackrel{~}{V}_L(F_0)B\right),`$ (2.25)
$`C_R(F_1)I`$ $`=`$ $`C_L(F_1)I,`$ (2.26)
$`V_R(F_0)I`$ $`=`$ $`V_L(F_0)I,`$ (2.27)
$`\stackrel{~}{V}_R(F_0)I`$ $`=`$ $`\stackrel{~}{V}_L(F_0)I,`$ (2.28)
where the function $`F_h(z)`$ satisfies $`F_h(1/z)=(z^2)^{1h}F_h(z)`$. Again, these equations hold if $`F_h(z)`$ is defined on the unit circle $`\left|z\right|=1`$. As in the case of eqs. (2.17) and (2.18), the function does not need holomorphicity. Here, the function in eq. (2.26) should behave like $`F_1(z)\mathrm{O}((zz_0)^ϵ),(ϵ>0)`$ near $`z_0=\pm i`$ in order that the state $`C_{L(R)}(F_1)I`$ has a well-defined Fock space expression. Because the ghost field $`c(z)`$ has a single pole at $`z=\pm i`$ on the identity state as seen in the next subsection . This condition is assured by imposing $`F_1(\pm i)=0`$ if the function is expandable in a Taylor series.<sup>9</sup><sup>9</sup>9Actually, we can expand the function as $`F_1(z)=O((zz_0)^2)`$ if $`F_1(\pm i)=0`$ holds.
Now, it can be easily shown that $`\mathrm{\Phi }_0`$ given by (2.5) is a classical solution:
$`e^{\mathrm{\Phi }_0}Q_\mathrm{B}e^{\mathrm{\Phi }_0}`$ $`=`$ $`e^{\stackrel{~}{V}_L(F)}Q_\mathrm{B}e^{\stackrel{~}{V}_L(F)}I`$ (2.29)
$`=`$ $`Q_\mathrm{B}I+[\stackrel{~}{V}_L(F),Q_\mathrm{B}]I+{\displaystyle \frac{1}{2!}}[\stackrel{~}{V}_L(F),[\stackrel{~}{V}_L(F),Q_\mathrm{B}]]I+\mathrm{}`$ (2.30)
$`=`$ $`V_L(F)I+{\displaystyle \frac{1}{4}}C_L(F^2)I,`$ (2.31)
where we have used (2.28) and $`[\stackrel{~}{V}_L(F),\stackrel{~}{V}_R(F)]=0`$ in the first equality, and eq. (2.17) in the last equality. We should note that the state $`C_L(F^2)I`$ is well-defined because $`F^2(\pm i)=0`$ due to $`F(1/z)=z^2F(z)`$. The $`\xi `$ zero mode is not contained in both operators $`V_L(F)`$ and $`C_L(F^2)`$ and the identity string field satisfies $`\eta _0I=0`$. As a result, we find that $`\eta _0(e^{\mathrm{\Phi }_0}Q_\mathrm{B}e^{\mathrm{\Phi }_0})=0`$ and then $`\mathrm{\Phi }_0`$ is a solution in open superstring field theory.
### 2.2 Fock space expressions
The operator $`\stackrel{~}{V}_L(F)`$ in the solution can be expressed in terms of integration with respect to $`\sigma `$ ($`z=e^{i\sigma }`$ on the contour $`C_{\mathrm{left}}`$) :
$`\stackrel{~}{V}_L(F)`$ $`=`$ $`{\displaystyle _{\frac{\pi }{2}}^{\frac{\pi }{2}}}{\displaystyle \frac{d\sigma }{2\pi }}{\displaystyle \frac{e^{i\sigma }F(e^{i\sigma })}{\sqrt{2}}}c(i\sigma )\xi (i\sigma ):e^{\varphi (i\sigma )}:\psi (i\sigma ),`$ (2.32)
where operators in the integrand are given by oscillator expansions:
$`c(i\sigma )={\displaystyle \underset{n𝐙}{}}c_ne^{in\sigma },\xi (i\sigma )={\displaystyle \underset{n𝐙}{}}\xi _ne^{in\sigma },\psi (i\sigma )={\displaystyle \underset{r𝐙+\frac{1}{2}}{}}\psi _re^{ir\sigma },`$ (2.33)
and $`:e^{\varphi (i\sigma )}:`$ is expanded in eq. (A.34). Using formulae: eqs. (A.28), (A.29) and (A.30) and computing straightforwardly, they are expressed on the identity state $`|I`$ as
$`c(i\sigma )|I`$ $`=`$ $`[c_1(2\mathrm{cos}\sigma )^1+c_0i\mathrm{tan}\sigma +c_1(1+(2\mathrm{cos}\sigma )^1)`$ (2.34)
$`+2{\displaystyle \underset{k1}{}}(c_{2k}i\mathrm{sin}2k\sigma +c_{(2k+1)}\mathrm{cos}(2k+1)\sigma )]|I,`$
$`\xi (i\sigma )|I`$ $`=`$ $`\left[\xi _0+2{\displaystyle \underset{k1}{}}\left(\xi _{2k}\mathrm{cos}2k\sigma +\xi _{(2k1)}i\mathrm{sin}(2k1)\sigma \right)\right]|I,`$ (2.35)
$`:e^{\varphi (i\sigma )}:|I`$ $`=`$ $`(2\mathrm{cos}\sigma )^{\frac{1}{2}}e^{_{k1}\left(\frac{\mathrm{cos}2k\sigma }{k}j_{2k}+\frac{2i\mathrm{sin}(2k1)\sigma }{2k1}j_{(2k1)}\right)}e^{\widehat{\varphi }_0}|I,`$ (2.36)
$`\psi (i\sigma )|I`$ $`=`$ $`{\displaystyle \underset{r,s1/2}{}}\left(\delta _{r,s}e^{is\sigma }+e^{ir\sigma }I_{rs}\right)\psi _s|I`$ (2.37)
$`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}(2\mathrm{cos}\sigma )^{\frac{1}{2}}{\displaystyle _0}{\displaystyle \frac{dz}{2\pi i}}z^{n1}{\displaystyle \frac{\sqrt{1+z^2}}{(1e^{i\sigma }z)(1+e^{i\sigma }z)}}\psi _{(n+\frac{1}{2})}|I`$
$`=`$ $`{\displaystyle \frac{\sqrt{2}}{(\mathrm{cos}\sigma )^{\frac{1}{2}}}}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}[\psi _{(2k+\frac{1}{2})}{\displaystyle \underset{m=0}{\overset{k}{}}}{\displaystyle \frac{(1)^{km+1}(2(km))!}{2^{2(km)}((km)!)^2(2(km)1)}}\mathrm{cos}(2m+1)\sigma `$
$`+\psi _{(2k+\frac{3}{2})}{\displaystyle \underset{m=0}{\overset{k}{}}}{\displaystyle \frac{(1)^{km+1}(2(km))!}{2^{2(km)}((km)!)^2(2(km)1)}}i\mathrm{sin}2(m+1)\sigma ]|I.`$
In computing $`\psi (i\sigma )|I`$, we have used a relation for the Neumann coefficients:
$`{\displaystyle \underset{r,s\frac{1}{2}}{}}I_{rs}y^{r\frac{1}{2}}z^{s\frac{1}{2}}={\displaystyle \frac{(h_I^{}(y))^{\frac{1}{2}}(h_I^{}(z))^{\frac{1}{2}}}{h_I(y)h_I(z)}}{\displaystyle \frac{1}{yz}}={\displaystyle \frac{\sqrt{1+y^2}\sqrt{1+z^2}}{(yz)(1+yz)}}{\displaystyle \frac{1}{yz}},`$ (2.38)
where $`h_I(z)=2z/(1z^2)`$ is the conformal map for the identity string field (see appendix A). Using the above expressions and the reality $`\mathrm{bpz}(|I)=(|I)^{}`$, we find a relation between the BPZ and hermitian conjugations:
$`\mathrm{bpz}(c(i\sigma )\xi (i\sigma ):e^{\varphi (i\sigma )}:\psi (i\sigma )|I)=(c(i\sigma )\xi (i\sigma ):e^{\varphi (i\sigma )}:\psi (i\sigma )|I)^{},`$ (2.39)
where we take a convention: $`c_n^{}=c_n,\xi _n^{}=\xi _n,j_n^{}=j_n2\delta _{n,0},\psi _r^{}=\psi _r`$. As a result, the reality condition for our solution $`\mathrm{bpz}(|\mathrm{\Phi }_0)=(|\mathrm{\Phi }_0)^{}`$ imposes $`(e^{i\sigma }F(e^{i\sigma }))^{}=e^{i\sigma }F(e^{i\sigma })`$ for the coefficient function in the integrand of (2.32), which is expanded as
$`e^{i\sigma }F(e^{i\sigma })`$ $`=`$ $`{\displaystyle \underset{n𝐙}{}}F_ne^{in\sigma }={\displaystyle \underset{n1}{}}F_n(e^{in\sigma }(1)^ne^{in\sigma })`$ (2.40)
$`=`$ $`2{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}\left(F_{2k}i\mathrm{sin}2k\sigma +F_{(2k1)}\mathrm{cos}(2k1)\sigma \right).`$
Here we have used $`F_n=(1)^nF_n`$ in the second equality, which follows from the condition $`F(1/z)=z^2F(z)`$ for the classical solution. The reality condition implies that $`F_{2k}`$ should be real and $`F_{(2k1)}`$ should be pure imaginary.
Putting the above expansions together, we obtain the explicit Fock space expression of the classical solution as follows:
$`|\mathrm{\Phi }_0`$ $`=`$ $`\stackrel{~}{V}_L(F)|I`$
$`=`$ $`\sqrt{2}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle _{\frac{\pi }{2}}^{\frac{\pi }{2}}}{\displaystyle \frac{d\sigma }{\pi }}\left(F_{2n}i\mathrm{sin}2n\sigma +F_{(2n1)}\mathrm{cos}(2n1)\sigma \right)`$
$`\times [c_1(2\mathrm{cos}\sigma )^1+c_0i\mathrm{tan}\sigma +c_1(1+(2\mathrm{cos}\sigma )^1)`$
$`+2{\displaystyle \underset{m1}{}}(c_{2m}i\mathrm{sin}2m\sigma +c_{(2m+1)}\mathrm{cos}(2m+1)\sigma )]`$
$`\times \left[\xi _0+2{\displaystyle \underset{l1}{}}\left(\xi _{2l}\mathrm{cos}2l\sigma +\xi _{(2l1)}i\mathrm{sin}(2l1)\sigma \right)\right]`$
$`\times \mathrm{exp}\left[{\displaystyle \underset{p1}{}}\left({\displaystyle \frac{\mathrm{cos}2p\sigma }{p}}j_{2p}+{\displaystyle \frac{2i\mathrm{sin}(2p1)\sigma }{2p1}}j_{(2p1)}\right)\right]e^{\widehat{\varphi }_0}`$
$`\times {\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}[\psi _{(2k+\frac{1}{2})}{\displaystyle \underset{q=0}{\overset{k}{}}}{\displaystyle \frac{(1)^{kq}(2(kq))!}{2^{2(kq)}((kq)!)^2(2(kq)1)}}\mathrm{cos}(2q+1)\sigma `$
$`+\psi _{(2k+\frac{3}{2})}{\displaystyle \underset{q=0}{\overset{k}{}}}{\displaystyle \frac{(1)^{kq}(2(kq))!}{2^{2(kq)}((kq)!)^2(2(kq)1)}}i\mathrm{sin}2(q+1)\sigma ]|I.`$
In particular, the integration with respect to $`\sigma `$ gives finite coefficients for each term of the form $`F_nc_m\xi _lj_{p_1}^{n_1}\mathrm{}j_{p_N}^{n_N}\psi _s|I`$ because both $`_{\frac{\pi }{2}}^{\frac{\pi }{2}}\frac{d\sigma }{2\pi }\left|\frac{\mathrm{sin}2k\sigma }{\mathrm{cos}\sigma }\right|`$ and $`_{\frac{\pi }{2}}^{\frac{\pi }{2}}\frac{d\sigma }{2\pi }\left|\frac{\mathrm{cos}(2k1)\sigma }{\mathrm{cos}\sigma }\right|`$,<sup>10</sup><sup>10</sup>10 Notice that coefficients of each term can be estimated as $`M\left|_{\frac{\pi }{2}}^{\frac{\pi }{2}}\frac{d\sigma }{2\pi }\frac{\mathrm{sin}2k\sigma }{\mathrm{cos}\sigma }\mathrm{sin}^{m_1}n_1\sigma \mathrm{}\mathrm{cos}^{p_1}q_1\sigma \mathrm{}\right|<M_{\frac{\pi }{2}}^{\frac{\pi }{2}}\frac{d\sigma }{2\pi }\left|\frac{\mathrm{sin}2k\sigma }{\mathrm{cos}\sigma }\right|`$ or $`M^{}\left|_{\frac{\pi }{2}}^{\frac{\pi }{2}}\frac{d\sigma }{2\pi }\frac{\mathrm{cos}(2k1)\sigma }{\mathrm{cos}\sigma }\mathrm{sin}^{m_1^{}}n_1^{}\sigma \mathrm{}\mathrm{cos}^{p_1^{}}q_1^{}\sigma \mathrm{}\right|<M^{}_{\frac{\pi }{2}}^{\frac{\pi }{2}}\frac{d\sigma }{2\pi }\left|\frac{\mathrm{cos}(2k1)\sigma }{\mathrm{cos}\sigma }\right|`$, where $`M,M^{}`$ are some finite positive constants. in which the numerators and the denominators come from the coefficients $`e^{i\sigma }F(e^{i\sigma })`$ and zero mode of ghost $`c`$ (i.e., $`c_1,c_0,c_1`$) respectively, are finite. We note that a factor $`(\mathrm{cos}\sigma )^{\frac{1}{2}}`$ in (2.37) is canceled by a factor $`(\mathrm{cos}\sigma )^{\frac{1}{2}}`$ in (2.36). Therefore, one can construct well-behaved solutions in the sense that coefficients of all modes in the Fock space become finite by taking appropriate $`F(z)`$. A sufficient condition is that only finite number of $`F_n`$s have nonzero value.
More concretely, the lowest few terms of the solution are computed as
$`|\mathrm{\Phi }_0`$ $`=`$ $`{\displaystyle \frac{\sqrt{2}}{\pi }}({\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^nF_{(2n1)}}{2n1}}c_1\xi _0\psi _{\frac{1}{2}}`$
$`+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^n4nF_{2n}}{4n^21}}((c_0\xi _0+c_1\xi _1+c_1\xi _0j_1)\psi _{\frac{1}{2}}+c_1\xi _0\psi _{\frac{3}{2}})+\mathrm{})e^{\widehat{\varphi }_0}|I.`$
In the above explicit expression, the first term implies the condensation of the massless vector field because it is expanded as $`c_1\xi _0\psi e^{\widehat{\varphi }_0}|Ic\xi e^\varphi \psi (0)|0+\mathrm{}`$ and $`c\xi e^\varphi \psi `$ is the vertex operator for massless vector with zero momentum . This coefficient constant for the lowest level can be rewritten as $`\frac{\sqrt{2}}{\pi }_{n=1}^{\mathrm{}}\frac{(1)^nF_{(2n1)}}{2n1}=_{C_{\mathrm{left}}}\frac{dz}{2\pi i}\frac{F(z)}{\sqrt{2}}`$.
### 2.3 background Wilson lines
We found that the classical solution involves the condensation of the massless vector field. This result suggests that our solution is related to a background Wilson line. In this subsection, we will discuss the vacuum energy of the classical solution, the theory expanded around the solution, and other characteristic features of the solution. Accordingly, we will show that the solution corresponds to a background Wilson line.
In order to evaluate the vacuum energy, it is convenient to use an alternative expression for the action:
$`S[\mathrm{\Phi }]`$ $`=`$ $`{\displaystyle \frac{1}{g^2}}{\displaystyle _0^1}𝑑t<<(\eta _0e^{t\mathrm{\Phi }}_te^{t\mathrm{\Phi }})(e^{t\mathrm{\Phi }}Q_\mathrm{B}e^{t\mathrm{\Phi }})>>.`$ (2.43)
$`=`$ $`{\displaystyle \frac{1}{g^2}}{\displaystyle _0^1}𝑑t<<(\eta _0\mathrm{\Phi })(e^{t\mathrm{\Phi }}Q_\mathrm{B}e^{t\mathrm{\Phi }})>>.`$ (2.44)
The equivalence of the actions (2.1) and (2.43) is proved in ref. . In general, the state $`\eta _0\mathrm{\Phi }`$ has no $`\xi `$ zero mode. For the solution (2.5), it is easily seen that
$`e^{t\mathrm{\Phi }_0}Q_\mathrm{B}e^{t\mathrm{\Phi }_0}=tV_L(F)I+{\displaystyle \frac{t^2}{4}}C_L(F^2)I,`$ (2.45)
and then the state $`e^{t\mathrm{\Phi }_0}Q_\mathrm{B}e^{t\mathrm{\Phi }_0}`$ also does not contain the $`\xi `$ zero mode. As a result, we find that the integrand in (2.43) becomes zero for the classical solution since there is no $`\xi `$ zero mode in the correlation function of the integrand.<sup>11</sup><sup>11</sup>11In the large Hilbert space, correlation functions are normalized as to be $`cc^2c\xi e^{2\varphi }0`$. Hence we confirm that the vacuum energy of the solution vanishes due to the ghost charge non-conservation in the large Hilbert space.
Let us consider the expansion of the string field around the solution (2.5). Generally, if we expand the string field $`\mathrm{\Phi }`$ around a classical solution $`\mathrm{\Phi }_0`$ as $`e^\mathrm{\Phi }=e^{\mathrm{\Phi }_0}e^\mathrm{\Phi }^{}`$, the action becomes
$`S[\mathrm{\Phi }]`$ $`=`$ $`S[\mathrm{\Phi }_0]+S^{}[\mathrm{\Phi }^{}],`$ (2.46)
$`S^{}[\mathrm{\Phi }^{}]`$ $`=`$ $`{\displaystyle \frac{1}{2g^2}}<<(e^\mathrm{\Phi }^{}Q_\mathrm{B}^{}e^\mathrm{\Phi }^{})(e^\mathrm{\Phi }^{}\eta _0e^\mathrm{\Phi }^{})`$ (2.47)
$`{\displaystyle _0^1}𝑑t(e^{t\mathrm{\Phi }^{}}_te^{t\mathrm{\Phi }^{}})\{(e^{t\mathrm{\Phi }^{}}Q_\mathrm{B}^{}e^{t\mathrm{\Phi }^{}}),(e^{t\mathrm{\Phi }^{}}\eta _0e^{t\mathrm{\Phi }^{}})\}>>,`$
where $`S[\mathrm{\Phi }_0]`$ corresponds to the vacuum energy and $`S^{}[\mathrm{\Phi }^{}]`$ has the same form as the original action (2.1) except the kinetic operator $`Q_\mathrm{B}^{}`$, which is defined as
$`Q_\mathrm{B}^{}\mathrm{\Psi }=Q_\mathrm{B}\mathrm{\Psi }+A_0\mathrm{\Psi }(1)^{\left|\mathrm{\Psi }\right|}\mathrm{\Psi }A_0,A_0=e^{\mathrm{\Phi }_0}Q_\mathrm{B}e^{\mathrm{\Phi }_0}\mathrm{for}{}_{}{}^{}\mathrm{\Psi }.`$ (2.48)
A proof is given in appendix B. For the new BRS charge, nilpotency holds automatically but $`\{Q_\mathrm{B}^{},\eta _0\}=0`$ is satisfied owing to the equation of motion, $`\eta _0A_0=0`$. For the solution (2.5), we find $`S[\mathrm{\Phi }_0]=0`$ as evaluated above. Substituting (2.5) into (2.48) and using (2.23), (2.24), (2.26), (2.27) and (2.31), we can write the new BRS charge as<sup>12</sup><sup>12</sup>12We can check nilpotency of the new BRS charge in terms of (2.18).
$`Q_\mathrm{B}^{}=Q_\mathrm{B}(V_L(F)+V_R(F))+{\displaystyle \frac{1}{4}}(C_L(F^2)+C_R(F^2)).`$ (2.49)
Using (2.17), the new BRS charge is rewritten as a similarity transformation from the original BRS charge:
$`Q_\mathrm{B}^{}=e^{\stackrel{~}{V}_L(F)+\stackrel{~}{V}_R(F)}Q_\mathrm{B}e^{\stackrel{~}{V}_L(F)\stackrel{~}{V}_R(F)}.`$ (2.50)
Here, we introduce the following half string operators,
$`X_{L(R)}(F)`$ $`=`$ $`{\displaystyle _{C_{\mathrm{left}(\mathrm{right})}}}{\displaystyle \frac{dz}{2\pi i}}F(z)X(z),`$ (2.51)
$`\mathrm{\Omega }_{L(R)}(F)`$ $`=`$ $`{\displaystyle _{C_{\mathrm{left}(\mathrm{right})}}}{\displaystyle \frac{dz}{2\pi i}}F(z)ic\xi \xi e^{2\varphi }X(z).`$ (2.52)
Using a similar procedure in the previous subsection, we can obtain (anti-)commutation relations between these operators in terms of their OPEs:
$`[X_{L(R)}(F),V_{L(R)}(F)]`$ $`=`$ $`i\sqrt{\alpha ^{}}C_{L(R)}(F^2),`$ (2.53)
$`[Q_\mathrm{B},X_{L(R)}(F)]`$ $`=`$ $`i2\sqrt{\alpha ^{}}V_{L(R)}(F),`$ (2.54)
$`\{Q_\mathrm{B},\mathrm{\Omega }_{L(R)}(F)\}`$ $`=`$ $`2\sqrt{\alpha ^{}}\stackrel{~}{V}_{L(R)}(F)iX_{L(R)}(F).`$ (2.55)
If $`F(z)`$ satisfies $`F(1/z)=z^2F(z)`$, it follows from (2.20) and (2.22) that
$`(X_R(F)A)B=A(X_L(F)B),X_R(F)I=X_L(F)I,`$ (2.56)
$`(\mathrm{\Omega }_R(F)A)B=(1)^{\left|A\right|}A(\mathrm{\Omega }_L(F)B),\mathrm{\Omega }_R(F)I=\mathrm{\Omega }_L(F)I.`$ (2.57)
Precisely speaking, $`X(z)`$ is not a primary field and we can not apply the formula (2.20) for the case that $`\sigma (z)=X(z)`$. However, it is directly shown that the equation (2.19) holds for $`X(z)`$ and then we can derive the same formula in which $`X(z)`$ behaves like a primary field with dimension 0 on the string vertex. The same holds for the formula associated with the identity string field.
In the theory expanded around the solution (2.5), we redefine the string field $`\mathrm{\Phi }^{}`$ as
$`\mathrm{\Phi }^{}`$ $`=`$ $`\mathrm{exp}\left({\displaystyle \frac{i}{2\sqrt{\alpha ^{}}}}X_L(F)I\right)\mathrm{\Phi }^{\prime \prime }\mathrm{exp}\left({\displaystyle \frac{i}{2\sqrt{\alpha ^{}}}}X_L(F)I\right)`$ (2.58)
$`=`$ $`\mathrm{exp}\left({\displaystyle \frac{i}{2\sqrt{\alpha ^{}}}}(X_L(F)+X_R(F))\right)\mathrm{\Phi }^{\prime \prime }.`$
Under this redefinition, the action of $`\mathrm{\Phi }^{}`$ is transformed to the exactly same form as the original action, because $`Q_\mathrm{B}^{}`$ is transformed to $`Q_\mathrm{B}`$:
$`Q_\mathrm{B}`$ $`=`$ $`\mathrm{exp}\left({\displaystyle \frac{i}{2\sqrt{\alpha ^{}}}}(X_L(F)+X_R(F))\right)Q_\mathrm{B}^{}\mathrm{exp}\left({\displaystyle \frac{i}{2\sqrt{\alpha ^{}}}}(X_L(F)+X_R(F))\right),`$ (2.59)
where use has been made of (2.53) and (2.54). This equivalence between the original and expanded actions suggests that $`\mathrm{\Phi }_0`$ is a pure gauge solution. Actually, we can represent the solution as a pure gauge form by using (2.55), (2.56) and (2.57):
$`e^{\mathrm{\Phi }_0}=\mathrm{exp}\left\{Q_\mathrm{B}\left({\displaystyle \frac{1}{2\sqrt{\alpha ^{}}}}\mathrm{\Omega }_L(F)I\right)\right\}\mathrm{exp}\left\{\eta _0\left({\displaystyle \frac{i}{2\sqrt{\alpha ^{}}}}\xi _0X_L(F)I\right)\right\}.`$ (2.60)
However, this is merely a locally pure gauge expression, because the operator $`X_L(F)`$ contains the zero mode $`\widehat{x}`$ which can not be defined globally if the direction is compactified. Hence, the classical solution turns out to be non-trivial.
In order to clarify the physical meaning of the solution, let us consider the case involving the Chan-Paton factor represented with indices $`(i,j)`$. The string field redefinition (2.58) can be generalized to<sup>13</sup><sup>13</sup>13The function $`F_i(z)`$ corresponds to $`\lambda _iF(z)`$ in ref. .
$`\mathrm{\Phi }_{ij}^{}`$ $`=`$ $`\mathrm{exp}\left({\displaystyle \frac{i}{2\sqrt{\alpha ^{}}}}X_L(F_i)I\right)\mathrm{\Phi }_{ij}^{\prime \prime }\mathrm{exp}\left({\displaystyle \frac{i}{2\sqrt{\alpha ^{}}}}X_L(F_j)I\right)`$ (2.61)
$`=`$ $`\mathrm{exp}\left({\displaystyle \frac{i}{2\sqrt{\alpha ^{}}}}X_L(F_i)+{\displaystyle \frac{i}{2\sqrt{\alpha ^{}}}}X_R(F_j)\right)\mathrm{\Phi }_{ij}^{\prime \prime },`$
where we take no sum with respect to $`(i,j)`$. Noting
$`{\displaystyle _{C_{\mathrm{left}}}}{\displaystyle \frac{dz}{2\pi i}}F(z)={\displaystyle _{C_{\mathrm{right}}}}{\displaystyle \frac{dz}{2\pi i}}F(z)\mathrm{for}F(1/z)=z^2F(z),`$ (2.62)
it is expanded as
$`\mathrm{\Phi }_{ij}^{\prime \prime }=\mathrm{exp}\left({\displaystyle \frac{i}{2\sqrt{\alpha ^{}}}}(f_if_j)\widehat{x}+\mathrm{}\right)\mathrm{\Phi }_{ij}^{},f_i={\displaystyle _{C_{\mathrm{left}}}}{\displaystyle \frac{dz}{2\pi i}}F_i(z),`$ (2.63)
where the abbreviation on the exponent denotes nonzero mode dependence. Consequently, this string field redefinition causes shift of momentum $`pp(f_if_j)/(2\sqrt{\alpha ^{}})`$ which is the same effect by background Wilson lines as shown in bosonic string field theory .
Thus, we conclude that the classical solution (2.5) corresponds to background Wilson lines, because the vacuum energy vanishes and the solution is represented as locally pure gauge form, and the theory expanded around the solution involves the momentum shift as expected from Wilson lines.
### 2.4 Ramond sector and supersymmetry
The action of the Ramond sector proposed in ref. is:
$`S_F={\displaystyle \frac{1}{2g^2}}<<(Q_\mathrm{B}\mathrm{\Xi })e^\mathrm{\Phi }(\eta _0\mathrm{\Psi })e^\mathrm{\Phi }>>,`$ (2.64)
where $`\mathrm{\Psi }`$ is a string field of the GSO(+) R sector which has ghost number 0 and picture number $`1/2`$. $`\mathrm{\Xi }`$ carries ghost number 0 and picture number $`1/2`$. Both $`\mathrm{\Psi }`$ and $`\mathrm{\Xi }`$ are Grassmann even. The total action is given by adding $`S_F`$ to the NS action (2.1), and a constraint is imposed on string fields as
$`Q_\mathrm{B}\mathrm{\Xi }=e^\mathrm{\Phi }(\eta _0\mathrm{\Psi })e^\mathrm{\Phi }.`$ (2.65)
The total action is invariant under the infinitesimal gauge transformation
$`\delta e^\mathrm{\Phi }`$ $`=`$ $`(Q_\mathrm{B}\delta \mathrm{\Lambda }_0)e^\mathrm{\Phi }+e^\mathrm{\Phi }(\eta _0\delta \mathrm{\Lambda }_1),`$ (2.66)
$`\delta \mathrm{\Psi }`$ $`=`$ $`\eta _0\delta \mathrm{\Lambda }_{3/2}+\mathrm{\Psi }(\eta _0\delta \mathrm{\Lambda }_1)(\eta _0\delta \mathrm{\Lambda }_1)\mathrm{\Psi },`$ (2.67)
$`\delta \mathrm{\Xi }`$ $`=`$ $`Q_\mathrm{B}\delta \mathrm{\Lambda }_{1/2}+(Q_\mathrm{B}\delta \mathrm{\Lambda }_0)\mathrm{\Xi }\mathrm{\Xi }(Q_\mathrm{B}\delta \mathrm{\Lambda }_0),`$ (2.68)
where $`\delta \mathrm{\Lambda }_P`$ denotes an infinitesimal parameter with the picture number $`P`$. The constraint (2.65) is unchanged under the transformation.
Variating the total action, we can derive the equations of motion to be
$`\eta _0(e^\mathrm{\Phi }Q_\mathrm{B}e^\mathrm{\Phi })`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\eta _0\mathrm{\Psi })e^\mathrm{\Phi }(Q_\mathrm{B}\mathrm{\Xi })e^\mathrm{\Phi }{\displaystyle \frac{1}{2}}e^\mathrm{\Phi }(Q_\mathrm{B}\mathrm{\Xi })e^\mathrm{\Phi }(\eta _0\mathrm{\Psi }),`$ (2.69)
$`\eta _0(e^\mathrm{\Phi }(Q_\mathrm{B}\mathrm{\Xi })e^\mathrm{\Phi })`$ $`=`$ $`0,`$ (2.70)
$`Q_\mathrm{B}(e^\mathrm{\Phi }(\eta _0\mathrm{\Psi })e^\mathrm{\Phi })`$ $`=`$ $`0.`$ (2.71)
Substituting the constraint (2.65) into these equations, we can obtain the equations of motion for the NS and R sectors, $`\eta _0(e^\mathrm{\Phi }Q_\mathrm{B}e^\mathrm{\Phi })=(\eta _0\mathrm{\Psi })^2`$ and $`Q_\mathrm{B}(e^\mathrm{\Phi }(\eta _0\mathrm{\Psi })e^\mathrm{\Phi })=0`$ .
Let us expand the string fields around a classical solution $`(\mathrm{\Phi },\mathrm{\Psi })=(\mathrm{\Phi }_0,\mathrm{\hspace{0.17em}0})`$ as ($`e^\mathrm{\Phi },\mathrm{\Psi })=(e^{\mathrm{\Phi }_0}e^\mathrm{\Phi }^{},\mathrm{\Psi }^{})`$. Then, the action of the R sector becomes
$`S_F={\displaystyle \frac{1}{2g^2}}<<(Q_\mathrm{B}^{}\mathrm{\Xi }^{})e^\mathrm{\Phi }^{}(\eta _0\mathrm{\Psi }^{})e^\mathrm{\Phi }^{}>>,`$ (2.72)
and the constraint is changed to
$`Q_\mathrm{B}^{}\mathrm{\Xi }^{}=e^\mathrm{\Phi }^{}(\eta _0\mathrm{\Psi }^{})e^\mathrm{\Phi }^{},`$ (2.73)
where $`Q_\mathrm{B}^{}`$ is the new BRS operator defined as (2.48) and $`\mathrm{\Xi }`$ is a superfluous string field redefined as $`\mathrm{\Xi }^{}=e^{\mathrm{\Phi }_0}\mathrm{\Xi }e^{\mathrm{\Phi }_0}`$. Like the NS sector, the expanded action and constraint in the R sector have the same structure as the original ones except the form of the BRS operator.
Now, we take the solution (2.5) as $`\mathrm{\Phi }_0`$ in the above expansion. In the expanded action, we redefine the string fields as
$`\mathrm{\Phi }^{}`$ $`=`$ $`\mathrm{exp}\left({\displaystyle \frac{i}{2\sqrt{\alpha ^{}}}}X_L(F)I\right)\mathrm{\Phi }^{\prime \prime }\mathrm{exp}\left({\displaystyle \frac{i}{2\sqrt{\alpha ^{}}}}X_L(F)I\right),`$ (2.74)
$`\mathrm{\Psi }^{}`$ $`=`$ $`\mathrm{exp}\left({\displaystyle \frac{i}{2\sqrt{\alpha ^{}}}}X_L(F)I\right)\mathrm{\Psi }^{\prime \prime }\mathrm{exp}\left({\displaystyle \frac{i}{2\sqrt{\alpha ^{}}}}X_L(F)I\right),`$ (2.75)
$`\mathrm{\Xi }^{}`$ $`=`$ $`\mathrm{exp}\left({\displaystyle \frac{i}{2\sqrt{\alpha ^{}}}}X_L(F)I\right)\mathrm{\Xi }^{\prime \prime }\mathrm{exp}\left({\displaystyle \frac{i}{2\sqrt{\alpha ^{}}}}X_L(F)I\right).`$ (2.76)
We can easily find that the new BRS charge is transformed to the original form in both of the total action and the constraint. As in the previous subsection, the classical solution has an effect only on the momentum in the string fields. The result indicates that the classical solution (2.5) corresponds to the background Wilson line in open superstring field theory including the NS and R sector.
Instead of using the action, we can see the effect of the classical solution in terms of the equation of motion. The original Berkovits’ equations of motion are given by $`\eta _0(e^\mathrm{\Phi }Q_\mathrm{B}e^\mathrm{\Phi })=(\eta _0\mathrm{\Psi })^2,Q_\mathrm{B}(e^\mathrm{\Phi }(\eta _0\mathrm{\Psi })e^\mathrm{\Phi })=0`$. Expanding the equations of motion around the classical solution, we find that the form of the equations is unchanged but the BRS charge is changed to $`Q_\mathrm{B}^{}`$ of (2.48). In the case of our solution $`\mathrm{\Phi }_0=\stackrel{~}{V}_L(F)I,\mathrm{\Psi }_0=0`$ (2.5), the expanded equations of motion are transformed back to the original ones by the field redefinitions: $`\mathrm{\Phi }^{}=e^{\frac{i}{2\sqrt{\alpha ^{}}}X_L(F)I}\mathrm{\Phi }^{\prime \prime }e^{\frac{i}{2\sqrt{\alpha ^{}}}X_L(F)I},\mathrm{\Psi }^{}=e^{\frac{i}{2\sqrt{\alpha ^{}}}X_L(F)I}\mathrm{\Psi }^{\prime \prime }e^{\frac{i}{2\sqrt{\alpha ^{}}}X_L(F)I}`$. This redefinition reflects the effect of the Wilson lines.
The equations of motion for the NS and R sectors have a fermionic gauge symmetry as follows :
$`\delta e^\mathrm{\Phi }`$ $`=`$ $`e^\mathrm{\Phi }(\eta _0\mathrm{\Psi }\mathrm{\Lambda }_{1/2}+\mathrm{\Lambda }_{1/2}\eta _0\mathrm{\Psi }),`$ (2.77)
$`\delta \mathrm{\Psi }`$ $`=`$ $`Q_\mathrm{B}\mathrm{\Lambda }_{1/2}+e^\mathrm{\Phi }Q_\mathrm{B}e^\mathrm{\Phi }\mathrm{\Lambda }_{1/2}+\mathrm{\Lambda }_{1/2}e^\mathrm{\Phi }Q_\mathrm{B}e^\mathrm{\Phi },`$ (2.78)
where $`\mathrm{\Lambda }_{1/2}`$ is a Grassmann odd parameter with the picture number $`1/2`$. This may include a global space-time supersymmetry, since the NS (R) string field is transformed to the R (NS) sector under the transformation. Actually, we can find the supersymmetry if we formally set a transformation parameter as
$`\mathrm{\Lambda }_{1/2}=\mathrm{\Omega }(ϵ)=ϵ_\alpha {\displaystyle _{C_{\mathrm{left}}}}{\displaystyle \frac{dz}{2\pi i}}\xi S_{(1/2)}^\alpha (z)I,`$ (2.79)
where $`S_{(1/2)}^\alpha `$ is a GSO(+) spin field with $`\varphi `$-charge $`1/2`$ and positive chirality and $`ϵ_\alpha `$ is a fermionic constant. We give the details of the supersymmetry in appendix C. Substituting (2.79) into eqs. (2.77) and (2.78), we can rewrite the transformation law as<sup>14</sup><sup>14</sup>14 We do not include a contribution from the first term in (2.78): $`\delta _\mathrm{M}(\eta _0\mathrm{\Psi })Q_\mathrm{B}\eta _0\mathrm{\Omega }(ϵ)=\frac{i}{2\pi }ϵ_\alpha (cS_{(1/2)}^\alpha (i)cS_{(1/2)}^\alpha (i))I`$. The elimination of this term is possible because the transformation $`\delta _\mathrm{M}`$ is a symmetry of the equations of motion: $`\delta _\mathrm{M}f_1=\{𝒪(i)I,\eta _0\mathrm{\Psi }\}=𝒪(i)\eta _0\mathrm{\Psi }𝒪(i)\eta _0\mathrm{\Psi }=0,\delta _\mathrm{M}f_2=\{𝒪(i)I,e^\mathrm{\Phi }Q_\mathrm{B}e^\mathrm{\Phi }\}=0`$, where $`f_1`$ and $`f_2`$ are given in eqs. (C.3) and (C.4) respectively and $`𝒪(i)\frac{i}{2\pi }ϵ_\alpha (cS_{(1/2)}^\alpha (i)cS_{(1/2)}^\alpha (i))`$.
$`\delta _ϵe^\mathrm{\Phi }=e^\mathrm{\Phi }𝒮(ϵ)\eta _0\mathrm{\Psi },\delta _ϵ(\eta _0\mathrm{\Psi })=\eta _0𝒮(ϵ)(e^\mathrm{\Phi }Q_\mathrm{B}e^\mathrm{\Phi }),`$ (2.80)
where the operator $`𝒮(ϵ)`$ is defined as
$`𝒮(ϵ)=ϵ_\alpha {\displaystyle \frac{dz}{2\pi i}\xi S_{(1/2)}^\alpha (z)}.`$ (2.81)
The operator $`𝒮(ϵ)`$ is an anti-derivation with respect to the star product. Now, we apply this transformation to the Wilson line solution, namely $`\mathrm{\Phi }_0`$ given by (2.5) and $`\mathrm{\Psi }_0=0`$:
$`\delta _ϵe^\mathrm{\Phi }`$ $`=`$ $`0,`$ (2.82)
$`\delta _ϵ(\eta _0\mathrm{\Psi })`$ $`=`$ $`\eta _0𝒮(ϵ)(e^{\mathrm{\Phi }_0}Q_\mathrm{B}e^{\mathrm{\Phi }_0})=ϵ_\alpha {\displaystyle \frac{dz}{2\pi i}S_{(1/2)}^\alpha (z)\left(V_L(F)\frac{1}{4}C_L(F^2)\right)I}`$ (2.83)
$`=`$ $`ϵ_\alpha {\displaystyle \frac{dz}{2\pi i}\{S_{(1/2)}^\alpha (z),V_L(F)\}I}=0,`$
where use has been made of $`\{v(y),S_{(1/2)}^\alpha (z)\}=0`$. We note that $`{\displaystyle \frac{dz}{2\pi i}S_{(1/2)}^\alpha I}=0`$ and $`𝒮(ϵ)I=0`$ because both $`S_{(1/2)}^\alpha (z)`$ and $`\xi S_{(1/2)}^\alpha (z)`$ are primary fields with conformal dimension 1. This result indicates that the Wilson line solution exactly preserves all global space-time supersymmetries.
### 2.5 symmetries and classical solutions
The half integration of $`F(z)`$ is related to a Wilson line parameter, which should be a gauge invariant quantity of the stringy gauge group. This relation suggests that other modes of $`F(z)`$ are redundant degrees of freedom under the gauge symmetry. In this subsection we would like to illustrate this point, namely the function $`F(z)`$ can be changed by an appropriate transformation except the half integration mode.
The total action including the R sector is invariant under the finite gauge transformation,
$`e^\mathrm{\Phi }^{}`$ $`=`$ $`e^{Q_\mathrm{B}\mathrm{\Lambda }_0}e^\mathrm{\Phi }e^{\eta _0\mathrm{\Lambda }_1},`$ (2.84)
$`\mathrm{\Psi }^{}`$ $`=`$ $`e^{\eta _0\mathrm{\Lambda }_1}\mathrm{\Psi }e^{\eta _0\mathrm{\Lambda }_1}+{\displaystyle \frac{1e^{\mathrm{ad}_{\eta _0\mathrm{\Lambda }_1}}}{\mathrm{ad}_{\eta _0\mathrm{\Lambda }_1}}}(\eta _0\mathrm{\Lambda }_{3/2}),`$ (2.85)
$`\mathrm{\Xi }^{}`$ $`=`$ $`e^{Q_\mathrm{B}\mathrm{\Lambda }_0}\mathrm{\Xi }e^{Q_\mathrm{B}\mathrm{\Lambda }_0}+{\displaystyle \frac{e^{\mathrm{ad}_{Q_\mathrm{B}\mathrm{\Lambda }_0}}1}{\mathrm{ad}_{Q_\mathrm{B}\mathrm{\Lambda }_0}}}(Q_\mathrm{B}\mathrm{\Lambda }_{1/2}).`$ (2.86)
This transformation can be obtained by performing $`N`$ times of the infinitesimal transformation given in eqs. (2.66), (2.67) and (2.68) and then taking the limit $`N\mathrm{}`$:
$`e^\mathrm{\Phi }^{}=\underset{N\mathrm{}}{lim}\left(1+{\displaystyle \frac{1}{N}}Q_\mathrm{B}\mathrm{\Lambda }_0\right)^Ne^\mathrm{\Phi }\left(1+{\displaystyle \frac{1}{N}}\eta _0\mathrm{\Lambda }_1\right)^N,`$ (2.87)
$`\mathrm{\Psi }^{}=\underset{N\mathrm{}}{lim}\left[\left(1{\displaystyle \frac{1}{N}}\mathrm{ad}_{\eta _0\mathrm{\Lambda }_1}\right)^N\mathrm{\Psi }+{\displaystyle \underset{k=0}{\overset{N1}{}}}\left(1{\displaystyle \frac{1}{N}}\mathrm{ad}_{\eta _0\mathrm{\Lambda }_1}\right)^k{\displaystyle \frac{1}{N}}\eta _0\mathrm{\Lambda }_{3/2}\right],`$ (2.88)
$`\mathrm{\Xi }^{}=\underset{N\mathrm{}}{lim}\left[\left(1+{\displaystyle \frac{1}{N}}\mathrm{ad}_{Q_\mathrm{B}\mathrm{\Lambda }_0}\right)^N\mathrm{\Xi }+{\displaystyle \underset{k=0}{\overset{N1}{}}}\left(1+{\displaystyle \frac{1}{N}}\mathrm{ad}_{Q_\mathrm{B}\mathrm{\Lambda }_0}\right)^k{\displaystyle \frac{1}{N}}Q_\mathrm{B}\mathrm{\Lambda }_{1/2}\right],`$ (2.89)
where we set $`\delta \mathrm{\Lambda }_P=\mathrm{\Lambda }_P/N`$ in (2.66), (2.67) and (2.68) and we have used the definition $`\mathrm{ad}_XY[X,Y]=XYYX`$.
Substituting $`Q_\mathrm{B}\mathrm{\Lambda }_0=\eta _0\mathrm{\Lambda }_1=\mathrm{{\rm Y}}_0`$ and $`\mathrm{\Lambda }_{3/2}=\mathrm{\Lambda }_{1/2}=0`$ into (2.84), (2.85) and (2.86), we find
$`e^\mathrm{\Phi }^{}=e^{\mathrm{{\rm Y}}_0}e^\mathrm{\Phi }e^{\mathrm{{\rm Y}}_0},\mathrm{\Psi }^{}=e^{\mathrm{{\rm Y}}_0}\mathrm{\Psi }e^{\mathrm{{\rm Y}}_0},\mathrm{\Xi }^{}=e^{\mathrm{{\rm Y}}_0}\mathrm{\Xi }e^{\mathrm{{\rm Y}}_0}.`$ (2.90)
In the large Hilbert space, any state $`\mathrm{{\rm Y}}_0`$ satisfying $`Q_\mathrm{B}\mathrm{{\rm Y}}_0=\eta _0\mathrm{{\rm Y}}_0=0`$ can be written as $`\mathrm{{\rm Y}}_0=Q_\mathrm{B}\mathrm{\Lambda }_0=\eta _0\mathrm{\Lambda }_1`$. Hence, the total action is invariant under the similarity transformation generated by $`\mathrm{{\rm Y}}_0`$ such that $`Q_\mathrm{B}\mathrm{{\rm Y}}_0=\eta _0\mathrm{{\rm Y}}_0=0`$.
The BRS charge $`Q_\mathrm{B}`$ corresponds to the external derivative in the WZW theory. It is easy to show that the “pure gauge connection” $`A_Q=e^\mathrm{\Phi }Q_\mathrm{B}e^\mathrm{\Phi }`$ is transformed as
$`A_Q^{}`$ $`=`$ $`e^{\eta _0\mathrm{\Lambda }_1}Q_\mathrm{B}e^{\eta _0\mathrm{\Lambda }_1}+e^{\eta _0\mathrm{\Lambda }_1}A_Qe^{\eta _0\mathrm{\Lambda }_1},`$ (2.91)
under the gauge transformation. If we perform the similarity transformation (2.90) on $`A_Q`$, we find $`A_Q^{}=e^{\mathrm{{\rm Y}}_0}A_Qe^{\mathrm{{\rm Y}}_0}`$ since $`Q_\mathrm{B}e^{\mathrm{{\rm Y}}_0}=0`$. This transformation law allows us to interpret (2.90) as a “global transformation” , which is a transformation of a subgroup of the stringy gauge group.
Now, let us consider a “global transformation” generated by the parameter,
$`\mathrm{{\rm Y}}_0=𝒯_L(f)I={\displaystyle _{C_{\mathrm{left}}}}{\displaystyle \frac{dz}{2\pi i}}f(z)T(z)I,f(1/z)=z^2f(z),`$ (2.92)
where $`T(z)`$ is the total energy momentum tensor. It is easily seen that $`\mathrm{{\rm Y}}_0`$ satisfies $`Q_\mathrm{B}\mathrm{{\rm Y}}_0=\eta _0\mathrm{{\rm Y}}_0=0`$ since $`[Q_\mathrm{B},T(z)]=[\eta _0,T(z)]=0`$. More explicitly, the gauge transformation parameters $`\mathrm{\Lambda }_0`$ and $`\mathrm{\Lambda }_1`$ satisfying $`Q_\mathrm{B}\mathrm{\Lambda }_0=\eta _0\mathrm{\Lambda }_1=\mathrm{{\rm Y}}_0`$ can be written as
$`\mathrm{\Lambda }_0=U_L(f)I,\mathrm{\Lambda }_1=\xi _0𝒯_L(f)I,`$ (2.93)
where $`U_L(f)`$ is defined as<sup>15</sup><sup>15</sup>15 $`T^\mathrm{m}(z)`$ is the matter energy momentum tensor. We note that $`\{Q_\mathrm{B},u(z)\}=T(z)`$.
$`U_L(f)`$ $`=`$ $`{\displaystyle _{C_{\mathrm{left}}}}{\displaystyle \frac{dz}{2\pi i}}f(z)u(z),`$ (2.94)
$`u(z)`$ $`=`$ $`T^\mathrm{m}c\xi \xi e^{2\varphi }(z)2bcc\xi \xi e^{2\varphi }(z)+c\xi ^2\xi e^{2\varphi }(z){\displaystyle \frac{3}{2}}^2c\xi \xi e^{2\varphi }(z)`$ (2.95)
$`{\displaystyle \frac{1}{2}}c\xi ^3\xi e^{2\varphi }(z)+c\xi \xi \left({\displaystyle \frac{1}{2}}(\varphi )^2+3^2\varphi \right)e^{2\varphi }(z).`$
Using some properties of $`𝒯_L(f)`$,<sup>16</sup><sup>16</sup>16The operator $`𝒯_L(f)`$ satisfies $`[𝒯_{L(R)}(f),𝒯_{L(R)}(g)]=𝒯_{L(R)}((f)gfg),[𝒯_L(f),𝒯_R(g)]=0,`$ (2.96) $`(𝒯_R(f)A)B=A(𝒯_L(f)B),𝒯_R(f)I=𝒯_L(f)I,`$ (2.97) as in the bosonic case . we can rewrite the global transformation as
$`\mathrm{\Phi }^{}=e^{𝒯(f)}\mathrm{\Phi },\mathrm{\Psi }^{}=e^{𝒯(f)}\mathrm{\Psi },\mathrm{\Xi }^{}=e^{𝒯(f)}\mathrm{\Xi },𝒯(f)={\displaystyle \frac{dz}{2\pi i}f(z)T(z)}.`$ (2.98)
If we apply an infinitesimal transformation on the classical solution (2.5), it changes to
$`\mathrm{\Phi }_0^{}`$ $`=`$ $`e^{𝒯(ϵ)}\mathrm{\Phi }_0=\stackrel{~}{V}_L(F^{})I,`$ (2.99)
$`F^{}(z)`$ $`=`$ $`F(z)(ϵ(z)F(z)),`$ (2.100)
where we have used the commutation relation $`[𝒯(f),\stackrel{~}{V}_L(F)]=\stackrel{~}{V}_L((fF))`$. We note that $`F^{}(z)`$ satisfies $`F^{}(1/z)=z^2F^{}(z)`$. Accordingly, we find that the function form of $`F(z)`$ are redundant under the gauge transformation. However, we cannot change the half integration mode of $`F(z)`$. Indeed, we find that, from eq. (2.100),
$`f={\displaystyle _{C_{\mathrm{left}}}}{\displaystyle \frac{dz}{2\pi i}}F(z)={\displaystyle _{C_{\mathrm{left}}}}{\displaystyle \frac{dz}{2\pi i}}F^{}(z),`$ (2.101)
where we have used $`ϵ(\pm i)=F(\pm i)=0`$ due to $`ϵ(1/z)=z^2ϵ(z)`$ and $`F(1/z)=z^2F(z)`$. These are consistent results with our expectation. The half integration mode of $`F(z)`$, which is to be a physical quantity, is invariant but other modes can be gauged away.
## 3 Marginal deformations and classical solutions
In the previous section, we have described a class of solutions which correspond to the Wilson lines. It turns out that they are based on algebra satisfied by $`u(1)`$ supercurrent $`𝐉(z,\theta )=\psi (z)+\theta \frac{i}{\sqrt{2\alpha ^{}}}X(z)`$. From this point of view, we can use the same method to construct classical solutions of superstring field theory which correspond to more general supercurrents or marginal deformations in the context of conformal field theory.<sup>17</sup><sup>17</sup>17 As a comparison, we discuss a similar generalization in the context of the Witten’s bosonic string field theory in appendix D.
Here we consider a supercurrent $`𝐉^a(z,\theta )=\psi ^a(z)+\theta J^a(z)`$ associated with a Lie algebra $`𝒢`$ in the matter sector ($`a=1,\mathrm{},\mathrm{dim}𝒢`$). In terms of component fields, we suppose that OPE is given by
$`\psi ^a(y)\psi ^b(z)`$ $``$ $`{\displaystyle \frac{1}{yz}}{\displaystyle \frac{1}{2}}\mathrm{\Omega }^{ab},`$ (3.1)
$`J^a(y)\psi ^b(z)`$ $``$ $`{\displaystyle \frac{1}{yz}}f_c^{ab}\psi ^c(z),`$ (3.2)
$`J^a(y)J^b(z)`$ $``$ $`{\displaystyle \frac{1}{(yz)^2}}{\displaystyle \frac{1}{2}}\mathrm{\Omega }^{ab}+{\displaystyle \frac{1}{yz}}f_c^{ab}J^c(z),`$ (3.3)
where $`f_c^{ab}`$ is the structure constant of $`𝒢`$ ($`f_c^{ab}=f_c^{ba},f_d^{ab}f_e^{cd}+f_d^{bc}f_e^{ad}+f_d^{ca}f_e^{bd}=0`$) and $`\mathrm{\Omega }^{ab}`$ is an invertible matrix<sup>18</sup><sup>18</sup>18 In the case of semi-simple Lie algebra, we can take $`\mathrm{\Omega }^{ab}`$ as the Killing form $`\gamma ^{ab}=f_d^{ac}f_c^{bd}`$. However, we have supposed the existence of invertible $`\mathrm{\Omega }^{ab}`$ in order to include the cases of non-semi-simple algebra after ref. . which satisfies
$`\mathrm{\Omega }^{ab}=\mathrm{\Omega }^{ba},f_c^{ab}\mathrm{\Omega }^{cd}+f_c^{ad}\mathrm{\Omega }^{cb}=0.`$ (3.4)
In this case, energy momentum tensor $`T(z)`$ and its super partner $`G(z)`$ are given by a general Sugawara construction :
$`T(z)`$ $`=`$ $`\mathrm{\Omega }_{ab}:(J^aJ^b+\psi ^a\psi ^b):(z)+{\displaystyle \frac{2}{3}}\mathrm{\Omega }_{ad}\mathrm{\Omega }_{be}f_c^{de}:(J^a:\psi ^b\psi ^c:+\psi ^a:(\psi ^bJ^cJ^b\psi ^c):):(z),`$ (3.5)
$`G(z)`$ $`=`$ $`2\mathrm{\Omega }_{ab}:J^a\psi ^b:(z)+{\displaystyle \frac{4}{3}}\mathrm{\Omega }_{ad}\mathrm{\Omega }_{be}f_c^{de}:\psi ^a:\psi ^b\psi ^c::(z),`$ (3.6)
where $`\mathrm{\Omega }_{ab}`$ is the inverse of $`\mathrm{\Omega }^{ab}`$: $`\mathrm{\Omega }^{ab}\mathrm{\Omega }_{bc}=\delta _c^a`$. In fact, they satisfy the following OPEs:
$`T(y)\psi ^a(z)`$ $``$ $`{\displaystyle \frac{1}{(yz)^2}}{\displaystyle \frac{1}{2}}\psi ^a(z)+{\displaystyle \frac{1}{yz}}\psi ^a(z),`$ (3.7)
$`T(y)J^a(z)`$ $``$ $`{\displaystyle \frac{1}{(yz)^2}}J^a(z)+{\displaystyle \frac{1}{yz}}J^a(z),`$ (3.8)
$`G(y)\psi ^a(z)`$ $``$ $`{\displaystyle \frac{1}{yz}}J^a(z),`$ (3.9)
$`G(y)J^a(z)`$ $``$ $`{\displaystyle \frac{1}{(yz)^2}}\psi ^a(z)+{\displaystyle \frac{1}{yz}}\psi ^a(z),`$ (3.10)
$`T(y)T(z)`$ $``$ $`{\displaystyle \frac{c}{2}}{\displaystyle \frac{1}{(yz)^4}}+{\displaystyle \frac{1}{(yz)^2}}2T(z)+{\displaystyle \frac{1}{yz}}T(z),`$ (3.11)
$`G(y)G(z)`$ $``$ $`{\displaystyle \frac{2c}{3}}{\displaystyle \frac{1}{(yz)^3}}+{\displaystyle \frac{1}{yz}}\mathrm{\hspace{0.17em}2}T(z),`$ (3.12)
$`T(y)G(z)`$ $``$ $`{\displaystyle \frac{1}{(yz)^2}}{\displaystyle \frac{3}{2}}G(z)+{\displaystyle \frac{1}{yz}}G(z),`$ (3.13)
and the central charge $`c`$ is given by $`c=\frac{3}{2}\mathrm{dim}𝒢f_d^{ac}f_c^{bd}\mathrm{\Omega }_{ab}`$ .
Let us consider the Berkovits’ open superstring field theory on the above CFT background. The action has the same form as the flat one:
$`S[\mathrm{\Phi }]`$ $`=`$ $`{\displaystyle \frac{1}{g^2}}{\displaystyle _0^1}𝑑t(\eta _0\mathrm{\Phi })(e^{t\mathrm{\Phi }}Q_\mathrm{B}e^{t\mathrm{\Phi }})`$ (3.14)
although $`T(z)`$ and $`G(z)`$ in the definition of the BRS operator<sup>19</sup><sup>19</sup>19 The BRS operator is given by the matter Virasoro operators $`T(z),G(z)`$ and ghosts $`(b,c,\varphi ,\xi ,\eta )`$ as: $`Q_\mathrm{B}`$ $`=`$ $`{\displaystyle \frac{dz}{2\pi i}\left[c(T\frac{1}{2}(\varphi )^2^2\varphi +\xi \eta )(z)+bcc(z)+\eta e^\varphi G(z)\eta \eta e^{2\varphi }b(z)\right]}.`$ (3.15) are given by (3.5) and (3.6). The star product among string fields in the action is constructed by LPP’s method in terms of conformal mappings and correlators in the above CFT. In order to make $`Q_\mathrm{B}`$ nilpotent, we assume that the total central charge in the matter sector is $`c=15`$. With this setup, we shall show that
$`\mathrm{\Phi }_0=\stackrel{~}{V}_L^a(F_a)I,`$ (3.16)
$`\stackrel{~}{V}_L^a(F_a)={\displaystyle _{C_{\mathrm{left}}}}{\displaystyle \frac{dz}{2\pi i}}F_a(z)\stackrel{~}{v}^a(z),F_a(1/z)=z^2F_a(z),`$ (3.17)
is a classical solution, where the operator $`\stackrel{~}{v}^a(z)`$ is given by the lowest component of supercurrent $`𝐉^a(z,\theta )`$ with dimension $`1/2`$ and appropriate ghost part:
$`\stackrel{~}{v}^a(z)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}c\xi e^\varphi \psi ^a(z),`$ (3.18)
and $`I`$ is the identity string field which is the identity element with respect to the star product. Noting $`\stackrel{~}{v}^a(z)`$ is a primary field with dimension $`0`$ and therefore satisfies
$`(\stackrel{~}{V}_L^a(F_a)I)B=\stackrel{~}{V}_R^a(F_a)IB=I(\stackrel{~}{V}_L^a(F_a)B)=\stackrel{~}{V}_L^a(F_a)B`$ (3.19)
for any string field $`B`$ where $`\stackrel{~}{V}_R^a(F_a)=_{C_{\mathrm{right}}}\frac{dz}{2\pi i}F_a(z)\stackrel{~}{v}^a(z)`$, we obtain
$`e^{\mathrm{\Phi }_0}Q_\mathrm{B}e^{\mathrm{\Phi }_0}`$ $`=`$ $`(e^{\stackrel{~}{V}_L^a(F_a)}Q_\mathrm{B}e^{\stackrel{~}{V}_L^a(F_a)})I=\left(V_L^a(F_a)+{\displaystyle \frac{1}{8}}\mathrm{\Omega }^{ab}C_L(F_aF_b)\right)I,`$ (3.20)
where
$`V_L^a(F_a)={\displaystyle _{C_{\mathrm{left}}}}{\displaystyle \frac{dz}{2\pi i}}F_a(z)v^a(z),v^a(z){\displaystyle \frac{1}{\sqrt{2}}}cJ^a(z)+{\displaystyle \frac{1}{\sqrt{2}}}\eta e^\varphi \psi ^a(z).`$ (3.21)
In the above computation, we have used the relations
$`[Q_\mathrm{B},\stackrel{~}{v}^a(z)]=v^a(z),`$ (3.22)
$`\stackrel{~}{v}^a(y)v^b(z){\displaystyle \frac{1}{yz}}{\displaystyle \frac{1}{4}}\mathrm{\Omega }^{ab}c(z),[\stackrel{~}{V}_L^a(f),V_L^b(g)]={\displaystyle \frac{1}{4}}\mathrm{\Omega }^{ab}C_L(fg),`$ (3.23)
which follow from (3.7), (3.9), (3.1) and (3.2). Because both $`V_L^a(F_a)`$ and $`C_L(F_aF_b)`$ in (3.20) do not include $`\xi _0`$, we conclude that $`\mathrm{\Phi }_0`$ (3.16) satisfies the equation of motion:
$`\eta _0(e^{\mathrm{\Phi }_0}Q_\mathrm{B}e^{\mathrm{\Phi }_0})=0.`$ (3.24)
By replacing $`F_a(z)`$ with $`tF_a(z)`$, we have $`\eta _0(e^{t\mathrm{\Phi }_0}Q_\mathrm{B}e^{t\mathrm{\Phi }_0})=0`$ ($`0t1`$), which implies that the value of the action at this solution is zero: $`S[\mathrm{\Phi }_0]=0`$ as we can easily check from eq. (3.14).
Around the solution $`\mathrm{\Phi }_0`$, using (3.20) and (B.19) and noting
$`A(V_L^a(F_a)I)=(1)^{|A|}(V_R^a(F_a)A)I=(1)^{|A|}V_R^a(F_a)A,`$ (3.25)
($`V_R^a(f)=_{C_{\mathrm{right}}}\frac{dz}{2\pi i}F_a(z)v^a(z)`$) for any string field $`A`$ because $`v^a(z)`$ is a primary field with dimension $`0`$, the new BRS operator $`Q_\mathrm{B}^{}`$ which is a derivation with respect to the star product becomes
$`Q_\mathrm{B}^{}B`$
$`=Q_\mathrm{B}B+[(V_L^a(F_a)+{\displaystyle \frac{1}{8}}\mathrm{\Omega }^{ab}C_L(F_aF_b))I]B(1)^{|B|}B[(V_L^a(F_a)+{\displaystyle \frac{1}{8}}\mathrm{\Omega }^{ab}C_L(F_aF_b))I]`$
$`=\left(Q_\mathrm{B}(V_L^a(F_a)+V_R^a(F_a))+{\displaystyle \frac{1}{8}}\mathrm{\Omega }^{ab}(C_L(F_aF_b)+C_R(F_aF_b))\right)B,`$ (3.26)
on any string field $`B`$, namely,
$`Q_\mathrm{B}^{}=Q_\mathrm{B}V^a(F_a)+{\displaystyle \frac{1}{8}}\mathrm{\Omega }^{ab}C(F_aF_b),`$ (3.27)
$`V^a(F_a)={\displaystyle \frac{dz}{2\pi i}F_a(z)v^a(z)},C(F_aF_b)={\displaystyle \frac{dz}{2\pi i}F_a(z)F_b(z)c(z)}.`$ (3.28)
We can directly check $`\{\eta _0,Q_\mathrm{B}^{}\}=0`$ and nilpotency $`Q_\mathrm{B}^2=0`$ noting $`\{Q_\mathrm{B},v^a(z)\}=0`$ from (3.22) and
$`v^a(y)v^b(z){\displaystyle \frac{\mathrm{\Omega }^{ab}}{4(yz)}}(cc\eta \eta e^{2\varphi })(z),\{V^a(F_a),V^b(F_b)\}={\displaystyle \frac{\mathrm{\Omega }^{ab}}{4}}\{Q_\mathrm{B},C(F_aF_b)\}.`$ (3.29)
We note that the above BRS operator $`Q_\mathrm{B}^{}`$ is obtained by replacing the matter Virasoro operators $`G(z),T(z)`$ in (3.15) with $`G(z)\frac{1}{\sqrt{2}}F_a(z)\psi ^a(z)`$, $`T(z)\frac{1}{\sqrt{2}}F_a(z)J^a(z)+\frac{1}{8}\mathrm{\Omega }^{ab}F_a(z)F_b(z)`$ respectively. In fact, if we define $`G^{}(z)=_rG_r^{}z^{r3/2}`$ and $`T^{}(z)=_nL_n^{}z^{n2}`$ as
$`G_r^{}`$ $`=`$ $`G_r{\displaystyle \frac{1}{\sqrt{2}}}{\displaystyle \underset{k}{}}F_{a,k}\psi _{rk}^a,`$ (3.30)
$`L_n^{}`$ $`=`$ $`L_n{\displaystyle \frac{1}{\sqrt{2}}}{\displaystyle \underset{k}{}}F_{a,k}J_{nk}^a+{\displaystyle \frac{1}{8}}\mathrm{\Omega }^{ab}{\displaystyle \underset{k}{}}F_{a,nk}F_{b,k},`$ (3.31)
where $`\psi ^a(z)=_r\psi _r^az^{r1/2},J^a(z)=_nJ_n^az^{n1},G(z)=_rG_rz^{r3/2},T(z)=_nL_nz^{n2}`$ and $`F_{a,n}=\frac{d\sigma }{2\pi }e^{i(n+1)\sigma }F_a(e^{i\sigma })`$, then, using OPEs among $`(\psi ^a,J^a,G,T)`$, we can check that they satisfy the super Virasoro algebra:
$`[L_m^{},L_n^{}]=(mn)L_{m+n}^{}+{\displaystyle \frac{c}{12}}(m^3m)\delta _{m+n,0},`$ (3.32)
$`\{G_r^{},G_s^{}\}=2L_{r+s}^{}+{\displaystyle \frac{c}{12}}(4r^21)\delta _{r+s,0},`$ (3.33)
$`[L_m^{},G_r^{}]=\left({\displaystyle \frac{m}{2}}r\right)G_{m+r}^{},`$ (3.34)
with the same central charge as original $`G(z),T(z)`$ system. Furthermore, let us define $`\psi ^a(z)=_r\psi _r^az^{r1/2}`$ and $`J^a(z)=_nJ_n^az^{n1}`$ by
$`\psi _r^a={\displaystyle \underset{k}{}}M_{b,k}^a\psi _{rk}^b,J_n^a={\displaystyle \underset{k}{}}M_{b,k}^a\left(J_{nk}^b{\displaystyle \frac{1}{2\sqrt{2}}}\mathrm{\Omega }^{bc}F_{c,nk}\right),`$ (3.35)
where $`M_{b,n}^a`$ is given by a path-ordered form:
$`M_b^a(\sigma )`$ $`=`$ $`{\displaystyle \underset{n}{}}M_{b,n}^ae^{in\sigma }=\left[𝐏\mathrm{exp}\left(i{\displaystyle _0^1}𝑑t\sigma A(t\sigma )\right)\right]_b^a`$ (3.36)
$`=`$ $`\delta _b^a+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}i^n\sigma ^n{\displaystyle _0^1}𝑑t_1{\displaystyle _0^{t_1}}𝑑t_2\mathrm{}{\displaystyle _0^{t_{n1}}}𝑑t_nA_{c_n}^a(t_n\sigma )A_{c_{n1}}^{c_n}(t_{n1}\sigma )\mathrm{}A_b^{c_2}(t_1\sigma ),`$
$`A_b^a(\sigma )`$ $``$ $`{\displaystyle \frac{1}{\sqrt{2}}}f_b^{ac}e^{i\sigma }F_c(e^{i\sigma }).`$ (3.37)
Noting the identities for invariant metric $`\mathrm{\Omega }^{ab}`$ (3.4) and the Jacobi identity for structure constants $`f_c^{ab}`$, we can show following relations:
$`i_\sigma M_b^a(\sigma )=M_c^a(\sigma )A_b^c(\sigma ),nM_{b,n}^a={\displaystyle \frac{1}{\sqrt{2}}}{\displaystyle \underset{k}{}}M_{d,k}^af_b^{dc}F_{c,nk},`$ (3.38)
$`M_c^a(\sigma )M_d^b(\sigma )\mathrm{\Omega }^{cd}=\mathrm{\Omega }^{ab},{\displaystyle \underset{k}{}}M_{c,k}^aM_{d,nk}^b\mathrm{\Omega }^{cd}=\mathrm{\Omega }^{ab}\delta _{n,0},`$ (3.39)
$`M_d^a(\sigma )M_e^b(\sigma )f_c^{de}=f_d^{ab}M_c^d(\sigma ),{\displaystyle \underset{k}{}}M_{d,k}^aM_{e,nk}^bf_c^{de}=f_d^{ab}M_{c,n}^d,`$ (3.40)
and we obtain commutation relations:
$`\{\psi _r^a,\psi _s^b\}={\displaystyle \frac{1}{2}}\mathrm{\Omega }^{ab}\delta _{r+s,0},[J_m^a,J_n^b]={\displaystyle \frac{1}{2}}\mathrm{\Omega }^{ab}m\delta _{m+n,0}+f_c^{ab}J_{m+n}^c,[J_n^a,\psi _r^b]=f_c^{ab}\psi _{n+r}^c,`$ (3.41)
$`[L_n^{},\psi _r^a]=\left({\displaystyle \frac{n}{2}}+r\right)\psi _{n+r}^a,[L_m^{},J_n^a]=nJ_{m+n}^a,`$ (3.42)
$`\{G_r^{},\psi _s^a\}=J_{r+s}^a,[G_r^{},J_n^a]=n\psi _{n+r}^a,`$ (3.43)
which are the same form as the original (unprimed) ones.
After all, by re-expanding the action (3.14) around a classical solution $`\mathrm{\Phi }_0`$ (3.16) as $`e^\mathrm{\Phi }=e^{\mathrm{\Phi }_0}e^\mathrm{\Phi }^{}`$, we obtain the action $`S^{}[\mathrm{\Phi }^{}]`$ with new BRS operator (3.27), which is realized by a replacement $`(\psi ^a,J^a,G,T)(\psi ^a,J^a,G^{},T^{})`$ in eqs. (3.30), (3.31) and (3.35) preserving algebra among them. This fact and vanishing vacuum energy: $`S[\mathrm{\Phi }_0]=0`$ suggest that the classical solution $`\mathrm{\Phi }_0`$ (3.16) might be a pure gauge solution in terms of superstring field theory. Indeed, the new BRS operator (3.27) can be rewritten as a similarity transform from the original one:
$`Q_\mathrm{B}^{}=e^{\stackrel{~}{V}^a(F_a)}Q_\mathrm{B}e^{\stackrel{~}{V}^a(F_a)},\stackrel{~}{V}^a(F_a)={\displaystyle \frac{dz}{2\pi i}F_a(z)\stackrel{~}{v}^a(z)}.`$ (3.44)
One might think that the original action $`S[\mathrm{\Phi }^{\prime \prime }]`$ could be reproduced by a field redefinition such as $`\mathrm{\Phi }^{\prime \prime }=e^{\stackrel{~}{V}^a(F_a)}\mathrm{\Phi }^{}=e^{\stackrel{~}{V}_L^a(F_a)I}\mathrm{\Phi }^{}e^{\stackrel{~}{V}_L^a(F_a)I}`$ in the re-expanded action $`S^{}[\mathrm{\Phi }^{}]`$. However, it is not so trivial because there is another derivation $`\eta _0`$ in the action and $`[\eta _0,\stackrel{~}{V}^a(F_a)]0`$.
In the following, we demonstrate that if the function $`F_a(z)`$ satisfies a condition, we can explicitly rewrite our solution in a pure gauge form and take an appropriate field redefinition around it, which recovers original action. Let us consider a particular pure gauge form and try to rewrite our solution to it. Noting the commutation relation $`[Q_\mathrm{B},J^a(z)(c\xi e^\varphi \psi ^a)(z)]=0`$, we find an identity:
$`J^a(z)`$ $`=`$ $`(c\xi e^\varphi \psi ^a)(z)+\{Q_\mathrm{B},\mathrm{\Omega }^a(z)\},`$ (3.45)
$`\mathrm{\Omega }^a(z)`$ $``$ $`{\displaystyle \frac{1}{2}}cc\xi \xi ^2\xi e^{3\varphi }\psi ^a(z)c\xi \xi e^{2\varphi }J^a(z),`$ (3.46)
which relates $`\stackrel{~}{v}^a`$ (3.18) to the current $`J^a`$ and we have
$`\stackrel{~}{V}_L^a(g_a)+J_L^a(g_a)=\{Q_\mathrm{B},\mathrm{\Omega }_L^a(g_a)\},`$ (3.47)
$`J_L^a(g_a)={\displaystyle _{C_{\mathrm{left}}}}{\displaystyle \frac{dz}{2\pi i}}{\displaystyle \frac{1}{\sqrt{2}}}g_a(z)J^a(z),\mathrm{\Omega }_L^a(g_a)={\displaystyle _{C_{\mathrm{left}}}}{\displaystyle \frac{dz}{2\pi i}}{\displaystyle \frac{1}{\sqrt{2}}}g_a(z)\mathrm{\Omega }^a(z),g_a(\pm i)=0.`$ (3.48)
The last condition is necessary to remove unwanted boundary contribution of the integration given by $`g_a(\pm i)\stackrel{~}{v}^a(\pm i)`$. We notice commutation relations
$`[\stackrel{~}{V}_L^a(F_a),\stackrel{~}{V}_L^b(G_b)]=0,[\stackrel{~}{V}_L^a(F_a),J_L^b(g_b)]={\displaystyle \frac{1}{\sqrt{2}}}f_c^{ab}\stackrel{~}{V}_L^c(F_ag_b),`$ (3.49)
and a kind of Hausdorff formula:
$`e^Ae^B`$ $`=`$ $`\mathrm{exp}\left(A+\mathrm{ad}_{\frac{A}{2}}(1+\mathrm{coth}(\mathrm{ad}_{\frac{A}{2}}))B+\mathrm{O}(B^2)\right),`$ (3.50)
where we have denoted $`\mathrm{ad}_XY=[X,Y]`$ and $`\mathrm{O}(B^2)`$ is quadratic and higher terms with respect to $`B`$. By substituting $`A=\stackrel{~}{V}_L^a(g_a)+J_L^a(g_a)`$ and $`B=\stackrel{~}{V}_L^a(F_a)`$ to the above formula, we obtain
$`e^{\stackrel{~}{V}_L^a(g_a)+J_L^a(g_a)}e^{\stackrel{~}{V}_L^a(F_a)}`$ $`=`$ $`e^{J_L^a(g_a)+\stackrel{~}{V}_L^a(g_aF_b((e^{}1)^1)_a^b)},`$ (3.51)
$`_a^b(z)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}f_a^{bc}g_c(z),`$ (3.52)
because $`\mathrm{O}(B^2)`$ in the exponent vanishes due to (3.49). If the second term in the exponent on the right hand side of the first line vanishes, we can compute as:
$`e^{Q_\mathrm{B}\mathrm{\Omega }_L^a(g_a)I}e^{\stackrel{~}{V}_L^a(F_a)I}=e^{(\stackrel{~}{V}_L^a(g_a)+J_L^a(g_a))}e^{\stackrel{~}{V}_L^a(F_a)}I=e^{J_L^a(g_a)}I=e^{\eta _0\xi _0J_L^a(g_a)I},`$ (3.53)
where we should impose $`g_a(1/z)=g_a(z)`$ to guarantee a relation such as eq. (2.20). The above calculation means that, by solving a differential equation with respect to $`g_a(z)`$:
$`F_b(z)=g_a(z)((e^{}1)^1)_b^a(z),`$ (3.54)
$`g_a(1/z)=g_a(z),g_a(\pm i)=0,`$ (3.55)
for a given $`F_a(z)`$ which specifies the classical solution $`\mathrm{\Phi }_0`$ (3.16), we obtain the form of gauge transformation from the trivial solution $`\mathrm{\Phi }=0`$ in superstring field theory:
$`e^{\mathrm{\Phi }_0}`$ $`=`$ $`e^{Q_\mathrm{B}(\mathrm{\Omega }_L^a(g_a))I}e^{\eta _0(\xi _0J_L^a(g_a))I}.`$ (3.56)
Note that the equation (3.54) is consistent with the conditions $`F_a(1/z)=z^2F_a(z)`$ and $`g_a(1/z)=g_a(z)`$. For an abelian $`𝒢`$, where $`_a^b(z)=0`$, it becomes a simple form $`g_a(z)=F_a(z)`$ because of $`((e^{}1)^1)_b^a(z)=\delta _b^a+\mathrm{O}()`$. Similarly, we can rewrite the new BRS operator (3.27) as
$`Q_\mathrm{B}^{}`$ $`=`$ $`e^\mathrm{\Lambda }Q_\mathrm{B}e^\mathrm{\Lambda },\mathrm{\Lambda }={\displaystyle \frac{dz}{2\pi i}\frac{1}{\sqrt{2}}g_a(z)J^a(z)}`$ (3.57)
using the above $`g_a`$. Thanks to $`[\eta _0,\mathrm{\Lambda }]=0`$, we can recover the original action by taking a field redefinition $`\mathrm{\Phi }^{\prime \prime }=e^\mathrm{\Lambda }\mathrm{\Phi }^{}=e^{\frac{1}{2}\lambda _a{\scriptscriptstyle {\scriptscriptstyle \frac{dz}{2\pi }}J^a(z)}+\mathrm{}}\mathrm{\Phi }^{}`$, which corresponds to a marginal deformation by the current $`J^a`$ . Its deformation parameter $`\lambda _a=\sqrt{2}i_\pi ^\pi \frac{d\sigma }{2\pi }g_a(e^{i\sigma })`$ is related to the function $`F_a`$ in our classical solution by eq. (3.54). However, a solution to the differential equation (3.54) which satisfies the conditions (3.55) does not necessarily exist. In fact, eq. (3.54) can be rewritten as,
$`{\displaystyle \frac{1}{\sqrt{2}}}F_b(z)T^b`$ $`=`$ $`e^{\frac{1}{\sqrt{2}}g_a(z)T^a}\left(e^{\frac{1}{\sqrt{2}}g_a(z)T^a}\right),`$ (3.58)
where $`T^a`$s are generator matrices of $`𝒢`$ such as $`[T^a,T^b]=f_c^{ab}T^c`$ and use has been made of $`(e^{}1)^1=_0^1𝑑te^t`$ and $`\delta (e^X)=_0^1𝑑te^{(1t)X}\delta Xe^{tX}`$. This equation can be solved by path-ordered form:
$`e^{\frac{1}{\sqrt{2}}g_a(e^{i\sigma })T^a}`$ $`=`$ $`𝐏e^{\frac{i}{\sqrt{2}}_0^1𝑑t\left(\sigma \mathrm{sgn}(\sigma )\frac{\pi }{2}\right)e^{it\sigma +i\mathrm{sgn}(\sigma )\frac{\pi }{2}(1t)}F_a(e^{it\sigma +i\mathrm{sgn}(\sigma )\frac{\pi }{2}(1t)})T^a},`$ (3.59)
where the path ordering denoted by $`𝐏`$ is taken as (3.36) (i.e., we put a matrix associated with larger $`t`$ to the right) and $`\mathrm{sgn}(\sigma )=+1(1)`$ for $`\sigma >0(\sigma <0)`$ is used. Here we have respected the second condition in (3.55): $`g_a(\pm i)=0`$ and solved separately on upper and lower half circle by taking the phase of $`z=e^{i\sigma }`$ as $`\pi \sigma \pi `$ on the unit circle. The property of the function $`F_a(z)`$ in our solution $`\mathrm{\Phi }_0`$ (3.16): $`F_a(1/z)=z^2F_a(z)`$ implies $`e^{i(\mathrm{sgn}(\sigma )\pi \sigma )}F_a(e^{i(\mathrm{sgn}(\sigma )\pi \sigma )})=e^{i\sigma }F_a(e^{i\sigma })`$ and then the above solution satisfies the first condition in (3.55): $`g_a(e^{i(\mathrm{sgn}(\sigma )\pi \sigma )})=g_a(e^{i\sigma })`$. In order to guarantee the continuity of $`g_a(e^{i\sigma })`$ at $`\sigma =0`$, which is needed for (3.47), there is a consistency condition for $`F_a(z)`$:
$`𝐏e^{\frac{i\pi }{2\sqrt{2}}_0^1𝑑te^{i\frac{\pi }{2}(1t)}F_a(e^{i\frac{\pi }{2}(1t)})T^a}=𝐏e^{\frac{i\pi }{2\sqrt{2}}_0^1𝑑te^{i\frac{\pi }{2}(1t)}F_a(e^{i\frac{\pi }{2}(1t)})T^a},`$ (3.60)
which is reduced to
$`{\displaystyle _{C_{\mathrm{left}}}}𝑑zF_a(z)=0`$ (3.61)
in the case of $`𝒢`$ : abelian. Namely, if $`F_a(z)`$ satisfies the condition (3.60) (or (3.61) for abelian $`𝒢`$), the solution (3.16) is rewritten in a pure gauge form (3.56) and induces a field redefinition generated by $`\mathrm{\Lambda }`$ in (3.57). Conversely, in the case that $`F_a(z)`$ breaks the condition (3.60), we cannot rewrite as (3.56) and we should know further informations about the supercurrent and its representation by specifying a model which realizes $`𝐉^a(z,\theta )`$ in order to find explicit relations between our solution and marginal deformation.
As an example of $`𝒢=u(1)^{10}`$, we take an ordinary flat background which is described by a supercurrent $`𝐉^\mu (z,\theta )=\psi ^\mu (z)+\theta \frac{i}{\sqrt{2\alpha ^{}}}X^\mu (z)`$. In this case, we can identify various quantities as follows:
$`T(z)={\displaystyle \frac{1}{4\alpha ^{}}}X^\mu X_\mu (z){\displaystyle \frac{1}{2}}\psi ^\mu \psi _\mu (z),G(z)={\displaystyle \frac{i}{\sqrt{2\alpha ^{}}}}X_\mu \psi ^\mu (z),`$ (3.62)
$`\mathrm{\Omega }^{\mu \nu }=2\eta ^{\mu \nu },\mathrm{\Omega }_{\mu \nu }={\displaystyle \frac{1}{2}}\eta _{\mu \nu },f_\rho ^{\mu \nu }=0,c={\displaystyle \frac{3}{2}}\mathrm{dim}(u(1)^{10})=15.`$ (3.63)
By taking functions $`F_\mu (z)`$ for a solution $`\mathrm{\Phi }_0`$ (3.16) as $`F_\mu (z)=\delta _{\mu ,9}F(z)`$ such as $`F(1/z)=z^2F(z)`$, we reproduce the solution in the previous section, which has turned out to correspond to the Wilson line. From the condition (3.61), non-vanishing Wilson line $`f=_{C_{\mathrm{left}}}\frac{dz}{2\pi i}F(z)0`$ implies non-existence of a function $`g_\mu (z)`$ which specifies globally defined gauge parameter of the form (3.56) and a field redefinition associated with (3.57). Instead, we have found other (local) expressions (2.60) and (2.58) using the integration of the current $`J^\mu (z)`$.
We comment on the analogy with arguments in the Witten’s bosonic string field theory (see appendix D). In both cases, we can construct a class of classical solutions (3.16), (D.11) based on (super-)currents, which have vanishing vacuum energy. The actions around the solution are also obtained by appropriate operator mappings (3.35), (D.15) which preserve current algebra in both cases. This fact suggests that the theory around the solutions is essentially the same as the original one and they might be gauge equivalent. Indeed, we can represent the solutions as a pure gauge form in both supersymmetric and bosonic string field theory if there exists a solution $`g_a(z)`$ to the differential equation (3.54), which is the same form as (D.21). In both cases, the change of BRS operator around the solution turns out to be absorbed by a field redefinition using $`g_a(z)`$. As was shown in the previous section for supersymmetric case and in ref. for bosonic case, we can obtain nontrivial solutions in a global sense by considering a compactified background with the non-vanishing Wilson lines $`_{C_{\mathrm{left}}}𝑑zF_\mu (z)0`$. In general case, there is a possibility that our classical solutions become nontrivial if $`F_a(z)`$ breaks the condition (3.60), although we cannot prove their non-triviality with respect to gauge transformation of string field theory at this stage because we have only investigated a particular pure gauge form: (3.56) or (D). We speculate that they will be rewritten as a locally pure gauge form using a kind of integration of the current in each model in the case that $`g_a(z)`$ does not exist.
## 4 Discussions
We constructed a class of analytic classical solutions in open superstring field theory, which is related to marginal deformations in conformal field theory. We showed that the resulting solutions can be represented using a well-defined Fock space expression, and the vacuum energy vanishes due to the ghost number non-conservation in the large Hilbert space. For the solution corresponding to background Wilson lines, we observe that the solution can be written as a locally pure gauge form, and the action expanded around the solution can be transformed locally back to the original action by a string field redefinition. The analytic classical solution enabled us to investigate gauge structure in the string field theory. We found that the half integration mode of the function in the solution is unchanged under the “global transformation”, but other modes can be gauged away. Space-time supersymmetry is realized on-shell in the theory and the solution is a supersymmetric solution. We note that the classical solutions in the present paper can be easily extended to that of the theory given in ref. including the GSO$`()`$ sector
We have extensively used formal properties of the identity string field to construct our solution and to investigate its structure. In general, the identity string field requires careful handling to evaluate some quantities of the form $`(𝒪_1I)(𝒪_2I)`$. Actually, in terms of oscillator representation, we encounter divergence from contractions of nonzero modes in computing $`I|\mathrm{}|I`$ as in the case of bosonic string field theory. Consequently, it is necessary to define appropriate regularization of the identity string field, which is not yet known, in both super and bosonic string field theory. However, we evaluated the vacuum energy at our solution $`\mathrm{\Phi }_0`$ using the zero-mode saturation rule for the $`\xi \eta `$ ghost system in the large Hilbert space. We hope to obtain information about a consistent regularization by comparing these calculations. The situation is different in bosonic string field theory, where we gave a formal proof (D.18) to show that the vacuum energy at the solution corresponding to the marginal deformation vanishes as in refs. . It is preferable to prove it more directly.
We briefly comment on the modified version of cubic superstring field theory . Setting the Ramond field to zero, the equation of motion is given by $`Y_2(Q_\mathrm{B}A+AA)=0`$, where $`Y_2`$ is the picture changing operator with picture number $`(2)`$. We note that a string field $`A_0=e^{\mathrm{\Phi }_0}Q_\mathrm{B}e^{\mathrm{\Phi }_0}`$, in which $`\mathrm{\Phi }_0`$ is our solution in the Berkovits’ theory, is a solution to the equation of motion. In fact, $`Q_\mathrm{B}A_0+A_0A_0=0`$ holds and $`A_0`$ itself has ghost number $`1`$ and picture number $`0`$ and is Grassmann odd in the small Hilbert space. Therefore, we find that $`A_0`$ is pure gauge with the gauge parameter $`J_L^a(g_a)I`$ if $`F_a(z)`$ in $`\mathrm{\Phi }_0`$ satisfies the condition (3.60). We notice that $`\mathrm{\Phi }_0`$ itself cannot be a gauge parameter in spite of the form of $`A_0`$ because $`\mathrm{\Phi }_0`$ is not in the small Hilbert space : $`\eta _0\mathrm{\Phi }_00`$.
For the Wilson line solution, the half integration of the function is unchanged under the global transformation. This is an obvious result because the half integration mode as the Wilson line should be a physical observable. However, we confirmed the invariance of the half integration mode merely for a part of the gauge symmetry, and it is difficult to prove the invariance for the whole gauge symmetry. More precisely, we have to relate the solution to a general gauge invariant quantity in string field theory. Although a gauge invariant quantity plays important roles in field theories, we have not yet understood completely how it can be constructed in string field theory. As gauge invariants, we know the action and some operators only . It is natural to ask how the Wilson loop operator is generalized in string field theory. This is an important open question.
In the theory expanded around the classical solution, the background of the theory can be changed from the unexpanded theory. If we choose $`su(2)`$ currents, the solution corresponds to the tachyon lump solution (see also appendix D). In the background, the boundary condition of a string coordinate is changed from the Neumann one to the Dirichlet one . In string field theory, we found this phenomenon indirectly by studying a gauge invariant operator . On the other hand, in ref. , it was proposed that string coordinates $`X^\mu (\sigma )`$ and its conjugate momenta $`P_\mu (\sigma )`$ are universal objects in string field theory, and the various backgrounds correspond to inequivalent representations of their canonical algebra. In the present case, it seems that the tachyon lump solution changes representation of universal coordinates, namely their boundary conditions. This subject was studied from the viewpoint of vacuum string field theory .
It is an important problem to construct an analytic solution representing tachyon condensation in the superstring field theory. In order to investigate this subject, we have to introduce Chan-Paton matrices to include both GSO$`(+)`$ and GSO$`()`$ sectors. Such a formulation were developed in ref. . We hope to obtain an analytic expression of the tachyon vacuum by an analogous construction to the bosonic case (i.e., a class of scalar solutions in refs. ).
## Acknowledgements
The authors would like to thank Yuji Igarashi and Katsumi Itoh for useful discussions. I. K. wishes to express his gratitude to Kazuki Ohmori for valuable comments.
## Appendix A Oscillator expression of the identity string field
We use the identity string field in constructing exact solutions in the framework of Berkovits’ open superstring field theory. The action is described in terms of the large Hilbert space in the NSR formalism. We shall give an explicit oscillator representation of the identity string field $`|I`$ in terms of modes of $`X^\mu ,\psi ^\mu ,b,c,\varphi ,\xi ,\eta `$ in the NS sector. In this paper, we formally regard the identity string field $`I`$ as the identity element with respect to the Witten $``$ product: $`AI=IA=A`$, which should be proved by LPP’s definition of string vertices and generalized gluing and resmoothing theorem at the critical dimension $`d=10`$ . Therefore, we define $`I|`$ as a 1-string LPP vertex using conformal mapping: $`h_I(z)=2z/(1z^2)`$ and CFT correlator in the large Hilbert space denoted by $`\mathrm{}`$ which is evaluated on the upper half plane:<sup>20</sup><sup>20</sup>20 In this paper, we implicitly use the doubling trick: a holomorphic field $`\sigma (z)`$ and antiholomorphic one $`\stackrel{~}{\sigma }(\overline{z})`$ in the upper half plane are combined into a holomorphic field defined in the whole complex plane with a boundary condition $`\sigma (z)=\stackrel{~}{\sigma }(\overline{z})`$ on the real axis: $`\mathrm{Im}z=0`$. $`I|A=h_I[𝒪_A(0)]`$, where $`|A=𝒪_A(0)|0`$ is an arbitrary state. In ref. , the integral expression for Neumann coefficients in $`X^\mu ,b,c`$ sector is given with this definition and the oscillator expression of the identity string filed is obtained by applying the conformal map $`h_I(z)`$, which is consistent with that in . The result is
$`|I_b`$ $`=`$ $`e^{E_{Xbc}}|p^\mu =0,`$ (A.1)
$`E_{Xbc}`$ $`=`$ $`{\displaystyle \underset{n1}{}}{\displaystyle \frac{(1)^n}{2n}}\alpha _n^\mu \alpha _{n\mu }+{\displaystyle \underset{n2}{}}(1)^nc_nb_n{\displaystyle \underset{k1}{}}(1)^k(2c_0b_{2k}+(c_1c_1)b_{2k1}),`$
which is the same as the identity string field in the Witten’s bosonic open string field theory if $`\mu `$ runs over $`0,1,\mathrm{},25`$.
In the same way, we can calculate the Neumann coefficients in the matter fermion sector ($`\psi ^\mu `$) as:
$`I_{rs}`$ $`=`$ $`I_{sr}={\displaystyle _0}{\displaystyle \frac{dy}{2\pi i}}y^{r\frac{1}{2}}{\displaystyle _0}{\displaystyle \frac{dz}{2\pi i}}z^{s\frac{1}{2}}{\displaystyle \frac{(h_I^{}(y))^{\frac{1}{2}}(h_I^{}(z))^{\frac{1}{2}}}{h_I(y)h_I(z)}}`$ (A.5)
$`=`$ $`\{\begin{array}{cc}\frac{r(2s1)}{r^2s^2}\left(\frac{1}{4}\right)^{\frac{r+s}{2}}\frac{\left(r\frac{1}{2}\right)!\left(s\frac{3}{2}\right)!}{\left[\left(\frac{1}{2}(r\frac{1}{2})\right)!\left(\frac{1}{2}(s\frac{3}{2})\right)!\right]^2}& (r\frac{1}{2}:\mathrm{even};s\frac{1}{2}:\mathrm{odd})\\ \frac{s(2r1)}{r^2s^2}\left(\frac{1}{4}\right)^{\frac{r+s}{2}}\frac{\left(r\frac{3}{2}\right)!\left(s\frac{1}{2}\right)!}{\left[\left(\frac{1}{2}(r\frac{3}{2})\right)!\left(\frac{1}{2}(s\frac{1}{2})\right)!\right]^2}& (r\frac{1}{2}:\mathrm{odd};s\frac{1}{2}:\mathrm{even})\\ 0& (\mathrm{otherwise})\end{array},`$
where we have used the expansion:
$`{\displaystyle \underset{r,s\frac{1}{2}}{}}(r^2s^2)I_{rs}y^{r\frac{1}{2}}z^{s\frac{1}{2}}=(y_yz_z)(y_y+z_z+1)\left({\displaystyle \frac{(h_I^{}(y))^{\frac{1}{2}}(h_I^{}(z))^{\frac{1}{2}}}{h_I(y)h_I(z)}}{\displaystyle \frac{1}{yz}}\right)`$
$`={\displaystyle \frac{y(1z^2)+z(1y^2)}{(1+y^2)^{\frac{3}{2}}(1+z^2)^{\frac{3}{2}}}}={\displaystyle \underset{k,l=0}{\overset{\mathrm{}}{}}}(2k+1)(4l+1)(\begin{array}{c}\frac{1}{2}\\ k\end{array})(\begin{array}{c}\frac{1}{2}\\ l\end{array})(y^{2k+1}z^{2l}+y^{2l}z^{2k+1}).`$ (A.10)
$`\left(\begin{array}{c}a\\ b\end{array}\right)=\frac{\mathrm{\Gamma }(a+1)}{\mathrm{\Gamma }(b+1)\mathrm{\Gamma }(ab+1)}`$ is the binomial coefficient. This formula for the coefficients $`I_{rs}`$ is consistent with that in .
As for the $`\varphi `$ sector, the formula for the Neumann coefficients is slightly different from that of $`X^\mu `$ because of the background charge $`Q=2`$ . One can compute explicitly by substituting $`h_I(z)`$ into the integrand:
$`𝒩_{mn}`$ $`=`$ $`{\displaystyle \frac{1}{mn}}{\displaystyle _0}{\displaystyle \frac{dy}{2\pi i}}y^m{\displaystyle _0}{\displaystyle \frac{dz}{2\pi i}}z^n{\displaystyle \frac{h_I^{}(y)h_I^{}(z)}{(h_I(y)h_I(z))^2}}={\displaystyle \frac{(1)^m}{m}}\delta _{m,n},(m,n1),`$ (A.11)
$`𝒩_{0n}`$ $`=`$ $`{\displaystyle \frac{1}{2n}}{\displaystyle _0}{\displaystyle \frac{dw}{2\pi i}}w^n_w\mathrm{log}(_y_w\mathrm{log}(h_I(y)h_I(w))|_{y=0})`$ (A.14)
$`=`$ $`\{\begin{array}{cc}\frac{(1)^{\frac{n}{2}}}{n}& (n:\mathrm{even})\\ 0& (n:\mathrm{odd})\end{array},(n1),`$
$`𝒩_{00}`$ $`=`$ $`0,`$ (A.15)
In the $`\xi \eta `$ sector, the Neumann coefficients for the identity string field are computed as<sup>21</sup><sup>21</sup>21 In general, $`N`$-string vertex is given by $`V_N|`$ $`=`$ $`{\displaystyle \underset{r}{}}{}_{r}{}^{}0|\xi _0^re^{_{r,s=1}^N_{n0,m1}\eta _n^rN_{nm}^{rs}\xi _m^s}(\eta _0^1+\mathrm{}+\eta _0^N),`$ (A.16) in this sector, which is obtained using the method in .
$`N_{mn}={\displaystyle _0}{\displaystyle \frac{dy}{2\pi i}}y^{m1}{\displaystyle _0}{\displaystyle \frac{dz}{2\pi i}}z^n{\displaystyle \frac{h_I^{}(z)}{h_I(y)h_I(z)}}=(1)^m\delta _{m,n},m,n1,`$ (A.17)
$`N_{0n}={\displaystyle _0}{\displaystyle \frac{dy}{2\pi i}}y^1{\displaystyle _0}{\displaystyle \frac{dz}{2\pi i}}z^n{\displaystyle \frac{h_I^{}(z)}{h_I(y)h_I(z)}}=\{\begin{array}{cc}2& (n:\mathrm{even})\\ 0& (n:\mathrm{odd})\end{array},n1,`$ (A.20)
where use has been made of the expansion
$`{\displaystyle \frac{h_I^{}(z)}{h_I(y)h_I(z)}}{\displaystyle \frac{1}{yz}}={\displaystyle \frac{y(1+z^2)+2z}{(1z^2)(1+yz)}}={\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}(2z^{2k+1}+(1)^ky^{k+1}z^k).`$ (A.21)
After all, the identity string field of the fermionic sector ($`\psi ^\mu ,\varphi ,\xi ,\eta `$) is given by
$`|I_f`$ $`=`$ $`e^{E_{\psi \varphi \xi \eta }}|q=0,`$ (A.22)
$`E_{\psi \varphi \xi \eta }`$ $`=`$ $`{\displaystyle \underset{r,s1/2}{}}{\displaystyle \frac{I_{rs}}{2}}\psi _r^\mu \psi _{s\mu }+{\displaystyle \underset{n1}{}}{\displaystyle \frac{(1)^n}{2n}}(j_n)^2{\displaystyle \underset{k1}{}}{\displaystyle \frac{(1)^k}{k}}j_{2k}+{\displaystyle \underset{n1}{}}(1)^n\eta _n\xi _n,`$
where the oscillators are given by
$`\psi ^\mu (z)={\displaystyle \underset{r}{}}\psi _r^\mu z^{r1/2},\{\psi _r^\mu ,\psi _s^\nu \}=\eta ^{\mu \nu }\delta _{r+s,0},`$ (A.23)
$`\varphi (z)=\widehat{\varphi }_0j_0\mathrm{log}z+{\displaystyle \underset{n0}{}}{\displaystyle \frac{1}{n}}j_nz^n,[j_0,\widehat{\varphi }_0]=1,[j_m,j_n]=m\delta _{m+n,0},`$ (A.24)
$`\xi (z)={\displaystyle \underset{n}{}}\xi _nz^n,\eta (z)={\displaystyle \underset{n}{}}\eta _nz^{n1},\{\xi _m,\eta _n\}=\delta _{m+n,0},`$ (A.25)
and the vacuum with $`\varphi `$-charge $`q`$ is defined by $`|q=e^{q\varphi }(0)|0=e^{q\widehat{\varphi }_0}|0`$.
Combining (A.1) and (A.22), the identity string field in open superstring field theory in the large Hilbert space is obtained:
$`|I`$ $`=`$ $`|I_b|I_f.`$ (A.26)
The index $`\mu `$ in the exponent of $`|I_b`$ runs over $`0,1,\mathrm{},9`$ and $`|p^\mu =0|q=0`$ is the conformal vacuum. $`|I`$ is Grassmann even and has both ghost and picture number $`0`$. BRS invariance $`Q_\mathrm{B}|I=0`$ follows from the construction of LPP vertex and $`\eta _0|I=0`$ can be checked directly. We can easily derive the following connection conditions of each oscillators on the identity string field $`|I`$ using the above explicit expression:<sup>22</sup><sup>22</sup>22 For a derivation using CFT, see ref. .
$`(\alpha _n^\mu +(1)^n\alpha _n^\mu )|I=0,(b_n(1)^nb_n)|I=0,`$ (A.27)
$`(c_{2k}+c_{2k}(1)^k2c_0)|I=0,(c_{2k+1}c_{(2k+1)}(1)^k(c_1c_1))|I=0,`$ (A.28)
$`\left(\psi _r^\mu {\displaystyle \underset{s1/2}{}}I_{rs}\psi _s^\mu \right)|I=0,(\xi _n(1)^n\xi _n)|I=0,(\eta _n+(1)^n\eta _n)|I=0,`$ (A.29)
$`(j_{2k}+j_{2k}(1)^k2)|I=0,(j_{2k1}j_{(2k1)})|I=0,(k1);j_0|I=0.`$ (A.30)
The identity string field $`|I`$ satisfies the reality condition: $`(|I)^{}=\mathrm{bpz}(|I)`$, where the BPZ conjugation is given by $`\mathrm{bpz}(\sigma _n)=(1)^{n+h}\sigma _n`$ for oscillators of a primary field $`\sigma (z)`$ with conformal dimension $`h`$ and $`\mathrm{bpz}(|p^\mu ;q)=(|p^\mu ;q)^{}`$ for zero mode part and use has been made of $`(1)^{r+s}I_{rs}=I_{rs}`$. We note that the identity string field $`|I`$ can be rewritten as
$`|I`$ $`=`$ $`{\displaystyle \frac{(2i)^{\frac{1}{4}}}{4i}}b(i\pi /2)b(i\pi /2):e^{\frac{1}{2}\varphi (i\pi /2)}::e^{\frac{1}{2}\varphi (i\pi /2)}:e^E^{}c_0c_1|p^\mu =0;q=1,`$ (A.31)
$`E^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{n1}{}}{\displaystyle \frac{(1)^n}{n}}\alpha _n^\mu \alpha _{n\mu }+{\displaystyle \underset{n1}{}}(1)^nc_nb_n`$ (A.32)
$`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{r,s1/2}{}}I_{rs}\psi _r^\mu \psi _{s\mu }+{\displaystyle \frac{1}{2}}{\displaystyle \underset{n1}{}}{\displaystyle \frac{(1)^n}{n}}j_nj_n+{\displaystyle \underset{n1}{}}(1)^n\eta _n\xi _n,`$
where we denoted
$`b(i\sigma )={\displaystyle \underset{n}{}}b_ne^{in\sigma },`$ (A.33)
$`:e^{q\varphi (i\sigma )}:=e^{\frac{1}{2}q(q+2)i\sigma }e^{q_{n1}\frac{1}{n}j_ne^{in\sigma }}e^{q\widehat{\varphi }_0}e^{iq\sigma j_0}e^{q_{n1}\frac{1}{n}j_ne^{in\sigma }}.`$ (A.34)
The extra factor $`e^{\frac{1}{2}q(q+2)i\sigma }`$ in the normal order form comes from the conformal factor under the map $`z=e^\rho `$.
## Appendix B Action around a classical solution
The Berkovits’ action for open superstring field theory in the NS sector is given by
$`S[\mathrm{\Phi }]`$ $`=`$ $`{\displaystyle \frac{1}{2g^2}}\overline{A}_{\eta _0}\overline{A}_Q{\displaystyle \frac{1}{2g^2}}{\displaystyle _0^1}𝑑tA_t\{A_Q,A_{\eta _0}\},`$ (B.1)
where $`A_{\eta _0},A_Q`$ and $`A_t`$ are defined by string field $`\mathrm{\Phi }(t)`$ parametrized by $`t`$ with boundary value $`\mathrm{\Phi }(1)=\mathrm{\Phi },\mathrm{\Phi }(0)=0`$ :
$`A_{\eta _0}=e^{\mathrm{\Phi }(t)}(\eta _0e^{\mathrm{\Phi }(t)}),A_Q=e^{\mathrm{\Phi }(t)}(Qe^{\mathrm{\Phi }(t)}),A_t=e^{\mathrm{\Phi }(t)}(_te^{\mathrm{\Phi }(t)}),`$ (B.2)
and $`\overline{A}_{\eta _0}=A_{\eta _0}|_{t=1},\overline{A}_Q=A_Q|_{t=1}`$. We usually take $`\mathrm{\Phi }(t)=t\mathrm{\Phi }`$ although the action $`S[\mathrm{\Phi }]`$ itself does not depend on this parameterization. We often denote $`\{A,B\}=AB+BA`$ and $`[A,B]=ABBA`$ and omit the symbol for the star product among string fields. We note that $`\eta _0,Q`$ and $`_t`$ are derivations with respect to the star product:
$`\eta _0(AB)`$ $`=`$ $`(\eta _0A)B+(1)^{|A|}A(\eta _0B),`$ (B.3)
$`Q(AB)`$ $`=`$ $`(QA)B+(1)^{|A|}A(QB),`$ (B.4)
$`_t(AB)`$ $`=`$ $`(_tA)B+A(_tB),`$ (B.5)
where $`(1)^{|A|}`$ is $`+1(1)`$ when $`A`$ is Grassmann even (odd) and have nilpotency $`\eta _0^2=Q^2=0`$ and (anti-)commutativity: $`\{\eta _0,Q\}=0,[_t,\eta _0]=[_t,Q]=0`$. The above WZW type action (B.1) can be rewritten in a rather simple form :
$`S[\mathrm{\Phi }]`$ $`=`$ $`{\displaystyle \frac{1}{g^2}}{\displaystyle _0^1}𝑑t(\eta _0A_t)A_Q.`$ (B.6)
Let us consider re-expansion of this action around $`\mathrm{\Phi }^{(0)}`$ with respect to $`\mathrm{\Phi }^{}`$ in the sense $`e^{\mathrm{\Phi }(t)}=e^{\mathrm{\Phi }^{(0)}(t)}e^{\mathrm{\Phi }^{}(t)}`$. The integrand of (B.6) can be rewritten as:
$`(\eta _0A_t)A_Q`$ $`=`$ $`(e^{\mathrm{\Phi }^{}(t)}(\eta _0A_t^{(0)})e^{\mathrm{\Phi }^{}(t)}+\eta _0A_t^{}+(\eta _0e^{\mathrm{\Phi }^{}(t)})A_t^{(0)}e^{\mathrm{\Phi }^{}(t)}+e^{\mathrm{\Phi }^{}(t)}A_t^{(0)}(\eta _0e^{\mathrm{\Phi }^{}(t)}))`$ (B.7)
$`\times (e^{\mathrm{\Phi }^{}(t)}A_Q^{(0)}e^{\mathrm{\Phi }^{}(t)}+A_Q^{})`$
$`=`$ $`(\eta _0A_t^{(0)})A_Q^{(0)}+(\eta _0A_t^{})A_Q^{}`$
$`+(\eta _0A_t^{})e^{\mathrm{\Phi }^{}(t)}A_Q^{(0)}e^{\mathrm{\Phi }^{}(t)}e^{\mathrm{\Phi }^{}(t)}(\eta _0e^{\mathrm{\Phi }^{}(t)})_tA_Q^{(0)}`$
$`+(\eta _0A_t^{(0)})e^{\mathrm{\Phi }^{}(t)}A_Q^{}e^{\mathrm{\Phi }^{}(t)}+(\eta _0e^{\mathrm{\Phi }^{}(t)})A_t^{(0)}e^{\mathrm{\Phi }^{}(t)}A_Q^{}+e^{\mathrm{\Phi }^{}(t)}A_t^{(0)}(\eta _0e^{\mathrm{\Phi }^{}(t)})A_Q^{}`$
$`+e^{\mathrm{\Phi }^{}(t)}(\eta _0e^{\mathrm{\Phi }^{}(t)})(QA_t^{(0)}),`$
where we denote $`A_t^{(0)}=e^{\mathrm{\Phi }^{(0)}(t)}(_te^{\mathrm{\Phi }^{(0)}(t)})`$, $`A_Q^{(0)}=e^{\mathrm{\Phi }^(0)(t)}(Qe^{\mathrm{\Phi }^{(0)}(t)})`$, $`A_t^{}=e^{\mathrm{\Phi }^{}(t)}(_te^{\mathrm{\Phi }^{}(t)})`$, and $`A_Q^{}=e^{\mathrm{\Phi }^{}(t)}(Qe^{\mathrm{\Phi }^{}(t)})`$ and use has been made of cyclic property:
$`A_1\mathrm{}A_{n1}\mathrm{\Phi }=\mathrm{\Phi }A_1\mathrm{}A_{n1},`$ (B.8)
$`A_1\mathrm{}A_{n1}(Q\mathrm{\Phi })=(Q\mathrm{\Phi })A_1\mathrm{}A_{n1},`$ (B.9)
$`A_1\mathrm{}A_{n1}(\eta _0\mathrm{\Phi })=(\eta _0\mathrm{\Phi })A_1\mathrm{}A_{n1},`$ (B.10)
and an identity $`[A_t^{(0)},A_Q^{(0)}]=QA_t^{(0)}_tA_Q^{(0)}`$. Using $`e^{\mathrm{\Phi }^{}(t)}(\eta _0A_t^{})e^{\mathrm{\Phi }^{}(t)}=_t(e^{\mathrm{\Phi }^{}(t)}(\eta _0e^{\mathrm{\Phi }^{}(t)}))`$, $`Q(e^{\mathrm{\Phi }^{}(t)}(\eta _0e^{\mathrm{\Phi }^{}(t)}))=e^{\mathrm{\Phi }^{}(t)}(\eta _0A_Q^{})e^{\mathrm{\Phi }^{}(t)}`$ and partial integrability:
$`Q(\mathrm{})=0,\eta _0(\mathrm{})=0,`$ (B.11)
we can simplify the last three lines of (B.7) as:
$`(\eta _0A_t^{})e^{\mathrm{\Phi }^{}(t)}A_Q^{(0)}e^{\mathrm{\Phi }^{}(t)}e^{\mathrm{\Phi }^{}(t)}(\eta _0e^{\mathrm{\Phi }^{}(t)})_tA_Q^{(0)}=_te^{\mathrm{\Phi }^{}(t)}(\eta _0e^{\mathrm{\Phi }^{}(t)})A_Q^{(0)},`$ (B.12)
$`(\eta _0A_t^{(0)})e^{\mathrm{\Phi }^{}(t)}A_Q^{}e^{\mathrm{\Phi }^{}(t)}+(\eta _0e^{\mathrm{\Phi }^{}(t)})A_t^{(0)}e^{\mathrm{\Phi }^{}(t)}A_Q^{}+e^{\mathrm{\Phi }^{}(t)}A_t^{(0)}(\eta _0e^{\mathrm{\Phi }^{}(t)})A_Q^{}`$
$`=A_t^{(0)}e^{\mathrm{\Phi }^{}(t)}(\eta _0A_Q^{})e^{\mathrm{\Phi }^{}(t)},`$ (B.13)
$`e^{\mathrm{\Phi }^{}(t)}(\eta _0e^{\mathrm{\Phi }^{}(t)})(QA_t^{(0)})=A_t^{(0)}e^{\mathrm{\Phi }^{}(t)}(\eta _0A_Q^{})e^{\mathrm{\Phi }^{}(t)}.`$ (B.14)
Then we have proved an identity
$`(\eta _0A_t)A_Q=(\eta _0A_t^{(0)})A_Q^{(0)}+(\eta _0A_t^{})A_Q^{}_te^{\mathrm{\Phi }^{}(t)}(\eta _0e^{\mathrm{\Phi }^{}(t)})A_Q^{(0)},`$ (B.15)
which implies the action (B.6) is rewritten for $`e^\mathrm{\Phi }=e^{\mathrm{\Phi }^{(0)}}e^\mathrm{\Phi }^{}`$ as
$`S[\mathrm{\Phi }]=S[\mathrm{\Phi }^{(0)}]+S[\mathrm{\Phi }^{}]+{\displaystyle \frac{1}{g^2}}e^\mathrm{\Phi }^{}(\eta _0e^\mathrm{\Phi }^{})\overline{A}_Q^{(0)},`$ (B.16)
where we have imposed ordinary boundary conditions $`\mathrm{\Phi }^{(0)}(1)=\mathrm{\Phi }^{(0)},\mathrm{\Phi }^{}(1)=\mathrm{\Phi }^{},\mathrm{\Phi }^{(0)}(0)=\mathrm{\Phi }^{}(0)=0`$ and denoted $`\overline{A}_Q^{(0)}=A_Q^{(0)}|_{t=1}`$. The last extra term of the above action can be rewritten as follows:
$`{\displaystyle \frac{1}{g^2}}e^\mathrm{\Phi }^{}(\eta _0e^\mathrm{\Phi }^{})\overline{A}_Q^{(0)}`$ (B.17)
$`=`$ $`{\displaystyle \frac{1}{g^2}}{\displaystyle _0^1}𝑑t_te^{\mathrm{\Phi }^{}(t)}(\eta _0e^{\mathrm{\Phi }^{}(t)})\overline{A}_Q^{(0)}={\displaystyle \frac{1}{g^2}}{\displaystyle _0^1}𝑑te^{\mathrm{\Phi }^{}(t)}(\eta _0A_t^{})e^{\mathrm{\Phi }^{}(t)}\overline{A}_Q^{(0)}`$
$`=`$ $`{\displaystyle \frac{1}{g^2}}{\displaystyle _0^1}𝑑t(\eta _0A_t^{})(e^{\mathrm{\Phi }^{}(t)}\overline{A}_Q^{(0)}e^{\mathrm{\Phi }^{}(t)}\overline{A}_Q^{(0)})A_t^{}(\eta _0\overline{A}_Q^{(0)}).`$
In the first equality, we have kept $`\overline{A}_Q^{(0)}`$ intact in $`t`$-integration and used $`e^{\mathrm{\Phi }^{}(t)}(\eta _0A_t^{})e^{\mathrm{\Phi }^{}(t)}=_t(e^{\mathrm{\Phi }^{}(t)}(\eta _0e^{\mathrm{\Phi }^{}(t)}))`$ again in the second equality. Using (B.6), (B.16) and (B.17), and imposing equation of motion for $`\mathrm{\Phi }^{(0)}`$: $`\eta _0\overline{A}_Q^{(0)}=0`$, we have obtained the action for $`\mathrm{\Phi }^{}`$ around a classical solution $`\mathrm{\Phi }^{(0)}`$ in the same form as the original one:
$`S^{}[\mathrm{\Phi }^{}]`$ $``$ $`S[\mathrm{\Phi }]S[\mathrm{\Phi }^{(0)}]={\displaystyle \frac{1}{g^2}}{\displaystyle _0^1}𝑑t(\eta _0A_t^{})A_Q^{}^{},`$ (B.18)
with $`A_Q^{}^{}=e^{\mathrm{\Phi }^{}(t)}(Q^{}e^{\mathrm{\Phi }^{}(t)})`$, where the new BRS operator $`Q^{}`$ is given by
$`Q^{}B`$ $`=`$ $`QB+\overline{A}_Q^{(0)}B(1)^{|B|}B\overline{A}_Q^{(0)},\overline{A}_Q^{(0)}=e^{\mathrm{\Phi }^{(0)}}(Qe^{\mathrm{\Phi }^{(0)}}).`$ (B.19)
We note that $`Q^{}`$ is a derivation with respect to the star product, nilpotency $`Q^2=0`$ holds automatically, and $`\{Q^{},\eta _0\}=0`$ is satisfied by equation of motion $`\eta _0\overline{A}_Q^{(0)}=0`$ for $`\mathrm{\Phi }^{(0)}`$. The above action can be rewritten in the ordinary WZW form again
$`S^{}[\mathrm{\Phi }^{}]`$ $`=`$ $`{\displaystyle \frac{1}{2g^2}}\overline{A}_{\eta _0}^{}\overline{A}_Q^{}^{}{\displaystyle \frac{1}{2g^2}}{\displaystyle _0^1}𝑑tA_t^{}\{A_Q^{}^{},A_{\eta _0}^{}\},`$ (B.20)
where $`A_{\eta _0}^{}=e^{\mathrm{\Phi }^{}(t)}(\eta _0e^{\mathrm{\Phi }^{}(t)}),\overline{A}_{\eta _0}^{}=A_{\eta _0}^{}|_{t=0},\overline{A}_Q^{}^{}=A_Q^{}^{}|_{t=1}`$, using the method in .
## Appendix C Supersymmetry in superstring field theories
First, we will show that the fermionic transformation (2.80), which is generated by (2.79), corresponds to global space-time supersymmetries.
For the parameter (2.79), the transformation law given in (2.77) and (2.78) becomes
$`\delta _ϵ\mathrm{\Phi }`$ $`=`$ $`{\displaystyle \frac{\mathrm{ad}_\mathrm{\Phi }}{1e^{\mathrm{ad}_\mathrm{\Phi }}}}\{\mathrm{\Omega }(ϵ),\eta _0\mathrm{\Psi }\}={\displaystyle \frac{\mathrm{ad}_\mathrm{\Phi }}{1e^{\mathrm{ad}_\mathrm{\Phi }}}}\left({\displaystyle \frac{dz}{2\pi i}ϵ_\alpha \xi S_{(1/2)}^\alpha (z)\eta _0\mathrm{\Psi }}\right),`$ (C.1)
$`\delta _ϵ(\eta _0\mathrm{\Psi })`$ $`=`$ $`\eta _0\{\mathrm{\Omega }(ϵ),e^\mathrm{\Phi }Q_\mathrm{B}e^\mathrm{\Phi }\}=\eta _0{\displaystyle \frac{dz}{2\pi i}ϵ_\alpha \xi S_{(1/2)}^\alpha (z)(e^\mathrm{\Phi }Q_\mathrm{B}e^\mathrm{\Phi })},`$ (C.2)
where the term $`\eta _0Q_\mathrm{B}\mathrm{\Omega }(ϵ)`$ is not included as is explained in footnote 14 in order to express it only in terms of a contour integration (2.81). These transformations are equivalent to (2.80). The equations of motion are given by
$`f_1`$ $``$ $`\eta _0(e^\mathrm{\Phi }Q_\mathrm{B}e^\mathrm{\Phi })+(\eta _0\mathrm{\Psi })^2=0,`$ (C.3)
$`f_2`$ $``$ $`e^\mathrm{\Phi }(Q_\mathrm{B}(e^\mathrm{\Phi }(\eta _0\mathrm{\Psi })e^\mathrm{\Phi }))e^\mathrm{\Phi }=Q_\mathrm{B}\eta _0\mathrm{\Psi }+\{e^\mathrm{\Phi }Q_\mathrm{B}e^\mathrm{\Phi },\eta _0\mathrm{\Psi }\}=0.`$ (C.4)
Note that the R sector string field $`\mathrm{\Psi }`$ is involved in the equations of motion through the particular form $`\eta _0\mathrm{\Psi }`$. We apply the transformation (2.80) to the string fields $`f_1`$ and $`f_2`$:
$`\delta _ϵf_1`$ $`=`$ $`\eta _0𝒮(ϵ)f_2,`$ (C.5)
$`\delta _ϵf_2`$ $`=`$ $`\{Q_\mathrm{B},𝒮(ϵ)\}f_1+[f_1,𝒮(ϵ)(e^\mathrm{\Phi }Q_\mathrm{B}e^\mathrm{\Phi })]+\{𝒮(ϵ)f_2,\eta _0\mathrm{\Psi }\}.`$ (C.6)
If $`f_1=f_2=0`$, we find $`\delta _ϵf_1=\delta _ϵf_2=0`$. Hence this symmetry is realized only on-shell.
Let us consider how massless fields are transformed by (2.80). The string fields contain massless fields as follows:
$`|\mathrm{\Phi }_A`$ $`=`$ $`{\displaystyle \frac{d^{10}p}{(2\pi )^{10}}(\stackrel{~}{A}_\mu (p)c\xi e^\varphi \psi ^\mu (0)+\stackrel{~}{B}(p)cc\xi \xi e^{2\varphi }(0))|p^\mu ,q=0},`$ (C.7)
$`|\mathrm{\Psi }_\lambda `$ $`=`$ $`{\displaystyle \frac{d^{10}p}{(2\pi )^{10}}\stackrel{~}{\lambda }_\alpha (p)\xi S_{(1/2)}^\alpha c(0)|p^\mu ,q=0},`$ (C.8)
where $`\stackrel{~}{A}_\mu (p),\stackrel{~}{B}(p)`$ and $`\stackrel{~}{\lambda }_\alpha (p)`$ denote Fourier transforms of gluon, auxiliary Nakanishi-Lautrup and gluino fields, respectively. We note that $`q`$ in $`|p^\mu ,q`$ implies a zero-mode momentum of $`\varphi `$. Applying the transformation (2.80) to these fields, we find
$`\delta _ϵ|\mathrm{\Phi }_A`$ $`=`$ $`{\displaystyle \frac{dz}{2\pi i}ϵ_\alpha \xi S_{(1/2)}^\alpha (z)\eta _0|\mathrm{\Psi }_\lambda }+\mathrm{}`$ (C.9)
$`=`$ $`{\displaystyle \frac{d^{10}p}{(2\pi )^{10}}(iϵ_\alpha (\mathrm{\Gamma }_\mu C)^{\alpha \beta }\stackrel{~}{\lambda }_\beta (p))c\xi e^\varphi \psi ^\mu (0)|p^\mu ,0}+\mathrm{},`$
$`\delta _ϵ(\eta _0|\mathrm{\Psi }_\lambda )`$ $`=`$ $`\eta _0{\displaystyle \frac{dz}{2\pi i}ϵ_\alpha \xi S_{(1/2)}^\alpha (z)Q_\mathrm{B}|\mathrm{\Phi }_A}+\mathrm{}`$ (C.11)
$`=`$ $`{\displaystyle }{\displaystyle \frac{d^{10}p}{(2\pi )^{10}}}ϵ_\alpha ({\displaystyle \frac{1}{4}}\sqrt{2\alpha ^{}}(p_\mu \stackrel{~}{A}_\nu (p)p_\nu \stackrel{~}{A}_\mu (p))(\mathrm{\Gamma }^{\mu \nu })_\beta ^\alpha `$
$`+(\sqrt{{\displaystyle \frac{\alpha ^{}}{2}}}p^\mu \stackrel{~}{A}_\mu (p)+\stackrel{~}{B}(p))\delta _\beta ^\alpha )S_{(1/2)}^\beta c(0)|p^\mu ,0+\mathrm{},`$
where $`(+\mathrm{})`$ denotes quadratic or higher order terms with respect to component fields. Here we have calculated the above results by using the OPEs,<sup>23</sup><sup>23</sup>23 We have used the convention in for spin fields and taken $`T^\mathrm{m}(z)=\frac{1}{4\alpha ^{}}X_\mu X^\mu (z)\frac{1}{2}\psi ^\mu \psi _\mu (z),G^\mathrm{m}(z)=\frac{i}{\sqrt{2\alpha ^{}}}\psi ^\mu X_\mu (z)`$ with $`X^\mu (y)X^\nu (z)2\alpha ^{}\eta ^{\mu \nu }\mathrm{log}(yz),\psi ^\mu (y)\psi ^\nu (z)\eta ^{\mu \nu }(yz)^1`$ in the matter sector and $`j_\mathrm{B}(z)=c(T^\mathrm{m}\frac{1}{2}(\varphi )^2^2\varphi +\xi \eta )(z)+bcc(z)+\eta e^\varphi G^\mathrm{m}(z)\eta \eta e^{2\varphi }b(z)+^2c(z)+(c\xi \eta )(z)`$ for the BRS current.
$`S_{(1/2)}^\alpha (y)S_{(1/2)}^\beta (z){\displaystyle \frac{1}{yz}}i(\mathrm{\Gamma }_\mu C)^{\alpha \beta }\psi ^\mu e^\varphi (z),`$ (C.12)
$`j_\mathrm{B}(y)c\xi e^\varphi \psi ^\mu e^{ip_\nu X^\nu }(z){\displaystyle \frac{1}{(yz)^2}}\sqrt{2\alpha ^{}}p^\mu ce^{ip_\nu X^\nu }(z)+{\displaystyle \frac{1}{yz}}({\displaystyle \frac{i}{\sqrt{2\alpha ^{}}}}cX^\mu +\eta ce^\varphi \psi ^\mu `$
$`+\sqrt{2\alpha ^{}}c(p_\nu \psi ^\nu \psi ^\mu +p^\mu (\varphi \xi \eta ))\alpha ^{}p^2cc\xi e^\varphi \psi ^\mu )e^{ip_\nu X^\nu }(z),`$ (C.13)
$`j_\mathrm{B}(y)cc\xi \xi e^{2\varphi }e^{ip_\nu X^\nu }(z){\displaystyle \frac{1}{(yz)^2}}ce^{ip_\nu X^\nu }(z)`$ (C.14)
$`+{\displaystyle \frac{1}{yz}}(c+2c(\varphi \xi \eta )+\sqrt{2\alpha ^{}}cc\xi p_\mu \psi ^\mu e^\varphi )e^{ip_\nu X^\nu }(z),`$
$`\xi S_{(1/2)}^\alpha (y)c(\varphi \xi \eta )(z){\displaystyle \frac{1}{yz}}\left({\displaystyle \frac{1}{2}}\xi S_{(1/2)}^\alpha c\right)(z),`$ (C.15)
$`\xi S_{(1/2)}^\alpha (y)c\psi ^\nu \psi ^\mu (z){\displaystyle \frac{1}{yz}}{\displaystyle \frac{1}{2}}\xi (\mathrm{\Gamma }^{\nu \mu }S_{(1/2)})^\alpha c(z).`$ (C.16)
From (C.9) and (C.11), we can read off the transformation law for massless fields:
$`\delta _ϵ\stackrel{~}{A}_\mu (p)`$ $`=`$ $`iϵ\mathrm{\Gamma }_\mu C\stackrel{~}{\lambda }(p)+\mathrm{},`$ (C.17)
$`\delta _ϵ\stackrel{~}{B}(p)`$ $`=`$ $`0+\mathrm{},`$ (C.18)
$`\delta _ϵ\stackrel{~}{\lambda }_\alpha (p)`$ $`=`$ $`{\displaystyle \frac{\sqrt{2\alpha ^{}}}{4}}(p_\mu \stackrel{~}{A}_\nu (p)p_\nu \stackrel{~}{A}_\mu (p))(ϵ\mathrm{\Gamma }^{\mu \nu })_\alpha \left(\sqrt{{\displaystyle \frac{\alpha ^{}}{2}}}p^\mu \stackrel{~}{A}_\mu (p)+\stackrel{~}{B}(p)\right)ϵ_\alpha +\mathrm{}.`$ (C.19)
For massless fields, the linearized equations of motion are calculated as
$`Q_\mathrm{B}\eta _0|\mathrm{\Phi }_A`$ $`=`$ $`{\displaystyle }{\displaystyle \frac{d^{10}p}{(2\pi )^{10}}}((\sqrt{2\alpha ^{}}p^\mu \stackrel{~}{A}_\mu (p)+2\stackrel{~}{B}(p))c\eta (0)`$ (C.20)
$`+(\alpha ^{}p^2\stackrel{~}{A}(p)+\sqrt{2\alpha ^{}}p_\mu \stackrel{~}{B}(p))cce^\varphi \psi ^\mu (0))|p^\nu ,0=0,`$
$`Q_\mathrm{B}\eta _0|\mathrm{\Psi }_\lambda `$ $`=`$ $`{\displaystyle }{\displaystyle \frac{d^{10}p}{(2\pi )^{10}}}(\alpha ^{}p^2\stackrel{~}{\lambda }_\alpha (p)S_{(1/2)}^\alpha cc(0)`$ (C.21)
$`+\sqrt{\alpha ^{}}ip_\mu \stackrel{~}{\lambda }_\alpha (p)(\mathrm{\Gamma }^\mu )_{\dot{\beta }}^\alpha S_{(1/2)}^{\dot{\beta }}\eta c(0))|p^\nu ,0=0,`$
where $`S_{(1/2)}^{\dot{\beta }}`$ is the GSO$`(+)`$ spin operator with $`\varphi `$-charge $`1/2`$, dimension $`0`$ and negative chirality. Consequently, we obtain the linearized equations of motion:
$`\stackrel{~}{B}(p)=\sqrt{{\displaystyle \frac{\alpha ^{}}{2}}}p^\mu \stackrel{~}{A}_\mu (p),(p^2\delta _\mu ^\nu p_\mu p^\nu )\stackrel{~}{A}_\nu (p)=0,`$ (C.22)
$`\stackrel{~}{\lambda }_\alpha (p)p^2=0,\stackrel{~}{\lambda }(p)\mathrm{\Gamma }^\mu p_\mu =0.`$ (C.23)
Under these on-shell conditions, the transformation laws (C.17), (C.18) and (C.19) become
$`\delta _ϵA_\mu =iϵ\mathrm{\Gamma }_\mu C\lambda ,\delta _ϵ\lambda ={\displaystyle \frac{i}{2}}\sqrt{{\displaystyle \frac{\alpha ^{}}{2}}}F_{\mu \nu }(ϵ\mathrm{\Gamma }^{\mu \nu }),(F_{\mu \nu }=_\mu A_\nu _\nu A_\mu ).`$ (C.24)
These are nothing but supersymmetry transformation of 10D supersymmetric Maxwell theory. Hence, the transformation law (2.80) contains space-time supersymmetries.
Finally, we would like to comment on supersymmetry in the cubic open superstring field theory and the modified cubic theory .
#### Cubic version
In the Witten’s open superstring field theory, there are fermionic gauge symmetry and global supersymmetry at least formally . Fermionic gauge symmetry, which is generated by Grassmann even gauge parameter $`\chi `$ in the Ramond sector with picture number $`1/2`$ and ghost number $`0`$, is given by
$`\delta _\chi A`$ $`=`$ $`\mathrm{\Psi }\chi \chi \mathrm{\Psi },`$ (C.25)
$`\delta _\chi \mathrm{\Psi }`$ $`=`$ $`Q_\mathrm{B}\chi +X(i)(A\chi \chi A),`$ (C.26)
where $`X(i)`$ is the picture changing operator at the midpoint, and $`A(\mathrm{\Psi })`$ denotes a Grassmann odd string field in the NS (R) sector with picture number $`1(1/2)`$ and ghost number $`1(1)`$. Formally, by taking
$`\chi `$ $`=`$ $`{\displaystyle _{C_{\mathrm{left}}}}{\displaystyle \frac{dz}{2\pi i}}ϵ_\alpha S_{(1/2)}^\alpha (z)I,`$ (C.27)
and omitting $`Q_\mathrm{B}\chi =\frac{1}{2\pi i}ϵ_\alpha (cS_{(1/2)}^\alpha (i)cS_{(1/2)}^\alpha (i))I`$ in $`\delta _\chi \mathrm{\Psi }`$, which is itself a symmetry of the action, the above gauge transformation becomes global space-time supersymmetry transformation:
$`\delta _ϵA`$ $`=`$ $`{\displaystyle \frac{dz}{2\pi i}ϵ_\alpha S_{(1/2)}^\alpha (z)\mathrm{\Psi }},`$ (C.28)
$`\delta _ϵ\mathrm{\Psi }`$ $`=`$ $`X(i){\displaystyle \frac{dz}{2\pi i}ϵ_\alpha S_{(1/2)}^\alpha (z)A}.`$ (C.29)
#### Modified cubic version
In modified version of cubic open superstring field theory, there are also fermionic gauge symmetry and global supersymmetry. Fermionic gauge symmetry, which is generated by Grassmann even gauge parameter $`\chi `$ in the Ramond sector with picture number $`1/2`$ and ghost number $`0`$, is given by
$`\delta _\chi A`$ $`=`$ $`X(i)(\mathrm{\Psi }\chi \chi \mathrm{\Psi }),`$ (C.30)
$`\delta _\chi \mathrm{\Psi }`$ $`=`$ $`Q_\mathrm{B}\chi +A\chi \chi A,`$ (C.31)
where $`A(\mathrm{\Psi })`$ denotes a Grassmann odd string field in the NS (R) sector with picture number $`0(1/2)`$ and ghost number $`1(1)`$. Formally, by taking
$`\chi `$ $`=`$ $`Y(i){\displaystyle _{C_{\mathrm{left}}}}{\displaystyle \frac{dz}{2\pi i}}Q_\mathrm{B}ϵ_\alpha \xi S_{(1/2)}^\alpha (z)I=Y(i){\displaystyle _{C_{\mathrm{left}}}}{\displaystyle \frac{dz}{2\pi i}}ϵ_\alpha W^\alpha (z)I,`$ (C.33)
$`W^\alpha (z)[Q_\mathrm{B},\xi S_{(1/2)}^\alpha (z)],`$
with the inverse picture changing operator $`Y(z)`$, which is global: $`Q_\mathrm{B}\chi =0`$, the above gauge transformation yields the global space-time supersymmetry transformation :<sup>24</sup><sup>24</sup>24 Similar formula can be found in . See also .
$`\delta _ϵA`$ $`=`$ $`{\displaystyle \frac{dz}{2\pi i}ϵ_\alpha W^\alpha (z)\mathrm{\Psi }},`$ (C.34)
$`\delta _ϵ\mathrm{\Psi }`$ $`=`$ $`Y(i){\displaystyle \frac{dz}{2\pi i}ϵ_\alpha W^\alpha (z)A}.`$ (C.35)
## Appendix D Classical solutions and marginal deformations in Witten’s bosonic open string field theory
In this section, we consider classical solutions of the Witten’s bosonic string field theory corresponding to marginal deformations, which are generalization of the previous ones investigated in <sup>25</sup><sup>25</sup>25 General arguments in bosonic string field theory are given in . and counterparts of the arguments in §3. We discuss a class of solutions using a current $`J^a`$ associated with a Lie algebra $`𝒢`$. We suppose the OPE among currents with adjoint indices of the form
$`J^a(y)J^b(z)`$ $``$ $`g^{ab}{\displaystyle \frac{1}{(yz)^2}}+{\displaystyle \frac{1}{yz}}f_c^{ab}J^c(z),`$ (D.1)
$`g^{ab}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(f_d^{ac}f_c^{bd}\mathrm{\Omega }^{ab}),`$ (D.2)
where $`f_c^{ab}`$ is the structure constant of $`𝒢`$ and $`\mathrm{\Omega }^{ab}`$ is a particular invertible invariant matrix such as (3.4). $`J^a`$ is a primary field with dimension 1 for the energy momentum tensor given by the Sugawara form:
$`T(z)`$ $`=`$ $`\mathrm{\Omega }_{ab}:J^aJ^b:(z),(\mathrm{\Omega }^{ab}\mathrm{\Omega }_{bc}=\delta _c^a).`$ (D.3)
In fact, we can show the OPEs
$`T(y)J^a(z)`$ $``$ $`{\displaystyle \frac{1}{(yz)^2}}J^a(z)+{\displaystyle \frac{1}{yz}}J^a(z),`$ (D.4)
$`T(y)T(z)`$ $``$ $`{\displaystyle \frac{c}{2}}{\displaystyle \frac{1}{(yz)^4}}+{\displaystyle \frac{1}{(yz)^2}}2T(z)+{\displaystyle \frac{1}{yz}}T(z),`$ (D.5)
where the central charge $`c`$ of the Virasoro algebra is given by $`c=\mathrm{dim}𝒢f_d^{ac}f_c^{bd}\mathrm{\Omega }_{ab}`$ . In the following, we assume that the background is described by the above CFT with $`c=26`$, and construct a classical solution of bosonic open string field theory of cubic form:
$`S[\mathrm{\Psi }]`$ $`=`$ $`{\displaystyle \frac{1}{g^2}}\left({\displaystyle \frac{1}{2}}\mathrm{\Psi }Q_\mathrm{B}\mathrm{\Psi }+{\displaystyle \frac{1}{3}}\mathrm{\Psi }\mathrm{\Psi }\mathrm{\Psi }\right)`$ (D.6)
on this background. Namely, the BRS operator in the kinetic term is
$`Q_\mathrm{B}`$ $`=`$ $`{\displaystyle \frac{dz}{2\pi i}\left(cT(z)+bcc(z)\right)}`$ (D.7)
where $`T(z)`$ is given by eq. (D.3) and interaction term is defined by the Witten $``$ product using conformal mappings and CFT correlators . With the above setup, we can show commutation relations
$`\{Q_\mathrm{B},cJ^a(z)\}=0,\{Q_\mathrm{B},c(z)\}=cc(z),`$ (D.8)
$`\{V_L^a(f),V_L^b(g)\}={\displaystyle \frac{1}{2}}g^{ab}\{Q_\mathrm{B},C_L(fg)\},\{V_L^a(f),C_L(g)\}=\{C_L(f),C_L(g)\}=0,`$ (D.9)
where $`V_L^a(f)=_{C_{\mathrm{left}}}\frac{dz}{2\pi i}\frac{1}{\sqrt{2}}f(z)cJ^a(z)`$ and $`C_L(f)=_{C_{\mathrm{left}}}\frac{dz}{2\pi i}f(z)c(z)`$ using similar method in . Then, noting $`cJ^a(z)`$ is a primary field with dimension $`0`$, we have
$`V_L^a(F_a)IV_L^b(F_b)I=V_L^a(F_a)V_L^b(F_b)I={\displaystyle \frac{1}{4}}g^{ab}\{Q_\mathrm{B},C_L(F_aF_b)\}I,`$ (D.10)
with $`F_a(1/z)=z^2F_a(z)`$. Using this relation, we can show that
$`\mathrm{\Psi }_0`$ $`=`$ $`V_L^a(F_a)I{\displaystyle \frac{1}{4}}g^{ab}C_L(F_aF_b)I,F_a(1/z)=z^2F_a(z),`$ (D.11)
satisfies the equation of motion: $`Q_\mathrm{B}\mathrm{\Psi }_0+\mathrm{\Psi }_0\mathrm{\Psi }_0=0`$. If we re-expand the action around this solution such as $`\mathrm{\Psi }=\mathrm{\Psi }_0+\mathrm{\Psi }^{}`$, we have $`S[\mathrm{\Psi }]=S[\mathrm{\Psi }_0]+S^{}[\mathrm{\Psi }^{}]`$ where the new action $`S^{}[\mathrm{\Psi }^{}]`$ is the same form as original one (D.6) except that the new BRS operator is given by $`Q_\mathrm{B}^{}A=Q_\mathrm{B}A+\mathrm{\Psi }_0A(1)^{|A|}A\mathrm{\Psi }_0`$, or more explicitly:
$`Q_\mathrm{B}^{}`$ $`=`$ $`Q_\mathrm{B}V^a(F_a){\displaystyle \frac{1}{4}}g^{ab}C(F_aF_b),`$ (D.12)
($`V^a(f)=\frac{dz}{2\pi i}\frac{1}{\sqrt{2}}f(z)cJ^a(z),C(f)=\frac{dz}{2\pi i}f(z)c(z)`$). Comparing this $`Q_\mathrm{B}^{}`$ with original one $`Q_\mathrm{B}`$ (D.7), we find the Virasoro operator $`T(z)`$ in the matter sector is replaced by $`T(z)\frac{1}{\sqrt{2}}F_a(z)J^a(z)\frac{1}{4}g^{ab}F_a(z)F_b(z)`$. In fact, if we define $`T^{}(z)=_nL_n^{}z^{n2}`$ as
$`L_n^{}`$ $`=`$ $`L_n{\displaystyle \frac{1}{\sqrt{2}}}{\displaystyle \underset{k}{}}F_{a,k}J_{nk}^a{\displaystyle \frac{1}{4}}g^{ab}{\displaystyle \underset{k}{}}F_{a,nk}F_{b,k},`$ (D.13)
with $`F_{a,n}=\frac{d\sigma }{2\pi }e^{i(n+1)\sigma }F_a(e^{i\sigma })`$, we obtain the Virasoro algebra with the same central charge $`c=26`$ as the original one:
$`[L_m^{},L_n^{}]=(mn)L_{m+n}^{}+{\displaystyle \frac{c}{12}}(m^3m)\delta _{m+n,0}.`$ (D.14)
Furthermore, by taking $`J^a(z)=_nJ_n^az^{n1}`$ as<sup>26</sup><sup>26</sup>26In the definition of $`M_{b,n}^a`$, we use path-ordered form in the same way as eq. (3.36). Therefore, we have similar formulae to eqs. (3.38), (3.39) and (3.40) by replacing $`\mathrm{\Omega }^{ab}`$ with $`2g^{ab}`$ in the bosonic case. Note that $`f_c^{ab}g^{cd}+f_c^{ad}g^{cb}=0`$.
$`J_n^a={\displaystyle \underset{k}{}}M_{b,k}^a\left(J_{nk}^b+{\displaystyle \frac{1}{\sqrt{2}}}g^{bc}F_{c,nk}\right),`$ (D.15)
$`{\displaystyle \underset{n}{}}M_{b,n}^ae^{in\sigma }=\left[𝐏\mathrm{exp}\left(i{\displaystyle _0^1}𝑑t\sigma A(t\sigma )\right)\right]_b^a,A_b^a(\sigma )={\displaystyle \frac{1}{\sqrt{2}}}f_b^{ac}e^{i\sigma }F_c(e^{i\sigma }),`$ (D.16)
the same commutation relations as the original one are recovered:
$`[J_m^a,J_n^b]=g^{ab}m\delta _{m+n,0}+f_c^{ab}J_{m+n}^c,[L_m^{},J_n^a]=nJ_{m+n}^a.`$ (D.17)
In the above, we have constructed a classical solution $`\mathrm{\Psi }_0`$ (D.11) and re-expanded around it. The obtained action $`S^{}[\mathrm{\Psi }^{}]`$ is also reproduced by replacing $`T(z)`$ with $`T^{}(z)`$ (D.13) in the original action $`S[\mathrm{\Psi }]`$ (D.6). This replacement is induced by the map $`J^aJ^a`$ (D.15) in terms of the current which preserves the algebra among $`(J^a(z),T(z))`$. At least formally, we can show that the vacuum energy vanishes at this solution :
$`S[\mathrm{\Psi }_0]={\displaystyle _0^1}𝑑t{\displaystyle \frac{d}{dt}}S[\mathrm{\Psi }_0(t)]={\displaystyle \frac{1}{g^2}}{\displaystyle _0^1}𝑑t{\displaystyle \frac{d}{dt}}\mathrm{\Psi }_0(t)(Q_\mathrm{B}\mathrm{\Psi }_0(t)+\mathrm{\Psi }_0(t)\mathrm{\Psi }_0(t))=0`$ (D.18)
where $`\mathrm{\Psi }_0(t)`$ is given by replacing $`F_a(z)`$ with $`tF_a(z)`$ in $`\mathrm{\Psi }_0`$ (D.11). These facts suggest that $`\mathrm{\Psi }_0`$ (D.11) may be a pure gauge solution. In fact, noting $`[Q_\mathrm{B},J^a(z)]=(cJ^a)(z)`$, we make an ansatz for the gauge parameter as
$`\mathrm{\Lambda }_LI={\displaystyle _{C_{\mathrm{left}}}}{\displaystyle \frac{dz}{2\pi i}}{\displaystyle \frac{1}{\sqrt{2}}}g_a(z)J^a(z)I;g_a(1/z)=g_a(z),g_a(\pm i)=0.`$ (D.19)
Using the OPE (D.1), we can compute its pure gauge form as
$`e^{\mathrm{\Lambda }_LI}Q_\mathrm{B}e^{\mathrm{\Lambda }_LI}`$
$`=`$ $`V_L^b(g_a((e^{}1)^1)_b^a)I{\displaystyle \frac{1}{4}}C(g_ag_b((e^{}1)^1)_c^ag^{cd}((e^{}1)^1)_d^b)I,`$
where we denoted $`_b^a(z)=\frac{1}{\sqrt{2}}f_b^{ac}g_c(z)`$. Comparing with the solution (D.11), we can get a gauge parameter by solving
$`F_b(z)`$ $`=`$ $`g_a(z)((e^{}1)^1)_b^a(z),`$ (D.21)
with respect to $`g_a(z)`$ for a given $`F_a(z)`$. For a solution $`g_a(z)`$, the new BRS operator $`Q_\mathrm{B}^{}`$ (D.12) can be rewritten as a similarity transformation from the original one:
$`Q_\mathrm{B}^{}=e^\mathrm{\Lambda }Q_\mathrm{B}e^\mathrm{\Lambda },\mathrm{\Lambda }={\displaystyle \frac{dz}{2\pi i}\frac{1}{\sqrt{2}}g_a(z)J^a(z)}.`$ (D.22)
In this case, we recover the original SFT action by performing a field redefinition such as $`\mathrm{\Psi }^{\prime \prime }=e^\mathrm{\Lambda }\mathrm{\Psi }^{}`$, which implies the effect of a marginal deformation by a current $`J^a`$. Its deformation parameter $`g_a(z)`$ is related to the classical solution of string field theory as (D.21). However, we cannot always obtain a solution $`g_a(z)`$ to the above differential equation (D.21) because of the boundary condition $`g_a(\pm i)=0`$, which is imposed by a partial integration in computing eq. (D). This situation is just the same as supersymmetric case in §3. Namely, consistency condition for the pure gauge form (D) is given by (3.60) because the differential equation (D.21) is the same as (3.54). Therefore, there is a possibility that a solution $`\mathrm{\Psi }_0`$ becomes nontrivial if $`F_a(z)`$ does not satisfy eq. (3.60).
As an example of $`𝒢=u(1)^{26}`$, we take the current $`J^\mu =\frac{i}{\sqrt{2\alpha ^{}}}X^\mu `$ on the flat background. In this case, we can identify as
$`T(z)={\displaystyle \frac{1}{4\alpha ^{}}}X^\mu X_\mu (z),f_c^{ab}=0,\mathrm{\Omega }^{\mu \nu }=2\eta ^{\mu \nu },c=\mathrm{dim}(u(1)^{26})=26,`$ (D.23)
using the above notation. If some directions are compactified to the torus, a solution $`\mathrm{\Psi }_0`$ (D.11) becomes nontrivial according to a nontrivial Wilson lines: $`_{C_{\mathrm{left}}}𝑑zF_\mu (z)0`$ . When one direction ($`X^{25}`$) is $`S^1`$-compactified at the critical radius $`R=\sqrt{\alpha ^{}}`$, we can regard the algebra as $`𝒢=u(1)^{25}\times su(2)`$ and identify as
$`J^1(z)=\sqrt{2}\mathrm{cos}({\displaystyle \frac{X^{25}}{\sqrt{\alpha ^{}}}})(z),J^2(z)=\sqrt{2}\mathrm{sin}({\displaystyle \frac{X^{25}}{\sqrt{\alpha ^{}}}})(z),J^3(z)={\displaystyle \frac{i}{\sqrt{2\alpha ^{}}}}X^{25}(z),`$ (D.24)
$`T(z)={\displaystyle \frac{1}{4\alpha ^{}}}X^{25}X^{25}(z)={\displaystyle \frac{1}{6}}:(J^1J^1+J^2J^2+J^3J^3):(z),`$ (D.25)
$`f_c^{ab}=\sqrt{2}iϵ_{abc},\mathrm{\Omega }^{ab}=6\delta ^{ab},c=\mathrm{dim}(su(2))34/6=1,`$ (D.26)
in the $`su(2)`$ sector. The corresponding solution $`\mathrm{\Psi }_0`$ was investigated in refs. . |
warning/0506/cond-mat0506370.html | ar5iv | text | # Isotropic-nematic interfacial tension of hard and soft rods: application of advanced grand canonical biased sampling techniques
## I Introduction
The aim of this paper is to present a computation of the interfacial tension $`\gamma _{\mathrm{IN}}`$ between the coexisting isotropic and nematic (IN) phase in suspensions of monodisperse hard rods via computer simulation. While the hard-rod fluid simplifies experimental reality, ignoring for example long-ranged interactions and polydispersity Chen and Gray (2002), it nevertheless captures the main mechanism of the IN phase transition and serves as a valuable model system. Experiments have shown that $`\gamma _{\mathrm{IN}}`$ is very small, typically in the range 10<sup>-3</sup>–10<sup>-4</sup> mN/m Chen and Gray (2002), which makes it difficult to extract $`\gamma _{\mathrm{IN}}`$ from simulation data. Simulation estimates of $`\gamma _{\mathrm{IN}}`$ are therefore rare, and have only been reported for ellipsoids McDonald et al. (2000); Akino et al. (2001), soft rods Vink and Schilling (2005), and lattice models Cleaver and Allen (1993). Theoretical estimates are more abundant Chen and Noolandi (1992); Koch and Harlen (1999); van der Schoot (1999); Velasco et al. (2002); Allen (2000a); Shundyak and van Roij (2001), but are usually obtained in the Onsager limit Onsager (1949) of infinite rod length ($`L/D\mathrm{}`$). The case of finite rod length is more difficult to describe theoretically, but has been addressed in Ref. Velasco et al., 2002 using density functional theory, and in Ref. van der Schoot, 1999 using a scaling relation. At the time of writing, no simulation estimate of $`\gamma _{\mathrm{IN}}`$ for the hard-rod fluid has been reported. Such an estimate would clearly be valuable to test theoretical predictions, and to see if the order of magnitude of $`\gamma _{\mathrm{IN}}`$ observed in experiments is reproduced.
Despite its simplicity, simulating the hard-rod fluid is not trivial Bolhuis and Frenkel (1997); Dijkstra et al. (2001). The bottleneck is the hard particle interaction, which complicates both molecular dynamics (MD) and Monte Carlo (MC) methods. In the case of MD, the discontinuous potential prevents the calculation of smooth forces. In the case of MC, equilibration times are long due to very low acceptance rates. An important improvement is the use of soft interactions, as was done for ellipsoids McDonald et al. (2000); Akino et al. (2001) and rods Al-Barwani and Allen (2000); Vink and Schilling (2005). By using soft interactions, the qualitative phase behavior is usually retained, but simulations become much more efficient. Moreover, MC simulations in the grand canonical ensemble become possible, enabling the investigation of IN coexistence via the probability distribution in the particle number density. This technique is well established in simulations of fluid-vapor coexistence Potoff and Panagiotopoulos (2000); Góźdź (2003); Virnau et al. (2004a); Müller and MacDowell (2000); Vink and Horbach (2004), and was recently extended to IN coexistence in suspensions of soft rods Vink and Schilling (2005). The advantage of grand canonical simulations is that the coexistence densities, as well as the interfacial tension, can be obtained.
Since coexisting phases are separated by a free energy barrier arising from the interfacial tension Wilding (2001), it is essential to use a biased sampling scheme to access regions of high free energy. In simulations of fluid-vapor coexistence, the bias is usually put on the density. While a density bias has also been used to simulate IN coexistence Vink and Schilling (2005), this choice is not optimal. In simulations that rely on standard MC moves, such as random translations and rotations of single particles, it is difficult to reach the nematic phase starting in the isotropic phase simply by increasing the density because the orientational degrees of freedom relax only very slowly Williams and Philipse (2003). This effect is called “jamming”, and it explains why the simulations of Ref. Vink and Schilling, 2005 were limited to rather small systems.
In this work, grand canonical MC simulations using a bias on the nematic order parameter are performed. As we will show, this approach is much less susceptible to jamming, and enables simulations of large systems. This in turn allows for accurate estimates of the interfacial tension in suspensions of soft rods. As an additional bonus, a bias on the nematic order parameter paves the way towards grand canonical simulations of hard rods, enabling a simulation estimate of $`\gamma _{\mathrm{IN}}`$ for the hard-rod fluid.
The outline of this paper is as follows: First, we introduce the liquid crystal model used in this work. The biased sampling scheme is described next. The results are presented in Sec. IV. We end with a summary and comparison to theoretical predictions in the last section.
## II Model and order parameters
We consider rods of elongation $`L`$ and diameter $`D`$. The simulations are performed in a three dimensional box of size $`L_\mathrm{x}\times L_\mathrm{y}\times L_\mathrm{z}`$ using periodic boundary conditions in all dimensions. In this work, we fix $`L_\mathrm{x}=L_\mathrm{y}`$, but we allow for elongation in the remaining dimension $`L_\mathrm{z}L_\mathrm{x}`$. Moreover, to avoid double interactions between rods through the periodic boundaries, we set $`L_\mathrm{x}>2L`$. The position of the center of mass of rod $`i`$ is denoted $`\stackrel{}{r}_i`$, and its orientation $`\stackrel{}{u}_i`$, with normalization $`|\stackrel{}{u}_i|=1`$. The interaction between two rods $`i`$ and $`j`$ is given by a pair potential of the form
$`v_{ij}(r)`$ $`=`$ $`\{\begin{array}{cc}ϵ\hfill & r<D,\hfill \\ 0\hfill & \mathrm{otherwise},\hfill \end{array}`$ (1)
with $`r`$ the distance between two line segments of length $`L`$, see also Ref. Vink and Schilling, 2005. The total energy is thus a function of the center of mass coordinates and the orientations of all rods
$$E(\stackrel{}{r}_1,\mathrm{},\stackrel{}{r}_N;\stackrel{}{u}_1,\mathrm{},\stackrel{}{u}_N)=\underset{i=1}{\overset{N}{}}\underset{j=i+1}{\overset{N}{}}v_{ij},$$
(2)
with $`N`$ the number of rods in the system (in the following we will drop the $`\stackrel{}{r}_i`$ and $`\stackrel{}{u}_i`$ dependences in our notation).
To investigate the IN transition, the density and the average rod alignment are used as order parameters. Since the density in the isotropic phase is lower than in the nematic phase, the rod number density $`\rho =N/V`$ may be used to distinguish the phases, with $`V`$ the volume of the simulation box. Following convention, we also introduce the reduced density $`\rho ^{}=\rho /\rho _{\mathrm{cp}}`$, with $`\rho _{\mathrm{cp}}=2/[\sqrt{2}+(L/D)\sqrt{3}]`$ the density of regular close packing of hard rods Bolhuis and Frenkel (1997).
In the nematic phase the rods are on average aligned, whereas in the isotropic phase the rods are randomly oriented. Therefore, the nematic order parameter may also be used to distinguish the phases. The latter quantity is defined in terms of the orientational tensor $`Q`$, whose components are given by
$$Q_{\alpha \beta }=\frac{1}{2N}\underset{i=1}{\overset{N}{}}\left(3u_{i\alpha }u_{i\beta }\delta _{\alpha \beta }\right),$$
(3)
with $`u_{i\alpha }`$ the $`\alpha `$ component ($`\alpha =x,y,z`$) of the orientation of rod $`i`$, and $`\delta _{\alpha \beta }`$ the Kronecker delta. In this work, the maximum eigenvalue $`S`$ of the orientational tensor is taken as nematic order parameter, being close to unity in the nematic phase, and close to zero in the isotropic phase. The eigenvector corresponding to $`S`$ is called the director, and it measures the preferred direction of the rods in the nematic phase.
## III Simulation method
We study IN coexistence via grand canonical MC simulations. In the grand canonical ensemble, the volume, the temperature $`T`$, and the chemical potential $`\mu `$ are fixed, while the number of rods in the simulation box fluctuates. Insertion and removal of rods are attempted with equal probability and accepted with appropriate Metropolis rules to be given later. The aim of grand canonical simulations is to measure the probability distribution in the number of particles $`P(N)`$. At the coexistence chemical potential, $`P(N)`$ becomes bimodal with two peaks of equal area. An example distribution is shown in Fig. 4, where we have plotted the logarithm of $`P(N)`$. The peak locations yield the coexistence densities; the average height of the peaks $`\mathrm{\Delta }\mathrm{\Omega }`$ in $`k_\mathrm{B}T\mathrm{ln}P(N)`$ is the free energy barrier separating the phases, with $`k_\mathrm{B}`$ the Boltzmann constant. In three dimensions using periodic boundary conditions and for sufficiently large systems, the barrier is related to the interfacial tension via $`\gamma _{\mathrm{IN}}=\mathrm{\Delta }\mathrm{\Omega }/(2L_\mathrm{x}^2)`$, where $`L_\mathrm{x}`$ is the lateral dimension of the simulation box Binder (1982).
In simulations, the free energy barrier presents a problem. Unless $`\mathrm{\Delta }\mathrm{\Omega }`$ is small, such as close to a critical point, simulations rarely cross the barrier, and spend most time in only one of the two phases. Biased sampling techniques are required to overcome the barrier. In general, these techniques aim to construct a weight function $`W(\zeta )`$ of some bias variable $`\zeta `$. The weight function is constructed such that a simulation using a modified potential $`E^{}(\zeta )=E+k_\mathrm{B}TW(\zeta )`$ yields a uniform probability distribution in the bias variable, with $`E`$ the potential of the original system. The grand canonical acceptance rules using the modified potential read as
$`A(N,\zeta _0N+1,\zeta _1)=`$
$`\mathrm{min}[1,{\displaystyle \frac{V}{N+1}}e^{\beta (\mathrm{\Delta }E\mu )W(\zeta _1)+W(\zeta _0)}],`$ (4)
$`A(N,\zeta _0N1,\zeta _1)=`$
$`\mathrm{min}[1,{\displaystyle \frac{N}{V}}e^{\beta (\mathrm{\Delta }E+\mu )+W(\zeta _1)W(\zeta _0)}],`$ (5)
for the insertion and removal of a single particle, respectively Landau and Binder (2000); Frenkel and Smit (2001). Here, $`\zeta _0`$ and $`\zeta _1`$ denote, respectively, the value of the bias parameter in the initial and final state, $`\mathrm{\Delta }E`$ is the potential energy difference between initial and final state given by Eq.(2), and $`\beta =1/(k_\mathrm{B}T)`$. For a properly constructed $`W(\zeta )`$, the biased simulation samples all states $`\zeta `$ with uniform probability. Once $`W(\zeta )`$ is known, the distribution $`P(N)`$ of the unbiased system can be constructed.
One is rather free in choosing the bias variable. The best choices are variables that change significantly when going from one phase to the other. For fluid-vapor transitions, a natural bias is the particle number density. In the case of IN coexistence, the density is still a valid variable because of the density gap between the isotropic and the nematic phase. This was used in Ref. Vink and Schilling, 2005 to study IN coexistence in suspensions of soft rods. Whether a bias on density in systems of elongated particles is efficient, depends on how easily a dense isotropic phase can rearrange itself to become nematic. In practice, the jamming effect limits density biased sampling to rather small systems and soft interactions. As it turns out, for IN transitions, a much more powerful bias variable is the nematic order parameter $`S`$. Note, however, that phase coexistence is defined in terms of $`P(N)`$. Therefore, in a simulation which biases on $`S`$, the distribution $`P(N)`$ must still be reconstructed. To this end, histograms in both the particle number $`N`$, as well as in $`S`$, have to be measured. In this section, we explain how the bias on $`S`$, and the subsequent reweighting in $`N`$ and $`S`$, are implemented. It is convenient to discuss the more straightforward procedure of a density bias first.
### III.1 Biased sampling on $`\rho `$
A convenient method to bias on density is Successive Umbrella Sampling (SUS) Virnau and Müller (2004). Here we describe the algorithm in its simplest form; refinements are given in the original reference. The choice of the sampling algorithm is not crucial. The general principles also apply to other schemes, such as conventional umbrella sampling Torrie and Valleau (1977), multicanonical sampling Berg and Neuhaus (1992), Wang-Landau sampling Wang and Landau (2001), or hyperparallel tempering Yan and de Pablo (2000).
The aim is to construct a function $`W(N)`$ of the number of particles such that a simulation using the modified potential $`E^{}(N)=E+k_\mathrm{B}TW(N)`$ yields a uniform distribution in $`N`$, with $`E`$ given by Eq.(2). The modified potential thus contains an explicit dependence on the bias variable $`N`$. Following Ref. Virnau and Müller, 2004, the particle number axis is divided into equally sized intervals called windows, starting with some minimum number of particles $`N_0`$. In the first window, the number of particles is confined to $`N_0NN_0+1`$, in the second window to $`N_0+1NN_0+2`$, and in the $`i`$-th window to $`N_0+i1NN_0+i`$. In this example, the window size equals a single particle but this choice is not essential: SUS works just as well using larger windows Virnau and Müller (2004). The choice of the window size is not completely arbitrary. Choosing the windows too large leads to poor sampling statistics at the window boundaries; choosing the windows too small runs the risk that certain relaxation pathways are cut-off. In practice, a compromise needs to be made.
The idea of SUS is to construct $`W(N)`$ by simulating the windows separately and successively. Starting in the first window ($`i=1`$), grand canonical MC moves are attempted (optionally combined with canonical moves such as translations and rotations), with the constraint that states outside the window bounds are rejected to fulfill detailed balance at the window boundaries. The relevant weights in the first window are $`W(N_0)`$ and $`W(N_0+1)`$, which we initially set to zero. We then record $`f_\mathrm{L}^1`$ and $`f_\mathrm{H}^1`$, counting the occurrence of the state with $`N_0`$ and $`N_0+1`$ particles, respectively. In this notation, the subscripts “L” and “H” refer to the “lower” and “higher” window bound, respectively, while the superscript refers to the window number. To obtain a uniform distribution in $`N`$, the ratio of the counts should be unity. This will generally not be the case, but is enforced by updating the weight of the higher window bound to
$$W_{\mathrm{new}}(N_0+i)=W_{\mathrm{old}}(N_0+i)+\mathrm{ln}(f_\mathrm{H}^i/f_\mathrm{L}^i),$$
(6)
leaving the weight of the lower bound $`W(N_0+i1)`$ unchanged, where $`i`$ is the window number. In case $`f_\mathrm{H}^i>f_\mathrm{L}^i`$, the effect of this modification is a lower insertion rate, see Eq.(III), and a higher removal rate, see Eq.(III), leading to a count ratio closer to one. The latter can be checked by performing a second simulation using the updated weight (the reasoning for $`f_\mathrm{H}^i<f_\mathrm{L}^i`$ is similar). In practice, it may occur that one of the counts is zero. It is then necessary to modify $`W(N_0+i)`$ by hand first, before starting the simulation. Note also that long simulation runs may be required to obtain the count ratio accurately.
Having simulated the first window, $`W(N_0)`$ and $`W(N_0+1)`$ are known. The choice $`W(N_0)=0`$ is arbitrary but has no physical consequences since it merely shifts the potential by a constant. Next, we consider window 2, where the number of particles is allowed to fluctuate between $`N_0+1`$ and $`N_0+2`$, with respective weights $`W(N_0+1)`$ and $`W(N_0+2)`$. An important optimization of Ref. Virnau and Müller, 2004 is to linearly extrapolate the known weights $`W(N_0)`$ and $`W(N_0+1)`$ to obtain an estimate for $`W(N_0+2)`$ (note that for the third and subsequent windows, quadratic extrapolation can be used). The simulations in the second window are then performed using the extrapolated estimate, and the respective counts, $`f_\mathrm{L}^2`$ and $`f_\mathrm{H}^2`$, of visiting the state with $`N_0+1`$ and $`N_0+2`$ particles, are recorded. Finally, the weight $`W(N_0+2)`$ is updated using Eq.(6), leaving the other weight $`W(N_0+1)`$ unchanged, and the next window is considered.
The above procedure is repeated until all windows of interest have been simulated, and the corresponding weight function $`W(N)`$ is constructed. The sought-for distribution in the number of particles $`P(N)`$ is trivially obtained via $`P(N)=Ce^{W(N)}`$, with $`C`$ a normalization constant Wilding (2001).
### III.2 Biased sampling on $`S`$
Next, we consider the extension to a bias on the nematic order parameter $`S`$. Here, the modified potential reads as $`E^{}(S)=E+k_\mathrm{B}TW(S)`$, with $`E`$ given by Eq.(2). The aim is to construct $`W(S)`$ such that a simulation using the modified potential samples all values of $`S`$ with uniform probability.
The windows are obtained by dividing the nematic order parameter axis into equally sized intervals of width $`\mathrm{\Delta }S`$. In the first window, the nematic order parameter is confined to $`0S<\mathrm{\Delta }S`$, in the second window to $`\mathrm{\Delta }S/2S<3\mathrm{\Delta }S/2`$, and in the $`i`$-th window to $`(i1)\mathrm{\Delta }S/2S<(i+1)\mathrm{\Delta }S/2`$, see Fig. 1. The windows thus partially overlap. To sample both the isotropic and the nematic phase, the sampling range should span from $`S=0`$ to $`S1`$. Note that $`S`$ is a continuous variable, whereas the density (expressed in the number of particles) is discrete. Therefore, a natural width for the windows does not exist, and one is forced to choose $`\mathrm{\Delta }S`$ rather arbitrarily. We found that $`\mathrm{\Delta }S0.0010.002`$ gives good results, which means that $`𝒪(10^3)`$ windows are required to sample the transition. A consequence of discretizing the nematic order parameter is that $`W(S)`$ is defined in steps of $`\mathrm{\Delta }S/2`$. Therefore, in the $`i`$-th window, $`W(S)`$ assumes only two distinct values
$`W(S)`$ $`=`$ $`\{\begin{array}{cc}W_{i1}\hfill & (i1)\mathrm{\Delta }S/2S<S_\mathrm{M}\hfill \\ W_i\hfill & S_\mathrm{M}S<(i+1)\mathrm{\Delta }S/2,\hfill \end{array}`$
with $`S_\mathrm{M}=i\mathrm{\Delta }S/2`$ the center of the window (note that $`i>0`$).
Starting in the first window ($`i=1`$), the relevant weights are $`W_0`$ and $`W_1`$, which are initially set to zero. While simulating the first window, we count the occurrence of states with $`0S<S_\mathrm{M}`$ ($`f_\mathrm{L}^1`$) and $`S_\mathrm{M}S<\mathrm{\Delta }S`$ ($`f_\mathrm{H}^1`$), with $`S_\mathrm{M}=\mathrm{\Delta }S/2`$. To obtain the distribution in the number of particles $`P(N)`$ (after all the quantity of interest) particle number histograms must also be stored (note that $`N`$ fluctuates freely in each window). In the first window, we thus record the probability distribution in the number of particles $`p_\mathrm{L}^1(N)`$ for states with $`0S<S_\mathrm{M}`$, and $`p_\mathrm{H}^1(N)`$ for states with $`S_\mathrm{M}S<\mathrm{\Delta }S`$. It is recommended to store the distributions unnormalized. This makes it more convenient to restart the simulations at a later stage in case higher precision is required. After simulating the first window, the weight of the higher window bound is updated to force a uniform distribution in $`S`$ using $`W_{i,\mathrm{new}}=W_{i,\mathrm{old}}+\mathrm{ln}(f_\mathrm{H}^i/f_\mathrm{L}^i)`$, while keeping $`W_{i1}`$ fixed, where $`i`$ is the window number.
In the second window ($`i=2`$), the relevant weights are $`W_1`$ and $`W_2`$. To simulate efficiently, the weights $`W_0`$ and $`W_1`$ of the previous window are extrapolated to estimate $`W_2`$. The extrapolated estimate is used while simulating the second window, and the counts $`f_\mathrm{L}^2`$ and $`f_\mathrm{H}^2`$ are recorded, as well as the distributions $`p_\mathrm{L}^2(N)`$ and $`p_\mathrm{H}^2(N)`$. After simulating the second window, $`W_2`$ is updated as before, and the next window is considered.
The above procedure is repeated up to some maximum number of windows $`w_{\mathrm{max}}`$ chosen well into the nematic phase. The remaining step is to combine the weights $`W_i`$ with the distributions $`p_\mathrm{L}^i(N)`$ and $`p_\mathrm{H}^i(N)`$ to obtain $`P(N)`$. Note that the upper region of window $`i`$ overlaps with the lower region of the next window $`i+1`$, see Fig. 1. More precisely, the distributions $`p_\mathrm{H}^i(N)`$ and $`p_\mathrm{L}^{i+1}(N)`$ stem from the same $`S`$ interval and are thus measured with the same probability by the sampling scheme. Therefore, these distributions may be combined $`\overline{p}_i(N)=p_\mathrm{H}^i(N)+p_\mathrm{L}^{i+1}(N)`$, and normalized such that $`_{N=0}^{\mathrm{}}\overline{p}_i(N)=1`$. The distribution $`P(N)`$ is simply a weighted sum of the above (normalized) $`\overline{p}_i(N)`$. Since $`k_\mathrm{B}TW_i`$ corresponds to a free energy, each $`\overline{p}_i(N)`$ contributes to $`P(N)`$ with a weight proportional to $`e^{W_i}`$. This leads to $`P(N)=C_{i=1}^{w_{\mathrm{max}}}\overline{p}_i(N)e^{W_i}`$, where the sum is over all windows, and normalization constant $`C^1=_{N=0}^{\mathrm{}}_{i=1}^{w_{\mathrm{max}}}\overline{p}_i(N)e^{W_i}`$.
### III.3 Bias on $`\rho `$ versus bias on $`S`$
Clearly, the discussed methods serve the same purpose: to measure the distribution $`P(N)`$ at coexistence. Density biased sampling is by far the easiest to implement. It has the additional advantage that the coexistence chemical potential need not be specified beforehand: once $`P(N)`$ has been measured at some chemical potential $`\mu _0`$, it can be extrapolated to any other chemical potential $`\mu _1`$ by using the equation
$$P(N|\mu _1)=P(N|\mu _0)e^{\beta (\mu _1\mu _0)N},$$
(7)
with $`P(N|\mu _\alpha )`$ the probability distribution $`P(N)`$ at chemical potential $`\mu _\alpha `$. Obviously, one should establish roughly beforehand the density at which the IN transition occurs, to avoid sampling large regions of irrelevant phase-space.
The situation is reversed when biasing on the nematic order parameter. In this case, the sampling range is always from $`S=0`$ to $`S1`$. However, to observe phase coexistence, it is essential to use a chemical potential that is rather close to the coexistence value. Of course, Eq.(7) still holds, but the range in $`\mu `$ over which one can extrapolate is much smaller, precisely because the bias is put on $`S`$ and not on $`\rho `$. An estimate of the coexistence chemical potential may be obtained in a density biased simulation of a small system, or via the Widom insertion algorithm Widom (1963); Frenkel and Smit (2001). This certainly makes biasing on $`S`$ more involved. Moreover, for each attempted MC move, $`S`$ in the final state must be determined, regardless of whether the move is accepted. It is therefore important to calculate $`S`$ efficiently. In particular, the $`𝒪(N)`$ loop of Eq.(3) should be eliminated, which can be done following the method outlined in Ref. Vink and Schilling, 2005.
## IV Results
An important conclusion of Ref. Vink and Schilling, 2005 is that the IN interfacial tension obtained from $`P(N)`$ may be prone to strong finite size effects. Away from any critical point, interfaces are the dominant source of finite size effects. The use of periodic boundary conditions leads to the formation of two interfaces. In small systems, the interfaces may interact and this will influence the estimate of $`\gamma _{\mathrm{IN}}`$. A convenient way to suppress interface interactions is to use an elongated simulation box with $`L_\mathrm{z}L_\mathrm{x}`$ Grossmann and Laursen (1993), in accord with Fig. 2. This forces an orientation of the interfaces perpendicular to the elongated dimension (since this minimizes the interfacial area), with a separation between the interfaces that is larger than it would be in a cubic system of the same volume. The absence of interface interactions is manifested by a pronounced flat region between the peaks in $`\mathrm{ln}P(N)`$. Note that a flat region is essential, but not sufficient, to extract $`\gamma _{\mathrm{IN}}`$ reliably. There may still be finite size effects in the lateral dimensions $`L_\mathrm{x}`$ and $`L_\mathrm{y}`$, arising for instance from capillary waves. Ideally, the lateral dimensions should be large enough to capture the long wavelength limiting form of the capillary spectrum Rowlinson and Widom (1982). To actually measure the capillary spectrum of the IN interface is demanding Akino et al. (2001). A more convenient approach sufficient for our purposes is to first establish a minimum elongation $`L_\mathrm{z}`$ in which interface interactions are suppressed, and to then check for finite size effects in the lateral dimensions by varying $`L_\mathrm{x}`$ and $`L_\mathrm{y}`$ explicitly.
An additional motivation to use large lateral dimensions is to stabilize the interfaces. The interfacial free energy is of order $`\gamma _{\mathrm{IN}}L_\mathrm{x}^2`$, and if this is small compared to $`k_\mathrm{B}T`$, the interfaces will generally not be stable. These issues are especially relevant for IN coexistence because $`\gamma _{\mathrm{IN}}`$ is very small. Therefore, in this section, we first perform MC simulations in the canonical ensemble to obtain an indication of the system size required to observe stable interfaces. Next, we present coexistence data obtained using the nematic order biased sampling scheme.
### IV.1 Interfacial profiles
We consider hard rods, i.e. $`ϵ\mathrm{}`$ in Eq.(1), of elongation $`L/D=15`$. The simulations are performed in the canonical ensemble, where the number of rods, the volume, and the temperature are fixed. The box dimensions are $`L_\mathrm{x}=L_\mathrm{y}=10L/3`$ and $`L_\mathrm{z}=20L`$. We set the overall density of the system to $`\rho ^{}=0.205`$, which is well inside the coexistence region Bolhuis and Frenkel (1997); Dijkstra et al. (2001), corresponding to ca. $`11,000`$ particles. An initial system is prepared containing two interfaces, with the director of the nematic phase aligned in the plane of the interface. This is the stable configuration, as confirmed by theory Chen and Noolandi (1992); Koch and Harlen (1999) and simulation Al-Barwani and Allen (2000); McDonald et al. (2000); Vink and Schilling (2005). The initial system is evolved with random rotations and translations of single rods, accepted with the standard Metropolis rules Newman and Barkema (1999); Landau and Binder (2000); Frenkel and Smit (2001). The system is equilibrated for $`10^6`$ sweeps, after which a snapshot is taken every 260 sweeps, up to a total of $`3\times 10^4`$ snapshots (one sweep corresponds to one attempted MC move per rod).
After equilibration, simulation snapshots schematically resemble Fig. 2. Note, however, that they contain far more particles than depicted in this simple sketch. The aim is to measure the density profile $`\rho (z)`$, and the nematic order parameter profile $`S(z)`$ along the elongated $`z`$-dimension, averaged over many different snapshots. The averages are taken with the center of mass of the snapshots shifted to the middle of the simulation box, with the constraint that the nematic phase is also located in the middle, in accord with Fig. 2. The constraint is necessary to remove ambiguity arising from cases where the isotropic phase is in the middle. Having shifted the center of mass, the density profile is obtained by binning the $`z`$-axis in steps of $`\mathrm{\Delta }z0.17L`$. The local density $`\rho (z)`$ in a single snapshot is given by $`n/v_\mathrm{B}`$, with $`n`$ the number of rods in the bin centered around $`z`$, and $`v_\mathrm{B}`$ the volume of a single bin. The density profiles are then averaged over all snapshots. Following Ref. Dijkstra et al., 2001, for the bin centered around $`z`$ in a single snapshot, we also define a local orientational tensor $`Q(z)`$, calculated using Eq.(3) considering only the rods inside the bin. The local orientational tensor elements are then averaged over all snapshots and $`S(z)=\mathrm{maxev}Q(z)`$.
The averaged profiles are shown in Fig. 3. The solid curves are hyperbolic tangent fits of the form $`A+B\mathrm{tanh}\left(\frac{zz_c}{w}\right)`$, which describe the data well. Note that the profiles are shifted with respect to each other. The magnitude of the shift, measured between the inflection points, equals $`\delta =0.37\pm 0.04L`$. This is consistent with theoretical predictions $`\delta =0.450.5L`$ Chen and Noolandi (1992); Shundyak and van Roij (2003), as well as $`\delta 0.33L`$ obtained in simulations of ellipsoids Allen (2000b). Note that the simulated profiles are broadened due to capillary waves Akino et al. (2001). Moreover, we observed considerable fluctuations in the amount of isotropic and nematic phase during the simulation, leading to large fluctuations in the interface positions along the elongated $`L_\mathrm{z}`$ dimension. The width of the averaged profile obtained by fixing the center of mass is therefore additionally broadened Binder and Müller (2000); Tepper and Briels (2002). Because of these effects, we cannot compare the interfacial width of the simulated profiles to theoretical predictions. More important for our purposes, however, is the observation that the interfaces are stable. For hard rods, the current system size thus seems sufficient to accommodate stable interfaces.
### IV.2 Comparison of $`\rho `$ and $`S`$ biased sampling
Having established the typical system size required to observe stable interfaces, biased sampling on the nematic order parameter is considered next. First, we show that density and nematic order biased sampling yield the same distribution $`P(N)`$. To this end, we consider a small system of soft rods with $`L/D=15`$ and $`\beta ϵ=2`$, in a simulation box of size $`L_\mathrm{x}=L_\mathrm{y}=2.1L`$ and $`L_\mathrm{z}=8.4L`$. The latter system was investigated in previous work using density biased sampling Vink and Schilling (2005). The corresponding coexistence chemical potential reads as $`\beta \mu 5.15`$. The nematic order biased sampling scheme is applied to the same system using the latter chemical potential and $`\mathrm{\Delta }S=0.002`$, see Fig. 4. Shown is the coexistence distribution $`P(N)`$ obtained using a bias on $`S`$ (solid curve), as well as using a bias on density (dashed curve, reproduced from Ref. Vink and Schilling, 2005). The agreement between both methods is strikingly confirmed, thereby justifying the approach of Sec. III.2. For small systems, the required CPU time is roughly equal for both methods. The data sets of Fig. 4 required ca. 700 CPU hours each, on 2.2 GHz Pentium machines.
### IV.3 Interfacial tension of soft rods
Next, we consider soft rods with $`L/D=10`$ and $`\beta ϵ=2`$. We aim to accurately measure the interfacial tension. To this end, large system sizes are required such that a bias on $`S`$ is essential. As explained before, the elongated $`L_\mathrm{z}`$ dimension of the simulation box must be large enough to accommodate non-interacting interfaces. At the same time, $`L_\mathrm{x}`$ and $`L_\mathrm{y}`$ must be large enough to suppress finite size effects in the lateral dimensions. We therefore consider two system sizes: $`L_\mathrm{x}=L_\mathrm{y}=3.5L;L_\mathrm{z}=10.5L`$ (system A), and $`L_\mathrm{x}=L_\mathrm{y}=4L;L_\mathrm{z}=14L`$ (system B), where the lateral dimensions are deliberately chosen to exceed those of Sec. IV.1. The simulations are performed using $`\mathrm{\Delta }S=0.001`$ and $`0.002`$, for system A and B, respectively. An initial estimate of the coexistence chemical potential was taken from previous work Vink and Schilling (2005).
In Fig. 5, the dependence of the nematic order parameter on the number of particle is shown, calculated using
$$S(N)=C\underset{i=1}{\overset{w_{\mathrm{max}}}{}}S_i\overline{p}_i(N)e^{W_i},$$
(8)
with $`S_i=i\mathrm{\Delta }S/2\mathrm{\Delta }S/4`$, and the remaining symbols defined as before. Analogous to fluid-vapor transitions Virnau et al. (2004a, b); MacDowell et al. (2004), five distinct regions can be distinguished. In region I, a single isotropic phase is observed. Region II corresponds to the transition from the bulk isotropic phase, to the phase with two parallel interfaces. The transition is characterized by the formation of a nematic droplet in an isotropic background, which grows with the density until it self-interacts through the periodic boundaries, ultimately leading to two parallel interfaces. In region III, the interfaces have formed and the system is at coexistence, schematically resembling Fig. 2. Increasing the density further leads to a growth of the nematic domain, at the expense of the isotropic domain. Region IV corresponds to the transition to the pure nematic phase, during which the system is characterized by an isotropic droplet in a nematic background. In region V, finally, a single nematic phase is observed.
In Fig. 6, we show the corresponding weight function $`W(S)`$ for both systems. The double-peaked structure is clearly visible. Note that the isotropic peak is significantly higher than the nematic peak. This indicates that the chemical potential used in the simulations is below the coexistence value. Since coexistence is defined by equal weight in the peaks of $`P(N)`$, and not in $`W(S)`$, Fig. 6 cannot be used to obtain the coexistence chemical potential. Instead, $`P(N)`$ must be constructed first, by combining $`W(S)`$ with the single window distributions $`\overline{p}_i(N)`$; Eq.(7) may then be used to extrapolate $`P(N)`$ to coexistence. The resulting coexistence chemical potential equals $`\beta \mu 7.13`$ for both systems. In Fig. 7, the logarithm of $`P(\rho ^{})`$ at coexistence is plotted for both systems, scaled with $`L^2/(2L_\mathrm{x}^2)`$, and the plateaus shifted to zero. In this way, the barrier directly reflects the interfacial tension $`\gamma _{\mathrm{IN}}`$, in units of $`k_\mathrm{B}T/L^2`$ Binder (1982). An important observation is that the peaks in both distributions are separated by a pronounced flat region. This shows that the elongated $`L_\mathrm{z}`$ dimension of the simulation box is sufficient. Moreover, the peak heights are similar, indicating that finite size effects in the lateral dimensions $`L_\mathrm{x}`$ and $`L_\mathrm{y}`$ are also small. Therefore, we conclude that the barrier in Fig. 7 accurately reflects the interfacial tension $`\gamma _{\mathrm{IN}}`$ for soft rods with $`L/D=10`$ and $`\beta ϵ=2`$. The resulting estimate reads as $`\gamma _{\mathrm{IN}}=0.49k_\mathrm{B}T/L^2=0.0049k_\mathrm{B}T/D^2`$.
In the nematic phase, ca. 6000 rods were simulated for system A, and 10,000 for system B. To obtain reliable results, a substantial investment in CPU time is thus required (ca. 3200 CPU hours were invested for system B). Since biased sampling schemes are easy to parallelize, results can typically be obtained within 1-2 weeks on a modern computer cluster. Accurate sampling is especially important around transitions II and IV, and this becomes increasingly difficult in large systems MacDowell et al. (2004). This may already be inferred from the scatter in the data of system B around transition II (arrow in Fig. 7). Transition IV, on the other hand, is sampled with surprisingly little difficulty. The likely explanation is that process II requires the formation of a nematic nucleus whose director is aligned in the $`xy`$-plane. Process IV, on the other hand, does not require any preferred orientation of the (isotropic) nucleus, and is therefore easier to sample.
### IV.4 Consequences for finite size extrapolation
An alternative method to obtain the interfacial tension is to measure $`\gamma _{\mathrm{IN}}(L_\mathrm{x})`$ in cubic systems with edge $`L_\mathrm{x}`$, and use the extrapolation equation of Binder Binder (1982)
$$\gamma _{\mathrm{IN}}(L_\mathrm{x})=\gamma _{\mathrm{IN}}+a/L_\mathrm{x}^2+b\mathrm{ln}(L_\mathrm{x})/L_\mathrm{x}^2,$$
(9)
to estimate $`\gamma _{\mathrm{IN}}`$. In principle, this approach enables estimates of $`\gamma _{\mathrm{IN}}`$ through an elimination of finite size effects, but it requires estimates over a range of values for which $`\gamma _{\mathrm{IN}}(L_\mathrm{x})L_\mathrm{x}^2/k_\mathrm{B}T1`$. In practice, however, one often tries to use Eq.(9) using data from smaller systems. In Ref. Vink and Schilling, 2005, this approach was applied to soft rods with $`L/D=10`$ and $`\beta ϵ=2`$, assuming $`b=0`$ in Eq.(9), leading to $`\gamma _{\mathrm{IN}}=0.0035k_\mathrm{B}T/D^2`$. This estimate differs profoundly from the one of the previous section, implying that finite size extrapolation must be used with care. The issue is investigated further in Fig. 8. Shown is $`\gamma _{\mathrm{IN}}(L_\mathrm{x})`$ as function of $`L_\mathrm{x}`$, where the open squares are data from Ref. Vink and Schilling, 2005, and closed squares data from larger systems obtained in this work. The horizontal line corresponds to the estimate of Fig. 7. Note that the data indeed approach the latter estimate. The curve is a fit to the open squares using Eq.(9) with $`b=0`$, which summarizes the result of Ref. Vink and Schilling, 2005. Clearly, the fit fails to capture the data of the larger systems. Allowing $`b`$ in Eq.(9) to be non-zero will obviously lead to a better fit, but the resulting $`\gamma _{\mathrm{IN}}`$ depends sensitively on the range over which the fit is performed, making this approach somewhat arbitrary. The problem partly stems from the difficulty in distinguishing $`a/L_\mathrm{x}^2`$ numerically from $`b\mathrm{ln}(L_\mathrm{x})/L_\mathrm{x}^2`$, since the range in $`L_\mathrm{x}`$ that can be sampled is rather small. Additionally, in small systems, the interface interactions may be strong. This will introduce corrections to Eq.(9), which may even yield non-monotonic behavior in $`\gamma _{\mathrm{IN}}(L_\mathrm{x})`$ Mon (1988); Berg et al. (1993); Hunter III and Reinhardt (1995). As a result, it is difficult to extract $`\gamma _{\mathrm{IN}}`$ via finite size extrapolation. In contrast, by using an elongated simulation box, and by explicitly checking for finite size effects in the lateral dimensions, $`\gamma _{\mathrm{IN}}`$ can be extracted reliably as shown in Fig. 7. This, we conclude in hindsight, should be the method of choice.
### IV.5 Interfacial tension of hard rods
Finally, we apply nematic order biased sampling to a system of hard rods with $`L/D=15`$, system size $`L_\mathrm{x}=L_\mathrm{y}=10L/3`$ and $`L_\mathrm{z}=10L`$, corresponding to ca. 6000 rods in the nematic phase. An initial estimate of the coexistence chemical potential was obtained via Widom insertion Widom (1963). The nematic order parameter is sampled with resolution $`\mathrm{\Delta }S=0.0025`$ to obtain $`W(S)`$. Combining $`W(S)`$ with the single window distributions $`\overline{p}_i(N)`$ and applying Eq.(7) yields for the coexistence chemical potential $`\beta \mu 5.58`$. The corresponding coexistence distribution is shown in Fig. 9.
Note that $`P(N)`$ for hard rods is prone to substantial statistical error. This is to be expected because the acceptance rate of grand canonical insertion for hard rods is only 0.004%, compared to 8% for soft rods. Nevertheless, the double-peaked structure is clearly visible. From the average peak locations, we obtain $`\rho _{\mathrm{ISO}}^{}=0.193`$ and $`\rho _{\mathrm{NEM}}^{}=0.220`$. The latter densities are consistent with the bulk plateaus in the density profile, indicated by the horizontal lines in Fig. 3. To further check the consistency of our results, an additional simulation in a smaller cubic system with $`L_\mathrm{x}=2.3L`$ was performed; the corresponding coexistence distribution is shown dashed in Fig. 9. Of course, this system is too small to extract the interfacial tension, but the peak positions, and hence the coexistence densities, agree well with those of the larger system. The agreement with bulk densities obtained via Gibbs ensemble simulations Dijkstra et al. (2001) and Gibbs-Duhem integration Bolhuis and Frenkel (1997) is better than 4%. The height of the free energy barrier of the larger system reads as $`\mathrm{\Delta }\mathrm{\Omega }=32\pm 3k_\mathrm{B}T`$, leading to an interfacial tension $`\gamma _{\mathrm{IN}}1.4k_\mathrm{B}T/L^2=0.0064k_\mathrm{B}T/D^2=0.096k_\mathrm{B}T/LD=0.10k_\mathrm{B}T/(L+D)D`$.
## V Discussion and summary
In this paper, we have presented methodic developments that allow for the estimation of the interfacial tension between isotropic and nematic phases in suspensions of rods. The problem is challenging because $`\gamma _{\mathrm{IN}}`$ is very small, and methods that work well for interfaces between isotropic phases become problematic, such as exploiting the anisotropy of the pressure tensor Rowlinson and Widom (1982), or analyzing the capillary wave spectrum (the latter requires very precise data from huge systems Akino et al. (2001)). The novelty of the present approach is to combine grand canonical MC simulations with a bias on the nematic order parameter, and obtain $`\gamma _{\mathrm{IN}}`$ from the grand canonical distribution $`P(N)`$. The advantage is that the problem of “jamming” is largely solved, enabling simulations of large systems.
The current approach also allows for grand canonical simulations of hard rods, enabling a direct comparison to theory. In the Onsager limit of infinite rod length, theoretical estimates of $`\gamma _{\mathrm{IN}}`$ typically range from 0.156 Shundyak and van Roij (2001) to 0.34 McMullen (1988), in units of $`k_\mathrm{B}T/LD`$. As expected, this exceeds the value for hard rods obtained in this work ($`\gamma _{\mathrm{IN}}0.096k_\mathrm{B}T/LD`$) because $`L/D=15`$ is still far from the Onsager limit. As shown by experiment Chen and Gray (2002) and theory Velasco et al. (2002), $`\gamma _{\mathrm{IN}}`$ increases with $`L/D`$. The latter theory is based on the Somoza-Tarazona density functional Somoza and Tarazona (1989) and its main findings are summarized in Fig. 10. Shown are the coexistence densities (left axis) and the interfacial tension (right axis) as function of the rod elongation $`L/D`$, where we have adopted the units of Ref. Velasco et al., 2002. Open and closed squares show the theoretical density of the isotropic and the nematic phase, respectively; the star and the cross are the corresponding simulation estimates of this work. Closed triangles are the theoretical interfacial tension, where the line serves to guide the eye; the open triangle represents the simulation estimate of $`\gamma _{\mathrm{IN}}`$. Theoretical estimates are reported for $`L/D=5,10,20`$, but unfortunately not for $`L/D=15`$. This makes a direct comparison difficult; interpolation of the theoretical results, however, seems in good agreement with our simulation results, as may be inferred from Fig. 10.
A typical rod dimension in experiments is $`L=150`$ nm and $`\gamma _{\mathrm{IN}}=0.00083`$ mN/m Chen and Gray (2002). For $`T=298`$ K, this length translates into $`0.00025`$ mN/m using our estimate of $`\gamma _{\mathrm{IN}}`$. Obviously, this estimate differs from the experimental one because the hard-rod fluid is a simplified model, but it is reassuring to see that the order of magnitude is confirmed.
The current biased sampling scheme thus seems well suited to simulate IN coexistence, even for hard interactions. Our scheme may also be useful for the application of transition path sampling Bolhuis et al. (2002) to anisotropic colloidal systems, since it can provide valuable starting paths; work along these lines is in progress. The remaining bottleneck is the low acceptance rate of grand canonical insertion. It remains a challenge to address this final problem. Since the overall density around the IN transition is low, it is anticipated that higher acceptance rates can be realized using smarter insertion schemes. To develop such schemes would be the subject of future work.
###### Acknowledgements.
We thank the Deutsche Forschungsgemeinschaft (DFG) for support (TR6/A5 and TR6/D5) and K. Binder, M. Müller, P. Virnau, P. van der Schoot, R. van Roij and M. Dijkstra for a careful reading of the manuscript and/or helpful suggestions. TS is supported by the Emmy Noether program of the DFG. Allocation of computer time on the JUMP at the Forschungszentrum Jülich is gratefully acknowledged. |
warning/0506/math0506273.html | ar5iv | text | # Nilpotency, almost nonnegative curvature, and the gradient flow on Alexandrov spaces
## 1. Introduction
Almost nonnegatively curved manifolds were introduced by Gromov in the late 70s \[Gro80\], with the most significant contributions to their study made by Fukaya and Yamaguchi in \[Yam91\] and \[FY92\]. Building on their ideas, in the present article we establish several new properties of these manifolds which yield, in particular, new topological obstructions to almost nonnegative curvature. Our techniques also provide simplified proofs of many results from \[FY92\].
A closed smooth manifold is said to be almost nonnegatively curved if it can Gromov–Hausdorff converge to a single point under a lower curvature bound. By rescaling, this definition is equivalent to the following one, which we will employ thruout this article.
###### Definition 1.0.1.
A closed smooth manifold $`M`$ is called almost nonnegatively curved if it admits a sequence of Riemannian metrics $`\{g_n\}_n`$ whose sectional curvatures and diameters have the following bounds
$$\mathrm{sec}(M,g_n)1/n\text{and}\mathrm{diam}(M,g_n)1/n.$$
Almost nonnegatively curved manifolds generalize almost flat as well as nonnegatively curved manifolds. The main source of examples comes from a theorem of Fukaya and Yamaguchi in \[FY92\]. It states that if $`FEB`$ is a fiber bundle over an almost nonnegatively curved manifold $`B`$ whose fiber $`F`$ is compact and admits a nonnegatively curved metric which is invariant under the structure group, then the total space $`E`$ is almost nonnegatively curved. Further examples are given by closed manifolds which admit cohomogeneity one actions of compact Lie groups (compare \[ST04\]).
In this work we combine collapsing techniques with a non-smooth analogue of the gradient flow of concave functions on Alexandrov spaces. This notion is based on the construction of gradient curves of $`\lambda `$-concave functions used in \[PP96\] and bears many similarities to the Sharafutdinov retraction \[Sha78\]. Gradient flows on general metric spaces were considered in \[AGS\]. The gradient flow on Alexandrov spaces plays a key role in the proofs of two of the three main results in this paper, and we believe that it should also prove useful for dealing with other problems related to collapsing under a lower curvature bound.
### 1.1.
Let us first briefly recall some previously known results:
Let $`M=M^m`$ be an almost nonnegatively curved $`m`$-manifold.
1. Gromov proved in \[Gro78\] that the minimal number of generators of the fundamental group $`\pi _1M`$ of $`M`$ can be estimated by a constant $`C_1(m)`$ depending only on $`m`$, and in \[Gro81\] that the sum of Betti numbers of $`M`$ with respect to any field of coefficients does not exceed some uniform constant $`C_2=C_2(m)`$.
2. In \[Yam91\], Yamaguchi showed that, up to a finite cover, the manifold $`M`$ fibers over a $`b_1(M;)`$-dimensionsal torus and that $`M`$ is diffeomorphic to a torus if $`b_1(M;)=m`$. (Here $`b_1(M;)`$ denotes the first Betti number with real coefficients.)
3. In \[FY92\], Fukaya and Yamaguchi proved that
1. $`\pi _1M`$ is almost nilpotent; that is, $`\pi _1M`$ contains a nilpotent subgroup of finite index
2. $`\pi _1M`$ is $`C_3(m)`$-solvable; that is, contains a solvable subgroup of index at most $`C_3(m)`$
4. If a closed manifold has negative Yamabe constant, then it cannot volume collapse with scalar curvature bounded from below (see \[Sch89, LeB01\]). In particular, no such manifold can be almost nonnegatively curved.
5. The $`\widehat{A}`$-genus of a closed spin manifold $`X`$ of almost nonnegative Ricci curvature satisfies the inequality $`\widehat{A}(X)2^{dimX/2}`$ (\[Gro82\], \[Gal83\]).
### 1.2.
Our first result concerns the relation between curvature bounds and the actions of the fundamental group on the higher homotopy groups.
Recall that an action by automorphisms of a group $`G`$ on an abelian group $`V`$ is called nilpotent if $`V`$ admits a finite sequence of $`G`$-invariant subgroups
$$V=V_0V_1\mathrm{}V_k=0$$
such that the induced action of $`G`$ on $`V_i/V_{i+1}`$ is trivial for any $`i`$. A connected CW-complex $`X`$ is called *nilpotent* if $`\pi _1X`$ is a nilpotent group that operates nilpotently on $`\pi _kX`$ for every $`k2`$.
Nilpotent spaces play an important role in topology since they enjoy some of the best homotopy-theoretic properties of simply connected spaces, like a Whitehead theorem or reasonable Postnikov towers. Furthermore, unlike the category of simply connected spaces, the category of nilpotent ones is closed under many constructions such as the based loop space functor or the formation of function spaces, and group-theoretic functors, like localization and completion, have topological extensions in this category.
###### Theorem A (Nilpotency Theorem).
Let $`M`$ be a closed almost nonnegatively curved manifold. Then a finite cover of $`M`$ is a nilpotent space.
It would be interesting to know whether the order of this covering can be estimated solely in terms of the dimension of $`M`$.
###### Example 1.2.1.
Let $`h:𝕊^3\times 𝕊^3𝕊^3\times 𝕊^3`$ be defined by
$$h:(x,y)(xy,yxy).$$
This map is a diffeomorphism with inverse given by
$$h^1:(u,v)(u^2v^1,vu^1).$$
The induced map $`h_{}`$ on $`\pi _3(𝕊^3\times 𝕊^3)`$ is given by the matrix $`A_h=\left(\begin{array}{cc}1& 1\\ 1& 2\end{array}\right)`$. Notice that the eigenvalues of $`A_h`$ are different from $`1`$ in absolute value. Let $`M`$ be the mapping cylinder of $`h`$. Clearly, $`M`$ has the structure of a fiber bundle
$$𝕊^3\times 𝕊^3M𝕊^1,$$
and the action of $`\pi _1M`$ on $`\pi _3M^2`$ is generated by $`A_h`$. In particular, $`M`$ is not a nilpotent space and hence, by Theorem A, it does not admit almost nonnegative curvature. This fact doesn’t follow from any previously known results.
### 1.3.
Our next main result provides an affirmative answer to a conjecture of Fukaya and Yamaguchi \[FY92, Conjecture 0.15\].
###### Theorem B ($`C`$-Nilpotency Theorem for $`\pi _1`$).
Let $`M`$ be an almost nonnegatively curved $`m`$-manifold. Then $`\pi _1M`$ is $`C(m)`$-nilpotent; that is, $`\pi _1M`$ contains a nilpotent subgroup of index at most $`C(m)`$.
Notice that Theorem B is new even for manifolds of nonnegative curvature.
###### Example 1.3.1.
For any $`C>0`$ there exist prime numbers $`p>q>C`$ and a finite group $`G_{pq}`$ of order $`pq`$ which is solvable but not nilpotent. In particular, $`G_{pq}`$ does not contain any nilpotent subgroup of index less than or equal to $`C`$.
Whereas none of the results mentioned so far excludes $`G_{pq}`$ from being the fundamental group of some almost nonnegatively curved $`m`$-manifold, Theorem B shows that for $`C>C(m)`$ none of the groups $`G_{pq}`$ can be realized as the fundamental group of such a manifold.
### 1.4.
In \[FY92\] Fukaya and Yamaguchi also conjectured that a finite cover of an almost nonnegatively Ricci curved manifold $`M`$ fibers over a nilmanifold with a fiber which has nonnegative Ricci curvature and whose fundamental group is finite. This conjecture was later refuted by Anderson \[And92\].
It is, on the other hand, very natural to consider this conjecture in the context of almost nonnegative sectional curvature. In fact, here Yamaguchi’s fibration theorem (\[Yam91\]) and the results of \[FY92\] easily imply that a finite cover of an almost nonnegatively curved manifold admits a map onto a nilmanifold whose homotopy fiber is a simply connected closed manifold.
From mere topology, it is, however, not clear whether this homotopy fibration can actually always be made into a genuine fiber bundle. Our next result shows that this is indeed true, and that for manifolds of almost nonnegative sectional curvature Fukaya’s and Yamaguchi’s original conjecture essentially does hold.
###### Theorem C (Fibration Theorem).
Let $`M`$ be an almost nonnegatively curved manifold. Then a finite cover $`\stackrel{~}{M}`$ of $`M`$ is the total space of a fiber bundle
$$F\stackrel{~}{M}N$$
over a nilmanifold $`N`$ with a simply connected fiber $`F`$. Moreover, the fiber $`F`$ is almost nonnegatively curved in the generalized sense as defined below.
###### Definition 1.4.1.
A closed smooth manifold $`M`$ is called almost nonnegatively curved in the generalized sense if for some nonnegative integer $`k`$ there exists a sequence of complete Riemannian metrics $`g_n`$ on $`M\times ^k`$ and points $`p_nM\times ^k`$ such that
1. the sectional curvatures of the metric balls of radius $`n`$ around $`p_n`$ satisfy
$$\mathrm{sec}(B_n(p_n))1/n;$$
2. for $`n\mathrm{}`$ the pointed Riemannian manifolds $`((M\times ^k,g_n),p_n)`$ converge in the pointed Gromov–Hausdorff distance to $`(^k,0)`$;
3. the regular fibers over $`0`$ are diffeomorphic to $`M`$ for all large $`n`$.
Due to Yamaguchi’s fibration theorem \[Yam91\], manifolds which are almost nonnegatively curved in the generalized sense play the same central role in collapsing under a lower curvature bound as almost flat manifolds do in the Cheeger–Fukaya–Gromov theory of collapsing with bounded curvature (see \[CFG92\]).
It is not known whether all manifolds which are almost nonnegatively curved in the generalized sense are almost nonnegatively curved. Clearly, if $`k=0`$, this definition reduces to the standard one. Moreover, it is easy to see that all results of the present article, as well as all results about almost nonnegatively curved manifolds mentioned earlier (except possibly for the ones concerning the $`\widehat{A}`$-genus and Yamabe constant), hold for manifolds which are almost nonnegatively curved in the sense of Definition 1.4.1.
### 1.5.
Let us describe the structure of the remaining sections of this article.
In section 2, after providing some necessary background from Alexandrov geometry, we introduce the gradient flow of the square of a distance function. It serves as one of the main technical tools in the proofs of theorem A and theorem B.
In section 3 we prove Theorem A by a direct application of the gradient flow technique.
In section 4 we prove Theorem B. The proof is also based on the gradient flow, but is more involved and employs further technical tools such as “limit fundamental groups” of Alexandrov spaces.
In section 5 we prove Theorem C. This section is completely independent from the rest of the article.
In section 6 we discuss some further open questions related to our results.
Acknowledgements. We would like to express our thanks to the following people for helpful conversations during the preparation of this work: I. Belegradek, V. Gorbunov, E. Formanek, I. Kapovitch, A. Lytchak, R. Matveyev, D. Robinson, D. Sullivan, B. Wilking, and Yu. Zarkhin.
We are also especially grateful to Yunhui Wu for noticing a mistake in the proof of Theorem 4.1.1 in the original version of this article.
## 2. Alexandrov geometry and the gradient flow
This section provides necessary background in Alexandrov geometry. The results of sections 2.1–2.3 are mostly repeated from \[PP96\], \[Pet95\] and \[Pet07\]. The reader may consult \[BGP92\] for a general reference on Alexandrov spaces.
### 2.1. $`\lambda `$-concave functions
###### Definition 2.1.1.
(for a space without boundary) Let $`A`$ be an Alexandrov space without boundary. A Lipschitz function $`f:A`$ is called $`\lambda `$-*concave* if for any unit speed minimizing geodesic $`\gamma `$ in $`A`$, the function
$$f\gamma (t)\lambda t^2/2$$
is concave.
If $`A`$ is an Alexandrov space with boundary, then its double $`\stackrel{~}{A}`$ is also an Alexandrov space (see \[Per91, 5.2\]). Let $`\text{p}:\stackrel{~}{A}A`$ be the canonical map. Given a function $`f`$ on $`A`$, set $`\stackrel{~}{f}=f\text{p}`$.
###### Definition 2.1.2.
(for a space with boundary) Let $`A`$ be an Alexandrov space with boundary. A Lipschitz function $`f:A`$ is called $`\lambda `$-*concave* if for any unit speed minimizing geodesic $`\gamma `$ in $`\stackrel{~}{A}`$, the function
$$\stackrel{~}{f}\gamma (t)\lambda t^2/2$$
is concave.
###### Remark 2.1.3.
Notice that the restriction of a linear function on $`^n`$ to a ball is not $`0`$-concave in this sense.
###### Remark 2.1.4.
In the above definitions, the Lipschitz condition is only technical. With some extra work, all results of this section can be extended to continuous functions.
### 2.2. Tangent cone and differential
Given a point $`p`$ in an Alexandrov space $`A`$, we denote by $`T_p=T_p(A)`$ the tangent cone at $`p`$.
If $`d`$ denotes the metric of an Alexandrov space $`A`$, let us denote by $`\lambda A`$ the space $`(A,\lambda d)`$. Let $`i_\lambda :\lambda AA`$ be the canonical map. The limit of $`(\lambda A,p)`$ for $`\lambda \mathrm{}`$ is the tangent cone $`T_p`$ at $`p`$ (see \[BGP92, 7.8.1\]).
###### Definition 2.2.1.
For any function $`f:A`$ the function $`d_pf:T_p`$ such that
$$d_pf=\underset{\lambda \mathrm{}}{lim}\lambda (fi_\lambda f(p))$$
is called the differential of $`f`$ at $`p`$.
It is easy to see that for a $`\lambda `$-concave function $`f`$ the differential $`d_pf`$ is defined everywhere, and that $`d_pf`$ is a $`0`$-concave function on the tangent cone $`T_p`$.
###### Definition 2.2.2.
Given a $`\lambda `$-concave function $`f:A`$, a point $`pA`$ is called critical point of $`f`$ if $`d_pf0`$.
### 2.3. Gradient curves
With a slight abuse of notation we will call elements of the tangent cone $`T_p`$ the “tangent vectors” at $`p`$. The origin of $`T_p`$ plays the role of the zero vector and is denoted by $`o=o_p`$. For a tangent vector $`v`$ at $`p`$ we define its absolute value $`|v|`$ as the distance $`|ov|`$ in $`T_p`$. For two tangent vectors $`u`$ and $`v`$ at $`p`$ we can define their “scalar product”
$$u,v=(|u|^2+|v|^2|uv|^2)/2=|u||v|\mathrm{cos}\alpha ,$$
where $`\alpha =\mathrm{}uov`$ in $`T_p`$.
For two points $`p,qA`$ we define $`\mathrm{log}_pq`$ to be a tangent vector $`v`$ at $`p`$ such that $`|v|=|pq|`$ and such that the direction of $`v`$ coincides with a direction from $`p`$ to $`q`$ (if such a direction is not unique, we choose any one of them). Given a curve $`\gamma (t)`$ in $`A`$, we denote by $`\gamma ^+(t)`$ the right and by $`\gamma ^{}(t)`$ the left tangent vectors to $`\gamma (t)`$, where, respectively,
$$\gamma ^\pm (t)T_{\gamma (t)},\gamma ^\pm (t)=\underset{\epsilon +0}{lim}\frac{\mathrm{log}_{\gamma (t)}\gamma (t\pm \epsilon )}{\epsilon }.$$
For a real function $`f(t)`$, $`t`$, we denote by $`f^+(t)`$ its right derivative and by $`f^{}(t)`$ its left derivative. Note that our sign convention (which is chosen to agree with the notion of right and left derivatives of curves) is not quite standard. For example,
$$\text{if}f(t)=t\text{then}f^+(t)1\text{and}f^{}(t)1.$$
###### Definition 2.3.1.
Given a $`\lambda `$-concave function $`f`$ on $`A`$, a vector $`gT_p(A)`$ is called a gradient of $`f`$ at $`pA`$ (in short: $`g=_pf`$) if
(i) $`d_pf(x)g,x`$ for any $`xT_p`$, and
(ii) $`d_pf(g)=g,g.`$
It is easy to see that any $`\lambda `$-concave function has a uniquely defined gradient vector field. Moreover, if $`d_pf(x)0`$ for all $`xT_p`$, then $`_pf=o`$ (here $`o`$ denotes the origin of the tangent cone $`T_p`$); otherwise,
$$_pf=d_pf(\xi )\xi $$
where $`\xi `$ is the (necessarily unique) unit vector for which the function $`d_pf`$ attains its maximum.
Moreover, for any minimizing geodesic $`\gamma :[a,b]U`$ parameterized by arclength, the following inequality holds:
(2.3.1)
$$\gamma ^+(a),_{\gamma (a)}f+\gamma ^{}(b),_{\gamma (b)}f\lambda (ba).$$
Indeed,
$`\gamma ^+(a),_{\gamma (a)}f+\gamma ^{}(b),_{\gamma (b)}f`$ $`d_{\gamma (a)}f(\gamma ^+(a))+d_{\gamma (b)}f(\gamma ^{}(b))=`$
$`=(f\gamma )^+|_a+(f\gamma )^{}|_b`$
$`\lambda (ba).`$
###### Definition 2.3.2.
A curve $`\alpha :[a,b]A`$ is called an $`f`$-*gradient curve* if for any $`t[a,b]`$
$$\alpha ^+(t)=_{\alpha (t)}f.$$
###### Proposition 2.3.3.
Given a $`\lambda `$-concave function $`f:A`$ and a point $`pA`$ there is a unique gradient curve $`\alpha :[0,\mathrm{})A`$ such that $`\alpha (0)=p`$.
Moreover, if $`\alpha `$ and $`\beta `$ are two $`f`$-gradient curves, then
$$|\alpha (t_1)\beta (t_1)||\alpha (t_0)\beta (t_0)|\mathrm{exp}(\lambda (t_1t_0))\text{for all}t_1t_0.$$
The gradient curve can be constructed as a limit of broken geodesics, made up of short segments with directions close to the gradient. The convergence, uniqueness, as well as the last inequality in Proposition 2.3.3 follow from inequality (2.3.1) above, while Corollary 2.3.5 below guarantees that the limit is indeed a gradient curve, having a unique right tangent vector at each point.
###### Lemma 2.3.4.
Let $`A_n\stackrel{\mathrm{GH}}{}A`$ be a sequence of Alexandrov spaces with curvature $`k`$ which Gromov–Hausdorff converges to an Alexandrov space $`A`$.
Let $`f_nf`$, where $`f_n:A_n`$ is a sequence of $`\lambda `$-concave functions converging to $`f:A`$.
Let $`p_np`$, where $`p_nA_n`$ and $`pA`$.
Then
$$|_pf|\underset{n\mathrm{}}{lim\; inf}|_{p_n}f_n|.$$
###### Corollary 2.3.5.
For any $`\lambda `$-concave function $`f`$ on $`A`$ the function
$$p|_pf|$$
is lower semicontinuous; that is, for any sequence of points $`p_nA`$, $`p_np`$, we have
$$|_pf|\underset{n\mathrm{}}{lim\; inf}|_{p_n}f|.$$
###### Proof of Lemma 2.3.4.
Fix an $`\epsilon >0`$ and choose $`q`$ near $`p`$ such that
$$\frac{f(q)f(p)}{|pq|}>|_pf|\epsilon .$$
Now choose $`q_nA_n`$ such that $`q_nq`$. If $`|pq|`$ is sufficiently small and $`n`$ is sufficiently large, the $`\lambda `$-concavity of $`f_n`$ then implies that
$$\underset{n\mathrm{}}{lim\; inf}\frac{d_{p_n}f_n(v_n)}{|v_n|}|_pf|2\epsilon \text{for}v_n=\mathrm{log}_{p_n}(q_n)T_{p_n}(A_n).$$
Therefore,
$$\underset{n\mathrm{}}{lim\; inf}|_{p_n}f_n||_pf|2\epsilon \text{for any}\epsilon >0;$$
that is,
$$\underset{n\mathrm{}}{lim\; inf}|_{p_n}f_n||_pf|.$$
###### Lemma 2.3.6.
Let $`f`$ be a $`\lambda `$-concave function, $`\lambda 0`$ and $`\alpha (t)`$ be an $`f`$-gradient curve, and let $`\overline{\alpha }(s)`$ be its reparameterization by arclength. Then $`f\overline{\alpha }`$ is $`\lambda `$-concave.
###### Proof.
$$(f\overline{\alpha })^+(s_0)=|_{\overline{\alpha }(s_0)}f|\frac{d_{\overline{\alpha }(s_0)}f(\mathrm{log}_{\overline{\alpha }(s_0)}(\overline{\alpha }(s_1))}{|\overline{\alpha }(s_1)\overline{\alpha }(s_0)|}$$
$$\frac{f(\overline{\alpha }(s_1))f(\overline{\alpha }(s_0))\lambda |\overline{\alpha }(s_1)\overline{\alpha }(s_0)|^2/2}{|\overline{\alpha }(s_1)\overline{\alpha }(s_0)|}$$
$$\frac{f(\overline{\alpha }(s_1))f(\overline{\alpha }(s_0))}{s_1s_0}\lambda |\overline{\alpha }(s_1)\overline{\alpha }(s_0)|/2.$$
Since $`\frac{|\overline{\alpha }(s_1)\overline{\alpha }(s_0)|}{(s_1s_0)}1`$ as $`s_1s_0+`$, it follows that $`f\overline{\alpha }`$ is $`\lambda `$-concave.
### 2.4. The Gradient Flow on Alexandrov Spaces
Let $`f`$ be a $`\lambda `$-concave function on an Alexandrov space $`A`$. Consider the map $`\mathrm{\Phi }_f^T:AA`$ defined as follows: $`\mathrm{\Phi }_f^T(x)=\alpha _x(T)`$, where $`\alpha _x:[0,\mathrm{})A`$ is the $`f`$-gradient curve with $`\alpha _x(0)=x`$. The map $`\mathrm{\Phi }_f^T`$ is called *$`f`$-gradient flow at time $`T`$*. From Proposition 2.3.3 it is clear that $`\mathrm{\Phi }_f^T`$ is an $`\mathrm{exp}(\lambda T)`$-Lipschitz map. Next we want to prove that this map behaves nicely under Gromov–Hausdorff-convergence.
###### Theorem 2.4.1.
Let $`A_n\stackrel{\mathrm{GH}}{}A`$ be a sequence of Alexandrov spaces with curvature $`k`$ which converges to an Alexandrov space $`A`$.
Let $`f_nf`$, where $`f_n:A_n`$ is a sequence of $`\lambda `$-concave functions and $`f:A`$.
Then $`\mathrm{\Phi }_{f_n}^T\mathrm{\Phi }_f^T`$.
Theorem 2.4.1 immediately follows from the following Lemma:
###### Lemma 2.4.2.
Let $`A_n\stackrel{\mathrm{GH}}{}A`$ be a sequence of Alexandrov spaces with curvature $`k`$ which converges to an Alexandrov space $`A`$.
Let $`f_nf`$, where $`f_n:A_n`$ is a sequence of $`\lambda `$-concave functions and $`f:A`$.
Let $`\alpha _n:[0,\mathrm{})A_n`$ be the sequence of $`f_n`$-gradient curves with $`\alpha _n(0)=p_n`$ and let $`\alpha :[0,\mathrm{})A`$ be the $`f`$-gradient curve with $`\alpha (0)=p`$.
Then $`\alpha _n\alpha `$.
###### Proof.
We may assume without loss of generality that $`f`$ has no critical points. (Otherwise consider instead the sequence $`A_n^{}=A_n\times `$ with $`f_n^{}(a\times x)=f_n(a)+x`$.)
Let $`\overline{\alpha }_n(s)`$ denote the reparameterization of $`\alpha _n(t)`$ by arc length. Since all $`\overline{\alpha }_n`$ are $`1`$-Lipschitz, we can choose a converging subsequence from any subsequence of $`\overline{\alpha }_n`$. Let $`\overline{\beta }:[0,\mathrm{})A`$ be its limit.
Clearly, $`\overline{\beta }`$ is also 1-Lipschitz and hence $`|\overline{\beta }^+|1`$. Therefore, by Lemma 2.3.4,
$`\underset{n\mathrm{}}{lim}f_n\overline{\alpha }_n|_a^b`$ $`=\underset{n\mathrm{}}{lim}{\displaystyle _a^b}|_{\overline{\alpha }_n(s)}f_n|𝑑s`$
$`{\displaystyle _a^b}|_{\overline{\beta }(s)}f|𝑑s`$
$`{\displaystyle _a^b}d_{\beta (t)}f(\beta ^+(t))𝑑t=f\beta |_a^b.`$
On the other hand, since $`\overline{\alpha }_n\overline{\beta }`$ and $`f_nf`$ we have $`f_n\overline{\alpha }_n|_a^bf\overline{\beta }|_a^b`$. Therefore, in both of these inequalities in fact equality holds.
Hence, $`|_{\overline{\beta }(s)}f|=lim_n\mathrm{}|_{\overline{\alpha }_n(s)}f_n|`$, $`|\overline{\beta }^+(s)|=1`$ and the directions of $`\overline{\beta }^+(s)`$ and $`_{\overline{\beta }(s)}f`$ coincide almost everywhere. This implies that $`\overline{\beta }(s)`$ is a gradient curve reparameterized by arc length. In other words, if $`\overline{\alpha }(s)`$ denotes the reparameterization of $`\alpha (t)`$ by arc length, then $`\overline{\beta }(s)=\overline{\alpha }(s)`$ for all $`s`$. It only remains to show that the original parameter $`t_n(s)`$ of $`\alpha _n`$ converges to the original parameter $`t(s)`$ of $`\alpha `$.
Notice that $`|_{\overline{\alpha }_n(s)}f_n|dt_n=ds`$ or $`dt_n/ds=ds/d(f_n\overline{\alpha }_n)`$. Likewise, $`dt/ds=ds/d(f\overline{\alpha })`$. Then the convergence $`t_nt`$ follows from the $`\lambda `$-concavity of $`f_n\overline{\alpha }_n`$ (see Lemma 2.3.6) and the convergence $`f_n\overline{\alpha }_nf\overline{\alpha }.`$
### 2.5. Gradient balls
Let $`A`$ be an Alexandrov space and let $`SA`$ be a subset of $`A`$. A function $`f:A`$ which can be represented as
$$f=\underset{i}{}\theta _i\frac{\mathrm{dist}_{a_i}^2}{2}\text{with }\theta _i0,\underset{i}{}\theta _i=1\text{and}a_iS$$
will be called cocos-function with respect to $`S`$ (where “cocos” stands for convex combination of squares of distance functions). A broken gradient curve for a collection of such functions will be called cocos-curve with respect to $`S`$.
For $`pA`$ and $`T,r_+`$, let us define “the gradient ball with center $`p`$ and radius $`T`$ with respect to $`B_r(p)`$”, $`\beta _T^r(p)`$, as the set of all end points of cocos-curves with respect to $`B_r(p)`$ that start at $`p`$ with total time $`T`$.
###### Lemma 2.5.1.
1. There exists $`T=T(m)_+`$ such that for any $`m`$-dimensional Alexandrov space $`A`$ with curvature $`1`$ and any $`qA`$ there is a point $`pA`$ such that
1. $`|pq|1`$, and
2. $`B_1(p)\beta _T^1(p)`$.
2. There exists $`T^{}=T^{}(m)_+`$ such that the following holds. Let $`A`$ be an Alexandrov space which is a quotient $`A=\stackrel{~}{A}/\mathrm{\Gamma }`$ of an $`m`$-dimensional Alexandrov space $`\stackrel{~}{A}`$ with curvature $`1`$ by a discrete action of a group of isometries $`\mathrm{\Gamma }`$. Let $`qA`$ and $`p=p(q)A`$ be as in part I above.
Then for any lift $`\stackrel{~}{p}\stackrel{~}{A}`$ of $`p`$ one has that $`B_1(\stackrel{~}{p})\beta _T^{}^1(\stackrel{~}{p})`$.
###### Proof.
The proof is similar to the construction of a strained point in an Alexandrov space (see \[BGP92\]).
Set $`\delta =10^m`$. Take $`a_1=q`$ and take $`b_1`$ to be a farthest point from $`a_1`$ in the closed ball $`\overline{B}_1(a_1)`$. Take $`a_2`$ to be a midpoint of $`a_1b_1`$ and let $`b_2`$ be a farthest point from $`a_2`$ such that $`|a_1b_2|=|a_1a_2|`$ and $`|a_2b_2|\delta |a_1b_1|`$, etc. On the $`k`$-th step we have to take $`a_k`$ to be a midpoint of $`a_{k1}b_{k1}`$ and $`b_k`$ to be a farthest point from $`a_k`$ such that $`|a_ib_k|=|a_ia_k|`$ for all $`i<k`$ and $`|a_kb_k|\delta |a_{k1}b_{k1}|`$.
After $`m`$ steps, take $`p`$ to be a midpoint of $`a_mb_m`$. We only have to check that we can find $`T=T(m)`$ such that $`\beta _T^1(p)B_1(p)`$.
Let $`t_i`$ be the minimal time such that $`B_{|a_ib_i|/\delta ^m}(p)\beta _{t_i}^1(p)`$. Then one can take $`T=t_1`$. Therefore it is enough to give estimates for $`t_m`$ and $`t_{k1}/t_k`$ only in terms of $`\delta `$ and $`m`$. Looking at the ends of broken gradient curves starting at $`p`$ for the functions $`\mathrm{dist}_p^2/2`$, $`\mathrm{dist}_{a_i}^2/2`$ and $`\mathrm{dist}_{b_i}^2/2`$, we easily see that $`t_n1/\delta ^m`$. Now, looking at the ends of broken gradient curves starting at $`B_{|a_{k1}b_{k1}|/\delta ^m}(p)`$ for the functions $`\mathrm{dist}_p^2/2`$, $`\mathrm{dist}_{a_i}^2/2`$ and $`\mathrm{dist}_{b_i}^2/2`$, we have that $`t_{k1}/t_k1/\delta ^m`$. Therefore $`t_11/\delta ^{m^2}=10^{m^3}`$. This finishes the proof of part (I).
For part (II), notice that
1. for any $`r,t>0`$ we have $`\beta _t^r(p)B_{re^t}(p);`$
2. if $`\beta _t^r(p)B_\rho (p)`$, then $`\beta _{t+\tau }^r(p)B_{\rho e^\tau }(p);`$
3. if $`\rho =|px|`$ and $`x\beta _t^{r+\rho }(p)`$, then $`\beta _\tau ^r(x)\beta _{t+\tau }^{r+\rho }(p)`$.
Take $`\epsilon =e^T/4`$ and apply part (I) of the lemma to $`\frac{1}{\epsilon }\stackrel{~}{A}`$ to find a point $`p^{}\stackrel{~}{A}`$ such that $`|\stackrel{~}{p}p^{}|\epsilon `$ and $`B_\epsilon (p^{})\beta _T^\epsilon (p^{})\stackrel{~}{A}`$. Then for some deck transformation $`\gamma `$ we have $`\gamma (p^{})\beta _T^\epsilon (\stackrel{~}{p})B_{\epsilon e^T}(\stackrel{~}{p})`$. Therefore it holds that $`\gamma (p^{})B_{1/2}(\stackrel{~}{p})`$. Hence, taking
$$T^{}=2T+1/\epsilon =2T+4e^T,$$
we obtain
$$\beta _T^{}^1(\stackrel{~}{p})\beta _{T+1/\epsilon }^\epsilon (\gamma (p^{}))B_1(\stackrel{~}{p}).$$
### 2.6. Short basis.
We will use the following construction due to Gromov.
Given an Alexandrov space $`A`$ with a marked point $`pA`$, and a group $`\mathrm{\Gamma }`$ acting discretely on $`A=(A,d)`$ one can define a short basis of the action of $`\mathrm{\Gamma }`$ at $`p`$ as follows:
For $`\gamma \mathrm{\Gamma }`$ define the norm of $`\gamma `$ by the formula $`|\gamma |=d(p,\gamma (p))`$. Choose $`\gamma _1\mathrm{\Gamma }`$ with the minimal norm in $`\mathrm{\Gamma }`$. Next choose $`\gamma _2`$ to have minimal norm in $`\mathrm{\Gamma }\backslash \gamma _1`$. On the $`n`$-th step choose $`\gamma _n`$ to have minimal norm in $`\mathrm{\Gamma }\backslash \gamma _1,\gamma _2,\mathrm{},\gamma _{n1}`$. The sequence $`\{\gamma _1,\gamma _2,\mathrm{}\}`$ is called a *short basis* of $`\mathrm{\Gamma }`$ at $`p`$. In general, the number of elements of a short basis can be finite or infinite. In the special case of the action of the fundamental group $`\pi _1(A,p)`$ on the universal cover of $`A`$ one speaks of the short basis of $`\pi _1(A,p)`$.
It is easy to see that for a short basis $`\{\gamma _1,\gamma _2,\mathrm{}\}`$ of the fundamental group of an Alexandrov space $`A`$ the following is true:
1. If $`A`$ has diameter $`d`$ then $`|\gamma _i|2d`$.
2. If $`A`$ is compact then $`\{\gamma _i\}`$ is finite.
3. For any $`i>j`$ we have $`|\gamma _i||\gamma _j^1\gamma _i|`$.
The third property implies that if $`\stackrel{~}{p}\stackrel{~}{A}`$ is in the preimage of $`p`$ in the universal cover $`\stackrel{~}{A}`$ of $`A`$ and $`\stackrel{~}{p}_i=\gamma _i(\stackrel{~}{p})`$, then
$$|\stackrel{~}{p}_i\stackrel{~}{p}_j|\mathrm{max}\{|\stackrel{~}{p}\stackrel{~}{p}_i|,|\stackrel{~}{p}\stackrel{~}{p}_j|\}.$$
As was observed by Gromov, if $`A`$ is an Alexandrov space with curvature $`\kappa `$ and diameter $`d`$, the last inequality implies that $`\mathrm{}\stackrel{~}{p}_i\stackrel{~}{p}\stackrel{~}{p}_j>\delta =\delta (\kappa ,d)>0`$. This yields an upper bound on the number of elements of a short basis in terms of $`\kappa ,d`$ and the dimension of $`A`$.
## 3. Nilpotency of almost nonnegatively curved manifolds
In this section we prove Theorem A.
### 3.1. Preliminary lemmas
Let $`M`$ be an almost nonnegatively curved manifold. Let us denote by $`M_n=(M,g_n)`$, $`n`$, a sequence of Riemannian metrics on $`M`$ such that $`\mathrm{sec}(M_n)1/n`$ and $`\mathrm{diam}(M_n)1/n`$. Let $`\stackrel{~}{M}M`$ be the universal covering and $`\stackrel{~}{M}_nM_n`$ be the universal Riemannian covering of $`M_n`$ (that is, $`\stackrel{~}{M}_n`$ is $`\stackrel{~}{M}`$ equipped with the pullback of the Riemannian metric $`g_n`$).
###### Key Lemma 3.1.1.
Given $`\epsilon >0`$ and $`r_2>r_1>0`$, let $`\stackrel{~}{M}_nB_{r_2}(p_n)B_{r_1}(p_n)`$. Then, for $`n`$ sufficiently large, there is a $`(1+\epsilon )`$-Lipschitz map
$$\mathrm{\Phi }_n:B_{r_2}(p_n)B_{r_1}(p_n)$$
which is homotopic to the identity on $`B_{r_2}(p_n)`$.
###### Proof.
Fix $`R>>r_2`$ (here $`R>1000(1+1/\epsilon )r_2`$ will suffice). Notice that as $`n\mathrm{}`$, we have that $`B_R(p_n)\stackrel{\mathrm{GH}}{}B_R^q`$. Choose a finite $`R/1000`$-net $`\{a_i\}`$ of $`B_R^q`$. Let $`a_{i,n}M_n`$ be sequences such that $`a_{i,n}a_n`$. Consider the sequence of functions $`f_n:M_n`$ with $`f_n=\mathrm{min}_i\{\text{dist}_{a_{i,n}}^2\}`$.
For large $`n`$, the functions $`f_n`$ are $`2`$-concave in $`B_R(p_n)`$, so that, in particular, the gradient flows $`\mathrm{\Phi }_{f_n}^T|_{B_{r_2}(p_n)}`$ are $`e^{2T}`$-Lipschitz. Moreover, if $`\xi _x`$ denotes the starting vector of a unit speed shortest geodesic from $`x`$ to $`p_n`$, then for any $`xB_{r_2}(p_n)\backslash B_{r_1}(p_n)`$ we have $`\xi _x,fR/2`$. Therefore, if $`T=2r_2/R`$, then $`\mathrm{\Phi }_{f_n}^T(B_{r_2}(p_n))B_{r_1}(p_n)`$. Thus $`\mathrm{\Phi }_n=\mathrm{\Phi }_{f_n}^{2r_2/R}`$ provides a $`4r_2/R`$-Lipschitz map $`B_{r_2}(p_n)B_{r_1}(p_n)`$, and it is $`(1+\epsilon )`$-Lipschitz if one chooses $`R`$ sufficiently large. ∎
For $`\gamma \pi _1M`$, set $`|\gamma |_n=|p\gamma (p)|_{\stackrel{~}{M}_n}`$, see 2.6.
###### Corollary 3.1.2.
Let $`M`$ be almost nonnegatively curved manifold. Let
$$h:\pi _1MAut(H^{}(\stackrel{~}{M},)/tor$$
be the natural action of $`\pi _1M`$ on $`H^{}(\stackrel{~}{M},)`$. Then there is a sequence of norms $`||||_n`$ on $`H^{}(\stackrel{~}{M},)/tor`$ such that the following holds. Given any $`\epsilon >0`$, there is $`n_+`$ such that for any $`\gamma \pi _1M`$ with $`|\gamma |_n2\mathrm{diam}(M_n)`$ we have $`h(\gamma )_n1+\epsilon `$.
###### Proof.
\[FY92, theorem 0.1\] and Yamaguchi’s fibration theorem \[Yam91\] imply that if $`n`$ is sufficiently large, for any fixed $`r_+`$ we have that for any $`p_n\stackrel{~}{M}_n`$ the inclusion map $`B_r(p_n)\stackrel{~}{M}_n`$ is a homotopy equivalence.
Let $`||||_{n,r}`$ denote the $`L_{\mathrm{}}`$-norm on differential forms on $`B_r(p_n)\stackrel{~}{M}_n`$.
Fix $`r_2>r_1>0`$. If $`\omega `$ is a differential form on $`B_{r_1}(p_n)M_n`$ and $`n`$ is sufficiently large, Lemma 3.1.1 implies that
$$\mathrm{\Phi }_n^{}(\omega )_{n,r_2}(1+\epsilon )\omega _{n,r_1}\text{and}2\mathrm{diam}(M_n)r_2r_1.$$
If now $`\omega `$ is a form on $`B_{r_2}(p_n)\stackrel{~}{M}_n`$ and $`\gamma \pi _1M`$ such that
$$|\gamma |_n=|p_n\gamma (p_n)|2\mathrm{diam}(M_n)r_2r_1,$$
then $`B_{r_1}(p_n)B_{r_2}(\gamma (p_n))\stackrel{~}{M}_n`$, whence
$$\mathrm{\Phi }_n^{}(\gamma ^{}(\omega ))_{n,r_2}(1+\epsilon )\gamma ^{}(\omega )_{n,r_1}(1+\epsilon )\omega _{n,r_2}.$$
Thus, for the induced norms on the de Rham cohomology of $`\stackrel{~}{M}`$ (and on its integral subspace $`H^{}(\stackrel{~}{M},)/tor)`$) we have
$$[\gamma ^{}(\omega )]_{n,r_2}(1+\epsilon )[\omega ]_{n,r_2}.$$
Therefore the sequence of norms $`||||_n=||||_{n,r_2}`$ satisfies the conditions of the Corollary. ∎
###### Lemma 3.1.3.
There exists a constant $`N=N(n,k)_+`$ such that the following holds. If $`G`$ is a subgroup of $`\mathrm{GL}(n,)`$ and $`S`$ is a set of generators of $`G`$ with $`\mathrm{\#}(S)k`$ such that the eigenvalues of each element of $`S^N`$ are all equal to $`1`$ in absolute value, then the same is true for the eigenvalues of all elements of $`G`$.
###### Proof.
Let $`B`$ be the set of all matrices in $`\mathrm{GL}(n,)`$ for which all of their eigenvalues are equal to 1 in absolute value. Since the characteristic polynomials of such matrices are uniformly bounded and have integer coefficients, there are only finitely many of them. Let $`\overline{B}`$ be the Zariski closure of $`B`$ in the set of all real $`n\times n`$ matrices. By the above, all elements of $`\overline{B}`$ satisfy that the absolute values of all of their eigenvalues are equal to 1.
Consider now the space $`V=^{kn^2}`$ of $`k`$-tuples of real $`n\times n`$ matrices.
Consider a collection of matrices $`(M_1,M_2,\mathrm{},M_k)V`$, where $`M_i\mathrm{GL}(n,)`$. Let $`F_k`$ be a free group on $`k`$ generators, generated by $`S=\{\gamma _1,\gamma _2,\mathrm{},\gamma _k\}`$, and let $`h:F_k\mathrm{GL}(n,)`$ be the homomorphism defined by $`h(\gamma _i)=M_i`$. The property that for any $`\gamma F_k`$ $`h(\gamma )`$ be an element of $`\overline{B}`$ then describes an algebraic subset $`A_\gamma V`$.
The intersection $`A=_{\gamma F_k}A_\gamma `$ is also algebraic, and therefore there is a finite number $`N=N(n,k)`$ such that for $`S^NF_k`$, $`A=_{\gamma S^N}A_\gamma `$. ∎
###### Lemma 3.1.4.
Let $`\mathrm{\Gamma }`$ be a subgroup of $`\mathrm{GL}(n,)`$ such that the eigenvalues of each element of $`\mathrm{\Gamma }`$ are equal to $`1`$ in absolute value. Then $`\mathrm{\Gamma }`$ contains a subgroup $`\mathrm{\Gamma }^{}`$ of finite index all of whose elements have eigenvalues equal to $`1`$.
###### Proof.
Let $`G`$ denote the Zariski closure of $`\mathrm{\Gamma }`$ in $`\mathrm{GL}(n,)`$ . Then $`G`$, being an algebraic group, is a Lie group with finitely many components. Let $`G_{}`$ be the identity component of $`G`$. By the same argument as in the proof of the previous lemma, the set of all characteristic polynomials of the elements of $`G`$ is finite. Therefore the characteristic polynomial of any element of $`G_{}`$ is identically equal to $`(x1)^n`$.
Therefore, the subgroup $`\mathrm{\Gamma }^{}=\mathrm{\Gamma }G_{}`$ satisfies all conditions of the Lemma. ∎
###### Remark 3.1.5.
As was pointed out to us by Yu. Zarkhin, one can alternatively take $`\mathrm{\Gamma }^{}`$ to be the kernel of the composition of the homomorphisms $`\mathrm{\Gamma }\mathrm{GL}(n,)\mathrm{GL}(n,/3)`$. In this way one obtains a bound
$$[\mathrm{\Gamma }:\mathrm{\Gamma }^{}]3^{n^2}.$$
To see that $`\mathrm{\Gamma }^{}`$ satisfies the conclusion of Lemma 3.1.4, one should notice that every element of $`\mathrm{\Gamma }`$ is a quasi-unipotent matrix since all its eigenvalues are roots of unity. The desired result then follows from the so-called Minkowski Lemma. Apply, for instance, \[SZ96, Th. 7.2\] for $`n=3,k=1`$ (so $`R(1,3)=1`$), where we take $`𝒪`$ to be the ring of $`n\times n`$ integer matrices.
### 3.2. Proof of Theorem A
Let $`M`$ be an almost nonnegatively curved manifold. Denote, as usual, by $`M_n=(M,g_n)`$, $`n`$, a sequence of Riemannian metrics on $`M`$ such that $`\mathrm{sec}(M_n)1/n`$ and $`\mathrm{diam}(M_n)1/n`$, by $`\stackrel{~}{M}`$ the universal covering of $`M`$, and by $`\stackrel{~}{M}_nM_n`$ the universal Riemannian covering of $`M_n`$.
After passing to a finite cover of $`M`$, by \[FY92\] we may assume that $`\pi _1M`$ is nilpotent.
Fix $`pM`$ and let $`\{\gamma _{i,n}\}`$ be a short basis of $`\pi _1(M_n,p)`$ (see 2.6). Then, if $`n`$ is sufficiently large, the short basis $`\{\gamma _{i,n}\}`$ has at most $`k=k(dimM)`$ elements and its elements satisfy $`|\gamma _{i,n}|_n2/n`$ for every $`i`$. Moreover, Corollary 3.1.2 implies that given $`\epsilon >0`$, for all large $`n`$ and every $`i`$ we have $`h(\gamma _{i,n})_n<1+\epsilon `$ and $`h(\gamma _{i,n}^1)_n<1+\epsilon `$.
Take $`N=N(k,m)`$ as in Lemma 3.1.3. One can choose $`\epsilon >0`$ so small that if $`p`$ is a polynomial with integer coefficients for which all of its roots have absolute values lying between $`1/(1+\epsilon )^N`$ and $`(1+\epsilon )^N`$, then all roots of $`p`$ have absolute values equal to $`1`$. This follows from the fact that the total number of integer polynomials all of whose roots are contained in a fixed bounded region is finite.
Set $`S_n:=\{\gamma _{i,n}\}`$. Then for any $`\gamma S_n^N`$ we have $`h(\gamma )_n<(1+\epsilon )^N`$ and $`h(\gamma ^1)_n<(1+\epsilon )^N`$. Therefore the absolute values of all eigenvalues lie between $`1/(1+\epsilon )^N`$ and $`(1+\epsilon )^N`$. Since the characteristic polynomial of $`h(\gamma )`$ has integer coefficients, the absolute values of all the eigenvalues of $`h(\gamma )`$ are in fact equal to $`1`$.
Apply now Lemma 3.1.3. It follows that for any $`\gamma \pi _1M`$ the absolute values of all eigenvalues of $`h(\gamma )`$ are equal to $`1`$.
Then Lemma 3.1.4 implies that after passing to a finite cover $`M^{}`$ of $`M`$, for any $`\gamma \pi _1M^{}`$ all eigenvalues of $`h(\gamma )`$ are equal to $`1`$. By Engel’s theorem, one can choose an integral basis of $`H^{}(\stackrel{~}{M},)`$ such that the action of $`\pi _1M`$ on $`H^{}(\stackrel{~}{M},)/tor`$ is given by upper triangular matrices.
Therefore, by passing to a finite cover $`M^{\prime \prime }`$ of $`M^{}`$, we can assume that the action of $`\pi _1M^{\prime \prime }`$ on $`H^{}(\stackrel{~}{M},)`$ (and on $`H_{}(\stackrel{~}{M},)`$) is nilpotent.
Recall (see, e.g., \[HMR75, 2.19\]) that a connected CW complex with nilpotent fundamental group is nilpotent if and only if the action of its fundamental group on the homology of its universal cover is nilpotent. Thus $`M^{\prime \prime }`$ is a nilpotent space, whence the proof of Theorem $`A`$ is complete.∎
## 4. C-nilpotency of the fundamental group
### 4.1.
In this section we will prove Theorem B. It will follow from the following somewhat stronger result.
###### Theorem 4.1.1.
For any integer $`m`$ there exist constants $`\epsilon (m)>0`$ and $`C(m)>0`$ such that the following holds. If $`M^m`$ is a closed smooth $`m`$-manifold which admits a Riemannian metric $`g`$ with $`\mathrm{sec}(M^m,g)>\epsilon (m)`$ and $`\mathrm{diam}(M^m,g)<1`$, then the fundamental group of $`M^m`$ is $`C(m)`$-nilpotent; that is, $`\pi _1M^m`$ contains a nilpotent subgroup of index $`C(m)`$.
###### Remark 4.1.2.
The proofs of Theorems A and C show that corresponding versions of those results also do hold when these theorems are reformulated in a fashion similar to Theorem 4.1.1.
By an argument by contradiction, Theorem 4.1.1 follows from the following statement:
Given a sequence of Riemannian $`m`$-manifolds $`(M_n,g_n)`$ such that
$$\mathrm{diam}(M_n,g_n)1/n\text{and}\mathrm{sec}(g_n)1/n,$$
for each $`n`$, one can find $`C`$ such that $`\pi _1M_n`$ is $`C`$-nilpotent for all sufficiently large $`n`$.
### 4.2. Algebraic lemmas
Recall that the group of outer automorphisms $`\mathrm{Out}(G)`$ of a group $`G`$ is defined as the quotient of its automorphism group $`\mathrm{Aut}(G)`$ by the subgroup of inner automorphisms $`\mathrm{Inn}(G)`$.
###### Lemma 4.2.1 (A characterization of $`C`$-nilpotent groups).
Let
$$\{1\}=G_{\mathrm{}}\mathrm{}G_1G_0=G$$
be a sequence of groups satisfying the following properties:
For any $`i`$ we have that
1. $`G_iG`$ is normal in $`G`$;
2. the image of the conjugation homomorphism $`h_i:G\mathrm{Out}(G_i/G_{i+1})`$ is finite of order at most $`C_i`$;
3. $`G_i/G_{i+1}`$ contains an abelian subgroup $`E_i`$ of index $`c_i`$.
Then $`G`$ contains a nilpotent subgroup $`N`$ of index at most
$$C=C(c_1,\mathrm{},c_{\mathrm{}},C_1,\mathrm{},C_{\mathrm{}})$$
where $`N`$ is of nilpotency class $`\mathrm{}`$.
###### Proof.
First of all, notice that property (i) assures that the objects described in parts (ii) and (iii) of the lemma are well-defined.
Set $`\mathrm{\Gamma }_i:=G_i/G_{i+1}`$.
Let $`H=\mathrm{ker}h_i`$. Notice that $`[G:H]_iC_i`$ and that the image of $`H`$ under the conjugation homomorphism $`f_i:G\mathrm{Aut}(\mathrm{\Gamma }_i)`$ lies in $`\mathrm{Inn}(\mathrm{\Gamma }_i)`$; that is, $`f_i|_H:H\mathrm{Inn}(\mathrm{\Gamma }_i)`$.
By passing to a subgroup we can assume that $`E_i\mathrm{\Gamma }_i`$ is normal of index $`C(c_i)`$ (we can take $`C(c_i)=c_i!`$).
By increasing $`E_i`$ if necessary we can assume that $`E_i`$ contains the center of $`\mathrm{\Gamma }_i`$.
Let $`Z_i`$ be the image of $`E_i`$ under the projection map $`\pi :\mathrm{\Gamma }_i\mathrm{Inn}(\mathrm{\Gamma }_i)`$. Clearly $`[\mathrm{Inn}(\mathrm{\Gamma }_i):Z_i]c_i`$ and $`Z_i\mathrm{Inn}(\mathrm{\Gamma }_i)`$ is normal.
Let
$$N=H\left(\underset{i}{}f_i^1(Z_i)\right),N_i=NG_i.$$
Then
$$[G:N]C=C(c_1,\mathrm{},c_{\mathrm{}},C_1,\mathrm{},C_{\mathrm{}})$$
and $`N`$ satisfies:
For any $`i`$
1. $`N_iN`$ is normal in $`N`$;
2. $`N_i/N_{i+1}`$ is in the center of $`N/N_{i+1}`$;
that is, $`N`$ is nilpotent of nilpotency length $`\mathrm{}`$.
Condition (i) is obvious so we only need to check (ii).
To see (ii) observe that by construction the image of the conjugation action $`N\mathrm{Aut}(\mathrm{\Gamma }_i)`$ lies in $`\mathrm{Inn}(\mathrm{\Gamma }_i)`$. Moreover, in fact it lies in $`\pi (A_i)`$ and as such it acts trivially on $`E_i`$. Lastly observe that $`N_i/N_{i+1}E_i`$.
Indeed, by construction, for any $`gN_i/N_{i+1}\mathrm{\Gamma }_i`$ there is $`aE_i`$ such that $`\pi (g)=\pi (a)`$. Therefore, $`g=az`$ for some $`z`$ in the center of $`\mathrm{\Gamma }_i`$. By our assumption on $`E_i`$ this means that $`gE_i`$.
Thus $`N`$ acts trivially on $`N_i/N_{i+1}`$ which means that $`N`$ is nilpotent and $`G`$ is $`C`$-nilpotent. ∎
###### Trivial Lemma 4.2.2 (A characterization of finite actions).
If $`S`$ is a finite set of generators of a group $`G`$ with $`S^1=S`$, and $`h:GH`$ is a homomorphism with $`|h(S^n)|<n`$ for some $`n>0`$, then $`h(S^n)=h(G)`$ and, in particular, $`|h(G)|<n`$.
Let now $`\mathrm{\Gamma }`$ be a group which acts discretely by isometries on an Alexandrov space $`A`$ with curvature $`1`$. Choose a marked point $`pA`$. Assume that $`\{\gamma _1,\gamma _2,\mathrm{},\gamma _n\}`$ is a finite short basis of $`\mathrm{\Gamma }`$ at $`p`$ (see 2.6), and that $`\theta |\gamma _i|1`$, where $`|\gamma |\stackrel{def}{=}|\mathrm{p}\gamma (\mathrm{p})|`$. Let $`\mathrm{\#}(R)`$ denote the number of elements $`\gamma \mathrm{\Gamma }`$ with $`|\gamma |R`$. The Bishop–Gromov inequality implies that
$$\mathrm{\#}(R)v_1^m(R)/v_1^m(\theta ),$$
where $`m=dimA`$ and $`v_1^m(r)`$ is the volume of the ball of radius $`r`$ in the $`m`$-dimensional simply connected space form of curvature $`1`$. Therefore, if $`\mathrm{\#}\mathrm{\#}(L)`$ denotes the number of homomorphisms $`h:\mathrm{\Gamma }\mathrm{\Gamma }`$ with norm $`L`$ (that is, the number of homomorphisms for which it holds that for any $`\gamma \mathrm{\Gamma }`$ one has that $`|h(\gamma )|L|\gamma |`$), then
(4.2.1)
$$\mathrm{\#}\mathrm{\#}(L)\mathrm{\#}(L)^n\left[\frac{v_1^m(L)}{v_1^m(\theta )}\right]^n.$$
### 4.3. The blow-up construction
For $`n\mathrm{}`$, the manifolds $`M_n`$ clearly converge to a point $`=:A_0`$.
Set $`M_{n,1}:=M_n`$ and $`\vartheta _{n,1}:=\mathrm{diam}M_{n,1}`$.
Rescale now $`M_{n,1}`$ by $`\frac{1}{\vartheta _{n,1}}`$ so that $`\mathrm{diam}(\frac{1}{\vartheta _{n,1}}M_{n,1})=1`$. Passing to a subsequence if necessary, one has that the manifolds $`\frac{1}{\vartheta _{n,1}}M_{n,1}`$ converge to $`A_1`$, where $`A_1`$ is a compact nonnegatively curved Alexandrov space with diameter $`1`$.
Now choose a regular point $`\overline{p}_1A_1`$, and consider distance coordinates around $`\overline{p}_1U_1^{k_1}`$, where $`k_1`$ is the dimension of $`A_1`$. The distance functions can be lifted to $`U_{n,1}\frac{1}{\vartheta _{n,1}}M_{n,1}`$.
Let $`M_{n,2}`$ be the level set of $`U_{n,1}^{k_1}`$ that corresponds to $`\overline{p}_1`$. Clearly, $`M_{n,2}`$ is a compact submanifold of codimension $`k_1`$. Set $`\vartheta _{n,2}:=\mathrm{diam}M_{n,2}`$.
Passing again to a subsequence if necessary, one has that the sequence $`\frac{1}{\vartheta _{n,2}}M_{n,2}`$ converges to some Alexandrov space $`A_2`$. As before, $`A_2`$ is a compact nonnegatively curved Alexandrov space with diameter 1. Set $`k_2:=k_1+dimA_2`$. If one now chooses a marked point in $`M_{n,2}`$, then, as $`n\mathrm{}`$, $`M_n/\vartheta _{n,2}`$ converges to $`A_2\times ^{k_1}`$, which is of some dimension $`k_2>k_1`$.
We repeat this procedure until, at some step, $`k_{\mathrm{}}=m`$.
As a result one obtains a sequence $`\{A_i\}`$ of compact nonnegatively curved Alexandrov spaces with diameter $`1`$ that satisfies
$$dimA_i=k_ik_{i1},\text{so that}\underset{i=1}{\overset{\mathrm{}}{}}dimA_i=m.$$
We also obtain a sequence of rescaling factors $`\vartheta _{n,i}=\mathrm{diam}M_{n,i}`$, and a nested sequence of submanifolds
$$\{p_n\}=M_{n,\mathrm{}}\mathrm{}M_{n,2}M_{n,1}=M_n,$$
which in turn induces a sequence of homomorphisms
$$\{1\}=\pi _1M_{n,\mathrm{}}\stackrel{ı}{}\mathrm{}\stackrel{ı}{}\pi _1M_{n,2}\stackrel{ı}{}\pi _1M_{n,1}=\pi _1M_n.$$
Let $`G_i:=G_i(n):=ı^i[\pi _1M_{n,i}]`$.
For $`n`$ sufficiently large, the subgroups $`G_i(n)`$ are those which are generated by elements of norm $`3\vartheta _{n,i}`$. Equivalently, if one takes a short basis $`\{\gamma _i\}`$ of $`G(n)`$, then $`G_i`$ is the subgroup generated by all elements $`\gamma _i`$ such that $`|\gamma _i|3\vartheta _{n,i}`$.
### 4.4. Limit fundamental groups of Alexandrov spaces.
We will now define the “limit” or “L-fundamental groups” of the Alexandrov spaces $`A_i`$ constructed above. This notion is similar to the notion of the fundamental group of an orbifold. However, we note in advance that the construction of the L-fundamental group does not only depend on the spaces $`A_i`$, but also on the chosen rescaled subsequence of $`M_n`$. In fact, the following construction shows that the limit fundamental group of $`A_i`$, $`\pi _1^L(A_i)`$, is isomorphic to $`\pi _1(M_{n,i},M_{n,i+1})`$ for all sufficiently large $`n`$. But, unlike $`\pi _1M_{n,i}`$, the groups $`\pi _1^LA_i`$ will not depend on $`n`$.
The limit fundamental groups of $`A_i`$. Consider the converging sequence
$$(M_n/\vartheta _{n,i},p_n)\stackrel{\mathrm{GH}}{}(A_i\times ^{k_{i1}},\overline{p}_i\times 0)$$
(here the interesting case is collapsing). Recall that $`\overline{p}_iA_i`$ is a regular point. Fix $`\epsilon >0`$ such that $`\text{dist}_{\overline{p}_i}`$ on $`A_i`$ does not have critical values in $`(0,2\epsilon )`$. Take a sequence $`R_n`$ which converges very slowly to infinity (here we will need $`R_n\vartheta _{n,i}/\vartheta _{n,i1}0`$ and $`R_n\mathrm{}`$).
Consider then a sequence of Riemannian coverings
$$\mathrm{\Pi }:(\stackrel{~}{B}_n,\stackrel{~}{p}_n)(B_{R_n}(p_n),p_n)$$
of $`B_{R_n}(p_n)M_n/\vartheta _{n,i}`$ with $`\pi _1(\stackrel{~}{B}_n,\stackrel{~}{p}_n)=\pi _1(B_\epsilon (p_n),p_n)`$, where $`B_\epsilon (p_n)M_n/\vartheta _{n,i}`$.
After passing to a subsequence if necessary, the sequence $`(\stackrel{~}{B}_n,\stackrel{~}{p}_n)`$ converges to a nonnegatively curved Alexandrov space $`\stackrel{~}{A}_i\times ^{k_{i1}}`$, where the space $`\stackrel{~}{A}_i`$ has the same dimension as $`A_i`$. Indeed, by construction it contains an isometric copy of $`B_\epsilon (p_{n,i})`$, and therefore
$$dim\stackrel{~}{A}_i+k_{i1}=dim\underset{i\mathrm{}}{lim}B_\epsilon (p_{n,i})=dimA_i+k_{i1}.$$
Let us show that for all sufficiently large $`n`$,
$$ı[\pi _1M_{n,i+1}]\mathrm{}\pi _1M_{n,i}.$$
Assume that $`\mathrm{\Pi }(\stackrel{~}{q}_n)=\stackrel{~}{p}_n`$ and that $`\stackrel{~}{q}_n\overline{q}_n\stackrel{~}{A}_i`$. Connect $`\overline{p}_n`$ and $`\overline{q}_n`$ by a geodesic which, by \[Pet98\], only passes regular points. Note that in a small neighborhood of this geodesic in $`M_n`$ we have two copies of $`M_{n,i+1}`$, near $`\stackrel{~}{p}_n`$ and $`\stackrel{~}{q}_n`$. Therefore, applying Yamaguchi’s Fibration Theorem in a small neighborhood of this geodesic, we can construct a diffeomorphism from $`M_{n,i+1}`$ to itself. This implies that for any loop $`\gamma `$ which after lifting connects $`\stackrel{~}{p}_n\stackrel{~}{q}_n`$, we have
$$\gamma ^1ı[\pi _1M_{n,i+1}]\gamma ı[\pi _1M_{n,i+1}];$$
that is,
$$ı[\pi _1M_{n,i+1}]\pi _1M_{n,i}$$
(for an alternative argument see also \[FY92\]).
This easily yields that $`A_i=\stackrel{~}{A}_i/\mathrm{\Gamma }_i`$, where $`\mathrm{\Gamma }_i`$ is a group of isometries which acts discretely on $`\stackrel{~}{A}_i`$. The group $`\mathrm{\Gamma }_i`$ is denoted by $`\pi _1^LA_i`$ (the *limit* or *L-fundamental group* of $`A_i`$). This group is clearly isomorphic to
$$\pi _1(M_{n,i},M_{n,i+1})=\pi _1M_{n,i}/ı[\pi _1M_{n,i+1}]$$
for all sufficiently large $`n`$, and the space $`\stackrel{~}{A}_i`$ will be called the *universal covering* of $`A_i`$.
Since $`\stackrel{~}{A}_i`$ is nonnegatively curved and $`A_i=\stackrel{~}{A}_i/\pi _1^LA_i`$ is compact, by Toponogov’s splitting theorem $`\stackrel{~}{A}_i`$ isometrically splits as $`\stackrel{~}{A}_i=K_i\times ^{s_i}`$, where $`K_i`$ is a compact Alexandrov space with curv$`0`$. Since $`\pi _1^LA_i`$ is a group of isometries that acts discretely on $`\stackrel{~}{A}_i`$, it follows that $`\pi _1^LA_i`$ is a virtually abelian group.
### 4.5. Final steps
Consider now the corresponding series
$$\{1\}=G_{\mathrm{}}(n)\mathrm{}G_1(n)G_0(n)=\pi _1M_n.$$
The theorem then follows from the following
###### Lemma 4.5.1.
For all sufficiently large $`n`$, the series
$$\{1\}=G_{\mathrm{}}(n)\mathrm{}G_1(n)G_0(n)$$
constructed above satisfies the assumptions of Lemma 4.2.1 for numbers $`C_i`$ and $`c_i`$ which do not depend on $`n`$.
We first prove the following
###### Sublemma 4.5.2.
Each subgroup $`G_i(n)`$ is normal in $`G(n)`$.
###### Proof.
We will show by reverse induction on $`k`$ that $`G_i(n)G_k(n)`$ for any $`ki`$. Let us assume that we already know that $`G_i(n)G_{k+1}(n)`$. Since
$$ı[\pi _1M_{n,k+1}]\pi _1M_{n,k},$$
we know that $`G_{k+1}(n)G_k(n)`$. Consider the covering
$$\mathrm{\Pi }_{k+1}:(\stackrel{~}{M}_{n,k+1},\stackrel{~}{p}_{n,k+1})(M_n,p_n)$$
with covering group $`\mathrm{\Gamma }_{k+1}(n)`$.
Clearly $`(\stackrel{~}{M}_{n,k+1},\stackrel{~}{p}_{n,k+1})\stackrel{\mathrm{GH}}{}^{s_i}`$ for some integer $`s_i`$. Applying Lemma 2.5.1, it follows that for any $`aG`$ with $`|a|<1`$ there is a cocos-curve $`\gamma `$ in $`\stackrel{~}{M}_{n,k+1}`$ with total time $`T`$ connecting $`\stackrel{~}{p}_n`$ and $`a(\stackrel{~}{p}_n)`$ in $`\stackrel{~}{M}_{n,k+1}`$. Then clearly $`\gamma ga`$ for some $`gG_{k+1}(n)`$. Let us denote by $`\mathrm{\Phi }^T:\stackrel{~}{M}_{n,i}\stackrel{~}{M}_{n,i}`$ the gradient flow corresponding to $`\gamma `$.
Let $`\gamma _j`$ be a loop from the short basis of $`G_i(n)`$. As was mentioned in 4.3, if $`n`$ is large, then $`\mathrm{length}\gamma _j3\vartheta _{n,i}`$. Let us denote by $`\stackrel{~}{\gamma }_j`$ a lift of $`\gamma _j`$ to $`\stackrel{~}{M}_{n,i}`$. Let $`\stackrel{~}{p}_{n,j}\stackrel{~}{M}_{n,i}`$ be its starting point. Since $`[\gamma _j]G_i(n)`$, we have that $`\stackrel{~}{\gamma }_j`$ is a loop in $`\stackrel{~}{M}_{n,i}`$. Consider then the loop $`\gamma _j^{}=\mathrm{\Pi }\mathrm{\Phi }^T\stackrel{~}{\gamma }_j`$. Clearly,
$$[\gamma _j]=a^1g^1[\gamma _j^{}]ga,\text{or}[\gamma _j^{}]=ga[\gamma _j]a^1g^1.$$
Now Proposition 2.3.3 implies that
$$\mathrm{length}(\gamma _j^{})\mathrm{exp}(2T)\mathrm{length}(\gamma _j).$$
Thus, for sufficiently large $`n`$,
$$ga[\gamma _j]a^1g^1G_i(n),$$
and since $`gG_i(n)G_{k+1}(n)`$ it follows that
$$a[\gamma _j]a^1G_i(n);$$
that is, $`G_i(n)G_k(n)`$.
###### Proof of Lemma 4.5.1.
The group
$$\pi _1^LA_i=\pi _1(M_{n,i},M_{n,i+1})=\pi _1M_{n,i}/ı[\pi _1M_{n,i}]$$
is virtually abelian. Let $`d_i`$ be the minimal index of an abelian subgroup of $`\pi _1^LA_i`$. The epimorphism $`ı^i:\pi _1M_{n,i}G_i`$ induces an epimorphism $`\pi _1^LA_iG_i(n)/G_{i+1}(n)`$. Therefore, $`G_i(n)/G_{i+1}(n)`$ is $`d_i`$-abelian for all large $`n`$.
Consider the covering $`\mathrm{\Pi }_i:\stackrel{~}{M}_{n,i}M_n`$ with covering group $`G_i(n)`$, and let $`\stackrel{~}{p}_{n,i}`$ be a preimage of $`p_n`$. Clearly $`(\stackrel{~}{M}_{n,i},\stackrel{~}{p}_{n,i})\stackrel{\mathrm{GH}}{}^{s_i}`$ for some integer $`s_i`$. Applying Lemma 2.5.1, it follows that for any $`aG(n)`$ with $`|a|<1`$ there is a cocos-curve $`\gamma `$ in $`\stackrel{~}{M}_{n,i}`$ which connects $`p`$ and $`a(p)`$. Then clearly $`\gamma ga`$ for some $`gG_i(n)`$. Let us denote by $`\mathrm{\Phi }^T:\stackrel{~}{M}_{n,i}\stackrel{~}{M}_{n,i}`$ the gradient flow corresponding to $`\gamma `$.
Let $`bG_i(n)`$ and $`\beta `$ be a loop representing $`b`$. Let us denote by $`\stackrel{~}{\beta }`$ a lift of $`\beta `$ to $`\stackrel{~}{M}_{n,i}`$. Let $`\stackrel{~}{p}_{n,i}\stackrel{~}{M}_{n,i}`$ be its starting point. Since $`[\beta ]G_i(n)`$, we have that $`\stackrel{~}{\beta }`$ is a loop in $`\stackrel{~}{M}_{n,i}`$.
Consider now the loop $`\beta ^{}=\mathrm{\Pi }\mathrm{\Phi }^T\stackrel{~}{\beta }`$. Clearly,
$$b=[\beta ]=a^1g^1[\beta ^{}]ga,\text{or}[\beta ^{}]=gaba^1g^1.$$
Proposition 2.3.3 then implies that
$$\mathrm{length}(\beta ^{})\mathrm{exp}(2T)\mathrm{length}(\beta )$$
Therefore, if $`h_a:G_i(n)/G_{i+1}(n)G_i(n)/G_{i+1}(n)`$ is induced by the conjugation $`baba^1`$, then for any $`aG(n)`$ there is $`gG_i(n)`$ such that $`|h_{ga}|\mathrm{exp}(2T)`$.
Let now $`\delta _i`$ be the minimal norm of the elements of $`\pi _1^LA_i`$, where $`\pi _1^LA_i`$ acts on $`\stackrel{~}{A}_i`$. Then (4.2.1) implies that the image of the action of $`G(n)`$ by conjugation in $`Out(G_i(n)/G_{i+1}(n))`$ is $`C_i`$-finite, where $`C_i`$ depends only on $`c_i`$, $`T`$, and $`\delta _i`$. ∎
### 4.6. Remark on nonfree actions
Theorem 4.1.1 can be reformulated as follows:
There exists a constant $`\epsilon (m)>0`$ such that if $`N^m`$ is a Riemannian manifold which admits a free discrete isometric action by a group $`\mathrm{\Gamma }`$ such that $`\mathrm{sec}(N)>\epsilon (m)`$ and $`\mathrm{diam}(N/\mathrm{\Gamma })<1`$, then $`\mathrm{\Gamma }`$ is $`C(m)`$-nilpotent.
As was pointed out to us by B. Wilking, in the above reformulation of Theorem 4.1.1 one can actually remove the assumption that the $`\mathrm{\Gamma }`$ action be free.
###### Corollary 4.6.1.
There exists a constant $`\epsilon (m)>0`$ such that if $`N^m`$ is a Riemannian manifold which admits a discrete isometric action by a group $`\mathrm{\Gamma }`$ such that $`\mathrm{sec}(N)>\epsilon (m)`$ and $`\mathrm{diam}(N/\mathrm{\Gamma })<1`$, then $`\mathrm{\Gamma }`$ is $`C(m)`$-nilpotent.
###### Proof.
Let $`\epsilon =\epsilon (m)`$ be as provided by Theorem 4.1.1 and suppose $`N`$ satisfies the assumptions of the corollary for this $`\epsilon `$. Let $`F`$ be the frame bundle of $`N`$. Then the action of $`\mathrm{\Gamma }`$ on $`N`$ lifts to a free isometric action on $`F`$. As was observed in \[FY92\], using Cheeger’s rescaling trick, $`F`$ can be equipped with a $`\mathrm{\Gamma }`$ invariant metric satisfying $`\mathrm{sec}(F)>\epsilon (m)`$ and $`\mathrm{diam}(F/\mathrm{\Gamma })<1`$. Since the induced action of $`\mathrm{\Gamma }`$ on $`F`$ is free, the claim of the corollary now follows from Theorem 4.1.1. ∎
## 5. Proof of the Fibration Theorem
### 5.1.
Let $`M`$ be an almost nonnegatively curved manifold. Denote by $`M_n=(M,g_n)`$ a sequence of Riemannian metrics on $`M`$ such that $`\mathrm{sec}(M_n)1/n`$ and $`\mathrm{diam}(M_n)1/n`$.
Let us denote by $`\stackrel{~}{M}`$ the universal cover of $`M`$ and by $`\stackrel{~}{M}_nM_n`$ the universal Riemannian covering of $`M_n`$ (that is, $`\stackrel{~}{M}`$ equipped with the pull back of the metric $`g_n`$ on $`M`$).
By \[FY92\], passing to a finite cover we may assume that $`\mathrm{\Gamma }=\pi _1M`$ is a nilpotent group without torsion. Hence, to prove the topological part of Theorem C, it is sufficient to show the following:
###### Theorem 5.1.1.
Let $`M`$ be a closed almost nonnegatively curved $`m`$-manifold such that $`\mathrm{\Gamma }=\pi _1M`$ is a nilpotent group without torsion. Then $`M`$ is the total space of a fiber bundle
$$FMN$$
where the base $`N`$ is a nilmanifold and the fiber $`F`$ is simply connected.
The assumption on $`\mathrm{\Gamma }`$ implies that we can fix a series
$$\mathrm{\Gamma }=\mathrm{\Gamma }_0\mathrm{\Gamma }_1\mathrm{\Gamma }_2\mathrm{}\mathrm{\Gamma }_{\mathrm{}}=\{1\}$$
such that $`\mathrm{\Gamma }_i`$ is normal in $`\mathrm{\Gamma }`$ and $`\mathrm{\Gamma }_i/\mathrm{\Gamma }_{i+1}`$.
Let us first us give an informal proof.
### 5.2. An informal proof of Theorem 5.1.1
We use induction to construct the bundles $`F_iM\stackrel{f_i}{}N_i`$, where each $`N_i`$ is a nilmanifold with $`\pi _1N_i=\mathrm{\Gamma }/\mathrm{\Gamma }_i`$ and $`\pi _1F_i\mathrm{\Gamma }_i`$. Since the base of induction is trivial, we are only interested in the induction step.
Fix $`pN_i`$, and let $`F_i(p)`$ be the fiber over $`p`$. For any sufficiently large $`n`$ choose a subgroup $`G_i=G_i(n)`$ such that $`\mathrm{\Gamma }_iG_i\mathrm{\Gamma }_{i+1}`$ and $`[\mathrm{\Gamma }_i:G_i]`$ is finite, but sufficiently large so that the cover $`\overline{F}_i(p)`$ of $`F_i(p)`$ corresponding to $`G_i`$ is Hausdorff close to a unit circle $`𝕊^1`$.
Construct now a bundle map $`\phi _p:\overline{F}_i(p)𝕊^1`$ by lifting distance functions from $`𝕊^1`$ (This can be done by a slight generalization of a construction in \[FY92\] and \[BGP92\]). Let $`\omega _p=d\phi _p`$.
Then $`\omega _p`$ is a closed integral non-degenerate one-form on $`F_i(p)`$. Since deck transformations are isometries, after averaging by $`_a`$, where $`a=[\mathrm{\Gamma }_i:G_i]`$, we can assume that $`\omega _p`$ is $`_a`$-invariant. Thus $`\omega _p`$ descends to a form on $`F_i(p)`$ which when integrated gives a bundle map $`F_i(p)`$ onto a small $`𝕊^1`$.
Note that altho this bundle is defined only up to rotations of $`𝕊^1`$, its fibers are well-defined.
Since $`\mathrm{\Gamma }_{i+1}`$ is normal in $`\mathrm{\Gamma }`$, the choice of the covering $`\overline{F}_i(p)`$ of $`F_i(p)`$ is unambiguous for all $`pN_i`$. By using a partition of unity on $`N_i`$ we can glue the forms $`\omega _p`$ into a global 1-form on $`M`$ which satisfies the following properties:
1. $`\omega |_{F(p)}`$ is closed and integral for any $`p`$;
2. $`\omega |_{F(p)}`$ is non-degenerate.
Integrating $`\omega `$ over the various $`F(p)`$’s we construct a continuous family of bundles $`F_p𝕊^1`$. The level sets partition each $`F(p)`$ and hence the whole $`M`$ into fibers of a fiber bundle, whose quotient space is then a circle bundle $`N_{i+1}`$ over $`N_i`$.∎
This gives a good idea of the proof. However, to make it precise some extra work has to be done. In particular, one has to be careful with the construction of $`\omega `$. To make this construction possible we have to keep track of how $`F(p)`$ was obtained. Namely, we have to use that the fiber $`F(p)`$ was obtained by a construction as in Yamaguchi’s fibration theorem (see \[Yam91\] or \[BGP92\]). This makes the induction proof quite technical.
We now proceed with the rigorous proof of Theorem 5.1.1.
### 5.3. Proof of Theorem 5.1.1
Let us denote by $`\stackrel{~}{M}_{n,i}`$ the Riemannian covering of $`M_n`$ with respect to $`\mathrm{\Gamma }_i`$.
For any choice of marked points $`p_n`$ we have that
$$(\stackrel{~}{M}_{n,i},p_n,\pi _1M)\stackrel{\mathrm{GH}}{}(^i,0,^i)$$
in equivariant Gromov–Hausdorff convergence, where $`^i`$ acts on itself by translations. Indeed, the limit space must be a nonnegatively curved simply connected Alexandrov space, and since $`\mathrm{diam}M_n0`$ we have that it possesses a transitive group action by a nilpotent group. Then Euclidean space, acting as a group of translations, is here the only choice, and it is easy to see that the dimension of the limit must be equal to $`i`$.
Therefore $`(\stackrel{~}{M}_n,p_n,\pi _1M)\stackrel{\mathrm{GH}}{}(^{\mathrm{}},0,^{\mathrm{}})`$, and we may also assume that for each $`i`$ we have that $`(\stackrel{~}{M}_n,p_n,\mathrm{\Gamma }_i)\stackrel{\mathrm{GH}}{}(^{\mathrm{}},0,^\mathrm{}i)`$, where $`^\mathrm{}i`$ is the coordinate subspace of $`^{\mathrm{}}`$ which corresponds to the first $`\mathrm{}i`$ elements of the standard basis.
Now, let us give a technical definition:
If $`R`$ is a Riemannian manifold, let us denote by $`\stackrel{~}{\mathrm{dist}}_p`$ the average of a distance function over a small ball around $`p`$. This enables us to work with the $`C^1`$ function $`\stackrel{~}{\mathrm{dist}}_p`$ instead of the Lipschitz function $`\mathrm{dist}_p`$.
###### Definition 5.3.1.
Let $`R_n\stackrel{\mathrm{GH}}{}R`$ be a sequence of Riemannian $`m`$-manifolds with curvature $`\kappa `$ which Gromov–Hausdorff converges to a Riemannian $`m^{}`$-manifold $`R`$, where $`m^{}m`$. A sequence of forms $`\omega _n`$ on $`R_n`$ is said to $`\epsilon `$-approximate a form $`\omega `$ on $`R`$, if
1. for any point $`pR`$ there is a neighborhood $`Up`$ which admits a distance chart $`f:U^m^{}`$,
$$f(x)=(\mathrm{dist}_{a_1}(x),\mathrm{dist}_{a_2}(x),\mathrm{},\mathrm{dist}_{a_m^{}}(x))$$
which is a smooth regular map, and
2. smooth lifts of $`f`$ to $`U_nR_n`$ give, for all large $`n`$, regular maps
$$f_n(x)=(\stackrel{~}{\mathrm{dist}}_{a_{1,n}}(x),\stackrel{~}{\mathrm{dist}}_{a_{2,n}}(x),\mathrm{},\stackrel{~}{\mathrm{dist}}_{a_{m^{},n}}(x))$$
with $`a_{i,n}M_n`$, $`a_{i,n}a_n`$ such that
$$|(f_nf^1)^{}(\omega )\omega _n|_{C^0}<\epsilon $$
for all sufficiently large $`n`$.
Theorem 5.1.1 now easily follows from the following lemma:
###### Lemma 5.3.2.
Given $`\epsilon >0`$ there is a sequence of one-forms
$$\{\omega _{1,n},\omega _{2,n},\mathrm{},\omega _{k,n}\}$$
on $`\stackrel{~}{M}_n`$ with the following properties:
1. For each $`i`$, $`\omega _{i,n}`$ is a $`\pi _1M`$-invariant form on $`\stackrel{~}{M}_n`$.
2. The forms $`\omega _{i,n}`$ $`\epsilon `$-approximate the coordinate forms $`dx_i`$ on $`^k`$. In particular, the forms $`\{\omega _{i,n}\}`$ are nowhere zero and almost orthonormal at each point.
3. If for any $`j<i`$ it holds that $`\omega _{j,n}(X)=\omega _{j,n}(Y)=0`$, then $`d\omega _{i,n}(X,Y)=0`$. In particular, for each $`i`$ and all sufficiently large $`n`$, the distribution corresponding to the system of equations
$$\omega _{j,n}(X)=0\text{for all}ji$$
defines on $`\stackrel{~}{M}_n`$ a foliation $`_{i,n}`$.
4. If $`\stackrel{~}{F}_{i,n}(x)\stackrel{~}{M}_n`$ denotes the fiber of the foliation $`_{i,n}`$ containing the point $`x\stackrel{~}{M}_n`$, then each $`\stackrel{~}{F}_{i,n}(x)`$ is $`\mathrm{\Gamma }_i`$-invariant; that is, for any $`\gamma \mathrm{\Gamma }_i`$ one has that $`\stackrel{~}{F}_{i,n}(x)=\stackrel{~}{F}_{i,n}(\gamma x)`$. Moreover, the action of $`\mathrm{\Gamma }_i`$ on $`\stackrel{~}{F}_{i,n}(x)`$ is cocompact for each $`i`$. In particular, $`_{i,n}`$ induces on $`M_n`$ the structure of a fiber bundle.
###### Proof.
We will construct these forms by induction. Assume that we have already constructed one-forms $`\omega _1,\omega _2,\mathrm{},\omega _{i1}`$ which meet all the required properties. They give a $`\pi _1M`$-invariant fibration of $`\stackrel{~}{M}_n`$ by submanifolds $`\stackrel{~}{F}_{i1,n}(x)`$ thru each point $`x\stackrel{~}{M}_n`$, with tangent spaces defined by the equations $`\omega _j(X)=0`$ for $`j=1,\mathrm{},i1`$.
Denote by $`\theta :[0,1]`$ a smooth monotone function which is equal to $`1`$ before $`0`$ and $`0`$ after $`1`$. Choose numbers $`\delta _n>0`$ slowly converging to $`0`$, and let $`\mathrm{\Theta }_{i,n}:\stackrel{~}{M}_n_+`$ be the function defined by
$$\mathrm{\Theta }_{i,n}(x)=\underset{yF_{i1,n}(x)}{\mathrm{min}}\{\theta (|p_ny|/\delta _n)\}.$$
Clearly $`\mathrm{\Theta }_{i,n}`$ is a continuous $`\mathrm{\Gamma }_{i1}`$-invariant function which is constant on each $`F_{i1,n}(x)`$. Moreover, for large $`n`$, $`\mathrm{\Theta }_{i,n}`$ has support in some $`C_i\delta _n`$-neighborhood of $`F_{i1,n}(p_n)`$, and is equal to $`1`$ in some $`c_i\delta _n`$-neighborhood of $`F_{i1,n}(p_n)`$.
Now let $`\phi :[0,1]`$ be a smooth nondecreasing function which is $`0`$ before $`1/2`$ and $`1`$ after $`3/2`$. Consider the form
$$\omega _{i,n}^{}=\mathrm{\Theta }_{i,n}d(\phi \stackrel{~}{\mathrm{dist}}_{\mathrm{\Gamma }_ia_{i,n}}),$$
where $`a_{i,n}\stackrel{~}{M}_n`$ is a sequence of points converging to $`e_i^{\mathrm{}}`$, and $`\stackrel{~}{\mathrm{dist}}_{\mathrm{\Gamma }_ia_{i,n}}`$ is the average of $`\mathrm{dist}_{\mathrm{\Gamma }_ix}`$ for $`x`$ in a small ball around $`a_{i,n}`$. The support of $`\omega _i^{}`$ has two components, one which contains $`p_n`$ (notice here that $`p_n0^{\mathrm{}}`$), and another which does not. (It follows from the construction that the limit of $`F_{i1,n}(p_n)`$ is a coordinate plane in $`^{\mathrm{}}`$).
Set $`\omega _{i,n}^{\prime \prime }:=\omega _{i,n}^{}`$ on the component of $`p_n`$, and let this form be $`0`$ otherwise. Clearly, $`\omega _{i,n}^{\prime \prime }`$ is then a continuous $`\mathrm{\Gamma }_i`$-invariant form whose restriction to $`\stackrel{~}{F}_{i1,n}(x)`$ is exact. Moreover, each level set of its integral over $`\stackrel{~}{F}_{i1,n}(x)`$ is $`\mathrm{\Gamma }_i`$-invariant.
By construction, the form $`\omega ^{\prime \prime }/|\omega ^{\prime \prime }|`$ is now (in the sense of definition 5.3.1) close to $`dx_i`$ at the points where $`|\omega ^{\prime \prime }|0`$. Take
$$\omega _{i,n}=c\underset{\gamma \mathrm{\Gamma }/\mathrm{\Gamma }_i}{}\gamma \omega ^{},$$
where the coefficient $`c`$ is chosen in such a way that $`|\omega _{i,n}(p_n)|=1`$. As $`\delta _n`$ is a sequence slowly converging to zero, we may assume that $`\mathrm{diam}(M_n)/\delta _n0`$. Therefore, $`\omega _{i,n}`$ is the form we need. ∎
Notice that the proof actually shows that the fibers in Theorem 5.1.1 are almost nonnegatively curved manifolds in the generalized sense with $`k=\mathrm{}`$. Therefore, the proof of Theorem C is complete.
## 6. Open questions
We would like to conclude this work by posing a number of related open questions.
### 6.1. Is the torsion contained in the center?
As was noted earlier, Theorem B is new even for manifolds of nonnegative curvature. For such manifolds it is known that their fundamental groups are almost abelian, and Fukaya and Yamaguchi conjectured the following (see \[FY92\]):
###### Conjecture 6.1.1 (Fukaya–Yamaguchi).
The fundamental group of a nonnegatively curved $`m`$-manifold is $`C(m)`$-abelian.
In this regard we would like to pose the following two conjectures:
###### Main Conjecture 6.1.2.
There is $`C=C(m)`$ such that if $`M^m`$ is almost nonnegatively curved then there is a nilpotent subgroup $`N\pi _1M`$ of index $`C`$ whose torsion is contained in its center (or, at least, whose torsion is commutative).
###### Conjecture 6.1.3.
If $`M^m`$ is almost nonnegatively curved, then the action of $`\pi _1M`$ on $`\pi _2M`$ is almost trivial, or maybe even $`C(m)`$-trivial; that is, there exists a finite index subgroup of $`\pi _1M`$ (or, respectively, a subgroup of index $`C(m)`$) which acts trivially on $`\pi _2M`$.
Conjecture 6.1.2 implies in particular that the fundamental groups of closed positively curved $`m`$-manifolds are $`C(m)`$-abelian.
In fact, as was pointed out to us by B. Wilking, if true, Conjecture 6.1.2 would also imply a positive answer to Conjecture 6.1.1. Indeed, if $`\mathrm{sec}(M)0`$, then the universal cover $`\stackrel{~}{M}`$ of $`M`$ is isometric to the product $`^n\times K`$, where $`K`$ is a compact Riemannian manifold and the $`\pi _1M`$ action on $`^n\times K`$ is diagonal. It follows from \[Wil00, Cor. 6.3\] that one can deform the metric on $`M`$ so that its universal cover is still isometric to $`^n\times K`$ and the induced action on $`K`$ is finite. By passing, as in the proof of Corollary 4.6.1, to the induced action on the frame bundle of $`K`$, one reduces the statement to Conjecture 6.1.2.
We tried to prove these conjectures by studying successive blow-ups of the collapsing sequence $`M_n`$ as done in Section 4.3.
We can prove Conjectures 6.1.2 and 6.1.3 in the case where all spaces $`A_i`$ which appear in the construction in Section 4.3 are closed Riemannian manifolds; see \[KPT\]. Moreover, we believe we have an argument to prove it if all $`A_i`$’s are Alexandrov spaces without boundary.
It seems that if we would have just a slightly better understanding of collapsing to a ray, then we could prove the conjectures. Here is the simplest related question which we cannot solve:
###### Question 6.1.4.
Let $`M_n=(𝕊^2\times ^2,g_n)`$ be a sequence of complete Riemannian manifolds with $`\mathrm{sec}(M_n)\epsilon _n`$, where $`\epsilon _n0`$ as $`n\mathrm{}`$. Assume that for a sequence of points $`p_nM_n`$ we have that $`(M_n,p_n)\stackrel{\mathrm{GH}}{}(_+,0)`$. Let $`q_nM_n`$ be a sequence of points such that $`|p_nq_n|=1`$ and such that there is a sequence of rescalings $`\lambda _n\mathrm{}`$ such that $`(\lambda _nM_n,q_n)\stackrel{\mathrm{GH}}{}𝕊^2\times 𝕊^1\times `$, where the latter space is equipped with the product of the canonical metrics.
1. Can it happen that $`(\lambda _nM_n,p_n)\stackrel{\mathrm{GH}}{}(_+,0)`$?
2. Is it true that the dimension of the Gromov–Hausdorff limit of $`(\lambda _nM_n,p_n)`$ is at least 3?
3. What are the possible limits of $`(\lambda _nM_n,p_n)`$?
Conjectures 6.1.1 and 6.1.2 are also related to the following conjecture of Rong (cf. \[Ron96b, Ron96a\]):
###### Conjecture 6.1.5 (Rong).
Positively curved $`m`$-manifolds have $`C(m)`$-cyclic fundamental groups.
This conjecture has been proved by Rong \[Ron96b\] under the additional assumption of a uniform upper curvature bound. We also believe that if one could carry out the above program for proving Conjecture 6.1.2, one would have a good shot at handling Rong’s Conjecture as well.
### 6.2. The simply connected case
So far we have only discussed manifolds with nontrivial fundamental groups. However, some of our arguments also work in a more general setting. We hope that it might be possible to use them to obtain new restrictions on simply connected almost nonnegatively curved manifolds as well as on collapsing with a lower curvature bound.
Let us indicate one possible approach to do so.
Let us denote by $`(F)`$ the space of self homotopy equivalences of a manifold $`F`$. Assume now that $`F`$ is simply connected and that $`\stackrel{~}{f}:𝕊^k\times FF`$ is a map such that $`\stackrel{~}{f}_u:FF`$ is homotopic to the identity for some (and therefore ANY) $`u𝕊^k`$. Then $`\stackrel{~}{f}`$ represents an element $`\alpha =[\stackrel{~}{f}]\pi _k((F))`$. Let $`g`$ be a Riemannian metric on $`F`$. Define
$$\mathrm{dil}_g(\stackrel{~}{f})=\underset{u𝕊^k}{\mathrm{max}}\mathrm{dil}_g(\stackrel{~}{f}_u),$$
where $`\mathrm{dil}_g(\stackrel{~}{f}_u)`$ stands for the optimal Lipschitz constant of $`\stackrel{~}{f}_u`$ with respect to $`g`$. For any $`\alpha \pi _k((F))`$ define
$$\mathrm{dil}_g(\alpha )=\underset{[h]=\alpha }{inf}\mathrm{dil}_g(h).$$
Finally, define
$$\mathrm{DIL}(\alpha )=\underset{g}{inf}\mathrm{dil}_g(\alpha )$$
over all Riemannian metrics $`g`$ on $`F`$ and
$$\mathrm{DIL}_+(\alpha )=\underset{g}{inf}\mathrm{dil}_g(\alpha )$$
over all Riemannian metrics $`g`$ on $`F`$ with $`\mathrm{diam}(F,g)1`$ and $`\mathrm{sec}(g)1`$.
Clearly, both $`\mathrm{DIL}(\alpha )`$ and $`\mathrm{DIL}_+(\alpha )`$ are homotopy invariants of $`\alpha `$.
Now suppose that $`M_n\stackrel{\mathrm{GH}}{}𝕊^{k+1}`$ is a sequence of Riemannian manifolds collapsing to a round sphere with $`\mathrm{sec}(M_n)k`$. By Yamaguchi’s fibration theorem, we have that $`M_n`$ is a fiber bundle over $`𝕊^{k+1}`$ with almost nonnegatively curved fiber $`F_n`$. This bundle is classified by an element $`\alpha `$ of $`\pi _k(Aut(F_n))`$ and by using our gradient flow technique we can estimate $`\mathrm{DIL}_+(\alpha )`$ (and hence $`\mathrm{DIL}(\alpha )`$) from above.
Therefore, if one could find examples of a simply connected $`F`$ and an $`\alpha `$ with arbitrary big $`\mathrm{DIL}_+(\alpha )`$, one would obtain new restrictions on collapsing to a sphere with curvature bounded from below, and probably more restrictions for the topological type of manifolds with lower curvature and upper diameter bounds in general. In fact, $`F`$ need not be simply connected as long as the total space of the bundle $`FM𝕊^{k+1}`$ is.
While we believe that finding examples with arbitrary large $`\mathrm{DIL}(\alpha )`$) is very difficult (and might even be impossible), we have several candidates to produce large $`\mathrm{DIL}_+(\alpha )`$.
On the other hand, the problem of finding examples of $`\alpha `$ with $`\mathrm{DIL}(\alpha )>1`$ seems quite interesting in its own right and might have other applications unrelated to collapsing.
Let us next describe some possible sources of examples with $`\mathrm{DIL}_+(\alpha )>1`$:
###### Example 6.2.1.
Obviously, if dil$`{}_{g}{}^{}(h)=1`$, then $`h_u`$ is a homotopy of isometries of $`(F,g)`$. Let $`G`$ be the isometry group of $`F`$. Then $`G`$ can be viewed as a subset of $`(F)`$. Therefore, if $`[h]0`$ in $`\pi _k((F))`$, then $`[h_u]0`$ in $`\pi _kG`$. Now $`G`$ is a compact Lie group, in particular, $`\pi _2G=0`$ (and even more generally $`\pi _{2n}G`$ is finite). On the other hand, there are spaces $`F`$ such that the space $`(F)`$ might have nontrivial second homotopy; for example, the canonical metric on $`F=\mathrm{SU}(6)/(\mathrm{SU}(3)\times \mathrm{SU}(3))`$ has nonnegative curvature, and it follows from \[OT97, Chapter 5, Example 4.14\], that $`\pi _2((F))`$ is nontrivial. Therefore, there is an $`\alpha \pi _2((F))`$ such that dil$`{}_{g}{}^{}(\alpha )>1`$ for any metric $`g`$ on $`F`$; we believe it should be true that $`\mathrm{DIL}_+(\alpha )>1`$. Still, it might happen that $`\mathrm{DIL}(\alpha )=1`$.
Another possible source of such manifolds is provided by the following example due to D. Sullivan.
###### Example 6.2.2.
Let $`N^7`$ be the total space of an $`𝕊^3`$ bundle over $`𝕊^4`$ with zero Euler class and nontrivial $`p_1`$. Clearly $`N^7`$ is rationally equivalent to $`𝕊^4\times 𝕊^3`$. In particular, its minimal model has no nontrivial derivations of degree $`1`$. Therefore, by \[Sul77, 13.3\], there exists a diffeomorphism $`f:NN`$ which is homotopic to the identity but such that the obstruction to it being diffeotopic to the identity is a nonzero element of $`H^3(N,)`$. Let $`M^8`$ be the mapping cylinder of $`f`$. Clearly $`M`$ is homotopy equivalent to $`N\times 𝕊^1`$ and hence it is spin with signature zero. On the other hand, by construction, $`p_1^2(M)0`$. Since the signature of $`M`$ is zero we must necessarily have that $`p_2(M)0`$ and hence $`\widehat{A}(M)0`$. In particular, by the Atiyah–Hirzebruch theorem, $`M`$ does not admit an $`𝕊^1`$ action and hence the corresponding element $`\alpha \pi _1((M))`$ has dil$`{}_{g}{}^{}(\alpha )>1`$ for any metric $`g`$ on $`M`$.
###### Remark 6.2.3.
As was mentioned in the introduction, it is known that a spin manifold $`X`$ of almost nonnegative Ricci curvature has $`\widehat{A}(X)2^{dimX/2}`$ (\[Gro82, page 41\], \[Gal83\]). Clearly, a finite cover of the manifold $`M`$ constructed above violates this restriction and therefore $`M`$ does not admit almost nonnegative Ricci curvature. However, it could possibly be almost nonnegatively curved in the generalized sense.
### 6.3. Further questions
Recall that a simply connected space $`C`$ is called *rationally elliptic* if it is homotopy equivalent to a finite CW-complex and
$$dim[\pi _{}(C,)]<\mathrm{}.$$
A conjecture of Grove–Halperin \[GH82\] says that simply connected nonnegatively curved manifolds are rationally elliptic. This conjecture was extended by Grove to include almost nonnegatively curved manifolds \[Gro02\]. Later, Totaro has proposed the following definition of rationally elliptic spaces which covers manifolds with infinite fundamental groups:
A connected topological space $`X`$ is *rationally elliptic* if it is homotopy equivalent to a finite CW complex, it has a finite covering which is a nilpotent space and its universal cover is rationally elliptic in the ordinary sense.
With this definition one can extend Grove’s conjecture to non simply connected manifolds as follows:
###### Conjecture 6.3.1.
Any almost nonnegatively curved manifold in the generalized sense is rationally elliptic.
Theorem A reduces this conjecture to the simply connected case which is undoubtedly the most difficult part of the problem.
It has been shown in \[PP06\] that if $`M`$ is a nilpotent closed manifold which admits a Riemannian metric with zero topological entropy, then its universal cover $`\stackrel{~}{M}`$ is rationally elliptic. Coupled with Theorem A this means that to prove Conjecture 6.3.1 it would be sufficient to show that a manifold with almost nonnegative curvature in the generalized sense admits a metric with zero topological entropy. However, we think that this might be wrong in general.
As was pointed out in the discussion in the Introduction before Theorem C, it already follows from Yamaguchi’s fibration theorem and \[FY92\] that a finite cover of an almost nonnegatively curved manifold maps onto a nilmanifold with homotopy fiber a simply connected closed manifold. While this is formally weaker than the statement of Theorem C, it would be interesting to have an answer to the following, purely topological, question:
###### Question 6.3.2.
Let $`M\stackrel{𝑓}{}N`$ be a map from a closed manifold $`M`$ to a nilmanifold $`N`$ such that the homotopy fiber of $`f`$ is a simply connected closed manifold. Is it true that after passing to a finite cover, the map $`f`$ becomes homotopic to a fiber bundle projection?
###### Question 6.3.3.
Is it true that manifolds which are almost nonnegatively curved in the generalized sense are almost nonnegatively curved?
In view of Theorems A and B it is also reasonable to pose the following question:
###### Question 6.3.4.
Is it true that almost nonnegatively curved $`m`$-manifolds $`M^m`$ are $`C(m)`$-nilpotent spaces?
It is clear from the proof of Theorems A and B that this is true if the universal cover of $`M^m`$ has torsion free integral cohomology.
In view of Theorem B it is moreover natural to raise the following question:
###### Question 6.3.5.
Can one give an explicit bound on $`C(m)`$ in Theorem B? |
warning/0506/gr-qc0506016.html | ar5iv | text | # Radiative spacetimes approaching the Vaidya metric
## 1 Introduction
The classic Vaidya metric (see also followed by reprints of the original Vaidya papers) is a spherically symmetric type D solution of the Einstein equations in the presence of pure radiation matter field which propagates at the speed of light. In various contexts this “null dust” may be interpreted as high-frequency electromagnetic or gravitational waves, incoherent superposition of aligned waves with random phases and polarisations, or as massless scalar particles or neutrinos. The Vaidya solution depends on an arbitrary “mass function” $`m(u)`$ of the retarded time $`u`$ which characterises the profile of the pure radiation (it is a “retarded mass” measured at conformal infinity). Various sandwiches and shells of null matter can thus be constructed that are bounded either by flat ($`m=0`$) or Schwarzschild-like ($`m=\text{const}0`$) vacuum regions. Due to this property such solutions have been extensively used as models of spherically symmetric gravitational collapse of a star, as an exterior solution describing objects consisting of heat-conducting matter, as an interesting toy model for investigation of singularities and their possible removal by quantum effects, for studies of various formulations of the cosmic censorship conjecture on both classical and quantum level, process of black-hole evaporation, and for other purposes (see, e.g., for more details and related references).
In fact, the Vaidya spacetime belongs to a large Robinson–Trautman class of expanding nontwisting solutions . Various aspects of this family have been studied in the last two decades. In particular, the existence, asymptotic behaviour and global structure of *vacuum* Robinson–Trautman spacetimes of type II with spherical topology were investigated , most recently in the works of Chruściel and Singleton . In these rigorous studies, which were based on the analysis of solutions to the nonlinear Robinson–Trautman equation for generic, arbitrarily strong smooth initial data, the spacetimes were shown to exist globally for all positive retarded times, and to converge asymptotically to a corresponding Schwarzschild metric. Interestingly, extension across the “Schwarzschild-like” event horizon can only be made with a finite order of smoothness. Subsequently, these results were generalized in to the Robinson–Trautman vacuum spacetimes which admit a nonvanishing *cosmological constant* $`\mathrm{\Lambda }`$. It was demonstrated that these cosmological solutions settle down exponentially fast to a Schwarzschild–(anti-)de Sitter solution at large times $`u`$. In certain cases the interior of a Schwarzschild–de Sitter black hole can be joined to an “external” cosmological Robinson–Trautman region across the horizon with a higher order of smoothness than in the corresponding case with $`\mathrm{\Lambda }=0`$. For the extreme value $`9\mathrm{\Lambda }m^2=1`$, the extension is smooth but not analytic (and not unique). The models with $`\mathrm{\Lambda }>0`$ also exhibit explicitly the cosmic no-hair conjecture under the presence of gravitational waves. On the other hand, when $`\mathrm{\Lambda }<0`$ the smoothness of such an extension is lower.
Our aim here is to further extend the Chruściel–Singleton analysis of the Robinson-Trautman vacuum equation by including matter, namely *pure radiation*. It was argued already by Bičák and Perjés that with $`\mathrm{\Lambda }=0`$ such spacetimes should generically approach the Vaidya metric asymptotically. We will analyze this problem in more detail, including also the possibility of $`\mathrm{\Lambda }0`$ in which case the Robinson–Trautman spacetimes containing pure radiation can be shown to approach the radiating Vaidya–(anti-)de Sitter metric.
## 2 The metric and field equations
In standard coordinates the Robinson–Trautman metric has the form
$$\mathrm{d}s^2=\left(K2r(\mathrm{ln}P)_{,u}2\frac{m}{r}\frac{\mathrm{\Lambda }}{3}r^2\right)\mathrm{d}u^22\mathrm{d}u\mathrm{d}r+2\frac{r^2}{P^2}\mathrm{d}\zeta \mathrm{d}\overline{\zeta },$$
(1)
where $`K=\mathrm{\Delta }(\mathrm{ln}P)`$ with $`\mathrm{\Delta }2P^2_\zeta _{\overline{\zeta }}`$ being the Gaussian curvature of the 2-surfaces $`2P^2\mathrm{d}\zeta \mathrm{d}\overline{\zeta }`$, $`m(u)`$ is the mass function, and $`\mathrm{\Lambda }`$ is the cosmological constant. When the function $`P(u,\zeta ,\overline{\zeta })`$ satisfies the fourth-order Robinson–Trautman field equation
$$\mathrm{\Delta }K+12m(\mathrm{ln}P)_{,u}4m_{,u}=2\kappa n^2,$$
(2)
the metric describes a spacetime (generally of the Petrov type II) filled with pure radiation field $`T_{\mu \nu }=n^2(u,\zeta ,\overline{\zeta })r^2k_\mu k_\nu `$, where $`\text{k}=_r`$ is aligned along the degenerate principal null direction (we use the convention $`G_{\mu \nu }+\mathrm{\Lambda }g_{\mu \nu }=\kappa T_{\mu \nu }`$). In particular, vacuum Robinson–Trautman spacetimes are given by $`n=0`$, in which case $`m`$ can be set to a constant by a suitable coordinate transformation . Vacuum spacetimes (1) — possibly with a nonvanishing $`\mathrm{\Lambda }`$ — thus satisfy the equation $`12m(\mathrm{ln}P)_{,u}=\mathrm{\Delta }K`$. These include the spherically symmetric Schwarzschild–(anti-)de Sitter solution which corresponds to $`P_0=1+\frac{1}{2}\zeta \overline{\zeta }`$. Indeed, replacing the complex stereographic coordinate $`\zeta `$ with angular coordinates by $`\zeta =\sqrt{2}e^{i\varphi }\mathrm{tan}(\theta /2)`$, we obtain $`2P_0^2\mathrm{d}\zeta \mathrm{d}\overline{\zeta }=\mathrm{d}\theta ^2+\mathrm{sin}^2\theta \mathrm{d}\varphi ^2`$, and $`K_0=\mathrm{\Delta }_0\mathrm{ln}(P_0)=1`$.
Here we will restrict ourselves to nonvacuum cases for which the dependence of the mass function $`m(u)`$ on the null coordinate $`u`$ is only caused by a *homogeneous* pure radiation with the density $`n^2(u)r^2`$. When the mass function $`m(u)`$ is decreasing, the field equation (2) can be naturally split into the following pair,
$`\mathrm{\Delta }K+12m(u)(\mathrm{ln}P)_{,u}`$ $`=`$ $`0,`$ (3)
$`2m(u)_{,u}`$ $`=`$ $`\kappa n^2(u).`$ (4)
In fact, it was demonstrated in that such a separation can always be achieved using the coordinate freedom. It is then possible to reformulate equation (3) by introducing a $`u`$-dependent family of smooth 2-metrics $`g_{ab}`$ on the submanifold $`r=\text{const}`$, $`u=\text{const}`$, such that $`g_{ab}=f(u,\zeta ,\overline{\zeta })^2g_{ab}^0`$, where $`g_{ab}^0(\zeta ,\overline{\zeta })`$ is the metric on a 2-dimensional sphere $`S^2`$. Since $`g_{ab}`$ is of the form $`2P^2\mathrm{d}\zeta \mathrm{d}\overline{\zeta }`$ in our case, we can write
$$P=fP_0,P_0=1+\frac{1}{2}\zeta \overline{\zeta },$$
(5)
and equation (3) becomes
$$\frac{f}{u}=\frac{1}{12m(u)}f\mathrm{\Delta }K,$$
(6)
where $`\mathrm{\Delta }`$ is the Laplace operator associated with the metric $`g_{ab}`$. Denoting $`\mathrm{\Delta }_0`$ and $`K_0=1`$ as the corresponding quantities related to $`g_{ab}^0`$, we obtain
$$\mathrm{\Delta }=f^2\mathrm{\Delta }_0,K=f^2(1+\mathrm{\Delta }_0(\mathrm{ln}f)).$$
(7)
## 3 Linear mass function
Let us first consider the simplest choice of $`m(u)`$ which, in fact, has been widely used in literature (see e.g. ): we will assume that the mass function is a *linearly* decreasing positive function
$$m(u)=\mu u,\mu =\text{const}>0,$$
(8)
on the interval $`[u_0,0]`$. Notice that for (8) the pure radiation field is uniform because equation (4) implies $`n=\sqrt{2\mu /\kappa }=\text{const}`$, independent of the retarded time $`u`$. The constant value $`u_0<0`$ localises an initial null hypersurface (that extends between the curvature singularity at $`r=0`$ and the conformal infinity $`r=\mathrm{}`$) on which an arbitrary sufficiently smooth *initial data* given by the function
$$f_0(\zeta ,\overline{\zeta })=f(u=u_0,\zeta ,\overline{\zeta }),$$
(9)
are prescribed, see Fig. 1.
### 3.1 Existence of the solutions
Now, the idea is to employ the Chruściel–Singleton results concerning the analysis of the Robinson–Trautman *vacuum* equation, in particular the existence and asymptotic behaviour of its solutions. In the vacuum case $`m`$ in equation (3) is constant, and the solution $`f(u,\zeta ,\overline{\zeta })`$ of the characteristic initial value problem (9) exists and is unique (in spite of the singularity at $`r=0`$). In the presence of pure radiation given by (8) it is possible to “eliminate” the variable mass function from the Robinson–Trautman field equation (6) mathematically by a simple reparametrisation
$$\stackrel{~}{u}=\mu ^1\mathrm{ln}(u),$$
(10)
cf. . Indeed, equation (6) is then converted to
$$\frac{\stackrel{~}{f}}{\stackrel{~}{u}}=\frac{1}{12}\stackrel{~}{f}\stackrel{~}{\mathrm{\Delta }}\stackrel{~}{K},$$
(11)
where $`\stackrel{~}{f}(\stackrel{~}{u},\zeta ,\overline{\zeta })=f(u(\stackrel{~}{u}),\zeta ,\overline{\zeta })`$, $`\stackrel{~}{K}=\stackrel{~}{f}^2(1+\mathrm{\Delta }_0\mathrm{ln}(\stackrel{~}{f}))`$, and $`\stackrel{~}{\mathrm{\Delta }}=\stackrel{~}{f}^2\mathrm{\Delta }_0`$. Notice that the transformation (10) moves the hypersurface $`u=0`$, on which the mass function $`m(u)`$ reaches zero, to $`\stackrel{~}{u}=+\mathrm{}`$.
Chruściel derived the following asymptotic expansion (as $`\stackrel{~}{u}\mathrm{}`$) for the function $`\stackrel{~}{f}`$ satisfying the evolution equation (11) for any smooth initial data $`\stackrel{~}{f}_0=f_0`$ on $`\stackrel{~}{u}_0=\mu ^1\mathrm{ln}(u_0)`$, namely
$`\stackrel{~}{f}`$ $`=`$ $`1+f_{1,0}e^{2\stackrel{~}{u}}+f_{2,0}e^{4\stackrel{~}{u}}+\mathrm{}+f_{14,0}e^{28\stackrel{~}{u}}`$
$`+f_{15,1}\stackrel{~}{u}e^{30\stackrel{~}{u}}+f_{15,0}e^{30\stackrel{~}{u}}+\mathrm{}`$
$`=`$ $`{\displaystyle \underset{i=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{j=0}{\overset{N_i}{}}}f_{i,j}\stackrel{~}{u}^je^{2i\stackrel{~}{u}},`$
where $`f_{i,j}`$ are smooth functions on $`S^2`$ such that $`f_{i,j}=0`$ for $`j>0`$, $`i14`$. The function $`\stackrel{~}{f}`$ thus exists and converges exponentially fast to 1, which means physically that the radiative Robinson–Trautman vacuum spacetimes approach asymptotically the Schwarzschild–(anti-)de Sitter solution as $`\stackrel{~}{u}\mathrm{}`$, see relation (5). In our case of pure radiation field (8) we employ the transformation (10) on expression (3.1) to obtain the following asymptotic expansion of $`f`$ as $`u0_{}`$,
$`f`$ $`=`$ $`1+f_{1,0}(u)^{2/\mu }+f_{2,0}(u)^{4/\mu }+\mathrm{}+f_{14,0}(u)^{28/\mu }`$
$`\mu ^1f_{15,1}\mathrm{ln}(u)(u)^{30/\mu }+f_{15,0}(u)^{30/\mu }+\mathrm{}`$
$`=`$ $`{\displaystyle \underset{i=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{j=0}{\overset{N_i}{}}}f_{i,j}[\mu ^1\mathrm{ln}(u)]^j(u)^{2i/\mu }.`$
As a result, for the initial data (9) the Robinson–Trautman type II spacetimes which contain uniform pure radiation field with the linear mass function (8) do exist in the whole region $`u_0u<0`$. It is also obvious that the function $`f`$ approaches $`1`$ as $`u0_{}`$ (where also $`m(u)0`$) according to (3.1). In other words, these spacetimes *approach the spherically symmetric Vaidya–(anti-)de Sitter metric* near $`u=0`$.
The global structure of such spacetimes is schematically indicated on Fig. 1. In fact, there are two possibly different conformal diagrams depending on the value of $`\mu `$: for $`\mu >1/16`$ there is a white hole singularity at $`r=0`$, for $`\mu 1/16`$ there is also a naked singularity, see e.g. for more details. At $`u=0`$ all of the mass $`m(u)`$ is radiated away, and we can attach Minkowski space (de Sitter space when $`\mathrm{\Lambda }>0`$, anti-de Sitter when $`\mathrm{\Lambda }<0`$; the presence of the cosmological constant would change the character of conformal infinity $``$ which would become spacelike or timelike, respectively) in the region $`u>0`$ along the hypersurface $`u=0`$. We will now investigate the smoothness of such an extension.
### 3.2 Extension of the metric across $`u=0`$
It follows from (3.1) that the smoothness of $`f`$ on $`u=0`$ is only finite. Depending on the value of $`\mu `$ two different cases have to be discussed separately: $`2/\mu `$ is an integer, and $`2/\mu `$ is a real non-integer positive number.
When $`2/\mu `$ is an integer then due to the presence of the $`\mathrm{ln}(u)`$ term associated with $`f_{15,1}0`$ the function $`f`$ is of the class $`C^{(30/\mu )1}`$. For $`\mu `$ very small, the integer number $`(30/\mu )1`$ is large so that $`f`$ becomes smoothly extendable to 1 across $`u=0`$ as $`\mu 0`$. This represents a naked-singularity Robinson–Trautman spacetime (see the right part of Fig. 1) unless $`\mu =0`$ which gives flat space everywhere. In the limiting case $`\mu =1/16`$ the function $`f`$ is of the class $`C^{479}`$. For the (white hole) Robinson–Trautman spacetimes given by $`\mu >1/16`$ the smoothness is lower. However, it is always at least $`C^{14}`$ because $`\mu 2`$ in this case.
In the generic case when $`2/\mu `$ is not an integer the function $`f`$ is only of the class $`C^{\{2/\mu \}}`$, where the symbol $`\{x\}`$ denotes the largest integer smaller than $`x`$. Again, with $`\mu 0`$ the function $`f`$ becomes smoothly extendable. For $`\mu <1/16`$ the function $`f`$ is at least of the class $`C^{32}`$, for $`\mu >2`$ it is not even $`C^1`$ but it remains continuous.
To investigate further the smoothness of the metric when approaching the hypersurface $`u=0_{}`$ which is the analogue of the Schmidt–Tod boundary of vacuum Robinson–Trautman spacetimes we should consider the conformal picture using suitable double-null coordinates. Such Kruskal-type coordinates for the Vaidya solution with linear mass function (8) were introduced by Hiscock , see also , and we will use this transformation only to replace the coordinate $`r`$ since the null coordinate $`u`$ is already appropriate. Introducing a new coordinate $`w`$ by
$$\mathrm{d}w=\frac{\mathrm{d}u}{u}\frac{2\mathrm{d}z}{z(2\mu z^2z+2)},\text{where}z=\frac{u}{r},$$
(14)
we put the Robinson–Trautman metric with linear mass function into the form
$`\mathrm{d}s^2=\left(K12{\displaystyle \frac{f_{,u}}{f}}r\right)\mathrm{d}u^2`$
$`\left(2r+u+2\mu {\displaystyle \frac{u^2}{r}}\right)\mathrm{d}u\mathrm{d}w+2{\displaystyle \frac{r^2}{P^2}}\mathrm{d}\zeta \mathrm{d}\overline{\zeta },`$ (15)
where $`r(u,w)`$. For the pure Vaidya metric characterized by $`f=1`$ and $`K_0=1`$ the first term vanishes identically so that the coordinates of (3.2) are indeed the Kruskal-type coordinates for the Vaidya spacetime with a linear mass function.
The smoothness of a general Robinson–Trautman metric (3.2) depends only on the smoothness of the metric coefficients $`g_{uu}`$ and $`g_{\zeta \overline{\zeta }}`$ (containing the function $`f`$) since the coefficient $`g_{uw}`$ tends to $`r`$ as $`u0`$. The smoothness of $`g_{\zeta \overline{\zeta }}`$ (for any finite $`r`$) and of $`K`$ is the same as of $`f`$, see (7). The function $`f_{,u}/f`$ is evidently one order less smooth than $`f`$. Consequently, for $`2/\mu `$ being integer or non-integer number, the metric (3.2) is of the class $`C^{(30/\mu )2}`$ or $`C^{\{2/\mu \}1}`$, respectively. We again observe that the spacetimes approaching the linear Vaidya metric with naked singularity (i.e., for small values of the parameter $`\mu `$) possess higher order of smoothness at $`u=0`$.
One might be worried about the invariance of our results, namely with respect to a rescaling of the null coordinate $`u(\widehat{u})`$ leading to a different smoothness of the function $`f`$ and of the metric. In order to change the smoothness on the hypersurface $`u=0`$ such rescaling must have a singular character there. But this would lead to a degeneracy of the metric coefficient $`g_{\widehat{u}w}`$ of the Vaidya metric, which is forbidden. Consequently, the above results are in this sense unique.
We would like to obtain analogous results concerning smoothness of the extension also for a non-zero value of the cosmological constant $`\mathrm{\Lambda }`$. Unfortunately, as far as we know, there is no *explicit* transformation of the Vaidya–de Sitter metric to the Kruskal-type coordinates even for the linear mass function (contrary to the Schwarzschild–de Sitter case ). However, it is possible to start with the Vaidya–de Sitter metric
$$\mathrm{d}s^2=h(u,r)\mathrm{d}u^22\mathrm{d}u\mathrm{d}r+r^2\mathrm{d}\mathrm{\Omega }^2,$$
(16)
where $`h(u,r)=1+2\mu ur^1\frac{\mathrm{\Lambda }}{3}r^2`$, and perform a coordinate transformation
$$\mathrm{d}w=g\mathrm{d}u+2\frac{g}{h}\mathrm{d}r,$$
(17)
where $`g(u,r)`$ is some function. We arrive at the double-null form for the metric
$$\mathrm{d}s^2=\frac{h}{g}\mathrm{d}u\mathrm{d}w+r^2(u,w)\mathrm{d}\mathrm{\Omega }^2.$$
(18)
Of course, we have to ensure that $`\mathrm{d}w`$ in (17) is a differential of the coordinate $`w`$. The integrability condition ($`\mathrm{d}^2w=0`$) gives the following quasilinear PDE,
$$h^2\frac{g}{r}2h\frac{g}{u}+4\frac{\mu }{r}g=0,$$
(19)
for the undetermined function $`g`$, which is difficult to solve analytically. The method of characteristic curves leads to the first-order ODE of the Abel type which has not yet been solved, but the existence of its solution is guaranteed. \[It is possible to apply the perturbative approach starting from the solvable case of the de Sitter metric ($`\mu =0`$) and then linearise the PDE in the parameter $`\mu `$. The result, however, can not be presented in a useful closed form.\] For our purposes it suffices to use a general argumentation: the coordinate $`u`$ is already suitably compactified and we are only determining the complementary null coordinate $`w`$ to obtain the Vaidya–de Sitter metric in the Kruskal-type coordinates (which is smooth on $`u=0`$). The corresponding Robinson–Trautman metric in these coordinates differs only by the term $`g_{uu}(u,r,\zeta ,\overline{\zeta })\mathrm{d}u^2`$ (which is absent in the Vaidya–de Sitter case in the double null coordinates), and by a different metric coefficient $`g_{\zeta \overline{\zeta }}=r^2f^2P_0^2`$, where $`r(u,w)`$ is finite and smooth when approaching the hypersurface $`u=0`$. The smoothness is thus not affected by the specific transformation (17) and it is the same as for the vanishing cosmological constant. This is different from vacuum spacetimes with $`m=\text{const}0`$ studied in because in the present case $`m0`$ near $`u=0`$, and the influence of $`\mathrm{\Lambda }`$ on the smoothness becomes negligible.
## 4 General mass function
The results obtained above can be considerably generalized. Inspired by a similar idea outlined in we may consider a reparametrisation on the null coordinate $`u`$ by
$$\stackrel{~}{u}=\gamma (u),$$
(20)
where $`\gamma `$ is an arbitrary continuous strictly monotonous function. We start with the evolution equation (11) for which the existence and uniqueness of solutions has been proven, and their general asymptotic behaviour (3.1) has been demonstrated. Now, by applying the substitution (20) in equation (11) we obtain
$$\frac{f}{u}=\frac{\dot{\gamma }}{12}f\mathrm{\Delta }K,$$
(21)
(where the dot denotes a differentiation) which is the evolution equation for the function $`f(u,\zeta ,\overline{\zeta })`$. This is exactly the Robinson–Trautman equation (6) for the mass function
$$m(u)=\frac{1}{\dot{\gamma }(u)}.$$
(22)
For a given smooth initial data on $`u_0`$ there thus exists the Robinson–Trautman spacetime (1), including the cosmological constant $`\mathrm{\Lambda }`$, with the mass function (22). To obtain a positive mass we consider a growing function $`\gamma (u)`$. Considering (4) this corresponds to a universe filled with homogeneous pure radiation
$$n^2(u)=\frac{2}{\kappa }\frac{\ddot{\gamma }}{\dot{\gamma }^2}.$$
(23)
For consistency the function $`\gamma `$ must be convex. An asymptotic behaviour of the function $`f`$ as $`\gamma (u)\mathrm{}`$ is easily obtained from the expansion (3.1) by substituting relation (20).
In particular, the linear mass function (8) discussed above is a special case of (22) for the transformation (20) of the form (10). More general explicit solutions can be obtained, e.g., by considering the power function
$$\gamma (u)=(u)^p,p>0,$$
(24)
which gives
$$m(u)=\frac{1}{p}(u)^{1+p},n^2(u)=\frac{2(p+1)}{\kappa p}(u)^p.$$
(25)
Both functions $`m`$ and $`n`$ approach zero as $`u0`$. Due to the theorems mentioned above, there exist Robinson–Trautman type II spacetimes in the region $`u<0`$ which approach the spherically symmetric Vaidya–(anti-)de Sitter metric as $`u0_{}`$ with the mass function and pure radiation given by (25). The asymptotic behaviour of such solutions is determined by expression (5) with
$$f=1+\underset{i=1}{\overset{\mathrm{}}{}}\underset{j=0}{\overset{N_i}{}}f_{i,j}(u)^{jp}\mathrm{exp}\left[2i(u)^p\right],$$
where $`f_{i,j}=0`$ for $`j>0`$ if $`i14`$. Interestingly, the function $`f`$ is now smooth on $`u=0`$ for any power coefficient $`p`$, but this still does not guarantee that the extension into flat region $`u>0`$ is analytic (see for a similar situation).
Another simple explicit choice is
$$\gamma (u)=M^1\mathrm{ln}\left[\mathrm{sinh}(u)\right],M>0,$$
(26)
which implies (see also )
$$m(u)=M\mathrm{tanh}(u),n^2(u)=\frac{2M}{\kappa \mathrm{cosh}^2u}.$$
(27)
In the region $`u<0`$ the mass function monotonically decreases from $`M`$ to zero, while the pure radiation field grows from zero to the value $`2M/\kappa `$ as $`u0`$. Let us note that in this case the integrated radiation density is finite on the interval $`(\mathrm{},0)`$, $`_{\mathrm{}}^0n^2(u)=2M/\kappa `$. The expansion near $`u=0_{}`$ is
$$f=1+\underset{i=1}{\overset{\mathrm{}}{}}\underset{j=0}{\overset{N_i}{}}f_{i,j}(M^1\mathrm{ln}\left[\mathrm{sinh}(u)\right])^j\mathrm{sinh}^{2i/M}(u).$$
If $`\mathrm{\hspace{0.17em}2}/M`$ is an integer then the function $`f`$ belongs to the class $`C^{(30/M)1}`$, otherwise it is of the class $`C^{\{2/M\}}`$.
## 5 Possible modifications and applications
The Robinson–Trautman pure radiation solutions in the region $`u_0u0`$ approaching the Vaidya metric near $`u=0`$, which can be extended (albeit non-smoothly) to flat Minkowski space in the region $`u0`$ as in Fig. 1, may be used for construction of various models of radiative spacetimes. For example, it is natural to further extend the solution “backwards” into the region $`u_1<uu_0`$ by the Robinson–Trautman *vacuum* solution with a constant mass $`m_0=m(u_0)`$, such that the function $`f`$ is continuous on $`u_0`$. This is shown in Fig. 2. In such a case the spacetime may describe the process of “evaporation” of a white hole (with a different character of the singularity at $`r=0`$ when $`\mu 1/16`$) with its mass decreasing from the value $`m_0`$ to zero. Let us emphasize that the region $`u<u_0`$ does not represent the Schwarzschild solution because the spacetime is *not spherically symmetric* there ($`f1`$). In fact, this is the region where the original Chruściel theorems on the behaviour of the Robinson–Trautman vacuum spacetimes with constant mass apply (cf. (11), (3.1)). However, the spacetime in this region can not be extended up to the past conformal infinity $`^{}`$ because the metric function $`f`$ diverges as $`u\mathrm{}`$.
In the presence of the cosmological constant $`\mathrm{\Lambda }`$ one obtains a family of exact spacetimes that describe evaporation of a white hole in the (anti-)de Sitter universe. In this case the schematic conformal diagram on Fig. 2 has to be modified in such a way that for all values of $`u`$ the conformal infinity $`^+`$ becomes timelike (for $`\mathrm{\Lambda }>0`$) or spacelike (for $`\mathrm{\Lambda }<0`$).
Another possible modification is to consider the “advanced” form of the spacetimes (which describes an ingoing flow) rather than the “retarded” form (corresponding to outgoing flow) employed above (see, e.g., for more details). This time-reversed form is obtained formally by a simple substitution $`uv`$ in the metrics and corresponding functions. The Robinson–Trautman metric thus reads
$$\mathrm{d}s^2=\left(K+2r(\mathrm{ln}P)_{,v}2\frac{m}{r}\frac{\mathrm{\Lambda }}{3}r^2\right)\mathrm{d}v^2+2\mathrm{d}v\mathrm{d}r+2\frac{r^2}{P^2}\mathrm{d}\zeta \mathrm{d}\overline{\zeta },$$
(28)
where $`m(v)`$ is an increasing mass function in $`v[0,v_0]`$. This is joined with flat Minkowskian region $`v<0`$, and extended to the region $`vv_0`$ by the corresponding Robinson–Trautman–(anti-)de Sitter black hole vacuum solution, see Fig. 3. It is a non-spherical generalization of the gravitational collapse of a shell of null dust forming a naked singularity — in these works the mass function was taken to be $`m(v)=\mu v`$ (with $`m(v)=0`$ for $`v0`$, and $`m(v)=M=\mu v_0`$ for $`vv_0`$). The metric function $`P`$ is now given by $`P=fP_0`$ where $`f`$ is analogous to (3.1),
$$f=\underset{i=0}{\overset{\mathrm{}}{}}\underset{j=0}{\overset{N_i}{}}f_{i,j}\left(\mu ^1\mathrm{ln}v\right)^jv^{2i/\mu },$$
(29)
so that the smoothness of the metric on the boundary $`v=0`$ depends on the parameter $`\mu `$. For $`v(v_0,v_1)`$ the spacetime is vacuum but not spherically symmetric. The metric diverges as $`v\mathrm{}`$. Our results can thus be interpreted in such a way that — at least within the Robinson–Trautman family of solutions — the model of collapse to a naked shell-focusing singularity which is based on the spherically symmetric Vaidya metric *is not stable* against perturbations.
## 6 Concluding remarks
In our contribution we have analyzed exact solutions of the Robinson–Trautman class which contain homogeneous pure radiation and a cosmological constant. This is a natural extension of previous works on properties of vacuum spacetimes of this family. We have demonstrated that these solutions exist for any smooth initial data, and that they approach the spherically symmetric Vaidya–(anti-)de Sitter metric. It generalizes previous results according to which vacuum Robinson–Trautman spacetimes approach asymptotically the spherically symmetric Schwarzschild–(anti-)de Sitter metric. We have investigated extensions of these solutions into Minkowski region, and we have shown that its order of smoothness is in general only finite. Finally, we suggested some applications of the results. For example, it follows that the model of gravitational collapse of a shell of null dust diverges as $`v\mathrm{}`$ which indicates that investigations of such process based on the spherically symmetric Vaidya metric are, in fact, not stable against “non-linear perturbations”, at least within the Robinson–Trautman family of exact solutions.
## Acknowledgements
We are grateful to Jiří Bičák for valuable comments, and Jerry Griffiths for reading the manuscript. |
warning/0506/hep-ph0506092.html | ar5iv | text | # DPNU-05-10 Pion Velocity near the Chiral Phase Transition Point in the Vector Manifestation #1#1footnote #1 Contribution to the proceedings of the YITP workshop on Non-equilibrium dynamics in the QCD phase transitions, February 22 - 24, 2005, YITP, Kyoto, Japan. Masayasu Harada (a), Mannque Rho (b) and Chihiro Sasaki (c) (a) Department of Physics, Nagoya University, Nagoya, 464-8602, Japan, (b) Service de Physique Théorique, CEA Saclay, 91191 Gif-sur-Yvette, France, (c) Gesellschaft für Schwerionenforschung (GSI), Planckstr. 1, 64291 Darmstadt, Germany
## 1 Introduction
It has been suggested by Cramer et al. that the recent result by the STAR collaboration at RHIC provides information on the pion velocity near the chiral phase transition in hot medium with the conclusion that the pions seen in Hanbury Brown-Twist interferometry are emitted from a chiral-symmetry restored phase . In this note, we wish to interpret the result of the analysis by Cramer et al. in terms of the vector manifestation scenario of hidden local symmetry theory .
The effective field theory based on the hidden local symmetry (HLS), which includes both pions and vector mesons as the dynamical degrees of freedom, implemented with the Wilsonian matching to determine the bare theory from the underlying QCD, leads to the vector manifestation (VM) of chiral symmetry in which the massless vector meson becomes the chiral partner of the pion at the critical point <sup>#2</sup><sup>#2</sup>#2 As studied in Ref. in detail, the VM is defined only as a limit with bare parameters approaching the VM fixed point from the broken phase.. This picture provides a strong support for Brown-Rho scaling which predicted that the mass of light-quark hadrons should drop in proportion to the quark condensate $`\overline{q}q`$. By now there are several experimental indications that this scenario is a viable one. The earliest one was the enhancement of dielectron mass spectra below the $`\rho /\omega `$ resonance observed at CERN SPS which has been satisfactorily explained by the dropping of the $`\rho `$ meson mass according to the Brown-Rho scaling . This explanation however is not unique as there are alternative – but not necessarily unrelated – mechanisms that can equally well describe the presently available data . A much more compelling evidence comes from the mass shift of the $`\omega `$ meson in nuclei measured by the KEK-PS E325 Experiment and the CBELSA/TAPS Collaboration , and also from that of the $`\rho `$ meson observed in the STAR experiment . These are clean signals manifested in a “pristine” environment unencumbered by a plethora of “trivial” effects.
In this note, we focus on the pion velocity near the critical temperature and make a prediction based on the vector manifestation (VM). The pion velocity is one of the important observable quantities in heavy-ion collisions as it controls the pion propagation in medium through a dispersion relation. Our prediction for $`v_\pi `$ is $`v_\pi (T_c)=0.830.99`$, which we should stress, is at strong variance with the result obtained in sigma models <sup>#3</sup><sup>#3</sup>#3By sigma models, we mean generically the chiral symmetry models that contain pions as the $`only`$ relevant long-wavelength degrees of freedom., i.e., $`v_\pi (T_c)=0`$ . We believe our result to be consistent with $`v_\pi (T)=0.65`$ of Cramer et al. extracted from the recent STAR data . What distinguishes our approach from that of sigma models is the intrinsic temperature and/or density dependence of the parameters of the HLS Lagrangian, that results from integrating out the high energy modes (i.e., the quarks and gluons above the matching scale) in medium . It is this intrinsic temperature and/or density dependence – which causes Lorentz symmetry breaking – that plays the essential role for realizing the chiral symmetry restoration in a consistent way and underlies the Brown-Rho scaling.
## 2 Vector manifestation of chiral symmetry
In this section, we start with the HLS Lagrangian at leading order including the effects of Lorentz non-invariance. Then we present the conditions satisfied at the critical point.
The HLS theory is based on the $`G_{\mathrm{global}}\times H_{\mathrm{local}}`$ symmetry, where $`G=SU(N_f)_\mathrm{L}\times SU(N_f)_\mathrm{R}`$ is the chiral symmetry and $`H=SU(N_f)_\mathrm{V}`$ is the HLS. The basic quantities are the HLS gauge boson $`V_\mu `$ and two matrix valued variables $`\xi _\mathrm{L}(x)`$ and $`\xi _\mathrm{R}(x)`$ which transform as $`\xi _{\mathrm{L},\mathrm{R}}(x)\xi _{\mathrm{L},\mathrm{R}}^{}(x)=h(x)\xi _{\mathrm{L},\mathrm{R}}(x)g_{\mathrm{L},\mathrm{R}}^{}`$, where $`h(x)H_{\mathrm{local}}\text{and}g_{\mathrm{L},\mathrm{R}}[SU(N_f)_{\mathrm{L},\mathrm{R}}]_{\mathrm{global}}`$. These variables are parameterized as <sup>#4</sup><sup>#4</sup>#4 The wave function renormalization constant of the pion field is given by the temporal component of the pion decay constant . Thus we normalize $`\pi `$ and $`\sigma `$ by $`F_\pi ^t`$ and $`F_\sigma ^t`$ respectively. $`\xi _{\mathrm{L},\mathrm{R}}(x)=e^{i\sigma (x)/F_\sigma ^t}e^{i\pi (x)/F_\pi ^t}`$, where $`\pi =\pi ^aT_a`$ denotes the pseudoscalar Nambu-Goldstone (NG) bosons associated with the spontaneous symmetry breaking of $`G_{\mathrm{global}}`$ chiral symmetry, and $`\sigma =\sigma ^aT_a`$ denotes the NG bosons associated with the spontaneous breaking of $`H_{\mathrm{local}}`$. This $`\sigma `$ is absorbed into the HLS gauge boson through the Higgs mechanism, and then the vector meson acquires its mass. $`F_\pi ^t`$ and $`F_\sigma ^t`$ denote the temporal components of the decay constant of $`\pi `$ and $`\sigma `$, respectively. The covariant derivative of $`\xi _L`$ is given by
$$D_\mu \xi _L=_\mu \xi _LiV_\mu \xi _L+i\xi _L_\mu ,$$
(1)
and the covariant derivative of $`\xi _R`$ is obtained by the replacement of $`_\mu `$ with $`_\mu `$ in the above where $`V_\mu `$ is the gauge field of $`H_{\mathrm{local}}`$, and $`_\mu `$ and $`_\mu `$ are the external gauge fields introduced by gauging $`G_{\mathrm{global}}`$ symmetry. In terms of $`_\mu `$ and $`_\mu `$, we define the external axial-vector and vector fields as $`𝒜_\mu =(_\mu _\mu )/2`$ and $`𝒱_\mu =(_\mu +_\mu )/2`$.
In the HLS theory it is possible to perform the derivative expansion systematically . In the chiral perturbation theory (ChPT) with HLS, the vector meson mass is to be considered as small compared with the chiral symmetry breaking scale $`\mathrm{\Lambda }_\chi `$, by assigning $`𝒪(p)`$ to the HLS gauge coupling, $`g𝒪(p)`$ . (For details of the ChPT with HLS, see Ref. .) The leading order Lagrangian with Lorentz non-invariance can be written as
$`=\left[(F_\pi ^t)^2u_\mu u_\nu +F_\pi ^tF_\pi ^s\left(g_{\mu \nu }u_\mu u_\nu \right)\right]`$
$`\times \text{tr}\left[\widehat{\alpha }_{}^\mu \widehat{\alpha }_{}^\nu \right]`$
$`+\left[(F_\sigma ^t)^2u_\mu u_\nu +F_\sigma ^tF_\sigma ^s\left(g_{\mu \nu }u_\mu u_\nu \right)\right]`$
$`\times \text{tr}\left[\widehat{\alpha }_{}^\mu \widehat{\alpha }_{}^\nu \right]`$
$`+[{\displaystyle \frac{1}{g_L^2}}u_\mu u_\alpha g_{\nu \beta }`$
$`{\displaystyle \frac{1}{2g_T^2}}(g_{\mu \alpha }g_{\nu \beta }2u_\mu u_\alpha g_{\nu \beta })]`$
$`\times \text{tr}\left[V^{\mu \nu }V^{\alpha \beta }\right],`$ (2)
where
$$\widehat{\alpha }_,^\mu =\frac{1}{2i}\left[D^\mu \xi _R\xi _R^{}D^\mu \xi _L\xi _L^{}\right].$$
(3)
Here $`F_\pi ^s`$ denote the spatial pion decay constant and similarly $`F_\sigma ^s`$ for the $`\sigma `$. The rest frame of the medium is specified by $`u^\mu =(1,\stackrel{}{0})`$ and $`V_{\mu \nu }`$ is the field strength of $`V_\mu `$. $`g_L`$ and $`g_T`$ correspond in medium to the HLS gauge coupling $`g`$. The parametric $`\pi `$ and $`\sigma `$ velocities are defined by
$$V_\pi ^2=F_\pi ^s/F_\pi ^t,V_\sigma ^2=F_\sigma ^s/F_\sigma ^t.$$
(4)
Now we approach the critical point of chiral symmetry restoration, which is characterized by the equality between the axial-vector and vector current correlators in QCD, $`G_AG_V0`$ for $`TT_c`$. This should hold even in the EFT side. In Ref. , it was shown that they are satisfied for any values of $`p_0`$ and $`\overline{p}`$ around the matching scale only if the following conditions are met: $`(g_{L,\mathrm{bare}},g_{T,\mathrm{bare}},a_{\mathrm{bare}}^t,a_{\mathrm{bare}}^s)(0,0,1,1)`$ for $`TT_c`$. This implies that at the bare level the longitudinal mode of the vector meson becomes the real NG boson and couples to the vector current correlator, while the transverse mode decouples. As shown in Ref. , $`(g_L,a^t,a^s)=(0,1,1)`$ is a fixed point of the RGEs and satisfied at any energy scale. Thus the VM condition is given by
$`(g_L,a^t,a^s)(0,1,1)\text{for}TT_c.`$ (5)
The vector meson mass is never generated at the critical temperature since the quantum correction to $`M_\rho ^2`$ is proportional to $`g_L^2`$. Because of $`g_L0`$, the transverse vector meson at the critical point, at any energy scale, decouples from the vector current correlator. The VM condition for $`a^t`$ and $`a^s`$ leads to the equality between the $`\pi `$ and $`\sigma `$ (i.e., longitudinal vector meson) velocities:
$`\left(V_\pi /V_\sigma \right)^4=\left(F_\pi ^sF_\sigma ^t/F_\sigma ^sF_\pi ^t\right)^2=a^t/a^s\stackrel{TT_c}{}1.`$ (6)
This can be easily understood from the point of view of the VM since the longitudinal vector meson becomes the chiral partner of pion. We note that this condition $`V_\sigma =V_\pi `$ holds independently of the value of the bare pion velocity which is to be determined through the Wilsonian matching.
## 3 Pion velocity near the critical temperature
One possible way to determine the bare parameters is the Wilsonian matching which is done by matching the axial-vector and vector current correlators derived from the HLS with those by the operator product expansion (OPE) in QCD at the matching scale $`\mathrm{\Lambda }`$ . The Wilsonian matching leads to the following conditions on the bare pion decay constants :
$`{\displaystyle \frac{F_{\pi ,\mathrm{bare}}^tF_{\pi ,\mathrm{bare}}^s}{\mathrm{\Lambda }^2}}={\displaystyle \frac{1}{8\pi ^2}}[(1+{\displaystyle \frac{\alpha _s}{\pi }})`$
$`+{\displaystyle \frac{2\pi ^2}{3}}{\displaystyle \frac{\frac{\alpha _s}{\pi }G^2_T}{\mathrm{\Lambda }^4}}+\pi ^3{\displaystyle \frac{1408}{27}}{\displaystyle \frac{\alpha _s\overline{q}q_T^2}{\mathrm{\Lambda }^6}}]`$
$`+{\displaystyle \frac{\pi ^2}{15}}{\displaystyle \frac{T^4}{\mathrm{\Lambda }^4}}A_{4,2}^\pi {\displaystyle \frac{16\pi ^4}{21}}{\displaystyle \frac{T^6}{\mathrm{\Lambda }^6}}A_{6,4}^\pi G_0,`$
$`{\displaystyle \frac{F_{\pi ,\mathrm{bare}}^tF_{\pi ,\mathrm{bare}}^s(1V_{\mathrm{bare}}^2)}{\mathrm{\Lambda }^2}}={\displaystyle \frac{32}{105}}\pi ^4{\displaystyle \frac{T^6}{\mathrm{\Lambda }^6}}A_{6,4}^\pi ,`$
where we use the dilute pion-gas approximation in order to evaluate the matrix element $`𝒪_T`$ in the low temperature region. From these conditions, we obtain the following matching condition to determine the deviation of the bare pion velocity from the speed of light in the low temperature region:
$$\delta _{\mathrm{bare}}1V_{\pi ,\mathrm{bare}}^2=\frac{1}{G_0}\frac{32}{105}\pi ^4\frac{T^6}{\mathrm{\Lambda }^6}A_4^\pi .$$
(8)
This implies that the intrinsic temperature dependence starts from the $`𝒪(T^6)`$ contribution.
As is discussed in Ref. , we should in principle evaluate the matrix elements in terms of QCD variables only in order for performing the Wilsonian matching, which is as yet unavailable from model-independent QCD calculations. Therefore, we make an estimation by extending the dilute gas approximation adopted in the QCD sum rule analysis in low temperature region to the critical temperature including all the light degrees of freedom expected in the VM. In the HLS/VM theory, both the longitudinal and transverse vector mesons become massless at the critical temperature. At the critical point, the longitudinal vector meson couples to the vector current whereas the transverse vector mesons decouple from the theory. Thus we assume that thermal fluctuations of the system are dominated near $`T_c`$ not only by the pions but also by the longitudinal vector mesons. We evaluate the thermal matrix elements of the non-scalar operators in the OPE, by extending the thermal pion gas approximation employed in Ref. to the longitudinal vector mesons that figure in our approach.
This is feasible since at the critical temperature, we expect the equality $`A_4^\rho (T_c)=A_4^\pi (T_c)`$ to hold as the massless longitudinal vector meson is the chiral partner of the pion in the VM. It should be noted that, although we use the dilute gas approximation, the treatment here is already beyond the low-temperature approximation because the contribution from vector meson is negligible in the low-temperature region. Since we treat the pion as a massless particle in the present analysis, it is reasonable to take $`A_4^\pi (T)A_4^\pi (T=0)`$. We therefore use
$$A_4^\rho (T)A_4^\pi (T)A_4^\pi (T=0)\text{for}TT_c.$$
(9)
Therefore from Eq. (8), we obtain the deviation $`\delta _{\mathrm{bare}}`$ as
$$\delta _{\mathrm{bare}}=1V_{\pi ,\mathrm{bare}}^2=\frac{1}{G_0}\frac{32}{105}\pi ^4\frac{T^6}{\mathrm{\Lambda }^6}\left[A_4^\pi +A_4^\rho \right].$$
(10)
This is the matching condition to be used for determining the value of the bare pion velocity near the critical temperature. Let us make a rough estimate of $`\delta _{\mathrm{bare}}`$. For the range of matching scale $`(\mathrm{\Lambda }=0.81.1\text{GeV})`$, that of QCD scale $`(\mathrm{\Lambda }_{QCD}=0.300.45\text{GeV})`$ and critical temperature $`(T_c=0.150.20\text{GeV})`$, we get
$$\delta _{\mathrm{bare}}(T_c)=0.00610.29.$$
(11)
Thus we obtain the $`bare`$ pion velocity as $`V_{\pi ,\mathrm{bare}}(T_c)=0.830.99`$.
We next consider the quantum and hadronic thermal corrections to the parametric pion velocity. It was proven in Ref. that the pion velocity is protected from renormalization by the VM. In the following, we show that this can be understood in terms of chiral partners: Away from $`T_c`$, the pion velocity receives hadronic thermal correction of the form :
$`v_\pi ^2(T)`$ $``$ $`V_\pi ^2N_f{\displaystyle \frac{2\pi ^2}{15}}{\displaystyle \frac{T^4}{(F_\pi ^t)^2M_\rho ^2}}`$ (12)
$`\text{for}T<T_c,`$
where the contribution of the massive $`\sigma `$ (i.e., the longitudinal mode of massive vector meson) is suppressed by the Boltzmann factor $`\mathrm{exp}[M_\rho /T]`$, and then only the pion loop contributes to the pion velocity. On the other hand, when we approach the critical temperature, the vector meson mass goes to zero due to the VM. Thus $`\mathrm{exp}[M_\rho /T]`$ is no longer the suppression factor. As a result, the hadronic correction in the pion velocity is absent due to the exact cancelation between the contribution of pion and that of its chiral partner $`\sigma `$. Similarly the quantum correction generated from the pion loop is exactly canceled by that from the $`\sigma `$ loop. Accordingly we conclude
$$v_\pi (T)=V_{\pi ,\mathrm{bare}}(T)\text{for}TT_c,$$
(13)
i.e., the pion velocity in the limit $`TT_c`$ receives neither hadronic nor quantum corrections due to the protection by the VM. This implies that $`(g_L,a^t,a^s,V_\pi )=(0,1,1,\text{any})`$ forms a fixed line for four RGEs of $`g_L,a^t,a^s`$ and $`V_\pi `$. When a point on this fixed line is selected through the matching procedure as explained in Ref. , that is to say when the value of $`V_{\pi ,\mathrm{bare}}`$ is fixed, the present result implies that the point does not move in a subspace of the parameters. Approaching the chiral symmetry restoration point, the physical pion velocity itself will flow into the fixed point. Finally thanks to the non-renormalization property, i.e., $`v_\pi (T_c)=V_{\pi ,\mathrm{bare}}(T_c)`$ given in Eq. (13), we arrive at the physical pion velocity at the chiral restoration:
$$v_\pi (T_c)=0.830.99,$$
(14)
close to the speed of light.
## 4 Conclusion
In this note, we studied, using the ChPT with HLS/VM, the pion velocity near the critical temperature. We exploited the non-renormalization property of the pion velocity to assure that it suffices to compute the $`bare`$ pion velocity at the matching scale to arrive at the $`physical`$ pion velocity at the chiral restoration temperature. We derived the matching condition on the $`bare`$ pion velocity and found that the pion velocity near $`T_c`$ is close to the speed of light, $`v_\pi ^{(\mathrm{VM})}(T)=0.830.99`$ and definitely far from the zero velocity. This is in a stark contrast to the result obtained from the chiral theory , wherein only the pion figures as the relevant degree of freedom near $`T_c`$, namely, $`v_\pi (T_c)=0`$. The drastic difference between the two approaches is not difficult to understand. In the HLS/VM approach, the $`\rho `$ meson becomes light as $`T_c`$ is approached from below and plays as important a role as the pion does. The effect of the massless vector meson cannot be approximated in chiral models by local operators in pseudoscalar fields.
By fitting the pion spectra observed by STAR at RHIC in terms of an optical potential that incorporates the dispersion relation of low-energy pions in nuclear matter, Cramer et al. deduced the in-medium pion velocity $`v_\pi (T)=0.65`$. The authors interpreted this result as an evidence for the pions being emitted from the chiral-symmetry restored phase. No error bars have been assigned to this value, so it is difficult to make a clear-cut assessment of what that value implies. It seems however difficult to identify it with what has been predicted by sigma models, namely, $`v_\pi =0`$. On the other hand, given that our prediction (14) based on HLS/VM is for the chiral limit, it seems reasonable to expect that the account of the explicit chiral symmetry breaking by quark masses would lower the velocity from (14), making it closer to the observed value, $`v_\pi =0.65`$. Whether or not it signals the chirally restored phase as interpreted in is not clear. However as argued in , there is an indication for a massless pion (in the same multplet with a scalar $`\sigma `$) just above $`T_c`$ with its velocity close to 1, it seems logical that $`v_\pi `$ stays near 1 – rather than near zero – as $`T_c`$ is approached from below as well as from above.
## Acknowledgments
We are grateful for discussions with Gerry Brown, Youngman Kim, Koichi Yamawaki. This work is supported in part by the JSPS Grant-in-Aid for Scientific Research (c) (2) 16540241, and by the 21st Century COE Program of Nagoya University provided by Japan Society for the Promotion of Science (15COEG01). |
warning/0506/astro-ph0506759.html | ar5iv | text | # Flaring and self-shadowed disks around Herbig Ae stars: simulations for 10 𝜇m interferometers
## 1 Introduction
Herbig Ae/Be stars (HAEBEs, Herbig 1960, for a more recent review see Natta et al. 2000) are intermediate mass pre-main-sequence stars, surrounded by material which is left from the star formation process. A sub-group of mostly late B and A-F type HAEBE stars (hereafter HAEs) show little or no optical extinction, and usually have low mass accretion rates as derived from radio analysis (Skinner et al. 1993) and the lack of significant veiling in optical spectra. There is ample evidence that the circumstellar material responsible for the large infrared excesses of these stars is located in a circumstellar disk (e.g. Mannings & Sargent 1997; Grady et al. 2001; Augereau et al. 2001; Eisner et al. 2003). Vink et al. (2002) show that the gaseous component has a disk-like geometry on scales of less than 0.1 AU.
Whereas the presence of these circumstellar disks seems firmly established, the structure of the disks is a matter of debate. Kenyon & Hartmann (1987) developed “flaring” disk models for T-Tauri stars, in which $`H/R`$ (the ratio of the disk surface height to the distance to the star) increases with increasing distance to the central star. The flaring disk model was refined by Chiang & Goldreich (1997, hereafter CG97) who introduced an optically thin surface layer responsible for the infrared emission features generally seen in circumstellar disks. Natta et al. (2001) and Dullemond et al. (2001) reconsider the CG97 model in the context of HAe stars, proposing that the innermost region of the disk has an increased scale height: the *“puffed-up inner rim”*. This configuration, which results from hydrostatic equilibrium at the directly irradiated inner rim, naturally explains the near-infrared bump commonly observed in the Spectral Energy Distribution (SED) of HAe systems (Natta et al. 2001).
Meeus et al. (2001) noted that based on the IR SED, HAEs can be divided into two main groups: “group I” sources that have a very strong, rising IR excess peaking around 60 $`\mu `$m, and “group II” sources displaying a more moderate IR excess, lacking the 60 $`\mu `$m bump. It was proposed that group I sources have a “flaring” geometry, allowing the disk to intercept and reprocess stellar radiation out to large stellocentric radii. In the outer disk of group II sources, on the other hand, $`H/R`$ is approximately constant, or decreasing with increasing distance to the star. The inner disk shields the outer disk from direct irradiation by the central star, hence the term “self-shadowed” disk. This substantially reduces the amount of radiation absorbed locally, leading to lower temperatures in the outer disk of a group II source.
Recent 2D modeling by Dullemond (2002), Dullemond & Dominik (2004a, henceforth DD04), has quantitatively confirmed that both flaring and self-shadowed disks are natural solutions of the equation of vertical hydrostatic equilibrium in passive circumstellar disks (see also section 2.4). These models form the basis of the current study. There is ample evidence that group I and group II disks indeed have a flaring and self-shadowed geometry, respectively (e.g. Grady et al. 2004, Dullemond & Dominik 2004b, Leinert et al. 2004). However, as this is not an observational study, we will consistently refer to the models as flaring/self-shadowed, rather than group I/II (which is by definition an SED classification).
With the advent of long ($`10^2`$ m) baseline infrared interferometry using large apertertures, it has now become possible to observe HAe disks in the thermal infrared with a spatial resolution of order $`10^2`$ arcsec. At the present, the number of baselines will be limited, and only interferometric amplitudes (no phases) are available. True aperture synthetis imaging of disks is therefore not (yet) possible. The interpretation of the measured visibility amplitudes, which contain information about the geometry of the disks, requires the use of disk models.
In the near-infrared Herbig Ae/Be star disks have been observed with long-baseline interferometers since a number of years (Tuthill et al. 2001; Millan-Gabet et al. 2001; Eisner et al. 2003, 2004, 2005). Up to recently these measurements were evaluated using extremely simplified models: Gaussian blobs, rings, ellipses etc. Such simple models made it possible to get a handle on the typical size and inclination of the emitting source, but did not go much further. For most sources the typical sizes were found to be in rough agreement with those predicted by the inner rim models, but at high accretion rates the observations deviate from predictions. This is explained by Akeson et al. (2005) as due to the emission from accretion inward of the inner dust rim, and by Monnier et al. (2005) as due to the protection of dust by optically thick gas, allowing the dust to survive closer to the star. In these, and other, recent papers the modeling of the data already starts to go well beyond the simple ring/ellipse models, using actual multi-dimensional radiative transfer calculations in the case of Akeson et al., and detailed accretion disk structure models in the case of Lachaume et al. (2003). In particular with the new phase-closure capabilities in the near-infrared at the IOTA and VLTI/AMBER interferometers such more advanced models are clearly of great use.
With the mid-infrared interferometry capabilities of the MIDI instrument on the VLTI it is now possible to study the structure of the disk at slightly larger scales than the inner rim. This is the region in which the self-shadowed and flaring disks would most clearly be distinguishable. A first set of measurements was published by Leinert et al. (2004), and a first tentative correlation between the shape of the disk (flaring/self-shadowed) and the visibility was found. Mid-infrared interferometry also has the interesting capability of measureming mineralogical properties of the dust as a function of stellocentric radius. First measurements of this kind (van Boekel et al. 2004a) have revealed the strong radius-dependence of the crystallinity of dust, as predicted by theoretical models. In the present paper, however, we will be mostly concerned with the first aspect of mid-infrared interferometry: measuring the geometry of the disk.
Based on the 2D disk models of Dullemond & Dominik (2004a) as well as the simpler models of Chiang & Goldreich (1997) we present calculations of the interferometric visibilities of HAe disks, to investigate if it is possible to distinguish between the various disk geometries predicted by these theoretical models. Since the mid-infrared probes the structures at somewhat larger scale than the inner rim (from 1 AU out to about 20 AU), this wavelength regime is more suited to our aims than the near-infrared. The MIDI instrument is, so far, the only instrument capable of doing such measurements for Herbig Ae/Be stars, so in our analysis we focus on the typical baselines and properties of the VLTI.
## 2 Modeling method
The goal of this study is to predict and compare interferometric visibilities of various disk models. Synthetic disk images are made using a ray tracing algorithm, where special care is taken to ensure all spatial scales in the disk are sufficiently resolved. Interferometric visibilities are calculated by Fourier transforming the images.
### 2.1 The inner rim
In the disk models considered here, the bright inner rim is treated in a highly simplified fashion: it is a sharp, “vertical wall”. As a consequence, when such a model is viewed not pole-on, the flux from the near side of the inner rim is strongly suppressed since the hot, irradiated rim surface is occulted by the cooler parts immediately outside the rim surface. The far side of the bright inner rim is in full view and is responsible for essentially all of the near-infrared excess. In this configuration, the total amount of near-infrared emission depends strongly on the disk inclination, suggesting that the observed strength of the “3 micron bump” in the SED is a measure of the latter. Observations of Herbig Ae stars, however, show that they all have rather similar near-infrared excesses, irrespective of their inclination (Natta et al. 2000, Dominik et al. 2003).
This indicates that the appearance of the inner rim is more smooth than the “vertical wall” used here. The processes that determine the shape of the inner rim are currently not yet understood. Isella & Natta (2005) recently showed that the dependence of the evaporation temperature on pressure naturally leads to a rounded-off inner rim. When such a disk is viewed at an inclination, both the near and the far side of the inner rim will be bright (although still the far side will be brighter). The bright inner rim will look like an inclined ring on the sky, rather than the “half ring” one obtains using the vertical wall model.
Realistic radiative transfer modeling of a rounded-off inner rim introduces various numerical complications. To avoid these difficulties we adopt the simplified vertical rim structure used in Dullemond & Dominik (2004a). For a slightly off-polar inclination we artificially circularize the disk emission to circumvent the near-side/far-side asymmetry of the rim. In this way we mimic the rounded-off shape of the rim without having to confront the numerical complexities of radiative transfer in extreme optical depth rounded-off rims. While the spatial resolution of current 10 $`\mu `$m interferometers is just sufficient to measure the diameter of the inner rim, observations at higher resolution are required to study details of the rim structure. We therefore believe that using this simplified approach is justified for our current purposes. To first approximation, inclination can be included by scaling the calculated spatial frequencies (or interferometric baselines) by a factor 1/cos($`\theta `$) along the minor axis of an inclined disk, where $`\theta `$ is the inclination of the disk. At high inclinations, this approximation brakes down.
### 2.2 Interferometry
Rather than images, an interferometer produces an interference signal called the interferometric “visibility” ($`V`$), which is the spatial coherence function of the intensity distribution of the source. The visibility is related to the intensity distribution through the van-Cittert-Zernike theorem, which states that the visibility is the Fourier transform of the intensity distribution of the source. For an introduction to long-baseline interferometry, we refer to Lawson (2000).
### 2.3 Spectrally resolved visibilities
Traditionally, visibility curves are represented as a function of the projected baseline $`B`$<sup>1</sup><sup>1</sup>1The spatial frequency of the observation is: $`k=B/\lambda `$, where $`\lambda `$ is the observing wavelength. The units of $`k`$ are cycles/radian if $`\lambda `$ and $`B`$ have the same units, at a specific wavelength (a $`V(B)`$ curve). Such a curve represents a number of visibility measurements at different projected baselines, which usually requires the use of multiple telescope pairs and/or moving telescopes.
A new possibility in the 10 $`\mu `$m region is the use of spectral dispersion with large wavelength coverage. When using an instrument that has this capability, one can obtain a whole “visibility curve” in one single measurement. Unlike the common $`V(B)`$ curve, the $`V(\lambda )`$ curve obtained this way holds many visibility values at only one baseline. The spatial resolution of the observation ($`B/\lambda `$) changes by almost a factor of 2 between 7.5 and 14 $`\mu `$m. Detailed modeling is required to interpret $`V(\lambda `$) curves. Most HAe stars show a prominent emission band between 8 and 12 $`\mu `$m, due to silicate dust. The shape of this emission band varies strongly, depending on chemical composition, particle size and lattice structure of the silicate grains. When simulating interferometric visibilities using disk models, one finds that the detailed shape of the visibility curve depends on the opacities used, i.e. on the dust properties. These are different from star to star, and vary within a disk as a function of distance to the star (van Boekel et al. 2004a). The visibilities measured in the silicate emission feature are a mixture of disk structure and mineralogy. Therefore, in order to deduce information on the disk structure, it is preferable to use measurements at wavelengths outside the silicate feature, which in practice means between about 12 and 13.5 $`\mu `$m. At 8 $`\mu `$m, it is also possible to sample the continuum emission, this is however more difficult since here the atmospheric transmission is rather poor.
### 2.4 The DD04 disk model
The disk models used in this work are described in DD04. These are 2D axisymmetric models in which the gas and dust density and temperature are given as a function of radius $`R`$ and polar coordinate $`\mathrm{\Theta }`$. The disk is assumed to be heated only by irradiation by the central star. A 2D continuum radiative transfer code is used to compute the entire temperature structure of the disk. The vertical density structure, for a given radial surface density distribution $`\mathrm{\Sigma }(R)`$, is computed by demanding vertical hydrostatic equilibrium. In this way the disk has a self-consistent temperature and density structure, from which images and SEDs can be computed using a ray-tracer. For this work we use the following stellar parameters: $`M_{}=2.5`$M, $`R_{}=2`$R and $`T_{}=\mathrm{10\hspace{0.17em}000}`$ K, which amounts to a stellar luminosity of $`L_{}=36`$L. All disk models in this work have a disk mass of 0.1 M, a gas-to-dust ratio of 100, a surface density distribution of $`\mathrm{\Sigma }R^{1.5}`$, and an outer radius of 200 AU. The DD04 models have an inner disk radius which is calculated self-consistently assuming an optically thick inner rim. The location of the inner rim is set by the dust evaporation temperature, which is about 1500 K for silicate dust. The CG97 model has an inner radius of 0.21 AU, which corresponds to the radius where the black-body temperature is 1500 K. For the dust opacities we use a simple model consisting only of small silicate grains (Laor & Draine 1993).
If the disk is optically thick enough, the disk has a flaring shape (DD04 and Dullemond 2002). When the optical depth is decreased, a flaring disk can turn into a self-shadowed disk and the SED changes from a group I to a group II shape. The flaring and self shadowed models shown here are the BL1 and the BL4 model from DD04. Both disks have a mass of 0.1 M. In the BL1 model all the dust mass is in 0.1 $`\mu `$m silicate grains. In the BL4 model 99.9% of the mass has been converted into 2 mm size grains located in the midplane, while only 0.1% remains in small 0.1 $`\mu `$m grains, thus strongly lowering the opacity of the disk. We stress however, that the flaring vs. self-shadowed behaviour of the disk depends on high vs. lower *optical depth*, and that dust coagulation is a possible mechanism to achieve lower optical depths.
In the proposed scheme, the outer disk of a self-shadowed source is shielded from direct stellar irradiation by its own inner disk. However, the outer disk receives near-IR radiation emitted by the hot innermost disk regions, and optical/UV radiation which is scattered by the diffuse inner disk atmosphere. Therefore, the temperature and scale height in a self-shadowed disk are still significantly larger than zero. Note that in a flaring disk, there is also a region just outward of the inner rim that is shielded from direct stellar radiation. Contrary to a self-shadowed disk, a flaring disk emerges from the shadow cast by the puffed up inner rim, at distances of a few AU from the star.
## 3 Results
The geometry of the disk models investigated in this work is schematically represented in Fig. 1. The outer disk geometries in the CG97 and DD04 flaring models are very similar, however the inner disk region is different. The DD04 flaring model has a puffed up inner rim, and an intermediate shadowed region (shaded light). The DD04 self-shadowed model has an inner disk structure that is very similar to the DD04 flaring model. The outer disk in the self-shadowed model never rises above the shadow cast by the puffed up inner rim, but is irradiated by the hot inner disk regions.
### 3.1 Spectral energy distributions
The emerging spectral energy distributions of the disk models are shown in Fig. 2. The infrared excess emission at far infrared (FIR) wavelengths ($``$60-100 $`\mu `$m) is clearly stronger in the flaring disk models than in the self-shadowed model. The FIR excess is significantly stronger in the DD04 flaring model than in the CG97 (also flaring) model. This can be traced to the simplifications made in the CG97/DDN01 models which do not take properly into account various 3D radiative transfer effects. In particular, these models do not account for the moderate ‘boosting’ of radiation toward the polar axis to compensate for the occultation in equatorial directions by the disk’s own flaring outer regions. In the DD04 models these effects are consistently taken into account by virtue of the full multi-dimensional radiative transfer treatment used in those models. The near-infrared excess around 2-3 $`\mu `$m which is very prominent in the DD04 model SEDs, is much less pronounced in the CG97 model. In the 10 $`\mu `$m region, the CG97 and DD04 flaring model SEDs are virtually identical in spectral shape, though the DD04 flaring model has a somewhat higher absolute flux level. The DD04 self-shadowed disk is fainter than the flaring models, and has a bluer continuum slope in the 10 $`\mu `$m region.
### 3.2 Radial intensity profiles
Fig. 3 shows the radial intensity profiles at 12.6 $`\mu `$m of a DD04 flaring and self-shadowed disk model, and the CG97 model. Both the DD04 flaring and self-shadowed model essentially exhibit three regimes:
* the “bright puffed up inner rim”, that causes a ringlike emission, contributing mainly between 0.5 and 0.8 AU from the star.
* a region just behind the inner rim (as seen from the star), where the dust temperatures are much lower than in the inner rim. From this “intermediate shadowed region”, relatively little radiation emerges (as can be seen in the cumulative flux distributions of the DD04 models in the lower panel of Fig. 3, which are nearly constant in this region).
* the “outer disk region”, whose main flux contribution arises between 3 and 20 AU from the central star. In a flaring disk model, this outer disk region is directly irradiated by the central star.
The “bright puffed up inner rim”, “intermediate shadowed region” and “outer disk region” are of course just different parts of the same physical structure, and the distinction made here serves merely to help the reader develop a qualitative understanding of how such geometries translate into interferometric visibilities. In Fig. 3 we have indicated the three regions discussed above. For the DD04 flaring model, the outer radius of the intermediate shadowed region can be well defined to be between 2 and 3 AU, where the slope of the cumulative flux distribution clearly increases. In the DD04 self-shadowed model, this radius is less clearly defined but evidently somewhat larger than in the DD04 flaring model. The contribution of the bright inner rim emission to the total system flux depends strongly on wavelength (for the self-shadowed model, the inner rim contributes more than 90% to the total flux at 6 $`\mu `$m, about 60% at 8 $`\mu `$m, about 35% at 13 $`\mu `$m and less than 5% at 30 $`\mu `$m), and is always higher in a self-shadowed model than in a flaring model harboring the same central star.
The CG97 model has, per definition, a flaring disk structure. Contrary to the DD04 flaring model however, it does *not* have a bright puffed up inner rim, and consequently it lacks an intermediate shadowed region.
### 3.3 Visibility curves
#### 3.3.1 Visibility as a function of baseline
Fig. 4 shows the predicted “classical” visibility curves $`V(B)`$ of the considered disk models, at several wavelengths. From top to bottom we show the predictions for the CG97 model, the DD04 flaring model and the DD04 self-shadowed model. To develop an intuitive understanding of how the characteristics of the emerging intensity distributions of the various models are reflected in their visibility curves, we will discuss the curves in Fig. 4. In this example, the star is put at a distance of 150 pc, typical for nearby Herbig stars.
At a baseline of 0 m all sources are of course unresolved and have a visibility of 1. The CG97 model shows a steady drop in visibility as the baseline is increased. The slope of the visibility curve changes gradually, reflecting that the radial intensity profile shown in Fig. 3 has no strong substructure. In this sense, the CG97 model is “scaleless”.
For the DD04 models, this is different. These models have essentially two scales: the bright inner rim, which emits between 0.5 and 0.8 AU, and the outer disk, which emits most of its flux between about 3 and 20 AU. In between lies the intermediate shadowed region, whence little flux emerges. This general picture is reflected in the visibility curves. Starting at 0 m, and increasing the baseline, we observe a steady drop in visibility as the outer disk gets more and more resolved. Note that at a baseline of 10 m the visibility is already significantly lower than 1, predicting that the largest modern day telescopes might marginally resolve the outer disks in such objects at 10 $`\mu `$m. For the HAe star HD 100546 this has indeed been observed (Liu et al. 2003, van Boekel et al. 2004b). At a baseline of about 30 m, the outer disk is mostly resolved while the inner rim is still essentially unresolved. Therefore, the visibility curves flatten at this point. The visibility level at this baseline (about 15% for the flaring model and 40% for the self-shadowed disk, at 9.8 $`\mu `$m) indicates the fraction of the total system flux that is emitted by the bright inner rim. At longer baselines, the bright inner rim itself becomes resolved by the interferometer, and the visibility gradually goes to its first null. Since the spatial resolution of the interferometer scales inversely with wavelength, whereas the apparent diameter of the inner rim hardly depends on the wavelength, zero visibility is reached first at the shortest wavelengths, and at longer baselines for the longer wavelengths.
#### 3.3.2 Visibility as a function of wavelength
Fig. 5 shows $`V(\lambda )`$ curves for the CG97, DD04 flaring and DD04 self-shadowed model, at a number of different baselines. Each of these curves represents a *single*, dispersed, visibility measurement (the spectral region inaccessible from the ground is shaded grey). Each curve in Fig. 5 can be regarded as a cut through Fig. 4 at a specific baseline, with a much denser wavelength sampling.
The overall trend for all curves is to show the highest visibilities at 8 $`\mu `$m, and lower visibilities at 13 $`\mu `$m. This is because the apparent size of the disks increases with wavelength more rapidly than the interferometric resolution decreases. There is generally a sharp decrease in visibility between 8 and 10 $`\mu `$m. There are two reasons for this. For the models with an inner rim one reason is that the emission from this rim dominates the spectrum below about 8 $`\mu `$m. The emission at 10 $`\mu `$m originates from more extended regions of the disk, resulting naturally in a lower visibility than the 8 $`\mu `$m emission. A second reason for the decline of the visibility between 8 and 10 $`\mu `$m – and for the gentle rise in visibility toward 13 $`\mu `$m in some models – is that the flux in the 10 $`\mu `$m silicate feature originates predominantly from the warm surface layers of the disk, while the flux outside the feature comes from the cooler regions below. The warm dust in the surface layer can radiate in the mid-infrared out to larger radii than the cooler dust in the disk interior. In other words: in the warm surface layers the Wien exponential cut-off in the mid-infrared takes effect at larger radii than in the disk interior. This explanation also holds for the CG97 models, which do not have an inner rim.
In Fig 6 we demonstrate the importance of the silicate resonance for the simulated visibility curves. We show the visibilities of the DD04 flaring model (full curves, see also the middle panel of Fig. 5). To calculate visibility curves of a model without silicate resonances, we removed the 10 and 20 micron silicate features from the opacity table prior to the ray tracing (dotted curves). The large influence of the opacity of the material on the resulting visibilities is evident. The SED of the model without silicate resonances is indicated by the full grey curve in Fig. 2.
The interpretation of the curves in Fig. 5 in terms of disk geometry is not straightforward, for several reasons. First, the emerging intensity distribution of the disk changes with wavelength. Second, the spatial resolution of the interferometer decreases by almost a factor of two between the short and long wavelength edges of the 10 $`\mu `$m atmospheric window (N-band). Third, how the disk intensity distribution (and thus the visibilities) changes with wavelength depends on the opacity of the dust, i.e. on mineralogy (Fig. 6). When we measure a single $`V(\lambda )`$ curve, what we see therefore is a mixture of disk geometry changing with wavelength, instrumental resolution changing with wavelength, and the mineralogy of the source.
As the mineralogy changes from star to star, it is difficult to obtain general diagnostics for disk *geometry* from a $`V(\lambda )`$ curve. The detailed interpretation of such measurements requires a model of each individual star, where both the spectrum (mineralogy) and the disk structure are fitted simultaneously. As a general diagnostic for disk structure it is thus preferable to measure visibilities in the continuum, where the visibilities do not depend strongly on mineralogy. We will come back to this in section 3.4.
However, the spectral capabilities of the interferometer develop their full strength, when one uses observations at several baselines in order to reconstruct an “image” of the disk<sup>2</sup><sup>2</sup>2In practice this will be easiest for disks that are seen not too far from pole on, so their image has a high degree of azimuthal symmetry. This limits the number of baselines needed, and requires the measurement of visibility amplitude only, since all phases will be approximately 0. With spectrally dispersed visibilities, the *spectrum* of the disk is then known immediately at all positions in the disk. It is then possible to study the mineralogy, size distribution and chemical composition of dust grains in the disk surface layer as a function of distance to the central star, providing crucial information about processing and radial mixing in disks.
If the intensity distribution in the disk is strongly centrally peaked like the models discussed in the present paper, the correlated flux obtained at a single, long baseline<sup>3</sup><sup>3</sup>3The visibility is the ratio of the correlated flux and the total flux ($`F_{\mathrm{cor}}/F_{\mathrm{tot}}`$). The correlated flux, or correlated spectrum in the case of a spectrally dispersed measurement, is the quantity an interferometer measures. can be directly interpreted as the spectrum of the innermost regions of the disk (in the correlated spectrum obtained with only one measurement, there is still an unknown spatial term mixed in, that typically introduces a slope in the spectrum. This however has little influence on the derived mineralogy). The outer disk spectrum can then be obtained as a difference between the integrated disk spectrum and the inner disk spectrum. Applying this method to the first spectrally resolved full N-band visibility measurements of HAe stars, it was demonstrated by van Boekel et al. (2004a) that the mineralogy in the disk can vary strongly with distance to the star.
### 3.4 Distinguishing between the various models
The goal of this study is to show how interferometric measurements can be used to distinguish between the various disk models. Clearly, the curves in Fig. 4 are different for the different models. However, one will typically not have continuous measured visibility curves at hand, but rather have samples at a few different baselines. Here we show that it is, at least in principle, possible to distinguish both between the CG97 and DD04 models on one hand, and between the DD04 flaring and DD04 self-shadowed on the other, using measurements at only two appropriately chosen baselines.
The distinction between the CG97 and DD04 models is based on the absence of an intermediate shadowed region in the former. At spatial scales corresponding to the intermediate shadowed region in the DD04 models, little flux emerges. Therefore, the visibility curves are relatively flat at the baselines corresponding to these spatial scales, they show a “plateau” (very prominent in the DD04 flaring model visibility at 12.6 $`\mu `$m in Fig. 4). The CG97 model does not have such specific spatial scales with much reduced emergent intensity, and therefore lacks the plateau in the visibility curve. The visibility curves of CG97 and DD04 models thus have a different *slope* at baselines corresponding to the scale of the intermediate shadowed region. Once this slope difference has been detected, DD04 flaring and DD04 self-shadowded disks can be distinguished by the relative contribution of the bright inner rim to the total system flux, which is much higher for a self-shadowed model. Note that the bright inner rim itself is virtually identical in both models, but the outer disk is much brighter in the flaring disk than in the self-shadowed case. Therefore, the predicted visibilities at our selected baselines are much lower for the flaring model. We recall that deducing properties about the disk structure is best done outside the silicate emission feature, which in practice favours the region between 12 and 13 $`\mu `$m.
A measurement at a specific baseline samples the corresponding angular scale, and the physical scale (in AU) associated with this baseline therefore depends linearly on the distance to the star. It is therefore convenient to introduce the “normalized baseline”
$$x\frac{B}{D}$$
(1)
where $`B`$ is the baseline in meter and $`D`$ is the distance to the star in parsec. Consider the emission of a small part of the disk, arising between $`R`$ and $`R+\mathrm{d}R`$ from the central star, i.e. a ring with angular diameter $`\theta =R/d`$ (where $`d`$ is the distance to the star). The visibility curve of such an *annulus* of emission is the zeroth order Bessel function:
$$V_{\mathrm{an}}(B)=J_0(\frac{\theta B}{\lambda })$$
(2)
where $`B`$ is the interferometric baseline, and $`\lambda `$ is the wavelength of observation. The visibility reaches the first null at a baseline of about
$$B_0158\frac{\lambda _{\mu \mathrm{m}}}{\theta _{\mathrm{mas}}}[\mathrm{m}]$$
(3)
where for convenience the wavelength and annulus angular diameter have been expressed in $`\mu `$m and milli-arcseconds, respectively. Our goal is to detect the effect of the intermediate shadowed region, i.e. the very low flux contribution from annuli between about 0.8 and 3 AU. To estimate which baseline is most sensitive to emission from an annulus with diameter $`\theta `$, let us take the baseline where the visibility has half its maximum value, $`V_{\mathrm{an}}=0.5`$:
$$B_{0.5}100\frac{\lambda _{\mu \mathrm{m}}}{\theta _{\mathrm{mas}}}[\mathrm{m}]$$
(4)
For an annulus of 2 AU radius, this corresponds to a normalized baseline of $`x=\lambda _{\mu \mathrm{m}}/40`$. Let us take the region between 1 and 2 AU from the star as characteristic for the intermediate shadowed region.
We find that in order to most clearly separate the specific DD04 models (with an intermediate shadowed region) used in this work from the CG97 models (without an intermediate shadowed region), the best choice for $`x`$ is:
$$\begin{array}{ccc}\hfill x_1& =& 0.038\left(\frac{\lambda }{10\mu \mathrm{m}}\right)\left(\frac{L_{}}{L_{}}\right)^{\frac{1}{2}}\left[\frac{\mathrm{m}}{\mathrm{pc}}\right]\hfill \\ \hfill x_2& =& 2x_1\hfill \end{array}$$
(5)
To allow the above formula to be applied to stars of different luminosity we have applied a simple scaling of $`x`$ with the square root of the luminosity. The visibility is now sampled at the baselines corresponding to $`x_1`$ and $`x_2`$ (for our 36 $`L_{}`$ star at 150 pc and at a wavelength of 12.6 $`\mu `$m these are 43 and 86 m, respectively, see also Fig 4), yielding visibilities $`V_1`$ and $`V_2`$. In Fig. 7 we plot the on the vertical axis the visibility measured on the longest of the two baselines ($`V_2`$). On the horizontal axis we plot the deduced slope of the visibility curve between $`x_1`$ and $`x_2`$. In addition to the CG97 model with an inner radius of 0.21 AU we plot CG97 models with inner radii of 0.02 and 0.5 AU, to illustrate the behaviour with varying inner radius. From this figure we see that:
* The CG97 models have a much steeper slope than the DD04 models in this baseline regime.
* The DD04 flaring model has a much lower visibility at the long baseline than the DD04 self-shadowed model (and on the short baseline as well).
We can conclude therefore, that it is possible to discriminate between the various models, with a very limited number of measurements. In the upper left corner of Fig. 7 we indicate the uncertainty in the measured location of a source in this diagram due to limited precision of the visibility measurments. We have here assumed a visibility accuracy of 1 %, which is what MIDI aims to achieve. The difference in the position of the models with and without a silicate feature (plotted with black and grey symbols, respectively) can be regarded as the extreme case for the uncertainty arising from mineralogy, since for each individual source, we *know* what the silicate feature looks like.
This analysis has been done in an idealized world where we have both assumed that the disks are pole-on, and that our models are a good representation of the true disk geometry. The optimum choice for $`x_1`$ and $`x_2`$ depends on the geometry of the disk and may therefore in reality be somewhat different than the values given in our recipe (equation LABEL:eq:optimum\_baseline). If the outer disk is smaller than we predict, the value of $`x_1`$ should be increased. If the bright inner rim is located at radii somewhat larger than predicted, the value of $`x_2`$ should be decreased. There is however evidence that the bright inner rim is located at radii somewhat smaller than predicted in our models (Eisner et al. 2003), and therefore this is not likely a reason for concern. In practice, measurements at more than two baselines ($``$5) are probably needed to unambiguously establish the nature of the sources.
## 4 Conclusions
We have presented model calculations of interferometric visibilities of circumstellar disks around Herbig Ae stars. We compare predictions for disks with (DD04) and without a bright inner rim (CG97). We show that it is possible to distinguish between both model possibilities using a small number of interferometric measurements in the 10 $`\mu `$m atmospheric window. Such measurements also allow to distinguish between flaring and self-shadowed disk models. This allows testing of the hypothesis that group I and group II sources correspond to flaring and self-shadowed disks, as is suggested by their spectral energy distributions.
###### Acknowledgements.
R. Lachaume and L. B. F. M. Waters are gratefully acknowleged for constructive comments on the manuscript. |
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