id
stringlengths 27
33
| source
stringclasses 1
value | format
stringclasses 1
value | text
stringlengths 13
1.81M
|
|---|---|---|---|
warning/0506/hep-ex0506016.html | ar5iv | text | # 𝑊 Boson Production and Mass at the Tevatron
## 1 Introduction
The properties of the $`Z`$ boson have been measured to very high precision at LEP. Naturally one wants to match this precision for the charged carriers of the electroweak interaction. Over the next few years the Tevatron is the only accelerator which can produce $`W`$ bosons. Measuring the properties of the $`W`$ boson to a very high precision is an important test of the Standard Model. From the measured $`W`$ cross section, one can infer an indirect measurement of the $`W`$ width and lepton universality. Since at the Tevatron the $`W`$ bosons are produced through quark anti-quark annihilation, a significant uncertainty for all direct electroweak measurements comes from the knowledge of the parton distributions inside the proton. The probability of finding a parton carrying a momentum fraction $`x`$ within the incoming proton is expressed in the parton distribution function (PDF). The measurement of the $`W`$ charge asymmetry provides important input on the ratio of the $`u`$ and $`d`$ quark components of the PDF and will help to further constrain parton distribution functions. The $`W`$ boson mass serves as a test of the Standard Model, but through radiative corrections is also sensitive to hypothetical new particles. Together with a precise measurement of the top quark mass, the $`W`$ boson mass constrains the mass of the Higgs boson, which has not yet been observed experimentally.
Both CDF and DØ are multi-purpose detectors. They consist of tracking systems surrounded by calorimeter and muon identification systems. CDF’s tracking system consists of a wire drift chamber (the Central Outer Tracker) and a 7-layer silicon microstrip vertex detector (SVXII) immersed in a 1.4 T magnetic field. A lead (iron) scintillator sampling calorimeter is used for measuring electromagnetic (hadronic) showers. DØ employs a silicon microstrip tracker (SMT) and a central fiber tracker (CFT), both located in a 2 T magnetic field. The sampling calorimeter consists of liquid argon and uranium.
Since the hadronic decay of the $`W`$ boson has an extremely large background originating from strongly interacting processes, CDF and DØ use the clean leptonic decays to study the $`W`$ boson. The signature is a high energy lepton with large missing transverse momentum originating from the neutrino, which does not interact with the detector. The momentum balance in the direction of the beam is unconstrained and as a result, the $`W`$ events are studied in the plane transverse to the beam. A typically used quantity is the transverse mass:
$$M_T=\sqrt{2p_T^lp_T^\nu (1cos(\mathrm{\Delta }\varphi ))},$$
(1)
which is similar to the invariant mass, just in the two transverse dimensions. If not otherwise stated, we restrict the lepton identification to the well instrumented central region of $`|\eta |<1`$.
$`Z`$ boson events are identified by two high energy leptons. These events have very low background.
## 2 Inclusive $`p\overline{p}W/Z+X`$ Cross Section Measurements
$`W`$ and $`Z`$ bosons are identified by their leptonic decays to electrons, muons and taus, from which the total rates $`\sigma \times `$Br($`Wl\nu `$) and $`\sigma \times `$Br($`Zll`$) are obtained. The cross section times branching ratio is calculated as follows:
$$\sigma \times Br(p\overline{p}W/Zll)=\frac{N_{cand}N_{bkg}}{AϵL}.$$
(2)
The $`W`$ and $`Z`$ boson cross sections have been measured by CDF with different datasets in different sub-detectors. Figure 1
shows a summary of the CDF and DØ cross section measurements in all leptonic decay modes. All measurements show good agreement with NNLO calculations, represented by the vertical band.
### 2.1 Lepton Universality in $`W`$ Decays
Lepton universality in $`W`$ decays can be tested by extracting the ratio of the electroweak couplings $`g_\mu /g_e`$ and $`g_\tau /g_e`$ from the measured ratio of $`Wl\nu `$ cross sections. The $`Wl\nu `$ couplings are related to the measured production cross section ratio $`U`$ as follows:
$$U=\frac{\sigma \times Br(Wl\nu )}{\sigma \times Br(We\nu )}=\frac{\mathrm{\Gamma }(Wl\nu )}{\mathrm{\Gamma }(We\nu )}=\frac{g_l^2}{g_e^2}$$
(3)
In this ratio, important systematic uncertainties cancel. The results obtained are:
$$\frac{g_\mu }{g_e}=0.998\pm 0.012$$
(4)
$$\frac{g_\tau }{g_e}=0.99\pm 0.04$$
(5)
where the largest systematic uncertainty comes from event selection efficiencies. Since these efficiencies are measured using the $`Zll`$ sample, the uncertainty will decrease as more $`Z`$ bosons are collected.
### 2.2 Indirect $`W`$ Width Determination
The ratio $`R`$ of the cross section measurements for $`W`$ and $`Z`$ bosons can be used to extract the total width of the $`W`$ boson. $`R`$ can be expressed as:
$$R=\frac{\sigma (p\overline{p}W)}{\sigma (p\overline{p}Z)}\frac{\mathrm{\Gamma }(Wl\nu )}{\mathrm{\Gamma }(Zll)}\frac{\mathrm{\Gamma }(Z)}{\mathrm{\Gamma }(W)}.$$
(6)
Using the very precise measurement of $`\mathrm{\Gamma }(Zll)`$/$`\mathrm{\Gamma }(Z)`$ from LEP and NNLO calculations of $`\sigma (p\overline{p}W)/\sigma (p\overline{p}Z)`$, together with the Standard Model prediction of $`\mathrm{\Gamma }(Wl\nu )`$ one can extract $`\mathrm{\Gamma }(W)`$ from equation 6.
Table 1 shows the values from CDF for two different datasets, together with the current world average (not including these measurements). The indirect width measurements show good agreement and have competitive uncertainties.
## 3 Direct $`W`$ Width Measurement
DØ has measured the $`W`$ boson width directly in the electron decay channel. The measurement uses an integrated luminosity of 177 pb<sup>-1</sup>. The width is determined by normalizing the predicted signal and background transverse mass distribution in the region of 50 GeV$`<`$$`M_T`$$`<`$100 GeV and then fitting the predicted shape to the candidate events in the tail region 100 GeV$`<`$$`M_T`$$`<`$200 GeV,
which is most sensitive to the width. Figure 2 shows the transverse mass distribution. The measurement yields $`\mathrm{\Gamma }_W=`$2011$`\pm `$93(stat)$`\pm `$107(syst) MeV, which is in good agreement with the world average, an improvement over the DØ Run 1 measurement, and competitive to the CDF Run 1 measurements in the muon and electron decay channels.
## 4 $`W`$ Charge Asymmetry
The $`W`$ bosons at the Tevatron are produced predominantly through annihilation of valence $`u`$ ($`d`$) and anti-$`d`$ (anti-$`u`$) quarks inside the proton and anti-protons for $`W^+`$ ($`W^{}`$) production. Since $`u`$ quarks carry, on average, a higher fraction of the proton momentum than $`d`$ quarks, a $`W^+`$ tends to be boosted in the proton direction, while a $`W^{}`$ is boosted in the anti-proton direction. This results in a non-zero forward-backward charge asymmetry, defined as:
$$A(y_W)=\frac{d\sigma (W^+)/dy_Wd\sigma (W^{})/dy_W}{d\sigma (W^+)/dy_W+d\sigma (W^{})/dy_W},$$
(7)
where $`y_W`$ is the rapidity of the $`W`$ bosons and $`d\sigma (W^\pm )/dy_W`$ is the differential cross section for $`W^+`$ or $`W^{}`$ boson production. However, because the $`p_Z`$ of the neutrino is unmeasured, $`y_W`$ cannot be directly determined, and we instead measure:
$$A(\eta _e)=\frac{d\sigma (e^+)/d\eta _ed\sigma (e^{})/d\eta _e}{d\sigma (e^+)/d\eta _e+d\sigma (e^{})/d\eta _e},$$
(8)
where $`\eta _e`$ is the electron pseudorapidity. Therefore, the observed asymmetry is a convolution of the aforementioned charge asymmetry and the Standard Model $`VA`$ couplings describing the $`We\nu `$ decays. A measurement of the forward-backward charge asymmetry is sensitive to the ratio of the $`u`$ and $`d`$ quark components of parton distribution functions. CDF has measured this asymmetry in the electron channel up to a pseudorapidity of $`|\eta |`$$`<`$2.5 using 170 pb<sup>-1</sup>.
Figure 3 shows the measured asymmetry corrected for the effects of charge misidentification and background contributions. The predictions using the latest CTEQ and MRST PDFs are overlaid. This measurement will provide important input for the next generation of PDFs.
## 5 $`W`$ Mass
Since its discovery in 1983, the $`W`$ boson mass has been measured with increasing precision. From an initial uncertainty of 5 GeV, the uncertainty of the $`W`$ mass has been reduced to 42 MeV from the LEP experiments and to 59 MeV from the Tevatron experiments. CDF has analyzed the first 200 pb<sup>-1</sup> of Run 2 data and estimated the corresponding $`W`$ boson mass uncertainty in the electron and muon decay channels. The uncertainty includes contributions from statistics, production and decay modeling, lepton energy scale and resolution, hadronic recoil and resolution, and backgrounds.
There are two important aspects to a precision $`W`$ mass measurement: Calibration of the detector to the highest possible precision, and simulation of the transverse mass spectrum, which cannot be predicted analytically. The simulation includes the production modeling and detector effects and produces transverse mass templates for a range of $`W`$ boson masses. Since backgrounds contaminate the signal, they are included in the templates. The $`W`$ mass is extracted from a maximum likelihood fit to the transverse mass spectrum.
### 5.1 Production Model
The uncertainty in the modeling of the $`W`$ boson production and decay results from parton distribution functions, QED radiative corrections, the transverse momentum of the $`W`$ boson and the $`W`$ boson width.
The parton distribution functions affect the $`W`$ mass through the limited acceptance of the detector for the $`W`$ decay lepton. The uncertainty has been determined using the set of 40 CTEQ6 PDFs, which explore the uncertainty on the 20 orthogonal eigenvector directions in parameter space. Each eigenvector direction corresponds to some linear combination of PDF parameters. The resulting uncertainty is $`\mathrm{\Delta }M_W(e,\mu )=\pm `$15 MeV. A cross check using the latest MRST PDF falls within this estimate.
The dominant higher-order QED effect on the $`W`$ boson mass is photon radiation off the final-state charged lepton. Additional QED uncertainties arise from multi-photon radiation, initial state radiation and radiation from interference terms, none of which are included in the simulation used to extract the $`W`$ boson mass. The uncertainty from QED corrections is $`\mathrm{\Delta }M_W(\mu )=\pm `$ 20 MeV in the muon channel and $`\mathrm{\Delta }M_W(e)=\pm 15`$ MeV in the electron channel.
The initial-state QCD radiation in vector boson production is constrained by a phenomenological parametrization of the $`Z`$ boson $`p_T`$ measurement from the previous collider run. The parameters are used for the modeling of the $`W`$ $`p_T`$ distribution and their uncertainties result in $`\mathrm{\Delta }M_W(e,\mu )=\pm `$13 MeV.
The uncertainty on the $`W`$ boson width affects the falling Jacobian edge and leads to $`\mathrm{\Delta }M_W(e,\mu )=\pm `$12 MeV.
### 5.2 Lepton Momentum/Energy Scale and Resolution
The lepton momentum measurement is based fundamentally on the calibration of the tracking wire chamber (COT). After the calibration of the track momentum and resolution using the muon decays of precisely known resonances, the energy scale of the electromagnetic calorimeter is calibrated using the ratio of calorimeter energy to track momentum ($`E/p`$) of electrons.
The quarkonium resonance decays $`J/\mathrm{\Psi }\mu \mu `$ and $`\mathrm{{\rm Y}}(1S)\mu \mu `$ are used to set the momentum scale (Figure 4).
The passive material in the simulation is tuned such that the reconstructed $`J/\mathrm{\Psi }`$ mass is constant as a function of mean track curvature. The measured momentum scale is the mean of the individual $`J/\mathrm{\Psi }`$ and $`\mathrm{{\rm Y}}(1S)`$ scales. The systematic uncertainty is taken as half the difference between the extracted scales which results in $`\mathrm{\Delta }M_W(e,\mu )=\pm `$15 MeV.
The track resolution is parametrized in the simulation by the individual hit resolution and by the hit multiplicity on the track. Muons from decays of $`Z`$ bosons are used to determine the resolution. The resulting uncertainty corresponds to $`\mathrm{\Delta }M_W(e,\mu )=\pm `$12 MeV. An additional uncertainty of $`\mathrm{\Delta }M_W(e,\mu )=\pm `$20 MeV is assigned for tracking chamber misalignments.
The $`E/p`$ distribution of electrons from $`W`$ boson decays is used to calibrate the electromagnetic energy scale of the calorimeter (Figure 5). The statistical uncertainty and the uncertainty from the momentum scale results in $`\mathrm{\Delta }M_W(e)=\pm `$35 MeV.
Additional energy scale uncertainties arise from the calibration of the detector passive material and from the calorimeter non-linearity. The passive material was measured during detector construction and is tuned using electrons from photon conversions. A final tuning uses the tail of the $`E/p`$ distribution which is sensitive to the amount of material modeled in the simulation. The uncertainty on the passive material results in $`\mathrm{\Delta }M_W(e)=\pm `$55 MeV. The calorimeter non-linearity is determined from the $`E_T`$ dependence of the energy scale. After applying a correction, the uncertainty on the slope results in $`\mathrm{\Delta }M_W(e)=\pm `$25 MeV.
The calorimeter resolution is parametrized as $`\sigma E_T/E_T`$=13.5%/$`\sqrt{E}_T`$$`\kappa `$, where $`\kappa `$ is determined from the width of the $`E/p`$ signal. The uncertainty on $`\kappa `$ results in $`\mathrm{\Delta }M_W(e)=\pm `$7 MeV.
The $`Z`$ boson masses are used as cross checks for the momentum and energy scales.
### 5.3 Recoil Scale and Resolution
The hadronic recoil is measured by summing over the energy in all calorimeter towers, excluding the lepton towers. The simulation removes an equivalent set of towers by subtracting the mean underlying event energy as measured from adjacent towers. The uncertainty on the measurement of this underlying event energy results in $`\mathrm{\Delta }M_W(e)=\pm `$15 MeV and $`\mathrm{\Delta }M_W(\mu )=\pm `$10 MeV uncertainties on the $`W`$ mass.
The hadronic recoil scale is the ratio of measured to true recoil. It is parametrized as a function of the true recoil and tuned using $`Z`$ events where both decay leptons are reconstructed and the $`Z`$ boson $`p_T`$ can be reconstructed precisely from the lepton momentum measurements. The uncertainty on the parametrization results in $`\mathrm{\Delta }M_W(e,\mu )=\pm `$20 MeV.
The recoil resolution model incorporates terms from the underlying event, which are modeled with generic inelastic collisions, and from hadronic jet resolution. The resolution uncertainty results in $`\mathrm{\Delta }M_W(e,\mu )=\pm `$42 MeV.
### 5.4 Backgrounds
The backgrounds in the $`W`$ boson data sample include $`W\tau \nu `$, $`Zll`$ where one lepton is outside the detector acceptance and not reconstructed, hadronic jets, where one jet mimics a lepton, cosmic rays, where one leg of the cosmic track is not reconstructed, and kaon decays, where the kaon track is misreconstructed, resulting in large apparent muon momentum. The background measurement uncertainties result in $`\mathrm{\Delta }M_W(e,\mu )=\pm `$20 MeV.
### 5.5 Mass Fits and Total Uncertainty
After including all measurement components, CDF obtains transverse mass (Figure 6)
and transverse energy distributions which are well modeled. Table 2 summarizes the uncertainties for the $`M_T`$ fits in the electron and muon decay channels.
For comparison the uncertainties from the previous collider run (Run 1b) are also included. The overall uncertainty is 76 MeV. The $`W`$ boson mass fit results are currently blinded with a constant offset. The offset will be removed when further cross checks have been completed.
## 6 Summary
The $`W`$ boson physics program at the Tevatron is very successful. CDF and DØ have measured the inclusive $`W`$ and $`Z`$ cross sections in all three leptonic decay channels, which show good agreement with NNLO calculations. From the cross section measurements, CDF has extracted competitive measurements on lepton universality and an indirect measurement of the $`W`$ boson width. DØ has measured the $`W`$ boson width directly in the electron channel with an uncertainty smaller than the Run 1 value. The new CDF $`W`$ charge asymmetry will help to further constrain the uncertainties of parton distribution functions, which affect all the aforementioned measurements. With the addition of 600 pb<sup>-1</sup> of data on tape, these measurements will further constrain the Standard Model.
CDF has determined the uncertainty on the $`W`$ boson mass with the first $``$200 pb<sup>-1</sup> of Run 2 data to be 76 MeV, which is lower than its Run 1 uncertainty of 79 MeV. With the additional data to come, Run 2 promises the world’s highest precision measurement of the $`W`$ boson mass, with an anticipated uncertainty of 30 MeV for 2 fb<sup>-1</sup>.
## 7 Acknowledgments
I would like to thank my colleagues from the CDF and DØ electroweak groups for their hard work and input to this talk. |
warning/0506/quant-ph0506260.html | ar5iv | text | # Random subspaces for encryption based on a private shared Cartesian frame
## I Introduction
Quantum information theory is concerned with implementing various communications tasks with a minimal use of resources Nie00 . In multi-party protocols, the most interesting resources are *nonlocal* ones, such as shared classical key or entanglement. Recently, it has become apparent that shared reference frames (SRFs) are another form of nonlocal resource that may be included in the accounting of any multi-party information-processing protocol. Heuristically, two parties are said to share a reference frame if there is a perfect correlation between the systems that define the bases of their respective local Hilbert spaces. For instance, if Alice and Bob each define their local Cartesian frame using classical gyroscopes in their labs, then they possess a shared reference frame if the rotation relating the frames defined by their gyroscopes is known to them.
Like entanglement, no amount of discussion between Alice and Bob will allow them to establish a shared reference frame; doing so requires a physical interaction between them that goes beyond the framework of classical information theory and, for that matter, the usual formalism of quantum information theory. For example, to establish a shared Cartesian frame between their respective labs, Alice and Bob may make use of a *pre-existing* frame such as the fixed stars or the Earth’s magnetic field. However, if no such shared frame exists *a priori*, then no amount of discussion will enable them to establish one; to do so, they must exchanges physical systems that carry some directional information such as spin-$`1/2`$ particles. Understanding the value of a shared reference frame as a new nonlocal resource and its relation to both private and quantum communication is therefore an important necessary step in the ongoing effort to understand the nature of information in physics.
Substantial progress has recently been made in this direction. Most research has focussed on determining the communication cost to establishing an SRF GisPop ; Per01 ; Bag01 ; Per01b ; Bag01b ; Lin03 ; Chi04 . There have also been several investigations into the impact of SRFs (or the lack thereof) on the efficiency with which one can perform a variety of bi-partite tasks, such as quantum and classical communication BRS03a ; Enk04 , key distribution Wal03 ; Boi03 , and the manipulation of entanglement Ver03 ; Bar03 ; BDSW04 ; Enk04 . Other work has studied the cryptographic consequences of the participants’ lack of an SRF KMP03 ; Ver03 ; HOT05 . Here we shall be interested instead in using *private* SRFs as a resource.
An SRF is private if the systems that define Alice and Bob’s local Hilbert space bases are not correlated with any other systems. In this case, the SRF can act as a new kind of key for private quantum and classical communication over a public channel. Consider a private shared Cartesian frame, as in BRS04 . Under the requirement that the communication have perfect fidelity and the privacy be perfect, it was found that for $`N`$ transmitted qubits, the number of private qubits that can be communicated asymptotically is $`\mathrm{log}_2N`$, and the number of private classical bits that can be communicated asymptotically is $`3\mathrm{log}_2N`$. This unusual factor of three relating the quantum and classical capacities is understood in terms of the details of the representation theory of SU(2) Ful91 , but should be contrasted with the “usual” factor of two that typically relates classical and quantum schemes Ben92 ; Amb00 . In the present paper, we ask the question of whether these capacities may be improved by allowing transmission with near-perfect rather than perfect privacy, as is usually considered in cryptography.
Note the following suggestive facts. $`2N`$ secret shared classical bits can be used in a one-time pad to encrypt $`2N`$ classical bits (cbits). However, if one asks how many private *qubits* can be transmitted using the secret key, the answer is a factor of two less; that same secret $`2N`$ cbit string can only encrypt $`N`$ qubits if perfect privacy is required Amb00 . If only near-perfect privacy is required, on the other hand, the number of secret cbits required per encrypted qubit shrinks from $`2`$ to $`1`$ asymptotically HLSW04 , so that the difference between encrypting cbits and qubits disappears.
A similar effect occurs in the domain of communication using shared entanglement. The superdense coding protocol Ben92 uses the transmission of $`N`$ qubits and the consumption of $`N`$ ebits to communicate $`2N`$ cbits with perfect privacy. (An ebit is a pure, maximally entangled state of two qubits.) The number of qubits that can be transmitted with perfect privacy and no errors is again a factor of two less, as implemented in the quantum Vernam cipher Leu02 . If either near-perfect privacy or near-perfect transmission is permitted, however, the superdense coding protocol can be extended to allow the transmission of nearly $`2N`$ *qubits* HHL04 ; HLW04 , again erasing the difference between sending classical and quantum data. The fact that the methods developed in BRS04 for communicating private classical data using a private shared reference frame made heavy use of superdense coding suggests that it may be possible, by relaxing the security conditions, to increase the private quantum capacity by a factor of three, from $`\mathrm{log}_2N`$ to nearly $`3\mathrm{log}_2N`$. We shall show that this is indeed the case.
## II Preliminaries
### II.1 Brief comment on notation
The symbol for a state (such as $`\phi `$ or $`\rho `$) also denotes its density matrix. A *pure* state is always denoted as a ket (e.g., $`|\phi `$) and the density matrix for a pure state $`|\phi `$ is written simply as $`\phi `$. We will also use the notation $`x_Ny_N`$ if $`lim_N\mathrm{}x_N/y_N=1`$. The term *irrep* denotes an irreducible representation of a group.
### II.2 Private quantum channels using private shared correlations
Whenever Alice and Bob have some private shared correlation, that is, one to which an eavesdropper Eve does not have access, Eve’s description of the systems transmitted along the channel is related to Alice’s description by a decohering superoperator, denoted by $``$ Amb00 ; BRS04 . Before discussing shared reference frames specifically, we begin by formalizing the notions of the private quantum and classical capacities of this decohering superoperator.
A $`\delta `$-*private quantum communication scheme for* $``$ consists of a completely positive, trace-preserving encoding $`𝒞`$, mapping message states on a logical Hilbert space $`_L`$ to encoded states on the Hilbert space $``$ of the transmitted system, such that (i) the operation $`𝒞`$ is invertible by Bob (who possesses the private shared correlations), allowing him to decode and recover states on $`_L`$ with perfect fidelity, and (ii) the encoding satisfies
$$(\phi )\rho _0_1\delta ,|\phi _L,$$
(1)
where $`\rho _0`$ is some fixed state on $``$, $`\rho \sigma _1\mathrm{Tr}|\rho \sigma |`$ is the trace distance between $`\rho `$ and $`\sigma `$, and $`\delta `$ is a security parameter. When $`\delta =0`$, the scheme is said to be *perfectly private*. The *private quantum capacity* of this channel, $`Q(,\delta )`$, is defined as $`Q(,\delta )=sup_𝒞\mathrm{log}_2\mathrm{dim}_L`$.
A $`\delta `$-*private classical communication scheme for* $``$ consists of a set $`\{\rho _i\}_{i=1}^m`$ of density operators on $``$ such that (i) the $`\{\rho _i\}`$ are orthogonal, so that Bob can distinguish these classical messages with certainty, and (ii) the encoding satisfies
$$[\rho _i]\rho _0_1\delta i,$$
(2)
where, again, $`\rho _0`$ is some fixed state in $``$ and $`\delta `$ is a security parameter. The *private classical capacity* of this channel, $`C(,\delta )`$, is defined to be $`C(,\delta )=sup_{\{\rho _i\}}\mathrm{log}_2|\{\rho _i\}|`$, where the supremum is over sets of density operators $`\{\rho _i\}`$ achieving $`\delta `$-privacy.
Given $`\delta `$-privacy, for any pair of quantum or classical messages, chosen with equal prior probabilities, the probability that Eve can distinguish these is bounded above by $`(1+\delta )/2`$, seen as follows. Suppose the message states are $`\varrho _1`$ and $`\varrho _2`$. These could be either an encoded pair of quantum messages, that is, $`𝒞(\varrho _{L,1})`$ and $`𝒞(\varrho _{L,2})`$ for some pair of density operators $`\varrho _{L,1}`$ and $`\varrho _{L,2}`$ on $`_L`$, or an encoded pair of classical messages, that is, orthogonal density operators. In either case, the optimal probability for Eve to distinguish $`(\varrho _1)`$ and $`(\varrho _2)`$ is given by $`\frac{1}{2}+\frac{1}{4}(\varrho _1)(\varrho _2)_1`$ Hel76 ; Fuc98 . Making use of the triangle inequality for the trace norm $`_1`$, we obtain
$`(\varrho _1)`$ $`(\varrho _2)_1`$
$`(\varrho _1)\rho _0_1+(\varrho _2)\rho _0_1`$
$`2\delta ,`$ (3)
where on the second line we have applied the definition of $`\delta `$-privacy. It follows that if the scheme is $`\delta `$-private, Eve’s probability of distinguishing the two messages is bounded above by $`(1+\delta )/2`$.
### II.3 Private quantum communication using a private shared Cartesian frame
We now determine the superoperator $``$ that describes Eve’s ignorance of Alice and Bob’s private shared Cartesian frame for states of $`N`$ spin-1/2 particles. (Our description applies equally well to any realization of a qubit that is entirely defined relative to some reference frame; another example is a single-photon polarization qubit.) The transmitted Hilbert space $``$ in this case is $`(^2)^N`$. This Hilbert space carries a tensor power representation $`R^N`$ of SU(2), by which an element $`\mathrm{\Omega }`$ SU(2) acts identically on each of the $`N`$ qubits. For simplicity, we restrict $`N`$ to be an even integer for the remainder of this paper, but our main results apply straightforwardly to all $`N`$. Then we can decompose
$$(^2)^N=\underset{j=0}{\overset{N/2}{}}_j,$$
(4)
where $`_j`$ is the eigenspace of total angular momentum with eigenvalue $`j`$.
Each subspace $`_j`$ in the direct sum can be factored into a tensor product $`_j=_{jR}_{jP}`$, such that SU(2) acts irreducibly on $`_{jR}`$ and trivially on $`_{jP}`$. Thus,
$$(^2)^N=\underset{j=0}{\overset{N/2}{}}_{jR}_{jP}.$$
(5)
The dimension of $`_{jR}`$ is
$$d_{jR}=2j+1,$$
(6)
and that of $`_{jP}`$ is
$$d_{jP}=\left(\genfrac{}{}{0pt}{}{N}{N/2j}\right)\frac{2j+1}{N/2+j+1}.$$
(7)
If Alice prepares $`N`$ qubits in a state $`\rho `$ and sends them to Bob, an eavesdropper Eve who is uncorrelated with the private SRF will describe the state as mixed over all rotations $`\mathrm{\Omega }`$ SU(2). Thus, the superoperator $``$ acting on a general density operator $`\rho `$ of $`N`$ qubits that describes the lack of knowledge of this private SRF is given by BRS03a
$$(\rho )=R(\mathrm{\Omega })^N\rho R^{}(\mathrm{\Omega })^Nd\mathrm{\Omega }.$$
(8)
The effect of this superoperator is best seen through use of the decomposition (5) of the Hilbert space. The action of the superoperator $``$ can be expressed in terms of this decomposition as
$$(\rho )=\underset{j=0}{\overset{N/2}{}}(𝒟_{jR}_{jP})(\mathrm{\Pi }_j\rho \mathrm{\Pi }_j),$$
(9)
where $`𝒟_{jR}`$ is the completely depolarizing superoperator on $`_{jR}`$, $`_{jP}`$ is the identity superoperator on $`_{jP}`$, and $`\mathrm{\Pi }_j`$ is the projector onto $`_j`$. The subsystems $`_{jP}`$ are called *decoherence-free* or *noiseless* subsystems Kni00 under the action of this superoperator; states encoded into these subsystems are completely protected from this decoherence. In contrast, $`_N`$ is completely depolarizing on each $`_{jR}`$ subsystem, and thus the $`_{jR}`$ are called *decoherence-full* subsystems BRS04 .
The largest decoherence-full subsystem occurs for $`j_{\mathrm{max}}=N/2`$ and has dimension $`2j_{\mathrm{max}}+1=N+1`$. As proven in BRS04 , this decoherence-full subsystem defines the optimally efficient perfectly secure private quantum communication scheme. Thus, given a private Cartesian frame and the transmission of $`N`$ qubits, Alice and Bob can with perfect privacy communicate $`Q(,0)=\mathrm{log}_2(N+1)\mathrm{log}_2N`$ qubits asymptotically.
In contrast, in that same paper it was shown that the private *classical* capacity using the private shared Cartesian frame was given by $`C(,0)3\mathrm{log}_2N`$. In Appendix A, we extend the result to show that $`C(,\delta )3(1+\delta )\mathrm{log}_2N+3`$ for $`\delta 1/2`$. The $`\delta `$-private classical capacity therefore does not change dramatically when $`\delta `$ is made non-zero.
### II.4 The working space $`^{}`$
To construct a “working” Hilbert space on which to investigate large random subspaces, we use the Hilbert space on which the states in the private classical communication scheme have support. This Hilbert space is constructed as follows. Note that for all $`j`$ strictly less than the maximum value $`N/2`$, the decoherence-free subsystem $`_{jP}`$ is always of *greater or equal* dimension than the decoherence-full subsystem $`_{jR}`$. Thus, we will employ irreps up to, but *not* including, $`j=N/2`$. Let $`j_{\mathrm{min}}<N/2`$ be some fixed irrep. Our working space $`^{}`$ will include elements from every irrep in the range $`j_{\mathrm{min}}j<N/2`$, that is, for $`jY,`$ where
$$Y=\{j_{\mathrm{min}},j_{\mathrm{min}}+1,\mathrm{},N/21\}.$$
(10)
For convenience, we denote the dimension of the decoherence-full subsystem of the $`j_{\mathrm{min}}`$ irrep by $`D`$, that is, $`D2j_{\mathrm{min}}+1`$. Choose a $`D`$dimensional subspace $`_{jR}^{}`$ of $`_{jR}`$ for every $`jY`$, and a subspace $`_{jP}^{}`$ of $`_{jP}`$ that is of dimension $`D_\alpha \frac{1}{\alpha }D,`$ for some parameter $`\alpha >1`$. Note that such subspaces always exist because $`\mathrm{dim}_{jR}=2j+1D`$ and $`\mathrm{dim}_{jP}\mathrm{dim}_{jR}`$ for all $`jY`$.
The Hilbert space of interest is then
$$^{}=\underset{jY}{}_{jR}^{}_{jP}^{},$$
(11)
with dimensionality $`Kdim^{}`$ given by
$`K`$ $`{\displaystyle \frac{1}{\alpha }}{\displaystyle \underset{jY}{}}D^2`$
$`={\displaystyle \frac{1}{\alpha }}(N/2j_{\mathrm{min}})(2j_{\mathrm{min}}+1)^2.`$ (12)
To maximize this dimension, we choose $`j_{\mathrm{min}}`$ to be the integer nearest to $`N/3`$. In this case, we have asymptotically
$$K\frac{2}{27}\frac{1}{\alpha }N^3.$$
(13)
(More precisely, $`K1`$ exceeds the righthand side for sufficiently large $`N`$, a result we will use later.)
The superoperator $``$ maps a state $`\phi `$ on $`^{}`$ to the state
$$(\phi )=\underset{jY}{}(I_{_{jR}}/d_{jR})\mathrm{Tr}_{jR}(\mathrm{\Pi }_j\phi \mathrm{\Pi }_j),$$
(14)
where $`I_{_{jR}}`$ is the identity on $`_{jR}`$. (Note that the state $`(\phi )`$ will, in general, have support outside of $`^{}`$.) To fully exploit the working space, we will pursue an encoding such that $`(\phi )`$ is close to maximally mixed on as large a subspace as possible. To this end, we define
$$\rho _0\underset{jY}{}(I_{_{jR}}/d_{jR})(I_{_{jP}^{}}/D_\alpha ),$$
(15)
where $`I_{_{jP}^{}}`$ is the identity on $`_{jP}^{}`$.
## III The main result
We wish to show that for fixed $`\delta `$, there exists a private quantum communication scheme for $``$ that scales as $`3\mathrm{log}_2N`$. This is achieved by encoding into particular subspaces of the working space $`^{}`$. Suppose that a subspace $`S^{}`$ of the appropriate dimensionality is drawn at random from some ensemble of subspaces of $`^{}`$. It is then sufficient to show that the probability that encoding in $`S`$ is not $`\delta `$-private is strictly less than 1, because this implies that there exist subspaces in the ensemble that do yield $`\delta `$-private schemes. Any such subspace can then constitute the logical Hilbert space $`_L`$ for such a scheme. In this case, the encoding map $`𝒞`$ is simply the embedding map, which takes states in $`S(^2)^N`$ to states in $`(^2)^N`$. Consequently, we leave the encoding map $`𝒞`$ implicit in the rest of the paper.
We shall consider the ensemble of subspaces that is generated by drawing uniformly at random from among all subspaces of $`^{}`$ of a given dimension. More precisely, we shall take $`S=US_0`$ where $`S_0`$ is a fixed subspace of $`^{}`$ and $`U`$ is a unitary on $`^{}`$ chosen according to the Haar measure $`dU`$. The condition that we require $`S`$ to satisfy in order to yield a $`\delta `$-private scheme is that for all $`|\phi S`$, $`(\phi )\rho _0_1\delta `$ for $`\rho _0`$ given by Eq. (15), so that from Eve’s perspective all the encoded states $`\phi `$ are near indistinguishable. This condition is equivalent to demanding that
$$\underset{|\phi S}{\mathrm{max}}(\phi )\rho _0_1\delta ,$$
(16)
where the maximization is over all pure states in $`S`$. The probability that $`S`$ fails to be $`\delta `$-private is therefore
$$\underset{S}{\mathrm{Pr}}\left(\underset{|\phi S}{\mathrm{max}}(\phi )\rho _0_1>\delta \right),$$
(17)
where we define the probability $`\mathrm{Pr}_S(g(S)>\delta )`$ that a randomly-chosen $`S`$ satisfies some inequality $`g(S)>\delta `$ by
$$\underset{S}{\mathrm{Pr}}\left(g(S)>\delta \right)_{\{U:g(US_0)>\delta \}}𝑑U.$$
(18)
For at least one of the $`S`$ to be $`\delta `$-private, we require that
$$\underset{S}{\mathrm{Pr}}\left(\underset{|\phi S}{\mathrm{max}}(\phi )\rho _0_1>\delta \right)<1,$$
(19)
for some $`\rho _0`$. The following theorem implies that such subspaces $`S`$, with dimension scaling in the desired fashion, do exist.
###### Theorem 1
For the decoherence map $``$ associated with lacking a reference frame for $`SU(2)`$, the condition
$$\underset{S}{\mathrm{Pr}}\left(\underset{|\phi S}{\mathrm{max}}(\phi )\rho _0_1>\delta \right)<1,$$
(20)
holds for sufficiently large $`N`$, where the probability is with respect to the unitarily invariant measure on subspaces $`S`$ of $`^{}`$, provided
$$\mathrm{log}_2dimS<3\mathrm{log}_2N+7/2\mathrm{log}_2\delta +C^{}$$
(21)
where $`C^{}`$ is a constant.
Consider, for example, $`1/\delta =\mathrm{polylog}(N)`$, i.e., a polynomial in $`\mathrm{log}(N)`$. Then we can find $`S^{}`$ with $`(\phi )\rho _0_1\delta `$ for $`|\phi S`$ such that
$$\mathrm{log}_2dimS3\mathrm{log}_2N,$$
(22)
recovering the same asymptotic rate as the classical private capacity. In this case, $`Q(,\delta )C(,0)`$.
We prove Theorem 1 via a sequence of lemmas. Our starting point is the following result, known as Levy’s Lemma MS86 :
###### Lemma 2 (Levy)
Let $`f:S^k`$ be a continuous real-valued function on the $`k`$-sphere with Lipschitz constant $`\eta `$ with respect to the Euclidean metric. Then, if $`x`$ is selected at random from $`S^k`$ according to the uniform measure,
$$\underset{x}{\mathrm{Pr}}\left(|f(x)M|>\gamma \right)<\mathrm{exp}_2\left(C(k1)\gamma ^2/\eta ^2\right),$$
(23)
where $`C>0`$ is a constant and $`M`$ is a median for $`f.`$
The function of interest is:
$$f(\phi )(\phi )\rho _0_1.$$
(24)
Note that the Hilbert space norm on $`^{}`$ is precisely the Euclidean norm if the Hilbert space is considered as the real vector space $`^{2K}`$. The following lemma bounds the Lipschitz constant of this function.
###### Lemma 3 (Lipschitz constant)
The Lipschitz constant of $`f(\phi )`$ is bounded above by $`2`$.
Proof. Using the triangle inequality gives
$`|f(\phi )f(\stackrel{~}{\phi })|`$ $`=\left|(\phi )\rho _0_1(\stackrel{~}{\phi })\rho _0_1\right|`$
$`(\phi )(\stackrel{~}{\phi })_1.`$ (25)
Because $``$ is a completely positive trace-preserving map, and the trace distance is non-increasing under such maps,
$$(\phi )(\stackrel{~}{\phi })_1\phi \stackrel{~}{\phi }_1.$$
(26)
Combining these inequalities with the fact that
$`|\phi |\stackrel{~}{\phi }_2^2`$ $`=`$ $`22\text{Re}\phi |\stackrel{~}{\phi }`$ (27)
$``$ $`1|\phi |\stackrel{~}{\phi }|^2`$ (28)
$`=`$ $`\left(\frac{1}{2}\phi \stackrel{~}{\phi }_1\right)^2,`$ (29)
we obtain
$$|f(\phi )f(\stackrel{~}{\phi })|\text{ }2|\phi |\stackrel{~}{\phi }_2,$$
(30)
which is the desired bound on the Lipschitz constant.
The next corollary, an immediate consequence of Levy’s Lemma, bounds the probability that Eve can distinguish a random state on $`^{}`$ from $`\rho _0`$ substantially better than she can distinguish states on average.
###### Corollary 4 (Concentration of $`f`$)
Let $`|\phi `$ be chosen at random from the uniform measure on the unit sphere in $`^{}`$ and $`M`$ a median for $`f`$. Then
$$\underset{\phi }{\mathrm{Pr}}\left(\left|f(\phi )M\right|>\gamma \right)\mathrm{exp}_2\left(\frac{C(K1)\gamma ^2}{2}\right).$$
(31)
Proof. Apply Levy to the function $`f(\phi )`$ of Eq. (24). In this case, $`k=2K1`$ and $`\eta 2`$.
Next, we relate the median of $`f`$ to its mean, which is easier to estimate. Write
$$𝔼_\phi ff(\phi )𝑑\nu (\phi )$$
(32)
for the expectation of $`f`$ with respect to the unitarily invariant measure $`d\nu (\phi )`$ on pure states in $`^{}`$. Let $`A_{}^{}`$ be the set of points $`|\phi `$ on the unit sphere for which $`f(\phi )M`$. By the definition of the median,
$$_A_{}f(\phi )𝑑\nu (\phi )M_A_{}𝑑\nu (\phi )=M\frac{1}{2}.$$
(33)
Letting $`A_<`$ be defined analogously, we get
$$𝔼_\phi f=_{A_<}f(\phi )𝑑\nu (\phi )+_A_{}f(\phi )𝑑\nu (\phi )\frac{M}{2},$$
(34)
because $`f(\phi )0`$.
###### Lemma 5 (Expectation of $`f`$)
The expectation value of $`f(\phi )`$ satisfies
$$𝔼_\phi f\frac{1}{\sqrt{\alpha }}.$$
(35)
The proof is supplied in Appendix B, but can be understood intuitively in terms of the action of $``$ on $`^{}`$. If the subspaces $`_{jP}^{}`$ were $`1`$-dimensional, then by virtue of the fact that the $`_{jR}`$ are decoherence-full, we would have complete decoherence on $`_{jR}_{jP}^{}`$. Because $`1/\alpha dim_{jP}^{}/dim_{jR}^{}`$, the larger the value of $`\alpha ,`$ the smaller the dimension of $`_{jP}^{}`$ relative to $`_{jR}^{}`$, and the less distinguishable on average are states on $`_{jR}_{jP}^{}`$ subsequent to the action of $``$. One might expect that states could be distinguished by their relative supports on the different irreps $`jY`$, because these supports are invariant under the action of $``$. However, the proof demonstrates that because we use only sufficiently large irreps, all encoded states will have similar supports on all irreps, and thus not be significantly more distinguishable than if a single irrep had been used.
We note that the proof of the lemma requires that, within each irrep $`jY`$, the dimension $`D_\alpha `$ of $`_{jP}^{}`$ be much smaller than the dimension $`D`$ of the decoherence-full subsystem $`_{jR}`$. For this reason, our result does not apply directly to cryptography using a U(1) phase reference, for which the decoherence-full subsystems are all one-dimensional. However, for any other reference frame that satisfies this condition, our results should be directly applicable.
We conclude that the median $`M`$ is upper bounded by $`2/\sqrt{\alpha }`$ which, using Corollary 4, leads to the result
$$\begin{array}{c}\underset{\phi }{\mathrm{Pr}}\left((\phi )\rho _0_1>\gamma +\frac{2}{\sqrt{\alpha }}\right)\hfill \\ \hfill \mathrm{exp}_2\left(\frac{C}{2}(K1)\gamma ^2\right).\end{array}$$
(36)
This inequality is sufficiently strong that we will be able to use it to conclude that large subspaces of $`^{}`$ have the property that the distinguishability of *all* states in the subspace are bounded.
###### Lemma 6 (Existence of good subspaces)
Let $`S_0^{}`$ be a fixed subspace and $`|\phi _0`$ a fixed state on $`S_0`$. Let $`S=US_0`$ be a random subspace obtained from $`S_0`$ using a Haar-distributed unitary $`U`$ on $`^{}`$. Then, for any $`\delta >0`$ and $`0<\epsilon <1/2`$,
$$\begin{array}{c}\underset{S}{\mathrm{Pr}}\left(\underset{|\phi S}{\mathrm{max}}f(\phi )>\delta \right)\hfill \\ \hfill \left(\frac{5}{\epsilon }\right)^{2dimS}\underset{U}{\mathrm{Pr}}\left(f(U|\phi _0)>\delta \epsilon \right).\end{array}$$
(37)
Proof. Fix an $`ϵ/2`$-net $`𝒩_0`$ for the unit sphere of a fixed subspace $`S_0`$ of $`^{}`$ with the Hilbert space norm. The net can be chosen such that the number of elements in the net satisfies $`|𝒩_0|(5/ϵ)^{2dimS_0}`$. (See HLW04 a proof of this fact.) By definition, given any $`|\phi S_0`$, there exists a state $`|\stackrel{~}{\phi }𝒩_0`$ such that $`|\phi |\stackrel{~}{\phi }_2ϵ/2`$. By Lemma 3, this implies that $`|f(\phi )f(\stackrel{~}{\phi })|ϵ`$.
Now choose a random subspace $`S=US_0`$ using a Haar-distributed unitary. This unitary $`U`$ maps the net $`𝒩_0`$ for $`S_0`$ into a net $`𝒩`$ for $`S`$. Let $`|\phi ^{}`$ be defined by
$$f(\phi ^{})=\underset{|\phi S}{\mathrm{max}}f(\phi ).$$
(38)
By definition, there exists a state $`|\stackrel{~}{\phi }^{}𝒩`$ such that $`|\phi ^{}|\stackrel{~}{\phi }^{}_2ϵ/2`$, and consequently $`|f(\phi ^{})f(\stackrel{~}{\phi }^{})|ϵ`$. It follows that if $`f(\phi ^{})>\delta `$, then $`f(\stackrel{~}{\phi }^{})>\delta ϵ`$. Therefore, if
$$\underset{|\phi S}{\mathrm{max}}f(\phi )>\delta ,\text{then}\underset{|\stackrel{~}{\phi }𝒩}{\mathrm{max}}f(\stackrel{~}{\phi })>\delta ϵ.$$
(39)
Finally, if $`x`$ implies $`y,`$ then $`\mathrm{Pr}(x)\mathrm{Pr}(y),`$ so we conclude that
$$\underset{S}{\mathrm{Pr}}\left(\underset{|\phi S}{\mathrm{max}}f(\phi )>\delta \right)\underset{U}{\mathrm{Pr}}\left(\underset{|\stackrel{~}{\phi }𝒩_0}{\mathrm{max}}f(U|\stackrel{~}{\phi })>\delta ϵ\right)$$
(40)
where $`\mathrm{Pr}_U`$ reminds the reader that we are varying over unitaries. We then have
$`\underset{U}{\mathrm{Pr}}(\underset{|\stackrel{~}{\phi }𝒩_0}{\mathrm{max}}`$ $`f(U|\stackrel{~}{\phi })>\delta ϵ)`$
$`{\displaystyle \underset{|\stackrel{~}{\phi }𝒩_0}{}}\underset{U}{\mathrm{Pr}}\left(f(U|\stackrel{~}{\phi })>\delta ϵ\right)`$
$`=\left|𝒩_0\right|\underset{U}{\mathrm{Pr}}\left(f(U|\stackrel{~}{\phi }_0)>\delta ϵ\right),`$ (41)
where the first inequality is the union bound for probabilities and the second line follows from the fact that the expression inside the sum over $`|\stackrel{~}{\phi }`$ is independent of $`|\stackrel{~}{\phi }`$ ($`|\stackrel{~}{\phi }_0`$ is an arbitrary state in $`^{}`$). Recalling that $`|𝒩_0|`$ $`(5/ϵ)^{2dimS_0}`$ establishes what we set out to prove.
Using the lemma together with Eq. (36), we obtain
$$\begin{array}{c}\underset{S}{\mathrm{Pr}}\left(\underset{|\phi S}{\mathrm{max}}(\phi )\rho _0_1>\delta \right)\hfill \\ \hfill \left(\frac{5}{\epsilon }\right)^{2dimS}\mathrm{exp}_2\left(\frac{C}{2}(K1)(\delta \epsilon \frac{2}{\sqrt{\alpha }})^2\right).\end{array}$$
(42)
If $`dimS`$ is chosen such that the right hand side is bounded away from $`1,`$ then the left hand side will also be bounded away from $`1,`$ and there will exist a $`\delta `$-private encoding into a subspace $`S.`$ We will therefore seek the largest value of $`dimS`$ that satisfies the inequality
$$\left(\frac{5}{\epsilon }\right)^{2dimS}<\mathrm{exp}_2\left(\frac{C}{2}(K1)(\delta \epsilon \frac{2}{\sqrt{\alpha }})^2\right),$$
(43)
or equivalently
$$dimS<\frac{\mathrm{ln}2}{\mathrm{ln}\left(\frac{5}{\epsilon }\right)}\frac{C}{4}(K1)(\delta \epsilon \frac{2}{\sqrt{\alpha }})^2.$$
(44)
Given that $`\mathrm{ln}x\sqrt{x}`$, we have $`1/\mathrm{ln}\left(\frac{5}{\epsilon }\right)1/\sqrt{\frac{5}{\epsilon }}`$ and any $`S`$ satisfying
$$dimS<\sqrt{\frac{\epsilon }{5}}\frac{C\mathrm{ln}2}{4}(K1)(\delta \epsilon \frac{2}{\sqrt{\alpha }})^2,$$
(45)
will also satisfy Eq. (44). Using the expression for $`K`$ in Eq. (13), it is sufficient to require that
$$dimS<\frac{C\mathrm{ln}2}{54\sqrt{5}}\frac{1}{\alpha }N^3(\delta \epsilon \frac{2}{\sqrt{\alpha }})^2\sqrt{\epsilon }$$
(46)
for sufficiently large $`N`$. If we choose $`\epsilon =\delta /3`$ and $`\alpha =36/\delta ^2`$, then this expression reduces to
$$dimS<\frac{C\mathrm{ln}2}{5832\sqrt{15}}N^3\delta ^{7/2}.$$
(47)
It is therefore possible to choose $`S`$ such that $`f(\phi )\delta `$ for all $`|\phi S`$ whenever
$$\mathrm{log}_2dimS<3\mathrm{log}_2N+7/2\mathrm{log}_2\delta +C^{}$$
(48)
where $`C^{}=\mathrm{log}_2[(C\mathrm{ln}2)/(5832\sqrt{15})]`$, completing the proof of Theorem 1.
## IV Discussion
We have seen that for fixed $`\delta >0`$, the $`\delta `$-private quantum capacity of a secret SU(2) reference frame is at least three times as large as its perfectly private quantum capacity. Indeed, the relaxation of the security requirement to $`\delta >0`$ causes the private quantum capacity to jump almost to the value of the perfectly private classical capacity, which is approximately $`3\mathrm{log}_2N`$, and within a factor of $`1+\delta `$ of the $`\delta `$-private classical capacity. In earlier work, a similar relaxation of the security condition in the quantum one time pad led to a doubling of the private quantum capacity of a shared secret key string HLSW04 as well as a similar doubling of the capacity of a maximally entangled state HHL04 ; HLW04 . The tripling of the capacity seen here, however, is unusual and reflects the particular structure of the tensor power representation of SU(2).
Because the private capacity of a shared reference frame is proportional to $`\mathrm{log}_2N`$ rather than $`N`$, however, the values of $`\delta `$ which provide an improvement over the perfectly private schemes is quite restricted. From Theorem 1, we see that for sufficiently large $`N`$,
$$Q(,\delta )3\mathrm{log}_2N+\frac{7}{2}\mathrm{log}_2\delta +C^{}$$
(49)
for some constant $`C^{}`$. In order to improve upon the perfectly private scheme, we require that $`Q(,\delta )>\mathrm{log}_2N`$, which implies that $`1/\delta O(N^{4/7})`$. In particular, our construction does not allow $`\delta `$ to be an exponentially decreasing function of $`N`$, which would obviously be more desirable for cryptographic applications.
Some questions remain about the optimality of the private quantum communication schemes we have presented here. In particular, our upper bounds on the private quantum capacity do not exclude the possibility that $`\delta `$ could be made to shrink exponentially with $`N`$ while maintaining a number of qubits sent scaling as $`3\mathrm{log}_2N`$. Also, we have not attempted to construct $`\delta `$-private classical communication schemes meeting the upper bound of Theorem 7 in Appendix A.
Finally, we note that a shared Cartesian frame is not the only possible form of a shared reference BRS04 , and it is useful to consider other practical examples such as a shared phase reference, shared direction, or reference ordering. These examples have different Hilbert-space structures arising from their group representation theory, and in general will result in different relations between their private classical and quantum capacities. We note that our technique should apply directly to cryptography using a reference frame for $`U(K)`$, with $`K2`$, because the Hilbert space structures for these groups satisfy the conditions required for our proof. Whether similar differences between perfectly-private and $`\delta `$-private capacities can be found for other reference frames is an open question.
###### Acknowledgements.
The authors gratefully acknowledge J. Emerson, D. Gottesman, and T. Rudolph for helpful discussions. SDB acknowledges support from the Australian Research Council, PH appreciates the support of the Canadian Institute for Advanced Research, and both PH and RWS are grateful for support from the Natural Sciences and Engineering Research Council of Canada.
## Appendix A $`\delta `$-private classical capacity
###### Theorem 7
For $`\delta 1/2`$, the $`\delta `$-private classical capacity satisfies $`C(,\delta )3(1+\delta )\mathrm{log}_2N+3`$.
Proof. Suppose we have a $`\delta `$-private classical communication scheme for $``$ consisting of $`m`$ states on $``$. If such a scheme exists, then there is also an $`m`$-state scheme using pure states, which we will label $`\{|\psi _i\}_{i=1}^m`$. We will use the privacy condition to find a small subspace of $``$ such that these states are almost entirely contained within the subspace. Combining the Holevo bound with the fact that the original states were all distinguishable will then lead to an upper bound on $`m`$, the number of states in the scheme.
Let $`_{jQ}`$ be the subspace of $`_{jP}`$ corresponding to the non-zero eigenvalues of $`\mathrm{Tr}_{jR}[\mathrm{\Pi }_j\psi _1\mathrm{\Pi }_j]`$ and let $`\mathrm{\Pi }_{jQ}`$ be the projector onto $`_{jQ}`$. It follows from the Schmidt decomposition for $`\mathrm{\Pi }_j|\psi _1`$ that $`dim_{jQ}\mathrm{min}(d_{jR},d_{jP})`$. Also let $`\mathrm{\Pi }^{}=_j\mathrm{\Pi }_{jR}\mathrm{\Pi }_{jQ}`$, where $`\mathrm{\Pi }_{jR}`$ is the projector onto $`_{jR}`$. Observe that for any $`\psi _i`$,
$`\mathrm{Tr}[\mathrm{\Pi }^{}\psi _i\mathrm{\Pi }^{}]=\mathrm{Tr}[(\mathrm{\Pi }^{}\psi _i\mathrm{\Pi }^{})]=\mathrm{Tr}[\mathrm{\Pi }^{}(\psi _i)\mathrm{\Pi }^{}]`$ (50)
because $``$ is trace-preserving and because projection by $`\mathrm{\Pi }^{}`$ commutes with $``$. By the privacy condition, however,
$`\delta `$ $``$ $`(\psi _1)(\psi _i)_1`$ (51)
$``$ $`2\{\mathrm{Tr}[\mathrm{\Pi }^{}(\psi _1)\mathrm{\Pi }^{}]\mathrm{Tr}[\mathrm{\Pi }^{}(\psi _i)\mathrm{\Pi }^{}]\}`$ (52)
$`=`$ $`2\{1\mathrm{Tr}[\mathrm{\Pi }^{}(\psi _i)\mathrm{\Pi }^{}]\}.`$ (53)
The second inequality holds because $`X_1=2\mathrm{max}_P\mathrm{Tr}[PX]`$ for traceless, Hermitian $`X`$, where the optimization is over projectors of all ranks. (See, for example, Nie00 .) Combining (53) with (50) shows that $`\mathrm{Tr}[\mathrm{\Pi }^{}\psi _i\mathrm{\Pi }^{}]1\delta /2`$. Thus the states $`\{|\psi _i\}_{i=1}^m`$ are essentially contained within the subspace defined by $`\mathrm{\Pi }^{}`$.
Now consider the set of states $`\{|\psi _i^{}\}_{i=1}^m`$, where
$$|\psi _i^{}=\frac{\mathrm{\Pi }^{}|\psi _i}{\sqrt{\mathrm{Tr}[\mathrm{\Pi }^{}\psi _i\mathrm{\Pi }^{}]}}.$$
(54)
Because $`|\psi _i|\psi _i^{}|^2=\mathrm{Tr}[\mathrm{\Pi }^{}\psi _i\mathrm{\Pi }^{}]`$, performing the measurement $`\{|\psi _j\psi _j|\}_{j=1}^m`$ on the set of states $`\{|\psi _i^{}\}_{i=1}^m`$ will correctly identify the state with probability at least $`1\delta /2`$. Assume a state $`|\psi _i^{}`$ is chosen from the uniform distribution. By Fano’s inequality CT91 ,
$$H(i|j)1+\frac{\delta }{2}\mathrm{log}_2m,$$
(55)
where $`H`$ is the Shannon conditional entropy function, which in turn implies that
$$I(i;j)(1\delta /2)\mathrm{log}_2m1,$$
(56)
where $`I`$ is the mutual information function. Because all the states $`|\psi _i^{}`$ are contained in the support of $`\mathrm{\Pi }^{}`$, however, the Holevo bound H73 implies that $`I(i;j)`$ is no more than the logarithm of $`\mathrm{rank}\mathrm{\Pi }^{}`$, which satisfies
$$\mathrm{rank}\mathrm{\Pi }^{}\underset{j}{}d_{jR}\mathrm{min}(d_{jR},d_{jP}).$$
(57)
In the case of a private shared SU(2) reference frame, for which $`d_{jR}=2j+1`$,
$$\mathrm{rank}\mathrm{\Pi }^{}(N/2+1)(N+1)^22N^3,$$
(58)
where the second inequality holds for all $`N2`$. This implies that
$`\mathrm{log}_2m`$ $``$ $`{\displaystyle \frac{3\mathrm{log}_2N+2}{1\delta /2}}`$ (59)
$``$ $`3(1+\delta )\mathrm{log}_2N+3,`$ (60)
provided $`\delta 1/2`$.
## Appendix B Proof of Lemma 5
The map $``$ depolarizes each of the systems $`_{jR}`$ but for the purposes of calculation, it is easier to simply discard them. In the proof, therefore, we will work with the space $`_P^{}=_{jY}_{jP}^{}`$, which has dimension $`d_Pdim_P^{}`$. Observe that if we introduce
$$(\rho )=\underset{jY}{}\mathrm{Tr}_{jR}(\mathrm{\Pi }_j\rho \mathrm{\Pi }_j),$$
(61)
which gives a normalized state on $`_P^{}`$, then
$$(\phi )\rho _0_1=(\phi )\varrho _0_1,$$
(62)
where $`\varrho _0=I_P/d_P`$ is the normalized identity operator on $`_P^{}`$.
Using $`X_1\sqrt{\mathrm{rank}X}X_2`$ gives
$$(\phi )\varrho _0_1\sqrt{d_P}(\phi )\varrho _0_2.$$
(63)
We therefore have
$`𝔼_\phi f`$ $`\sqrt{d_P}{\displaystyle (\phi )\varrho _0_2𝑑\nu (\phi )}`$
$`=\sqrt{d_P}{\displaystyle \sqrt{\mathrm{Tr}\left[(\phi )^2(\phi )/d_P+I_P/d_P^2\right]}𝑑\nu (\phi )}.`$ (64)
Using the normalization $`\mathrm{Tr}(\phi )=1`$ and the concavity of the square root function, this expression reduces to
$$𝔼_\phi f\sqrt{d_P\mathrm{Tr}[(\phi )^2]𝑑\nu (\phi )1}.$$
(65)
It therefore suffices to evaluate
$$\begin{array}{c}\mathrm{Tr}[(\phi )^2]𝑑\nu (\phi )\hfill \\ \hfill =\mathrm{Tr}\left[\left(\underset{jY}{}\mathrm{Tr}_{jR}\left(\mathrm{\Pi }_j\phi \mathrm{\Pi }_j\right)\right)^2\right]𝑑\nu (\phi ).\end{array}$$
(66)
Because $`\mathrm{\Pi }_j`$ has the form $`\mathrm{\Pi }_j=\mathrm{\Pi }_{jR}\mathrm{\Pi }_{jP}`$, where $`\mathrm{\Pi }_{jR}`$ and $`\mathrm{\Pi }_{jP}`$ are the projectors onto $`_{jR}`$ and $`_{jP}`$ respectively, $`\mathrm{Tr}_{jR}(\mathrm{\Pi }_j\phi \mathrm{\Pi }_j)`$ and $`\mathrm{Tr}_{kR}(\mathrm{\Pi }_k\phi \mathrm{\Pi }_k)`$ have orthogonal supports, implying that
$$\mathrm{Tr}\left[\left(\underset{jY}{}\mathrm{Tr}_{jR}\left(\mathrm{\Pi }_j\phi \mathrm{\Pi }_j\right)\right)^2\right]=\mathrm{Tr}\left[\underset{jY}{}\left(\mathrm{Tr}_{jR}\left(\mathrm{\Pi }_j\phi \mathrm{\Pi }_j\right)\right)^2\right].$$
(67)
To evaluate the resulting integral, fix bases $`\{|m\}_{m=1}^D`$ and $`\{|l\}_{l=1}^{D_\alpha }`$ for the spaces $`_{jR}^{}`$ and $`_{jP}^{}`$ respectively. (Note that we identify bases labelled by different values of $`j`$.) Also let $`|\phi _0=|j_0m_0l_0`$ for some fixed values of $`j_0`$, $`m_0`$ and $`l_0`$. Using Eq. (67) and making use of the invariance of the measure, we can expand the integral of Eq. (66) as
$$_{\mathrm{U}(K)}\mathrm{Tr}[(U\phi _0U^{})^2]𝑑U=\underset{jY}{}\underset{m,m^{}=1}{\overset{D}{}}\underset{l,l^{}=1}{\overset{D_\alpha }{}}_{\mathrm{U}(K)}U_{jml,j_0m_0l_0}U_{jml^{},j_0m_0l_0}^{}U_{jm^{}l^{},j_0m_0l_0}U_{jm^{}l,j_0m_0l_0}^{}𝑑U,$$
(68)
which can be evaluated using the identity (see, for example, AL03 )
$$_{\mathrm{U}(K)}U_{ij}U_{kl}^{}U_{mn}U_{pq}^{}𝑑U=\frac{1}{K^21}\left\{\delta _{ij,kl}\delta _{mn,pq}+\delta _{ij,pq}\delta _{kl,mn}\frac{1}{K}\delta _{ik}\delta _{jq}\delta _{ln}\delta _{mp}\frac{1}{K}\delta _{ip}\delta _{jl}\delta _{km}\delta _{nq}\right\}.$$
(69)
We obtain
$`{\displaystyle _{\mathrm{U}(K)}}\mathrm{Tr}[(U\phi _0U^{})^2]𝑑U`$ $`={\displaystyle \underset{jY}{}}{\displaystyle \underset{m,m^{}=1}{\overset{D}{}}}{\displaystyle \underset{l,l^{}=1}{\overset{D_\alpha }{}}}{\displaystyle \frac{1}{K(K+1)}}\left\{\delta _{l,l^{}}+\delta _{m,m^{}}\right\}`$ (70)
$`={\displaystyle \frac{_{jY}(D^2D_\alpha +D_\alpha ^2D)}{K(K+1)}}.`$ (71)
Substituting this back into the expression for $`𝔼_\phi f`$ yields
$$𝔼_\phi f\sqrt{\frac{d_P}{K(K+1)}\left(\underset{jY}{}(D^2D_\alpha +D_\alpha ^2D)\right)1}.$$
(72)
Recalling that $`d_P=_{jY}D_\alpha `$ and $`K=_{jY}D_\alpha D`$, we get $`𝔼_\phi f\sqrt{D_\alpha /D}`$. Because $`D_\alpha \frac{1}{\alpha }D`$, we have the desired inequality:
$$𝔼_\phi f\sqrt{\frac{1}{\alpha }}.$$
(73) |
warning/0506/nucl-th0506079.html | ar5iv | text | # Semiclassical description of shell effects in finite fermion systems
## I Introduction
The Nilsson model sgn gave rise to my first steps towards scientific research. My teacher in theoretical physics was Kurt Alder, who had been a member of the theoretical ’Coulomb excitation crew’ coulex at the Niels Bohr Institute during the years when collective nuclear motion was explored and single-particle motion in deformed nuclei was studied. When I asked Alder about a possible subject for my diploma work, he gave me a copy of Sven Gösta’s famous paper and asked me to write a program for computing the deformed single-particle wavefunctions; he had become weary of interpolating their expansion coefficients in Nilsson’s tables. I set about to diagonalize the Nilsson Hamiltonian using, however, not the spherical basis but that of the deformed harmonic oscillator, i.e., the asymptotic Nilsson states. The results were bound to be the same as those of Nilsson and were therefore not published mbdip – but through this exercise I became initiated into the shell structure of nuclei. I felt very privileged when I later came to know the great human being behind the famous model.
In this paper I want to review the semiclassical description of shell-effects in finite fermion systems using the periodic orbit theory. After a brief reminder about trace formulae, I will discuss some of their applications to systems in four different branches of physics. Since all results have been published elsewhere, I will not reproduce here any figures, but just discuss in words some of the most important results and conclusions.
## II Periodic orbit theory and gross-shell effects
### II.1 Semiclassical trace formulae
The periodic orbit theory (POT) was initiated by M. Gutzwiller in a series of publications culminating in his seminal paper in 1971 that contains the semiclassical trace formula gutz . It relates the quantum spectrum $`\{E_i\}`$ (which we here assume to be discrete, although the inclusion of a continuum is possible) of a hermitian Hamiltonian $`\widehat{H}`$ to the periodic orbits of the corresponding classical Hamiltonian $`H(𝐪,𝐩)`$. The quantum-mechanical level density, defined as the sum of Dirac delta functions peaked at the levels $`E_i`$, can be decomposed into a smooth part $`\stackrel{~}{g}(E)`$ and an oscillating part $`\delta g(E)`$:
$`g(E)={\displaystyle \underset{i}{}}\delta (EE_i)=\stackrel{~}{g}(E)+\delta g(E).`$ (1)
The smooth part contains by definition the average level density which usually is a monotonously increasing function of $`E`$ and can be obtained in the extended Thomas-Fermi (ETF) model (see, e.g., book , Ch. 4). The oscillating part can be expressed by the semiclassical trace formula
$`\delta g_{sc}(E){\displaystyle \underset{po}{}}𝒜_{po}(E)\mathrm{cos}[S_{po}(E)/\mathrm{}\sigma _{po}\pi /2].`$ (2)
The sum is over all periodic orbits $`(po)`$ of the classical system, $`S_{po}(E)=𝐩𝑑𝐪`$ are their action integrals, the amplitudes $`𝒜_{po}(E)`$ depend on their stabilities and degeneracies, and $`\sigma _{po}`$ are called the Maslov indices. The sum in (2) is an asymptotic one, correct to leading order in $`1/\mathrm{}`$, and in non-integrable systems it is hampered by convergence problems chaos . For isolated orbits, Gutzwiller expressed gutz the amplitudes $`𝒜_{po}(E)`$ in terms of their periods and stability matrices.
The trace formula (2) was later generalized to billiard systems bablo and to systems with continuous symmetries struma ; bertab ; crli , including integrable systems. A relativistic trace formula for spin 1/2 particles was derived in boke , and a nonrelativistic trace formula for particles with arbitrary spin $`s`$ in plet . In all cases, the trace formula has the same general form (2), but the amplitudes $`𝒜_{po}(E)`$ take different forms. For isolated orbits, their $`\mathrm{}`$ dependence is given by a factor $`\mathrm{}^1`$, while for orbits appearing in $`f`$-fold degenerate families, the amplitudes go like $`\mathrm{}^{(1+f/2)}`$.
In integrable and mixed-dynamical systems, periodic orbits can change their stability under the variation of a control parameter (e.g., the energy $`E`$, a potential parameter, or an additional external field) and thereby undergo bifurcations. In such situations, the amplitudes $`𝒜_{po}`$ diverge at the bifurcation points. The same happens also in limits where continuous symmetries are broken (or restored), since hereby the $`\mathrm{}`$ dependence of the $`𝒜_{po}`$ changes discontinuously. The remedy to remove these (unphysical!) divergences is to go beyond the stationary-phase approximation for the integration(s) used in the derivation of the semiclassical trace formula. This has, besides bablo ; struma ; bertab , been developed most systematically in ozoha for symmetry breaking and bifurcations, and in crpert for symmetry breaking in weakly perturbed integrable systems, leading in all cases to local uniform approximations with finite amplitudes $`𝒜_{po}`$. Global uniform approximations which yield finite amplitudes at symmetry-breaking and bifurcation points, and far from them go over into the standard (extended) Gutzwiller trace formula, were developed for the breaking of U(1) symmetry in toms , for some cases of U(2) and SO(3) symmetry breaking in hhuni , for the symmetry breaking U(3) $``$ SO(3) in boys , and for various types of bifurcations in ssun . (Details and further references may be found in book , Ch. 6.3.)
For interacting finite fermion systems described in the mean-field approximation (i.e., in Hartree-Fock or density functional theory), one can also obtain semiclassical trace formulae for the oscillating parts of the total binding energy $`E_b`$ and the particle number $`N`$. Hereby one writes, similarly to (1), $`E_b=\stackrel{~}{E}_b+\delta E`$ and $`N=\stackrel{~}{N}+\delta N`$. The average quantities $`\stackrel{~}{E}_b`$ and $`\stackrel{~}{N}`$ are taken from the ETF model, and for the oscillating parts one finds book ; struma
$`\delta E_{sc}`$ $``$ $`{\displaystyle \underset{po}{}}𝒜_{po}(\lambda )\left({\displaystyle \frac{\mathrm{}}{T_{po}}}\right)^2\mathrm{cos}[{\displaystyle \frac{S_{po}(\lambda )}{\mathrm{}}}\sigma _{po}{\displaystyle \frac{\pi }{2}}],`$
$`\delta N_{sc}`$ $``$ $`{\displaystyle \underset{po}{}}𝒜_{po}(\lambda )\left({\displaystyle \frac{\mathrm{}}{T_{po}}}\right)\mathrm{sin}[{\displaystyle \frac{S_{po}(\lambda )}{\mathrm{}}}\sigma _{po}{\displaystyle \frac{\pi }{2}}],`$ (3)
both to be evaluated at the Fermi energy $`\lambda (N)`$ for a given number of particles $`N`$. The periodic orbits are ideally those of the classical counterpart of the self-consistent mean field, which for practical purposes often is taken as a shell-model type potential.
### II.2 Coarse-graining and finite temperatures
Our present emphasis in the use of POT is not the full quantization of the spectra of finite fermion systems, but on the semiclassical description of their gross-shell structure. For this purpose we coarse-grain the quantum spectrum by a convolution of the level density (1) with a normalized Gaussian of width $`\gamma `$:
$`g_{qm}(E,\gamma )={\displaystyle \frac{1}{\gamma \sqrt{\pi }}}{\displaystyle \underset{i}{}}e^{(EE_i)^2/\gamma ^2}.`$ (4)
The coarse graining of the trace formula (2) gives, using the stationary-phase approximation for the convolution integral, an extra exponential factor in the trace formula:
$`\delta g_{sc}(E,\gamma )`$ $``$ $`{\displaystyle \underset{po}{}}𝒜_{po}(E)e^{(\gamma T_{po}/2\mathrm{})^2}\times `$ (5)
$`\mathrm{cos}[{\displaystyle \frac{S_{po}(\lambda )}{\mathrm{}}}\sigma _{po}{\displaystyle \frac{\pi }{2}}].`$
The same exponential factor appears in the trace formulae (3) for $`\delta E_{sc}(\lambda ,\gamma )`$ and $`\delta N_{sc}(\lambda ,\gamma )`$. It suppresses the contributions from orbits with larger periods $`T_{po}`$. A similar suppression of longer orbits and the overall amplitude of the shell effects occurs at finite temperatures. E.g. in a grand-canonical system at temperature $`T`$, the oscillating part of the free Helmholtz energy has the trace formula
$`\delta F_{sc}(\lambda ,T)`$ $``$ $`{\displaystyle \underset{po}{}}𝒜_{po}(\lambda )\left({\displaystyle \frac{\mathrm{}}{T_{po}}}\right)^2{\displaystyle \frac{\tau _{po}}{\mathrm{Sinh}(\tau _{\mathrm{po}})}}\times `$ (6)
$`\mathrm{cos}[{\displaystyle \frac{S_{po}(\lambda )}{\mathrm{}}}\sigma _{po}{\displaystyle \frac{\pi }{2}}].`$
Here $`\tau _{po}=k_BT\pi T_{po}(\lambda )/\mathrm{}`$ and $`k_B`$ is the Boltzmann contant. In both situations the gross-shell effects are dominated by the shortest periodic orbits of the system struma .
In mixed-dynamical and integrable systems, orbits with different degrees $`f`$ of degeneracy can coexist. The gross-shell structure then results from a competition between the periods and the degeneracies of the shortest orbits ozotom . Whereas their coarse-grained semiclassical amplitudes decrease with growing period $`T_{po}`$, they increase with growing degeneracy $`f`$ due to their dependence $`\mathrm{}^{(1+f/2)}`$ already mentioned above.
In arbitrary spherical three-dimensional systems, the most degenerate orbits undergo both radial and angular oscillations with rational frequency ratios (cf. bertab ; boys ; bm ):
$`\omega _r:\omega _\varphi =n:m,|n|,|m|.`$ (7)
(For physical reasons, only pairs of integers $`n,m`$ with equal signs are allowed; orbits with negative $`n,m`$ correspond to the time-reversed of the orbits with positive $`n,m`$). The orientiations of these ’rational tori’ can be rotated about three Euler angles without changing their shapes, periods or actions; therefore they appear in three-fold degenerate families ($`f=3`$). Their existence for arbitrary ratios $`n:m`$ depends, however, on the form of the radial potential $`V(r)`$ (cf. boys ; arita ) and, in general, on the energy $`E`$. The special orbits with angular momentum $`L=0`$ and $`L=L_{max}`$, corresponding to librating ’diameter’ and rotating ’circle’ orbits, respetively, form families with only two-fold degeneracy ($`f=2`$), since one of the three Euler rotations does not change their orientations. The Coulomb potential $`V(r)=\alpha /r`$ has an extra dynamical symmetry, leading to O(4); here all orbits have $`n:m=1:1`$. Spherical harmonic oscillators $`V(r)=ar^2`$ also have an extra dynamical symmetry, leading in three dimensions to U(3); here all orbits have $`n:m=2:1`$. In each of these two special potentials, all periodic orbits (including the diameters and circles) form one family with degeneracy $`f=4`$. Quantum-mechanically, the dynamical symmetries reflect themselves in an accidental extra degeneracy of the eigenvalue spectrum $`\{E_i\}`$. The semiclassical trace formulae for $`\delta g(E)`$ of these systems, added to their ETF expressions for $`\stackrel{~}{g}(E)`$, reproduce the exact quantum-mechanical level densities according to (1) \[see book , Eqs. (3.144) and (3.69) for their explicit analytical expressions\].
## III POT for finite fermion systems
### III.1 Nuclei
1. Ground-state deformations. Strutinsky et al. struma ; strdo were the first to extend the POT to systems with continuous symmetries and to apply it to the study of gross-shell effects in nuclei. In strdo they studied the periodic orbits in a spheroidal cavity with axis ratio $`\eta `$ as a model for the mean field of a deformed nucleus. They plotted the shell-correction energy $`\delta E`$, obtained quantum-mechanically using Strutinsky’s shell-correction method strut with the spectra of realistic deformed Woods-Saxon potentials fuhil with the same spheroidal deformations (including spin-orbit interaction), versus particle number $`N`$ and deformation $`\eta `$. The slopes of the valleys in these deformation energy surfaces $`\delta E(N,\eta )`$ corresponding to the ground-state deformations could then be correctly reproduced by the condition that the actions $`S_{po}`$ of the shortest and most degenerate periodic orbits in the spheroidal cavity (with $`f=2`$) sphero be constant. Contributions of orbits with $`f=1`$ were negligible. (The Fermi energies corresponding to the spherical magic numbers had to be adjusted, as no spin-orbit interaction was included in the cavity model.) Although this was only a qualitative result, it proved the correctness of the concept to interpret quantum-mechanical gross-shell effects semiclassically in terms of short periodic orbits of the corresponding classical system. The use of a cavity with infinitely steep walls for the mean field hereby justifies itself through the short range and the saturating property of the effective nucleon-nucleon interaction, leading to steep walls of the self-consistent Hartree-Fock potentials or their approximations by Woods-Saxon type shell-model potentials.
2. Left-right asymmetry of fission barriers. A prominent manifestation of shell effects is the ’double-humped’ fission barrier of nuclei in the actinide region sven . One particular aspect is that of the onset of a left-right asymmetry of the fissioning nuclear shapes which eventually leads to the asymmetric mass distributions of the fission fragments. Since Sven Gösta Nilsson and his group, and other scientists in Lund, were much involved in the study of this shell effect, I may dwell a little on its history.
The mixing of pairs of single-particle states with opposite parities in a spheroidal harmonic-oscillator potential was studied earlier leeing as a possible mechanism leading to ’pear-shaped’ nuclei. In 1962, S. A. E. Johansson johan took this question up and investigated the possibility of octupole-deformed fission barriers. Since Strutinsky’s shell-correction method strut did not yet exist at that time, no realistic fission barriers could be obtained qantum-mechanically with the Nilsson model. Johansson showed, however, that at the typical deformations of the actinide fission barriers predicted by the liquid-drop model (LDM) ldm , the mixing of single-particle orbits of the type used in leeing leads to an instability of the barrier against octopole shapes. Using the shell-correction method with the Nilsson model, P. Möller and S. G. Nilsson asymn obtained in 1970 the instability of the outer fission barrier against a suitable mixture of $`ϵ_3`$ and $`ϵ_5`$ deformations. Thus, the onset of the fission mass asymmetry was clearly a quantum-mechanical shell effect that could not be explained by the classical LDM model. In a detailed microscopical study, C. Gustafsson, P. Möller, and S. G. Nilsson gumni showed a year later that those pairs of single-particle states, which are most sensitive to the left-right asymmetric shapes and hence responsible for their onset, have their wavefunction nodes and extrema on parallel planes at and near the waist-line of the fissioning nucleus perpendicular to its symmetry axis.
30 years later, in a Lund-Regensburg-Dresden collaboration chaofi , this effect was studied semiclassically. The POT had been used in the same collaboration brs for cavities with the $`(c,h,\alpha )`$ shapes of fuhil , for which the mass asymmetry of the outer fission barrier of actinide nuclei had also been obtained quantum-mechanically in asyplb . The shortest periodic orbits here are families with $`f=1`$, having the axial U(1) symmetry; they are simply the diagonal, triangular and square-shaped orbits in the circular planes perpendicular to the nuclear symmetry axis. The semiclassical trace formula (with a uniform approximation for the bifurcation occurring at the onset of the neck, where the orbits in the central plane become unstable and give birth to two new parallel planes with stable orbits) was shown brs to yield realistic deformation enery surfaces $`\delta E_{sc}(c,h,\alpha )`$ in the region of the outer fission barrier, predicting its instability against the asymmetry parameter $`\alpha `$ in good agreement with the old quantum-mechanical results asyplb . (Again, the spin-orbit interaction was omitted; the Fermi energy as the only parameter was adjusted to yield the fission isomer minimum at the correct deformation.) Similarly to strdo , the valley of steepest descent through the deformation energy surface, leading over a left-right asymmetric outer saddle, is obtained by the stationary condition $`\delta S_{po}=0`$ for the shortest orbits. In chaofi , it was shown that these orbits are situated in a very small regular island of a dominantly chaotic phase space. An approximate Einstein-Brillouin-Keller (EBK) quantization of the linearized classical motion in these regular islands reproduced rather precisely the quantum-mechanical energies of those diabatic single-particle states with opposite parity which are most sensitive to the $`\alpha `$ deformations and hence quantum-mechanically responsible for the $`\alpha `$ instability of the outer barrier. Furthermore, their wavefunctions were found to have their nodes and extrema precisely in the planes near the nuclear waist-line that contain the shortest periodic orbits responsible semiclassically for the asymmetry effect.
This application of the POT represents an interesting example for the classical-to-quantum correspondence of the interplay between chaos and order: a tiny regular island in an almost chaotic phase space causes a quantum shell effect in an interacting many-body system with observable consequences in the form of the mass asymmetry of the fission fragments.
### III.2 Metal clusters
Metal clusters are interesting finite fermion systems which allow one to study the transition from atoms over molecules towards condensed matter deheer . In the simplest theoretical description, the so-called ’jellium model’, the ions are replaced by a structureless but deformable positive background and the systems of $`N`$ interacting valence electrons in the external jellium potential are studied braclu . In neutral clusters, the eletrostatic long-range forces cancel and the valence electrons are only bound by the short-ranged exchange and correlation effects. Neutral metal clusters therefore have much in common with nuclei. One difference to nuclei is that there is no measurable spin-orbit interaction in most metal clusters. The magic numbers $`N_i`$ of the smallest spherically stable clusters correspond to those of a harmonic oscillator ($`N_i=2,8,20,40`$) deheer . For the analysis of early experimental abundance spectra of small sodium clusters, the Nilsson model without spin-orbit term was successfully employed clem to interpret the regions between the spherical shell closures in terms of prolate and oblate deformations. In a self-consistent mean-field description eka1 ; genz , the average potential of clusters with $`N\text{ }>\text{ }80`$ valence electrons has steep walls like heavy nuclei, and therefore cavity models provide again a good approximation.
One early result of POT was the observation by Balian and Bloch bablo that the coarse-grained level density of a spherical cavity exhibits a pronounced beating pattern: a rapid regular oscillation, reflecting the shell structure of the spectrum, modulated by a slow oscillation reaching over some 13 - 14 shells. From the trace formula derived in bablo one sees that the beat comes about by the interference of the shortest periodic orbits of highest degeneracy ($`f=3`$), which here are the triangle ($`n:m=3:1`$) and square ($`4:1`$) orbits. (The diameter orbit with degeneracy $`f=2`$ can be neglected.) The rapid shell oscillations are determined by the average length of these two orbits, while the period of their amplitude modulation, the so-called ‘super-shell’ oscillation, is given by the difference of their lengths.
The numbers of fermions needed to reach the first super-shell node is of the order of $`8001000`$. Super-shells can therefore not be seen in nuclei. Neutral metal clusters, however, can be made arbitrarily large. This inspired Nishioka et al. nishi to study the super-shell structure in Woods-Saxon potentials fitted to self-consistent mean fields eka1 ; bra89 , and to predict that it should be observable in metal clusters. Indeed, the super-shells were experimentally observed, for the first time in supersonic beams of hot sodium clusters klavs . Their abundance in an adiabatically expanding beam is dominated by their stability against evaporation of single atoms and exhibits pronounced peaks at the spherically magic numbers. The larger ones, $`N_i=92,138,192,264,\mathrm{}`$, are almost exactly those of a spherical cavity. By a suitable extraction of the oscillating part of the abundance spectra, the super-shell beat could clearly be exhibited klavs , with its first node appearing around $`N900`$.
Plotting the cube roots of the magic numbers, $`N_i^{1/3}`$ which are proportional to the r.m.s. radii of the magic clusters, against the shell number $`i`$, one obtains a straight line with a slope $`s=N_{i+1}^{1/3}N_i^{1/3}`$. The experimental value of this slope is $`s_{exp}=0.61\pm 0.01`$; it was confirmed by later experiments, also with other types of metal clusters (their different Wigner-Seitz radii do not affect the value of $`s`$, as long as they are single-valenced) deheer ; braclu . Around $`i14`$, there is a slight discontinuity in the plot $`s(i)`$ before it continues again with the same slope. This is due to the phase change of the rapid shell oscillations when passing through the first super-shell node. The value of this slope predicted by the POT of bablo ; nishi is $`s_{pot}=0.603`$, that of the self-consistent quantum-mechanical mean-field calculations is $`s_{mf}=0.61`$, both in perfect quantitative agreement with experiment.
This provides another example of the good agreement of POT with both experiment and quantum mechanics. An easily readable account of the super-shells in metal clusters is given in sciam .
Unfortunately, it is very difficult to extract spectroscopic data on the electronic single-particle states in metallic clusters. The most direct access to their shapes is given by the so-called ’Mie plasmons’, the collective dipole oscillations of the valence electrons against the ions eka2 ; bra89 (which are the origin of the colors in stained-glass windows krvol ). Similarly to the giant-dipole resonances in nuclei, their splitting gives evidence of the average deformations of the clusters. This gave rise to a series of theoretical investigations of cluster deformations using deformed shell models (see braclu ). For applications of the POT one employs hereby most comfortably the spheroidal cavity model used also in struma . I refer to frauen for a detailled account containing also a nice application of the POT to the semiclassical interpretation of moments of inertia. The effects of weak magnetic fields on spherical metal clusters were studied in kaoclu . In pash , spheroidal cavities with the lowest multipole deformations $`ϵ_2`$, $`ϵ_3`$, and $`ϵ_4`$ were studied, and the perturbative trace formula of crpert (cf. also peter ) was used to predict their ground-state deformations. The results were in very good agreement with those of quantum-mechanical shell-correction calculations using the spectra of the same cavities.
### III.3 Semiconductor nanostructures
Semiconductor heterostructures can be used to construct two-dimensional systems of quasi-free electrons on the nanometer scale. With the help of external metallic gates or lithography, the electrons can further be laterally confined to form so-called quantum dots, quantum channels, quantum wires, antidot superlattices, etc. nano ; reimat . Applying a perpendicular magnetic field $`B`$, one can measure the magneto-resistance of such devices. Under suitable experimental circumstances, both the mean free path and the phase coherence length can be made larger than the sizes of these structures, so that quantum interference still takes place while the dimensions are large enough to allow for a semiclassical description. Nanostructues are therefore ideal tools to study the interplay between classical and quantum mechanics.
Weiss et al. weiss measured the resistance of andidot superlattices and found oscillations which can be explained classically by the commensurability of cyclotron orbits with the superdot lattice: when an electron is trapped in a cyclotron orbit that fits around 1, 2, 4, 9, etc. antidots, it does not contribute to the conductance and hence a peak is seen in the magneto-resistance (see weri for an easily readable account). In weak $`B`$ fields at very low temperatures, some rapid $`B`$-periodic oscillations could be observed. They could be interpreted semiclassically by the interference of different trapped periodic orbits of comparable lenghts; the linear response of the system to the $`B`$ field was hereby described by a semiclassical version of the Kubo theory yielding a trace formula for the conductance kubosc . These so-called ’Aharonov-Bohm (AB) oscillations’ can be easily understood in a perturbative approach crpert in which the effect of the magnetic field is taken into account to lowest order only in the actions $`S_{po}=𝐩𝑑𝐪`$, while the shapes of the orbits are left unchanged. Under the canonical substitution $`𝐩𝐩(e/c)𝐀`$, where $`𝐀`$ is the vector potential with $`𝐁=\times 𝐀`$, the action of a periodic orbit changes like
$`S_{po}S_{po}(e/c){\displaystyle 𝐀𝑑𝐪}=S_{po}(e/c)\mathrm{\Phi }_{po},`$ (8)
where $`\mathrm{\Phi }_{po}=𝐁𝑑𝐅_{po}`$ is the magnetic flux through the area $`F_{po}`$ surrounded by the orbit. Consequently, the perturbed semiclassical trace formula is modified only by a factor $`\mathrm{cos}(e\mathrm{\Phi }_{po}/\mathrm{}c)`$ containing the Aharonov-Bohm phase which causes the $`B`$-periodic oscillations.
Similar AB oscillations have also been measured in the magneto-conductance of a circular quantum dot containing some $`1200`$ to $`2000`$ electrons pelrei . They could be qualitatively well explained qdot by the perturbed level density $`\delta g(E,B)`$ of a two-dimensional circular billiard. This is an integrable system whose trace formula is analytically known disk . The conductance oscillations as functions of the radius of the quantum dot (regulated experimentally by the applied gate voltage) were well reproduced by the average length of the shortest orbits (here: diameters and triangles), while the period of the AB oscillations according to (8) was well reproduced by the area enclosed by the triangular orbit qdot .
An analytical trace formula for the circle billiard in arbitrarily strong transverse magnetic fields $`B`$ was given in blaqd , and the magnetization of quantum dots was studied semiclassically in kaoqd .
In a mesoscopic semiconductor channel with two antidots, the magneto-conductance was measured in goul and also found to exhibit AB oscillations. These could be well explained jphd using the semiclassical Kubo formula kubosc with a suitably modeled two-dimensional confinement potential for the channel (with antidots), which represents a mixed-dynamical system with a rather chaotic phase space. Plotting the maxima of the experimental AB oscillations versus magnetic field $`B`$ and gate voltage (which regulates the radii of the antidots), one obtains a grid of lines exhibiting some characteristic displacements goul . At first sight, these might be attributed to missing flux units. Quantum-mechanical calculations kirc reproduced these displacements but could not explain the physics behind them. In the semiclassical calculations the displacements could, in fact, be attributed to bifurcations of some of the trapped periodic orbits jphd ; chan .
### III.4 Trapped dilute atomic gases
I mention only briefly the finite fermion systems produced by confining diluted fermionic atom gases, e.g. in magneto-optic traps grimm . Hartree-Fock calculations for $`N`$ harmonically trapped atoms with a short-ranged repulsive interaction yylund yielded shell effects $`\delta E(N)`$ in their total binding energies $`E_b(N)`$ which remind about the super-shells discussed above. More details are given in magnus ; let me just emphasize here that the origin of the beating shell structure is different here from that in a spherical cavity bablo . The self-consistent mean field can be modeled by a perturbed harmonic oscillator $`V(r)=ar^2+ϵr^4`$. For such potentials it was shown recently boys that the gross-shell structure is dominated by the two-fold degenerate diameter and circle orbits, whose interference explains the super-shells found in yylund . The shortest three-fold degenerate orbits have frequency ratios $`n:m7:3`$ and contribute only to finer details of the quantum spectrum at relatively high energies. |
warning/0506/hep-ex0506006.html | ar5iv | text | # STUDIES OF BEAUTY AT H1 AND ZEUS
## 1 Introduction
At HERA (DESY, Hamburg) electrons (positrons) are brought into collision with protons at a centre-of-mass energy of 318 GeV (300 GeV up to 1997). From 1996 to 2000 the two experiments H1 and ZEUS collected about $`100\mathrm{pb}^1`$, on which the results presented here are based. The measurement of beauty production still is a rich testing ground for QCD, as several hard scales can be chosen, depending on the process to be described. The beauty mass always provides one hard scale, but due to the presence of additional scales, calculations are not straight forward. In HERA ep collisions, beauty quarks are predominantly produced in boson gluon fusion; a photon from the lepton and a gluon from the proton collide to produce a $`\mathrm{b}\overline{\mathrm{b}}`$ pair. This process introduces several scales relevant for the calculation of the cross sections:
| mass of the b quark | $`m_b`$ | $`5\mathrm{GeV}`$ | |
| --- | --- | --- | --- |
| transverse momentum of the b quark | $`p_T^b`$ | $``$ typically | a few GeV |
| virtuality of the exchanged photon | $`Q^2`$ | $`1\mathrm{GeV}^2`$ | $`\text{Photoproduction (}\gamma \text{p})`$ |
| | | $`1\mathrm{GeV}^2`$ | $`\text{Deep Inelastic Scattering (DIS)}`$ |
The optimal scheme for Next to Leading Order (NLO) QCD predictions depends on the dominant scale. If the virtuality, $`Q^2`$, and the transverse momentum of the b quark squared, $`(p_T^b)^2`$, are of the order of the b mass squared, $`m_b^2`$, threshold effects due to the b mass need to be considered and the so called massive scheme (used in the Fixed Flavour Number Scheme - FFNS) is used. Once $`Q^2`$ or $`(p_T^b)^2`$ are much larger than $`m_b^2`$, massless schemes (used in the Zero-Mass Variable Flavour Scheme - ZM-VFNS) can be used. Here the b quark is treated as a massless parton in the hard scatter. In the Variable Flavour Number Scheme (VFNS) the massive and massless approaches are combined. Differential cross sections are compared to massive NLO predictions by FMNR for $`\gamma `$p and HVQDIS for DIS. For references to all schemes, NLO predictions and MC generators please refer to ref. $`^{\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?}}`$.
## 2 Measurements in $`\mu `$+jet(s)
In events with a jet and a muon, two properties of the b hadrons were exploited to tag beauty events: the large b mass and the long lifetime. The b mass leads to a relatively large transverse momentum of the muon relative to the closest jet, $`p_T^{\mathrm{rel}}`$. The long lifetime of the b was used by measuring the distance of closest approach in the transverse plane, $`\delta `$, of the muon track to the beam spot or main vertex. A sign was added to $`\delta `$, which was set positive if the muon track crosses the jet axis downstream of the beam spot, else it was set negative.
### 2.1 Beauty in $`\gamma `$p in $`\mu `$+jets
H1 has measured differential cross sections in events with a muon and 2 jets $`^\mathrm{?}`$ in the photoproduction regime for $`Q^2<1`$ GeV<sup>2</sup> using an integrated luminosity of 50 pb<sup>-1</sup>. To tag beauty events a simultaneous fit to both the $`p_T^{\mathrm{rel}}`$ and the signed $`\delta `$ distribution was performed. The cross sections obtained are given in fig. 1. ZEUS performed a similar measurement $`^\mathrm{?}`$ relying on $`p_T^{\mathrm{rel}}`$ only in a wider muon eta range of $`|\eta ^\mu |<2.5`$.
Differential cross sections in $`p_T^\mu `$ and $`\eta ^\mu `$ are shown in fig. 1 together with the massive NLO QCD prediction (FMNR) which has been corrected, using LO+PS MC simulations, to describe the measured hadron level properties. Reasonable agreement of the NLO prediction with data is observed. The data are a bit steeper in $`p_T^\mu `$ in the H1 measurement, see fig. 1a. H1 and ZEUS agree within errors, as shown in fig. 1a and b.
### 2.2 Beauty in DIS in $`\mu `$+jet
Beauty cross sections were also measured for $`Q^2>1`$ GeV<sup>2</sup> using the same techniques as described above. H1 has used a $`Q^2`$ range of $`2<Q^2<100\text{GeV}^2`$. The H1 cross sections were obtained for events with $`p_T^\mu >2.5`$ GeV, $`0.75<\eta ^\mu <1.15`$, $`0.1<y<0.7`$, $`p_{T\mathrm{jet}}^{\mathrm{Breit}}>6`$ GeV and $`|\eta _{\mathrm{jet}}^{\mathrm{lab}}|<2.5`$ where the index ’Breit’ refers to the Breit frame. The $`\eta ^\mu `$ cross section is shown in fig. 2a. ZEUS has performed a similar measurement in DIS (fig. 2b) using data in the $`Q^2`$ range of $`1<Q^2<1000\text{GeV}^2`$ and defining the kinematic range as $`p_T^\mu >2.0`$ GeV, $`1.6<\eta ^\mu <1.3`$, $`0.05<y<0.7`$, $`E_{T\mathrm{jet}}^{\mathrm{Breit}}>6`$ GeV and $`2.0<\eta _{\mathrm{jet}}^{\mathrm{lab}}<2.5`$ $`^\mathrm{?}`$. The $`p_T^\mu `$ spectrum from both H1 and ZEUS is slightly steeper than the prediction (not shown). An interesting feature is seen in the $`\eta ^\mu `$ distribution in fig. 2: for both H1 and ZEUS (kinematic ranges differ) the cross section rises towards higher $`\eta ^\mu `$, a trend that is not well reproduced by the NLO predictions.
## 3 Inclusive measurements using impact parameters
Events with heavy flavour mesons can contain several tracks with large impact parameters. To inclusively determine the beauty contribution, events were selected in H1 by inclusive lifetime tagging, using all $`p_T>0.5`$ GeV tracks with reasonable precision vertex tracker information. For these tracks the impact parameter significance $`S_{1(2)}=\delta _{1(2)}/\sigma _{\delta _{1(2)}}`$ was calculated for the track with the largest (second largest) impact parameter $`\delta `$ (with $`|\delta |<0.1\text{cm}`$). The sign of the impact parameter was chosen relative to the direction of the closest jet if present or from the hadronic final state, also requiring $`sign(S_1)=sign(S_2)`$. This method yields high statistics and a good separation of beauty from charm and light flavour processes.
### 3.1 Beauty contribution to F<sub>2</sub>
To determine the beauty contribution to the proton structure function, double differential beauty cross sections in bins of $`x`$ and $`Q^2`$ were measured using inclusive lifetime tagging. The resulting $`\mathrm{F}_2^{\mathrm{b}\overline{\mathrm{b}}}`$ is shown in fig. 3. Good agreement with both the H1 PDF 2000 fit $`^\mathrm{?}`$ and MRST03 $`^\mathrm{?}`$ was found.
The integrated beauty cross section for $`Q^2>150\text{GeV}^2`$ and $`0.1<y<0.7`$ has been determined to be $`\sigma ^{b\overline{b}}=55.4\pm 8.7\pm 12`$ pb. This is compatible with several predictions: H1 PDF 2000 (massless, ZM-VFNS) $`\sigma ^{b\overline{b}}=52`$ pb, NLO (massive$``$massless, VFNS) $`\sigma ^{b\overline{b}}=47`$ pb, and NLO (massive, FFNS) $`\sigma ^{b\overline{b}}=37`$ pb.
## 4 Conclusion and Outlook
NLO QCD predictions generally agree with the data, though they are somewhat below the data in some kinematic regions. H1 and ZEUS measurements with muons + jet(s) in $`\gamma `$p and DIS agree. The measurement of the beauty contribution to F<sub>2</sub> at high $`Q^2`$ using inclusive lifetime tagging is in good agreement with NLO.
The upgraded HERA II detectors are giving promising results. ZEUS presented a study of beauty in photoproduction based on the first 20 pb<sup>-1</sup> taken in 2003/4 using the new ZEUS vertex detector $`^\mathrm{?}`$. The impact parameter distribution, as shown in fig. 4a, has been validated by comparing the difference of the positive and negative parts of the impact parameter distribution with MC which in turn has been normalised using the $`p_T^{\mathrm{rel}}`$ method, as shown in fig. 4b.
## References |
warning/0506/cond-mat0506572.html | ar5iv | text | # Ultracold atoms in optical lattices with random on-site interactions
## Abstract
We consider the physics of lattice bosons affected by disordered on-site interparticle interactions. Characteristic qualitative changes in the zero temperature phase diagram are observed when compared to the case of randomness in the chemical potential. The Mott-insulating regions shrink and eventually vanish for any finite disorder strength beyond a sufficiently large filling factor. Furthermore, at low values of the chemical potential both the superfluid and Mott insulator are stable towards formation of a Bose glass leading to a possibly non-trivial tricritical point. We discuss feasible experimental realizations of our scenario in the context of ultracold atoms on optical lattices.
Ultracold atomic gases have attracted large interest during the last years, in particular due to the experimental achievement of quantum degeneracy both for bosonic BEC and fermionic Fermi gases. Very recently, the extraordinary experimental developments in the manipulation of ultracold atoms have opened the path towards a new fascinating research area, namely the analysis of strongly-correlated atomic gases. In this sense, remarkable experimental results have been already reported, including experiments on the so-called BEC-BCS crossover in degenerate Fermi gases BCSBEC , and experiments on ultracold gases in optical lattices, such as e.g. the observation of the Mott insulator-to-superfluid transition Greiner ; ETH and the experimental realization of a Tonks-Girardeau gas Paredes .
Up to now the experiments performed in optical lattices have considered a defect-free square lattice potential. However, recent proposals have discussed new scenarios which depart from this standard. In particular, various lattice geometries are achievable by standard laser techniques, and should lead to fascinating physics Kagome ; triangular . On the other hand, disorder may be induced in the laser potential in different ways. Localized impurities can be created, leading to Kondo-like physics Cirac . Randomness can be produced by means of additional incommensurable lattices, or by laser speckles BG , and may lead to Anderson localization Anderson and Bose glass phases Fisher . First experimental steps towards random potentials have already been achieved Inguscio . Additionally, randomness may be induced in the hopping rates in the optical lattice, leading to Fermi glasses, spin glasses and quantum percolation Sanpera .
One of the most fascinating possibilities for the control of atomic gases is provided by the manipulation of the interatomic interactions by means of Feshbach resonances Feshbach . Using a suitably adjusted magnetic field a resonance is induced between a molecular state and the unbound states of the incoming atoms, leading to important modifications of the scattering properties. The $`s`$-wave scattering length undergoes a very large change in a relatively narrow window of values of the applied magnetic field, becoming $`\pm \mathrm{}`$ at the resonance. This extreme sensitivity of the scattering properties at the verge of a Feshbach resonance is crucial for our discussion.
In this Letter, we consider the physics of lattice bosons subject to a novel kind of disorder, namely bounded disorder in the strength of the interatomic interactions Adrian . In particular, we show how this type of disordered system can be realized in the context of ultracold atoms near a Feshbach resonance. We discuss the consequences that the disorder in the interaction strength has on the zero-temperature phase diagram of the system by means of a strong-coupling expansion (SCE) and quantum Monte Carlo (QMC) simulations, and contrast our findings to the case of randomness in the chemical potential Fisher .
The above scenario can be realized in a gas of ultracold bosons confined to an optical lattice using state-of-the-art experimental techniques on atom chips Schmiedmayer . In the following, we consider a one-dimensional configuration, where the atoms are assumed to be strongly confined in the other two dimensions (by e.g. additional optical confinement). However, qualitatively similar results as those discussed in this Letter also hold in higher dimensions. We assume that the gas is brought at the verge of a Feshbach resonance by an off-set magnetic field, where, as discussed previously, slight modifications of the magnetic field lead to large variations of the scattering properties. We furthermore consider the bosonic gas to be close to a magnetic wire, inducing a spatially random magnetic field Schmiedmayer , which can be considered sufficiently weak, such that the variations in the potential energy can be considered negligible when compared to other energy scales. However, since the off-set field sets the system at the verge of a Feshbach resonance, the slight random variations of the magnetic field lead to a spatially random variation of the local interatomic interactions (see Fig. 1). We assume that the variations of the scattering properties are bounded, i.e. the spatial variations of the magnetic field do not induce a crossing of the resonance.
In order to avoid significant van der Waals losses, the atoms should be placed at a distance $`>0.5\mu `$m from the surface of the atom chip, which restricts the characteristic wavelength of the variations of the magnetic field at the wire to $`\mathrm{\Delta }x1\mu `$m. The inter-site separation in the optical lattice is approximately $`\lambda =400`$nm, but can be increased to $`1\mu `$m by setting an angle between the counter-propagating lasers in such a way that $`\mathrm{\Delta }x\lambda `$. The variation of the magnetic field at the wire can be adapted to typical Feshbach resonance widths, which are of the order of $`110`$ mG. Fig. 1 shows an example of the variation of the on-site interaction for a typical experimental setup.
Under these conditions, the system is described by a Bose-Hubbard Hamiltonian with random interatomic interactions:
$$\widehat{H}=J\underset{<ij>}{}(b_i^{}b_j+h.c.)+\underset{i}{}\frac{U_i}{2}n_i(n_i1)\underset{i}{}\mu _in_i,$$
(1)
where $`b_i`$ ($`b_i^{}`$) is the annihilation (creation) operator for bosons in the $`i`$-th lattice site, $`n_i=b_i^{}b_i`$, $`J`$ is the hopping constant, $`\mu _i=\mu `$ is the local chemical potential, and $`U_i`$ denotes the local value of the coupling constant for the interatomic interactions. We assume $`U_i`$ uniformly distributed inside the range $`U(1ϵ)U_iU(1+ϵ)`$ footnote .
In order to assess the effects of randomness in the on-site interactions and to contrast it to the case of a random chemical potential, we shortly review the latter case, where $`\mu \mathrm{\Delta }\mu _i\mu +\mathrm{\Delta }`$, at a fixed $`U_i=U`$ ($`ϵ=0`$Fisher : The phase diagram in the strongly interacting regime can be obtained from considering first the trivial limit, $`J=0`$: One finds that if $`\mu `$ falls within an interval $`(n1)U+\mathrm{\Delta }\mu nU\mathrm{\Delta }`$, with an integer $`n`$, the ground-state of the system is an incompressible Mott insulator with $`n`$ bosons on each lattice site. For values of $`\mu `$ outside these intervals, i.e. for $`nU\mathrm{\Delta }<\mu <nU+\mathrm{\Delta }`$, the system enters an insulating, but compressible phase commonly referred to as Bose glass, in which the occupation per site varies between $`n`$ and $`n+1`$. This physics is maintained up to moderate values of $`J`$, leading to a phase diagram with three distinguished phases: superfluid (SF) at large $`J`$, Bose glass (BG), and the Mott insulator (MI) regions. Recent quantum Monte Carlo calculations have shown that the transition from SF to MI always occurs through an intermediate BG phase Svistunov , as suggested in Ref. Fisher .
A characteristic difference to the case of random $`U`$ at a fixed value of $`\mu _i=\mu `$ ($`\mathrm{\Delta }=0`$), is already obtained in the limit $`J=0`$: Now compressible phases occur for $`n(1ϵ)<\mu /U<n(1+ϵ)`$, in which the occupation per site varies between $`n`$ and $`n+1`$, whereas outside these regions the system enters into a MI. In contrast to the case of randomness in the chemical potential, the extent of compressible regions between the MI thus increases with the occupation number $`n`$. The disorder in $`U`$ thus leads to an increased destabilization of MI states with increasing $`n`$. Eventually, MI regions with an occupation of $`n(1+ϵ)/2ϵ`$ particles per site disappear. Therefore, given a disorder strength $`ϵ`$, only a finite number of MI regions remain stable. This selective destruction of MI states cannot be accessed in the case of a random chemical potential, where all the MI lobes vanish once $`\mathrm{\Delta }U/2`$. More generally, if both $`U_i`$ and $`\mu _i`$ have a bounded random distribution, the $`n`$-th lobe disappears for $`(2n1)ϵ+2\mathrm{\Delta }/U=1`$.
In order to analyze the phase diagram for finite hopping $`J>0`$, we first estimate the extent of the MI regions using the SCE Monien . This method allows for a quantitative calculation of the boundaries between compressible and incompressible phases for sufficiently low $`J`$, particularly in low dimensions, where the coordination number is small.
From the SCE one obtains the energy gap for adding or removing a particle from the MI, and performs an expansion in the small parameters $`J/U`$, $`\mathrm{\Delta }/U`$ and $`ϵ`$. For a given value of $`J`$, these gaps determine the boundaries of the MI lobes, which are obtained in first order as
$`{\displaystyle \frac{\mu _u}{U}}`$ $`=`$ $`2(n+1){\displaystyle \frac{J}{U}}{\displaystyle \frac{\mathrm{\Delta }}{U}}+n(1ϵ)`$ (2)
$`{\displaystyle \frac{\mu _l}{U}}`$ $`=`$ $`2n{\displaystyle \frac{J}{U}}+{\displaystyle \frac{\mathrm{\Delta }}{U}}+(n1)(1+ϵ)`$ (3)
for the upper ($`u`$) and lower ($`l`$) boundaries of the MI with integer filling $`n`$. In particular we find, that the lower boundary of the first MI lobe ($`n=1`$) is not affected by the disorder in $`U`$, in contrast to the random-$`\mu `$ case Monien . This reflects the observation, that no compressible phase emerges below the $`n=1`$ MI at $`J=0`$ for the random-$`U`$ case, in contrast to the case of random $`\mu `$.
From equating the two boundaries one obtains an (under-) estimate for the largest extent of the MI phase with filling $`n`$:
$$\frac{J_c}{U}=\frac{12\mathrm{\Delta }/U(2n1)ϵ}{2(2n+1)}.$$
(4)
Calculating higher-order terms in the strong-coupling expansion improves the quantitative description of the MI lobes, but does not affect our qualitative conclusions. In particular, all the higher order terms in the random interaction are at least of order $`(J/U)^2ϵ`$. Considering that the critical hopping $`J_c/U1`$, even relatively large disorder strengths may be approximated using the first order shift of the phase boundaries. For the figures presented below, we used expansions up to third order and also considered finite size effects Monien . The latter, as shown below, are rather significant, and converge only very slowly to the thermodynamic limit when the number of lattice sites considered increases.
While the SCE allows to estimate the extent of the Mott-insulating phases quantitatively, it does not provide information about the complete phase diagram, in particular concerning the nature of the compressible phases and the presence of SF and BG phases. In order to obtain a more complete phase diagram for bosons with random on-site interactions we have performed QMC simulations using the stochastic series expansion method sse with directed loop updates directedloop1 ; directedloop2 . We performed simulations for periodic chains with $`L=200`$ sites averaged over typically 200 disorder realizations for each data point. In the simulations, the temperature $`T`$ was chosen low enough to obtain ground state properties of the finite systems footnoteT . In order to reliably identify the various phases of the system, we measured the global compressibility $`\kappa =N/\mu `$ from the fluctuations of the particle number $`N`$, and the winding number fluctuations $`W^2`$, from which the superfluid density is obtained as $`\rho _S=LT/(2J)W^2`$ rho\_S . Fig. 2 shows the disorder averaged values of both observables as functions of $`J/U`$ along a cut of constant $`\mu /U=0.51`$ for $`ϵ=0.25`$, which at $`J=0`$ belongs to the MI region with $`n=1`$. At low values of $`J`$, an extended MI region with $`n=1`$ is clearly identified by vanishing values of both $`\kappa `$ and $`\rho _S`$. For large value of $`J/U`$, both $`\kappa `$ and $`\rho _S`$ take on finite values, thus identifying the large-$`J`$ region as a SF phase. In the intermediate region, for $`0.078J/U0.133`$, the system shows a finite compressibility, but no superfluid response. This regime is thus characterized as a disorder-induced BG phase, separating the MI and SF region at this value of $`\mu /U=0.51`$.
Performing similar scans along different lines in parameter space, we obtain QMC estimates of the various phase boundaries, leading to the phase diagram of Fig. 3. Here, we combined the QMC results for $`ϵ=0.25`$ with third-order SCE results for the extent of the MI phases in the thermodynamic limit (TDL). Within the chosen resolution, the finite size corrections of the phase boundaries obtained by QMC show only a weak dependence on the system size. An accurate determination of the MI boundaries from QMC in the disordered case scalettar would require significantly larger system sizes due to rare regions of delocalized bosons Monien . In fact, performing a third-order SCE for a finite chain of $`L=200`$ sites compares well to the estimated extend of the MI regions from QMC (Fig. 3) in the strong-coupling regime, $`J/U0.1`$, and indicates the rather strong finite size effects in the disordered case.
Our numerical calculations confirm the qualitative picture considered above: The MI lobes clearly shrink, not only in the $`\mu `$-direction but also in $`J`$, and eventually disappear for sufficiently large values of $`\mu `$. In the case considered in Fig. 3, where $`ϵ=0.25`$, only the MI lobes with $`n=1`$ and $`2`$ remain stable. Furthermore, we observe that the lower boundary of the $`n=1`$ MI does not vary significantly with the disorder strength up to the critical hopping $`J_c`$. The reduced relevance of disorder in the interaction strength in the low-density region is also reflected by the complete absence of a BG phase for $`\mu <0`$. Therefore, a dilute compressible lattice boson gas remains superfluid even in the presence of disordered interactions. This is in clear contrast to the case of random $`\mu `$, where the BG extends down to arbitrary low but finite densities. The absence of a BG for $`\mu <0`$ indicates a tricritical point for the SF, MI, and BG phases along the lower boundary of the $`n=1`$ MI. However, our calculations cannot reliably distinguish whether the BG phase disappears exactly at a tricritical point at $`\mu =J=0`$, or if this occurs for finite values of $`\mu ,J>0`$. This interesting question thus remains open for future studies.
To summarize, in this Letter we analyzed the novel physics occurring when the on-site interaction in a Bose-Hubbard Hamiltonian acquires a bounded random character. We discussed how this exciting possibility can be realized in ultracold atomic gases in optical lattices at the verge of a Feshbach resonance, when they are brought in the vicinity of weak spatially random magnetic fields generated by magnetic wires. As for the case of random chemical potentials, a phase diagram with superfluid, Bose-glass and Mott-insulator phases is predicted. However, important differences can be found when comparing the random $`U`$ and random $`\mu `$ cases: For random interactions, the Mott-insulator phases become progressively narrower (both in the $`\mu `$ and $`J`$ direction) for larger occupation numbers, and eventually disappear beyond a given filling. Furthermore, the Bose-glass phase disappears for low chemical potentials and hopping, and hence a tricritical point occurs, although our numerical calculations cannot discern whether this occurs at $`\mu =J=0`$ or very small but finite values.
Our calculations were performed in absence of an overall (usually harmonic) confining potential, whereas the necessary commensurability for Mott-insulator phases can be achieved in practice only in the presence of inhomogeneous potentials. However, if the confining potential is sufficiently shallow, a local chemical potential may be considered, and local phase occur, as thoroughly demonstrated by means of QMC calculations Rigol ; Wessel for the pure case. Of course, the finite size of these different regions could significantly affect the boundaries between the Bose glass and superfluid, once the localization length becomes smaller than the spatial extent of the various regions. Typical experiments on the Mott-insulator phase of atoms in optical lattices Greiner ; ETH rely on the observation of the broadening of interference fringes in time-of-flight pictures (provided by the insulating character of the phase), and the monitoring of the opening of the Mott-gap in the excitation spectrum by means of tilting experiments. In these experiments, although other spatial phases are expected, the Mott-insulator region is the largest one and dominates the experimental detection. Similar experiments could be performed to detect the shrinking of the Mott-insulator phases as a function of the disorder in $`U`$ predicted in this paper.
Fruitful conversations with A. Kantian, M. Lewenstein, and P. Öhberg are acknowledged. This work was supported by the Alexander von Humboldt Foundation and by NIC at FZ Jülich. H.G. thanks the Studienstiftung des deutschen Volkes for support. |
warning/0506/astro-ph0506061.html | ar5iv | text | # Probing dark matter caustics with weak lensing
## 1 Introduction
The nature of dark matter, which constitutes about $`30\%`$ of the mass of the Universe, remains largely unknown. Results from cosmic microwave background explorers and large-scale galaxy surveys suggest that dark matter is cold with little velocity dispersion (e.g. Spergel et al. 2003; Tegmark et al. 2004). If so, then its evolution is mainly governed by its self-gravity and expressed by Jeans-Vlassov-Poisson equations. The collisionless nature of dark matter predicts the formation of multi-stream regions bounded by very high density manifolds known as caustics.
In a cold dark matter Universe, caustics can form on large scales of many megaparsecs, manifesting in the filamentary structure of galaxy distribution (Shandarin & Zeldovich 1989) and also at smaller scales of a few parsecs or kiloparsecs in dark matter haloes. This is well illustrated in some pionering numerical works (Melott & Shandarin 1989). Due to their abundance, the rich observational and numerical data and their high density contrasts, haloes are likely areas for the caustics.
Analytic models for the formation and evolution of dark matter haloes are still rare and most works are based on the selfsimilar accretion model (Gott 1975; Gunn 1977; Fillmore & Goldreich 1984; Bertschinger 1985). In this model, first proposed to explain the rotation curve of galaxies, haloes form by temporally self-similar collapse of dark matter shells onto an initially over-dense perturbation. Dark matter shells initially expand until they reach their turnaround radii where they separate from the background expansion and collapse. After collapse they re-bounce and collapse again and the density profile settles asymptotically into a power-law which is convolved with singular spikes, namely with caustics.
In the following we shall refer to such caustics as ”outer caustics”. They are suggestive of the sharp stellar shells observed around giant ellipticals which can arise in the merger of galaxies (Malin & Carter 1980; Quinn 1984; Fort et al. 1986). The main observational difference between the merger and selfsimilar spherical infall is that the former predicts that the caustics are interleaved with the caustic radii alternating on opposite sides of the galaxy and the latter predicts concentric spherical shells.
The spherically symmetric model has been extended to consider infalling matter with angular momentum and calculated the properties of an additional kind of ”inner caustics” with torus-like topology (Sikivie 1998, 1999). Such accretion with angular momentum is more relevant for galaxy-size haloes than for clusters of galaxies dominated by radial infall.
Here, we focus on the outer caustics of cluster-size haloes and argue that they can be reasonably approximated by selfsimilar infall models. In the original version of this model, dark matter is absolutely cold, i.e. with zero velocity dispersion, and caustics are infinitely thin concentric spherical shells with diverging densities. However, realistically, dark matter has a small velocity dispersion and these shells have finite thicknesses. The thicknesses of the caustics would however remain very small (due to the coldness of dark matter) and thus they would contain very little mass in spite of their significant density. Various characteristics of the caustics such as their density profile, their thickness and their approximate maximum density for a low velocity dispersion dark matter medium have been recently evaluated (Mohayaee & Shandarin 2006, hereafter MS05).
Detection of dark matter caustics remains a challenging problem both for 3-dimensional numerical simulations and for observations. Dark matter annihilation in the caustics has already been studied (See e.g. Sikivie et al. 1997; Hogan 2001; Mohayaee & Shandarin 2006; Pieri & Branchini 2005). Since the flux of the annihilation products, e.g. $`\gamma `$-rays or antiproton flux, depends on the square of the local density<sup>1</sup><sup>1</sup>1See e.g. Donato et al. (2004) for the antiproton flux and Bertone & Merritt (2005) for a general recent review., caustics with their sharply-peaked densities would be the likely places for significant emissivity. Such kind of observations are promising for the eventual detection of dark matter in the caustics. However, the major problem with dark matter detection through gamma-ray emission is the severe background contamination of the signal and the low signal-to-noise ratio of present day observations.
It has been shown that rotation curves of galaxies might be sensitive to the presence of inner caustics and claimed a marginal detection based on an ensemble average over 32 rotation curves (Kinney & Sikivie 2000). Inner caustic rings are located in the plane of the disk and are likely to modify rotation curves more efficiently than outer caustics with spherical symmetry. Moreover, rotation curves only probe the dark matter potential of galactic haloes where a large amount of tracers (gas or stars) is available. Therefore, outer caustics which are located well beyond the observable tracers cannot leave a detectable imprint on rotation curves.
Gravitational lensing provides a promising alternative tool. Since lensing probes the projected density profile with no regards to the nature or dynamical state of the deflecting mass, it should be sensitive to the caustics. The lensing properties of dark matter caustics and particularly their efficiency in magnifying and/or producing multiple images of background sources has already been investigated (Charmousis et al. 2003).
Inner caustics may be dense enough to produce substantial magnification and small separation multiple images as in micro-lensing by compact objects. High magnification events due to such caustics may explain the anomaly in flux ratios observed in multiply-imaged quasars which are hardly reproduced by a smooth halo or even subhalos (see e.g Dalal & Kochanek 2002; Kochanek & Dalal 2004).
Conversely, it has been found that outer caustics are inefficient in magnifying distant sources and would yield at most a few percent net magnification or shear (Charmousis et al. 2003). As we shall detail below, the selfsimilarity of dark matter accretion implies that outer caustics occur at the same radius provided physical radii are properly rescaled by the halo mass. Hence the tiny lensing signal of caustics could show up statistically by averaging over many rescaled haloes. Such a statistical approch is much more complicate for the lensing-based detection of inner caustic rings because of the random orientation of rings (or angular momenta).
In this work, we consider the lensing properties of the outer caustics only and we propose the weak-lensing effect as a potential way to detect caustics. We demonstrate that the caustics will produce sharp variations in the projected surface mass density around haloes. Depending on the height and width of caustics, gravitationally-distorted background galaxies will experience local variations of shear.
If the aforementioned universal property of haloes is fullfilled, then the rapid progress in X-ray and lensing observations of cluster of galaxies may offer one possibility to observe dark matter caustics. Deep and wide surveys such as CFHTLS will cover fields of view of a few hundreed square degrees and will provide us with a useful material for the detection of dark matter caustics through the capability of mass to coherently stretch the image of the background galaxies. They will provide us with a large enough number of galaxy clusters<sup>2</sup><sup>2</sup>2$`5`$ per square degree, (Hennawi & Spergel 2005) to achieve the required level of signal-to-noise ratio. Wide field spatial surveys will be even more powerful for the investigation of the lensing properties of dark matter haloes and their associated caustics. An important requirement for the detection of caustics is the measurement of virial radii of the clusters which can be provided by X-rays (e.g. Arnaud et al. 2005).
Throughout this paper, we assume an Einstein-de Sitter Universe but our results should be qualitatively similar in a concordance $`\mathrm{\Lambda }`$CDM model. The role of dark energy becomes important at low redshifts ($`0.2`$) which we expect to occur well after the formation of the typical dark matter haloes we consider here. Furthermore, once a particle turns around and collapses, it separates from the background expansion and its subsequent motion should not be affected by the $`\mathrm{\Lambda }`$ term.
The paper is organized as follows. We review the three-dimensional and projected properties of self-similar haloes and caustics in Sect. 2. We derive the lensing signal for a single halo, compare it to the noise level of fiducial observations and estimate the number of haloes required to achieve a significant signal-to-noise ratio in Sect. 3, where we also examine the ability of weak lensing to constrain the velocity dispersion of dark matter particles. We summarize, discuss the prospectives for future works and conclude in Sect. 4.
## 2 Dark matter caustics
### 2.1 Tri-dimensional key equations
In an Einstein-de Sitter Universe a spherical overdensity expands and then turns around to collapse. After collapse and at late times, the fluid motion becomes selfsimilar: its form remains unchanged when its length is re-scaled in terms of the radius of the shell that is currently at the “turn around” and is falling onto the galaxy for the first time. Physically, selfsimilarity arises because gravity is scale-free and because mass shells outside the initial overdensity are also bound and turn around at successively later times. Self-similar solutions give power-law density profiles whose exact scaling properties depend on the central boundary conditions and on whether the fluid is collisionless or collisional (Gott 1975; Gunn 1977; Fillmore & Goldreich 1984; Bertschinger 1985). The density profile obeys a power-law on the scale of the halo which provides an explanation of the flattening of the rotation curves of the galaxies. However, on smaller scales the density profile contains many spikes (i.e. caustics) of infinite density. The position and the time of formation of these caustics are among the many properties that have been studied in the selfsimilar infall model (Bertschinger 1985).
In the presence of a small velocity dispersion the maximum density and thickness of the caustic shells and their density profiles have been evaluated in the framework of a selfsimilar collapse model (MS05).
The global halo density profile, asymptotically reached in this process, is well-approximated by
$$\frac{\rho _{\mathrm{halo}}(\lambda )}{\rho _H}\frac{2.8\lambda ^{9/4}}{(1+\lambda ^{3/4})^2},$$
(1)
where $`\lambda =r/r_{ta}`$ with $`r_{ta}`$ the present turnaround radius of the halo and $`\rho _H=3H^2/8\pi G`$ the background density. The turnaround radius can be easily computed using the virial radius. Within the virial radius, $`r_{\mathrm{vir}}=r_{200}`$, the mean density, $`\rho _{\mathrm{vir}}`$, is, by definition, $`200`$ times the background density. Thus, using the density profile (1), we obtain the following relationship between the turnaround and the virial radii: $`r_{\mathrm{ta}}\mathrm{\hspace{0.17em}4}r_{\mathrm{vir}}`$.
For a perfectly cold dark matter medium, the density profile close to a caustic at $`\lambda _k`$ is (Bertschinger 1985)
$$\frac{\rho _0(\lambda )}{\rho _H}=\frac{G_k}{\sqrt{\lambda _k\lambda }};\sigma =0.$$
(2)
with
$$G_k=\frac{\pi ^2}{4\sqrt{2\lambda _k^{\prime \prime }}}\frac{e^{2\xi _k/3}}{\lambda _k^2},$$
(3)
where the values of the various quantities $`\xi _k`$, $`\lambda _k`$, $`\lambda _k^{\prime \prime }`$ \[and $`\mathrm{\Lambda }_k`$ which will appear in the coming expression (6)\] are given in Table LABEL:table:caustic-parameters \[see MS05 for a detailed description of these parameters\].
When the temperature of particles is not strictly zero and the velocity of particles is distributed according to the distribution function $`f(v)`$, caustic positions are shifted by a small value $`\delta \lambda `$ and the caustic density is modified as (MS05):
$$\rho _\sigma (\lambda )=dv\rho _0[\lambda \delta \lambda (v)]f(v);\sigma 0$$
(4)
In this work, we choose a top-hat velocity distribution function (MS05). Then, the density close to the $`k`$-th caustic in the halo is given by
$$\frac{\rho _\sigma (\lambda )}{\rho _H}=\frac{G_k}{\mathrm{\Delta }_k}\{\begin{array}{cc}\sqrt{\lambda _k^+\lambda }\sqrt{\lambda _k^{}\lambda }\hfill & \text{ for }\lambda <\lambda _k^{},\hfill \\ \sqrt{\lambda _k^+\lambda }\hfill & \text{ for }\lambda _k^{}<\lambda <\lambda _k^+,\hfill \\ 0\hfill & \text{ for }\lambda >\lambda _k^+,\hfill \end{array}$$
(5)
where $`\lambda _k^{}=\lambda _k\mathrm{\Delta }_k`$ and $`\lambda _k^+=\lambda _k+\mathrm{\Delta }_k`$ and the thickness of the $`k`$-th caustic, $`\mathrm{\Delta }_k`$, in nondimensional coordinate is given by
$$r_{\mathrm{ita}}\mathrm{\Delta }_k=\mathrm{\Delta }r_k=\frac{(3\pi )^{2/3}}{4}e^{5\xi _k/9}|\mathrm{\Lambda }_k|t\sigma (t)$$
(6)
and $`r_{\mathrm{ita}}`$ is the initial turnaround radius and $`\sigma (t)`$ is the value of the velocity dispersion of dark matter particles at time $`t`$ which is that at decoupling re-scaled by the expansion factor<sup>3</sup><sup>3</sup>3Hereafter, velocity dispersion is given at the present time, $`z=0`$. For instance, neutralinos have $`\sigma 0.03`$ cm/s and axions have $`\sigma 10^7`$ cm/s.. The standard spherical collapse model yields a relation $`r_{\mathrm{ita}}=2r_{\mathrm{vir}}`$ for a constant overdensity \[see e.g. chapter 5.10 of (Padmanabhan 2002)\], However, the real value of the initial turnaround radius would be lower due to the continuous accretion by the halo. Thus the physical thickness of the caustic, $`\mathrm{\Delta }r_k`$, depends only on its position in the halo and the nature of dark matter.
The maximum density at each caustic position can also be approximated by the expression (see Table LABEL:table:caustic-parameters and also MS05)
$$\rho _{\mathrm{max}}=\frac{2G_ke^{5\xi _k/18}}{(3\pi )^{1/3}\sqrt{\mathrm{\Lambda }_k}}\sqrt{\frac{r_{\mathrm{ta}}}{t\sigma (t)}}\rho _H$$
(7)
Expressions (5), (6) and (7) provide us with a sufficient basis for the evaluation of magnification and shear due to dark matter caustics. Although in the rest of this work we directly integrate expression (5) and never use (7), values for the latter are given in Table LABEL:table:caustic-parameters to demonstrate the density contrast of each caustic with respect to its host halo.
Although the selfsimilar model might seem naive, it provides a good approximation for outer caustics in galaxy-cluster haloes which are not significantly disrupted by merger and substructures and to a good approximation are spherical. Furthermore, fluctuations caused by large scale structure would already be averaged out in the statistical evaluation of the shear.
### 2.2 Projected densities
Since we are concerned with the lensing properties of caustics, we have to calculate the projected density profile for the caustics and for the halo. The Abel integral relates the 3-dimensional density ($`\rho `$) and the 2-dimensional density profiles ($`\mathrm{\Sigma }`$) by
$$\mathrm{\Sigma }(\lambda )=2r_{ta}_\lambda ^{\mathrm{}}\frac{\rho (\lambda ^{})\lambda ^{}\mathrm{d}\lambda ^{}}{\sqrt{\lambda ^2\lambda ^2}}.$$
(8)
We numerically integrate the above expression, for the density profiles (2) and (5). The halo surface mass density can be evaluated analytically if we neglect the $`(1+\lambda ^{3/4})^2`$ term in equation (1), i.e. at small scales $`\lambda 1`$, yielding the approximate projected halo profile
$$\mathrm{\Sigma }_{\mathrm{halo}}(\lambda )r_{ta}7.56\lambda ^{5/4}\rho _H.$$
(9)
Let us define the mean projected density enclosed by the cylinder of radius $`\lambda `$:
$$\overline{\mathrm{\Sigma }}(\lambda )=\frac{2}{\lambda ^2}_0^\lambda \mathrm{\Sigma }(\lambda ^{})\lambda ^{}d\lambda ^{}.$$
(10)
Thus, under the hypothesis of small $`\lambda `$, the mean projected density for the halo component is $`\overline{\mathrm{\Sigma }}_{\mathrm{halo}}(\lambda )=\frac{8}{3}\mathrm{\Sigma }(\lambda )`$.
Figure 1 shows the 3D (left) and 2D (right) density profiles for the halo and the caustics. We consider the case of a perfectly cold dark matter with very peaked caustics and the case in which they are smoothed out by a finite velocity dispersion $`\sigma 60\mathrm{km}\mathrm{s}^1`$. This value is very high for most cold dark matter models but its extremity illustrates well the peak dilution and shows that caustics survive a large thermal softening. When considering the projected density, instead of spikes we rather see stairs that come from the analytic $`1/\sqrt{r_kr}`$ singularity of caustics. Consequently even with singular caustics, the projected density profile is not peaked. The implications for lensing are discussed in the next section.
## 3 Weak lensing
### 3.1 General equations
The fundamental quantity for gravitational lensing is the lens potential $`\psi (\theta )`$ at the angular position $`\theta `$ which is related to the surface mass density $`\mathrm{\Sigma }(\theta )`$ projected into the lens plane through :
$$\psi (\theta )=\frac{4G}{c^2}\frac{D_\mathrm{l}D_\mathrm{s}}{D_{\mathrm{ls}}}\mathrm{d}^2\theta ^{}\mathrm{\Sigma }(\theta ^{})\mathrm{ln}|\theta \theta ^{}|,$$
(11)
where $`D_\mathrm{l}`$, $`D_\mathrm{s}`$ and $`D_{\mathrm{ls}}`$ are angular distances to the lens, to the source and between the lens and the source respectively. The deflection angle $`\alpha =\psi `$ relates a point in the source plane $`\beta `$ to its image(s) in the image plane $`\theta `$ through the lens equation $`\beta =\theta \alpha (\theta )`$. The local relation between $`\beta `$ and $`\theta `$ is the Jacobian matrix $`A_{ij}=\beta _i/\theta _j`$ :
$$A_{ij}=\delta _{ij}\psi _{,ij}=\left(\begin{array}{cc}1\kappa \gamma _1& \gamma _2\\ \gamma _2& 1\kappa +\gamma _1\end{array}\right).$$
(12)
with the convergence $`\kappa (\theta )=\mathrm{\Sigma }(\theta )/\mathrm{\Sigma }_{\mathrm{crit}}`$ directly related to the surface mass density via the critical density
$$\mathrm{\Sigma }_{\mathrm{crit}}=\frac{c^2}{4\pi G}\frac{D_{\mathrm{ls}}}{D_\mathrm{l}D_\mathrm{s}},$$
(13)
and the 2-component shear $`\gamma =\gamma _1+i\gamma _2`$ in complex notation. The convergence satisfies the Poisson equation $`\mathrm{\Delta }\psi =2\kappa `$. In the weak lensing regime ($`\gamma 1`$), an elliptical object in the source plane with complex ellipticity $`e_s`$ is mapped into an elliptical image with a different ellipticity $`e=e_s+\gamma `$. We refer the reader to the reviews of Mellier (1999) and Bartelmann & Schneider (2001) for detailed accounts of weak lensing.
For a circularly symmetric lens, $`\gamma `$ is oriented tangentially to the lens center and its amplitude at radius $`r`$ is $`\gamma (r)=(\overline{\mathrm{\Sigma }}(r)\mathrm{\Sigma }(r))/\mathrm{\Sigma }_{\mathrm{crit}}`$. Since sources are randomly oriented, the tangential component of the observed galaxies is an unbiased estimator of $`\gamma `$. When averaging the estimate of $`\gamma `$ within an aperture of solid angle $`\mathrm{\Omega }`$, containing $`N=\mathrm{\Omega }n`$ galaxies ($`n`$ is the number density of sources), the noise dispersion of $`\gamma `$ is $`\sigma _\gamma =\frac{\sigma _e}{\sqrt{N}}`$ where $`\sigma _e0.3`$ is the intrinsic dispersion of source ellipticities (along one component).
### 3.2 Shear measurement
In order to be consistent with our calculations of Sec. 2.2, we define a pseudo-shear: $`\mathrm{\Gamma }(\lambda )=(\overline{\mathrm{\Sigma }}\mathrm{\Sigma })/\rho _Hr_{ta}`$ and the corresponding noise level $`\mathrm{\Gamma }_N=\mathrm{\Sigma }_{\mathrm{crit}}\sigma _\gamma /\rho _Hr_{ta}`$. For an EdS cosmology, and considering an annulus of inner and outer radii $`\lambda _1`$ and $`\lambda _2`$ respectively, it is straightforward to write $`\mathrm{\Gamma }_N`$ in this useful form:
$$\begin{array}{c}\mathrm{\Gamma }_N(\lambda _1,\lambda _2)=\mathrm{\hspace{0.17em}2.16}\frac{D_{os}}{D_{ls}}\left(\frac{5\mathrm{Mpc}}{r_{ta}}\right)^2(1+z_l)^3\times \hfill \\ \hfill \sqrt{\frac{30\mathrm{arcmin}^2}{n}}\left(\frac{\sigma _e}{0.3}\right)\frac{1}{\sqrt{\lambda _2^2\lambda _1^2}}.\end{array}$$
(14)
This expression seems to indicate that the noise level will be lower for nearby haloes. However, it hides the fact that low redshift haloes require a very wide sky coverage for the outermost caustics ($`r_{\mathrm{ta}}/3`$) to fit the field of view of the observation. Consequently, intermediate redshift haloes ($`z0.20.5`$) are the most interesting targets. In addition, nearby (and thus large angular scale) clusters suffer from noise due to large-scale structure (LSS) fluctuation integrated along the line-of-sight and unrelated to the halo we are considering (Hoekstra 2003). For scales $`15`$ arcmin, the smearing of the shear profile by LSS is minimised.
Figure 2 shows $`\mathrm{\Gamma }`$ as a function of $`\lambda `$ for the same values of the thermal velocity dispersions $`\sigma =0\mathrm{km}\mathrm{s}^1`$ and $`\sigma =60\mathrm{km}\mathrm{s}^1`$. Comparing these curves, one can see that the sawtooth patterns due to caustics survive significantly high temperatures. Next, we consider the noise level for a fiducial halo at redshift $`z_l=0.3`$ and a turnaround radius $`r_{ta}=5\mathrm{Mpc}`$ which is a typical value for clusters (upper green “binned” curve). With a single halo the detection of caustics is impossible. If we are able to stack the signal from a few tens of clusters (upper blue-binned curve), the noise level will be low enough to be sensitive to caustics as a whole. However, we would detect a smooth contribution which is indistinguishable from the halo itself.
Sawtooth patterns cannot be confused with the effect of substructures since the contribution of the latter would be averaged out over the azimuthal angle (most outer caustics have spherical symmetry) and once rescaled, caustics always appear at the same place within haloes. This is not the case for substructures which can appear at any radius inside the host halo.
Consequently, the right way to probe the existence of caustics is to measure the $`\mathrm{\Gamma }`$ signal in excess/default relative to the extreme value of $`\sigma `$, e.g. $`\sigma 300\mathrm{km}\mathrm{s}^1`$ as taken here. For this purpose, one needs $`100`$ stacked clusters. With the corresponding noise level, it would be possible to test the thermal smoothing of caustics and put constraints on $`\sigma `$. However, the sensitivity is poor and only upper limits can be put on $`\sigma `$ with a realistic number of haloes. For instance with $`N=100`$ (resp. 250) clusters, we could achieve a limit $`\sigma <170\mathrm{km}\mathrm{s}^1`$ (resp. $`40\mathrm{km}\mathrm{s}^1`$) at a $`95.4\%`$ confidence level.
When considering galaxies instead of clusters with turnaround radii about $`10`$ times smaller, the number of haloes required to achieve the same detection level is $`10^4`$ times higher. So using a few $`10^5`$ galaxies between $`z0.10.5`$ would yield the same results.
This required level of signal can easily be achieved with a wide and moderately deep survey like the ongoing CFHTLS. Typically, a square-degree field of view will contain a few such cluster-size haloes with $`r_{\mathrm{ta}}5\mathrm{Mpc}`$, a few thousand elliptical galaxies with $`r_{\mathrm{ta}}1\mathrm{Mpc}`$ and tens of thousands of spiral galaxies with $`r_{\mathrm{ta}}500\mathrm{kpc}`$. The total coverage of the CFHTLS wide survey is 170 square degree and will contain a large enough number of clusters/galaxies. Furthermore the wide fields of view are well-suited to measure shear up to the outermost caustic of clusters ($`1\mathrm{Mpc}=5\mathrm{}`$ at $`z=0.3`$ and $`4\mathrm{}`$ at $`z=0.5`$).
Spatial observations provide an improvement on shear measurements: (i) the intrinsic dispersion in source ellipticities is lowered $`\sigma _e0.2`$ and (ii) the density of sources is increased $`n50`$ arcmin<sup>-2</sup>. Hence the total number of haloes required to achieve the same detection level can be lowered by a factor of three. Future wide spatial surveys like SNAP or DUNE will provide the required sky coverage.
### 3.3 Halo-stacking issues
So far, we have considered the observational difficulties encountered by the dispersion in intrinsinc source ellipticities. However, practical complications such as the difficult signal stacking process should also be overcome before a detection of the caustics can be achieved. To do so, high-precision measurements of the location of the center and the turnaround radius of each halo are required. Otherwise, imperfect alignment/rescaling would tend to blur the caustics spikes and reduce the sensitivity. The center of a well-relaxed cluster coincides within one arc-second with the center of the halo. Hence the only scaling factor is the turnaround radius that can be related to the virial radius \[see the sentence following expression (1)\].
Figure 3 shows the blurring of caustics due to an imperfect knowledge of the turnaround radius of each stacked halo. We consider a dispersion around the true value of 3% and 10%. To properly detect caustics, one needs a precise estimate of $`r_{\mathrm{ta}}`$. Provided dark matter is cold enough, $`\sigma `$ a few $`\mathrm{km}\mathrm{s}^1`$, which is a realistic prescription, a reasonable number of stacked haloes, $`200`$, can overcome the loss of a few percents of relative precision caused by the errors in the determination of the turnaround radius.
Consequently the number of haloes necessary for detection needs to be increased in order to achieve the necessary level of signal-to-noise ratio. Moreover, additional external information for the determination of $`r_{\mathrm{ta}}`$ such as X-rays or dynamical observations in addition to weak lensing data are indispensable for a convincing detection of caustics. The present-day state-of-the-art X-ray estimates for the scale radius is $`10\%`$ for nearby clusters ($`z<0.2`$) (Arnaud et al. 2005; Pointecouteau et al. 2005) and it seems that a similar precision can be reached at higher redshifts ($`0.4<z<0.7`$) (Kotov & Vikhlinin 2005). In addition shear, which is used for caustic detection, also provides constraints on the halo density profile. The virial radius of some clusters presenting strong and weak lensing features can be measured with good accuracy ($`\mathrm{\Delta }r_{200}/r_{200}3`$%) (Broadhurst et al. 2005; Gavazzi 2005). Consequently, future large cluster samples with X-rays and lensing data of present-day quality will provide us with the necessary precision to probe dark matter caustics.
## 4 Discussion & conclusion
Although the Liouville theorem claims that singularities will survive, they are likely to develop a complex topology in the course of evolution of structures under self-gravity. It is not clear to what degree the merging processes and substructures will smear out the caustics or complicate their geometries. Here, we have considered the outer caustics of cluster-size haloes which are expected to have suffered far less from mergers and due to their relatively large separations are unlikely to have been washed away by the dispersive-like effect of the substructures. These caustics are the ones that contribute most to the lensing signal for the following reasons. Their amplitude relative to the background smooth halo component is more important as compared to the inner caustics. The inner caustics smear most from thermalisation and imperfect stacking and finally, the noise level increases toward the center of haloes thus giving more weight to the outer caustics. Hence, we have focused on caustics of cluster-size haloes and have argued that they can be reasonably approximated by a selfsimilar spherical accretion model, though the triaxiality of haloes is well established in numerical simulations (Jing & Suto 2002). However, caustic patterns should exist in triaxial matter distributions. The singularities will have the same elliptical symmetry and may be properly stacked from one halo to another by choosing a subsample of apparently circular projected haloes or by using the shear azimuthal variations to constrain the halo ellipticity. We expect this effect to be comparable to the uncertainty in the halo scaling radius (either turnaround or virial).
We have shown that the existence of dark matter caustics could be probed by properly stacking the weak lensing signal of a reasonable number of haloes. The main observational limitation is perhaps the precise estimation of the turnaround radius, $`r_{\mathrm{ta}}`$, of superimposed haloes but we have shown that the loss of a few percent relative accuracy in the determination of $`r_{\mathrm{ta}}`$ (or asphericity) can be compensated for by stacking about $`200`$ haloes.
Although the sensitivity is low, a detection of caustics provides an upper bound for the temperature of dark matter, thus excluding hot dark matter models. The sensitivity is not sufficient to distinguish between various cold dark matter candidates (like axions or neutralinos) since for most of the corresponding velocity dispersions, the shear signal would be similar.
However, a detection of caustics would be a strong argument for the existence of cold dark matter since alternative models like MOND could not explain such density singularities and at most could serve in place of a smooth halo (namely provide an equivalent effective gravitational potential).
Putting constraints on the velocity dispersion of dark matter is a challenging topic in modern cosmology since it offers the possibility to pin-down an actual physical parameter of dark matter. In this paper, we investigated the possibility of such a measurement with the weak lensing effect of dark matter caustics. The implicit observational hypothesis is that we can have a selfsimilar geometrical description of the caustic shell distribution which depends on a single characteristic scaling parameter: i.e. the virial radius.
Wide field surveys such as the ongoing CFHTLS accompanied by X-ray observations can provide the required statistics for a successful detection of caustics. The number of haloes required to be superimposed will be lowered by a further factor of 3 for future space-based experiments like SNAP or DUNE.
Here, we have used the first and most common caustic singularity, that for which the density profile falls with inverse square-root of the distance from the caustic. Caustics of higher-order singularities can also appear in collisionless media and have already been classified (Arnold 1986). It remains a challenging task to generalize our work to haloes with non-vanishing eccentricity where higher-rank caustics would occur and to examine if they can modify the magnification properties of the lensed images and account for anomalous image flux ratios.
###### Acknowledgements.
Special thanks go to Sergei Shandarin for an ongoing collaboration on dark matter caustics. We also thank Francis Bernardeau, Ed Bertschinger, Monique Arnaud, Gary Mamon and Etienne Pointecouteau for many helpful discussions and comments. R.G. is supported at Toulouse from a CNRS postdoc contract #1019 and at Oxford by a grant from the Leverhulme Trust. R.M. is supported by European Gravitational Observatory grant at the school of astronomy, university of Cardiff, UK. |
warning/0506/hep-ex0506055.html | ar5iv | text | # Evidence for 𝜅 Meson Production in 𝐽/𝜓→𝐾̄^∗(892)⁰𝐾⁺𝜋⁻ Process
## Abstract
Based on 58 million BESII $`J/\psi `$ events, the $`\overline{K}^{}(892)^0K^+\pi ^{}`$ channel in $`K^+K^{}\pi ^+\pi ^{}`$ is studied. A clear low mass enhancement in the invariant mass spectrum of $`K^+\pi ^{}`$ is observed. The low mass enhancement does not come from background of other $`J/\psi `$ decay channels, nor from phase space. Two independent partial wave analyses have been performed. Both analyses favor that the low mass enhancement is the $`\kappa `$, an isospinor scalar resonant state. The average mass and width of the $`\kappa `$ in the two analyses are 878 $`\pm `$ 23$`{}_{55}{}^{}{}_{}{}^{+64}`$ MeV/$`c^2`$ and 499 $`\pm `$ 52$`{}_{87}{}^{}{}_{}{}^{+55}`$ MeV/$`c^2`$, respectively, corresponding to a pole at $`(841\pm 30_{73}^{+81})i(309\pm 45_{72}^{+48})`$ MeV/$`c^2`$.
M. Ablikim<sup>1</sup>, J. Z. Bai<sup>1</sup>, Y. Ban<sup>15</sup>, J. G. Bian<sup>1</sup>, X. Cai<sup>1</sup>, H. F. Chen<sup>21</sup>, H. S. Chen<sup>1</sup>, H. X. Chen<sup>1</sup>, J. C. Chen<sup>1</sup>, Jin Chen<sup>1</sup>, Y. B. Chen<sup>1</sup>, S. P. Chi<sup>2</sup>, Y. P. Chu<sup>1</sup>, X. Z. Cui<sup>1</sup>, Y. S. Dai<sup>23</sup>, Z. Y. Deng<sup>1</sup>, Q. F. Dong<sup>19</sup>, S. X. Du<sup>1</sup>, Z. Z. Du<sup>1</sup>, J. Fang<sup>1</sup>, S. S. Fang<sup>2</sup>, C. D. Fu<sup>1</sup>, C. S. Gao<sup>1</sup>, Y. N. Gao<sup>19</sup>, S. D. Gu<sup>1</sup>, Y. T. Gu<sup>4</sup>, Y. N. Guo<sup>1</sup>, Y. Q. Guo<sup>1</sup>, Z. J. Guo<sup>20</sup>, F. A. Harris<sup>20</sup>, K. L. He<sup>1</sup>, M. He<sup>16</sup>, Y. K. Heng<sup>1</sup>, H. M. Hu<sup>1</sup>, T. Hu<sup>1</sup>, G. S. Huang<sup>1</sup><sup>a</sup>, X. P. Huang<sup>1</sup>, X. T. Huang<sup>16</sup>, M. Ishida<sup>10</sup>, S. Ishida<sup>14</sup>, X. B. Ji<sup>1</sup>, X. S. Jiang<sup>1</sup>, J. B. Jiao<sup>16</sup>, D. P. Jin<sup>1</sup>, S. Jin<sup>1</sup>, Yi Jin<sup>1</sup>, T. Komada<sup>14</sup>, S. Kurokawa<sup>8</sup>, Y. F. Lai<sup>1</sup>, G. Li<sup>2</sup>, H. B. Li<sup>1</sup>, H. H. Li<sup>1</sup>, J. Li<sup>1</sup>, R. Y. Li<sup>1</sup>, S. M. Li<sup>1</sup>, W. D. Li<sup>1</sup>, W. G. Li<sup>1</sup>, X. L. Li<sup>9</sup>, X. Q. Li<sup>13</sup>, Y. L. Li<sup>4</sup>, Y. F. Liang<sup>17</sup>, H. B. Liao<sup>6</sup>, C. X. Liu<sup>1</sup>, F. Liu<sup>6</sup>, Fang Liu<sup>21</sup>, H. H. Liu<sup>1</sup>, H. M. Liu<sup>1</sup>, J. Liu<sup>15</sup>, J. B. Liu<sup>1</sup>, J. P. Liu<sup>22</sup>, R. G. Liu<sup>1</sup>, Z. A. Liu<sup>1</sup>, F. Lu<sup>1</sup>, G. R. Lu<sup>5</sup>, H. J. Lu<sup>21</sup>, J. G. Lu<sup>1</sup>, C. L. Luo<sup>12</sup>, F. C. Ma<sup>9</sup>, H. L. Ma<sup>1</sup>, L. L. Ma<sup>1</sup>, Q. M. Ma<sup>1</sup>, X. B. Ma<sup>5</sup>, Z. P. Mao<sup>1</sup>, T. Matsuda<sup>11</sup>, X. H. Mo<sup>1</sup>, J. Nie<sup>1</sup>, H. P. Peng<sup>21</sup>, N. D. Qi<sup>1</sup>, H. Qin<sup>12</sup>, J. F. Qiu<sup>1</sup>, Z. Y. Ren<sup>1</sup>, G. Rong<sup>1</sup>, L. Y. Shan<sup>1</sup>, L. Shang<sup>1</sup>, D. L. Shen<sup>1</sup>, X. Y. Shen<sup>1</sup>, H. Y. Sheng<sup>1</sup>, F. Shi<sup>1</sup>, X. Shi<sup>15</sup><sup>b</sup>, H. S. Sun<sup>1</sup>, J. F. Sun<sup>1</sup>, S. S. Sun<sup>1</sup>, Y. Z. Sun<sup>1</sup>, Z. J. Sun<sup>1</sup>, K. Takamatsu<sup>8</sup>, Z. Q. Tan<sup>4</sup>, X. Tang<sup>1</sup>, Y. R. Tian<sup>19</sup>, G. L. Tong<sup>1</sup>, T. Tsuru<sup>8</sup>, K. Ukai<sup>8</sup>, D. Y. Wang<sup>1</sup>, L. Wang<sup>1</sup>, L. S. Wang<sup>1</sup>, M. Wang<sup>1</sup>, P. Wang<sup>1</sup>, P. L. Wang<sup>1</sup>, W. F. Wang<sup>1</sup><sup>c</sup>, Y. F. Wang<sup>1</sup>, Z. Wang<sup>1</sup>, Z. Y. Wang<sup>1</sup>, Zhe Wang<sup>1</sup>, Zheng Wang<sup>2</sup>, C. L. Wei<sup>1</sup>, D. H. Wei<sup>1</sup>, N. Wu<sup>1</sup>, X. M. Xia<sup>1</sup>, X. X. Xie<sup>1</sup>, B. Xin<sup>9</sup><sup>a</sup>, G. F. Xu<sup>1</sup>, Y. Xu<sup>13</sup>, K. Yamada<sup>14</sup>, I. Yamauchi<sup>18</sup>, M. L. Yan<sup>21</sup>, F. Yang<sup>13</sup>, H. X. Yang<sup>1</sup>, J. Yang<sup>21</sup>, Y. X. Yang<sup>3</sup>, M. H. Ye<sup>2</sup>, Y. X. Ye<sup>21</sup>, Z. Y. Yi<sup>1</sup>, G. W. Yu<sup>1</sup>, J. M. Yuan<sup>1</sup>, Y. Yuan<sup>1</sup>, S. L. Zang<sup>1</sup>, Y. Zeng<sup>7</sup>, Yu Zeng<sup>1</sup>, B. X. Zhang<sup>1</sup>, B. Y. Zhang<sup>1</sup>, C. C. Zhang<sup>1</sup>, D. H. Zhang<sup>1</sup>, H. Y. Zhang<sup>1</sup>, J. W. Zhang<sup>1</sup>, J. Y. Zhang<sup>1</sup>, Q. J. Zhang<sup>1</sup>, X. M. Zhang<sup>1</sup>, X. Y. Zhang<sup>16</sup>, Yiyun Zhang<sup>17</sup>, Z. P. Zhang<sup>21</sup>, Z. Q. Zhang<sup>5</sup>, D. X. Zhao<sup>1</sup>, J. W. Zhao<sup>1</sup>, M. G. Zhao<sup>13</sup>, P. P. Zhao<sup>1</sup>, W. R. Zhao<sup>1</sup>, Z. G. Zhao<sup>1</sup><sup>d</sup>, H. Q. Zheng<sup>15</sup>, J. P. Zheng<sup>1</sup>, Z. P. Zheng<sup>1</sup>, L. Zhou<sup>1</sup>, N. F. Zhou<sup>1</sup>, K. J. Zhu<sup>1</sup>, Q. M. Zhu<sup>1</sup>, Y. C. Zhu<sup>1</sup>, Y. S. Zhu<sup>1</sup>, Yingchun Zhu<sup>1</sup><sup>e</sup>, Z. A. Zhu<sup>1</sup>, B. A. Zhuang<sup>1</sup>, X. A. Zhuang<sup>1</sup>.
(BES Collaboration)
<sup>1</sup> Institute of High Energy Physics, Beijing 100049, People’s Republic of China
<sup>2</sup> China Center for Advanced Science and Technology(CCAST), Beijing 100080, People’s Republic of China
<sup>3</sup> Guangxi Normal University, Guilin 541004, People’s Republic of China
<sup>4</sup> Guangxi University, Nanning 530004, People’s Republic of China
<sup>5</sup> Henan Normal University, Xinxiang 453002, People’s Republic of China
<sup>6</sup> Huazhong Normal University, Wuhan 430079, People’s Republic of China
<sup>7</sup> Hunan University, Changsha 410082, People’s Republic of China
<sup>8</sup> KEK, High Energy Accelerator Research Organization, Ibaraki 305-0801, Japan
<sup>9</sup> Liaoning University, Shenyang 110036, People’s Republic of China
<sup>10</sup> Meisei University, Tokyo 191-8506, Japan
<sup>11</sup> Miyazaki University, Miyazaki 889-2192, Japan
<sup>12</sup> Nanjing Normal University, Nanjing 210097, People’s Republic of China
<sup>13</sup> Nankai University, Tianjin 300071, People’s Republic of China
<sup>14</sup> Nihon University, Chiba 274-8501, Japan
<sup>15</sup> Peking University, Beijing 100871, People’s Republic of China
<sup>16</sup> Shandong University, Jinan 250100, People’s Republic of China
<sup>17</sup> Sichuan University, Chengdu 610064, People’s Republic of China
<sup>18</sup> Tokyo Metropolitan College of Technology, Tokyo 140-0011, Japan
<sup>19</sup> Tsinghua University, Beijing 100084, People’s Republic of China
<sup>20</sup> University of Hawaii, Honolulu, HI 96822, USA
<sup>21</sup> University of Science and Technology of China, Hefei 230026, People’s Republic of China
<sup>22</sup> Wuhan University, Wuhan 430072, People’s Republic of China
<sup>23</sup> Zhejiang University, Hangzhou 310028, People’s Republic of China
<sup>a</sup> Current address: Purdue University, West Lafayette, IN 47907, USA
<sup>b</sup> Current address: Cornell University, Ithaca, NY 14853, USA
<sup>c</sup> Current address: Laboratoire de l’Accélératear Linéaire, Orsay, F-91898, France
<sup>d</sup> Current address: University of Michigan, Ann Arbor, MI 48109, USA
<sup>e</sup> Current address: DESY, D-22607, Hamburg, Germany
In the field of hadron spectroscopy, whether the low mass iso-scalar scalar meson, the $`\sigma `$, exists or not had been an important but controversial problem for many years. Recently, its evidence has been reported - not only in $`\pi \pi `$ scattering, but also in various production processes. The $`\sigma `$ meson with a mass around 600 MeV/$`c^2`$ and a broad width around 500 MeV/$`c^2`$ is, now, widely accepted .
The evidence for the $`\sigma `$ meson suggests the possibility of a $`\sigma `$ nonet, {$`\sigma (600)`$, $`\kappa (900)`$, $`f_0(980)`$, $`a_0(980)`$ }, either in an extended linear sigma model realizing chiral symmetry or in a unitarized meson model . The $`\kappa `$ has been observed in analyses on $`K\pi `$ scattering phase shifts by several groups using a unitarized meson method , an interfering amplitude method considering a repulsive background suggested by chiral symmetry, and a nonperturbative method - with an effective chiral Lagrangian. The observed mass and width values are scattered in the ranges from 700 to 900 MeV/$`c^2`$ and 550 to 650 MeV/$`c^2`$, respectively, depending on the model used. Recently, a rather wider width around 800 MeV/$`c^2`$ was reported for the $`\kappa `$ in the analysis of $`K\pi `$ scattering phase shifts with a T-matrix method including a prescription for zero suppression. Also recently, an analysis of LASS data on $`\pi K`$ scattering phase shifts using a unitarization method combined with chiral symmetry has found the $`\kappa `$ with a slightly lighter pole mass . However, some authors have found no evidence for the $`\kappa `$ -. A criticism has been presented for the existence of the $`\kappa `$ based on unitarity and universality arguments, similar to the case of the $`\sigma `$ .
Here, it is to be noted that in $`\pi \pi `$ and $`K\pi `$ scattering, effects due to $`\sigma `$ and $`\kappa `$ production are, as a result of chiral symmetry, largely cancelled by those due to non-resonant background scattering, while the cancellation mechanism does not necessarily work in the production process . Therefore, it is more suitable for the investigation of $`\sigma /\kappa `$ mesons to use the $`\pi \pi /K\pi `$ production process, which is parameterized independently of the scattering process .
Evidence for the $`\kappa `$ has been reported, recently, in the production process by the E791 experiment at Fermilab in the analysis of $`D^+K^{}\pi ^+\pi ^{}`$ . Preliminary $`\kappa `$ results have been reported in the analysis of the $`K\pi `$ system produced in $`J/\psi \overline{K}^{}(892)^0K^+\pi ^{}`$ with BESI data. The FOCUS experiment has presented evidence for the existence of a coherent $`K\pi `$ S-wave contribution to $`D^+K^{}\pi ^+\mu ^+\nu `$. CLEO has seen no evidence for the $`\kappa `$ in $`D^0K^{}\pi ^+\pi ^0`$. Preliminary results on the $`\kappa `$ have also been reported in analyses of BESII data -.
Here we report analyses of $`\overline{K}^{}(892)^0K^+\pi ^{}`$ in $`J/\psi K^+K^{}\pi ^+\pi ^{}`$ to study the $`\kappa `$. Partial wave analyses (PWA analyses) have been performed in $`J/\psi \overline{K}^{}(892)^0K^+\pi ^{}`$ using 58 million $`J/\psi `$ decays obtained with BESII at the BEPC (Beijing Electron Positron Collider) storage ring. The BESII detector is described in detail elsewhere .
In the event selection, four charged tracks with zero net charge are required for each event. Every charged track should have a good helix fit in the Main Drift Chamber (MDC), and the polar angle $`\theta `$ of each track in the MDC must satisfy $`|\mathrm{cos}\theta |<0.8`$. The event must originate near the collision point; tracks must satisfy $`\sqrt{x^2+y^2}2`$ cm, $`z20`$ cm, where $`x`$, $`y`$ and $`z`$ are space coordinates of the point of closest approach of tracks to the beam axis. Particle identification is performed using combined TOF and dE/dx information, and two identified kaons and two identified pions are required.
For a neutral track, it is required that it should have hits in the Barrel Shower Counter (BSC), the number of layers hit should be greater than one, the shower starts before layer 6, and its deposited energy is more than 50 MeV. The angle between the photon emission direction and the shower development direction of the track in the BSC should be less than $`30^{}`$. An event associated with a neutral track(s) is rejected.
Surviving events are fitted kinematically (4C kinematic fits) under the hypotheses $`J/\psi K^+K^{}\pi ^+\pi ^{}`$, $`\pi ^+\pi ^{}\pi ^+\pi ^{}`$, $`K^+K^{}K^+K^{}`$, $`K^+\pi ^{}\pi ^+\pi ^{}`$, and $`\pi ^+K^{}\pi ^+\pi ^{}`$. It is required that $`\chi _{4C}^2(K^+K^{}\pi ^+\pi ^{})`$ be less than 40. The sum of $`\chi ^2`$ of the 4C kinematic fit, TOF, and dE/dx for $`K^+K^{}\pi ^+\pi ^{}`$ is required to be less than those for the $`\pi ^+\pi ^{}\pi ^+\pi ^{}`$, $`K^+K^{}K^+K^{}`$, $`K^+\pi ^{}\pi ^+\pi ^{}`$, and $`\pi ^+K^{}\pi ^+\pi ^{}`$ hypotheses.
In order to remove background from $`J/\psi \varphi \pi ^+\pi ^{}`$, events with $`|M_{K^+K^{}}1.02|<0.02`$ GeV/$`c^2`$ are vetoed, and to remove $`J/\psi \overline{K}^{}(892)^0K_S^0`$ background, events with $`|M_{\pi ^+\pi ^{}}0.497|<0.04`$ GeV/$`c^2`$ and $`R_{xy}>0.8`$ cm are vetoed, where $`M_{K^+K^{}}`$ and $`M_{\pi ^+\pi ^{}}`$ are the invariant masses of $`K^+K^{}`$ and $`\pi ^+\pi ^{}`$, and $`R_{xy}`$ is the decay length of $`K_S`$ transverse to the beam axis.
Fig. 1a shows the scatter plot of $`M_{K^+\pi ^{}}`$ versus $`M_{K^{}\pi ^+}`$ after all above requirements. Two bands of $`\overline{K}^{}(892)^0`$ and $`K^{}(892)^0`$ corresponding to $`J/\psi \overline{K}^{}(892)^0K^+\pi ^{}`$ and $`J/\psi K^{}(892)^0K^{}\pi ^+`$, respectively, are clearly seen. There is an accumulation of events in the region where the two bands cross, which is not seen in the scatter plot of Monte Carlo $`J/\psi anything`$ events (shown in Fig. 1b) obtained with the Lund-charm generator , which is described below. This accumulation comes from $`J/\psi `$ decaying to $`\overline{K}^{}(892)^0`$ (or $`K^{}(892)^0`$) against a low mass enhancement.
The invariant mass spectrum of $`K^{}\pi ^+`$ is shown in Fig. 2a, where the $`\overline{K}^{}(892)^0`$ peak is clearly seen. The requirement 0.812 GeV/$`c^2`$ $`<M_{K^{}\pi ^+}<`$ 0.972 GeV/$`c^2`$ is imposed on the $`J/\psi K^+K^{}\pi ^+\pi ^{}`$ sample to select $`J/\psi \overline{K}^{}(892)^0K^+\pi ^{}`$ events, the total number of which is 24674. The PWA analyses are applied to the selected $`J/\psi \overline{K}^{}(892)^0K^+\pi ^{}`$ sample.
After the $`\overline{K}^{}(892)^0`$ mass requirement, the $`K^+\pi ^{}`$ invariant mass distribution, shown as the solid histogram in Fig. 3a, has a clear $`K^{}(892)^0`$ peak, a peak around 1430 MeV/$`c^2`$, and a broad enhancement in the low mass region. The $`K^{}\pi ^+`$ mass distribution of the charge conjugate channel $`J/\psi K^{}(892)^0K^{}\pi ^+`$, denoted as crosses in Fig. 3a, shows the same structures.
Fig. 3c is the Dalitz plot of $`J/\psi \overline{K}^{}(892)^0K^+\pi ^{}`$, and the insertion is that of its charge conjugate channel. In the Dalitz plot, the two diagonal bands correspond to the low mass enhancement and the peak around 1430 MeV/$`c^2`$ in the $`K^+\pi ^{}`$ invariant mass spectrum, and the horizontal band corresponds to $`J/\psi K_1(1400)K`$ and $`K_1(1270)K`$ with $`K_1(1400)`$ and $`K_1(1270)`$ decaying to $`\overline{K}^{}(892)^0\pi ^{}`$. The clear top diagonal band in the Dalitz plot indicates that the low mass enhancement observed in this decay does not come from phase space since phase space events would be uniformly distributed in the Dalitz plot.
Though this low mass enhancement overlaps with the narrow $`K^{}(892)^0`$, it can be clearly seen due to its broad structure. The spectrum of $`K^{}\pi ^+`$ mass recoiling against the low mass enhancement ($`M_{K^+\pi ^{}}<0.8`$ GeV/$`c^2`$) is shown in Fig. 2b, where a clear $`\overline{K}^{}(892)^0`$ peak can be seen. This means that the low mass enhancement is produced dominantly through $`J/\psi `$ decays associated with the $`\overline{K}^{}(892)^0`$.
$`K^{}(892)^0`$ signals are clearly recognized in the $`K^+\pi ^{}`$ spectrum against $`\overline{K}^{}(892)^0`$ side-band events (the dark shaded histogram) in Fig. 3a. The side-bands events are taken from the $`K^{}\pi ^+`$ mass ranges directly neighboring to $`\overline{K}^{}(892)^0`$ with 80 MeV/$`c^2`$ widths. After side-band subtraction, the $`K^{}(892)^0`$ peak is suppressed appreciably in the invariant mass spectrum of $`K^+\pi ^{}`$, as is shown in the same figure (the dotted histogram). This means that the $`K^{}(892)^0`$ peak mostly comes from background processes. The main part of the broad low mass enhancement remains after side-band subtraction, which indicates that the broad low mass structure does not come from $`J/\psi `$ decay processes which do not contain $`\overline{K}^{}(892)^0`$ in the final states.
The main background channels for $`J/\psi \overline{K}^{}(892)^0K^+\pi ^{}`$ are studied through Monte Carlo simulation. More than 20 $`J/\psi `$ decay channels, including $`J/\psi \gamma K^+K^{}\pi ^+\pi ^{}`$, $`\gamma \pi ^+\pi ^{}\pi ^+\pi ^{}`$, $`K^{}(892)^\pm K^{}`$, $`K^{}(892)^0K_S`$, and $`\overline{K}^{}(892)^0K_S`$ decays are generated using uniform phase space generators, and no peak is produced in the $`K^+\pi ^{}`$ mass spectrum. This means that the low mass broad structure in the $`K^+\pi ^{}`$ mass spectrum does not come from these background channels. From the Monte Carlo simulation, we also see that the background level in $`J/\psi \overline{K}^{}(892)^0K^+\pi ^{}`$ is about 1/7 of that in $`J/\psi K^+K^{}\pi ^+\pi ^{}`$. Therefore, we think $`J/\psi \overline{K}^{}(892)^0K^+\pi ^{}`$ is a good place to study the $`\kappa `$. The background level ranges from 10 to 15 % in this analysis.
We also study 58 million Monte Carlo $`J/\psi anything`$ events which are generated using the Lund-charm model . The generator is developed for simulating $`J/\psi `$ inclusive decay. In the models, charmonium decays into hadrons via quarks and gluons are simulated. Quarks and gluons shower development and their hardronization are handled by the LUND string model. The model reproduces the main properties of hadronic events from $`J/\psi `$ inclusive decay. The process, $`J/\psi K^{}(892)\kappa `$ is not included in the generator. Using the same selection criteria on this Monte Carlo sample as those for data, the scatter plot of $`M_{K^+\pi ^{}}`$ versus $`M_{K^{}\pi ^+}`$ (Fig. 1b), shows no accumulation of events around the region where the $`K^{}(892)^0`$ and $`\overline{K}^{}(892)^0`$ bands cross. There is no corresponding broad low mass enhancement in the invariant mass spectrum of $`K^+\pi ^{}`$ (or $`K^{}\pi ^+`$) recoiling against $`\overline{K}^{}(892)^0`$ (or $`K^{}(892)^0`$ for the charge conjugate channel) shown as the shaded area in Fig. 3b (or in the insertion of Fig. 3b), and there is no diagonal band of the low mass enhancement in the Dalitz plot of $`J/\psi \overline{K}^{}(892)^0K^+\pi ^{}`$ (or $`J/\psi K^{}(892)^0K^{}\pi ^+`$). Because the generator of this Monte Carlo simulation does not contain the process $`J/\psi K^{}(892)\kappa `$, the Monte Carlo shows different structures at lower $`K\pi `$ mass region from those of data. This difference just means that the low mass enhancement is not due to backgrounds coming from other $`J/\psi `$ decay channels.
In the PWA of $`J/\psi \overline{K}^{}(892)^0K^+\pi ^{}`$, the $`\overline{K}^{}(892)^0`$ is treated as a stable particle. After integrating over the $`K^{}\pi ^+`$ angular information, no interference effect is expected between the charge conjugate processes, $`J/\psi \overline{K}^{}(892)^0K^+\pi ^{}`$ and $`K^{}(892)^0K^{}\pi ^+`$, if the solid angle coverage of the detector is $`4\pi `$. We examined the interference effect taking into account the detector acceptance and the width of the $`\overline{K}^{}(892)^0`$ using Monte Carlo simulation and find that the interference effect in the cross region of the two $`K^{}(892)^0`$ bands in the scatter plot is negligibly small, and that the interference between $`\overline{K}^{}(892)^0K^+\pi ^{}`$ and $`\rho K\overline{K}`$ (or $`K_J^{}(1430)K\pi `$) is also negligible. We ignore these interferences in the PWA analyses.
We consider the $`K^{}(892)^0K\pi `$ channel in $`J/\psi `$ decays, with well defined quantum states recoiling against the $`K^{}(892)^0`$, is suitable to study the $`\kappa `$ in the $`K\pi `$ system. The $`K^{}(892)^0K\pi `$ system simplifies the PWA, and backgrounds can be treated easily. To simplify the analysis, one of the charge conjugate states, $`\overline{K}^{}(892)^0K^+\pi ^{}`$, is used, and it has enough events for this analysis. As mentioned above, the analysis of this channel is not affected by the charge conjugate state, $`K^{}(892)^0K^{}\pi ^+`$, since the interference between them is found to be negligibly small.
Four decay mechanisms considered in our analyses are shown in Fig. 4. These are $`J/\psi K^{}R_{K\pi }`$ (Fig. 4a), $`J/\psi KR_{K^{}\pi }`$ (Fig. 4b), $`J/\psi \pi R_{K^{}K}`$ (Fig. 4c) and $`J/\psi K^{}K\pi `$ (Fig. 4d), where $`R_{K\pi }`$, $`R_{K^{}\pi }`$ and $`R_{K^{}K}`$ are intermediate resonant states decaying into $`K\pi (\kappa ,K_0^{}(1430),K_2^{}(1430),K_2^{}(1920))`$, $`K^{}\pi (K_1(1270),K_1(1400))`$ and $`K^{}K(b_1(1235))`$, respectively. (Here, $`K^{}`$ denotes $`K^{}(892)^0`$.) The independent PWA analyses by Method A and Method B have been performed on the $`\overline{K}^{}(892)^0K^+\pi ^{}`$ channel. The same data set is used in both analyses.
Method A is based on the covariant helicity amplitude analysis . The maximum likelihood method is utilized in the fit. The fit is performed by using Breit-Wigner parameterizations with an $`s`$-dependent width, $`\rho (s)`$, (See eq. (1) below.) for the low mass enhancement and with constant widths for the other intermediate states. Fits are also performed using two other parameterizations of constant width and of a width for the $`\kappa `$ using a unitarization approach with chiral symmetry as follows;
$`BW`$ $``$ $`1/(M^2si\sqrt{s}\mathrm{\Gamma }_\kappa (s)),\mathrm{\Gamma }_\kappa (s)=g_\kappa ^2k_K/8\pi s,`$ (1)
$`BW`$ $``$ $`1/(M^2siM\mathrm{\Gamma }_{const}),`$ (2)
and
$`BW`$ $``$ $`1/(M^2si\sqrt{s}\mathrm{\Gamma }_\kappa (s)),\mathrm{\Gamma }_\kappa (s)=\alpha _\kappa k_K.`$ (3)
Eq. (1) is for Method A, eq. (2) is for Method A-1, and eq. (3) is for Method A-2.
The broad low mass enhancement is fit by a $`0^+`$ resonance. The possibility that its spin-parity is $`1^{}`$, $`2^+`$, $`\mathrm{}`$ is excluded by at least 10$`\sigma `$. If the kappa is removed from the fit, the log-likelihood becomes worse by 294. So, its statistical significance is above 20$`\sigma `$. This iso-spinor scalar resonance is considered to be the $`\kappa `$ particle. The above parameterizations are tried for the $`\kappa `$. Though these parameterizations have different behavior, the quality of the fits given by them is almost the same. This is because there are many resonances with interferences between them, and changes in one can be compensated by changes in the others while keeping the total contribution unchanged. Though mass and width parameters given by these parameterizations are somewhat different, the shapes of the $`\kappa `$ obtained by these parameterizations are similar, and the pole positions are close to each other.
The biggest peak at about 1430 MeV/$`c^2`$ in the $`K^+\pi ^{}`$ spectrum is relatively complex, containing $`0^+`$, $`1^{}`$, and $`2^+`$ components. The $`0^+`$ and $`2^+`$ components are identified to be $`K_0^{}(1430)`$ and $`K_2^{}(1430)`$, respectively, and their masses and widths determined from the fit are consistent with PDG values . A $`1^{}`$ is included in the fit in this region. Changes by removing it are considered in the systematic uncertainties. In the higher $`K^+\pi ^{}`$ mass region, a broad resonance is needed. It is found that different treatments of it have little influence on the mass and width of the $`\kappa `$.
$`K_1(1270)`$ and $`K_1(1400)`$ are used to fit the enhancement near threshold of the $`\overline{K}^{}(892)^0`$$`\pi ^{}`$ spectrum, and $`b_1(1235)`$ is used to improve the fit in the $`\overline{K}^{}(892)^0K^+`$ spectrum. Because the mass of $`b_1(1235)`$ is below $`\overline{K}^{}(892)^0K`$ threshold, only the tail of $`b_1(1235)`$ affects the fit of the $`\overline{K}^{}(892)^0K`$ spectrum. Changes caused by removing the $`b_1(1235)`$ are included in the systematic uncertainties.
In the PWA, the backgrounds from the charge conjugate channel $`J/\psi K^{}(892)^0K^{}\pi ^+`$ and from $`J/\psi \overline{K}^{}(892)^0K_S`$ and $`K_S\pi ^+\pi ^{}`$ where $`\pi ^+`$ is misidentified with $`K^+`$ are fitted by non-interfering amplitudes, and the background from other $`J/\psi `$ decay channels, including $`J/\psi \rho K\overline{K}`$ are fitted by non-interfering phase space.
The uncertainty of the background level is considered, and its influence on the mass and width of the $`\kappa `$ is estimated and is put into the systematic errors.
Since there exist broad $`K^+\pi ^{}`$ resonances around 1430 MeV and those broad resonances and other background processes $`K^{}(892),K_S`$, and PS contribute in the $`\kappa `$ region , and because of the complicated interferences between $`\kappa `$ and other resonances, different solutions are obtained in the fits. These solutions give almost the same fit quality and almost the same pole position of $`\kappa `$. Uncertainties coming from the multi-solutions on the $`\kappa `$ parameters are included in the systematic errors.
The fit obtained in the analysis is shown by the solid histogram in Fig. 5a for the $`K^+\pi ^{}`$ invariant mass spectrum. The data are shown by crosses with error bars. The contribution of the $`\kappa `$ is shown by the dark shaded histogram. The fit of the $`\overline{K}^{}(892)^0\pi ^{}`$ invariant mass spectrum is shown by the solid histogram in Fig. 5c. The fit for the angular distribution of the whole $`K^+\pi ^{}`$ mass region is shown by the solid histogram in Fig. 5e, and that for the angular distribution of the $`K^+\pi ^{}`$ mass below 1.0 GeV/$`c^2`$ is shown by the histogram in Fig. 5g.
The values for Breit-Wigner parameters of mass and width and those for the pole position for the $`\kappa `$ obtained in Method A are given in Table 1. Systematic uncertainties in the mass and width of the $`\kappa `$ include the changes from 1$`\sigma `$ variations of the masses and widths of the other resonances, different treatments of background, and uncertainties from the fitting.
Mass and width parameters of intermediate resonances and of those background processes in the fit of Method A are tabulated in Table 2. Errors shown in the table are only statistical.
Method B uses the VMW (Variant Mass and Width) method, a covariant field theoretical approach consistent with generalized unitarity . In this method, the total amplitude is expressed as a coherent sum of respective amplitudes, corresponding to the relevant processes of strong interactions among all color-singlet hadrons. As the bases of the $`S`$-matrix for the strong interaction, a residual interaction of QCD, all unstable (or resonant) as well as stable hadrons which are to be color-singlet bound states of quarks, anti-quarks and gluons are to be included. The propagator of a resonant particle is given by the conventional Feynman propagator with substitution of i$`ϵ`$ by $`i\sqrt{s}\mathrm{\Gamma }(s)`$.
For the scalar $`K\pi `$-resonant particles, $`R_{K\pi }`$’s are the $`\kappa `$ and $`K_0^{}(1430)`$. The Lagrangian of strong interaction, $`_\mathrm{S}`$, describing the process in Fig. 4a is taken to be the most simple form. This form and the corresponding decay amplitude $`_S`$ are given by
$`_S`$ $`=`$ $`{\displaystyle \underset{R=\kappa ,K_0^{}}{}}(\xi _R\psi _\mu K_\mu ^{}R+g_RRK\pi ),`$
$`_S`$ $`=`$ $`S_{h_\psi h_K^{}}{\displaystyle \underset{R=\kappa ,K_0^{}}{}}r_Re^{i\theta _R}\mathrm{\Delta }_R(s_{K\pi }),`$
$`\mathrm{\Delta }_R(s_{K\pi })`$ $`=`$ $`{\displaystyle \frac{m_R\mathrm{\Gamma }_R}{m_R^2s_{K\pi }i\sqrt{s}_{K\pi }\mathrm{\Gamma }_R(s_{K\pi })}},`$ (4)
where $`\mathrm{\Delta }_R(s_{K\pi })`$ is the Breit-Wigner formula with $`\mathrm{\Gamma }_R(s_{K\pi })=pg_R^2/(8\pi s_{K\pi })`$, describing the decay of $`R=\kappa `$ and $`K_0^{}(1430)`$, and $`S_{h_\psi h_K^{}}`$ is the helicity amplitude $`S_{h_\psi h_K^{}}ϵ_\mu ^{(h_\psi )}\stackrel{~}{ϵ}_\mu ^{(h_K^{})}`$. $`S_{h_\psi h_K^{}}r_Re^{i\theta _R}`$ describes the $`S`$-matrix element $`{}_{out}{}^{}RK^{}|J/\psi _{in}^{}`$, where $`e^{i\theta _R}`$ parametrizes the rescattering phase of $`{}_{out}{}^{}RK^{}|`$. This form of $`_S`$ is consistent with the generalized unitarity of the $`S`$-matrix.
The decay amplitudes through the tensor $`R_{K\pi }`$, $`R_{K^{}\pi }`$, and $`R_{K^{}K}`$, denoted as $`_D`$, $`_{K_1}`$, and $`_{b_1}`$, respectively, are obtained in a similar manner. The direct $`K\pi `$ production amplitude is taken to be $`_{\mathrm{dir}}=S_{h_\psi h_K^{}}r_{K\pi }e^{i\theta _{K\pi }}`$, which is from $`_{\mathrm{dir}}\xi _{K\pi }\psi _\mu K_\mu ^{}K\pi `$. The total amplitude $``$ is given by the sum of all these amplitudes,
$``$ $`=`$ $`_S+_D+_{K_1}+_{b_1}+_{\mathrm{dir}}.`$ (5)
We also consider the background processes coming from $`J^P=1^{}`$ $`K^{}(892)`$ and $`K^{}(1410)`$ (decaying into $`K^+\pi ^{}`$), from $`K_S`$, and from phase space $`K^{}(892)K\pi `$, which are described by the amplitudes incoherent with the above $``$. The details are described elsewhere . The treatments of resonances and background processes in this method are the same as those in Method A, except for the $`K^{}(1410)\overline{K}^{}(892)`$ and $`K^{}(892)K\pi `$ processes, as explained below.
The least $`\chi ^2`$ method is used for the fitting of the mass distribution of $`K^+\pi ^{}`$, that of $`\overline{K}^{}(892)^0\pi ^{}`$, and the $`K^+`$ angular distribution in the $`K^+\pi ^{}`$ system. The results obtained in this analysis are shown by the solid histograms for the $`K^+\pi ^{}`$ invariant mass spectrum in Fig. 5b and for the $`\overline{K}^{}(892)^0\pi ^{}`$ mass spectrum in Fig. 5d. The data for the analysis are shown by crosses with error bars. The contributions of the $`\kappa `$ are shown by the dark shaded histograms superimposed on Figs. 5b and 5d. The results for the $`K^+`$ angular distributions in the whole $`K^+\pi ^{}`$ mass region and below 1.0 GeV/$`c^2`$ are shown by the solid histograms in Figs. 5f and 5h, respectively, and that for the mass region above 1.0 GeV/$`c^2`$ is inserted in Fig. 5f. The contributions of the $`\kappa `$ are also shown by the dark shaded histograms in the figures.
The values for Breit-Wigner parameters of mass and width and pole position for the $`\kappa `$ obtained in the analysis of Method B are given in Table 1. Uncertainties are estimated on the same items as in Method A. Method B takes the direct $`\overline{K}^{}(892)^0K^+\pi ^{}`$ process to be coherent and phase space background contribution to be incoherent. The contribution of the latter is estimated using the results obtained in Method A. The uncertainties of it are included in the systematic errors of the $`\kappa `$ parameters. The $`\overline{K}^{}(892)^0K^{}(1410)`$ amplitude is taken as an incoherent background process. The $`K^{}(892)^0`$ events are considered to be associated with $`K^{}\pi ^+`$ and/or $`\overline{\kappa }`$ which are in the $`\overline{K}^{}(892)^0`$ region. No interference is expected between charge conjugate states. A coherent $`\overline{K}^{}(892)^0K^{}(1410)`$ amplitude is also examined, and the difference is also included in the uncertainty for the $`\kappa `$ parameters. Uncertainties for the $`\kappa `$ parameters contain also the change from 1$`\sigma `$ variations of the masses and widths of the other resonances and uncertainties of the fit. Mass and width parameters of intermediate resonances and of background processes in the fit by Method B are also tabulated in Table 2. Errors shown in the table are only statistical.
The $`\chi ^2/`$d.o.f value is 1.10. We examined also the fit without the $`\kappa `$ resonance and obtained the value to be 2.83. In the latter, the parameters for the $`K_S`$ background are fixed.
The results obtained in the two analyses reproduce the data well and are in good agreement with each other. Both fits favor strongly that the low mass enhancement of the $`K^+\pi ^{}`$ system is a resonance. The scalar resonance is considered to be the $`\kappa `$ which is necessary in both fits. The average values for Breit-Wigner parameters of masses and widths for the $`\kappa `$ (given in the third row in Table 1) are obtained from Methods A and B,
$$M_\kappa =878\pm 23_{55}^{+64}\text{MeV}/c^2,\mathrm{\Gamma }_\kappa =499\pm 52_{87}^{+55}\text{MeV}/c^2,$$
where the first term errors are statistical ones. The second ones show total uncertainties, taking the largest values between systematic uncertainties of Methods A and B. The average values are in good agreement with those obtained in the analysis of $`K\pi `$ scattering phase shifts . The $`\kappa `$ parameters obtained are also consistent with those obtained in the analysis of $`D^+K^{}\pi ^+\pi ^+`$ in the E791 experiment with $`M_\kappa =797\pm 19\pm 43`$ MeV/$`c^2`$ and $`\mathrm{\Gamma }_\kappa =410\pm 43\pm 87`$ MeV/$`c^2`$. The relative contribution for the kappa normalized for the total event number ranges between 0.08 and 0.25 taking the effects coming from the multi-solutions and systematic uncertainties in the analyses into account.
Recent analyses of $`J/\psi 1^{}0^{}`$ decays and $`0^{}0^{}`$ decays show the large amount of strong phases between relevant amplitudes. This fact suggests that, in the relevant $`J/\psi K^{}(892)K\pi `$ decay, the $`K\pi `$ system is not isolated out of the final three-body system, and accordingly in the decay amplitude the phase of the pure $`K\pi `$-scattering amplitude is not directly observed experimentally.
In Methods A and B, the phase of the total $`J/\psi K^{}(892)K\pi `$ amplitude comes from a sum of respective contributions of the $`S`$-matrix elements with the final systems, $`K^{}(892)\kappa `$, $`K^{}(892)K_0^{}(1430)`$, $`K^{}(892)K\pi `$, etc. We obtained the $`\kappa `$ resonance with width, $`\mathrm{\Gamma }_\kappa 500`$MeV/$`c^2`$ in the Breit-Wigner parameterization, which is consistent with the generalized unitarity of S-matrix. This behavior is also consistent with the result of analysis of $`D^+K^{}\pi ^+\pi ^+`$ process in the E791 experiment .
In conclusion, we have shown that the low mass enhancement in the invariant mass spectrum of $`K^+\pi ^{}`$ in the $`J/\psi \overline{K}^{}(892)^0K^+\pi ^{}`$ decays comes neither from phase space, nor from other $`J/\psi `$ decay processes. The angular distribution of $`K`$ in the $`K\pi `$ rest frame in the low mass region shows S-wave decay. Two independent analyses for the process, by the covariant helicity amplitude method and by the VMW method, have been performed, providing a cross check with each other. They reproduce the data well, and the results are in good agreement. The low mass enhancement is well described by the scalar resonance $`\kappa `$, which is highly required in the analyses. Parameter values for BW mass and width of the $`\kappa `$, averaged from those obtained by these two methods, are 878 $`\pm `$ 23$`{}_{55}{}^{}{}_{}{}^{+64}`$ MeV/$`c^2`$ and 499 $`\pm `$ 52$`{}_{87}{}^{}{}_{}{}^{+55}`$ MeV/$`c^2`$, respectively. They are in good agreement with those obtained in the analysis on the $`K\pi `$ scattering phase shifts. The pole position is determined to be $`(841\pm 30_{73}^{+81})i(309\pm 45_{72}^{+48})`$ MeV/$`c^2`$ from the average values.
Acknowledgements
The BES collaboration thanks the staff of BEPC for their hard efforts. This work is supported in part by the National Natural Science Foundation of China under contracts Nos. 10491300, 10225524, 10225525, 10425523, the Chinese Academy of Sciences under contract No. KJ 95T-03, the 100 Talents Program of CAS under Contract Nos. U-11, U-24, U-25, and the Knowledge Innovation Project of CAS under Contract Nos. U-602, U-34 (IHEP), the National Natural Science Foundation of China under Contract No. 10225522 (Tsinghua University), the Department of Energy under Contract No.DE-FG02-04ER41291 (U Hawaii), the Core University Program of Japan Society for the Promotion of Science, JSPS under Contract No. JR-02-B4, the fund for the international collaboration and exchange of RIQS, Nihon-U. |
warning/0506/nucl-th0506049.html | ar5iv | text | # Perfect Fluidity of the Quark Gluon Plasma Core as Seen through its Dissipative Hadronic Corona
## I Introduction
One of the most intriguing experimental findings at the Relativistic Heavy Ion Collider (RHIC) in Brookhaven National Laboratory (BNL) is the large magnitude of the elliptic flow parameter $`v_2`$ Ackermann:2000tr ; Adcox:2002ms ; Back:2002gz in comparison with the smaller values observed at lower collision energies (for results at Super Proton Synchrotron (SPS) energies, see Refs. Alt:2003ab ; Agakichiev:2003gg ; Aggarwal:2004ub ). The magnitude of $`v_2`$ and in particular its transverse momentum $`p_T`$ and mass $`m`$ dependences at RHIC were found to be close to predictions based on ideal, non-dissipative hydrodynamics simulations around midrapidity ($`\eta \stackrel{<}{}\mathrm{\hspace{0.25em}1}`$), in the low transverse momentum region ($`p_T\stackrel{<}{}\mathrm{\hspace{0.25em}1}`$ GeV/$`c`$), and up to semicentral collisions ($`b\stackrel{<}{}\mathrm{\hspace{0.25em}5}`$ fm) Kolb:2000fh ; Hirano:2001eu . This result has led to the recent BNL announcement BNL about the discovery of the near perfect fluidity of the strongly coupled/interacting quark gluon plasma (sQGP) Lee:2005gw ; Gyulassy:2004zy ; Shuryak:2004cy produced in ultra-relativistic nuclear reactions at RHIC.
Until RHIC data, “perfect fluidity” was never observed nor expected to apply theoretically in high energy hadronic or nuclear reactions due to nonvanishing viscous dissipation namiki . Especially, since the discovery of asymptotic freedom in QCD, the prevailing paradigm has been the expectation of large viscosities in a weakly coupled/interacting QGP (wQGP) at very high densities. In addition, it is well established Stoecker:2004qu that the hadronic resonance gas phase of QCD matter is highly dissipative. The discovery of elliptic flow at RHIC consistent with nearly perfect fluidity is therefore an experimental and theoretical surprise. Hence a new name, sQGP, has been adopted to characterize the observed strong coupling properties of the QGP near the critical temperature $`T_c`$ 160–170 MeV that keep viscous effect to a minimum at RHIC.
In this paper, we present the case for the following physical interpretation of RHIC data based on current hydrodynamic analyses: (1) the high density core part of matter produced in relativistic heavy ion collisions, i.e. the sQGP, must expand as a nearly perfect fluid despite of its higher viscosity, (2) the perfect fluidity of the sQGP core is a consequence from a large jump of the entropy density at the critical temperature, $`T_c`$, i.e. deconfinement, and not from some anomalous reduction of its viscosity, (3) viscous effects on its hadronic corona are necessarily large despite its smaller viscosity, and (4) ideal inviscid hydrodynamics should not be applied to the hadronic corona which requires a nonequilibrium transport description.
In Sec. II, we discuss why we expect a surprising monotonic increase of the viscosity of QCD matter through the critical temperature and emphasize the important role played by the rapidly varying viscosity to entropy ratio in connection with perfect fluidity of the sQGP phase. In Sec. III, we discuss different assumptions for the hadronic matter in the hydrodynamic models to clarify what are open issues in the current hydrodynamic approaches. In Sec. IV, the time evolution of the transverse energy per particle is discussed. The mean transverse energy is found to be the key to distinguish the model assumptions in the hadron phase. Results from the hydrodynamic simulations are reviewed in Sec. V. We will show how the perfect fluid description for the hadronic matter in chemical equilibrium in the conventional hydrodynamic simulations leads to accidental reproduction of $`p_T`$ spectrum and $`v_2(p_T)`$. In order to understand analytically the role of chemical freezeout on the transverse dynamics, we employ a blast wave model and give a dynamical meaning to this model in Sec. VI. Finally, summary of this study and an outlook are presented in Sec. VII.
## II Viscosity and Entropy in QCD
Weak coupling perturbative QCD (pQCD) estimates Hosoya:1983xm ; Danielewicz:1984ww ; Thoma:1991em of the viscosity of a wQGP were based on basic kinetic theory relations
$`\eta _{\mathrm{wQGP}}`$ $``$ $`{\displaystyle \frac{4}{15}}ϵ_{\mathrm{SB}}(T)\lambda _{\mathrm{tr}}{\displaystyle \frac{1}{5}}{\displaystyle \frac{T}{\sigma _{\mathrm{tr}}}}{\displaystyle \frac{s_{\mathrm{SB}}(T)}{n_{\mathrm{SB}}(T)}},`$
$`{\displaystyle \frac{\eta _{\mathrm{wQGP}}}{s_{\mathrm{SB}}}}`$ $``$ $`{\displaystyle \frac{T\lambda _{\mathrm{tr}}}{5}}`$ (1)
where (in $`\mathrm{}=c=k_B=1`$ units), $`ϵ_{\mathrm{SB}}(T)=3P_{\mathrm{SB}}(T)=\frac{3}{4}Ts_{\mathrm{SB}}(T)3Tn_{\mathrm{SB}}(T)K_{\mathrm{SB}}T^4`$ is the energy density, pressure, entropy density, and number density of an ideal Stefan-Boltzmann (SB) gas of quarks and gluon characterized by the constant $`K_{\mathrm{SB}}=\frac{\pi ^2}{30}[2(N_c^21)+\frac{7}{8}12N_f]12`$–15 for $`N_c=3,N_f=2`$–3. The key microscopic dynamical quantity in Eq. (II) is the transport mean free path $`\lambda _{\mathrm{tr}}=1/(n_{\mathrm{SB}}\sigma _{\mathrm{tr}})`$ which is controlled in pQCD by the Debye screened transport cross section Danielewicz:1984ww ; Molnar:2001ux
$`\sigma _{\mathrm{tr}}`$ $`=`$ $`{\displaystyle 𝑑\sigma _{\mathrm{el}}\mathrm{sin}^2\theta _{\mathrm{cm}}}`$ (2)
$`=`$ $`{\displaystyle \frac{8\pi \alpha _s^2}{\widehat{s}}}(1+z)\left[(2z+1)\mathrm{ln}\left(1+{\displaystyle \frac{1}{z}}\right)2\right],`$
where $`z=\mu ^2/\widehat{s}`$ and $`\widehat{s}17T^2`$ is the mean partonic Mandelstam variable. Perturbatively, the screening mass squared varies as $`\mu ^24\pi \alpha _sT^2`$. For numerical estimates we take $`\alpha _s(T)4\pi /[18\mathrm{ln}(4T/T_c)]`$, so that $`\alpha _s(T_c)0.5`$. In the range $`T_c<T<5T_c`$ ($`0.5>\alpha _s(T)>0.23`$), the perturbative transport cross section $`\sigma _{\mathrm{tr}}<2`$ mb remains much smaller than typical hadronic cross sections $`\sigma _H1020`$ mb Danielewicz:1984ww ; Gavin:1985ph ; Muronga:2003tb .
An important dimensionless measure of how imperfect or dissipative a fluid may be given by the ratio of viscosity to entropy density, $`\eta /s`$ noteoncausal . This is most easily seen via the Navier-Stokes equation in (1+1)-dimensional boost invariant hydrodynamics Hosoya:1983xm ; Danielewicz:1984ww ; Gavin:1985ph . In the perfect (Euler) fluid limit, the proper energy density decreases with proper time, $`\tau `$, due to longitudinal expansion $`dV=\pi R^2d\tau `$ and $`PdV`$ work via $`dϵ/d\tau =(ϵ+P)/\tau =sT/\tau `$ with a solution $`T=T_0(\tau _0/\tau )^{1/3}`$ for massless particles Bjorken:1982qr . However, shear and bulk viscosity reduce the ability of the system to perform useful work by adding a term $`(4\eta /3+\zeta )/\tau ^2`$ to the right hand side. Neglecting bulk viscosity, $`\zeta `$, that vanishes in equilibrium for massless partons, $`dϵ/d\tau =(sT/\tau )[14(\eta /s)/3T\tau ]`$. This shows that for the earliest times consistent with the uncertainty principle Danielewicz:1984ww , $`\tau 1/3T`$, the cooling of the plasma due to both expansion and work is canceled if $`\eta /s>1/4`$. The ability of the system to convert internal energy into collective flow is thus severely impaired at early times if $`\eta /s`$ and $`\zeta /s`$ are not very small. In fact, in order to reproduce the observed elliptic flow at RHIC, numerical solutions to covariant parton transport equation Molnar:2001ux and blast wave analysis with viscous corrections Teaney:2003pb showed that $`\eta /s`$ had to be less than about 0.2 during the first 3 fm/$`c`$.
The viscosity to entropy ratio in the weakly coupled QCD plasma on the other hand is
$$\left(\frac{\eta }{s}\right)_{\mathrm{wQGP}}=\frac{3}{5}\frac{T}{\sigma _{\mathrm{tr}}}\frac{1}{K_{\mathrm{SB}}T^3}\frac{0.071}{\alpha _s(T)^2\mathrm{ln}[1/\alpha _s(T)]}$$
(3)
This ratio is not small ($`\eta /s=0.35,0.48,0.58,0.66`$ for $`T/T_c=1,1.5,2.0,2.5`$) indicating that the wQGP is expected to be a rather “poor fluid” with large dissipative corrections.
The analytic dependence on $`\alpha _s`$ in Eq. (3) reproduces well the approximate temperature dependence implied by Eq. (2) if we assume the perturbative variation of the screening scale $`\mu `$. Lattice QCD data Kaczmarek:1999mm indicate that $`\mu `$ (2.0–2.5)$`T`$ is not far from the perturbative estimate extrapolated into the physical $`g>1`$ region and that $`\alpha _s(T)<0.5`$ above $`T_c`$ is also reasonable. However, the underlying simple gas kinetic approximation for the viscosity is only rigorously valid in the $`g1`$ region.
Formally, by increasing $`\alpha _s>0.5`$, it would seem that the right hand side of Eq. (3) could be made to be as small as we like if we ignore the $`\mathrm{ln}(1/\alpha _s)`$ singularity. However, by the Heisenberg uncertainty principle, the transport mean free path cannot be localized to less than $`\mathrm{\Delta }x1/p1/3T`$. This leads to a quantum kinetic lower bound on the viscosity for ultrarelativistic gases Danielewicz:1984ww :
$$\frac{\eta }{s_{\mathrm{SB}}}\stackrel{>}{}\frac{1}{15}$$
(4)
with an undetermined multiplicative factor on the order of unity.
A special quantum field theoretic determination of a viscosity lower bound was found recently for infinitely coupled supersymmetric Yang-Mills (SYM) gauge theory using the Anti de-Sitter Space/Conformal Field Theory (AdS/CFT) duality conjecture Kovtun:2004de :
$$\left(\frac{\eta }{s}\right)_{\mathrm{SYM}}=\frac{1}{4\pi }.$$
(5)
This bound is obtained in the dual $`N_c=\mathrm{}`$ and $`g^2N_c=\mathrm{}`$ limits of the special $`𝒩=4`$ conformal SYM schalm . It is interesting to note that this analytic SYM bound is close to the simple kinetic theory uncertainty principle bound in Eq. (4). It has been conjectured Kovtun:2004de that $`1/4\pi `$ in Eq. (5) is the universal minimal viscosity to entropy ratio even for QCD. In that case, the viscosity of the sQGP could be up to a factor of $`1/2\pi `$ smaller than of a wQGP. It is then tempting to conclude that the sQGP must have anomalously small viscosity if perfect fluid behavior is observed. However, as we show below, the sQGP viscosity is actually very close to that of ordinary hadronic matter just below $`T_c`$.
To develop this argument further, we first digress to recall that the entropy density in the $`N_c1,g^2N_c1`$ limits of $`𝒩=4`$ SYM is given by Gubser:1998nz
$$s_{\mathrm{SYM}}=\left[\frac{3}{4}+\frac{0.6}{(g^2N_c)^{3/2}}+O\left(\frac{1}{N_c^2}\right)\right]\frac{4}{3}K_{\mathrm{SYM}}T^3.$$
(6)
where the Stefan-Boltzmann constant for $`𝒩=4`$ SYM is $`K_{\mathrm{SYM}}=\pi ^2(N_c^21)/239.5`$ is about 3 times greater than $`K_{\mathrm{SB}}`$ of our QCD world schalm . What is especially remarkable about Eq. (6) is that, at infinitely strong coupling, the entropy density is only reduced by $`25\%`$ from its non-interacting SB value. On the other hand, the viscosity in this extreme limit is reduced about an order of magnitude from the weak coupling value and limited only by the quantum (Heisenberg uncertainty) bound on the effective scattering rate. Current lattice data on the QCD viscosity near $`T_c`$ Nakamura:2004sy are with large numerical error bars between these weak and super strong coupling limits but the relatively small deviation of the lattice entropy density from the SB limit is consistent with Eq. (6).
The AdS/CFT lower bound (5) together with the assumed universal $`3/4`$ reduction of the SB entropy density implies that the absolute value of the sQGP viscosity at $`T_c`$ would be
$$\eta _{\mathrm{sQGP}}(T)\eta _{\mathrm{SYM}}(T)=\frac{K_{\mathrm{SB}}T^3}{4\pi }T_c^3\left(\frac{T}{T_c}\right)^3$$
(7)
where we used a fact that for QCD $`K_{\mathrm{SB}}12`$–15 is accidentally close to $`4\pi `$. The monotonic increase of $`\eta _{\mathrm{SYM}}(T)`$ is illustrated by the dashed curve in Fig.1.
The effective transport cross section via Eq. (II) at $`T_c160`$ MeV is in this case
$$\sigma _{\mathrm{tr}}^c\frac{4}{5}\frac{T_c}{\eta _c}12\mathrm{mb}.$$
(8)
Here $`\eta _cT_c^3=0.106`$ GeV/fm<sup>2</sup> at $`T_c=0.16`$ GeV. See Ref. Peshier:2005pp for an independent estimate of the transport cross section in the sQGP phase leading to similar $`\sigma _{\mathrm{tr}}(T)`$ near $`T_c`$.
While there is no consensus yet on what physical mechanisms could enforce the minimal viscosity bound in the sQGP Molnar:2001ux ; Shuryak:2003xe ; Xu:2004mz , we take as empirical fact that the sQGP viscosity must be close (within a factor of two) to the minimal (uncertainty) bound, Eq. (7). Our central assumption is that local thermal equilibrium is maintained in the sQGP core with minimal dissipative deviations and with the equation of state and hence speed of sound as predicted by QCD. Alternate scenarios, with arbitrary equations of state with higher speed of sound that in principle could compensate the higher dissipation and viscosity in a wQGP will not be considered here. In this connection we also emphasize the importance of fixing sQGP initial conditions with Color Glass Condensate or saturating gluon distributions constrained by the global entropy observables Gyulassy:2004zy ; Hirano:2004rs . With fixed initial conditions and equation of state, the remaining degrees of dynamical freedom are reduced to the dissipation corrections discussed in this section for the sQGP phase and the dynamical constraints on its dissipative hadronic corona discussed in the subsequent sections.
Note that the effective transport cross section in the sQGP $`\sigma _{\mathrm{tr}}^c`$ just above $`T_c`$ is remarkably close to the hadron resonance gas transport cross section just below $`T_c`$ Gavin:1985ph ; Muronga:2003tb . However, due to the $`1/T^2`$ scaling at $`T2T_c`$, the effective transport cross section in the sQGP would already drop to $`3`$ mb while preserving the (uncertainty principle) lower bound Eq. (5).
In contrast to the novel sQGP phase above $`T_c`$, for $`T<T_c`$, matter is well known to be in the confined hadron resonance gas (HRG) phase where the kinetic theory viscosity Danielewicz:1984ww ; Gavin:1985ph is
$$\eta _{\mathrm{HRG}}\frac{T}{\sigma _H}\eta _c\frac{T}{T_c},$$
(9)
as illustrated by the solid curve below $`T_c`$ in Fig. 1. Because the hadronic transport cross sections are typically $`\sigma _H1020`$ mb, the combination of Eqs. (7) and (9) shows that we should not expect a large variation of the absolute value of the matter viscosity across $`T_c`$ if the minimal $`\eta /s`$ holds above $`T_c`$. In Ref. Gavin:1985ph , Gavin found that for a pion gas with P-wave $`\rho `$ and D-wave $`f^0`$ resonance interactions, the thermal averaged transport cross section and viscosity from his Fig. 3 are $`(\eta /\eta _c,\sigma _H[\mathrm{mb}])(0.66,20),(0.9,17),(1.1,15)`$ for $`T=180,160,140`$ MeV. In Ref. Muronga:2003tb , Muronga used the UrQMD resonance gas Monte Carlo to compute $`\eta (T)/\eta _c0.75,1.1,1.9`$ for $`T=0.14,0.16,0.18`$ GeV. In both studies Gavin:1985ph ; Muronga:2003tb numerical estimates are thus consistent with Eq. (9) for $`T<T_c`$. For nonvanishing baryon density, see recent estimates in Ref. Muroya:2004pu .
In the sQGP phase, the minimal viscosity, Eq. (7), is predicted to grow cubically with $`T`$ beyond $`T_c`$. However, at $`TT_c`$ asymptotic freedom predicts that it would grow even more rapidly as the sQGP transforms gradually into a wQGP. An interpolation formula between these phases can be constructed as
$$\eta (T)T_c^3\{\begin{array}{cc}(T/T_c)^1,\hfill & T<T_c\hfill \\ (T/T_c)^3[1+w(T)\mathrm{ln}(T/T_c)]^2,\hfill & T>T_c\hfill \end{array}$$
(10)
The extra squared logarithmic growth of the viscosity at $`TT_c`$ is expected from Eq. (3). To be consistent with the perturbative wQGP at $`TT_c`$ one should take
$$\frac{w^2(TT_c)}{4\pi }=\frac{9\beta _0^2}{80\pi ^2K_{\mathrm{SB}}}\frac{1}{\mathrm{ln}1/\alpha _s(T)}.$$
(11)
With $`K_{\mathrm{SB}}=12`$–15 and $`\beta _0=112N_f/39`$–10, a possible scenario may be that $`w1`$ already near $`T_c`$. This possibility is shown in Fig. 1 by the solid curve above $`T_c`$ which would imply sQGP$``$ wQGP already above $`(23)T_c`$. In fact, current lattice data on the evolution of screening scales in the QGP phase suggest that hadronic scale correlations may persist only up to $`T3T_c`$ Asakawa:2003re ; Datta:2003ww . A value $`w(T>2T_c)1`$, is also not inconsistent with current lattice results Nakamura:2004sy . We note that future measurements of elliptic flow in A+A collisions at LHC with $`\sqrt{s_{NN}}=5500`$ GeV will be able to test experimentally if such a precocious onset of dissipative wQGP dynamics occurs.
The approximate continuity of the viscosity across $`T_c`$ indicated in Fig. 1 is to be understood to hold within a factor on the order of unity. What changes rapidly at $`T_c`$ is not the viscosity of QCD matter but rather its entropy density due to the deconfinement of the quark and gluon degrees of freedom.
For a hadronic resonance gas charactered by a speed of sound $`c_H^2`$, the entropy density $`s(T)=s_H(T/T_c)^{1/c_H^2}`$ with decreasing temperature decreases much more rapidly than does the viscosity for typical $`c_H^21/6`$–1/3. Just beyond $`T_c`$ –possibly only up to several times $`T_c`$– it is postulated that the sQGP phase may exist with $`\eta /s<0.2`$ but with viscosity close to ideal gas.
Summarizing the discussion up to now, we expect that $`\eta `$ varies smoothly near $`T_c`$ as in Fig. 1 while the ratio $`\eta /s`$ may have a significant discontinuity due to the rapid onset of deconfinement as a function of $`T`$ as shown in Fig. 2. We therefore propose the following interpolation formula for the temperature dependence of the $`\eta /s`$ ratio with a $`T`$ independent constant $`w`$
$$\frac{\eta (T)}{s(T)}\frac{1}{4\pi }\{\begin{array}{cc}\left(\frac{s_Q}{s_H}\right)\left(\frac{T}{T_c}\right)^{11/c_H^2},\hfill & T<T_c\hfill \\ \left[1+w\mathrm{ln}(T/T_c)\right]^2,\hfill & T>T_c\hfill \end{array}$$
(12)
with the negative discontinuity
$$\left[\frac{\eta }{s}\right]_{T_c}=\frac{\eta (T_c)}{s(T_c)}|_Q\frac{\eta (T_c)}{s(T_c)}|_H=\frac{1}{4\pi }\left(\frac{s_Q}{s_H}1\right).$$
(13)
We illustrate Eq. (12) in Fig. 2. Note that it is the entropy jump $`s_Q/s_H3`$–10 that causes a drop of $`\eta /s`$ across $`T_c`$. Since the HRG$``$ sQGP transition with dynamical quarks appears from the lattice QCD to be only a rapid crossover, the discontinuity is understood to be spread out over a temperature range $`\mathrm{\Delta }T_c/T_c0.1`$.
## III Imperfect Fluidity of the Hadronic Corona
In the last section we presented the case that the $`\eta /s`$ ratio may be small enough above $`T_c`$ in the sQGP core for perfect fluidity to arise during the critical first $`3`$ fm/$`c`$, while the spatial azimuthal asymmetry of the matter produced in non-central reactions is large enough to induce collective elliptic flow. However, during the subsequent $`10`$ fm/$`c`$ evolution after hadronization of the sQGP core, the whole system evolves as HRG corona. In the HRG, $`\eta /s`$ is too high for local equilibrium to be maintained due to its small entropy density compared with sQGP. Nevertheless, the data on $`v_2(p_T)`$ seem to agree very well with some hydrodynamic predictions based on the assumption that local equilibrium is maintained until thermal freeze-out. However, various assumptions about the hadro-chemical evolution are known to lead to significantly different predictions for the differential elliptic flow. In this section we focus on the question of the validity of the application of hydrodynamics to analyze the entire evolution in A+A at RHIC.
The key problem that we now address is the role of final state hadronic interactions in possibly modifying conclusions inferred about the prefect fluidity of the sQGP core. In order that the sQGP elliptic flow signature of perfect fluidity survives during the evolution through the extended hadronic “corona” we must study how longitudinal flow, transverse radial flow, as well as the elliptic deformation of the transverse flow may evolve in hadronic matter.
Several puzzling features suggest the importance of investigating more closely the distortions that can be caused by final state hadronic interactions involving hadro-chemical and thermal freezeout. In ideal hydrodynamics it is well known that while the central rapidity region is well reproduced by hydrodynamics, this is not the case in forward/backward rapidity regions Hirano:2001eu . Hydrodynamics also strongly overestimates $`v_2`$ at energies below $`\sqrt{s_{NN}}=`$200 GeV as well as in the most peripheral collisions where initially a larger fraction of the transverse elliptic interaction region starts out in the hadronic phase. All these data point to the fact that the dynamics in the hadronic corona is not at all ideal.
Another important issue in ideal hydrodynamics as well as in other dynamical models is the so-called HBT puzzle Heinz:2002un . In spite of the apparent success of hydrodynamic description for elliptic flow, hydrodynamics fails to reproduce the experimental data of the HBT radii Adler:2001zd ; Adcox:2002uc ; Back:2004ug . The best current description of hadron freezeout consistent with the HBT puzzle involves an assumption of a highly dissipative resonance gas dynamics and transport Lin:2002gc .
As compiled in Fig. 20 in Ref. Adcox:2004mh , some hydrodynamic results reproduce neither $`v_2(p_T)`$ nor $`p_T`$ spectrum. This immediately raises the following two questions:
(Q1) Are hydrodynamic results consistent with each other at RHIC energies? What assumptions lead to the differences among hydrodynamic predictions?
(Q2) How robust is the statement that hydrodynamic description at RHIC data is good at low $`p_T`$?
The differences arise from the treatment of hadron phase dynamics. The treatment of the sQGP phase is essentially the same in all approaches so far: Parton chemical equilibrium and inviscid hydrodynamics are assumed in the sQGP phase. There are, however, three generic classes of approaches to the evolution of the hadronic corona in the literature.
Chemical equilibrium model (CE). Most of the hydrodynamic calculations so far are based on the assumption that the hadron phase is a perfect fluid in both hadro-chemical and thermal equilibrium. With this assumption, the centrality, transverse momentum, and/or (pseudo)rapidity dependences of elliptic flow are studied Kolb:2000fh ; Hirano:2001eu . However, it is known that the yields of heavy hadrons (essentially all hadrons except for pions) are smaller in CE than data since the numbers of particle are counted on the hypersurface at thermal freezeout within this approach. At relativistic collisional energies, thermal freezeout temperature $`T^{\mathrm{th}}`$ is smaller than chemical freezeout temperature $`T^{\mathrm{ch}}`$ Heinz:1997za ; Shuryak:1997xb . So the numbers of heavy particles are suppressed due to the Boltzmann factor and, eventually, lead to discrepancy from the data. CE therefore fails to account for the observed particle abundance systematics Braun-Munzinger:2003zd .
Partial chemical equilibrium model (PCE). In Refs. Arbex:2001vx ; Hirano:2002ds ; Teaney:2002aj ; Kolb:2002ve , chemical freezeout being separated from thermal freezeout is taken into account in the hydrodynamic simulations toward simultaneous reproduction of particle ratios and particle spectra. Below $`T^{\mathrm{ch}}`$, one introduces chemical potential for each hadron associated with the conserved number. Inelastic processes only through strong interactions are supposed to be equilibrated in the hadron phase, e.g. $`\mu _\rho =2\mu _\pi `$, $`\mu _\mathrm{\Delta }=\mu _N+\mu _\pi `$, and so on. Note that the conserved pion number within this approach is $`\stackrel{~}{N}_\pi =N_\pi +_Rb_RN_R`$. Here $`b_R`$ is the effective branching ratio and $`N_i`$ is the number of $`i`$-th hadron Bebie:1991ij . For details, see Refs. Hirano:2002ds ; Teaney:2002aj ; Kolb:2002ve . This particular model does not reproduce $`v_2(p_T;m)`$ nor $`p_T`$ spectra at RHIC as shown in Fig. 20 in Ref. Adcox:2004mh . Note that a model is called “chemical freezeout (CFO)” when the number of hadrons $`N_i`$ instead of $`\stackrel{~}{N}_i`$ is conserved. This means even inelastic scatterings through strong interaction cease to happen.
Hadronic cascade model (HC). One can utilize a hadronic cascade model just after the phase transition between the QGP and hadron phases Bass:2000ib ; Teaney:2000cw . This approach dynamically describes both chemical and thermal freezeouts without assuming explicit freezeout temperatures. Viscous effects are effectively taken into account through the non-zero mean free path among the particles (see, e.g. Eqs. (II) and (9)).
Hydrodynamic results from the above three classes of hadro-dynamical models are summarized in Table 1. For reviews of hydrodynamic results at the RHIC energies, see also Huovinen:2003fa ; Kolb:2003dz ; Hirano:2004ta . As long as the space-time evolution of the hadron phase is concerned, the approach HC seems to be the most realistic one among the three classes. The second best model should be the model PCE from the experimental data of particle ratios and spectra. The models CE and HC reproduce $`v_2(p_T)`$ for pions and protons in low $`p_T`$ regions, whereas the model PCE fails. The failure of PCE for $`v_2(p_T)`$ is particularly perplexing since the spatial azimuthal asymmetry is mostly gone by the time hadronization is competed. In the following sections, we shall show why the differences between the models CE and PCE appears and why the result from the perfect fluid model CE eventually looks similar to the one from the highly dissipative model HC. The key quantities to understand these differences are found to be the temporal behavior of the mean transverse momentum/energy and the ratio of the particle number to the entropy.
## IV Temporal behavior of transverse energy
In this section, we briefly review the time evolution of total transverse energy and transverse energy per particle in relativistic heavy ion collisions within a framework of the Bjorken solution for longitudinal expansion Bjorken:1982qr .
Assuming the Bjorken scaling solution, the time evolution of the entropy density is represented by $`s(\tau )=s_0\tau _0/\tau `$. Here $`s_0`$ is the initial entropy density and $`\tau _0`$ is the initial time. As long as a perfect fluid is considered, the entropy per unit rapidity is conserved
$`{\displaystyle \frac{dS}{dy}}`$ $`=`$ $`{\displaystyle d^2x_{}\tau s(\tau )}=A_{}\tau _0s_0,`$ (14)
where $`A_{}`$ is the transverse cross section of a nucleus. Here we assume a smooth function for the equation of state (EOS). For EOS with $`P=c_s^2ϵ`$ $`(0<c_s^2<1/3)`$, the time evolution of energy density becomes
$`ϵ(\tau )`$ $`=`$ $`ϵ_0\left({\displaystyle \frac{\tau _0}{\tau }}\right)^{1+c_s^2}`$ (15)
where $`ϵ_0`$ is the energy density at $`\tau _0`$. Thus the transverse energy per unit rapidity
$`{\displaystyle \frac{dE_T}{dy}}`$ $`=`$ $`{\displaystyle d^2x_{}\tau ϵ(\tau )}=A_{}ϵ_0\tau _0\left({\displaystyle \frac{\tau _0}{\tau }}\right)^{c_s^2}`$ (16)
decreases with proper time due to $`PdV`$ work in spite of the conservation of entropy Gyulassy:1983ub ; Ruuskanen:1984wv . In the following discussion in this section, we consider three EOS models for pions and see time evolution of total transverse energy and transverse energy per pion.
Massless Pions. The number density $`n`$ is proportional to the entropy density $`s`$ for massless pions: $`n=(d/\pi ^2)\zeta (3)T^3`$, $`s=d(4\pi ^2/90)T^3`$. Inserting $`d=3`$ and $`\zeta (3)1.2`$, one obtains $`s3.6n`$. From Eq. (14), the number of pions per unit rapidity $`dN/dy`$ is also conserved. Therefore the mean transverse mass, which is identical to the mean transverse momentum in the massless pion case, $`m_T=(dE_T/dy)/(dN/dy)(dE_T/dy)`$ decreases with $`\tau `$ from Eq. (16).
Massive Pions in Chemical Equilibrium. The proportionality between the number density and the entropy density is approximately valid only for ultra-relativistic particles ($`Tm`$). In relativistic heavy ion collisions, the typical temperature is around the order of the pion mass. So pions are no longer relativistic particles in this particular situation. The number density and the entropy density for pions are evaluated in the usual prescription of statistical physics. Thus the ratio of the number $`N`$ and the entropy $`S`$ increases with temperature due to the finite mass of pions as shown in Fig. 3. Note that the volume $`V`$ is canceled and that $`N/S`$ equals to the ratio of their densities $`n/s`$. Therefore $`m_T`$ can increase with proper time (or with decreasing temperature of the system) even as $`dE_T/dy=m_TdN/dy`$ decreases from Eq. (16). This “local reheating” can occur because the total transverse energy is distributed among a smaller number of pions at lower temperature. The resulting temporal behavior of $`m_T`$ is determined through competition of how rapidly $`dE_T/dy`$ and $`dN/dy`$ decrease with proper time. As we will see in the next section, $`p_T`$ increases with proper time in hydrodynamic simulations with chemical equilibrium EOS.
Chemically Frozen Massive Pions. When the system expands strongly, inelastic collisions can cease to happen. So one can think about a situation in which the system keeps only thermal equilibrium through elastic scattering. Analyses based on statistical models and blast wave models show that thermal freezeout temperature $`T^{\mathrm{th}}`$ is smaller than chemical freezeout temperature $`T^{\mathrm{ch}}`$ at AGS, SPS and RHIC energies. Moreover, $`T^{\mathrm{ch}}`$ is found to be close to the (pseudo)critical temperature $`T_c`$. This indicates that the hadron phase in relativistic heavy ion collisions is in thermal equilibrium, not in chemical equilibrium Heinz:1997za ; Shuryak:1997xb . Usually, the term “chemical equilibrium” is associated with a state equilibrated among a finite number ($`>1`$) of compositions in the system. Here we simply use the term “chemical freezeout” in spite of one hadronic species, i.e. pions. This means the number of pions per unit rapidity is fixed below $`T^{\mathrm{ch}}`$. The entropy is also conserved as long as a perfect fluid is considered, so the ratio of the number density and the entropy density is a constant of motion similar to the case for massless pions. It is interesting to mention that the entropy is not generated even in the evolution of chemically frozen fluids. Entropy production originates from the source term in the rate equation for chemical non-equilibrium processes. The number of hadron is, however, conserved after chemical freezeout. It is easy to show the conservation of entropy $`_\mu S^\mu =0`$ from the conservation of energy and momentum $`_\mu T^{\mu \nu }=0`$, where $`T^{\mu \nu }=(ϵ+P)u^\mu u^\nu Pg^{\mu \nu }`$, and the conservation of the number of hadrons $`_\mu N_i^\mu =0`$. In this sense, one needs to distinguish “chemical freezeout” from “chemical non-equilibrium”. In the chemical non-equilibrium process, the system is approaching to chemical equilibrium state, i.e. the maximum entropy state through inelastic scattering. Entropy is certainly generated during this process. Contrary to this, chemical freezeout means that the system suddenly leaves from the chemical equilibrium state by keeping particle ratios due to the strong expansion.
Figure 3 shows comparison of the ratio of pion number $`N`$ and its entropy $`S`$ between chemical equilibrium EOS and chemically frozen EOS. Here chemical freezeout temperature is taken as being $`T^{\mathrm{ch}}=170`$ MeV. Similar to the massless pion case, $`m_T`$ in the chemical freezeout case decreases with decreasing decoupling temperature. As long as the Bjorken scaling solution for longitudinal expansion is considered, transverse expansion does not spoil the above discussion: $`PdV`$ work done in the transverse direction is simply converted into the kinetic energy of fluid elements. The resultant slope of $`p_T`$ spectrum for pions should become softer at lower decoupling (thermal freezeout) temperature. Quantitatively, the $`p_T`$ slope is insensitive to the choice of $`T^{\mathrm{th}}`$ since $`dE_T/dy\tau ^{c_s^2}`$ changes only gradually in the late stage. The universal behavior of the $`p_T`$ slope is already confirmed in the hydrodynamic simulations with chemically frozen (or partial chemical equilibrium) EOS Hirano:2002ds ; Kolb:2002ve and will be mentioned in the next section.
From the above consideration, the key quantity which governs the transverse dynamics, particularly the time evolution of mean transverse mass, within ideal hydrodynamics is found to be the ratio of the number $`N`$ and the entropy $`S`$. It is commonly expected that the behavior of the mean transverse energy/momentum is governed by the stiffness of the EOS. But it is not always true since the sound velocity of chemical freezeout EOS (or simply the slope of $`P`$ as a function of $`ϵ`$) is almost the same as that of chemical equilibrium EOS as shown in Ref. Hirano:2002ds . Interestingly, whether the mean transverse energy increases or decreases with the proper time is governed by $`N/S`$ and the longitudinal work, not the stiffness of EOS. We will also mention these hydrodynamic results in the next section.
To summarize this section, $`m_T`$ decreases with proper time for massless pions or chemically frozen massive pions, while it can increase for massive pions in chemical equilibrium. We emphasize here that these conclusions are obtained from quite basic assumptions: the first law of thermodynamics ($`PdV`$ work) and the Bjorken scaling solution.
## V Results from Hydrodynamic Simulations
We compare the hydrodynamic results from the model PCE with the ones from the model CE with respect to the EOS, space-time evolution, $`p_T`$ spectra, and $`v_2(p_T)`$. Hydrodynamic simulations have already been performed for Au+Au collisions at $`\sqrt{s_{NN}}=130`$ GeV and the essential results have already been obtained in Ref. Hirano:2002ds . In this section, we make an interpretation of these hydrodynamic results obtained so far. In particular, we emphasize that the temporal behavior of the mean $`p_T`$ is the key to understand the difference of the results between these two models. For further details of the hydrodynamic model, see also Ref. Hirano:2002ds .
### V.1 Equation of state and space-time evolution
Chemical freezeout does not change EOS, i.e. pressure as a function of energy density $`P(ϵ)`$, so much Hirano:2002ds ; Teaney:2002aj . This means that the difference of chemical composition in the hadron phase does not affect the space-time evolution of energy density significantly. However, at a fixed temperature, the energy density in chemically frozen (or partial chemical equilibrium) hadronic matter is considerably larger than the one for hadronic matter in chemical equilibrium. This is due to the fact that the large resonance population keeps in the system during the expansion after chemical freezeout and that the mass terms significantly contributes to the energy density. Therefore the space-time evolution of temperature field does change significantly while $`P(ϵ)`$ remains essentially unchanged. The temperature of the chemically frozen system cools down more rapidly than that of the chemical equilibrium system. This leads to the reduction of life time of fluids, radial flow, and HBT radii ($`R_{\mathrm{long}}`$ and $`R_{\mathrm{out}}`$) Hirano:2002ds . Longitudinal size $`R_{\mathrm{long}}`$, which relates with the lifetime of a fluid through the gradient of longitudinal flow velocity $`(dv_z/dz)^1\tau _f`$ Makhlin:1987gm , can be interpreted by the effect of chemical freezeout on the life time of a fluid in hydrodynamics. Nevertheless, $`R_{\mathrm{side}}`$ and $`R_{\mathrm{out}}`$ are still inconsistent with data.
### V.2 $`p_T`$ spectra and elliptic flow
In hydrodynamic simulations with chemical equilibrium, thermal freezeout temperature $`T^{\mathrm{th}}`$ is an adjustable parameter. In order to fix $`T^{\mathrm{th}}`$, one usually utilizes the experimental data of $`p_T`$ spectra for hadrons. Reduction of $`T^{\mathrm{th}}`$ leads to generation of additional radial flow. Generically, the resulting $`p_T`$ spectra at $`T^{\mathrm{th}}=T_1`$ becomes harder than the ones at $`T^{\mathrm{th}}=T_2>T_1`$ even though temperature, i.e. the inverse slope of the momentum distribution in the local rest frame decreases. Contrary to this behavior, when chemical freezeout is appropriately taken into account in hydrodynamic simulations, the $`p_T`$ slope becomes insensitive to $`T^{\mathrm{th}}`$ compared to the one in the model CE Hirano:2002ds .
To confirm these behaviors, we perform hydrodynamic simulations again for a particular choice of impact parameter $`b=5`$ fm and see the average transverse momentum $`p_T`$. Details of the hydrodynamic models used here can be found in Ref. Hirano:2002ds ; Hirano:2003pw . Figure 4 shows $`p_T`$ for pions as a function of $`T^{\mathrm{th}}`$ at midrapidity $`y=0`$. In chemical equilibrium, $`p_T`$ increases with decreasing $`T^{\mathrm{th}}`$. On the contrary, $`p_T`$ gradually decreases with decreasing $`T^{\mathrm{th}}`$ when early chemical freezeout is taken into account. Even in the case that a full hydrodynamic simulation without boost invariant ansatz is performed and that the contribution from resonance decays is included, the temporal behavior of $`p_T`$ as already discussed in Sec. IV is seen in Fig. 4. It should be emphasized here that increase of $`p_T`$ in chemical equilibrium is a direct consequence of neglecting chemical freezeout in hydrodynamic calculations and, more definitely, of the experimental results of particle ratios. One has made full use of this unrealistic behavior to reproduce $`p_T`$ spectra at the cost of hadron ratios in the conventional hydrodynamic calculations.
In chemical equilibrium, the slope of elliptic flow parameter $`dv_2(p_T)/dp_T`$ for pions is insensitive to $`T^{\mathrm{th}}`$ and is consistent with the experimental data. See also Figs. 8 and 11 in Ref. Hirano:2002ds . This is apparently plausible since the elliptic flow is a self-quenching phenomenon and is sensitive to the early stage of the collisions Ollitrault:1992bk . On the other hand, $`dv_2(p_T)/dp_T`$ for pions in the model PCE increases with decreasing $`T^{\mathrm{th}}`$ and is ended with overprediction of the experimental data when $`T^{\mathrm{th}}`$ is chosen so as to reproduce the proton $`p_T`$ spectrum and $`v_2(p_T)`$ in the low $`p_T`$ region. This means that $`v_2(p_T)`$ for pions varies in the late hadronic stage ($`\tau \stackrel{>}{}`$ 10 fm/$`c`$).
We have to be careful in understanding the difference between integrated elliptic flow $`v_2`$ and differential elliptic flow $`v_2(p_T)`$. $`v_2`$ reflects the momentum anisotropy of bulk matter, while $`v_2(p_T)`$ represents how total $`v_2`$ distributes among various particles with different $`p_T`$. As shown in Fig. 5, the integrated $`v_2`$ varies only weakly with decrease of $`T^{\mathrm{th}}`$ (and hence proper time $`\tau `$) in both cases. This is consistent with the self-quenching picture of elliptic flow again.
The slope of $`v_2(p_T)`$, on the other hand, can be evaluated approximately by $`dv_2(p_T)/dp_Tv_2/p_T`$ since $`v_2(p_T)`$ for pions is almost a linear function of $`p_T`$ noteonv2slope . In chemical equilibrium (CE), the moderate increase of $`v_2`$ ($``$13% increase as $`T^{\mathrm{th}}`$ decreases from 160 MeV to 100 MeV) is approximately canceled by the increase of $`p_T`$ ($``$22% from 160 MeV to 100 MeV). Eventually, the ratio remains almost constant (or even decreases slightly) as shown in Fig. 6. The CE predictions for the differential elliptic flow work because without chemical freezeout the slope of $`v_2(p_T)`$ stalls accidentally and reproduces the experimental data.
On the contrary, in PCE the numerator $`v_2`$ increases by $``$20% and the denominator $`p_T`$ decreases by $``$10%. The resultant ratio $`v_2/p_T`$ therefore increases with decreasing $`T^{\mathrm{th}}`$ as shown in Fig. 6. This is the reason why the slope of $`v_2(p_T)`$ in PCE increases even in the late stage after the spatial azimuthal asymmetry has reversed sign.
It is now easy to understand why $`v_2(p_T)`$ at the SPS energies is so close to the one at the RHIC energy (see, e.g. Fig. 17 in Ref. Adams:2005dq ). This is due to the correlated change of both $`v_2`$ and $`p_T`$ from SPS to RHIC energies: The increase of the average $`v_2`$ is compensated for by the increase of $`p_T`$. The slopes of $`v_2(p_T)`$ therefore vary surprisingly weakly with the beam energy.
Summarizing the discussion, differential elliptic flow is sensitive to the late hadronic stage in ideal hydrodynamic calculations with early chemical freezeout since the slope of $`v_2(p_T)`$ is determined by the mean $`p_T`$, i.e. radial flow. The apparent consistency of $`v_2(p_T)`$ between RHIC data and results based on conventional CE hydrodynamic simulations is therefore fortuitous.
## VI Analytic Approach
In order to understand the effect of chemical freezeout on the transverse momentum dependence of elliptic flow analytically, we employ the blast wave model discussed in Ref. Pratt:2004zq . In the framework of the blast wave model, one can choose the radial flow parameter and temperature independently to reproduce the slope of $`p_T`$ spectra. However, these two parameters are correlated in actual hydrodynamic simulations. After chemical freezeout, the mean $`p_T`$ decreases with decreasing $`T^{\mathrm{th}}`$, i.e. $`dp_T/dT>0`$, as already discussed in the previous sections. On the contrary, the mean $`p_T`$ increases with decreasing $`T^{\mathrm{th}}`$, i.e. $`dp_T/dT<0`$, in chemical equilibrium. So we can constrain the average flow velocity as a function of temperature through the condition $`dp_T/dT=0`$. We call the obtained radial flow the critical radial flow, $`v_r^{\mathrm{crit}}(T)`$. The critical radial flow ensures that the mean $`p_T`$ is a constant of motion. Qualitatively, $`v_r(T)<v_r^{\mathrm{crit}}(T)`$ corresponds to the chemically frozen system, while $`v_r(T)>v_r^{\mathrm{crit}}(T)`$ corresponds to the system under chemical equilibrium. In the following in this section, we assume only pions which are dominant particles in a fluid element.
### VI.1 Momentum Distribution
The invariant momentum spectrum is given by the Cooper-Frye formula Cooper:1974mv in the Boltzmann approximation:
$`E{\displaystyle \frac{dN}{d^3p}}`$ $`=`$ $`{\displaystyle \frac{g}{(2\pi )^3}}{\displaystyle p^\mu 𝑑\sigma _\mu \mathrm{exp}\left(\frac{p^\mu u_\mu \mu }{T}\right)}.`$ (17)
Here $`g`$ is the degree of freedom, $`d\sigma ^\mu `$ is the element of freezeout hypersurface. $`p^\mu `$ is the four momentum measured in the laboratory system. $`\mu `$ and $`T`$ are, respectively, chemical potential and temperature at freezeout. Four fluid velocity ($`u^\mu u_\mu =1`$) can be parametrized as
$`u^\mu `$ $`=`$ $`\mathrm{cosh}\rho (\mathrm{cosh}\eta _f,\mathrm{sinh}\rho \mathrm{cos}\varphi ,\mathrm{sinh}\rho \mathrm{sin}\varphi ,\mathrm{sinh}\eta _f).`$
Here $`\rho `$ is the transverse rapidity and $`\eta _f`$ is the longitudinal rapidity which is to be identified with the space-time rapidity $`\eta _s`$ in the Bjorken boost invariant solution Bjorken:1982qr . One can also write $`\mathrm{cosh}\rho =\sqrt{1+u_{}^2}`$ and $`\mathrm{sinh}\rho =u_{}`$, respectively. Energy of a particle in the local rest frame becomes
$`p^\mu u_\mu =m_T\mathrm{cosh}(y\eta _f)\sqrt{1+u_{}^2}p_Tu_{}\mathrm{cos}(\varphi _p\varphi ).`$
Here $`y`$ is the momentum rapidity and $`\varphi _p`$ is the azimuthal angle of the momentum. According to Ref. Pratt:2004zq , we also make the same ansatz $`u_{}(\varphi )=u_{}(1+\epsilon \mathrm{cos}2\varphi )`$ for azimuthal dependence of radial flow and take terms up to the first order in $`\epsilon `$
$`p^\mu u_\mu `$ $`=`$ $`m_T\mathrm{cosh}(y\eta _f)\sqrt{1+u_{}^2}\left[1+{\displaystyle \frac{u_{}^2}{1+u_{}^2}}\epsilon \mathrm{cos}2\varphi \right]`$ (19)
$``$ $`p_Tu_{}\mathrm{cos}(\varphi _p\varphi )(1+\epsilon \mathrm{cos}2\varphi ).`$
Assuming the matter suddenly freezes out at $`\tau _f`$,
$`p^\mu d\sigma _\mu =EdV=m_T\mathrm{cosh}y\times rdrd\varphi \tau _fd\eta _s.`$ (20)
Note that $`u_{}(\varphi )`$ is supposed to include all possible azimuthal anisotropic effects and that $`\varphi `$ dependences of $`T`$ and $`\mu `$ are neglected as in the hydrodynamic approach. Then Eq. (17) reduces to
$`{\displaystyle \frac{dN}{m_Tdm_Td\varphi _pdy}}`$ $``$ $`{\displaystyle 𝑑\varphi m_T\mathrm{cosh}ye^A},`$ (21)
$`A`$ $`=`$ $`{\displaystyle \frac{m_T\mathrm{cosh}(y\eta _f)\sqrt{1+u_{}^2}}{T}}`$ (22)
$`+`$ $`{\displaystyle \frac{p_Tu_{}\mathrm{cos}(\varphi _p\varphi )}{T}}+{\displaystyle \frac{\mu }{T}}`$
$``$ $`\epsilon {\displaystyle \frac{m_Tu_{}^2\mathrm{cosh}(y\eta _f)\mathrm{cos}2\varphi }{T\sqrt{1+u_{}^2}}}`$
$`+`$ $`\epsilon {\displaystyle \frac{p_Tu_{}\mathrm{cos}2\varphi \mathrm{cos}(\varphi _p\varphi )}{T}}.`$
### VI.2 Azimuthal Anisotropy
By using the above momentum distribution, one can calculate $`v_2(m_T)`$ (or $`v_2(p_T)`$)
$`v_2(m_T)={\displaystyle \frac{𝑑\varphi _p\mathrm{cos}2\varphi _p\frac{dN}{m_Tdm_Td\varphi _p}}{𝑑\varphi _p\frac{dN}{m_Tdm_Td\varphi _p}}}.`$ (23)
Thus we obtain the same equation as Eq. (33) in Ref. Pratt:2004zq
$`v_2(m_T)`$ $`=`$ $`{\displaystyle \frac{\epsilon }{J_0}}\left({\displaystyle \frac{m_Tu_{}^2}{T\sqrt{1+u_{}^2}}}J_E+{\displaystyle \frac{p_Tu_{}}{T}}J_p\right),`$ (24)
$`J_0`$ $`=`$ $`2K_1I_0,`$ (25)
$`J_E`$ $`=`$ $`\left(K_0+{\displaystyle \frac{K_1}{z_E}}\right)I_2,`$ (26)
$`J_p`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(I_1+I_3\right)K_1.`$ (27)
Here $`K_i`$ and $`I_i`$ are modified Bessel functions. It is always understood that the argument of $`K_i`$ ($`I_i`$) is $`z_E=m_T\sqrt{1+u_{}^2}/T`$ ($`z_p=p_Tu_{}/T`$). Detailed calculations can be found in Appendix A.
One can also obtain an analytic expression for the slope of $`v_2(p_T)`$
$`{\displaystyle \frac{dv_2}{dp_T}}`$ $`=`$ $`{\displaystyle \frac{\epsilon }{J_0^2}}\{{\displaystyle \frac{u_{}^2}{T\sqrt{1+u_{}^2}}}[{\displaystyle \frac{p_T}{m_T}}J_E`$
$``$ $`m_T({\displaystyle \frac{z_E}{p_T}}{\displaystyle \frac{J_E}{z_E}}+{\displaystyle \frac{z_p}{p_T}}{\displaystyle \frac{J_E}{z_p}})]`$
$`+`$ $`{\displaystyle \frac{u_{}}{T}}J_p+{\displaystyle \frac{p_Tu_{}}{T}}({\displaystyle \frac{z_E}{p_T}}{\displaystyle \frac{J_E}{z_E}}+{\displaystyle \frac{z_p}{p_T}}{\displaystyle \frac{J_p}{z_p}})\}.`$
Here,
$`{\displaystyle \frac{z_E}{p_T}}`$ $`=`$ $`{\displaystyle \frac{\sqrt{1+u_{}^2}}{T}}{\displaystyle \frac{p_T}{m_T}},`$ (29)
$`{\displaystyle \frac{z_p}{p_T}}`$ $`=`$ $`{\displaystyle \frac{u_{}}{T}},`$ (30)
$`{\displaystyle \frac{J_E}{z_E}}`$ $`=`$ $`\left(K_1+{\displaystyle \frac{K_2}{z_E}}\right)I_2,`$ (31)
$`{\displaystyle \frac{J_E}{z_p}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(K_0+{\displaystyle \frac{K_1}{z_E}}\right)(I_1+I_3),`$ (32)
$`{\displaystyle \frac{J_p}{z_E}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(K_0+{\displaystyle \frac{K_1}{z_E}}\right)(I_1+I_3),`$ (33)
$`{\displaystyle \frac{J_p}{z_p}}`$ $`=`$ $`{\displaystyle \frac{1}{4}}K_1(I_0+2I_2+I_4).`$ (34)
Note that one can replace a higher order modified Bessel function with lower order functions.
### VI.3 Incorporation of Transverse Dynamics
From discussion in Sec. IV, we find $`dm_T/dT<0`$ (or $`dm_T/d\beta >0`$) for chemical equilibrium pions and $`dm_T/dT>0`$ (or $`dm_T/d\beta <0`$) for chemically frozen pions. These features are quite generic for ideal Bjorken fluids of pions. Obviously, the analytic approach is just a parametrization at freezeout and contains almost no information about the time evolution of the system. For example, the analytic approach does not tell us anything about how $`u_{}`$ increases with decrease of temperature. In this subsection, we try to give a dynamical meaning to the blast wave approach discussed in the previous subsection.
The transverse mass distribution in the analytic approach is (see Eq. (25)) Schnedermann:1992hp
$`{\displaystyle \frac{dN}{m_Tdm_T}}`$ $``$ $`m_TK_1I_0.`$ (35)
Thus one obtain the mean transverse mass
$`m_T={\displaystyle \frac{𝑑m_Tm_T^3K_1I_0}{𝑑m_Tm_T^2K_1I_0}}`$ (36)
and its derivative with respect to the inverse temperature
$`{\displaystyle \frac{dm_T}{d\beta }}`$ $`=`$ $`{\displaystyle \frac{1}{\left(𝑑m_Tm_T^2K_1I_0\right)^2}}`$ (37)
$`\times \{\left[{\displaystyle }dm_1m_1^3({\displaystyle \frac{dK_1}{d\beta }}I_0+K_1{\displaystyle \frac{dI_0}{d\beta }})\right]`$
$`\times \left({\displaystyle 𝑑m_2m_2^2K_1I_0}\right)\left({\displaystyle 𝑑m_1m_1^3K_1I_0}\right)`$
$`\times \left[{\displaystyle }dm_2m_2^2({\displaystyle \frac{dK_1}{d\beta }}I_0+K_1{\displaystyle \frac{dI_0}{d\beta }})\right]\},`$
$`{\displaystyle \frac{dK_1}{d\beta }}`$ $`=`$ $`{\displaystyle \frac{K_1}{\beta }}+{\displaystyle \frac{d\rho }{d\beta }}{\displaystyle \frac{dK_1}{d\rho }}`$ (38)
$`=`$ $`m_T{\displaystyle \frac{dK_1(z_E)}{dz_E}}\left(\mathrm{cosh}\rho +{\displaystyle \frac{d\rho }{d\beta }}\beta \mathrm{sinh}\rho \right)`$
$`=`$ $`m_T\left(K_0+{\displaystyle \frac{K_1}{z_E}}\right)`$
$`\times \left(\mathrm{cosh}\rho +{\displaystyle \frac{d\rho }{d\beta }}\beta \mathrm{sinh}\rho \right),`$
$`{\displaystyle \frac{dI_0}{d\beta }}`$ $`=`$ $`p_T{\displaystyle \frac{dI_0(z_p)}{dz_p}}\left(\mathrm{sinh}\rho +{\displaystyle \frac{d\rho }{d\beta }}\beta \mathrm{cosh}\rho \right)`$ (39)
$`=`$ $`p_TI_1\left(\mathrm{sinh}\rho +{\displaystyle \frac{d\rho }{d\beta }}\beta \mathrm{cosh}\rho \right).`$
The numerator of Eq. (37) reduces to
$`[\mathrm{numerator}\mathrm{of}\mathrm{Eq}.(\text{37})]`$
$`=`$ $`{\displaystyle 𝑑m_1𝑑m_2m_1^2m_2^2(m_2m_1)(K_1I_0)}_{m_T=m_1}`$
$`\times `$ $`[(m_T\mathrm{cosh}\rho +m_T\mathrm{sinh}\rho {\displaystyle \frac{d\rho }{d\beta }}\beta )K_0I_0`$
$`+\left({\displaystyle \frac{1}{\beta }}+{\displaystyle \frac{d\rho }{d\beta }}\mathrm{tanh}\rho \right)K_1I_0`$
$`(p_T\mathrm{sinh}\rho +p_T\mathrm{cosh}\rho {\displaystyle \frac{d\rho }{d\beta }}\beta )K_1I_1\left]\right|_{m_T=m_2}.`$
The second term in the square bracket vanishes due to the antisymmetry ($`m_1m_2`$). Thus,
$`[\mathrm{numerator}\mathrm{of}\mathrm{Eq}.(\text{37})]`$
$`=`$ $`{\displaystyle 𝑑m_1𝑑m_2m_1^2m_2^2(m_2m_1)K_1(z_{E,1})I_0(z_{p,1})}`$
$`\times `$ $`[m_2(\mathrm{cosh}\rho +\mathrm{sinh}\rho {\displaystyle \frac{d\rho }{d\beta }}\beta )K_0(z_{E,2})I_0(z_{p,2})`$
$``$ $`\sqrt{m_2^2m^2}(\mathrm{sinh}\rho +\mathrm{cosh}\rho {\displaystyle \frac{d\rho }{d\beta }}\beta )K_1(z_{E,2})I_1(z_{p,2})],`$
where $`z_{E,i}=z_E_{m_T=m_i}`$ and $`z_{p,i}=z_p_{p_T=\sqrt{m_i^2m^2}}`$. The sign of $`dm_T/d\beta `$ is determined by this equation. One can obtain the relation between $`\rho `$ and $`\beta `$ by solving an equation $`dm_T/d\beta =0`$:
$`\beta {\displaystyle \frac{d\rho }{d\beta }}[F(\rho )\mathrm{sinh}\rho G(\rho )\mathrm{cosh}\rho ]`$ (42)
$`+`$ $`[F(\rho )\mathrm{cosh}\rho G(\rho )\mathrm{sinh}\rho ]=0,`$
where,
$`F(\rho )`$ $`=`$ $`{\displaystyle 𝑑m_1𝑑m_2(m_2m_1)m_1^2m_2^3}`$ (43)
$`\times `$ $`K_1(z_{E,1})I_0(z_{p,1})K_0(z_{E,2})I_0(z_{p,2}),`$
$`G(\rho )`$ $`=`$ $`{\displaystyle 𝑑m_1𝑑m_2(m_2m_1)m_1^2m_2^2\sqrt{m_2^2m^2}}`$ (44)
$`\times `$ $`K_1(z_{E,1})I_0(z_{p,1})K_1(z_{E,2})I_1(z_{p,2}).`$
The temperature dependence of transverse rapidity $`\rho =\rho (\beta )`$ is obtained for a given “initial” condition $`(\beta _0,\rho _0)`$. This particular radial flow ensures the mean transverse mass becomes a constant and is an upper limit of average radial flow in chemical freezeout EOS for massive pions. We call this solution the critical radial flow $`v_r^{\mathrm{crit}}=\mathrm{tanh}\rho (\beta )`$. One can parametrize the temperature dependence of radial flow by introducing a parameter $`\alpha `$ within the analytic approach which embodies the transverse dynamics of the chemically frozen/equilibrated pion fluid:
$`v_r(T)=v_r(T^{\mathrm{ch}})+\alpha [v_r^{\mathrm{crit}}(T)v_r(T^{\mathrm{ch}})],`$ (45)
where $`v_r(T^{\mathrm{ch}})=\mathrm{tanh}\rho _0(\beta _0)`$ is an initial condition for Eq. (42). Although the exact value of $`\alpha `$ needs much more involved dynamical calculation, radial flow qualitatively corresponds to the chemically frozen system for $`0<\alpha <1`$ and to the chemical equilibrium system for $`\alpha >1`$. It should be mentioned here that the temperature dependence of average radial flow can be described to some extent without solving hydrodynamic equations. Note that $`\alpha `$ should be taken as being a moderate value so that the total energy of the system (per unit rapidity) does not increase due to generation of radial flow.
Figure 7 shows temperature dependences of radial flow. The solid line shows the critical radial flow $`v_r^{\mathrm{crit}}`$ which results from $`dm_T/dT=0`$. This is obtained by solving Eq. (42) with an initial condition $`(T,v_r)=(170,0.25)`$ which is consistent with a value at RHIC energies Hirano:2002ds . The critical flow distinguishes the system of chemical equilibrium from that of chemical freezeout: $`v_r(T)`$ for CE (CFO) should be located above (below) this line since $`dm_T/dT<0`$ for CE ($`dm_T/dT>0`$ for CFO). Dashed line ($`\alpha =1.2`$) shows an example of radial flow in the analytic model CE, while dotted line ($`\alpha =0.6`$) shows the one in the analytic model CFO. These results look very similar to the results from real hydrodynamic simulations as shown in Fig. 5 in Ref. Hirano:2002ds .
By using these profiles for radial flow, we calculate $`v_2(p_T)`$ for pions below the chemical freezeout temperature. Hydrodynamic analysis tells us that the integrated $`v_2`$ is saturated within first 3–4 fm/$`c`$ just after the collision and insensitive to the late hadronic stage Kolb:2003dz . However, within our analytic approach, there is no dynamical mechanism which saturates the integrated $`v_2`$ in the late hadron stage. Therefore $`v_2(p_T)/v_2`$ is the quantity to be compared with the results from full hydrodynamic simulations.
In Fig. 8, $`v_2(p_T)/v_2`$ for the analytic models CE and CFO are represented. The thick solid line shows the result at $`T=T^{\mathrm{ch}}=170`$ MeV. The overall slope up to 1 GeV/$`c`$ is gradually increasing with decreasing the temperature in the analytic model CFO, while the slope is decreasing in the analytic model CE. These results clearly show that the slope of $`v_2(p_T)`$ can vary in the late hadronic stage and that the temporal behavior of average transverse momentum $`p_T`$ and radial flow $`v_r`$ is the key to understand the shape of $`v_2(p_T)`$.
## VII Conclusion and Outlook
In this paper we showed that the differential elliptic flow observable $`v_2(p_T)`$, which is a critical component for the interpretation of RHIC data in terms of perfect fluidity of the sQGP core, is sensitive to the degree of hadro-chemical equilibrium in late time evolution of the hadronic corona. If local equilibrium hydrodynamics is applied to the hadronic corona below $`T_c`$, an inevitable logical impasse arises when confronting all the data on (1) hadron abundances (2) radial flow and (3) differential elliptic flow. In CE hydrodynamics (2) and (3) can be reproduced at the expense of (1). In PCE (1) is enforced at the expense of (2) and (3) as summarized in Table I. We presented a simple analytic blast wave model to explain these nonintuitive consequences of hadro-chemical (non)equilibrium in (P)CE implementations of hydrodynamics.
In CE both the average transverse momentum per hadron $`p_T`$ and average $`v_2`$ increase with proper time in the hadronic phase in a way that accidentally preserves the slope of differential elliptic flow $`dv_2(p_T)/dp_Tv_2/p_T`$ in agreement with the data. In PCE, $`p_T`$ decreases due to the basic Bjorken longitudinal cooling. The main point is that in PCE the hadronic yields are fixed at $`T^{\mathrm{ch}}`$ and the compensating CE “local reheating” mechanism (the conversion of heavy resonance mass back into internal energy which “mimics” a sort of dissipative effect) is absent. This is why PCE fails to describe the proton radial flow data. In addition, the slight increase of the average $`v_2`$, as in CE, with proper time cannot be compensated for in PCE. Therefore, the slope of differential elliptic flow $`dv_2(p_T)/dp_Tv_2/p_T`$ continues to grow in PCE during the hadronic phase, which leads to disagreement with RHIC data.
The subtle interplay among (1) longitudinal expansion work, (2) maintenance of hadronic abundance yields, (3) the long time development of radial flow, and (4) the differential azimuthal asymmetric elliptic flow provides a formidable dynamical constraint on the dynamics of the hadronic corona. Only by abandoning ideal hydrodynamics in the hadronic corona, have nonequilibrium hadron cascade (HC) models been able to deal with the interplay of the above hadron dynamics in a way consistent with present RHIC data. As discussed in Sec. II, this approach is natural since the viscosity to entropy ratio in a hadronic resonance gas below $`T_c`$ is too large to support even local thermal equilibrium. By relaxing both thermal and hadro-chemical equilibrium assumptions, the hybrid QGP hydrodynamics plus hadron cascade model in Teaney:2000cw has been able to account for all three major low $`p_T`$ observables as summarized in Table I. The effect of viscosity in the hadron phase Teaney:2003pb substitutes for the “local reheating” in the CE model and compensates the small growth of the average $`v_2`$ in PCE to preserve the slope of $`v_2(p_T)`$ for pions. The slope of $`v_2(p_T)`$ is thus found to stall at the SPS energy from the hybrid model analysis Teaney:2002aj . In the classical transport approach, both $`p_T`$ Gyulassy:1997ib and $`v_2`$ Zhang:1999rs do not vary significantly when the mean free path among the particles becomes comparable with the typical gradient length scales. Moreover, the shear viscous effect changes the momentum distribution function Dumitru:2002sq and reduces the slope of $`v_2(p_T)`$ slightly Teaney:2003pb . These are the reasons why the slope of $`v_2(p_T)`$ does not changed significantly in the hadronic transport stage.
So how robust is the statement that hydrodynamic description at RHIC works remarkably well? We emphasized that the behavior of $`v_2`$ differs from that of $`v_2(p_T)`$: The integrated elliptic flow does not develop so much in the late hadronic stage in which either the inviscid, chemical (non-)equilibrium fluid or the dissipative gas is assumed, whereas the differential elliptic flow depends largely on these assumptions. The large magnitude of integrated $`v_2`$ observed at RHIC is reproduced only when a small $`\eta /s`$ is assumed Molnar:2001ux . Therefore the large $`v_2`$ developed in the early stage is obtained from the evolution of the sQGP core, which as discussed in Sec. II must have near minimal viscosity $`\eta _{\mathrm{SYM}}T^3`$. Even though the minimal sQGP viscosity is larger than the viscosity of the HRG corona, the core exhibits near perfect fluid behavior due to the deconfinement of almost all the QCD degrees of freedom. The near perfect fluidity of the sQGP core is therefore a signal of deconfinement. On the other hand, the breaking of local and hadro-chemical equilibrium in the hadronic corona is critical for this interpretation of RHIC data. If inviscid ideal hydrodynamics were valid in both sQGP and HRG phases, the crucial $`v_2(p_T)`$ would be sensitive to the hadronic thermal freezeout dynamics and not only to the equation of state of sQGP matter.
Perhaps most surprising in connection with the hydrodynamics robustness question is the important role played by hadro-chemical freezeout at $`T^{\mathrm{ch}}T_c`$ that is implied by the extensive systematics of observed hadron abundances Braun-Munzinger:2003zd . Without this constraint the different hadro-chemical results with CE and PCE would preclude a conclusion about the perfect fluidity of the sQGP core as well as the highly dissipative, imperfect fluidity of its hadronic resonance corona.
Despite the success of the hybrid HC approach, there exist open technical questions that must be still investigated. One important issue is the violation of energy-momentum conservations at the boundary between the QGP and hadron phases Bugaev:2002ch . The Cooper-Frye prescription Cooper:1974mv is employed to obtain the particle distribution just after hadronization which is to be used as an initial condition in the sequential cascade calculation. The violation is expected to be small when radial flow is large. Nevertheless there always exists in the space-like hypersurface $`d\sigma ^\mu `$ in-coming particles which contributes to the number of particles negatively. This negative contribution is omitted in the actual calculations, which causes the violation of energy-momentum conservation. Proper treatment of the boundary condition may lead to change the dynamics in the QGP phase Bugaev:2002ch . In this connection, the approximate continuity of the viscosity from the sQGP to the HRG phase discussed in Sec. II minimizes this interface problem since there is no discontinuity of the stress tensor including the viscous correction at $`T_c`$.
Another important future problem is the rapidity dependence of elliptic flow. The (3+1)-dimensional ideal hydrodynamic calculations Hirano:2001eu ; Hirano:2002ds have not been able to reproduce the observed pseudorapidity dependence of $`v_2`$ in forward rapidity region at RHIC Back:2002gz . The forward rapidity region at RHIC is similar to the midrapidity region at SPS in the sense that local particle density $`(1/S)dN_{\mathrm{ch}}/dy`$ is similar. Heinz and Kolb Heinz:2004et proposed a “thermalization coefficient” to describe enhanced nonequilibrium effects in the low particle density region defined from the experimental data $`v_2/\epsilon `$ as a function of $`(1/S)dN_{\mathrm{ch}}/dy`$ Alt:2003ab . To address correctly the rapidity and beam energy dependence taking into account the highly viscous nature of the hadronic corona, a new hybrid model must be developed in which the (3+1)-D hydrodynamic model for the sQGP core is combined with the HRG transport approach for the dissipative hadronic corona. Hirano and Nara Hirano:2004rs ; Hirano:2003pw ; Hirano:2002sc have already developed a dynamical model to describe three different aspects of relativistic heavy ion collisions in one consistent framework: Color glass condensate initial conditions for high energy colliding nuclei, hydrodynamic evolution in 3D space for long wave length components of produced matter, and quenching jets for short wave length components. A further unified framework by combining a hadronic cascade model with the above one HiranoNara will be necessary to understand more quantitatively the dynamics and the properties of QCD matter produced in relativistic heavy ion collisions at all SPS, RHIC and LHC energies.
We emphasize that continued work toward such a unified dynamical framework will be essential to further test our physical interpretation of RHIC data - as outlined in the introduction and analyzed in sections II-VI - that “perfect fluidity” of the higher viscosity sQGP core and “imperfect fluidity” of the lower viscosity HRG corona taken together with hadro-chemical equilibrium near $`T_c160170`$ MeV already provide a compelling set of signatures for QCD deconfinement at RHIC.
###### Acknowledgements.
We thank C. Ogilvie for discussions on compiled hydrodynamic results on transverse momentum dependences of spectra and elliptic flow in the PHENIX white paper. Valuable discussions with K. Bugaev, P. Huovinen, D. Molnar, K. Schalm and D. Teaney are also acknowledged. We also thank P. Huovinen for cereful reading of the manuscript. This work was supported in part by the United States Department of Energy under Grant No. DE-FG02-93ER40764.
## Appendix A Derivation of Eq. (24)
Let us recall the formulae for modified Bessel functions
$`2K_1(z)`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑x\mathrm{cosh}x\mathrm{exp}(z\mathrm{cosh}x),`$ (46)
$`2\pi I_2(z)`$ $`=`$ $`{\displaystyle _\pi ^\pi }𝑑\varphi \mathrm{cos}2\varphi \mathrm{exp}(z\mathrm{cos}\varphi ).`$ (47)
The numerator of Eq. (23) becomes
$`{\displaystyle 𝑑\varphi _p\mathrm{cos}2\varphi _p𝑑\varphi 𝑑y\mathrm{cosh}y\mathrm{exp}(A)}`$ (48)
$`=`$ $`{\displaystyle 𝑑\varphi 2K_1(B)\times 2\pi \mathrm{cos}2\varphi I_2(C)},`$
where,
$`B={\displaystyle \frac{m_T\sqrt{1+u_{}^2}}{T}}+\epsilon {\displaystyle \frac{m_Tu_{}^2\mathrm{cos}2\varphi }{T\sqrt{1+u_{}^2}}},`$
$`C={\displaystyle \frac{p_Tu_{}}{T}}+\epsilon {\displaystyle \frac{p_Tu_{}\mathrm{cos}2\varphi }{T}}.`$
We here assume $`\epsilon `$ is small, expand the modified Bessel function, and take the first order term with respect to $`\epsilon `$,
$`4\pi {\displaystyle 𝑑\varphi \left(K_1+\epsilon K_1^{}\frac{m_Tu_{}^2\mathrm{cos}2\varphi }{T\sqrt{1+u_{}^2}}\right)}`$ (49)
$`\times `$ $`\left(I_2+\epsilon I_2^{}{\displaystyle \frac{p_Tu_{}\mathrm{cos}2\varphi }{T}}\right)`$
$``$ $`4\pi ^2\epsilon \left(K_1^{}I_2{\displaystyle \frac{m_Tu_{}^2}{T\sqrt{1+u_{}^2}}}+K_1I_2^{}{\displaystyle \frac{p_Tu_{}}{T}}\right).`$
Let us also recall some useful formulae for the derivatives of modified Bessel functions,
$`K_n^{}(z)`$ $`=`$ $`K_{n1}(z){\displaystyle \frac{n}{z}}K_n(z),`$ (50)
$`I_n^{}(z)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[I_{n1}(z)+I_{n+1}(z)\right].`$ (51)
Then the numerator of Eq. (23) is proportional to
$`4\pi ^2\epsilon `$ $`[(K_0+{\displaystyle \frac{K_1}{z}})I_2{\displaystyle \frac{m_Tu_{}^2}{T\sqrt{1+u_{}^2}}}`$ (52)
$`+{\displaystyle \frac{1}{2}}(I_1+I_3)K_1{\displaystyle \frac{p_Tu_{}}{T}}].`$
Here the arguments of $`K`$ and $`I`$ are, respectively, $`z_E=m_T\sqrt{1+u_{}^2}/T`$ and $`z_p=p_Tu_{}/T`$. An analogous calculation can be done for the denominator of Eq. (23). The result is proportional to $`4\pi ^2\times 2K_1I_0`$. |
warning/0506/astro-ph0506149.html | ar5iv | text | # The near-UV pulse profile and spectrum of the pulsar PSR B0656+14Based on observations with the NASA/ESA Hubble Space Telescope, obtained at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc. under NASA contract NAS5-26555.
## 1 Introduction
The middle-aged pulsar PSR B0656+14 (PSR J0659+14
14) was first discovered in the radio by Manchester et al. (1978). The pulsar period, $``$0.385 s, spin-down age, $``$1.1$`\times 10^5`$ years, and several other pulsar parameters are given in Table 1. This is one of the brightest isolated neutron stars (NSs) in the X-ray sky (Córdova et al. 1989), and also the optical counterpart is brighter ($`V25`$) than most other middle-aged pulsars (Caraveo et al. 1994). This made it possible to study the optical pulsations of PSR B0656+14 (Shearer et al. 1997; Kern et al. 2003), which has previously only been convincingly done for a few much younger pulsars.
Owing to its relative brightness and proximity, PSR B
0656+14 is one of the most intensively studied isolated NSs in different wavelength bands. The combined phase integrated spectrum of the pulsar, from radio to gamma-rays (cf. Koptsevich et al. 2001), shows that its radiation consists of two parts. One component is non-thermal emission from the pulsar magnetosphere described by a power-law spectrum with different spectral indices in different spectral ranges. The other component is thermal blackbody like emission from the whole surface of the cooling NS and from much smaller and significantly hotter polar caps of the pulsar.
While the thermal emission from the whole surface dominates in the extreme UV (EUV) and soft X-rays (e.g., Edelstein 2000; Greiveldinger et al. 1996; Zavlin & Pavlov 2004; De Luca et al. 2005), the emission at longer wavelengths is mainly of non-thermal origin and has a negative spectral slope (Pavlov et al. 1996; Pavlov et al. 1997; Koptsevich et al. 2001; Komarova et al. 2003). Ground-based and HST UV to near-IR photometry suggest that a main spectral change from the near-IR to EUV occurs somewhere between the B band and the far-UV (FUV).
Here we report on time-resolved HST STIS NUV-MAMA prism spectroscopy of PSR B0656+14 in the near-UV (NUV). This is the first time the pulse profile of PSR B0656+14 has been revealed in the UV. The detection of UV pulsations from another middle-aged pulsar, Geminga, and the younger Vela-pulsar was very recently reported by Kargaltsev et al. (2005) and Romani et al. (2005), but was previously only obtained for the Crab pulsar (NUV, Gull et al. 1998; FUV, Sollerman et al. 2000). In this paper we study the properties of the NUV emission of PSR B0656+14 and compare them with available data in other wavelength ranges. We describe our observations and data reduction in Sect. 2. The results are discussed in Sect. 3, and we then draw our conclusions in Sect. 4.
## 2 Observations and data analysis
### 2.1 Observations
Time-resolved spectroscopy of PSR B0656+14 in the near-UV ($`17003400`$ Å) range was performed with the HST/STIS during two visits in fall 2001 using the UV PRISM, the 52″$`\times `$0$`.^{\prime \prime }`$5 slit and the NUV MAMA detector in TIMETAG recording mode. The PRISM mode was used because of the very low pulsar flux. The relative time resolution was 125$`\mu `$s. Accurate target coordinates were obtained from earlier WFPC2 images. The target acquisition was performed using a B=17.9 field star, U0975-04374987, which is positioned 13$`.^{\prime \prime }`$9 away from the pulsar (Fig. 1). We requested and received two visits, with the first being on 1 September and the second on 16 November. The slit position angle was PA=$`135`$°.3 east of north during the first visit and $`122`$°.45 during the second visit. These position angles allowed us to avoid any contamination of the pulsar flux by a faint extended background object only $``$1″ north of our target (Koptsevich et al. 2001; Komarova et al. 2003) as well as to exclude other nearby objects (Fig. 1). The two visits also allowed initial data evaluation and checks to ensure proper astrometry had been used. Details of the individual science observations are listed in Table 2.
The initial visit, three orbits in length, began with normal acquisition plus three STIS/CCD clear aperture images, each with 120 s exposure time, for a control image of the pulsar position. Figure 1 shows the composite field image with the 52″$`\times `$0$`.^{\prime \prime }`$5 aperture boundaries inscribed. An internal lamp-illuminated image of the 52″$`\times `$0$`.^{\prime \prime }`$5 slit was then recorded, and then the pulsar was viewed through the same aperture by the CCD to confirm that the pulsar was well-centered in the slit. This also precisely determines the pulsar position for the prism dispersion wavelength reference. The remainder of the first visit was allocated to three science exposures with the PRISM and NUV MAMA detector (Table 2). Internal WAVECAL exposures were interspersed between science exposures to ensure good reference for wavelength scale. The total UV PRISM science exposure time in the first visit was 6790 s.
The second visit was five orbits long. CCD images of the pulsar viewed through the slit were performed at the beginning of the first and the third orbits to check for any target or image drift internal to the telescope, but no such drift was detected. Five TIMETAG exposures, one during each orbit, were recorded during 12760 s (Table 2). The total TIMETAG science exposure time for the two visits was thus 19.55 ks. Additional information can be obtained through the HST archive<sup>1</sup><sup>1</sup>1archive.stsci.edu under GO program 9156<sup>2</sup><sup>2</sup>2presto.stsci.edu/observing/phase2-public/9156.pro.
### 2.2 Data reduction
#### 2.2.1 Extraction of the Spectrum
All science exposures for each visit were co-added. The spectrum was then extracted and calibrated using the STIS CALSTIS pipeline tasks with a 7 pixel wide extraction slit (0$`.^{\prime \prime }`$013/pix). The backgrounds were extracted with a 7 pixels wide slit centered 20 pixels above and below the center of the pulsar spectrum. The associated wavecals recorded between science exposures were used to adjust the zeropoint of the wavelength scale. The two extracted spectra were then registered and averaged with weights scaled to the total exposure times for each visit. As the PRISM dispersion increases with decreasing wavelength and the detector sensitivity falls, we binned the flux below 2000 Å in 4 pixel bins. Between 2000 Å and 2300 Å, the binning was for 2 pixels and above 2300 Å, the data were left at the original PRISM resolution. This resulted in bins from $`20`$ Å in the blue to $`60`$ Å towards the red. The detector sensitivity and prism dispersion cut the sensitivity considerably below 2000 Å, but in the range $`20003200`$ Å the pulsar is detected with S/N $`3`$ per bin. The NUV MAMA detector is the dominating source of the background and the total sky+detector background was $``$700 and $``$3600 counts for the 1st and 2nd visits, while the signal+background was $``$1200 and $``$5000 counts, respectively.
#### 2.2.2 Generation of the Pulse Profiles
The pulse profile analysis was done with STIS Instrument Development Team tools designed originally for testing the TIMETAG mode of the MAMA detectors (Gull et al. 1998). These tools are very similar in function to the STSDAS/STIS routines. Both sets of software have been used on the Crab pulsar data by Gull et al. (1998) and by Sollerman et al. (2000) and provide very similar results. The pulse profile was determined for each visit by concatenating the time-tagged events for each PRISM science exposure. Event times were corrected for the Earth’s motion around the Sun and for HST’s motion around the Earth using the definitive orbital ephemeris extracted from the HST archive. The orbital position was tabulated at one minute intervals and was determined by cubic-spline interpolation for each event.
The pulse period was determined from the data stream by the following process. All events within a rectangular patch centered on the pulsar spectrum were selected. Likewise, we selected all events in background patches above and below the pulsar spectrum. The pulse period was determined by maximizing the RMS of the pulse profile of the target events (see Sollerman et al. 2000). For the two visits, the computed periods are:
Visit 1: 0.384902 s
Visit 2: 0.384903 s
This is in excellent agreement to within the last digit with the radio pulsar period (Table 1). Using the ephemeris from Kern et al. (2003) we obtain:
Visit 1: 0.3849028 s
Visit 2: 0.3849031 s
at the epochs of our observations. The measaured period is also in agreement with earlier radio observations presented by Taylor et al. (1993) and Chang & Ho (1999), and there is thus no evidence for glitches in PSR B0656+14 over the last 16 years. This also ensures us that the obtained pulse profile is correct. To increase the S/N, the data were binned into phase bins at intervals 1/16th that of the pulse period. The background data were likewise binned, but showed no structure with phase. We then subtracted the combined background average bin count. The TIME-TAG mode provided relative time resolution of 125 $`\mu `$s, but absolute time determination is not known to better than one second due to design limitations of the STIS computer timing updates. Internal clock stability is sufficient to ensure good time stability from orbit to orbit, even up to several days apart, but not for a ten week interval. We were thus forced to use the pulse profile shapes in the NUV and in the optical to phase align the two visits and to determine the NUV pulse position with respect to the radio pulse (see Sect. 3.1 and 3.2).
## 3 Results and Discussion
### 3.1 The NUV pulse profile
The background subtracted spectral integrated pulse profiles obtained from the data of visits 1 and 2 are shown in the top and middle panels of Fig. 2, respectively. While the pulsations in the first visit are only significant at the $`4\sigma `$ level, the significance of the pulsations in the second, longer visit is $`8\sigma `$. The pulse profile contains two sub-pulses whose maxima are separated by $`0.5`$ in phase
and are connected by an emission “bridge”. The shape of the first sub-pulse is rather symmetric, while the second sub-pulse is characterized by a linear flux increase in its leading part as well as by a sharp edge in the trailing part. This sharp drop and the bridge can also be seen in the first visit profile. We have used these features to align both pulse profiles in phase and combined them. The result is shown in the bottom panel of Fig. 2. The significance of the pulsations increased in the combined profile to $`9.2\sigma `$.
We cannot exclude the presence of a persistent (or a weakly variable) flux component from the pulsar in the NUV range. As shown in Fig. 2, the estimated pulse fraction Pf (the ratio of the number of counts above the minimum of the profile to the total number of counts) is in the range $`58\%82\%`$ representing a $`1\sigma `$ confidence interval based on count statistics. The persistent flux contribution into the pulse integrated flux is thus within the $`1842\%`$ range. The lower boundary of this uncertainty range is compatible with the published upper limit $`16`$% in the optical range (Kern et al. 2003).
### 3.2 Pulse profiles from the radio to gamma-rays
The shape of the observed NUV double pulse profile is very similar to what is seen in the optical range by Kern et al. (2003). This supports a single emission mechanism being responsible for the emission in both ranges. We note that the Kern et al. (2003) pulse profile was actually very different from the optical pulse profile for PSR B0656+14 earlier reported by Shearer et al. (1997). Kern et al. (2003) discussed whether this large discrepancy could be due to differences in pass bands, or due to a pulse profile variation with time. Our NUV profile clearly supports the findings of Kern et al. (2003) and disfavours the explanations to accomodate the earlier profile by Shearer et al. (1997).
The optical pulse profile by Kern et al. (2003) was obtained with absolute timing which allowed a phase alignment between the optical and the radio pulse profiles. The same was done for X-ray data obtained with the Chandra and XMM-Newton observatories (Pavlov et al. 2002; Zavlin & Pavlov 2004; De Luca et al. 2005), and for gamma-ray data obtained with EGRET (Ramanamurthy et al. 1996). We have used the similarity of the optical and NUV profiles to align our data with the radio pulse and to compare its absolute position in phase and morphology with the available data in a wide spectral range, from radio through gamma-rays (Fig. 3).
In the radio (Gould & Lyne 1998) the pulsar has a sharp single pulse with a width of about 0.2 of the NS rotational period (Fig. 3a). In the soft X-rays it instead displays shallow (Pf$`1024`$%) pulsations with a sine-like light curve. It has a single broad maximum per pulsar period, typical for the thermal emission from the surface of a NS with non-uniform temperature distribution over the surface (Fig. 3d). This is completely different from the optical and NUV profiles (Figs. 3b–c) and underline the difference in the emission mechanisms dominating in the optical-NUV (non-thermal) and in the soft X-rays (thermal). However, in the smooth light curves in soft X-rays one can discover small, but significant, features approximately coinciding in phase with the positions of the NUV sub-pulses (Figs. 3c and 3d–g). This can be more easily seen in the zoomed examples of the X-ray profiles (Fig. 4, see also Pavlov et al. 2002 and Zavlin & Pavlov 2004). With increasing photon energy, where the thermal emission is assumed to be dominated by the pulsars hot polar caps, the maximum of the X-ray profile is shifted toward the position of the 2nd NUV sub-pulse (Fig. 3g). However, this emission can not contribute significantly to the NUV because the hot spot area is too small ($`R_{HS}2`$ km, De Luca et al. 2005). In the hard X-rays (Fig. 3h), where the spectrum is dominated by a non-thermal magnetospheric component, both the position and the morphology of the profile is similar to that of the 2nd NUV sub-pulse. The counterpart of the 1st NUV sub-pulse is not observed at the low S/N available
in the hard X-ray range, while it coincides in position with the gamma-ray pulse marginally detected with the EGRET (Ramanamurthy et al. 1996). The pulsar has not yet been reliably detected in the intermediate energy range between the hard X-rays and gamma-rays. Available RXTE observations obtained in the 15-250 keV range provide only an upper limit on the pulsar flux and a very uncertain pulse profile (Chang & Ho 1999).
To summarize, the observed profiles suggest that the pulsed emission outside the radio range is a sum of at least two components. One component, likely non-thermal emission, displays a two peak pulse profile which is clearly resolved in the optical and in the NUV. The other component, likely thermal emission, has a shallow broad sine-like profile with a single pulse per period. This component could potentially be partially responsible for the persistent and bridge emission seen in the NUV, since its maximum coincides with the NUV inter-pulse position. The counterparts of the optical-NUV sub-pulses can be resolved from the thermal background in soft X-rays. The 2nd NUV sub-pulse counterpart is clearly resolved in hard X-rays, while the 1st sub-pulse may be seen in gamma-rays. The coincidence of the position and morphology of the 2nd NUV sub-pulse and the X-ray profile supports a strong correlation of the non-thermal emission in these ranges, which has also been deduced from phase averaged spectral analysis (e.g., Koptsevich et al. 2001). The radio pulse is centered at the off-pulse emission in all other ranges.
#### 3.2.1 Thermal nature of the bridge and persistent emission?
The variation of the thermal flux with the NS rotation may be understood in terms of anisotropic heat conduction through the strongly magnetized envelope of a cooling NS. The regions of the star with radial magnetic fields are hotter than the regions with tangential fields (e.g., Greenstein & Hartke 1983; Page 1995; Shibanov & Yakovlev 1996; Potekhin & Yakovlev 2001). The observed flux modulation in the soft X-rays is in agreement with a $``$14% modulation expected from a middle-aged cooling NS with a dipole magnetic field of a few times $`10^{12}`$ G (e.g., Potekhin & Yakovlev 2001). A single maximum per period with a relatively symmetric pulse profile suggests that only one magnetic pole is visible in the X-rays during the pulsar rotation. This is also compatible with the pulsar viewing geometry derived from the radio (Everett & Weisberg 2001) and optical (Kern et al. 2003) polarization measurements. However, the $`0.5`$ phase lag between the maxima in the radio and soft X-rays suggests that the hollow polar cap magnetospheric cone assumed to be responsible for the radio and optical emission must be strongly curved with height. Alternatively, the magnetic field surface structure is different from a simple centered dipole configuration. The latter is supported by the significant phase lag (De Luca et al. 2005) between the pulse maximum of the cool blackbody component from the whole surface of the NS and the hot blackbody component from the hot spots (Figs. 3d and 3g).
To investigate if the NUV bridge flux is dominated by the thermal emission we have applied a simplified model (Greenstein & Hartke 1983) assuming that the NUV and soft X-rays fluxes ($`F`$) vary with rotational phase ($`\mathrm{\Phi }`$) as $`F=F_{}cos^2(\pi \mathrm{\Phi }+\alpha )+F_{}sin^2(\pi \mathrm{\Phi }+\alpha )`$, where $`F_{}`$, $`F_{}`$ are the fluxes at the minimum and maximum of the pulse profile, respectively. These are assumed to originate from the surface regions with predominantly radial and tangential magnetic fields, and $`\alpha `$ accounts for a small shift ($`\mathrm{\Delta }\mathrm{\Phi }0.1`$) of the profile minimum from the zero phase.
Figure 4 shows that this model provides a qualitative fit (dashed curves) to the soft X-ray profiles and to the smooth part of the NUV profile including the bridge. It also brings out the non-thermal counterparts of the NUV sub-pulses in X-rays. However, assuming a blackbody spectrum of the thermal emission with a maximum in the soft X-rays (e.g., Zavlin & Pavlov 2004), and that the temperature $`TF^{1/4}`$ in this range, we can estimate from the X-ray model light curves with $`F_{}/F_{}1.25`$ that the effective temperature variations averaged over the visible surface of the NS are $`5`$%. The smooth flux variation over the NUV is larger: $`F_{}/F_{}2`$. Taking into account that any thermal NUV is in the Rayleigh-Jeans (RJ) part of the blackbody spectrum where $`TF`$, we obtain much larger temperature variations ($``$50%) from the NUV fit. This contradiction with the X-rays can hardly be accounted for by any variations of the visible emitting area during the NS rotation ($`10\%`$, De Luca et al. 2005). An extrapolation of the estimated X-ray temperature variations to the NUV (Fig. 4, dotted horisontal curve), shows that the pulsed bridge emission in the NUV can not be due to this mechanism, but the resulting NUV emission is compatible with the observed persistent flux. This means that the bridge emission is probably not thermal. It could instead result from overlapping of the two non-thermal sub-pulses.
To further investigate if the NUV persistent flux is thermal, we can compare its relative contribution to the total flux, $`30`$%, with the expected contribution from the thermal cool blackbody component provided by extrapolation of the combined two blackbody $`+`$ power-law fit obtained in X-rays for the phase averaged spectrum (e.g., Koptsevich et al. 2001) to the NUV. Such a model is further discussed in Sect. 3.3 and shown in Fig. 5, and yields $`25`$% thermal contribution at 2800 Å, where the PRISM throughput peaks. This is consistent with the unpulsed fraction in the NUV observations. For comparison, the expected blackbody contribution in the optical range at 5000 Å is several times smaller, $`7`$%. In this region, Kern et al. (2003) were not able to detect any persistent flux at $`16`$%. In this interpretation, we are actually resolving the surface of the NS during the off-pulse in the NUV. Future observations, allowing for phase resolved spectroscopy in the NUV, FUV, and optical as well as detailed comparison with the phase resolved spectroscopy in the X-rays, are necessary to confirm such an interpretation.
### 3.3 The near-UV spectrum
The combined phase averaged spectrum of the pulsar is shown in Fig. 5. The spectrum is rather noisy at the edges of the observed range, while in the range $`21503200`$ Å, where the NUV-MAMA/PRISM has its maximal throughput, the S/N is $`3`$ per bin. The spectral energy distribution is rather flat with the flux apparently increasing with increasing frequency. A power-law (PL; $`F_\nu \nu ^{\alpha _\nu }`$) fit to the spectrum results in a spectral index of $`\alpha _\nu =0.1\pm 0.6`$. This is indicated by the flat dotted line in Fig. 5. The slope may be slightly positive, but a negative index is also consistent with the data.
To check the absolute flux calibration we compared the NUV spectrum with the available broadband NUV photometric data obtained with the HST/FOC in the F342W and F195W bands (Pavlov et al. 1997). Within the uncertainties the spectral data are consistent with the photometric magnitudes.
The interstellar color excess $`E(BV)`$ toward PSR B0656+14 has been estimated to be in the range of $`0.010.05`$ mag with a likely value of $`E(BV)=0.03`$ mag (Pavlov et al. 1996; Kurt et al. 1999). The dereddened \[$`E(BV)=0.03`$\] spectral index is $`\alpha _\nu =0.35\pm 0.5`$. This is shown by the solid line in Fig. 5. Dereddening the spectrum by this amount gives an integrated flux in the range $`17003400`$ Å of $`3.4\pm 0.3\times 10^{15}`$erg s<sup>-1</sup> cm<sup>-2</sup>. At a distance of 288 pc this corresponds to the luminosity $`L_{\mathrm{NUV}}=3.4\times 10^{28}`$ erg s<sup>-1</sup>, assuming isotropy of the emission. Even for an extinction of $`E(BV)=0.05`$ mag, the one sigma limit of the spectral index remains smaller than unity. This is significantly lower than the index of the RJ part of the blackbody spectrum ($`\alpha _\nu =2`$), and strongly favors a non-thermal component for the bulk of the emission in the NUV range. To state this in a different way, to reach the RJ spectral index would require an extinction of $`E(BV)=0.2`$, which is much outside the suggested range. Even the maps of Schlegel et al. (1998) only indicate $`E(BV)=0.09`$ throughout the entire Galaxy along this line-of-sight.
The observed NUV spectrum is actually compatible with the extension of the absorbed combined spectral fit, two blackbodies (BB) $`+`$ power-law (PL), of the X-ray data (the hot blackbody component does not contribute significantly in the NUV and is not shown in Fig. 5). Although such an extrapolation is uncertain (see, e.g., Kargaltsev et al. 2005), we note that the shape of the dereddened sum BB+PL (Fig. 5, thick dashed line) in the NUV is very similar to the observed NUV spectrum, although the PL clearly dominates in most of the NUV range. A non-thermal origin for the bulk of the NUV emission is also consistent with our analysis of the pulse profile. The change-over from the PL-component to the BB-slope seems to occur at the high frequency boundary of the observed NUV-range (Fig. 5). If this is true, we expect a higher spectral slope and a larger fraction of the persistent flux in the FUV range.
A possible absorption feature is hinted at $``$2400 Å (15.10 in Log $`\nu `$). The formal significance is low ($`2\sigma `$, see Fig. 5), but a flux depression around this wavelength is present in both visits. If confirmed by future observations, such a feature could be interpreted as an electron/positron cyclotron line formed at $`B5\times 10^8`$ G. Assuming the dipole surface field obtained from the pulsar spin-down rate (Table 1) this places the absorbing plasma at 220 km above the NS surface. This is roughly consistent with the minimum height ($`350`$ km) for the optical emission estimated from the assumed critical synchrotron frequency by Kern et al. (2003). The faintness of the absorption feature suggests that the obscuring magnetospheric plasma is optically thin and cannot affect the radiation in X-rays. No spectral features or depressions have been detected in the soft X-rays. This disfavors interpretations of the phase lags between the maxima of the cool and hot BB profiles and the radio pulse based on existence of a magnetospheric “blanket” affecting the pulse shape (e.g., De Luca et al. 2005 and references therein). These phase lags could instead be the result of a strong inhomogeneity of the surface magnetic field of the NS. On the other hand, the spectral energy distribution of PSR B0656+14 in the optical/near-IR range appears to be non-monotonic and shows potential depressions (Koptsevich et al. 2001; Komarova et al. 2003). These features are deeper than the possible absorption line in the NUV. This could be explained by increasing the optical depth of the absorbing magnetospheric plasma with the wavelength. Optical spectroscopy and timing observations of the pulsar in different optical bands are needed to study the nature of the dips and their relation to the marginal absorption line in the NUV.
## 4 Conclusions
We have detected the near-UV pulsations of PSR B0656
+14 using the STIS NUV-MAMA onboard the HST. The pulse period we derive is in agreement with previous estimates. The pulse profile is double-peaked with an inter-pulse bridge and is very similar to the profile detected in the optical range by Kern et al. (2003). Sharp pulses and a high NUV pulse fraction, $`70\pm 12`$%, suggest a non-thermal origin of the emission.
We have compared the NUV pulse profile with pulse profiles from the radio to gamma-rays. The first NUV sub-pulse is in phase with the marginal gamma-ray pulse, while the second NUV sub-pulse is similar in both shape and phase with the pulse detected in the hard X-rays. This favors a single origin for the non-thermal pulsed part of the emission from the optical to the X-ray range, as previously has also been seen from the broadband photometric observations in the optical and NIR (e.g., Koptsevich et al. 2001). The NUV sub-pulse counterparts can also be resolved from the shallow soft X-ray profiles dominated by thermal emission from the surface of the NS.
A simple model for the thermal X-ray flux, in terms of an anisotropic heat distribution of the NS surface, does not explain the pulsating NUV bridge emission, which is therefore likely non-thermal. However, any persistent flux in the NUV could contribute at the 30$`\pm 12`$% level, which is potentially higher than the upper limit $`16`$% measured in the optical (Kern et al. 2003). This is in qualitative agreement with the thermal flux from the abovementioned model, and we therefore argue that the non-pulsed NUV emission is due to the RJ part of the thermal emission component.
We have also measured the NUV spectrum of PSR B0656+14. The phase averaged spectrum is basically flat, $`\alpha _\nu =0.35\pm 0.5`$. This slope is consistent with contributions from both a thermal component from the whole surface of the cooling NS and from a magnetospheric power-law component extrapolated from the X-rays. The main spectral change from the power-law dominating magnetospheric emission to the thermal RJ from the NS surface is likely to occur at the boundary between the NUV and FUV ranges. Preliminary reports on FUV emission from PSR B0656+14 indicate that the FUV is indeed thermal (Pavlov et al. 2004). This also agrees with our model where the non-pulsed NUV flux is thermal. We therefore predict a higher spectral slope and a lower pulse fraction in the FUV.
We note that a similar situation applies to the middle-
-aged Geminga pulsar (Kargaltsev et al. 2005, their Figs. 7 and 10). However, the younger Vela pulsar (Shibanov et al. 2003; Romani et al. 2005) and the Crab pulsar (Sollerman 2003) display a relatively flat spectrum from the near-IR to UV range. Moreover, the non-thermal emission of the young Crab pulsar and the Crab twin PSR B0540-69 (Boyd et al. 1995) display pulse-profiles that also does not change significantly from the optical to the UV. The same seems to be true for Geminga (Kargaltsev et al. 2005) and for Vela (Romani et al. 2005). This could indicate that a unique mechanism, which apparently does not strongly depends on pulsar age, drives the non-thermal pulsed emission in the optical and UV.
###### Acknowledgements.
We are grateful for help from Alexei Kop-
-tsevich in the initial phase of this study. Part of this research has made use of the database of published pulse profiles maintained by the European Pulsar Network. This work was supported by NASA and The Royal Swedish Academy of Sciences. YAS were supported by the RFBR (grants 03-02-17423, 03-07-90200 and 05-02-16245) and RLSS programme 1115.2003.2. The research of PL is further sponsored by the Swedish Research Council. PL is a Research Fellow at the Royal Swedish Academy supported by a grant from the Wallenberg Foundation. DL was supported in part by funding through GO proposal 9156. TG and DL were supported in part through the STIS Guaranteed Time Observations resources. |
warning/0506/quant-ph0506200.html | ar5iv | text | # Solving Satisfiability Problems by the Ground-State Quantum Computer
## I Introduction
A quantum computer has been expected to outperform its classical counterpart in some computation problems. For example, the well-known Shor’s factoring algorithmShor and Grover’s algorithmGrover accelerate exponentially and quadratically compared with the classical algorithms, respectively. It is a challenge to find whether a quantum computer outperforms on other classically intractable problemsFarhi1 ; Hogg , which cannot be solved classically in polynomial time of $`N`$, the number of the input bits.
Especially interesting are the NP-complete problemsNP-Complete , which include thousands of problems, such as the Traveling Salesman problemEC and the satisfiability (SAT) problems. All NP-complete problems can be transformed into each other in polynomial steps. If one of the NP-complete problems can be solved in polynomial time by an algorithm even in the worst case, then all NP-complete problems can be solved in polynomial time. However, it is widely believed that such a classical algorithm doesn’t exist. In this paper we will discuss quantum algorithm for solving SAT problems. A $`K`$-SAT problem deals with $`N`$ binary variables submitted to $`M`$ clauses with each clause $`C_i`$ involving $`K`$ bits, and the task is to find $`N`$-bit states satisfying all clauses. When $`K>2`$, $`K`$-SAT is NP-Complete, and some instances become classically intractable when the parameter $`\alpha =M/N`$, as $`M,N\mathrm{}`$, approaches the threshold $`\alpha _c(K)`$3SAT ; nature ; SAT ; NP .
Due to the properties of quantum mechanics, it’s hard to design quantum algorithms directly from intuition. In the present paper, we will study the properties of the ground-state quantum computer(GSQC), and show that the special property of the GSQC naturally leads to algorithm for solving SAT problems. Although we cannot determine whether or not this algorithm solves the NP-complete problems in polynomial time, we try to shed light on the complexity of the NP-complete problems.
In the following sections, at first we introduce the idea of the ground-state quantum computerMizel1 ; Mizel2 ; Mizel3 and its energy gap analysisours , then demonstrate the particular property of the GSQC, which provides a direct approach to solving SAT problems, and finally an example, an algorithm for solving the 3-bit Exact Cover problem, is given.
## II Ground-State Quantum Computer and its Energy Gap
A standard computer is characterized by a time-dependent state $`|\psi (t_i)=U_i|\psi (t_{i1}),`$ where $`t_i`$ denotes the instance of the $`i`$-th step, and $`U_i`$ represents for a unitary transformation. For a GSQC, the time sequence is mimicked by the spatial distribution of its ground-state wavefunction $`|\psi _0`$. As proposed by Mizel et.al.Mizel1 , the time evolution of a qubit may be represented by a column of quantum dots with multiple rows, and each row contains a pair of quantum dots. State $`|0`$ or $`|1`$ is represented by finding the electron in one of the two dots. It is important to notice that only one electron exists in a qubit. The energy gap, $`\mathrm{\Delta }`$, between the first excited state and the ground state determines the scale of time cost.
### II.1 Hamiltonians of GSQC
A GSQC is a circuit of multiple interacting qubits, whose ground state is determined by the summation of the single qubit unitary transformation Hamiltonian $`h^j(U_j)`$, the two-qubit interacting Hamiltonian $`h(CNOT)`$, the boost Hamiltonian $`h(B,\lambda )`$ and the projection Hamiltonian $`h(|\gamma ,\lambda )`$.
The single qubit unitary transformation Hamiltonian has the form
$`h^j(U_j)=ϵ[C_{j1}^{}C_{j1}+C_j^{}C_j(C_j^{}U_jC_{j1}+h.c.)],`$ (1)
where $`ϵ`$ defines the energy scale of all Hamiltonians, $`C_j^{}=\left[c_{j,0}^{}c_{j,1}^{}\right]`$, $`c_{j,0}^{}`$ is the electron creation operator on row $`j`$ at position $`0`$, and $`U_j`$ is a two dimension matrix representing the unitary transformation from row $`j1`$ to row $`j`$. The boost Hamiltonian is
$`h^j(B,\lambda )`$ $`=`$ $`ϵ[C_{j1}^{}C_{j1}+{\displaystyle \frac{1}{\lambda ^2}}C_j^{}C_j`$ (2)
$`{\displaystyle \frac{1}{\lambda }}(C_j^{}C_{j1}+h.c.)],`$
which amplifies the wavefunction amplitude by the large value number $`\lambda `$ compared with the previous row at $`|\psi _0`$. The projection Hamiltonian is
$`h^j(|\gamma ,\lambda )`$ $`=`$ $`ϵ[c_{j1,\gamma }^{}c_{j1,\gamma }+{\displaystyle \frac{1}{\lambda ^2}}c_{j,\gamma }^{}c_{j,\gamma }`$ (3)
$`{\displaystyle \frac{1}{\lambda }}(c_{j,\gamma }^{}c_{j1,\gamma }+h.c.)],`$
where $`|\gamma `$ is the state to be projected to on row $`j`$ and to be amplified by $`\lambda `$ at $`|\psi _0`$. The interaction between qubit $`\alpha `$ and $`\beta `$ can be represented by $`h(CNOT)`$:
$`h_{\alpha ,\beta }^j(CNOT)`$ (4)
$`=`$ $`ϵC_{\alpha ,j1}^{}C_{\alpha ,j1}C_{\beta ,j}^{}C_{\beta ,j}+h_\alpha ^j(I)C_{\beta ,j1}^{}C_{\beta ,j1}`$
$`+c_{\alpha ,j,0}^{}c_{\alpha ,j,0}h_\beta ^j(I)+c_{\alpha ,j,1}^{}c_{\alpha ,j,1}h_\beta ^j(N).`$
where for $`c_{a,b,\gamma }^{}`$, its subscription $`a`$ represents for qubit $`a`$, $`b`$ for the number of row, $`\gamma `$ for the state $`|\gamma `$. With only $`h^j(U_j)`$ and $`h_{\alpha ,\beta }^j(CNOT)`$, its ground state isMizel2 :
$`|\psi _0^j`$ $`=`$ $`[1+c_{\alpha ,j,0}^{}c_{\alpha ,j1,0}(1+C_{\beta ,j}^{}C_{\beta ,j1})`$ (5)
$`+c_{\alpha ,j,1}^{}c_{\alpha ,j1,1}(1+C_{\beta ,j}^{}NC_{\beta ,j1})]`$
$`\times {\displaystyle \underset{a\alpha ,\beta }{}}(1+C_{a,j}^{}U_{a,j}C_{a,j1})|\psi ^{j1}.`$
All above mentioned Hamiltonians are positive semidefinite, and are the same as those in Mizel1 ; Mizel2 ; Mizel3 . Only pairwise interaction is considered.
The input states are determined by the boundary conditions applied upon the first rows of all qubits, which can be Hamiltonian $`h^0=E(I+_ia_i\sigma _i)`$ with $`\sigma _i`$ being Pauli matrix and $`_ia_i^2=1`$. For example, with $`h^0=E(I+\sigma _z)`$, $`|\psi _0`$ on the first row is $`|1`$; with $`h^0=E(I\sigma _x)`$, it is $`\left(|0+|1\right)`$. If $`E`$ is large enough, for example, at $`E10ϵ`$, the energy gap will saturate and become independent of the magnitude of $`E`$ ours .
To implement an algorithm, on final row of each qubit a boost or a projection Hamiltonian is applied so that $`|\psi _0`$ concentrates on the position corresponding to the final instance in the standard paradigm, hence measurement on the GSQC can read out the desired information with appreciable probability. With boost Hamiltonian or projection Hamiltonian on last rows, the ground-state wavefunction amplitude on those rows will be $`\lambda `$ of that on their neighboring rows.
By observing the expression Eq.(5), it’s easy to find that, for two interacting qubits, the ground-state wavefunction has the formours
$`\left(|\psi _{upstream}^{control}+|\psi _{downstream}^{control}\right)|\psi _{upstream}^{target}`$
$`+|\psi _{downstream}^{control}|\psi _{downstream}^{target},`$ (6)
where each qubit is divided by the interacting Hamiltonian as two parts, and the part with boundary Hamiltonian $`h^0`$ is called as upstream, and the other part is called downstream. In this paper, we always use this definition when upstream or downstream is mentioned.
### II.2 Energy Gap of GSQC
Now we briefly introduce how to find the scale of the energy gap of a GSQC. For details, please find in ours .
With multiple interacting qubits, one needs to evaluate on each qubit the parameter $`1/x`$, the overall amplitude of lowest excited state on top rows of this qubit before meeting the first interacting Hamiltonian, assuming that on the top rows of this qubit the lowest energy excited state is orthonormal to $`|\psi _0`$ while states on all other qubits remain the same as the corresponding ground state with only magnitude changed. The energy gapours is given by the minimum parameter $`1/x`$ as
$`\mathrm{\Delta }ϵ(1/x)_{min}^2.`$ (7)
The rule of estimating $`1/x`$ is as followingours : With each qubit ended with either a projection or a boost Hamiltonian containing the same (for simplicity) amplifying factor $`\lambda 1`$, when estimating $`1/x`$ for a qubit, say qubit $`A`$, ($`i`$) at first $`x`$ is set to 1; ($`ii`$) the boost Hamiltonian, $`not`$ the $`projection`$ Hamiltonian, on qubit $`A`$ itself increases $`x`$ by multiplication of $`\lambda `$; ($`iii`$) if qubit $`A`$ directly interacts with another qubit, say qubit $`B`$, by Hamiltonian $`h_{AB}`$, then we determine, excluding qubit $`A`$, on the qubit $`B`$ the ground-state wavefunction amplitude ratio of the upstream part (with respect to $`h_{AB}`$) over its final row, $`\frac{1}{x_B}`$, contributions to $`\frac{1}{x_B}`$ are found one by one according to Eq.(6): if the upstream part of qubit $`B`$ doesn’t coexist with the states on final rows of any one qubit, except for qubit $`A`$, then $`x_B`$ should be multiplied by a $`\lambda `$; ($`iv`$) finally, the value of $`1/x`$ on qubit $`A`$ should be multiplied by $`\frac{1}{x_B}`$, or $`\mathrm{\Pi }_i\frac{1}{x_B^i}`$ if more than one qubit directly interact with qubit $`A`$.
According to the above rule, the energy gap $`\mathrm{\Delta }`$ of single qubit with length $`n`$ and ended with boost Hamiltonian $`h(B,\lambda )`$ scales as $`ϵ/\lambda ^2`$ as $`\lambda n`$; when ended with projection Hamiltonian $`h(|\gamma ,\lambda )`$, $`\mathrm{\Delta }`$ is independent of $`\lambda `$. For two $`n`$-row qubits interacting by $`h(CNOT)`$, $`\mathrm{\Delta }ϵ/\lambda ^4`$ as $`\lambda n`$ if both qubits ended with $`h(B,\lambda )`$ or one with $`h(B,\lambda )`$ and the other with $`h(|\gamma ,\lambda )`$. Numerical calculations confirm these results. The Fig.(1b) and Fig.(2) in ours are two examples on how to apply the above rule on complicated circuits.
Complicated GSQC circuit may have exponentially small energy gap, like the circuit in Fig.(1b) of ours , and assembling the GSQC circuit directly following the algorithm for the standard paradigm, such as quantum Fourier transform, leads to exponentially small energy gap. In order to avoid such small gap, the teleportation boxes are introduced on each qubit between two control Hamiltoniansours . Fig.(1) shows how the CNOT interacting qubits is modified by inserting teleportation boxes on each qubit’s upstream and downstream part. The teleportation boxes make all qubits short (the longest qubit has length 8), on the other hand, for arbitrary GSQC circuit they make the energy gap only polynomially small $`\mathrm{\Delta }ϵ/\lambda ^8`$ours if all boost and projection Hamiltonians have the same amplifying factor $`\lambda `$. To determine magnitude of $`\lambda `$, one only needs to count the total number of qubits in the circuit, say $`L`$, which is proportional to the number of control operation in an algorithm, then the probability of finding all electrons on final rows is $`P(1C/\lambda ^2)^L`$ with $`C`$ being 8, the maximum length of qubit. In order to have appreciable $`P`$, we set $`\lambda L^{1/2}`$, hence $`\mathrm{\Delta }ϵ/L^4`$. The details can be found in ours .
### II.3 Energy Gap When Projecting Small Fraction of a State
In the previous section the rule for finding scale of the energy gap is under the assumption that when a projection Hamiltonian $`h(|\gamma ,\lambda )`$ is applied, $`|a|/\sqrt{|a|^2+|b|^2}`$ is appreciable for the ground state on row just before the projection Hamiltonian:
$`a|\gamma +b|\stackrel{~}{\gamma },`$ (8)
where $`|\gamma =|0`$ ($`|1`$) and $`|\stackrel{~}{\gamma }=|1`$ ($`|0`$). The ground-state wavefunction concentrates on the last row, hence the first excited state wavefunction cannot have appreciable weight there because otherwise $`\psi _1|\psi _00`$. When evaluate $`1/x`$ on a qubit, the projection Hamiltonian on the qubit itself doesn’t contribute to $`1/x`$. For example, concerning a single qubit, as shown in Fig.(2), with only identical transformations $`h(I)`$ and ended by $`h(|0,\lambda )`$, if $`h^0=E(I\sigma _x)`$ so that $`|\psi _0`$ on the first row is $`|0+|1`$, then the energy gap $`\mathrm{\Delta }`$ is almost independent of $`\lambda `$, as shown in the top line of Fig.(3).
However, if in Eq.(8) $`|a|/\sqrt{|a|^2+|b|^2}1`$, then $`\mathrm{\Delta }`$ depends on $`\lambda `$ until $`\lambda `$ reaching $`\sqrt{|a|^2+|b|^2}/|a|`$. This is because when $`\lambda <\sqrt{|a|^2+|b|^2}/|a|`$, the ground-state wavefunction has little weight on the last row, and the first excited state concentrates there, hence $`1/x`$ is small, leading to small energy gap. When $`\lambda >\sqrt{|a|^2+|b|^2}/|a|`$, ground state wavefunction has large part on the last row, then just like the above situation, energy gap is not further affected by increasing $`\lambda `$.
To confirm the above analysis, we numerically calculate the energy gap of a 6-row single qubit ended with projection Hamiltonian, as shown in Fig.(2).
The boundary Hamiltonian is $`h^0=10ϵ(I+\alpha \sigma _z\sqrt{1\alpha ^2}\sigma _x)`$, all other Hamiltonians except for that at final row are $`h^j(I)`$ with $`j=1,2,3,4`$, and on the final row there is a projection Hamiltonian $`h(|0,\lambda )`$. By tuning $`\alpha `$, we can determine what fraction of wavefunction is projected from the 5th row to the last row. At $`\alpha =0,0.9,0.99,0.999,0.9999,0.99999`$, on the 5th row the ground state wavefunctions are $`a|0+b|1`$ with $`a/b=1,0.23,0.071,0.022,0.0071,0.0022.`$ Fig.(3) shows that the energy gap is $`\mathrm{\Delta }ϵ/\lambda ^2`$ as $`\lambda <|\sqrt{|a|^2+|b|^2}/a|`$, and when $`\lambda >|\sqrt{|a|^2+|b|^2}/a|`$, $`\mathrm{\Delta }`$ becomes independent on $`\lambda `$. The independent $`\mathrm{\Delta }`$ is proportional to $`ϵ|a/\sqrt{|a|^2+|b|^2}|^2`$.
In order to make the ground-state wavefunction concentrate on the last row so that measurement corresponds to the desired state, $`\lambda `$ must be larger than $`|\sqrt{|a|^2+|b|^2}/a|`$, Thus the energy gap is determined by the fraction of state been projected. If $`|a/\sqrt{|a|^2+|b|^2}|`$ is exponentially small, which may happen in certain case, then the energy gap is exponentially small. Fortunately, this doesn’t happen to the GSQC implement of Quantum Fourier Transform, there all projection Hamiltonians are applied to teleportation circuit, and $`|a/b|=1`$. However, it plays a role in the algorithm presented in the following section.
For multiple interacting qubits, if $`|a/\sqrt{|a|^2+|b|^2}|1`$ in Eq.(8), the rule of finding energy gap needs modification: With all qubits ended with either a projection or a boost Hamiltonian containing the same amplifying factor $`\lambda 1`$, when estimating $`1/x`$ for any qubit, say qubit $`A`$, ($`i`$) at first $`x`$ is set to 1; ($`ii`$) the boost Hamiltonian, $`or`$ the $`projection`$ Hamiltonian, on qubit $`A`$ itself increases $`x`$ by multiplication of $`\lambda `$ or Min$`(\lambda ,|\sqrt{|a|^2+|b|^2}/a|)`$; ($`iii`$) if qubit $`A`$ directly interacts with another qubit, say qubit $`B`$ by Hamiltonian $`h_{AB}`$, then we determine, excluding qubit $`A`$, on the qubit $`B`$ the amplitude ratio of the upstream part (divided by $`h_{AB}`$) over its final row, $`1/x_B`$, and contribution to $`1/x_B`$ from other qubits are found one by one according to Eq.(6): if the upstream part of qubit $`B`$ doesn’t coexist with the states on final rows of a qubit, except for qubit $`A`$, then $`x_B`$ should be multiplied by $`\lambda `$ (ended with boost Hamiltonian) or $`\lambda |a^{\prime \prime }/\sqrt{|a^{\prime \prime }|^2+|b^{\prime \prime }|^2}|`$ (ended with projection Hamiltonian); ($`iv`$) finally, the value of $`1/x`$ on qubit $`A`$ should be multiplied by $`1/x_B`$ or $`\mathrm{\Pi }_i1/x_B^i`$ if more than one qubit directly interact with qubit $`A`$.
It is easy to find that when $`|b/a|`$, $`|b^{\prime \prime }/a^{\prime \prime }|1`$ and $`\lambda 1`$, we get the same result as the previous subsection. After $`1/x`$’s on all qubits being evaluated, the minimum $`1/x`$ gives the energy gap scale as
$`\mathrm{\Delta }ϵ(1/x)_{min}^2.`$
## III Quantum Algorithm by GSQC
There are some interesting properties for the GSQC. Although it was shownLloyd that, concerning on time cost, a quantum computer composed of (time varying) local Hamiltonians is equivalent to standard circuit quantum computer, GSQC provides some insights to design quantum algorithm for certain problems. For example, the projection Hamiltonian, which corresponds to measurement in standard paradigm, can amplify the probability at a particular state. Here we are not claiming that the GSQC is more powerful than standard quantum computer, however, the GSQC does provide a direct approach for certain problem, as shown below is the algorithm for the SAT problems.
At first we give the simplest example, considering that qubit $`i`$ $`CNOT`$ controls an ancilla qubit that is at the right side in Fig.(4), and their boundary Hamiltonians make the ground state on their first rows are $`|0+|1`$ and $`|0`$, respectively. On last rows the ground state is $`|0|0+|1|1`$. If we apply a boost Hamiltonian on qubit $`i`$ and a projection Hamiltonian $`h(|0,\lambda )`$ on the ancilla qubit, then at the ground state the state on final rows becomes $`|0|0`$. The large value of $`\lambda `$ makes sure that there is large probability to find two electrons on the final rows of the two qubits at the ground state. So by choosing projected state on the ancilla qubit, we can have the selected state $`|0`$ on qubit $`i`$, and prevent the other state $`|1`$ from reaching its final row. If qubit $`i`$ entangles with other qubit, such as $`|0|\alpha +|1|\beta `$, the entanglement of $`|0|\alpha `$ will not be affected. Thus we call circuit in Fig.(4) a filter for the clause $`i=0`$.
Another example makes more sense. Lets consider a SAT problem with clauses, each of which involves two qubits, say qubit $`i`$ and $`j`$, and requires $`i+j=1`$. We can implement this clause by the GSQC circuit in Fig.(5). In this figure there are three qubits: qubit $`i`$, qubit $`j`$ and an ancilla qubit that is at the left side in the figure. It’s easy to find that if on the first row $`|i=\alpha _i|0+\beta _i|1,|j=\alpha _j|0+\beta _j|1`$ and the ancilla qubits at $`|0`$, then at the ground state on the final rows of the three qubits the state is $`|i|j|ancilla=(\beta _i\alpha _j|1|0+\alpha _i\beta _j|0|1)|1`$, which satisfies the clause. Thus circuit in Fig.(5) filters out states not satisfying this simple clause and lets through those satisfying states. It is important to note that at the beginning if the satisfying states entangle with other qubits not showing in the figure, these entanglements keep untouched.
The property of GSQC brings up new quantum algorithm naturally. Here we present one to solve the SAT problems as shown in Fig.(6), a GSQC circuit to solve a 3-SAT problem with only 9 bits. It’s easy to be extended to $`N`$-bit $`K`$-SAT problems. Each clause is implemented by a “filter box”, and the circuit inside each filter box makes sure that on rows immediately below it the ground state satisfies the clause $`C_i`$, or we can say those unsatisfying states are filtered out. This can be realized by projection and boost Hamiltonians like in Fig.(4) and Fig.(5).
In Fig.(6), the initial state on the top rows of qubit from 1 to 9 is $`(|0+|1)(|0+|1)\mathrm{}(|0+|1)`$, which is enforced by the boundary Hamiltonians, $`h^0=E(I\sigma _x)`$; the clause involving qubit 1, 2 and 3 is implemented by filter box 1, the clause involving qubit 2, 3 and 4 implemented by filter box 4, the clause involving qubit 3, 4 and 8 implemented by filter box 6, etc.
When all constraints are implemented, at ground state the states measured on the final rows of the $`N`$ qubits should be superposition of all states satisfying all constraints. No backtracking is needed.
Concerning energy gap, unlike the circuit for quantum Fourier transform, in which the energy gap is determined by the number of control operationours , the SAT problems is more complicated to evaluate because it might involve the situation to project a very small fraction of state as shown in section II.3. For example, if one constructs a GSQC for the Grover’s search problem with one condition to find a unique satisfying state from $`2^N`$ states, then he will find that there is an ancilla qubit containing such unnormalized state
$`|0|\text{satisfying}+{\displaystyle \underset{i=1}{\overset{2^N1}{}}}|1|\text{unsatisfying}^{(i)}`$ (9)
before the projection Hamiltonian $`h(|0,\lambda )`$. In order to amplify the amplitude of the correct state on the final row, it requires $`\lambda 2^{N/2}`$. Its energy gap is hence less than $`2^N`$, which is consistent with the limit set by many other worksGrover ; Farhi2 ; Bennett .
## IV Example: the 3-Bit Exact Cover Problem
Up to now the filters, Fig.(4) and Fig.(5), we have given are trivial, and now we give an example on how to implement a filter for a serious problem. We focus on the 3-bit Exact Cover problemEC , an instance of SAT problem, which belongs to NP-complete. Following is the definition of the 3-bit Exact Cover problem:
There are $`N`$ bits $`z_1,z_2,\mathrm{},z_N`$, each taking the value 0 or 1. With $`O(N)`$ clauses applied to them, each clause is a constraint involving three bits: one bit has value 1 while the other two have value 0. The task is to determine the $`N`$-bit state satisfying all the clauses.
### IV.1 GSQC Circuit for the 3-bit Exact Cover Problem
The algorithm is implemented by the circuit in Fig.(6). Each filter box, in our algorithm, involves three qubits, say qubit $`i,j`$ and $`k`$, which are represented by gray dot columns in Fig.(7). We add two ancilla qubits: qubit 1 and qubit 2, which are represented by dark dot columns. Qubit $`i,j`$ and $`k`$ at the first row are in the state $`(|1+|0)`$ if they have not experienced any clause yet, and the two ancilla qubits are in the states $`|\widehat{0}`$ and $`|\stackrel{~}{0}`$ on top rows by selecting proper boundary Hamiltonians, where $`|\widehat{\gamma }`$ corresponds to the state of ancilla qubit 1, and $`|\stackrel{~}{\gamma }`$ to the state of ancilla qubit 2.
Inside the dashed triangle of Fig.(7), after the first $`CNOT`$, we obtain state $`|\widehat{1}|1+|\widehat{0}|0`$; after the second $`CNOT`$: $`|\widehat{1}|1|0+|\widehat{0}|0|0+|\widehat{0}|1|1+|\widehat{1}|0|1`$; after the third $`CNOT`$:
$`|\widehat{1}\left(|1|0|0+|0|1|0+|0|0|1+|1|1|1\right)`$
$`+`$ $`|\widehat{0}\left(|1|1|0+|0|1|1+|1|0|1+|0|0|0\right).`$
Immediately below the triangle, if the system stays at the ground state, if electron in ancilla qubit 1 is measured to be on the row labeled by $`X`$ and at state $`|\widehat{1}`$, and if the three electrons on qubit $`i,j,k`$ are all found on the rows labeled by $`X`$, then the three-qubit state satisfies the clause except for $`|1|1|1`$.
The ancilla qubit $`2`$, starting at state $`|\stackrel{~}{0}`$, experiences $`CNOT`$ gates controlled by qubits $`j`$ and $`k`$, and $`R(\pm \pi /4)`$ transformations, defined in Tof as $`R_y(\pm \pi /4)`$, as shown within the dotted pentagon in Fig.(7). All those transformations happened inside the dotted pentagon are equivalent to a Toffoli gate except for some unimportant phasesTof : if both qubits $`j`$ and $`k`$ are in state $`|1`$, then the ancilla qubit $`2`$ reverses to state $`|\stackrel{~}{1}`$, otherwise, it remains at state $`|\stackrel{~}{0}`$. After this nearly Toffoli transformation, if at ground state electrons in qubit $`j,k`$ and ancilla qubit $`2`$ are found on rows labeled by $`Y`$, and if ancilla qubit 2 is at $`|\stackrel{~}{0}`$, then the three qubits will be at $`|\stackrel{~}{0}(|0|0+|1|0+|0|1)`$. Thus if at ground state all electrons are found on rows immediately below both the dashed triangle and the dotted pentagon, and if ancilla qubit 1 is at $`|\widehat{1}`$ and ancilla qubit 2 at $`|\stackrel{~}{0}`$, then the three qubits $`i,j,k`$ satisfy the clause:
$`|\widehat{1}|\stackrel{~}{0}\left(|1|0|0+|0|1|0+|0|0|1\right).`$ (10)
In order to make the satisfying states pass through the filter box with large probability, we add projection Hamiltonians and boost Hamiltonians as shown in the lower part of Fig.(7). The projection Hamiltonians on final rows of the two ancilla qubits limit and amplify the amplitude of the states we prefer: ancilla qubit 1 at $`|\widehat{1}`$, and ancilla qubit 2 at $`|\stackrel{~}{0}`$. If a qubit does not experience any more clause, it will end with a boost Hamiltonian, otherwise, its quantum state will be teleported to a new qubit through teleportation box, not shown in Fig.(7), and the new qubit experiences more clauses. Thus the projection Hamiltonians on two ancilla qubits and boost Hamiltonians on the three qubits make sure that the ground-state wavefunction concentrates on the final rows in Fig.(7) with state at Eq.(10).
Noting that in the filter box all the three qubits $`i,j,`$ and $`k`$ always act as control qubits, thus the entanglements of these three qubits with other qubits not involved in this particular clause still keep the same. When adding a clause, the resulted states satisfying this clause will also satisfy all previous applied clauses. Thus unlike classical algorithm, no backtracking is needed.
### IV.2 Energy Gap Without Projecting Small Fraction of State
In this subsection, we assume applying each clause does decrease the number of satisfying state $`gradually`$, or equivalently, the projection Hamiltonian in the two ancilla qubits in each filter box, Fig.(7), does project appreciable part of state on the second last row. This assumption may not be correct in many SAT problems, especially close to $`\alpha _c`$.
In the circuit of Fig.(6), if there is at least one solution, and all electrons are simultaneously found on the final rows of all qubits, then the reading of the $`N`$-bit state satisfies all clauses.
In order to keep the energy gap from being too small, like in ours , on every qubit teleportation boxes are inserted between two control Hamiltonians, thus the total number of qubits increases while the energy gap $`\mathrm{\Delta }ϵ/\lambda ^8`$ if all the boost and the projection Hamiltonians have the same value of amplifying factor $`\lambda `$.
For one clause, or a filter box, it needs 10 teleportation boxes (each teleportation box adds two more qubits) on the original five-qubit circuit, noting that on the end of qubit $`i,j`$ and $`k`$ in Fig.(7) teleportation boxes are needed because more clause will be added. Thus adding one more filter box means adding 20 more qubits. The number of clause for a NP hard 3-bit Exact Cover problem is about the same order as the number of bits $`N`$3SAT , say $`\alpha N`$ with $`\alpha `$ being $`O(1)`$, then there are about $`20\alpha N`$ qubits and each of them ends with either a projection or a boost Hamiltonian. Probability of finding all electrons at the final rows is approximately
$`P\left(1C/\lambda ^2\right)^{20\alpha N},`$ (11)
where $`C=8`$, the length of the longest qubitours . It is assumed that, at ground state, in each filter box the ancilla qubit 1 and 2 have appreciable probability in $`|1`$ and $`|0`$ states, respectively, before projection Hamiltonians. Later we will address the situation when this assumption is violated.
In order to make the probability independent of number of bits $`N`$, we take $`\lambda ^2=DN`$, where $`D`$ is an arbitrary number. Then as $`N`$ becomes large, we obtain
$`P\left(1C/(DN)\right)^{20\alpha N}e^{20\alpha C/D},`$ (12)
and energy gap isours
$`\mathrm{\Delta }ϵ/\lambda ^8ϵ/(D^4N^4),`$ (13)
from which one can estimate time cost.
To make the GSQC circuit at ground state, we can use adiabatic approach: first we set $`\lambda =1`$ for boost and projection Hamiltonian on final rows of all qubits, and replace the single qubit Hamiltonian between the first two rows of all qubits by a boost Hamiltonian
$`h^{}(B,\lambda ^{})=ϵ[{\displaystyle \frac{1}{\lambda ^2}}C_1^{}C_1+C_2^{}C_2{\displaystyle \frac{1}{\lambda ^{}}}(C_1^{}C_2+h.c.)],`$ (14)
so that the wavefunction amplitude of the first row is boosted as $`\lambda ^{}1`$. Now in the ground state the electrons concentrate at the first rows as $`1/\lambda ^{}0`$, thus the ground state is easy to be prepared, and the energy gap $`\mathrm{\Delta }ϵ/n^2`$ with $`n=8`$ being the length of the longest qubit. The next step is turning the quantity $`1/\lambda ^{}`$ to 1 adiabatically, during which the energy gap remains at $`ϵ/n^2`$ and the ground-state wavefunction spreads to other rows from the first row. The third step is turning $`1/\lambda `$ from 1 to $`1/\sqrt{DN}`$ adiabatically. In this process the energy gap decreases monotonically from $`ϵ/n^2`$ to what we obtained above: $`ϵ/D^4N^4`$, and the ground-state wavefunction concentrates on the final rows of all qubit as we wish. Thus the scale of time cost is about $`T1/\mathrm{\Delta }^2N^8`$Farhi0 , local adiabatic approach may reduce the time cost furtherlocal .
### IV.3 Energy Gap for SAT Problems
Above analysis is under the assumption that the number of satisfying states gradually decreases as the clauses are implemented one by one. There is a situation that might hurt our algorithm: after adding one more clause, the number of satisfying states drops dramatically. Just like what happens to Grover’s search algorithm, in which the number of satisfying states drops from $`2^N`$ to 1, and as shown in Eq.(9), our algorithm involves a projection Hamiltonian on an ancilla qubit to project an exponentially small fraction of a state, thus the energy gap evaluation in the above subsection becomes invalid.
Does this happen to the general SAT problems? In nature it was suggested that close to the threshold $`\alpha _c`$ computational complexity might be related with the forming of a backbone, each of a subset of bits has average value close to 1 or 0 in the subspace of satisfying states. The existence of the backbone means that most satisfying states contain the state represented by the backbone, and if adding one more clause kicks out the states consistent with the backbone from satisfying subspace, the number of satisfying states drops dramatically, and this corresponds to projecting a small fraction of state.
Performance of our algorithm is not affected by forming of backbone, however, as more clauses applied, the disappearance of the already existed backbone in the satisfying subspace surely hurts. There is a criterion determining efficiency of our algorithm: the ratio $`S_j/S_{j+1}`$, with $`S_j`$ being the number of solutions when the $`j`$th clause is applied, and $`S_{j+1}`$ the number of solutions when the $`(j+1)`$th clause is applied. For example, $`S_0/S_1=8/3`$ for 3-bit Exact Cover problem. If $`S_j/S_{j+1}1`$, on the ancilla qubit of the $`(j+1)`$th filter box, the probability of finding electron on its final row will be $`p(1CS_j/(\lambda ^2S_{j+1}))`$. To make sure of appreciable probability of finding all electrons on the final row of all qubits, an overhead factor $`\sqrt{S_j/S_{j+1}}`$ for $`\lambda `$ on the ancilla qubit is needed, hence the amplifying factor in the projection Hamiltonian on the ancilla qubit should be $`\lambda \sqrt{S_j/S_{j+1}}`$. According to the analysis in Sec.II.3, the energy gap might be also determined by the parameter $`S_j/S_{j+1}`$. Because in a filter box, the ancilla qubit will end after the projection Hamiltonian, which should be at the position of qubit 8 or qubit 10 in Fig.(1) without the dotted line following. According to the rule described in section II.3, the parameter $`1/x`$ on this ancilla qubit should be
$`{\displaystyle \frac{1}{x}}={\displaystyle \frac{1}{\lambda ^2Min(\lambda ,\sqrt{\frac{S_j}{S_{j+1}}})}}.`$ (15)
The energy gap thus is
$`\mathrm{\Delta }=Min({\displaystyle \frac{ϵ}{\lambda ^8}},{\displaystyle \frac{ϵS_{j+1}}{\lambda ^4S_j}}).`$ (16)
If this ratio $`S_{j+1}/S_j`$ happens to be exponentially small, then our algorithm cannot solve the SAT problem in polynomial time. We cannot know in advance what $`S_{j+1}/S_j`$ is, however, we might be able to identify backbone by trials, and then choose proper order to implement clauses so that $`S_{j+1}/S_j`$ always can be kept not too small. However, if the NP-Complete problem means that one can never avoid an exponentially small $`S_{j+1}/S_j`$, then the quantum algorithm cannot solve NP-Complete problem in polynomial time.
## V Conclusion
In conclusion, we have demonstrated that a ground state quantum computer can solve a general SAT problem. A specific example, the 3-bit Exact Cover problem, is given. We show that a 3-bit Exact Cover problem can be solved by the quantum algorithm described here, and the time cost is related with the number of bits $`N`$ and the parameter $`S_{j+1}/S_j`$. If $`S_{j+1}/S_j`$ stays only polynomially small, then the presented algorithm can solve this SAT problem in polynomial time. It will be interesting if one finds the equivalent algorithm by standard paradigm.
I would like to thank A. Mizel for helpful discussion. This work is supported in part by Jun Li Foundation, the NSF under grant # 0121428, and ARDA and DOD under the DURINT grant # F49620-01-1-0439. |
warning/0506/astro-ph0506493.html | ar5iv | text | # Monthly Notices of the Royal Astronomical Society: LaTeX 2ε style guide for authors
## 1 Introduction
There has been much work done surveying young open clusters for very low mass stars (VLMS) and brown dwarfs (BDs) (e.g. Bouvier et al. 1998, Dobbie et al. 2002). The focus has generally been to identify BDs and measure cluster mass functions (MFs) in a variety of environments to establish if the initial mass function (IMF) is universal or not – the answer to which has important implications for our understanding of star and BD formation. The study of older clusters is also important if we are to understand how the dynamical evolution of the cluster affects the shape of their MF – over time VLMS and BDs are expected to be preferentially ejected (eg. de la Fuente Marcos & de la Fuente Marcos 2000). For the Hyades, Dobbie et al. (2002) searched for low-mass (0.1–0.06M) members in a 10.5 square degree survey, but found only one (previously known) stellar member. They estimate that 4–5 Hyads should have been found in this mass range if the Hyades and Pleiades MFs are identical, and conclude that dynamical evaporation of low-mass Hyads is the most likely explanation for the deficit. However, studies of other clusters with similar ages are clearly needed if we are to properly address this question.
Late cluster members are also ideal for testing ultra-cool atmosphere models of late M and L spectral types (T<sub>eff</sub> $``$ 2500K Leggett et al. 2001). At such low T<sub>eff</sub>, atmospheric dust grains condense out of the gas phase, and strongly affect both colour and spectral properties. However, dust grain properties depend not only on temperature but on gravity ($`g`$)and metallicity (\[M/H\]) which cannot be measured with confidence for late M and L field dwarfs. But with well constrained age (and hence radii and masses inferred from evolutionary models) late M and L cluster members’ $`g`$, T<sub>eff</sub> and \[M/H\] can be known. Indeed, spectroscopic studies of such objects with a range of different (but known) ages would allow us to empirically isolate the T<sub>eff</sub>, $`g`$, and \[M/H\] dependence of both narrow and broadband spectral features.
The Praesepe cluster has an age of (0.9$`\pm `$0.5) Gyrs, lies at a distance of $``$ 170 pc, has a near solar metallicity (Hambly et al. 1995a) and zero reddening (Crawford & Barnes 1969). Large scale proper motion studies of Praesepe include Hambly et al. (1995b) (hereafter HSHJ) covering 19 sq. degs down to 0.1M, and Adams et al. (2002) covering the whole cluster to similar depth. A number of smaller but deeper Praesepe surveys have also been carried out. Pinfield et al. (1997) (P97 hereafter) covered $``$ 1 sq. deg down to $`I_c`$ = 21. Pinfield et al. (2000) (P00 hereafter) covered $``$ 5 sq. degs down to I$``$19.5. Magazzu et al. (1998) covered 800 sq. arcminutes down to I=21.2, and discovered one candidate member of M9 spectral type (Roque Praesepe1; RPr1 hereafter). If confirmed as a Praesepe member, its estimated mass would be M=0.063–0.084$`M_{}`$ making it a possible BD. NIR characterisation of the faint candidates has been presented by Hodgkin et al. (1999) and Pinfield et al. (2003) (P03 hereafter).
In this paper, we present the results of a deep 2.6 square degree optical survey of the cluster, and use near-infrared follow-up measurements to refine photometric membership status. We then combine our candidates with others from the literature, and identify likely unresolved binaries. We also recover the “M dwarf gap” in the cluster – a dearth of M7-8 dwarfs previously noted in the Pleiades and other clusters (Dobbie et al. 2002; P03). We derive the luminosity function, use state-of-the-art evolutionary models to estimate candidate masses and determine the cluster mass function. We also derive proper motions for several previously identified cluster candidate members. Finally, we discuss planned future work.
## 2 WFC IZ SURVEY
### 2.1 Observations
Deep I and Z band images of Praesepe were obtained with the INT/WFC during the nights 2001/12/24-28. The WFC instrument consists of 4 thinned 2048 x 4196 pixel CCDs (0.333 arcseconds/pixel), and operates at the prime focus of the INT, covering a 0.29 sq. deg field of view.
The conditions for $``$ 35 percent of the time were photometric with modest humidity and seeing of $``$ 1 arcsec. During the run, we observed 9 pointings through the I<sub>rgo</sub> and Z<sub>rgo</sub> filters with exposure times of 20 minutes per band. In total, 2.6 square degrees of the cluster were surveyed. Our coverage is shown in Figure 1.
### 2.2 IZ reduction
Data reduction was performed using standard IRAF<sup>1</sup><sup>1</sup>1IRAF is distributed by the National Optical Astronomy Observatories, which are operated by the Association of of Universities for Research in Astronomy, Inc., under cooperative agreement with the National Science Foundation. routines. The images were bias subtracted, and non-linearity accounted for prior to flat-fielding. Fringe maps, which were constructed in both bands by median filtering images acquired throughout the entire run, were used to remove the effects of the interference between night sky lines in the CCD substrate.
Sources were identified in the I-band images using IRAF’S DAOFIND at a detection level of $`3\sigma `$, and matched to Z\- counterparts using a 4 pixels search radius of the I-band source position (taking into account any systematic shifts between the I\- and Z\- images). Aperture photometry was measured for each source using the PHOT routine with a 7 pixel aperture size and a sky value computed from within a circular annulus.
To calibrate the photometric data, standard stars drawn from Landolt (1992) were observed, and zero points and extinction coefficients determined. In order to calibrate the 7 images obtained in non-photometric conditions, we re-observed these fields with shorter (45s) integrations in the I- and Z- bands on the night of 23rd February 2002 during photometric conditions. Landolt standard observations were used to calibrate the magnitudes of $``$ 50 of the brighter stars in each of these shorter integrations, and of these, the $``$ 40 stars that were not saturated in our longer integrations were used as secondary standards to determine zero points for these longer exposure images.
### 2.3 WFC Results
Figure 2 plots the frequency of the total number of sources detected per 0.1 magnitude interval in the I\- and Z-bands. The dotted lines are straight line fits to the I- and Z- band histograms in the magnitude ranges I=17-20.5 and Z=17-19.5. Survey completeness limits for the I and Z bands were taken to be the magnitudes at which the binned star counts drop below 90 percent of the expected frequency, and are shown as vertical dashed lines. In this way, the survey was estimated to be 90 percent complete to I =21.3, Z=20.5.
Our WFC survey resulted in a catalogue of $``$ 15000 objects. The resultant I,I-Z colour-magnitude diagram (CMD) is shown in Figure 3. In order to separate candidate cluster members from background sources, it was necessary to locate the cluster sequence. To do this, we used 0.5 Gyr DUSTY (Chabrier et al. 2000) and NEXTGEN (Baraffe et al. 1998) model isochrones (dashed and dotted lines respectively), transformed onto our CMD. Predicted I-Z colours were calculated using the relations given in Dobbie et al. (2002), where we assume a distance modulus of (m-M)<sub>0</sub> = 6.16 (Pinsonneault et al. 1998) and zero extinction. To further highlight the location of the cluster sequence, we also identified HSHJ proper motion members and optical/NIR candidates from P03 that were recovered in our survey. We choose not to include objects that fall within 1 degree of the centre of a suspected sub-cluster, discussed in $`\mathrm{\S }`$9. These are highlighted in the figure as plus signs and triangles respectively. Further, we overplot candidates from P03 that are shown to be astrometric cluster members in section $`\mathrm{\S }`$10, and are shown as filled stars.
Also plotted in Figure 3 are the expected magnitudes of several M field dwarfs as they would appear at the cluster distance. Praesepe members are sufficiently old to have finished their contraction stage, implying that cluster UCDs will have the same luminosity as field stars of the same T<sub>eff</sub>. The magnitudes for these are taken from Cossburn et al. (1998), and are shown as squares in the CMD.
Using the photometric and astrometric members, combined with the offset field star positions and the synthesized model tracks as a guide on the CMD, we have defined a cut-off line to separate background sources from potential cluster members. By taking into account photometric uncertainties (indicated on the right hand side of the figure), we have positioned our cut-off line such that it is at least 1-sigma (in I-Z uncertainty) blueward of these photometric guides. Note that for I$`>`$19, the offset field stars are the closest objects to the cut-off line, and the proper motion cluster members are rather redder than the cut-off. This is as expected, since Praesepe low-mass members (Teff$`<`$2500K) will not be fully contracted and we expect lower Teffs as a function of magnitude (e.g. 100K lower for Mi 14 at 0.5Gyr age), and thus slightly larger I-Z colours (e.g. 0.1 larger; Dobbie et al. 2002). The use of the offset field stars thus provides a conservative way to define our cluster region on the CMD. We are thus confident that our selection will not miss any genuine cluster members. We further constrained the selection of candidates using a bright limit (brighter than which 2MASS is sensitive to cluster members which will have been previously identified by Adams et al. (2002) and HSHJ), and a faint completeness limit (both indicated on the CMD). All stellar-like objects redward of the cut-off line, and between the bright and faint limits were selected as IZ candidates, and are shown as filled circles in the figure. We identified 320 cluster candidates in this way from $`I_c`$ = 17.5-21.3. Visual inspection of each revealed 95 to be spurious (mostly diffraction spikes and cosmic ray hits), leaving 225 genuine candidate members.
## 3 NIR FOLLOW UP
Near-infrared follow-up measurements of the 225 WFC candidate members were made using the Fast Track Imager (UFTI) instrument on the United Kingdom Infrared Telescope (UKIRT). UFTI was windowed (512x512 pixels) for faster readout providing a 46.5 arcsecond field of view, with each integration consisting of a 5-point dither pattern. K-band photometry of the 120 brightest candidates was measured during flexibly scheduled queue observing in photometric conditions from the 19th to 24th of February 2003. K-band photometry of the remaining 105 fainter candidates, as well as longer J-, H\- and K-band integrations were obtained by us in photometric conditions during our UKIRT observing run on 7-10 March 2003.
Our observing strategy was to initially measure K-band photometry (S/N$``$20) for all candidates, with exposure times ranging between 6s to 60s depending upon source brightness. At this stage we considered a target as a potential cluster member if it lay near to the NEXTGEN isochrone for K$``$15.5, had I-K$`>`$3.3 for K=15.5–16.5, or had I-K$`>`$3.8 for K$`>`$16.5. We then obtained longer integration J, H and K measurements of 12 of the faintest potential cluster members, where we used exposure times of 24s-120s, aiming for a S/N of $``$ 50. The optical and NIR photometry of all WFC sources flagged as potential candidate members is presented in Table 7.
The images were de-biased, dark subtracted, flat fielded and finally combined into mosaics using the UKIRT software pipeline ORACDR (Bridger et al. 2000). The STARLINK package GAIA was then used to measure the photometry, using apertures matched to the seeing (0.6–1 arcseconds). Individual aperture corrections were calculated using UKIRT faint standards measured either side of the candidate, and instrumental magnitudes were then transformed onto the MKO system using airmass curves derived from these standards.
## 4 ADDITIONAL PRAESEPE CANDIDATES
As well as the new WFC candidates, we also consider Praesepe candidate members from several other surveys. Proper motion candidates from HSHJ define the bright main sequence down to I<sub>c</sub> $``$ 17.5. We have transformed the HSHJ photographic I<sub>N</sub> magnitudes onto the Cousins system using Bessell (1986). Sources from P97 and P00 (which we will refer to as Riz and Iz candidates respectively) were also considered in the present analysis. NIR photometry of the Riz and Iz candidates was taken from P03, and is on the MKO system. We also considered the RPr1 candidate from Magazzu et al. (1998) (the M98 survey hereafter). The NIR photometry of this object has been transformed onto the MKO system using the transformations of Hawarden et al. (2001).
In order to obtain NIR measurements of as many of the HSHJ stars as possible, we cross matched their positions with the 2MASS All Sky catalog. 2MASS photometry typically provides a SNR$``$10 for J = 15.8, H = 15.1, and Ks = 14.3. We identified NIR counterparts within a search radius of 3 arcseconds from each source. Of the 459 HSHJ candidates that have astrometric membership probabilities $``$ 70 percent (from HSHJ), 447 had a corresponding 2MASS counterpart. 3 of these were flagged by 2MASS as having potentially contaminated photometry. These, and the 12 unmatched HSHJ candidates, are noted in Table 8. We transformed the 2MASS photometry onto the MKO system using equations in Carpenter (2001).
## 5 PHOTOMETRIC MEMBERSHIP AND BINARITY
### 5.1 Membership from the IK CMD
Figure 4 shows the K,I-K CMD for Praesepe. Iz and Riz candidates from P97 and P03 are shown as triangles and squares respectively. The new WFC candidates are shown as diamonds. Large diamonds represent WFC candidates that were observed using longer K, J and H integrations. All candidates with photometry which we deem inconsistent with cluster membership (see below) are displayed as empty symbols. RPr1 is shown as an asterisk.
Over-plotted on the CMD are 0.5 Gyr NEXTGEN (solid line) and Dusty (dotted line) model isochrones from the Lyon Group (Baraffe et al. 1998, Chabrier et al. 2000), where we have assumed a distance modulus of (m-M)<sub>o</sub> = 6.16 and zero reddening.
The large number of HSHJ and Iz sources define the sequence well down to K$``$15.5, and it is clear from the figure that the NEXTGEN isochrone fits the cluster sequence fairly well down to this magnitude (out to I-K$``$3.3; as was noted by P03). We therefore adjudged the membership status of all candidates with I-K $``$ 3.3 based on their proximity to the NEXTGEN isochrone. Candidates which appeared to be too blue, despite making allowance for the cluster depth (+/-0.15 mags; Holland et al. 2000) and the photometric uncertainties, were rejected.
All 444 HSHJ objects with 2MASS photometry have I-K $``$ 3.3. Using the NEXTGEN isochrone as a test for cluster membership, we adjudge 18 of these to be non-members. Furthermore, three HSHJ candidates appear significantly redder than the isochrone, with photometry indicative of foreground field dwarfs. However, we caution that these sources are all found on the same Schmidt plate in HSHJ suggesting that their redness could be due to calibration systematic errors, so we choose not to reject them as candidates. In Table 8 we list these 3 HSHJ candidates, the photometric non-members and the HSHJ candidates for which no 2MASS counterpart was found or a 2MASS contamination flag was logged. All other HSHJ candidates have 2MASS photometry consistent with cluster membership.
For I-K $``$ 3.3 (T<sub>eff</sub> $``$ 2500K) the cluster sequence is not obviously defined by the candidates themselves. Also it is not immediately apparent if either the NEXTGEN or DUSTY isochrones are appropriate. This is the T<sub>eff</sub> range when atmospheric dust grains are expected to condense out of the gas phase (Jones & Tsuji 1997), so we might expect the NEXTGEN isochrone to be inadequate. However, the DUSTY isochrone is clearly too blue when one moves fainter than K=15.5, and this isochrone is thus not a viable alternative to guide candidate member selection. Consequently, we chose not to rely on the models to define our Praesepe sequence in this range.
Instead, we defined a cluster sequence region in which to identify candidates as likely members. In order to do this we began by defining the top of the equal mass unresolved binary sequence in the K,I-K CMD, which we assume consists of the brightest candidates in this colour range – these sources then define the bright limit of our cluster sequence region. The single star sequence will be 0.75 magnitudes fainter than the equal mass binary sequence. We allow for the broadening of both the single and binary sequences by cluster depth effects and assume photometric uncertianties of +/-0.2 magnitudes are applicable to both single and binary members. Therefore, as illustrated by the shaded zone in Figure 4, our bright and faint selection criteria are separated by 1.45 magnitudes. We adjudge candidates with I-K $`>`$ 3.3 lying outwith this region as non-members.
To summarise, of the 225 WFC candidates with K-band measurements, 90 look like non members, and 36 have I-K colours consistent with cluster membership (84 percent contamination amongst the candidates from the I, I-Z CMD). Of the 12 WFC candidates for which we measured long integration J, H, and K magnitudes, 3 are now flagged as non members, because although red (and initially considered interesting) they lie below our cluster sequence region in the CMD. Our IK membership criteria are given in Table 7, and all IK members (including those from the other surveys considered) are listed in Table 9.
### 5.2 Identifying Unresolved Binarity
Possible unresolved cluster binaries were identified in the I, I-K CMD by dividing our cluster sequence region in two, using a division that is mid-way between the region’s bright and faint limits. This division is shown in Figure 4 as a dashed line. Candidate cluster members above this line will lie closer to the equal mass binary sequence than to the single star sequence, and have been duly flagged as binary candidates (these sources are circled in Figure 4). Unresolved binarity amongst the Riz and Iz candidates (as well as the HSHJ sample) has been previously addressed by P03. We have nothing significant to add to their analysis of candidate members with I-K$`<`$3.3. However, our current sample of redder candidates is larger and goes fainter than P03. Therefore we have re-assessed the unresolved binarity of any Riz and Iz candidates with I-K$``$3.3 (as well as the WFC candidates) using our current approach. Unresolved probable binaries are flagged in the last column of Table A3.
### 5.3 The JK CMD
The K,J-K CMD is shown in Figure 5. All faint candidates with accurate J-band measurements are plotted. We only include HSHJ sources with K- uncertainties $`<`$0.06 magnitudes, and we do not show Iz candidates with photometry from 2MASS due to the large photometric uncertainties. NEXTGEN and DUSTY model isochrones for 0.5Gyr age are shown as solid and dotted line respectively. It can be seen that the NEXTGEN isochrone agrees well with the cluster members to K $``$ 15.5 (as with the IK CMD). However, for K$`>`$15.5 the candidates move significantly red-ward of the NEXTGEN isochrone. In the main, these fainter candidates lie slightly above the DUSTY isochrone. However, the NIR colour trend shows the rapid increase in J-K that the DUSTY models predict, although it occurs at a slightly brighter magnitude.
Three candidates with J-band photometry were assigned as non-members in §5 (two of which appear on the plot as open symbols). WFC30 and WFC22 lie slightly below the dusty isochrone but have J-K colours that are comparable with the majority of faint candidate members. This is consistent with these two sources being background stars. WFC57 has a very red colour (J-K=1.91) and lies redward of the plotting area. Although this is suggestive of a late L dwarf, the I-K colour of 3.88 is not red enough to be a late L dwarf (cf. I-K=4–5). Consequently we believe this source is a red galaxy. Iz4 is also highlighted in Figure 5, and appears to look like a bright foreground field dwarf, but the possibility of it being an unresolved binary cannot be ruled out. Thus at this stage we do not discard it as a contaminant.
The approximate M/L transition colour is indicated at the bottom of the plot (estimated using colours from Leggett et al. 1998 and Knapp et al. 2004), and shows that most of our faint Praesepe candidate members should be very late M dwarfs. At least two however (WFC 11 and WFC60), have colours consistent with being early-mid L spectral types. These are labelled in Figure 5. WFC60 was flagged as an unresolved binary in §5.2, and in the JK CMD it is at least 0.5 magnitudes brighter than candidate members with similar J-K colours. This candidate could be an unresolved binary L dwarf in the cluster.
### 5.4 Contamination
In this section we estimate the expected level of field contamination at the faint end of the Praesepe sequence. We do this for three K- band magnitude ranges: 15.5-16.0, 16.0-16.5 and 16.5-17.0. In the cluster member region of the IK CMD, it can be seen that these ranges correspond to I-K colours of $``$ 3.5, 4.0 and 4.4 respectively.
These colours relate to spectral types $``$ M6, M8 and L0 respectively. The absolute magnitudes of field stars with these spectral types (from Dahn et al. 2002 and Knapp et al. 2004) means those at 160-200 pc will overlap the cluster sequence in a CMD.
The sky area covered for each magnitude range varies, since the different surveys considered have different photometric depth. All of the surveys considered (excepting the HSHJ survey) contribute to the brightest magnitude range. However, the Iz survey does not contribute to the two fainter ranges. All except the HSHJ and Iz surveys contribute to the middle magnitude range, and only the WFC and M98 surveys contribute to the faintest magnitude range. After accounting for overlap between certain survey areas, we find that the brightest magnitude range covers a volume of $``$ 1.2 x 10<sup>3</sup> pc<sup>3</sup>, and the two fainter ranges each cover a volume of 1.0 x 10<sup>3</sup> pc<sup>3</sup>.
Luminosity functions predict late M and early L field dwarfs densities of $`\mathrm{\Phi }`$ = $``$ $`2.5`$ x 10<sup>-3</sup> and $`1.9`$ x $`10^3`$ stars $`pc^3`$ (Kirkpatrick et al. (1994) and Cruz et al. (2003) respectively). We hence expect a field contamination of $``$ 3 late M dwarfs and $``$ 2 early L dwarfs amongst the candidates.
While clearly the levels of contamination amongst our candidate members with the late M spectral types is fairly low ($``$20%), it is entirely possible that both of our candidate L type members could be field stars. A more rigorous assement of their membership status can be performed by measuring their proper motions.
### 5.5 J-H,H-K 2 colour diagram
Figure 6 shows our candidates in the J-H,H-K 2-colour diagram. The symbols are the same as in the CMDs. NEXTGEN (dashed lines) and DUSTY (solid lines) model colours are over-plotted, with the surface gravity values shown. The bluer synthetic H-K colours for H-K$`<`$0.2 are thought to be due to inaccuracies with water vapour modelling in the H band. Approximate late M and L spectral class locations are also indicated (from Leggett et al. (1998) and Knapp et al. 2004).
Visual inspection of the diagram shows that, as expected, the 2 likely background stars (WFC 22 and 30) have colours consistent with late M dwarfs. WFC57, probably a galaxy, has NIR colours similar to a late L dwarf. The possible unresolved L dwarf binary (WFC60) and the even redder single candidate member (WFC11) have NIR colours consistent with being early-mid L dwarfs. Similarly, Riz117 appears too blue for an early L dwarf, and is therefore considered as a probable non-member. It is shown in Figure 6 as an open square.
## 6 THE PRAESEPE SEQUENCE
### 6.1 M dwarf gap
In Figure 7 we show a zoomed in portion of the Praesepe IK CMD with non-members omitted. The four binary tracks (shown with lines joined by plus signs) will be discussed in §8.1.
It can be seen that there is a dearth of cluster candidates from K=15.3-15.5. Dobbie et al. (2002) and P03 presented photometric evidence for a similar paucity of M7-M8 spectral types in other clusters, as well as in star forming regions and the field. This paucity has been referred to as the “M dwarf gap” (P03). The fact that the M dwarf gap has been observed in both young associations such as the Pleiades and $`\theta `$ Orionis as well as amongst older populations (M4 and the field), suggests that its origins are due to a sharp local increase in the slope of the magnitude-mass relation for late-M spectral types ($`T_{eff}2700K`$). At this T<sub>eff</sub>, models predict that dust grains begin to condense in the outer atmospheric layers (e.g. Tsuji et al. 1996), and as Dobbie et al. (2002) suggests, grain opacities could be responsible for the gap feature. The Dusty models do not consider grains with sizes greater than $``$ 0.24 $`\mu m`$, and consequently predict negligible grain opacities in the NIR, and no gap. Grain formation models however, suggest that grains could have sizes up to tens of microns (Cooper et al. 2003), which could result in a significant opacity increase due to Mie scattering.
We have measured the extent of Praesepe’s “M dwarf gap” in the I-,J-,H- and K- bands. For each band, the gap size was calculated as the magnitude difference between the faintest candidate above the gap, and the brightest candidate below, ensuring that only the candidates with full IJHK photometry were used. The results are tabulated (with the corresponding values for the Pleiades from P03) in Table 1. For all passbands, quoted errors were derived by combining the photometric uncertainties of the two candidates that were used to measure the gap sizes.
It can be seen that the gap is smaller for Praesepe than for the Pleiades in all the bands considered, although there is a small overlap if the uncertainties are taken into account. This is consistent with the interpretation just mentioned, since (as Dobbie et al. (2002) explains) one expects a shallower luminosity-T<sub>eff</sub> relation for older populations (such as Praesepe), because the radius-mass relation will be steeper for younger populations (like the Pleiades), where objects are still in the relatively early stages of contraction. The relative extent of the Praesepe M dwarf gap in the four photometric bands considered is comparable for both clusters. We will, however, discuss more subtle colour changes across the gap in the next section.
### 6.2 The faint single star sequence
In order to study how the bulk properties of the Praesepe candidates change for $`K>`$15, we split the Praesepe candidate members into the following K-band magnitude bins: bin A (15.0 – 15.3) represents the candidates slightly brighter than the M dwarf gap (discussed in the last section). Bins B and C (15.5–15.75 and 15.75–16.00 respectively), are below the gap. Bin D covers all the fainter candidates (except WFC 11) down to K =16.8. These 4 bins are indicated by dashed lines in Figure 7.
In order to determine an “observational photometric sequence”, we used the NEXTGEN model isochrone above the gap region in the CMD (since this model agrees well with the data in this range). To bridge the M dwarf gap we have used our measured gap sizes (see Table 1) to provide colour and magnitude changes from K=15.3 – 15.5. Below the gap (bins B–D), we have averaged the colours and magnitudes of our single star candidate members in each bin. Our “observed cluster sequence” is shown as a solid line in Figure 7.
When comparing the observed sequence to the model isochrones, it is apparent that for I-K=3.3–4 the NEXTGEN model is still quite reasonable. However, for I-K$`>`$4 the candidates appear to drop slightly below the NEXTGEN isochrone, moving closer to the DUSTY isochrone. P03 and Jameson et al. (2002) pointed out a similar transition from the NEXTGEN to DUSTY isochrones for the Pleiades cluster sequence from I-K=4.1–4.5.
We trace the colour changes of the single source candidates (ie. avoiding sources flagged as unresolved binaries) through the 4 K-band magnitude bins by indicating where the bin members lie (with boxes) on a zoomed in version of the 2-colour diagram (shown in Figure 8). The boxes were constructed such that single candidate members from each of the K-band magnitude bins are all contained within the respective box dimensions in the diagram.
Bin A, which is populated by candidates above the M dwarf gap, appears to be slightly bluer in (J-H) than the NEXTGEN model colours, lying closer to the DUSTY model colours. We expect no dust condensation at these T<sub>eff</sub>s, so this difference probably results from incompleteness in the NEXTGEN H<sub>2</sub>O opacities (eg. see Jones et al. 2003). Bin B is slightly redder in H-K than Bin A, but has similar J-H colours. This colour change could be caused by the addition of new grain opacities, consistent with the appearance of atmospheric dust across the M dwarf gap. However, note that no colour changes were observed across the M dwarf gap in the Pleiades (see P03). If the M dwarf gap does result from the first appearance of atmospheric dust grains, then they must have different effects for different surface gravities. This would not be unexpected since we would expect grain growth and rain-out to be affected by both surface gravity and pressure. The fainter bins (C and D) have colours consistent with the Dusty models.
## 7 THE MASSES OF LOW MASS CANDIDATES
We have estimated the masses of all single Praesepe candidates that appear in our candidate sequence region (see Figure 4). The exception to this is Iz117, which was identified in $`\mathrm{\S }`$5.5 as a likely non-member. Masses were estimated using the NEXTGEN and Dusty models (with (m-M)<sub>0</sub>=6.16) for isochrone ages of 0.5 and 1 Gyrs, and solar metallicity. Martín et al. (2000) note that mass estimates from J- band magnitudes are the least sensitive to the choice of model, either NEXTGEN or DUSTY. Therefore, we choose to derive masses from this passband where possible. Where J- band photometry was absent, masses were estimated from the K- band magnitudes.
Linear interpolation was used in between the isochrone mass points. The DUSTY mass estimates are slightly higher than the NEXTGEN ones, and, rather more obviously, the younger mass estimates are lower than the older ones. The mass estimates are given in Table 9. Estimates taken from the K- band are slightly higher than if taken from the J- band.
It can be seen from the table that we have identified at least 15 single VLMS candidates (with masses $``$ 0.075 M). The five faintest candidates could be sub-stellar or stellar depending on the age or model we adopt. For example, four of these are substellar if a 0.5 Gyr cluster age is assumed, regardless of which model is used.
## 8 PRAESEPE BINARIES
### 8.1 Binary mass ratios
In Figure 7 we have over-plotted 6 binary tracks, which we use to estimate mass ratios of our Praesepe binaries. We calculated these tracks by combining single star points from the cluster sequence. Each binary track begins at some single star point, and the affect of adding an unresolved companion is simulated by combining the photometric brightness of this single star point with other single star points from lower mass points in the sequence (beginning with the lowest mass point on the sequence, and then increasing the companion mass until an equal mass binary is reached, 0.75 magnitudes above the original single star point).
When estimating binary mass ratios (q), one must account for the cluster depth which can make sources appear brighter without the need for a binary companion, as well as photometric uncertainty. The tidal radius of Praesepe is 12.1 pc (Holland et al. 2000) giving a depth effect $``$ $`\pm `$ 0.15 magnitudes. Photometric uncertainty is typically $`\pm `$0.2 in I-K. Mass ratios were estimated by comparing unresolved binary candidates to the binary tracks in the CMD, and allowing for both these forms of uncertainty. The results are tabulated in Table 2.
The q values range from 0.4–1.0. This represents a rather larger range of q values than was found for Pleiades BDs by P03 (q=0.7–1.0). However, the photometric uncertainties are slightly larger for the Praesepe candidates, and the very low mass unresolved binary populations identified in the two clusters cannot be shown to be significantly different (in their q range) by this analysis.
### 8.2 The binary fraction
P03 determined binary fractions (BFs; defined as the number of binary systems divided by the total number of systems in some photometric range) for Praesepe in 4 colour bins for 1 $``$ I-K $``$ 3.6, corresponding to an approximate mass range 1.0 - 0.09 M. We extend this to lower masses by considering two redder bins from I-K =3.3-3.9 and 3.9-4.6. For the first bin the Iz survey is not complete over the full single star sequence. Therefore, to avoid underestimating the relative number of single star sources we only considered Riz and WFC candidates. With 5 unresolved binary and 11 single source candidates, this bin has a BF$``$31%. For the second bin, the Riz survey is not complete over the full single star sequence. However, the Riz candidates are all contained within the WFC survey area, and would have been discovered by this survey had they not already been identified. We therefore considered all sources in this bin, resulting in 3 unresolved binary and 7 single source candidates, or a BF$``$30%. These results are summarised in Table 3, where uncertainties have been calculated assuming binomial statistics (Burgasser et al. 2003). Note that the low number statistics give large uncertainties, and that the BF for the redder bin is a lot more uncertain even than this, since several of these candidates could be non-members (see Sections 5.4 and 10.4).
The Praesepe BFs are displayed in Figure 9, with the 4 higher mass BFs from P03 over-plotted. Also shown are the BFs from Close et al. (2003) and Burgasser et al. (2003) for VLMS and BD binaries in the field with separations $`>`$1AU. Although our lowest mass point is too uncertain to accurately constrain the BF, the higher mass point shows that the BF is decreasing as one goes below 0.1M. Even if our expected level of contamination ($``$3 late M dwarfs) removes only single star candidates (thus pushing up the BF), the BF would only go up to 39$`{}_{12}{}^{}{}_{}{}^{+14}`$%, which would still suggest a BF decrease. This decreasing BF trend appears to be consistent with the lower BFs seen for field sources. P03 however, observed an increasing unresolved BF into the Pleiades BD regime, reaching 50$`{}_{11}{}^{}{}_{}{}^{+10}`$ for 0.07M. This suggests that different clusters could have different very low mass BFs. Such differences could result from different cluster environments, and in this context the distribution of different cluster BFs could provide an important observational test of VLMS and BD formation theories. However, before any firm conclusions are drawn, it is clearly important to improve the robustness of the measured cluster BFs by confirming (or not) cluster membership from proper motions.
## 9 LUMINOSITY AND MASS FUNCTION
We divided our cluster members into 3 magnitude ranges covering K=15–15.8, 15.8–16.3 and 16.3–17. For each magnitude range, cluster candidates were counted in 0.5 degree rings from the centre, out to a radius of 2.5 degrees for the brightest bin, and out to 2.0 degrees for the 2 fainter bins. We did not consider candidates found in the region of the sub-cluster identified by Holland et al. (2000), 3 pc from the centre of Praesepe. This sub-cluster could be a smaller older open cluster that has collided with Praesepe, and we avoided candidates within 1 degree of the sub-cluster centre. The positions of candidates in the magnitude bins K=15–15.8, 15.8–16.3 and 16.3–17.0 are shown in figure 10 as crosses, squares and asterisks respectively. It can be seen from the figure that there is an over-density of faint objects in the sub-cluster region. It is clearly important to isolate (and ignore) this region when measuring the LF and MFs of the main cluster, since the sub-cluster is thought to be older than the rest of Praesepe (Holland et al. 2000), and it would otherwise be possible for faint low-mass sub-cluster stars to masquerade as Praesepe BDs.
Ring counts were then converted into surface densities by dividing by the appropriate survey areas within each ring (ignoring the sub-cluster region), and expected levels of field star contamination (see Section 5.4) were subtracted off the surface density profiles assuming a uniform spatial distribution over the cluster. We then assumed that the cluster could be represented by a King surface density distribution (King 1962) (as has been done previously by Pinfield et al. 2000, Jameson et al. 2002 and Holland et al. 2000). This function takes the form:
$$f_s=k\left\{\frac{1}{\sqrt{1+(r/r_c)^2}}\frac{1}{\sqrt{1+(r_t/r_c)^2}}\right\}^2$$
(1)
Here, $`f_s`$ is the surface density, r the radius from the cluster centre, k a normalization constant and $`r_c`$ and $`r_t`$ the core and tidal radii respectively. Assuming $`r_t`$ = 12.1 pc (Holland et al. 2000), and using Poisson uncertainties for the surface density profiles, we obtained the best fit value of $`r_c`$ using the whole sample, and found $`r_c`$=7.56 pc. This was then assumed to be constant for each magnitude bin, so we proceeded to minimize the $`\chi ^2`$ statistic of the King profile to the 3 observed density distributions. We then determined the total number of cluster stars in each bin using the surface integral (out to $`r_t`$) of equation 1;
$$n=\pi r_c^2K\left\{ln(1+x_t)\frac{(3\sqrt{1+x_t}1)(\sqrt{1+x_t}1)}{1+x_t}\right\}$$
(2)
where $`x_t=(r_t/r_c)^2`$. Our results are summarised in Table 4. The resulting luminosity function ($`\mathrm{\Phi }_K`$) is shown in figure 11 as filled diamonds. Also shown is the luminosity function from HSHJa (converted into $`\mathrm{\Phi }_K`$), although note that the faintest three HSHJ points (open circles) are lower limits due to photometric incompleteness.
The mass function then follows directly from the luminosity function:
$$N(m)=\mathrm{\Phi }_K\frac{dK}{dm}.$$
(3)
We estimated $`\frac{dK}{dm}`$ for our 3 magnitude ranges using both the 0.5 Gyr and 1.0 Gyr NEXTGEN and DUSTY isochrones to define the mass limits of each bin, and derived the cluster mass function per 0.1M interval (see Table 4).
The log-normal Pleiades mass function as derived by Adams et al. (2001) is displayed on Figure 12 as a dashed line, and is normalised to the Praesepe function at 0.5 M. An upper limit Hyades mass function point from Dobbie et al. (2002) in the range 0.08-0.06 M is also shown for comparison.
It can be seen from Figure 11 that both the Praesepe mass function derived here and the Hyades mass function falls below that of the Pleiades at $``$0.1M. The Hyades and Praesepe are thought to be of a similar age ($``$ 0.6 Gyr), compared to the younger 0.13 Gyr Pleiades cluster (Burke et al. 2004). This age gap has lead Dobbie et al. (2002) to suggest that the paucity of low-mass Hyads is caused by dynamical evolution of the environment. In this scenario, the lowest-mass members of an association are, over time, preferentially ejected from the cluster – a theory supported by the N–body simulations of de la Fuente Marcos & de la Fuente Marcos (2000) in which a 0.6 Gyr open cluster has retained at best only $``$ 17 percent of the BDs that were originally bound to the association. Evidently, our Praesepe mass function is in strong agreement with this scenario.
Despite the qualitative similarities, however, between the mass functions of Praesepe and Hyades, there are to date no Hyad substellar candidates, in comparison to the Praesepe candidates proposed in this work. As a result, the Hyades low-mass MF upper limit is only marginally consistent with our lowest mass Praesepe point, which suggests that the Praesepe MF may be somewhat higher than that of the Hyades in this mass range. Clearly, differing dynamical processes would not be an issue here since both clusters are of a similar age and can be assumed therefore to be at a similar stage of evolution. If this is a real feature of the MF, rather than as a consequence of low-number statistics, we speculate that this difference can perhaps be explained by a substellar mass function sensitivity to natal molecular cloud properties such as density, metallicity, temperature etc. Indeed, the simulations of Elmegreen (2000) suggest that the shape of the MF below the observed 0.1 M turnover is dependent upon the physical cloud properties, in contrast to the universal Salpeter IMF that characterises the intermediate to high mass domain. We stress, though, that it is imperative to increase the known substellar population in a variety of associations before this hypothesis can be fully tested, and that the possibility of a common substellar MF in both environments cannot be ruled out.
## 10 ASTROMETRY
We have mentioned, in previous sections, how important proper motions are in confirming cluster membership. The peculiar motion of Praesepe is large compared to that of typical field stars, and the intrinsic velocity dispersion of the cluster is small (HSHJb). With a base-line of a few years, it is thus possible to measure the proper motion of candidates with sufficient accuracy to separate true cluster members from field star contamination. In this section we describe how we have measured proper motions for a subset of the Praesepe candidates we consider in this work.
### 10.1 Epoch data
First epoch data for the astrometric analysis of Praesepe RIZ candidates was taken in 1993 using the 2.5 m Isaac Newton Telescope (INT), providing a baseline of $``$ 8.9 years. For Roque PR1, we use the original image acquired in February 1996, also taken with the INT, with a baseline of $``$ 5.9 years. We use our WFC images as 2nd epoch data. The images in which RIZ 21 and 23 were recovered were acquired on the third night of the observing run, which could not be photometrically calibrated due to poor conditions. Consequently, these two candidates are not plotted in Figure 3. However, we do not expect this to affect the astrometric measurements. For all images, we use only the I band for the astrometry.
### 10.2 Measurements
The process of deriving the astrometry of the cluster candidates is as follows. The pixel positions of the candidates on both epoch images ($`x_{1st}`$,$`y_{1st}`$),($`x_{2nd}`$,$`y_{2nd}`$) are first accurately determined using the IRAF:CENTRE routine. The positions of presumed background point sources were also measured in both epochs, to be used as astrometric references. IRAF:XYXYMATCH was used to calculate the transformation function which maps the two epoch images. GEOMAP and GEOXYTRAN were then used to transpose the second epoch coordinates of the candidates into the first epoch coordinate system. Any reference sources lying more than a Gaussian 3 $`\sigma `$ from the median of absolute deviation were rejected and the transform redetermined.
The relative proper motions along the x and y axes $`\mu _{x,y}`$ of the CCD then follows from the pixel displacement:
$$\mu _X=\frac{(\mathrm{\Delta }X)}{\mathrm{\Delta }t}=\frac{X_{2nd}X_{1st}}{\mathrm{\Delta }t}$$
(4)
$$\mu _Y=\frac{(\mathrm{\Delta }Y)}{\mathrm{\Delta }t}=\frac{Y_{2nd}Y_{1st}}{\mathrm{\Delta }t}$$
(5)
The uncertainties of the transformations $`\sigma _{x,y}`$ were found to be $``$ 0.2 pixels, corresponding to proper motion errors of $``$ 10-15 masyr<sup>-1</sup>, in both RA and Dec direction.
### 10.3 Membership probabilities
The method of establishing membership probabilities is that outlined by Sanders (1971), and utilizes the following likelihood function:
$$L=\underset{i=1}{\overset{N}{}}\mathrm{\Phi }(\mu _x,\mu _y)$$
(6)
where
$$\mathrm{\Phi }(\mu _x,\mu _y)=\frac{n_c}{2\pi \sigma ^2}\mathrm{\Phi }_c(\mu _x,\mu _y)+\frac{1n_c}{2\pi \mathrm{\Sigma }_x\mathrm{\Sigma }_y(1\rho ^2)}\mathrm{\Phi }_f(\mu _x,\mu _y)$$
(7)
here, $`n_c`$ is the normalized number of cluster stars, and $`\rho `$ is a correlation coefficient. Following the method of Wang et al. (1995), $`\sigma `$ follows from the proper motion uncertainties $`\sigma _{x,y}`$ and an intrinsic variance $`\sigma _0`$ :
$$\sigma ^2=\sigma _0^2+\sigma _{x,y}^2$$
(8)
The distribution functions of the cluster and field $`\mathrm{\Phi }_c,\mathrm{\Phi }_f`$ are evaluated as:
$$\mathrm{\Phi }_c(\mu _x,\mu _y)=exp\left\{\frac{1}{2}\left[\left(\frac{\mu _x\mu _{xc}}{\sigma }\right)^2+\left(\frac{\mu _y\mu _{yc}}{\sigma }\right)^2\right]\right\}$$
(9)
and
$$\begin{array}{c}\mathrm{\Phi }_f(\mu _x,\mu _y)=exp\{\frac{1}{2(1\rho ^2)}[\left(\frac{\mu _x\mu _{xf}}{\mathrm{\Sigma }_x}\right)^2\hfill \\ \hfill \frac{2\rho (\mu _x\mu _{yf})(\mu _y\mu _{yf})}{\mathrm{\Sigma }_x\mathrm{\Sigma }_y}+\left(\frac{\mu _y\mu _{yf}}{\mathrm{\Sigma }_y}\right)^2]\}\end{array}$$
(10)
here, $`\sigma `$ is the proper motion dispersion, $`\mu _{xc}`$,$`\mu _{yc}`$ are the cluster mean proper motions, and $`\mu _{xf}`$,$`\mu _{yf}`$ are the field mean proper motions with standard deviations $`\mathrm{\Sigma }_x`$ and $`\mathrm{\Sigma }_y`$.
The free parameters were found by Wang et al. (1995) using a conjugate gradient method using a large Praesepe census, and are reproduced in Table 5.
### 10.4 Astrometric results
Proper motions and membership probabilities of RIZ objects from P97 that have JHK colours consistent with cluster memberships are shown in Table 6 and illustrated in a vector point diagram (VPD) in Figure 13 as filled stars with error bars. We were unable to measure the proper motion of RIZ 4,11,18 and 24 due to insufficient shared spatial coverage between epochs, which would have resulted in unreliable transformations.
Proper motions from Wang et al. (1995) are overplotted as small circles, or large circles where membership probabilities were evaluated to be greater than 70 percent. These help to clarify the cluster position at $`\mu _{x,y}`$ -33,-16 $`masyr^1`$. All stars that were used as references for the astrometric transformation are also shown on the VPD as asterisks. As expected, they cluster around the origin of the diagram, since they were selected on the basis of having a negligible motion.
All of the RIZ candidates have an evaluated membership probability greater than 70 percent, and are therefore consistant with cluster membership within the astrometric uncertainties. Our astrometric analysis of Roque Pr1, shown as a triangle on the VPD, suggests that this object is most likely a non-cluster member. RPr1 is more likely to be a field star in the line of sight to the cluster.
## 11 CONCLUSIONS & FUTURE WORK
In this paper, we have presented the results of a deep optical survey (I$``$21.3) and NIR follow-up of 2.6 square degrees of the Praesepe open cluster. We have used our census of cool red members to place constraints on the bottom end of the cluster mass function. Our results indicate that the dynamical evolution of the cluster is likely leading to the evaporation of the lowest mass members. However, the relative depletion of the very-low mass stellar and substellar members may be less than that reported for the similarly aged Hyades cluster, suggesting that perhaps the cluster MF is sensitive to natal molecular cloud properties.
To test the robustness of these arguments, it is imperative to ascertain beyond doubt the validity of the candidates’ claims to cluster membership. We have done this using proper motion measurements for several low-mass Praesepe candidates, but our priority is to continue this for all. Only then can we probe with confidence the substellar MF. Furthermore, a bona-fide census would strengthen the comparison between low-mass binary fractions in this cluster compared to earlier clusters such as the Pleiades. Our results suggest a difference between the two environments, but its cause or quantitative nature cannot be elaborated upon without further census expansion.
The photometric sample we have compiled have revealed a paucity of candidates close to the region where the main sequence begins to deviate away from the NEXTGEN track. This feature has been observed in younger regions too, and would appear to originate because of a steepening in the mass-luminosity relation at late M spectral type. We propose that further measurements, both photometric (e.g L- band) and spectroscopic, be made of candidates which are located just above and below the gap in both this cluster and others to isolate the underlying physics, and to test the sensitivity of dust formation in the atmosphere of early cluster L dwarfs to the fundamental properties of their association.
Finally, to further the results presented in this work it is of utmost importance to increase the known number of VLMS and BDs in young SFRs and open clusters. In the near future, this need will primarily be met by the advent of automated deep surveys conducted over wide areas. One such effort is UKIDSS, the collective name for a series of large-scale IR surveys using the WFCAM instrument on the UK Infrared Telescope. The Galactic Cluster Survey (GCS) in particular will survey 10 associations over the next few years on an unprecedented scale, uncovering many more cluster BDs that will go a long way to addressing some of the issues raised in this work.
## 12 Acknowledgments
RJC acknowledges financial support from The Particle Physics and Astronomy Research Council (PPARC) of the UK. PDD is sponsored by a PPARC postdoctoral grant. This publication makes use of data products from the Two Micron All Sky Survey, which is a joint project of the University of Massachusetts and the Infrared Processing and Analysis Center/California Institute of Technology, funded by the National Aeronautics and Space Administration and the National Science Foundation. The United Kingdom Infrared Telescope is operated by the Joint Astronomy Centre on behalf of the UK Particle Physics and Astronomy Research Council. Some of the data reported here were obtained as part of the UKIRT Service Programme.
## Appendix A Tables of photometry, membership and binarity criteria, and mass estimates |
warning/0506/cond-mat0506317.html | ar5iv | text | # Dependence of the superconducting effective mass on doping in cuprates
## 1 Introduction
The superconducting condensate density ($`n_s`$) characterizes the order parameter of this macroscopic quantum commonwealth. The way to cuprate high-T<sub>c</sub> superconductivity includes (necessarily) a doping treatment. The nature of the reorganizations in the physical-chemical basis of the material and the variation of the superconducting properties on the doping scale become of primarily interest. At present there remain nevertheless some debatable aspects in the cuprate superconductivity until the pairing mechanism itself. Correspondingly this concerns also the behaviour of the superconducting condensate density on doping, e.g. .
An essential property which determines numerous applications of superconductors is the penetration depth. It contains $`n_s`$ with the superconducting condensate effective mass $`m`$ in the combination $`m/n_s`$. The doping dependence of the superconductivity playground CuO<sub>2</sub> plane $`m_{ab}`$ in cuprates is practically unknown. Usually one supposes a constant value $`mxm_0`$, with $`x`$ between, say, 2 and 5 ($`m_0`$ is the free electron mass). The main aim of the present contribution is the calculation of $`m_{ab}`$ on the whole scale of the hole doping ($`p`$) in association with other cuprate superconductivity characteristics.
We use a very simple, partly postulative model of a ”typical” cuprate superconductor which uses only general knowledge on the system. The model has been started by Ref. and developed in \[5-7\]. The following comparison of the outcome of the model for various properties with the observations is expected to illuminate in some extent the background physics.
The model supports on the two-component scenario of cuprate superconductivity which states the essential functioning of the doping-created defect subsystem besides the itinerant one. This means that the background electron spectrum shows essential dynamics under doping. Bare normal state gaps are assumed between the mentioned subsystems and supposed to be quenched by progressive doping. Overlap dynamics of the bands appears now as a novel source of critical doping concentrations. The nature of the minimal quasiparticle excitation energies changes with doping in accordance with the chemical potential position. This explains naturally the presence of pseudogaps in the model. The pseudogaps appear as precursors on the doping scale (not on the energetic one) to the superconducting gaps and survive for $`T>T_c`$ as normal state gaps.
The pairing interaction is supposed to be of interband pair transfer type between the itinerant and defect states. This mechanism is seemingly the most effective in serving high transition temperatures in a simple way .
The calculated phase diagram of cuprate energetic characteristics ($`T_c`$, superconducting- and pseudogaps, the condensation energy) agree qualitatively well with the experimental results. Controversial statements in the literature on interrelations and coexistence of various gaps in distinct doping regions have been elucidated. Compounds with two or one pseudogap in the charge channel can described.
## 2 The physical model
We proceed with the description of the physical content of the model background used. It enables the understanding of the behaviour of the characteristics investigated. The cuprate electron spectrum created and reorganized by doping is chosen as follows. The mainly oxygen itinerant valence band ($`\gamma `$) top fixes the energy zero and this band extends until $`\xi =D`$. The states in it are normalized to $`1c`$, where $`c`$ is a measure of the doped hole concentration ($`p`$). The corresponding scaling for a given case must be made by joining a characteristic concentration on the phase diagram. There is a huge (e.g. \[13-17\]) amount of appointments that doping creates new defect (midband) states near the top of the valence band in the charge-transfer gap. Extended doping brings them to merge with the valence band. The functioning of the defect subsystem is anisotropic in the momentum space . At least the ”hot” ($`\pi ,0`$)-type and ”cold” $`(\frac{\pi }{2},\frac{\pi }{2})`$-type regions must be distinguished. Accordingly we introduce two defect system subbands characterized by energy intervals $`d_1\alpha c`$ and $`d_2\beta c`$, i.e. they expand down from energies at $`c=0`$. The overlap of these bands with the valence band is reached at $`c_\alpha =d_1\alpha ^1`$ and $`c_\beta =d_2\beta ^1`$. Note that the infrared manifestations of the defect subsystem are suppressed in favour of the free carriers (Drude peak appearance) with progressive doping . We take $`d_1`$ and $`d_2`$ to be positive: the optical charge-transfer gap is reduced by doping . The choice $`c_\beta <c_\alpha `$ accounts for that the lowest doping-created states belong to the cold subsystem. The weight of the defect states is taken to be $`c/2`$, cf. .
The $`2D`$ (CuO<sub>2</sub> planes) densities of states in the bands read $`\rho _\alpha =(2\alpha )^1`$, $`\rho _\beta =(2\beta )^1`$, $`\rho _\gamma =(1c)D^1`$. There are the following different arrangements of the bands and the chemical potential ($`\mu `$). At very underdoping $`c<c_\beta `$, $`\mu _1=d_2\beta c`$ is connected with the ”cold” $`\beta `$-band and charge carriers become concentrated here, cf. . For $`c>c_\beta `$, $`\mu _2=(d_2\beta c)[1+2\beta (1c)D^1]^1`$ intersects both ($`\beta ,\gamma `$)-bands. The Fermi level shifts into the valence band, cf. . The overlap of the narrow defect $`\beta `$-band with the itinerant band leads to the formation of two sheets of the Fermi surface. The one at $`(\frac{\pi }{2},\frac{\pi }{2})`$ is holelike with the dominating itinerant contribution and the other with a tendency to form an electronlike ”flat band” with lowering $`\mu `$. This behaviour is in agreement with the observation of two Fermi-sheets and of a single ”hole barrel” at $`(\frac{\pi }{2},\frac{\pi }{2})`$ and a flatband at ($`\pi ,0`$) .
For the expressed dopings larger than $`c_0`$, determined by $`d_1\alpha c_0=\mu _2`$, the role of the ($`\pi ,0`$)-type region increases essentially as also Refs. state. Now $`\mu _3=[\alpha d_2+\beta d_12\alpha \beta c][\alpha +\beta +(1c)2\alpha \beta D^1]^1`$ intersects all three overlapping bands. For extended overdoping $`c>c_1`$, where $`c_1`$ is defined by $`d_2\beta c=\mu _3`$, the chemical potential $`\mu _4=(d_1\alpha c)[1+2\alpha (1c)D^1]^1`$ falls out of the cold defect band. The chemical potential is not affected significantly by pairing and its general trend agrees with the results of special investigations \[27-29\].
## 3 Necessary formulas
Our basic Hamiltonian with the coupling ($`W`$) of itinerant and defect subsystems by the pair-transfer interaction reads
$$H=\underset{\sigma ,\stackrel{}{k},s}{}ϵ_\sigma (\stackrel{}{k})a_{\sigma ,\stackrel{}{k},s}^+a_{\sigma ,\stackrel{}{k},s}+W\underset{\sigma ,\sigma ^{}}{}{}_{}{}^{}\underset{\stackrel{}{k},\stackrel{}{k}^{}}{}\underset{\stackrel{}{q}}{}a_{\sigma \stackrel{}{k}}^+a_{\sigma (\stackrel{}{k}+\stackrel{}{q})}^+a_{\sigma ^{}(\stackrel{}{k}^{}+\stackrel{}{q})}a_{\sigma ^{}\stackrel{}{k}^{}}.$$
(1)
Here $`ϵ_\alpha =\xi _\sigma \mu `$, $`s`$ is the spin index, $`\sigma `$ counts the bands and $`\stackrel{}{q}`$ is the pair momentum with the components from the same bands. Various aspects of the work with (1) at $`\stackrel{}{q}=0`$ for calculating the superconductivity energetic characteristics can be followed in . The superconductivity gap parameters are defined as
$`\mathrm{\Delta }_\gamma `$ $`=`$ $`2W{\displaystyle \underset{\stackrel{}{k},\tau }{}}{}_{}{}^{\tau }<a_{\tau \stackrel{}{k}}a_{\tau \stackrel{}{k}}>,`$ (2)
$`\mathrm{\Delta }_\tau `$ $`=`$ $`2W{\displaystyle \underset{\stackrel{}{k}}{}}<a_{\sigma \stackrel{}{k}}a_{\sigma \stackrel{}{k}}>,`$
where $`^\tau `$ means the integration with the densities $`\rho _{\alpha ,\beta }`$ over the corresponding energy intervals of the defect system subbands ($`\tau =\alpha `$, $`\beta `$; $`\mathrm{\Delta }_\alpha =\mathrm{\Delta }_\beta `$). The nongapped nature of the cold subsystem can be accounted by the multiplication of $`\mathrm{\Delta }_\alpha `$ with the suitable $`d`$-symmetry factor.
The gap equation reads ($`\theta =k_BT`$)
$`\mathrm{\Delta }_\sigma =W{\displaystyle \underset{\stackrel{}{k},\tau }{}}{}_{}{}^{\tau }\mathrm{\Delta }_{\tau }^{}(\stackrel{}{k})E_\tau ^1(\stackrel{}{k})th{\displaystyle \frac{E_\tau (\stackrel{}{k})}{2\theta }}`$ (3)
$`\mathrm{\Delta }_\tau =W{\displaystyle \underset{\stackrel{}{k}}{}}\mathrm{\Delta }_\gamma (\stackrel{}{k})E_\gamma ^1(\stackrel{}{k})th{\displaystyle \frac{E_\gamma (\stackrel{}{k})}{2\theta }}`$
with the usual form of the quasiparticle energies $`E_\sigma (\stackrel{}{k})=\sqrt{ϵ_\sigma ^2(\stackrel{}{k})+\mathrm{\Delta }_\sigma ^2(\stackrel{}{k})}`$. The density of the paired carriers is
$$n_s=\frac{1}{2}\left[\underset{\stackrel{}{k}}{}\frac{\mathrm{\Delta }_\sigma ^2(\stackrel{}{k})}{E_\sigma ^2(\stackrel{}{k})}th^2\frac{E_\gamma (\stackrel{}{k})}{2\theta }+\underset{\stackrel{}{k}}{}{}_{}{}^{\tau }\frac{\mathrm{\Delta }_\tau ^2(\stackrel{}{k})}{E_\tau ^2(\stackrel{}{k})}th^2\frac{E_\tau (\stackrel{}{k})}{2\theta }\right].$$
(4)
The free energy corresponding to the Hamiltonian (1) has been calculated in (see also ), where the paired carrier effective mass isotope defect has been investigated. The ”soft” order parameter with the critical behaviour at $`T_c`$ is characterized by the effective mass
$$m_{ab}=\frac{1}{2}\frac{(\eta _\alpha +\eta _\beta +\eta _\gamma )(\delta _\alpha +\delta _\beta +\delta _\gamma )}{(\eta _\alpha +\eta _\beta )\delta _\gamma m_\gamma ^1+\eta _\gamma (\delta _\alpha m_\alpha ^1+\delta _\beta m_\beta ^1)},$$
(5)
where $`m_\sigma =2\pi h^2\rho _\sigma V^1`$, and $`V=a^2`$ for the CuO<sub>2</sub> plaquette. For $`\mu s`$ being not too close to limiting energies $`\mathrm{\Gamma }_{0\sigma }`$ and $`\mathrm{\Gamma }_{c\sigma }`$ of the bands the following formulas for the quantities entering (5) can be used
$$\eta _\sigma =W\rho _\sigma \mathrm{ln}\left[\left(\frac{2\gamma }{\pi }\right)^2\theta _c^2|\mathrm{\Gamma }_{0\sigma }\mu ||\mathrm{\Gamma }_{c\sigma }\mu |\right],$$
(6)
$$\delta _\sigma =\frac{7}{2}\zeta (3)W\rho _\sigma |\mu \mathrm{\Gamma }_{0\sigma }|(\pi \theta _c)^2,$$
(7)
when $`\mu `$ is located in the integration region ($`\zeta (x)`$ is the zeta-function; $`\gamma =\mathrm{exp}(0.577)`$).
If $`\mu `$ lies out of the band $`\delta _\sigma =0`$ and
$$\eta _\sigma =W\rho _\sigma \mathrm{ln}\left|\frac{\mathrm{\Gamma }_{c\sigma }\mu }{\mathrm{\Gamma }_{0\sigma }\mu }\right|.$$
(8)
According to (5) the supercarrier effective mass depends on the position of the spectral components and the chemical potential at a given doping. The mixed nature of the excitations in the multiband model is reflected in the expression (5). With the approximations used $`m_{ab}`$ is temperature independent.
## 4 Calculated doping dependences
The calculations of $`T_c`$, superconductivity gaps, etc. have been made by numerical integration using a plausible parameter set of Ref. . The doping dependences of the cuprate energetic characteristics are illustrated in . Here we represent the zero-temperature superfluid density and the condensation energy in Fig.1 vs. hole doping $`p=0.28c`$ (it has been taken $`p=0.16`$ for the $`T_c(max)`$). The condensation energy is represented by the thermodynamic critical field as
$$H_{c0}=\sqrt{4\pi [(\rho _\alpha +\rho _\beta )\mathrm{\Delta }_\alpha ^2+\rho _\gamma \mathrm{\Delta }_\gamma ^2]}.$$
(9)
The bell-like curves of $`T_c`$, $`\mathrm{\Delta }_\sigma `$ , $`n_s`$ and $`H_{c0}`$ show the similar behaviour in agreement with the results of . Our calculation of the $`\xi _{ab}`$ correlation length vs $`p`$ has given a valley-profile like curve . The second critical field $`H_{c2}(0)`$ calculated from $`\xi _{ab}`$ has also a well expressed maximum . In this manner the strength of the pairing and the phase coherence develop and vanish simultaneously in our model. Analogous conclusion has been done in a number of recent investigations \[1,2,32-35\]. The $`T_c`$ maximum peak on the doping scale is a result of the electron spectrum doping-driven dynamics which brings all the band components into overlap at the Fermi energy. Bands overlap dynamics appears as a novel source of critical doping concentrations. In the normal state at $`c_\beta `$ and $`c_0`$ insulator-to-metal (in the cold and hot subsystems, respectively) are expected to appear. This is in agreement with the observations . At $`c_0`$ also the large pseudogap vanishes, cf. , due to this overlap. The following overdoped regime corresponds to higher carrier concentrations but to smaller superconducting carrier concentrations. The scatterings which cause pairing by the interband mechanism are reduced here.
It must be also mentioned that the bare itinerant – defect gaps are not manifested in the superconducting density curve because the interband nature of the pairing.
The $`n_s(T)`$ dependence is illustrated in Fig.2 and describes also the penetration depth
$$\lambda =\left[\frac{xm_0c^2a^2l}{4\pi e^2n_s(T)}\right]^{1/2}$$
(10)
dependence on temperature. Here $`m_{ab}`$ is expressed through the free electron mass $`m_0`$ as $`m_{ab}=xm_0`$ and $`l`$ is the c-axis lattice constant.
The doping-dependence curve of the paired carrier effective mass is given in Fig.3. To the authors knowledge it has not been obtained earlier. It is seen that $`m_{ab}`$ cannot be assumed to be constant on the whole doping scale. However, in the actual region of remarkable $`T_c`$ values, a rough estimation near $`x3`$, as often supposed, can be considered as acceptable for estimations. It must be stated that the parameter set used can serve only plausible estimations on the quantitative level.
The large values of $`x`$ at very underdoping correspond to the large effective mass of the narrowest defect $`\beta `$-band. Starting from $`c_\beta `$ the contribution of the wide valence band carriers is continuously added. The kink at $`c_0`$ corresponds to the simultaneous action of the $`\alpha \beta `$ carriers. The essential decrease in $`m_{ab}`$ after reaching $`c_1`$ is connected with the vanished contribution of the heaviest $`\beta `$-carriers. The superconducting carrier effective mass reflects the structure of the electron spectrum and reproduces the well known trend of the superconducting collective towards the normal Fermi liquid behaviour with doping in cuprates.
The essential ratio $`n_s(0)/x`$ determining the inverse penetration depth squared is shown vs doping in Fig.4. The superfluid density bell-like dependence dominates over the effective mass changes. As the result, $`\lambda ^2`$ is characterized by a curve with a well expressed maximum on the doping scale. This is in agreement with the experimental findings . The peaked behaviour of $`n_s(0)/m_{ab}`$ is a natural consequence of the interband superconductivity. This answers the question rised in Ref. about the decrease of $`n_s`$ at overdoping. The sublinear $`T_c`$ vs $`n_s`$ (or $`\lambda ^2`$) plot at underdoping, which has been observed in a number of investigations , is reproduced also by the present model – Fig.5.
Having in mind the results obtained for the energetic characteristics of cuprates \[3-6\] and the coherence length , the results of the present work, and also for the effect of photodoping , one can conclude that the model under consideration, despite of its simplicity, is able to describe the cuprate charge-channel associated properties on the whole doping scale.
The spin-channel effects are out of the scope of the model at present. However, there seems to be some correspondence with the models like where the magnetic properties become explained. The ability of the authors model to reproduce qualitatively the behaviour of various cuprate superconducting characteristics on doping is expected to stimulate to fill it in with precised suppositions and quantitative refinements.
## Acknowledgement
This work was supported by Estonian Science Foundation grant No 6540.
## References
* \] C. Bernhard et al., Phys. Rev. Lett. 86 (2001) 1614.
* \] J.L. Tallon et al., Phys. Rev. B 68 (2003) 180501(R).
* \] N. Kristoffel, P.Rubin, Physica C 356 (2001) 171; Solid State Commun. 122 (2002) 265.
* \] N. Kristoffel, P. Rubin, Eur. Phys. J. B 30 (2002) 495.
* \] N. Kristoffel, P. Rubin, Physica C 402 (2004) 257.
* \] N. Kristoffel, P. Rubin, Proc. Estonian Acad. Sci. Phys. Math. 54 (2005) 98; cond.-mat./0408574v1 (2004).
* \] N. Kristoffel, T. Örd, P. Rubin, cond.-mat./0504431v1 (2005).
* \] K.A. Müller, Physica C 341-348 (2000) 11.
* \] D. Mihailovic, K.A. Müller, in Proceedings of the NATO ASI Materials Aspects of High-T<sub>c</sub> Superconductivity, Kluwer, Dordrecht, 1997, p.1.
* \] N. Kristoffel, P. Konsin, T. Örd, Riv. Nuovo Cim. 17 (1994) 1.
* \] H. Suhl, B.T. Matthias, L.R. Walker, Phys. Rev. Lett. 3 (1959) 552.
* \] V.A. Moskalenko, Fiz. Met. Metalloved. 8 (1959) 503.
* \] A. Ino et al., Phys. Rev. B 65 (2002) 094504.
* \] Y. Ando et al., Phys. Rev. Lett. 87 (2001) 017001.
* \] H. Ihara, Physica C 364-365 (2001) 289.
* \] C.C. Homes et al. Phys. Rev. B 67 (2003) 184516.
* \] A.A. Borisov , V. A. Gavrichkov, S. G. Ovchinnikov, Modern Phys. Lett. B 17 (2003) 479.
* \] T. Timusk, B. Statt, Rep. Progr. Phys. 62 (1999) 61.
* \] P. Calvani, phys. stat. sol. b 237 (2003) 194.
* \] Y. G. Zhao et al., Phys. Rev. B 63 (2001) 132507.
* \] H. Romberg et al., Phys. Rev. B 42 (1990) 8768.
* \] P.V. Bogdanov et al., Phys. Rev. B 64 (2001) 180505.
* \] J. Mesot et al., Phys. Rev. B 63 (2001) 224516.
* \] Y. Yoshida et al., Phys. Rev. B 63 (2001) 220501.
* \] A.D. Gromko et al., Phys. Rev. B 68 (2003) 174520.
* \] S. Sugai et al., Phys. Rev. B 68 (2003) 184504.
* \] A. Ino et al., Phys. Rev. Lett. 79 (1997) 2101.
* \] T. Tohyama, S. Maekawa, Phys. Rev. B 67 (2003) 092509.
* \] N. Harima et al., Phys. Rev. B 67 (2003) 172501.
* \] T. Örd, N. Kristoffel, phys. stat. sol. b 216 (1999) 1049.
* \] T. Schibauchi et al., Phys. Rev. Lett. 86 (2001) 5763.
* \] D.L. Feng et al., Science 289 (2000) 277.
* \] T. Schneider, Physica B 326 (2003) 289.
* \] M.R. Trunin, Yu. A. Nefyodov, A.F. Shevchun, Phys. Rev. Lett. 92 (2004) 067006.
* \] R. H. He et al., Phys. Rev. B 69 (2004) 220502(R)
* \] F. Venturini et al., Phys. Rev. Lett. 89 (2002) 107003.
* \] X.F. Sun, K. Segawa, Y. Ando, Phys. Rev. Lett. 93 (2004) 107001.
* \] Y.J. Uemura et al., Phys. Rev. Lett. 66 (1991) 2665.
* \] N. Kristoffel, P. Rubin, Physica C 418 (2005) 49.
* \] A.S. Alexandrov, P.P. Edwards, Physica C 331 (2000) 97.
Figure captions
Fig. 1. The superfluid density (full line) and the thermodynamic critical field (dashed line) on the doping scale.
Fig. 2. The dependence of the superconducting density on temperature for the doping $`p=0.14`$.
Fig. 3. The supercarrier effective mass dependence on doping.
Fig. 4. The inverse penetration depth squared vs doping representated by $`n_s(0)/x`$; $`m_{ab}=xm_0`$.
Fig. 5. The Uemura type plot. |
warning/0506/hep-th0506174.html | ar5iv | text | # 1 Introduction
## 1 Introduction
During the past year a number of objects of have been discovered in $`𝒩=1`$ supersymmetric gauge theories, such as boojums at the interface of vortices and domain walls and an axionic string . Besides the new axionic string a number of other strings are known to exist in this theory, such as the torsion ($`_n`$) charged generalizations of Douglas-Shenker strings in pure SYM, the nontorsion ($``$) charged and BPS strings of Refs. and the non-BPS strings of Refs. which are thought to be unstable and so carry no conserved charges. This wide variety of known strings leads one to wonder whether even more strings await discovery.
In this note we will use M5-brane constructions of various $`𝒩=1`$ and $`𝒩=2`$ supersymmetric gauge theories to find a classification of the conserved charges carried by the unconfined and stable massive matter in these theories in terms of the topology of the configuration. We will consider matter that arises from the dimensional reduction of M2-branes, although some instantons come instead from momentum modes about nontrivial cycles in the spacetime. M2-branes can only end on M5-branes, and so, generalizing Witten’s classification of Douglas-Shenker strings in pure super Yang-Mills , matter charges correspond to worldvolumes of M2-branes with boundaries on the M5-brane. We will argue that the topological charges are valued in the relative homology groups $`\text{H}_k(M,\mathrm{\Sigma })`$ where $`M`$ is the internal spacetime manifold, $`\mathrm{\Sigma }`$ is the embedding of the M5 in the internal space, and $`k=1,2`$ and $`3`$ for strings, particles and certain instantons respectively. If a string carries a conserved charge then there is no matter in the theory on which that string can end, and so the corresponding string charge will necessarily not be screened. There may be matter on which two strings carrying conserved charges can end, but such matter does not screen the charge of the individual strings, instead it screens the difference in charges of the two strings.
The charge groups will first be calculated for pure super Yang-Mills with a polynomial superpotential that breaks the gauge symmetry from $`U(N)`$ to a product of $`U(N_i)`$’s. Here we will see that the only topologically stable (not screened) strings are Douglas-Shenker strings which are charged under a single torsion group $`_K`$ where $`K`$ is the confinement index of Ref. . In addition there will be various ’t Hooft-Polyakov monopoles and dyons, as well as W bosons which will usually be confined. The fact that these particles are confined by torsion charged strings means that there will be finite combinations of particles that will be unconfined, leading to a rich spectrum of stable electric, magnetic and mixed glueballs with masses proportional to the various Higgs VEVs.
When fundamental matter is added the Douglas-Shenker strings become unstable, they may decay via the nucleation of quark-antiquark pairs and correspondingly the relative homology group that classifies them vanishes. The strings may still be quite long-lived if a large bare mass is given to the quarks, and they are still physically relevant as they confine the quarks into baryons. One may try to save some Douglas-Shenker strings by making some of the bare quark masses degenerate, in which case the relative homology group becomes nontrivial, or more precisely becomes dependent upon the way the topology at infinity is treated. However under an arbitrarily small perturbation this degeneracy is destroyed, and so we claim that the strings may decay. That is, the physics determines the kind of homology used. On the other hand one may make the degeneracy stable against small fluctuations by gauging the flavor symmetry, and we will see that in some such theories there are vacua with stable Douglas-Shenker strings. However in cascading theories and more generally in baryonic vacua it appears as though the corresponding homology group is trivial and so no strings carry the topological charge classified by this group.
In the presence of fundamental matter there are supersymmetric configurations in which the M5-brane may have two connected components, or even more for some theories that in the UV already have several gauge groups. We will see that the homology group classifying strings contains a number of $``$ factors equal to the number of components minus one. For example, in SQCD with a FI term one finds, as has been demonstrated already in Refs. in three and four dimensions, stable BPS vortices. If one introduces a superpotential then these vortices are generically no longer BPS, although supersymmetries may appear in the worldsheet theory that were not present in the bulk theory . However for any superpotential polynomial in the adjoint chiral multiplets the M5-brane will remain disconnected if the FI term is nonzero, and so even when they are non-BPS we claim that these vortices will be absolutely stable.
On the other hand vortices of the type studied in Refs. correspond to connected M5-branes and so may decay via monopole-antimonopole creation. An FI term is not required for the stability of these vortices. Instead we will see that vortices created from superpotentials alone will be stable whenever every color is locked to a flavor by a nonvanishing meson VEV, or equivalently when all monopoles are confined by either 0 or 2 vortices. We will also argue that the axionic vortices of Ref. come in two varieties and both are always unstable.
In this analysis we treat the spacetime and the M5-brane embedding classically. This supergravity limit may be different from the limit in which the dimensionally-reduced theory of interest is obtained. Thus calculations of tensions and lifetime should be independently verified in the theories in question when they are not protected by nonrenormalization theorems. However the topological charges computed in this note, at least in the examples that we know, survive the change of limits despite the fact that many of the stable objects are not BPS.
In section 2 we argue that in general the relative homology of the M5-brane embedding is a group of conserved charges and we review an exact sequence which will be used repeatedly throughout the paper to calculate the relative homologies in examples. In section 3 we describe a simple example, M-theory compactified on a 2-torus on which three M5-branes are wrapped. This gives a dimensional reduction of 5-dimensional $`U(3)`$ pure super Yang-Mills, which consists of $`𝒩=4`$ 4-dimensional $`U(3)`$ super Yang-Mills coupled to a $`U(3)`$ lower-form gauge theory with a zero-form connection and a one-form field strength. The various M2-branes connecting the M5-branes correspond to the stable objects known to exist in the gauge theory plus those of the 1-form theory. Next in section 4 we extend this analysis to $`𝒩=2`$ pure super Yang-Mills, or more precisely flavorless MQCD, by compactifying M-theory on a single circle and considering an M5-brane embedding given by the logarithm of the corresponding Seiberg-Witten curve . A superpotential is added, softly breaking the supersymmetry to $`𝒩=1`$, in section 5 and a number of examples are considered. Finally in section 6 flavored matter is added, and both global and local flavor symmetries are considered. Gubser-Herzog-Klebanov strings are described in the IIA theory, but the calculation of conserved vortex charges shows that they do not carry any conserved charge of the kind classified here.
## 2 Charges from Relative Homology
Consider M-theory on an 11-dimensional spacetime of the form $`^{1,d}\times M`$ where the gauge theory of interest lives on the spacetime $`^{1,d}`$ and $`M`$ will be referred to as the internal space. All of the results in this paper extend equally well to the case in which $`^{1,d}`$ is replaced by an arbitrary ($`d`$+1)-dimensional manifold. We will be interested in the low energy theory of an M5-brane that fills the gauge theory $`^{1,d}`$ and sweeps out a $`(5d)`$-manifold $`\mathrm{\Sigma }M`$. In general $`\mathrm{\Sigma }`$ and even $`M`$ may depend upon the position in the physical spacetime $`^{1,d}`$, for example SQCD and SYM contain domain walls that separate regions with topologically distinct embeddings. In this note we will consider a region of spacetime in which the topology of the embedding of the M5 in the internal space is constant.
We will now describe a classification of topological charges carried by M2-branes whose boundaries lie entirely on the M5-brane. M2-brane boundaries can only lie on M5-branes because the supergravity 7-form current $`G_4+C_3G_4`$ is gauge-invariant and in particular globally defined. This means that it is annihilated by the square of the exterior derivative. Identifying $`dG_4`$ and $`dG_4`$ with the M5-brane and M2-brane charge densities $`\rho _5`$ and $`\rho _2`$ respectively the gauge-invariance of the 7-form implies
$$0=d^2(G_4+C_3G_4)=d\rho _2+G_4\rho _5+C_3d\rho _5$$
(2.1)
where the gauge invariance of $`G_4`$ implies that $`d\rho _5`$ vanishes. The boundary of an M2-brane is characterized by the nonvanishing of $`d\rho _2`$, but Eq. (2.1) implies that when $`d\rho _2`$ is nonzero $`\rho _5`$ is also nonzero and so there must be an M5-brane at the M2’s boundary. Of course this argument does not apply at the end of the world where derivatives are not defined, but we will restrict our attention to spacetimes with no boundaries.
If the spacetime is geometrically a product of $`^{1,d}\times M`$, as it will be for the gauge theories of interest far from a domain wall, we may deform the worldvolume of each M2-brane into a number of components each of which is a product of a submanifold of the internal space and a submanifold of the gauge theory spacetime. For example a diagonal line can be deformed into a straight horizontal and a straight vertical line without changing its topology. We will see that such composites of products correspond to composite objects in the gauge theory, for example, monopoles with vortices attached. We will then classify each product manifold separately.
Strings of tension $`T`$ correspond to M2-branes that extend in two gauge theory spacetime directions while the third is a line segment in $`M`$ of length $`T`$ bounded by the M5-brane. Similarly particles of mass $`M`$ correspond to M2-branes that extend in one gauge theory spacetime direction while the other two sweep out an area $`M`$ surface bounded by a collection of curves in the M5. M2-branes that correspond to instantons are extended in no gauge theory spacetime directions and an internal 3-manifold bounded by a collection of surfaces in the M5.
However there are also composite configurations, in which, for example, the 2-dimensional sheet corresponding to a particle has one boundary not on the M5. Such a boundary is a 1-dimensional line corresponding to a vortex which, as the M2 is everywhere 3-dimensional, continues into a spacetime direction. That is, while the M2 can only end on the M5 it may, now that we have distorted it into product submanifolds, have corners in which it turns into a spacetime direction. In the dimensionally reduced theory this configuration appears to be a particle with a vortex ending on it, in other words this particle is confined and the vortex charge is screened. In general a particle is confined by a number of vortices equal to the number of components of the boundary of the corresponding sheet that are not on the M5-brane, which in SQCD may be 0, 1 or 2 . We are classifying unconfined objects, and so we are searching for a group of charges that corresponds to surfaces with boundaries only on the M5.
If a given set of vortices confines a particle then this set of vortices is unstable and can decay via nucleation of that particle and its antiparticle. The lifetime is exponential in the mass of the particle, so such vortices may last for quite awhile, but because their lifetime is finite they will carry no conserved charge. Similarly instantons may be “confined” by particles if the corresponding 3-manifolds have boundaries off of the M5, and the particles may decay via the confined instanton. Thus the group of conserved charges in dimension $`k`$ consists of the k-manifolds with boundaries lying entirely along the M5 (the unconfined) quotiented by those that are themselves boundaries (the unstable). This is the definition of the relative homology group $`\text{H}_k(M,\mathrm{\Sigma })`$.
The group $`\text{H}_3(M,\mathrm{\Sigma })`$, which classifies instantons, is trivial for 4-dimensional SQCD, where $`\mathrm{\Sigma }`$ is the logarithm of the corresponding Seiberg-Witten curve and $`M=^6\times S^1`$. However if this theory is dimensionally reduced on a spacetime circle then the particles of the 4-dimensional theory may wrap this circle and their dimensional reductions will be instantons that carry a charge classified by $`\text{H}_3(M,\mathrm{\Sigma })`$. Such wrapping configurations need not be instantonic, but may also correspond to a particle-antiparticle pair that is created, circumnavigates the circle, and then annihilates leaving only a quantized flux residue behind. If the particle was confined by a string in the 4-dimensional theory, then the dimensional reduction of the string to 3-dimensions will be an unstable particle that decays by the dimensional reduction of the above process.
The relative homology groups will be calculated using the long exact sequence for relative homology
$$\mathrm{}\stackrel{j_{}^{k+1}}{}\text{H}_{k+1}(M,\mathrm{\Sigma })\stackrel{_{}^{k+1}}{}\text{H}_k(\mathrm{\Sigma })\stackrel{i_{}^k}{}\text{H}_k(M)\stackrel{j_{}^k}{}\text{H}_k(M,\mathrm{\Sigma })\stackrel{_{}^k}{}\mathrm{}$$
(2.2)
where $`i_{}`$ is the map induced by the inclusion $`i:\mathrm{\Sigma }M`$, $`j`$ is the quotient on chains by the image of the M5-brane, and $`_{}`$ is the map induced by the boundary map.
## 3 $`𝒩=4`$ U(3) Gauge Theory With Strings
A 4-dimensional $`𝒩=4`$ $`U(3)`$ gauge theory coupled to a lower-form gauge theory with a zero-form $`U(3)`$ connection and one-form field strength may be engineered from M-theory on $`^{1,8}\times T^2`$ with 3 M5-branes wrapping the $`T^2^5\times T^2`$ and extended along $`^{1,3}`$. Matter corresponds to M2-branes extending between any two of the M5’s and wrapping some subtorus of the $`T^2`$. The ordinary gauge theory dyons come from M5’s wrapping a circle $`S^1T^2`$ and extending along a strip whose boundary consists of a line, the particle trajectory, along two M5’s. There are also M2-branes that wrap the entire torus and a line segment between two branes. These are the electrically charged instantons of the one-form theory. Finally there are M2-branes that do not wrap the torus at all but fill a 3-dimensional strip which is a line segment between two M5’s crossed with a 2-dimensional surface along the M5’s. This 2-dimensional surface is the worldsheet of the ’t Hooft-Polyakov string which is magnetically charged under the abelian gauge 1-form field strength. Note that the particles are not mutually BPS with the strings and instantons, although the strings and instantons may be mutually BPS with each other.
The masses of the particles, tensions of the strings and actions of the instantons are proportional to the corresponding Higgs VEVs, which are the distances between pairs of M5’s multiplied by the volumes of the wrapped subtorii. The instanton, particle and string each transform in the adjoint of U(3) and correspondingly each has 6 massive components corresponding to the roots of U(3), two of which are independent corresponding to the two simple roots. Thus the charge group is $`^2`$ for the instanton and string, and $`^4`$ for the particle since there will be a W boson and also a monopole for each simple root.
If two of the M5’s are coincident then some of the objects become massless, but under a slight deformation they will be massive. The classification scheme proposed in this note only applies to configurations that are stable under slight deformations, otherwise, for example, we would have predicted the existence of stable Douglas-Shenker strings in the case in which two M5’s are degenerate. Homology classes that are unstable under small perturbations of the configuration do not correspond to stable objects, instead these objects delocalize out of existence as one approaches transitions in the moduli space. Thus we will restrict our attention to the case in which no two M5’s are coincident.
The relative homology groups are easily computed from the exact sequence (2.2). The homology of the internal space $`M`$ is just that of the torus $`T^2`$, since the remaining $`^5`$ is contractible.
$$\text{H}_0(M)=\text{H}_2(M)=,\text{H}_1(M)=^2.$$
(3.1)
The M5 wraps the torus, and so the homology of each component of the M5 will also be that of the torus, and that of the whole M5 will then be that of the torus cubed
$$\text{H}_0(\mathrm{\Sigma })=\text{H}_2(\mathrm{\Sigma })=^3,\text{H}_1(\mathrm{\Sigma })=^6.$$
(3.2)
Each M5 component wraps the torus once, and so the $`i_{}`$ maps will be a copy of the identity on each component. In particular, on $`\text{H}_0`$ and $`\text{H}_2`$ it will be a 3 by 1 matrix with all entries equal to 1 while on $`\text{H}_1`$ it will be a $`6`$ by $`2`$ matrix made of three copies of the $`2`$ by $`2`$ identity identity matrix. $`j_{}`$ is the zero map as it quotients chains on the torus by themselves.
This leaves three short exact sequences
$$0\stackrel{j_{}^{k+1}}{}\text{H}_{k+1}(M,\mathrm{\Sigma })\stackrel{_{}^{k+1}}{}\text{H}_k(T^2)^3\stackrel{i_{}^k}{}\text{H}_k(T^2)\stackrel{j_{}^k}{}0$$
(3.3)
which yields
$$\text{H}_{k+1}(M,\mathrm{\Sigma })=\text{H}_k(T^2)^2.$$
(3.4)
In particular
$$\text{H}_1(M,\mathrm{\Sigma })=\text{H}_3(M,\mathrm{\Sigma })=^2,\text{H}_2(M,\mathrm{\Sigma })=^4$$
(3.5)
matching the expectations from the field theory.
We have only classified topologically stable objects charged under the adjoint of the gauge group, but the bulk fields lead to additional objects that are $`U(3)`$ singlets. For example the first compactified circle leads to a $`U(1)`$ gauge symmetry on the remaining 5-dimensional theory, complete with magnetic and electric charges which are realized by D6 and D0-branes respectively. Curiously these D0-branes, the $`U(1)`$ electric charges, are also the instantons of the $`U(3)`$ gauge theory. Reducing on the second circle leads to a second $`U(1)`$ gauge symmetry with a new set of electric and magnetic charges, while the old $`U(1)`$, similarly to the $`U(3)`$, is decomposed into an ordinary two-form Maxwell theory plus an abelian one-form theory. The electric charges of the one-form theory are now the instantons of the $`U(3)`$ gauge theory. Each D6-brane that does not wrap the second circle, at least in BPS configurations, leads to a flavor of quark matter in the theory. The other branes do not lead to objects that exist in the pure gauge theory, but rather to topological charges characteristic of one-form and mixed one-form-two-form theories. While these objects all carry topological charges, they do not correspond to M2-branes ending on the M5-brane and so are missed by the charge classification scheme of this note.
## 4 $`𝒩=2`$ Super Yang-Mills
We next turn our attention to $`𝒩=2`$ super Yang-Mills with gauge group $`U(N)`$, whose bare Lagrangian in $`𝒩=1`$ superspace takes the form
$$=d^2\theta \frac{1}{2e^2}\text{Tr}_{N_c}(W^\alpha W_\alpha )+h.c.+d^2\theta d^2\overline{\theta }\frac{2}{e^2}\text{Tr}_{N_c}(\mathrm{\Phi }{}_{}{}^{}e_{}^{V}\mathrm{\Phi }e^V).$$
(4.1)
Here $`\mathrm{\Phi }`$ and $`W`$ are chiral and vector superfields respectively, both transforming in the adjoint of the gauge group.
Following the construction in Ref. this theory is a sector of the theory obtained by embedding an M5-brane in the spacetime $`^{1,9}\times S^1`$ such that the M5 covers the physical spacetime $`^{1,3}`$ and the embedding in the internal $`^6\times S^1`$ is given by the logarithm of the Seiberg-Witten curve
$$y^2=P_N^2(v)\mathrm{\Lambda }^{2N}$$
(4.2)
where $`P_N`$ is a polynomial of order $`N`$
$$P_N(v)=\frac{1}{2}\underset{i=1}{\overset{N}{}}(v\varphi _k)$$
(4.3)
constructed from the eigenvalues $`\varphi _k`$ of the adjoint scalar.
The internal space $`M`$ is $`^6\times S^1`$, which we parametrize with three complex coordinates $`v,w`$ and $`s`$ and one real coordinate $`x^7`$. $`v`$, $`w`$ and $`x^7`$ parametrize the $`^5^6\times S^1`$ while the complex coordinate $`s=x^6+ix^{10}`$ parametrizes the remaining $`\times S^1`$. $`x^10`$ is a periodic coordinate parametrizing the M-theory circle. For future use we define the complex coordinate $`t=exp(s)`$ which is valued in $`^{}`$. These coordinates are related to the Seiberg-Witten curve (4.2) via the change of coordinates $`t=y+P_N(v)`$. In all the M5-brane embedding is the Riemann surface in $`M`$ given by the equations
$$t^22P_N(v)t+\mathrm{\Lambda }^{2N}=w=x^7=0.$$
(4.4)
Dimensionally reducing to type IIA this configuration becomes two parallel NS5-branes with $`N`$ D4-branes stretched between them. The effective worldvolume theory of the D4’s contains the gauge theory of interest.
The $`𝒩=2`$ theory enjoys a $`(2N2)`$-dimensional moduli space of vacua in which the dyon masses are smooth functions of the coordinates. In particular configurations with massless dyons are codimension two in the moduli space and so, as the moduli space is connected, will not satisfy our stability criteria. Thus it will suffice to consider vacua in which all dyons are massive. In such vacua the M5-brane consists of $`N`$ separated tubes which each wrap the M-theory circle once. All of the tubes connect at large and small $`x^6`$ to form a genus $`N1`$ Riemann surface $`\mathrm{\Sigma }_{N1}`$ with two punctures at $`x^6=\pm \mathrm{}`$. The integral homology groups of this Riemann surface are
$$\text{H}_0(\mathrm{\Sigma }_{N1})=,\text{H}_1(\mathrm{\Sigma }_{N1})=^{2N1},\text{H}_2(\mathrm{\Sigma }_{N1})=0$$
(4.5)
where the $`2N1`$ generators of $`\text{H}_1`$ are the $`N`$ $`A`$-cycles that wrap the tubes and the $`N1`$ dual $`B`$-cycles that run down the $`i`$th tube and back up the $`i+1`$st.
The relative homology groups may be calculated using the long exact sequence
$`0`$ $`=`$ $`\text{H}_2(^6\times S^1)\stackrel{j_{}^2}{}\text{H}_2(^6\times S^1,\mathrm{\Sigma }_{N1})\stackrel{_{}^2}{}\text{H}_1(\mathrm{\Sigma }_{N1})=^{2N1}`$
$`\stackrel{i_{}^1}{}`$ $`\text{H}_1(^6\times S^1)=\stackrel{j_{}^1}{}\text{H}_1(^6\times S^1,\mathrm{\Sigma }_{N1})\stackrel{_{}^1}{}\text{H}_0(\mathrm{\Sigma }_{N1})=`$
$`\stackrel{i_{}^0}{}`$ $`\text{H}_0(^6\times S^1)=\stackrel{j_{}^0}{}\text{H}_0(^6\times S^1,\mathrm{\Sigma }_{N1})=0.`$
The inclusion $`i_{}^0:`$ is degree one, as the embedding maps a point to a single point, and so the kernel of $`i_{}^0`$ is trivial. The kernel of $`i_{}^0`$ is the image of $`_{}^1`$ and so $`_{}^1`$ is the zero map and $`j_{}^1`$ is onto
$$i_{}^0=1,_{}^1=0,\text{H}_1(^6\times S^1,\mathrm{\Sigma }_{N1})=\text{Image}(j_{}^1).$$
(4.6)
Of the $`2N1`$ cycles that generate $`\text{H}_1(\mathrm{\Sigma }_{N1})`$, the $`N`$ $`A`$-cycles each wrap the M-theory circle, which generates $`\text{H}_1(^6\times S^1)`$, while a basis of $`N1`$ $`B`$-cycles is chosen such that no $`B`$-cycle wraps the M-theory circle. Therefore $`i_{}^1`$ is a $`(2N1)`$-dimensional column vector consisting of $`N`$ $`1`$’s and $`N1`$ $`0`$’s. In particular, every element of $`\text{H}_1(^6\times S^1)`$ is in the image of $`i_{}^1`$, reflecting the fact that there is some cycle in the Riemann surface that wraps the M-theory circle any given number of times. Thus $`i_{}^1`$ is onto, so $`j_{}^1`$ must be the zero map. However we have seen the $`j_{}^1`$ is onto, thus the group of vortex charges vanishes
$$i_{}^1(a_i,b_i)=\underset{i=1}{\overset{N}{}}a_i,j_{}^1=0,\text{H}_1(^6\times S^1,\mathrm{\Sigma }_{N1})=\text{Image}(j_{}^1)=0$$
(4.7)
and there are no topologically stable vortices in the pure $`𝒩=2`$ theory.
The map $`j_{}^2`$ has a trivial domain, and so it has a trivial image. Thus the map $`_{}^2:\text{H}_2(^2\times S^1,\mathrm{\Sigma }_{N1})^{2N1}`$ is into. The image of $`_{}^2`$ is the kernel of $`i_{}^1`$, which consists of all $`N1`$ $`B`$-cycles plus the $`(N1)`$-dimensional subspace of $`A`$-cycles that wrap the M-theory circle zero times. The fact that $`_{}^2`$ is into implies that this image is also its domain, the group of particle charges $`\text{H}_2(^2\times S^1,\mathrm{\Sigma }_{N1})`$. Thus there are $`2N2`$ independent M2-brane worldvolumes yielding particles, of which $`N1`$ are bounded by the $`B`$-cycles while $`N1`$ are bounded by pairs of $`A`$-cycles with opposite orientations. These are just a basis of the ’t Hooft-Polyakov monopoles and the W bosons. The entire group of conserved charges,
$`\text{H}_2(^2\times S^1,\mathrm{\Sigma }_{N1})=^{2N2}`$ (4.8)
consists of all combinations of W bosons and ’t Hooft-Polyakov monopoles. Every such combination has a BPS bound, given by the areas of complex surfaces (when they exist) bounded by the corresponding cycles. Physically most of these BPS bounds are never saturated, and the various monopoles and W bosons often repel. However the topological classification is insensitive to this repulsion, it includes non-BPS topologically stable configurations.
Note that there is an analogous configuration to the vortex of the $`𝒩=4`$ 1-form gauge theory, an M2 which extends between two tubes but whose boundary consists of two points on the Riemann surface, rather than a cycle. However unlike the $`𝒩=4`$ case this configuration does not yield a nontrivial homology class, as it may be pushed in the $`\pm x^6`$ direction until it is lies on the M5-brane where it may dissolve. We will argue that this unstable vortex, rotated by 90 degrees, is an $`𝒩=2`$ version of the unstable axionic vortex of Ref. . There is no finite-volume configuration analogous to the $`𝒩=4`$ 1-form instanton. These two configurations, which correspond to the physics of the 1-form theory and not the 2-form gauge theory of interest, reappear in some vacua when we compactify the $`x^6`$ direction to introduce bifundamental matter. The gauge theory instantons arise from momentum modes about the M-theory circle, and not M2-branes, and so are missed by this classification.
If one adds $`N_f2N2`$ semi-infinite tubes that wrap the M-theory circle, corresponding to the Seiberg-Witten curve
$$y^2=\frac{1}{4}\underset{k=1}{\overset{N}{}}(v\varphi _k)^2+\mathrm{\Lambda }^{2n_cn_f}\underset{j=1}{\overset{N_f}{}}(v+m_j)$$
(4.9)
the gauge theory will include $`N_f`$ flavors of fundamental matter hypermultiplets with bare masses $`m_j`$. Correspondingly $`\text{H}_1(\mathrm{\Sigma })`$ will gain $`N_f`$ new generators. $`i_{}^1`$ is already onto, and so these extra generators imply an extra $`^{N_f}`$ factor in the group of stable particles $`\text{H}_2(M,\mathrm{\Sigma })`$ corresponding to the quarks. Linear combinations of all of the particles are now generically flavored dyons.
## 5 $`𝒩=1`$ Super Yang-Mills
### 5.1 The $`𝒩=1`$ Curve
We will now softly break the $`𝒩=2`$ supersymmetry to $`𝒩=1`$ by introducing a superpotential $`𝒲(\mathrm{\Phi })`$ that is degree $`m+1`$ polynomial in the adjoint chiral multiplet $`\mathrm{\Phi }`$
$``$ $`=`$ $`{\displaystyle d^2\theta \frac{1}{2e^2}\text{Tr}_{N_c}(W^\alpha W_\alpha )}+h.c.`$
$`+{\displaystyle \frac{2}{e^2}}\text{Tr}_{N_c}(\mathrm{\Phi }{}_{}{}^{}e_{}^{V}\mathrm{\Phi }e^V)+{\displaystyle d^2\theta \sqrt{2}\text{Tr}_{N_c}𝒲(\mathrm{\Phi })}+h.c..`$
This corresponds to replacing replacing the embedding condition $`w=0`$ of Eq. (4.4) with
$$w^22𝒲{}_{}{}^{}(v)w\stackrel{~}{f}_{m1}(v)=0$$
(5.2)
for some degree $`m1`$ polynomial $`\stackrel{~}{f}_{m1}`$ that captures quantum corrections to the superpotential.
Eq. (5.2) is quadratic and so we have not only deformed but also doubled our original Riemann surface, that is each point $`v`$ of the $`𝒩=2`$ curve has split into two different points representing the two roots of $`w(v)`$ in (5.2). However in Ref. the authors have argued that $`𝒩=1`$ supersymmetry requires that all odd-degree roots of $`w^2`$ are also odd-degree roots of $`y^2`$, and vice versa. Therefore any loop will encircle an even number of roots, and so change sheets an even number of times. Thus no path connects the two sheets, and we may throw one sheet away, leaving the Riemann surface undoubled. If we do not throw this sheet away we will not obtain a deformation of the $`𝒩=2`$ theory, but rather two coupled copies of such a deformation.
Reducing to IIA one obtains two NS5-branes, again at different $`x^6`$ positions. The NS5 on the left is at $`w=0`$, while that on the right, which is often named NS5$`^{}`$, is at $`w=𝒲{}_{}{}^{}(v)`$. The roots of $`w(v)`$ are thus the critical points of the superpotential, while those of $`P_N(v)`$ still correspond to the expectation values of the adjoint scalars. $`𝒩=1`$ supersymmetry implies that the roots of $`P_N(v)`$ are all within a distance of order $`\mathrm{\Lambda }`$ of the roots of $`w(v)`$. In field theory this corresponds to the fact that the eigenvalues of the adjoint scalars are at extrema of the superpotential up to quantum corrections of order $`\mathrm{\Lambda }`$.
From the point of view of the brane cartoon this condition reflects the fact that supersymmetry requires the D4’s to proceed along the $`x^6`$ direction without bending in the $`w`$-plane, and so they may only connect the NS5-branes when they both occupy the same $`w`$ coordinate. The NS5 on the left is always at $`w=0`$, so D4’s may only be placed at points $`v`$ such that the NS5$`^{}`$ is also at $`0=w=𝒲{}_{}{}^{}(v)`$. The distribution of eigenvalues among the critical points of the superpotential leads to a classical breaking of the gauge symmetry
$$U(N)\underset{i=1}{\overset{k}{}}U(N_i)$$
(5.3)
where $`i`$ runs over the $`k`$ critical points, which are taken to be separated by a distance much greater than $`\mathrm{\Lambda }`$ to avoid exotic vacua such as those of Ref. . Only a finite set of points $`\{\varphi _k\}`$ on the Coulomb branch satisfy the $`𝒩=1`$ supersymmetry requirement, corresponding to the existence of $`N_i1`$ independent massless dyons in the $`U(N_i)`$ subsector. In particular this means that the above stability condition no longer excludes vacua with massless states, as it did in $`𝒩=2`$. In field theory terms this stability is caused by a nontrivial gluino condensate that obstructs the deformations away from these points. It is the topology of the condensate field that stabilizes the vortex.
The $`N_i`$ D4-branes at each critical point of the superpotential lift to a single tube of M5-brane which wraps the M-theory circle $`N_i`$ times. This tube cannot split into smaller tubes, as the distance between these tubes would give a mass to all of the dyons of their respective gauge groups which would be incompatible with the known gluino condensate. Thus the $`𝒩=1`$ super Yang-Mills M5-brane is a genus $`k1`$ Riemann surface $`\mathrm{\Sigma }_{k1}`$. However the soliton spectrum is not identical to the $`N=k1`$ case of the $`𝒩=2`$ curve because the embedding of $`\mathrm{\Sigma }_{k1}`$ into the internal space $`M`$ is topologically inequivalent as a result of the fact that the $`A`$-cycles now wrap the M-theory circle $`N_i`$ times instead of just once. In addition the embedding of the $`B`$-cycles in the spacetime now will sometimes wrap the M-theory circle, in which case the cycle will not bound a disk in $`M`$ and the corresponding monopole will be confined. This is in contrast with the $`𝒩=2`$ case, where it was always possible to concatenate a $`B`$-cycle that encircles the M-theory circle $`k`$ times with $`k`$ $`A`$-cycles to construct a new $`B`$-cycle with winding number zero. To understand these phenomena we will consider some special cases. The case $`N_i=1`$ is identical to that of $`𝒩=2`$.
### 5.2 $`U(N)`$ Super Yang-Mills
This case has been analyzed by Witten in Ref. . If, up to corrections of order $`\mathrm{\Lambda }`$, all of the VEVs of the adjoint scalar are at the same critical point of the superpotential, for example if the superpotential is quadratic and so only has one critical point, then the gauge symmetry is classically unbroken, although quantum mechanically it will be dynamically broken.
The Riemann surface is genus zero, with no $`B`$-cycles and a single $`A`$-cycle which wraps the M-theory circle $`N`$ times. The Witten index of this theory is $`N`$ and correspondingly there are $`N`$ distinct Riemann surfaces. For example in the limit in which the adjoint chiral multiplet mass, which is the quadratic term in the superpotential, goes to infinity one finds
$$v^n=t,w=\zeta v^1$$
(5.4)
where $`\zeta `$ is an $`n`$th root of unity that labels the vacuum. The choice of root of unity corresponds to a choice of how to identify the $`N`$ $`x^{10}=0`$ paths on each side of the tube, or more precisely to cyclic permutations of the identifications of the points $`x^{10}=0`$ at $`x^6=+\mathrm{}`$ and $`x^6=\mathrm{}`$. While the choice of identifications will affect the set of topological charges in the later examples, it will not have any effect in the $`k=1`$ case considered in this subsection.
This Riemann surface is homeomorphic to that of the $`U(1)`$ case of the $`𝒩=2`$ theory, however the embedding is inequivalent as the A-cycle now wraps the M-theory circle $`N`$ times. Thus the exact sequence is identical to Eq. (4) but now $`i_{}^1`$ is multiplication by the number $`N`$
$$i_{}^1:\text{H}_1(\mathrm{\Sigma }_0)=\text{H}_1(^6\times S^1)=:jNj.$$
(5.5)
In particular $`i_{}^1`$ is no longer onto, the image consists only of integers that are multiples of $`N`$.
Physically this means that an M2-brane that winds around the M-theory circle $`N`$ times can unwind by touching the M5, opening and letting the two ends slide around the $`A`$-cycle one time relative to each other. However a loop that winds less than $`N`$ times, while it may still open and slide around the $`A`$-cycle, only changes its winding number by $`N`$ and so can never get its winding number to zero and disappear. Thus these winding M2-branes carry stable charges classified by $`_N`$.
To see this from the exact sequence, recall that $`_{}^1`$ is the zero map. This implies that the group of conserved vortex charges is
$$\text{H}_1(^6\times S^1,\mathrm{\Sigma }_0)=\text{Image}(j_{}^1)=\frac{\text{H}_1(^6\times S^1)}{\text{Ker}(j_{}^1)}=\frac{\text{H}_1(^6\times S^1)}{\text{Image}(i_{}^1)}=\frac{}{N}=_N.$$
(5.6)
These are the charges of the Douglas-Shenker strings. The relative homology groups correspond to lines at constant $`x^{10}`$ connecting points on the $`A`$-cycle that are at equal values of the M-theory circle coordinate $`x^{10}`$.
The rest of the sequence is unchanged and so, as in the $`𝒩=2`$ theory, there are no topologically stable, unconfined particles. For example one may follow a cylinder that was a W boson in the $`𝒩=2`$ $`U(N)`$ case as the $`N`$ $`A`$-cycles merge and one finds that the two $`A`$-cycles wrapped by the ends of this cylinder no longer close . Instead the circles have been broken, and the two ends are still at the same $`x^{10}`$ coordinate but now are separated by $`1/N`$ of the $`A`$-cycle. While $`G_4+C_3G_4`$ gauge-invariance does not allow an M2-brane to break without all boundaries laying on the M5, there is a composite current-conserving configuration that contains this broken cylinder. One may attach each broken end of the cylinder to a strip that continues into one of the gauge theory spacetime directions. That is, one may attach the W boson to two Douglas-Shenker strings. Thus W bosons have disappeared from the spectrum because they are confined. The two strings have opposite orientations, and so the total string charge emitted from the confined W boson is zero. Thus in this example Douglas-Shenker strings will not be able to decay via the nucleation of pairs of W bosons, and so their corresponding charges are not screened. In fact if a W-boson is inserted in a string worldsheet the $`_N`$ vortex charge is the same on both sides of the W. Note that the relative charges of the two strings that end on a given W-boson depends on the direction from which the strings approach, as this determines the relative orientations of the M2-branes in the gauge theory directions.
### 5.3 $`U(N+1)U(N)\times U(1)`$ Super Yang-Mills
This configuration is identical to that of the previous subsection except that now there is an extra tube of M5 that wraps the M-theory circle once. The $`v`$ position (adjoint scalar VEV) of this new tube is a different critical point of the superpotential than the point used by the original tube. Now there will be a second $`A`$-cycle, which goes around the new tube once, on which $`i_{}^1`$ acts by multiplication by one.
There are now massive W bosons corresponding to cylindrical M2-branes whose two boundary circles, one of which is on each $`A`$-cycle, each wrap the M-theory circle once. As in the previous example, the boundary on the $`U(N)`$ $`A`$-cycle fails to close, and to make it close a vortex must be inserted. Thus each W boson is confined by a single $`U(N)`$ Douglas-Shenker string and in turn each Douglas-Shenker strings may decay via the nucleation of a W-anti-W pair. The lifetime of such strings is then exponential in the W mass, which is proportional to the distance between the two critical points of the superpotential. Notice that $`N`$ W bosons may come together to form an unconfined glueball, as the $`N`$ vortices confining them may annihilate each other. This glueball corresponds to a cylindrical M2 whose boundary circles wrap the $`U(N)`$ $`A`$-cycle once and the $`U(1)`$ $`A`$-cycle $`N`$ times. As both winding numbers are integral, both boundaries close.
We may construct the $`B`$-cycle by connecting the two paths that connect the tubes to paths that run down the two tubes at constant values of $`x^{10}`$. A choice of $`B`$-cycle exists such that $`x^{10}`$ is constant along the whole cycle and so the $`B`$-cycle is a boundary in the internal space. We will refer to the bounded disk as the monopole or $`(1,0)`$-dyon.
Summarizing, the set of conserved particle charges is identical to the case of a $`U(2)`$ $`𝒩=2`$ theory, although the minimum charge of an unconfined W boson is higher in the $`𝒩=1`$ case. The set of particle charges is generated by the charge 1 monopole and a charge $`N`$ bound state of W bosons. No decay modes for particles have been found, and so the particle charge group is expected to be $`^2`$. Similarly the vortices can decay via W boson pair production, and so the vortex charge group is expected to be the trivial group $`0`$.
This may be compared against the relative homology groups obtained using the long exact sequence. Everything is as before, except now the inclusion map is
$$i_{}^1:\text{H}_1(\mathrm{\Sigma }_1)=^3\text{H}_1(^6\times S^1)=:(a_1,a_2,b)Na_1+a_2.$$
(5.7)
That is, the inclusion map is matrix multiplication by the column vector $`(N,1,0)`$. Here $`a_1,a_2`$ and $`b`$ are the $`U(N)`$ $`A`$-cycle, the $`U(1)`$ $`A`$-cycle and the $`B`$-cycle respectively. As in the $`𝒩=2`$ case $`i_{}^1`$ is onto and so $`j_{}^1`$ is the zero map, which again implies that there are no nontrivial topological charges for vortices
$$\text{H}_1(^6\times S^1,\mathrm{\Sigma }_1)=0.$$
(5.8)
The kernel of $`i_{}^1`$ consists of all triplets $`(a_1,a_2,b)^3=\text{H}_1(\mathrm{\Sigma }_1)`$ such that $`Na_1+a_2=0`$. This is just $`^2`$, as each choice of $`a_1`$ and $`b`$ yields precisely one possible value of $`a_2`$, namely $`a_2=Na_1`$. $`j_{}^2`$ is the zero map and so $`_{}^2`$ is an isomorphism between this kernel and the particle charge group $`\text{H}_2(^6\times S^1,\mathrm{\Sigma }_1)`$. Therefore the group of conserved particle charges is
$$\text{H}_2(^6\times S^1,\mathrm{\Sigma }_1)=\text{Image}(_{}^2)=\text{Ker}(i_{}^1)=\{a_1,a_2,b|Na_1+a_2=0\}=^2$$
(5.9)
in line with the above field theory expectations.
The $`U(N)`$ sector has a Witten index of $`N`$, and so there are $`N`$ different equivalent vacua, which we will parametrize by a number $`r`$. In the $`r`$th vacuum the $`B`$-cycle will wrap the M-theory circle $`r`$ times, so the monopole will be confined by a charge $`r`$ string. The $`(1,r)`$ dyon, on the other hand, will not be confined, nor will a bound state of $`\text{lcm}(r,N)`$ monopoles where lcm is the least common multiple. The charge group of the strings is independent of the naming convention of the dyons, and so it is the same for all vacua. In the general case the confinement pattern will depend nontrivially on the choice of $`r`$. In Ref. this $`r`$-dependent vortex decay via monopole pair creation was refered to as magnetic screening.
### 5.4 $`U(6)U(4)\times U(2)`$ Super Yang-Mills
Next we will consider a case in which the gauge symmetry is classically broken from $`U(6)`$ to $`U(4)\times U(2)`$ (see Figure 4), where the two components reside at two extrema of, for example, a cubic superpotential.
As in all $`𝒩=1`$ pure Yang-Mills theories there will be a further dynamical symmetry breaking to an abelian group, and it is this symmetry breaking which leads to the existence of discrete vacua. In this case the vacua are labeled by two numbers $`r_1`$ and $`r_2`$, the first of which runs from 1 to 4 and the second from 1 to 2. In addition there are two kinds of Douglas-Shenker strings, those of the $`U(4)`$ and those of the $`U(2)`$, which we will call 4-strings and 2-strings respectively.
At energies much lower than the separation between the two critical points of the superpotential, this theory reduces to two decoupled pure super Yang-Mills theories with gauge groups $`U(4)`$ and $`U(2)`$. Thus bound states of 4 4-strings or 2 2-strings decay rapidly. There are W bosons whose masses are proportional to the separation between the critical points. These correspond to cylinders whose two bounding circles each travel a distance $`2\pi `$ in $`x^{10}`$ along one of the two $`A`$-cycles. The two $`A`$-cycles now wrap the M-theory circle $`2`$ and $`4`$ times respectively, thus neither of the circles that bounds the cylinder closes. To make a gauge-invariant M2 configuration we may attach each end of the cylinder to one of the Douglas-Shenker strings of the corresponding $`A`$-cycle. Thus the W bosons that are charged under the bifundamental of the $`U(4)\times U(2)`$ are confined by a 2-string and a 4-string. This means that the nucleation of a pair of such W bosons may cause a 2-string to be turned into a 4-string, or vice versa. As a result the charge $`(1,1)`$ in the maximal vortex charge group $`_4\times _2`$ is not conserved. The largest possible conserved vortex charge group is then
$$_2=\frac{_4\times _2}{\{(0,0),(1,1),(2,0),(3,1)\}}.$$
(5.10)
One may then ask whether the $`_2`$ charges themselves are conserved. The W bosons appear to conserve them, but monopole pair creation may not. A charge 1 monopole is confined by $`r_1`$ 4-strings and $`r_2`$ 2-strings, for a total charge of $`r_1+r_2_2`$. Therefore monopoles are confined by a vortex with nontrivial charge in the remaining $`_2`$ when $`r_1=r_2+1`$ mod 2, and so in this case all Douglas-Shenker strings are unstable and monopoles are confined into charge 2 magnetic baryons, which in turn may be bound to W bosons or other monopoles if it is confined by the $`_4`$ group that was broken by the W bosons. On the other hand if $`r_1=r_2`$ mod 2 then the $`_2`$ string quantum number is preserved. For example in the vacuum $`r_1=0,`$ $`r_2=0`$ the monopole is unconfined. On the other hand in the vacuum $`r_1=2,`$ $`r_2=0`$ the monopole is confined by a charge 2 4-string which carries no $`_2`$ charge, while a charge 2 monopole is unconfined, as is a bound state of a charge 1 monopole and two W bosons.
This spectrum can be read from the usual exact sequence, again with $`\text{H}_1(\mathrm{\Sigma }_1)=^3`$ but now with the induced inclusion map
$$i_{}^1:\text{H}_1(\mathrm{\Sigma }_1)=^3\text{H}_1(^6\times S^1)=:(a_1,a_2,b)4a_1+2a_2+(r_1r_2)b.$$
(5.11)
which is onto precisely when $`r_1r_2`$ is odd. Therefore the vortex charge group $`\text{H}_1(^6\times S^1,\mathrm{\Sigma }_1)`$ is $`_2`$ when $`r_1r_2`$ is even and $`0`$ when $`r_1r_2`$ is odd, while the particle charge group $`\text{H}_2(^6\times S^1,\mathrm{\Sigma }_1)`$ is again $`^2`$.
### 5.5 The General Case
In general the superpotential may classically break the gauge symmetry to a product of groups $`U(N_i)`$, $`1ik`$. By now we have seen that the exact sequence calculation of the topological charges always leads to a particle charge group
$$\text{H}_2(^6\times S^1,\mathrm{\Sigma }_{k1})=^{2k2}$$
(5.12)
where $`k1`$ of the generators are elementary W bosons corresponding to simple roots of $`U(k)`$ and the other $`k1`$ are the elementary ’t Hooft-Polyakov monopoles. There can be no torsion elements as $`j_{}^2`$ is the trivial map and there is no torsion in
$$\text{H}_1(\mathrm{\Sigma }_{k1})=^{2k1}.$$
(5.13)
In addition $`_{}^1`$ is the zero map and so the group of topological charges carried by particles is determined by the map
$$i_{}^1:^{2k1}:(a_1\mathrm{}a_k,b_1\mathrm{}b_{k1})\underset{i}{}N_ia_i+(r_{i+1}r_i)b_i$$
(5.14)
whose image is
$$\text{Image}(i_{}^1)=\mathrm{gcd}(N_i,r_{i+1}r_i).$$
(5.15)
This image of $`i_{}^1`$ is the kernel of $`j_{}^1`$ which determines the group of particle charges
$$\text{H}_1(^6\times S^1,\mathrm{\Sigma }_{k1})=\frac{\text{H}_1(^6\times S^1)}{\text{Ker}(j_{}^1)}=\frac{}{\mathrm{gcd}(N_i,r_{i+1}r_i)}=_{\mathrm{gcd}(N_i,r_{i+1}r_i)}.$$
(5.16)
In Figure 5 we see how the B-cycle can affect the stability of Douglas-Shenker strings.
Notice that the cardinality $`t`$ of the group $`\text{H}_1(^6\times S^1,\mathrm{\Sigma }_{k1})`$ is the confinement index, which is the minimal charge of an unconfined monopole, computed in . The charge $`t`$ monopole is confined by $`t`$ vortices, which carry a charge of $`t_t`$, which is equivalent to zero modulo $`t`$ and so the monopole is unconfined.
## 6 $`𝒩=1`$ $`U(N)`$ Theory with Flavored Matter
### 6.1 Super QCD: A Global Flavor Symmetry
In general the inclusion of matter leads to an instability of all Douglas-Shenker strings and often to the existence new, nontorsion, stable vortices. However we will see that sometimes theories with a gauged flavor symmetry can have stable Douglas-Shenker strings in certain vacua in which all bifundamental quarks have a bare mass.
We begin with the simplest case of fundamental matter that transforms in the fundamental representation of a global $`U(N_f)`$ flavor symmetry. The corresponding $`𝒩=2`$ theory is described by the Seiberg-Witten curve (4.9) which again describes an M5-brane extending along the gauge theory directions $`^{1,3}`$ and a Riemann surface $`\mathrm{\Sigma }M`$. In the IIA reduction this corresponds to the addition of $`N_f`$ semi-infinite D4-branes, which we will call flavor branes, that extend from the NS5$`^{}`$ to $`x^6=+\mathrm{}`$. If two flavor branes are coincident then there is an enhanced flavor symmetry, but arbitrarily small perturbations will separate the branes at sufficiently large $`x^6`$ and so a classification of stable configurations requires that all of the bare quark masses, that is all of the flavor brane positions on the $`v`$-plane, are distinct.
Softly breaking the supersymmetry to $`𝒩=1`$ with a superpotential, there is now a choice of two supersymmetric configurations for each of the $`N`$ color D4-branes that extended between the NS5’s before the deformation. Each color brane may either move along $`v`$ to an extremum of the superpotential as in the previous section, or else it may connect to a flavor brane. We will ignore configurations in which a color brane touches both a flavor brane and the NS5$`^{}`$. If a color brane located at some $`v`$ is connected to a flavor brane on one end then the two NS5’s no longer need to be coincident in the $`w`$-plane, as the D4-branes at $`v`$ no longer connect them. The minimum distance between the NS5$`^{}`$ and the locked color-flavor D4-brane pair is the VEV of the meson field corresponding to mesons built from quarks that extend between the now locked color and flavor brane. These meson VEVs Higgs the color symmetry on the locked branes, and so the remaining classical gauge symmetry group $`U(N_1)\times \mathrm{}\times U(N_k)`$ has a number of components equal to the number of minima of the superpotential at which unlocked color branes reside. Thus $`k`$ again is the number of tubes of the Riemann surface which connect the two NS5’s, and so $`k1`$ is again the genus of $`\mathrm{\Sigma }`$.
While the genus of the Riemann surface is the same as in the unflavored case, the number of punctures when one deletes the points at infinity is now $`N_f+2`$ instead of just $`2`$ as in pure super Yang-Mills. This changes the first homology group of the Riemann surface to
$$\text{H}_1(\mathrm{\Sigma })=^{2k+N_f1}$$
(6.1)
where the new $`N_f`$ generators are loops that circle the flavor branes. The inclusion map $`i_{}^1`$ multiplies these new generators by the number of coincident flavor branes at each bare mass, which is one. Therefore the inclusion map is onto and again $`j_{}^1`$ is the zero map,
$$i_{}^1(a_1\mathrm{}a_k,b_1\mathrm{}b_{k1},c_1\mathrm{}c_{N_f})=\underset{i}{}N_ia_i+(r_{i+1}r_i)b_i+c_i,j_{}^1=0.$$
(6.2)
However the triviality of $`j_{}^1`$ no longer implies that the group of vortex charges is trivial, because as we will now explain $`j_{}^1`$ is no longer onto as $`_{}^1`$ is no longer necessarily the zero map.
The fact that flavor branes extend to $`x^6=+\mathrm{}`$ with $`x^7/x^60`$ gives a fixed notion of the $`x^6`$ direction, and so in particular a relative rotation of the two NS5-branes on the $`x^6x^7`$ plane may no longer be absorbed into a coordinate redefinition. Correspondingly in the presence of fundamental matter a Fayet-Iliopoulos (FI) term may no longer be eliminated by a field redefinition. The FI term, $`r`$, corresponds to the $`x^7`$ coordinate of the NS5$`^{}`$. All of the color D4’s extend in the $`x^6`$ direction but must remain at constant $`x^7`$ to preserve supersymmetry, otherwise they would not be parallel to the flavor D4-branes. Thus when $`r0`$ all of the color branes extend from the NS5-brane to $`x^6=+\mathrm{}`$ without touching the NS5$`^{}`$, and so every color is locked. This does not imply that all flavors are locked, unlocked flavors remain at a constant $`x^7=r`$ as they extend from the NS5$`^{}`$ to $`x^6=\mathrm{}`$. As every brane occupies a constant $`x^7`$ position, if the two NS5’s are at different $`x^7`$ positions then the entire configuration is disconnected. In general a high energy gauge symmetry consisting of $`j`$ gauge group components may be engineered using $`j+1`$ NS5’s and if the FI terms are independent then the entire configuration will consist of $`j+1`$ independent components.
We will now restrict attention to configurations with only two NS5’s and so at most two components. Thus $`\text{H}_0(\mathrm{\Sigma })`$ will be $`^2`$ if all colors are locked, but if any color is unlocked then the corresponding D4 connects the two components of the M5 and so $`\text{H}_0(\mathrm{\Sigma })=`$. The spacetime is still connected and so $`\text{H}_0(^6\times S^1)=`$ as in pure super Yang-Mills. This means that the map $`i_{}^0`$ is no longer into when all colors are locked, instead $`i_{}^0=(1,1)`$, and so as suggested above $`_{}^10`$. On the contrary the image of $`_{}^1`$ is the kernel of $`i_{}^0`$, which is the additional $``$ component. This $``$ must then appear in its domain, the group of vortex charges
$$\text{H}_1(^6\times S^1,\mathrm{\Sigma })=\{\begin{array}{cc}\text{if all colors are locked}\hfill & \\ 0\text{if any color is unlocked.}\hfill & \end{array}$$
(6.3)
Thus while the Douglas-Shenker vortices may decay via the nucleation of fundamental matter, new nontorsion vortices arise in some cases. Such vortices have been studied for example in Refs. , where it was seen that monopoles in a theory with an FI term are confined by two vortices. This means that such a vortex cannot decay by monopole-antimonopole nucleation as monopoles are merely kinks in vortex worldvolumes. Rather than terminating a vortex, these kinks transmute a vortex into a different type of vortex that carries the same conserved charge.
The vortices studied in Refs. , in theories with a superpotential rather than an FI term, are not topologically stable as some of the colors are unlocked. As a result there are monopoles confined by a single vortex whose pair creation causes these vortices to decay (In Figure 6 we have an example of this kind).
This does not imply that an FI term is required for stability. The vortices of those theories would have been stable had one not included the unlocked color, in which case the 1-monopoles would have been removed from the spectrum. Alternately one could have added a new flavor with a different bare mass and changed the adjoint scalar VEV of the unlocked color to the bare mass of the new flavor, locking them together. The resulting ’t Hooft-Polyakov monopoles would have been confined by two vortices and thus would not have led to magnetic screening. Note that unlike the vortices of Ref. , such vortices are not be BPS, but the topological charge is nonetheless conserved. In conclusion, matter that transforms in the fundamental of a global flavor symmetry screens Douglas-Shenker string charge, but when the M5-brane is disconnected a new kind of vortex is stable.
The group of particles is also augmented as a result of the extra $`^{N_f}`$ in Eq. (6.1). As the dimension of $`\text{H}_1(^6\times S^1)`$ is still $``$, the extra dimensions must come from the group of particle charges
$$\text{H}_2(^6\times S^1,\mathrm{\Sigma })=^{2k+N_f2}.$$
(6.4)
The extra $`N_f`$ independent particles are the quarks with a given color. A quark of a different color will not carry an independent charge, as its charge is carried by a boundstate of a quark of the given color and a W boson.
### 6.2 A Local Flavor Symmetry
The Douglas-Shenker strings of super QCD are unstable because each quark is defined by a single string, and so a pair of quarks produced on the string leads to a gap which grows until the string has disappeared. Topologically this reflects the fact that $`i_{}^1`$ is onto because each flavor brane wraps the M-theory circle once and so $`i_{}^1`$ multiplies the quark charges by one.
The situation is quite different if the flavor symmetry is gauged and confining. In this case there will be two species of Douglas-Shenker strings, those of the color group and those of the flavor group. A quark will be confined by one of each, and so quark pair-productions, like W boson pair-production in SYM, will only cause vortices to be transmuted into different types of vortices rather than to disappear altogether. For example a confining $`U(N)`$ color symmetry and $`U(M)`$ local flavor symmetry lead to $`_N`$ charged and $`_M`$ charged vortices. A $`_N`$ vortex becomes a $`_M`$ vortex when one crosses a quark, and so in particular the charge $`M`$ $`_N`$ vortex is unstable. As in the case of $`U(M)\times U(N)`$ super Yang-Mills, the group of conserved charges is then
$$\text{H}_1(^6\times S^1,\mathrm{\Sigma }_0)=_{\mathrm{gcd}(M,N)}$$
(6.5)
where $`\mathrm{\Sigma }_0`$ is the lift of three NS5-branes, one pair connected by $`N`$ D4’s and the other by $`M`$ D4’s. Topologically this is a result of
$$i_{}^1(a_1,a_2)=Ma_1+Na_2$$
(6.6)
whose cokernel is $`_{\mathrm{gcd}(M,N)}`$, just as in the case of pure SYM with low energy gauge group $`U(M)\times U(N)`$ where the W bosons play the role played by the quarks here.
This is not to suggest that gauging the flavor symmetry always stabilizes some of the Douglas-Shenker strings. For example if any classically unbroken gauge subgroup $`U(N_c)`$ is in its baryonic root vacua then locally $`\mathrm{\Sigma }`$ factorizes into two funnels which intersect at $`2N_cN_f`$ distinct points. The loop which begins at one intersection, travels along one sheet to the next, and then returns along the other sheet wraps the M-theory circle precisely once. Thus the image of the corresponding generator of $`\text{H}_1(\mathrm{\Sigma })`$ will be the element $`1\text{H}_1(^6\times S^1)`$. This means that $`i_{}^1`$ is again onto and no vortices carry any conserved topological charge. The confined particle must correspond to an M2 that has a boundary that encircles this loop, although it may then encircle a second loop which gives it a net winding number around the M-theory circle of zero. If the net winding number is zero then the M2 could have disk topology which would allow the particle to be a hypermultiplet . This theory contains flavored magnetic monopoles, which in the Seiberg dual IR free theory become unconfined quarks, and so do not connect to strings.
When the flavor symmetry is gauged, for example in the Klebanov-Strassler theory in which $`x^6`$ is periodic, the UV theory may be strongly coupled and in fact potentially continues to cascade. However the above confined particles exist at each Seiberg duality, as the above cycle may be constructed each time the two M5-brane sheets pass each other. Thus the vortices confine particles at the energy scale of each step in the cascade, and the vortices will be broken by the lightest particles which come from the last step of the cascade.
### 6.3 Gubser-Herzog-Klebanov Axionic Vortices
Gubser, Herzog and Klebanov have conjectured the existence of a new kind of vortex in the Klebanov-Strassler $`SU(N+M)\times SU(N)`$ theory. This corresponds to a D1-brane on the conifold which extends in the gauge theory directions. It does not extend in any internal directions, and so after T-dualizing an internal direction to arrive in IIA it will be a D2-brane extending along the T-dual circle, $`x^6`$. Lifting to M-theory one finds an internal spacetime of $`^5\times T^2`$, as in the $`𝒩=4`$ example above. As usual in the IIA reduction there are two NS5-branes, which are at two different coordinates on the $`x^6`$ circle.
There are now two inequivalent embeddings of the D2-brane extended along $`x^6`$ and two gauge theory directions. It may either wrap all of $`x^6`$, or else it may extend between two NS5’s. In the first case it will be T-dual to a D-string, while in the second it will be T-dual to a half D-string at the conifold singularity which blows up into a D3 when the conifold is resolved. Both of these branes are non-BPS, as they share three common directions with the D4-branes. In fact both kinds of brane attract the D4-brane, and once they become coincident will dissolve into the D4. In the D4 worldvolume theory these are magnetic flux tubes. The first brane carries a unit of magnetic flux in both gauge groups while the second carries flux in only one. These flux tubes will smear out of existence unless they are confined for some reason, for example if they really are vortices in an unexpected condensate field in the D4 worldvolume. In the large $`N`$ IIB description perhaps the large $`N`$ limit, which places them far from the horizon’s D3-branes that are T-dual to the D4’s here, leads to a long lifetime for these strings.
The possibility that these D2’s may dissolve into the D4’s follows from the topological classification above. There is no extra $``$ component in the group of vortex charges that could stabilize these axionic string charges. In fact the M5-branes of these models are connected, and so there are no vortices at all carrying topological charges of the type classified in this note.
At distance scales much smaller than the size of the smallest $`A`$-cycle, corresponding to the $`SU(K)`$ theory at the bottom of the cascade, one may ignore the D4 and so miss the instability of these vortices. The decay is caused by the nucleation of degrees of freedom from the $`SU(K+M)\times SU(M)`$ UV completion and so involves the nucleation of particles of mass equal to the UV cutoff. Similarly the characteristic lifetime of the vortices will be exponential in the UV cutoff. This is in contrast with the case of ordinary $`SU(K)`$ pure super Yang-Mills, which is asymptotically free. The $`SU(K)`$ theory at the bottom of the cascade is not asymptotically free. The fact that both $`SU(K)`$ theories are in the same universality class does not preclude the existence of stable vortices in one theory and not the other, as the topological stability of a configuration depends on the UV physics.
## 7 Conclusions
In field theories that can be engineered from M5-branes the relative homology groups of the embedding of the M5 yield conserved charges. In particular we have reproduced the confinement index formula of Ref. and we have found the charges corresponding to torsion Douglas-Shenker strings as well as BPS Hanany-Tong strings. Strings obtained, for example, by softly breaking $`𝒩=2`$ SQCD to $`𝒩=1`$ with a mass for the adjoint chiral multiplet are found, as expected, to be unstable. However the conservation of the homology charge led us to conjecture that such vortices may be stabilized by adding a correction to the superpotential that leads the corresponding M5 to be disconnected. Further we have seen that the M5 always is disconnected, and so the vortices are always stable, in supersymmetric backgrounds with a nonzero FI term and any superpotential. In both the superpotential and FI cases the disconnected M5 only preserves supersymmetry if all colors are locked to flavors.
Noticably absent from the spectrum of stable vortices is the GHK axionic string. We have T-dualized the conifold realization of this string to a IIA brane cartoon. In doing so we have found that there are in fact two species of this string and that both can be continuously deformed into the M5-brane, where no charge that we have seen prevents them from smearing into oblivion.
The topological charge construction considered here, and first proposed in Ref. , applies to a wide variety of theories. For example the index of confinement can easily be calculated for theories with various kinds of fundamental and bifundamental matter. Usually the index will be trivial in these cases. However we have seen that in some theories in which the flavor group is gauged and confining some of the Douglas-Shenker strings are stable and so the confinement index will not be trivial.
The construction itself is much more general than this, allowing treatment of gauge theories in various numbers of dimensions and even of higher (lower) form gauge theories as in the $`𝒩=4`$ example above. The only critical assumption was the validity of the supergravity approximation, and in particular the classical treatment of the geometry. This is a weaker condition than the supersymmetry of the configuration, for example we have found non-BPS torsion strings and also dyons that cannot satisfy the BPS bound. However we have not found a criterion that allows one to determine when this approximation applies.
One may attempt to use relative homology to classify the charges of objects in the presence of a domain wall, which is a 3-dimensional cobordism that interpolates between two M5 embeddings $`\mathrm{\Sigma }`$. Objects with a boundary on the domain wall will be bound to the wall.
## Acknowledgement
The work of JE is partially supported by IISN - Belgium (convention 4.4505.86), by the “Interuniversity Attraction Poles Program – Belgian Science Policy” and by the European Commission RTN program HPRN-CT-00131, in which he is associated to K. U. Leuven. Stefano is supported as well. |
warning/0506/astro-ph0506220.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Cosmic rays (CRs) were discovered almost a century ago, yet their origin is unknown. There is sufficient evidence that at least part of their spectrum ($`E<10^{18}eV`$) originates in the Galaxy, while particles of higher energies are thought to come to us from extragalactic sources. One simple theoretical argument is that particles with such high energies (usually referred to as ultra high energy CRs, UHECR) could not be confined to the Galaxy, on account of their large gyroradius. Second, their (power law- $`E^{2.7}`$) spectrum is harder than that of (presumably) galactic particles at $`E<10^{18}eV`$ ($`E^{3.1})`$, which is consistent with their extra-galactic origin. The last point becomes clear if we turn to the other break on the overall CR spectrum, the one at $`E10^{15}eV`$, commonly known as the “knee”, shown in Fig.1. The spectrum above this energy steepens (from $`E^{2.7}`$ to $`E^{3.1}`$), so that the premise of their extragalactic origin would require an explanation of why the galactic part of the spectrum terminates, while the extragalactic part appears *exactly* at the knee energy. As we shall see, to explain the cosmic ray power law spectrum between the “knee” at$`10^{15}eV`$ and the “ankle” at $`10^{18}eV`$ in terms of acceleration within the Galaxy is one of the most serious challenges of contemporary acceleration mechanisms and one of the main motivations of this study.
The explanation of the spectrum beyond the ankle (the highest energy event observed so far is $`310^{20}eV`$) poses a major challenge to fundamental physics. Given the distance of possible accelerators (at least a few tens of Mpc), particles of such high energy should have experienced significant losses through their interaction with the cosmic microwave background radiation (the so called Greisen, Zatsepin, Kusmin or GZK cut-off) while propagating over long distances. Note that UHECR are sub-atomic particles with the energy of a well-hit baseball. This energy exceeds (by at least three orders of magnitude) that achievable by all existing terrestrial accelerators ( e.g., Large Hadron Collider, LHC-$`10^{17}eV`$). Another fundamental aspect of the problem of CR origin is that they are the ultimate receptacle of a significant portion of the gravitational energy in the Universe. Indeed, star formation from the gravitational collapse of the primordial gas with subsequent SN (supernova) explosions and their blast waves results in CR acceleration. On the observational side, there is an evidence (in the form of both synchrotron and inverse Compton radiation) that electrons of energies up to $`100TeV`$ are accelerated in the supernova shock waves.
## 2 Acceleration mechanism
The leading CR acceleration mechanism, namely the diffusive shock acceleration (DSA, also known as the 1st -order Fermi mechanism) was proposed originally by Fermi in Ref., and in its modern form by a number of authors in the late seventies . The mechanism is basically simple – particles gain energy by bouncing between converging upstream and downstream regions of the flow near a shock wave such as that from an SNR (supernova remnant) shock. This mechanism requires magnetic fields. First, the field binds particles to the accelerator (shock wave), in the direction perpendicular to the field. Confinement in the direction *along* the field lines is, in turn, ensured by particles themselves through the generation of Alfven waves by accelerated particles streaming ahead of the shock. This occurs via Doppler resonance $`\mathrm{\Omega }=k(p/m)\mu `$, where $`\mathrm{\Omega }`$ and $`p`$ are the (nonrelativistic) gyrofrequency and momentum, $`k`$ is the wave number, $`m`$ is the particle mass, and $`\mu `$ is the cosine of its pitch angle. These waves, in turn, scatter particles in pitch angle (at the rate $`\nu \mathrm{\Omega }(mc/p)\left(\delta B/B_0\right)^2`$, where $`c`$ is the speed of light) back and force so that they can gain energy by repeatedly crossing the shock. Particle self-confinement along the field is thus diffusive and the diffusivity $`\kappa c^2/\nu `$ is inversely proportional to the fluctuation energy $`\delta 𝐁^2`$, as the fluctuating field is responsible for pitch-angle scattering. However, the mean field $`B_0`$ ultimately determines the acceleration rate and the particle maximum energy since it sets the work done by the induced electric field on the particles. Indeed, $`E_{max}(e/c)u_sB_0R_s`$, where $`u_s`$ and $`R_s`$ are the speed and typical size of the shock wave, such as the radius of the SNR shock. The fluctuating part, $`\delta B`$, is typically assumed to be saturated at most at the level $`\delta BB_0`$, which provides pitch-angle scattering at the rate of gyrofrequency, and thus limits the particle mean free path (m.f.p.) along the field to a distance of the order of gyroradius (the so called Bohm diffusion limit). An important thing to keep in mind is that, due to the resonance condition $`kp=const`$, confinement of particles of higher energies requires that longer waves need be excited.
The most critical test of this mechanism is the requirement that it accelerates galactic CRs to the energy of $`10^{15}`$ eV over the life-time of supernova remnant shocks. Even with the above “optimistic” estimates of the turbulence level, the mechanism passes this test at best only marginally. If the turbulence level is lower, then the maximum energy should be reduced proportionally. There are indeed a number of phenomena which may either reduce the turbulence level , or which can shift the turbulence spectrum (in wave number) away from resonance with the high energy particles and therefore cause their losses .
Another reason for concern about this mechanism, at least in its standard (Bohm limit) version, is its seeming inability to explain acceleration of particles beyond $`10^{15}`$eV. As was discussed above, the cosmic ray spectrum has only a break at this energy, and continues to about $`10^{18}`$ eV where the extragalactic component is believed to start dominating the spectrum.
One approach to this problem is to generate a fluctuating component $`\delta B`$ significantly exceeding the unperturbed field $`B_0`$ . Physically, such generation is possible since the free energy source is the pressure gradient of accelerated particles, which in turn may reach a significant fraction of the shock ram energy. Specifically, the wave energy density $`\left(\delta B/B_0\right)^2`$ may be related to the partial pressure $`P_c`$ of CRs that resonantly drives these waves through the relation
$$\left(\delta B/B_0\right)^2M_AP_c/\rho u_s^2$$
(1)
where $`M_A=u_s/V_A1`$ is the Alfven Mach number and $`\rho u_s^2`$ is the shock ram pressure. Of course, when $`\delta B/B_0`$ exceeds unity, particle dynamics, and thus their confinement and acceleration rates, are very difficult to assess if the turbulence spectrum is sufficiently broad. The numerical studies by showed that at least in the case of an MHD (magnetohydrodynamic) description of the background plasma and rather narrow wave (and particle energy) band, the amplitude of the principal mode can reach a few times that of the background field. Moreover, the authors of Ref. argue that in the case of efficient acceleration, field amplification may be even stronger, reaching a $`mG`$ ($`10^3`$ Gauss) level from the background of a few $`\mu G`$ ($`10^6`$ Gauss) ISM field, thus providing acceleration of protons up to $`10^{17}`$ eV in SNRs.
Recently, the authors of Ref. approached this problem from a different perspective. They considered a Kolmogorov turbulent cascade to small scales assuming the waves are generated by efficiently accelerated particles on the long-wave part of the spectrum. They obtained a particle maximum energy similar to that of .
Apart from the excitation of magnetic fluctuations during acceleration process, there is yet another aspect of the CR– magnetic field connection discussed in the literature. Zweibel points out that since CRs were already present in young galaxies (observed at high redshifts), magnetic fields of appreciable strength must also have been there at that time. She emphasizes, however, that the approximate equipartition between the CR and magnetic field energy established by the current epoch in our Galaxy is not required for the acceleration mechanism and presumably results from the fact that they both have a common energy source, namely the supernovae. Indeed, as we discussed already, the magnetic field strength merely determines the *maximum* energy of accelerated particles, given the time available for acceleration and the size of the accelerator. If the latter are sufficient then the total energy of accelerated particles can, and in most of the DSA models *does,* exceed that of the magnetic field. The latter, in turn, remains unchanged, apart from the conventional compression at the shock and the MHD fluctuations discussed earlier.
In this paper we discuss the possibility of a *different* scenario, in which the magnetic field may absorb a significant part of the shock energy as a result of the acceleration process, which may in fact be strongly enhanced. The mechanism of such enhancement is based on the transfer of magnetic energy to longer scales, which we call *inverse cascade* for short, even though specific mechanisms of such transfer may differ from what is usually understood as a cascade in MHD turbulence. This transfer is limited only by some outer scale $`L_{out}`$ such as the shock precursor size $`\kappa (p_{max})/u_sr_g(p_{max})c/u_sr_g(p_{max})`$. This approach is in contrast to the above discussed models , which operate on the generated magnetic fields with the scale lengths of the order of the Larmor radius $`r_g(p_{max})`$ of the highest energy particles and smaller. The advantage of the inverse cascade for the acceleration is that the turbulent field at the outer scale $`\delta B(L_{out})B_{rms}`$ (which necessarily must have long autocorrelation time) can be obviously regarded as an “ambient field” for accelerated particles of all energies. If $`B_{rms}B_0`$, then the acceleration can be enhanced by a factor $`B_{rms}/B_0`$. Note that the resonance field $`\delta B(r_g)`$ may remain smaller than $`B_{rms}`$, so that standard arguments about Bohm diffusion apply, and it is less likely that this field will be rapidly dissipated by nonlinear processes, such as induced scattering on thermal protons, not included in the enhanced acceleration model .
As it should be clear from the above, an adequate description of the acceleration mechanism must include *both* particle and wave dynamics on an equal footing. In fact the situation is even more difficult, since the acceleration process turns out to be so efficient that the pressure of accelerated particles markedly modifies the structure of the shock (both the overall shock compression and the flow profile).
## 3 Accelerated Particles and plasma flow near the shock front
The transport and acceleration of high energy particles (CRs) near a CR modified shock is usually described by the diffusion-convection equation. It is convenient to use a distribution function $`f(p)`$ normalized to $`p^2dp`$.
$$\frac{f}{t}+U\frac{f}{x}\frac{}{x}\kappa \frac{f}{x}=\frac{1}{3}\frac{U}{x}p\frac{f}{p}$$
(2)
Here $`x`$ is directed along the shock normal, which for simplicity, is assumed to be the direction of the ambient magnetic field. The two quantities that control the acceleration process are the flow profile $`U(x)`$ and the particle diffusivity $`\kappa (x,p)`$. The first one is coupled to the particle distribution $`f`$ through the equations of mass and momentum conservation
$`{\displaystyle \frac{}{t}}\rho +{\displaystyle \frac{}{x}}\rho U=0`$ (3)
$`{\displaystyle \frac{}{t}}\rho U+{\displaystyle \frac{}{x}}\left(\rho U^2+P_\mathrm{c}+P_\mathrm{g}\right)=0`$ (4)
where
$$P_c(x)=\frac{4\pi }{3}mc^2_{p_{inj}}^{\mathrm{}}\frac{p^4dp}{\sqrt{p^2+1}}f(p,x)$$
(5)
is the pressure of the CR gas and $`P_\mathrm{g}`$ is the thermal gas pressure. The lower boundary $`p_{inj}`$ in momentum space separates CRs from the thermal plasma that is not directly involved in this formalism but rather enters the equations through the magnitude of $`g`$ at $`p=p_{inj}`$, which specifies the injection rate of thermal plasma into the acceleration process. The particle momentum $`p`$ is normalized to $`mc`$. The spatial diffusivity $`\kappa `$, induced by pitch angle scattering, prevents particle streaming away from the shock, thus facilitating acceleration by ensuring the particle completes several shock crossings.
The system (2-5) indicate a marked departure from the test particle theory. Perhaps the most striking result of the nonlinear treatment is the bifurcation of shock structure (in particular shock compression ratio) in the parameter space formed by the injection rate, shock Mach number and particle maximum momentum .
## 4 Wave dynamics in the CR shock precursor
The transformation of magnetic energy to longer scales, while bearing certain characteristics of the conventional turbulent dynamo problem, is still rather different from it, in its conventional form. First, this process should take place in the strongly compressible fluid near the shock. Second, the Alfven wave turbulence is generated by accelerated particles via Cerenkov emission, and thus is strongly coupled to them. Third, the shock precursor itself is unstable to emission of acoustic waves. The latter phenomenon is known as the Drury instability and will be discussed later. Acoustic waves, in turn, interact with particle generated Alfvenic turbulence, stimulating decay instability ( i.e., “inverse cascade”).
The spatial structure of an efficiently accelerating shock, i.e., the shock that transforms a significant part of its energy into accelerated particles, is very different from that of the ordinary shock, Fig.2. Most of the shock structure consists of a precursor formed by accelerated CRs diffusing ahead of the shock. If the CR diffusivitity $`\kappa (p)`$ depends linearly on particle momentum $`p`$ (as in the Bohm diffusion case), then, at least well inside the precursor, the velocity profile $`U(x)`$ is approximately a *linear* function of $`x`$, where $`x`$ points along the shock normal . Ahead of the shock precursor, the flow velocity tends to its upstream value, $`U_1`$, while on the downstream side it undergoes a conventional plasma shock transition to its downstream value $`U_2`$ (all velocities are taken in the shock frame). This extended CR precursor (of the size $`L_{CR}\kappa (p_{max})/U_1`$) is the place where we expect turbulence is generated by the CR streaming instability and where it cascades to longer wavelengths.
### 4.1 Alfvenic turbulence
The growth rate of the ion-cyclotron instability is positive for the Alfven waves traveling in the CR streaming direction i.e., upstream, and it is negative for oppositely propagating waves. The wave kinetic equation for both types of waves can be written in the form
$$\frac{N_k^\pm }{t}+\frac{\omega ^\pm }{k}\frac{N_k^\pm }{x}\frac{\omega ^\pm }{x}\frac{N_k^\pm }{k}=\gamma _k^\pm N_k^\pm +C^\pm \{N_k^+,N_k^{}\}$$
(6)
Here $`N_k^\pm `$ denotes the number of quanta propagating in the upstream and downstream directions, respectively. Also, $`\omega ^\pm `$ are their frequencies, $`\omega ^\pm =kU\pm kV_A`$, where $`V_A`$ is the Alfven velocity. The linear growth rates $`\gamma ^\pm `$ are nonzero only in the resonant part of the spectrum, $`kr_g(p_{max})>1`$. In the most general case, the last term on the r.h.s. represents nonlinear interaction of different types of quanta $`N_k^+`$ and $`N_k^{}`$ and, if compressibility effects are present, also interactions between the same type. As seen from this equation, the coefficients in the wave transport part of this equation (l.h.s.) depend on the parameters of the medium through $`U`$ and $`V_A`$, which in turn, may be subjected to perturbations. This usually results in parametric phenomena . We will concentrate on the acoustic type perturbations (which may be induced by Drury instability), so that we can write for the density $`\rho `$ and velocity $`U`$
$$\rho =\rho _0+\stackrel{~}{\rho };U=U_0+\stackrel{~}{U}$$
The variation of the Alfven velocity $`\stackrel{~}{V}_A=V_AV_{A0}`$
$$\stackrel{~}{V}_A\frac{1}{2}V_A\frac{\stackrel{~}{\rho }}{\rho _0}.$$
For simplicity, we assume that the plasma $`\beta <1`$ (which is not universally true in the shock environment) and neglect the variation of $`U`$ compared to that of $`V_A`$ in eq.(6). The above perturbations of $`V_A`$ in turn induce perturbations of $`N_k^\pm `$, so we can write
$$N_k^\pm =N_k^\pm +\stackrel{~}{N}_k^\pm $$
Our goal is to obtain an evolution equation for the averaged number of plasmons $`N_k^\pm `$. Averaging eq.(7) we have
$$\frac{}{t}N_k^\pm +\left(U\pm V_A\right)\frac{}{x}N_k^\pm kU_x\frac{}{k}N_k^\pm +\frac{}{k}kV_A\frac{\stackrel{~}{\rho }_x}{\rho _0}\stackrel{~}{N}_k^\pm =\gamma _k^\pm N_k^\pm +C\left(N_k^\pm \right)$$
(7)
Here the index $`x`$ stands for the $`x`$-derivatives. To calculate the correlator $`\frac{\stackrel{~}{\rho }_x}{\rho _0}\stackrel{~}{N}_k^\pm `$in the last equation, we expand the r.h.s. of eq.(6) retaining only the main linear part in $`\stackrel{~}{N}`$
$$\gamma _k^\pm N_k^\pm +C^\pm \{N_k^+,N_k^{}\}\mathrm{\Delta }\omega _k^\pm \stackrel{~}{N}_k^\pm $$
(8)
The time scale separation between the l.h.s. and r.h.s. of eq.(6) suggests that to lowest order, the linear growth $`\gamma ^+`$ rate is approximately balanced by the local nonlinear term $`C^+`$. Likewise, the linear damping of the backward waves $`\gamma ^{}`$may be balanced by their nonlinear growth and conversion of the forward waves $`C^{}`$. Generally, the $`\mathrm{\Delta }\omega _k^\pm `$ in eq.(8) is a $`2\times 2`$ matrix operator. If the wave collision term is quadratic in $`N`$, then $`\mathrm{\Delta }\omega _k^\pm \gamma _k^\pm `$. This balance can be established only for the resonant waves ($`\gamma ^\pm 0`$), whereas our primary focus will be on the extended longwave interval $`k<1/r_g(p_{max})`$ for which $`\gamma 0`$. In this domain, cascading from the generation region $`k>1/r_g(p_{max})`$ takes place and the refraction (last) term on the l.h.s. of eq.(7) plays a dominant role along, with the nonlinear term on the r.h.s.
To calculate the refraction term we write eq.(6), linearized with respect to $`\stackrel{~}{N}_k^\pm `$, as:
$$L^\pm \stackrel{~}{N}_k^\pm =kV_A\frac{\stackrel{~}{\rho }_x}{2\rho _0}\frac{}{k}N_k^\pm $$
(9)
where
$$L^\pm =\frac{}{t}+\left(U\pm V_A\right)\frac{}{x}kU_x\frac{}{k}+\mathrm{\Delta }\omega _k^\pm $$
Solving eq.(9) for $`\stackrel{~}{N}_k^\pm `$, from eq.(7) we thus have the following equation for $`N_k^\pm `$
$$\frac{}{t}N_k^\pm +U\frac{}{x}N_k^\pm kU_x\frac{}{k}N_k^\pm \frac{}{k}D\frac{}{k}N_k^\pm =\gamma _k^\pm N_k^\pm +C\left(N_k^\pm \right)$$
(10)
Here we introduced a diffusion operator for the Alfven waves in $`k`$ space due to *random* refraction by the acoustic perturbations $`\stackrel{~}{\rho }`$ (via the density dependence of $`V_A`$), i.e.,
$$D_k=\frac{1}{4}k^2V_A^2\frac{\stackrel{~}{\rho }_x}{\rho _0}L^1\frac{\stackrel{~}{\rho }_x}{\rho _0}$$
(11)
$`D_k`$ is an example of the well-known phenomenon of induced diffusion. Transforming to Fourier space, we first represent $`\stackrel{~}{\rho }`$ as
$$\stackrel{~}{\rho }=\underset{q}{}\rho _qe^{iqxi\mathrm{\Omega }_qt}$$
and note that due to the local Galilean invariance of $`L`$, we can calculate its Fourier representation in the reference frame moving with the plasma at the speed $`U(x)`$ as:
$$L_{k,q}^\pm =\pm iqV_A+\mathrm{\Delta }\omega _k^\pm kU_x\frac{}{k}$$
(12)
Then, eq.(11) can be re-written as:
$$D_k=\frac{1}{2}k^2V_A^2\underset{q}{}q^2\left|\frac{\rho _q}{\rho _0}\right|^2\mathrm{}L_{k,q}^1$$
(13)
The last (wave refraction) term on the r.h.s. of eq.(12) can be estimated as $`U_1^2/\kappa (p_{max})`$, which is the inverse acceleration time and can be neglected as compared to the frequencies $`qV_A`$ and $`\mathrm{\Delta }\omega `$. Hence, for $`\mathrm{}L_{k,q}^{\pm 1}`$ we have:
$$\mathrm{}L_{k,q}^{\pm 1}\frac{\mathrm{\Delta }\omega _k^\pm }{q^2V_A^2+\mathrm{\Delta }\omega _k^{\pm 2}}$$
For further convenience, we introduce here the number of phonons
$$N_q^s=\frac{W_q}{\omega _q^s}$$
where $`W_q`$ is the energy density of acoustic waves (with $`\omega _q^s=qC_s`$ as their frequency).
$$W_q=C_s^2\frac{\rho _q^2}{\rho _0}$$
For $`D_k`$ in eq.(10) we thus finally have
$$D_k=\frac{k^2V_A^2}{2C_s^2\rho _0}\underset{q}{}q^2\omega _q^s\frac{\mathrm{\Delta }\omega _k^\pm }{q^2V_A^2+\mathrm{\Delta }\omega _k^{\pm 2}}N_q^s$$
Note that $`D_k`$ represents the rate at which the wave vector of the Alfven wave random walks due to stochastic refraction. Of course, such a random walk necessarily must generate larger scales (smaller $`k`$), thus in turn facilitating the confinement (to the shock) of higher energy particles. Thus, confinement of higher energy particles is a natural consequence of Alfven wave refraction in acoustic wave generated density perturbations.
### 4.2 Acoustic turbulence
Unlike the Alfvenic turbulence that originates in the shock precursor from accelerated particles, there are two separate sources of acoustic perturbations. One is related to parametric processes undergone by the Alfven waves in the usual form of a decay of an Alfven wave into another Alfven wave and an acoustic wave. The other source is the pressure gradient of CRs, which directly drives instability. The latter leads to emission of sound waves due to the Drury instability. By analogy with eq.(7) we can write the following wave kinetic equation for the acoustic waves:
$$\frac{}{t}N_q+U\frac{}{x}N_qqU_x\frac{}{k}N_q=\left(\gamma _q^d+\gamma _D\right)N_q+C\left\{N_q\right\}$$
Here $`\gamma _D`$ is the Drury instability growth rate and $`\gamma _q^d`$ is that of the decay instability. We first consider the decay instability of Alfven waves. Note, however, that the combination of Drury instability and decay instability can lead to generation of mesoscale fields at a faster than – exponential rate, by coupling together the Drury and decay instability processes.
#### 4.2.1 Decay Instability
The mechanism of this instability is the growth of the density (acoustic) perturbations due to the action of the ponderomotive force from the Alfven waves. This force can be regarded as a radiative pressure term appearing in the hydrodynamic equation of motion for the sound waves (written below in the comoving plasma frame)
$$\frac{V}{t}=\frac{1}{\rho _0}\frac{}{x}\left(C_s^2\stackrel{~}{\rho }+P_{rad}\right)$$
Eliminating velocity by making use of continuity equation
$$\frac{\stackrel{~}{\rho }}{t}+\rho _0\frac{V}{x}=0,$$
we obtain
$$\frac{^2\stackrel{~}{\rho }}{t^2}=\frac{^2}{x^2}\left(C_s^2\stackrel{~}{\rho }+P_{rad}\right)$$
(14)
The Alfven wave pressure can be expressed through their energy
$$P_{rad}=\underset{k}{}\omega _k\left(\stackrel{~}{N}_k^++\stackrel{~}{N}_k^{}\right)$$
Using the relation (9) between the density perturbations and the Alfven waves and separating forward and backward propagating sound waves $`\rho ^\pm `$, we can obtain from eq.(14) the following dispersion relation for the acoustic branch
$$\omega ^2q^2C_s^2=q^2\underset{k}{}\frac{\omega _k}{2\rho _0}iqkV_AL_{k,q}^{\pm 1}\frac{}{k}N_k^\pm $$
or on writing $`\omega =\pm qC_s+i\gamma ^\pm `$, we have the following growth rate of acoustic perturbations
$$\gamma ^\pm =\frac{q^2}{4\rho _0}\frac{V_A}{C_s}\underset{k}{}k\omega _kL_{k,q}^{\pm 1}\frac{}{k}N_k^\pm $$
Note that the instability requires an inverted population of Alfven quanta. As they are generated by high energy resonant particles in a finite domain of $`k`$ space, such an inversion clearly can occur.
#### 4.2.2 Drury Instability
This instability also leads to efficient generation of acoustic waves and it is driven by the pressure gradient of the CRs in the shock precursor. The growth rate has been calculated in Ref. (see also ), and can be written in the form:
$$\gamma _D^\pm =\frac{\gamma _CP_C}{\rho \kappa }\pm \frac{P_{Cx}}{C_s\rho }\left(1+\frac{\mathrm{ln}\kappa }{\mathrm{ln}\rho }\right)$$
(15)
Here $`P_C`$ and $`P_{Cx}`$ are the CR pressure and its derivative, respectively, and $`\gamma _C`$ is their adiabatic index. For an efficiently accelerating shock $`\gamma _C4/3`$. Note that we have omitted a term $`U_x`$ which is related to a simple compression of wave number density by the flow and should be generally incorporated into the r.h.s. of eq.(14) (see ). This term is smaller by a factor $`C_s/U`$ than the second (destabilizing) term. The first term is damping caused by CR diffusion, calculated earlier by Ptuskin .
## 5 Mechanisms of transfer of magnetic energy to larger scales
As it follows from the above considerations, there are a variety of nonlinear processes that can lead to the transfer of magnetic energy (generated by accelerated particles in form of the resonant Alfven waves) to longer scales. First, as it can be seen from eq.(10) (last term on the l.h.s.) scattering of the Alfven waves in $`k`$ -space due to acoustic perturbations transfers magnetic fluctuations away from the resonant excitation region to smaller (but also to larger) $`k`$. Second, the nonlinear interaction of Alfven waves and magnetosonic waves represented by the wave collision term on the r.h.s. can be responsible for such process. It is also well known that in the presence of nonzero magnetic helicity there is a strong inverse cascade of magnetic energy . Finally, even in the frame work of the weak turbulence, induced scattering of Alfven waves on thermal protons leads to a systematic decrease in the energy of quanta which, given the dispersion law, means again transformation to longer waves.
Returning to the wave refraction process on acoustic perturbations generated by the Drury instability, it is important to emphasize the following. As is seen from the instability growth rate, eq.(15), it is proportional to the gradient of $`P_c`$. As the latter should be increased as a result of instability, through a better confinement and faster acceleration, this will reinforce the instability, possibly triggering “explosive” growth. This can significantly contribute to mechanisms of regulation of $`P_c`$ discussed earlier in .
## 6 Conclusions and Discussion
We have considered a number of possible mechanisms for generation of large scale magnetic field in front of strong astrophysical shocks. All these mechanisms are immediate results of the particle acceleration process. Such generation is necessary to further accelerate particles well beyond the “knee” energy at $`10^{15}eV`$, as is suggested by observations of the CR background spectrum and the wide consensus on their SNR shock origin. The fact that accelerated CRs constitute an ample reservoir for the turbulence which is required to further accelerate them provides a logical basis for our approach. Indeed, the scenario proposed here may be viewed as a self-regulating enhanced acceleration process, which ultimately forces the energetic particle pressure gradient to its marginal point for Drury instability.
A new description of the instability of Alfven wave spectrum to acoustic modulations is given. Along with the Drury instability this is shown to provide an efficient mechanism for transformation of magnetic energy to longer scales. It should be also emphasized that the theoretical analysis of CR acceleration is a challenging problem in plasma wave physics, particle kinetics and shock hydrodynamics. It should, and indeed must, include a self-consistent description of particle transport, wave dynamics, and shock structure.
## Acknowledgements
This work was supported by the U.S. DOE under Grant No. FG03-88ER53275. |
warning/0506/quant-ph0506222.html | ar5iv | text | # Classicality in discrete Wigner functions
## I Introduction
Continuous-variable Wigner functions $`W(q,p)`$ have been used for a long time to represent quantum systems in phase space Wigner32 ; HilleryOSW84 . The Wigner function $`W(q,p)`$ is an alternative complete description of quantum states which behaves almost like a phase-space probability density. Not only is it real-valued and normalized but it also yields the correct value of the probability density for the quadrature $`a\widehat{Q}+b\widehat{P}`$ when integrated along the phase-space line $`aq+bp`$. However, unlike probability densities, the Wigner function can assume negative values for some quantum states. This negativity of the Wigner function has been considered a defining signature of non–classicality (or quantum coherence and interference) Zurek03 ; PazZ00 .
In quantum information science we usually deal with systems with a space of states with a finite dimension $`d`$. For example, for a system of $`n`$ qubits the dimension of the (Hilbert) space of states is $`d=2^n`$. For such systems, various discrete analogues of the Wigner function have been proposed Buot74 ; HannayB80 ; CohenS86 ; Feynman87 ; Wootters87 ; GalettiP88 ; GibbonsHW04 and used to investigate a variety of interesting problems connected with quantum computation such as the phase-space representation of quantum algorithms MiquelPS02 , separability PittengerR05 , quantum state tomography Wootters04 ; PazRS04 , teleportation KoniorczykBJ01 ; Paz02 , decoherence in quantum walks LopezP03 , and error correction PazRS04b . Here we shall concentrate on a class of discrete Wigner functions $`W`$ introduced recently by Gibbons et al. GibbonsHW04 . This elegant approach seems to be a potentially powerful tool to establish connections between phase-space techniques and problems in quantum information and foundations of quantum mechanics.
In this paper we study the set of states with non-negative discrete Wigner functions $`W`$ for all functions in the class proposed by GibbonsHW04 , and the group of unitaries that preserve non-negativity of $`W`$. Our first result is a complete characterization of the set of quantum states having non-negative discrete Wigner functions $`W`$. This is done by proving a conjecture presented by one of us in Galvao05 (a related discussion in a somewhat different context, using concepts in high-dimensional geometry appeared in Bengtsson04 ; BengtssonE04 ). Our proof is elementary and constructive, and shows that the only pure states with non-negative $`W`$ are stabilizer states, i.e. simultaneous eigenstates of generalized Pauli operators Gottesman97 . We then study the group of unitaries which preserve non-negativity of $`W`$, and prove that they form a subgroup of the Clifford group. This means such states and unitaries are classical in the sense of allowing for an efficient classical simulation scheme using the stabilizer formalism.
The paper is organized as follows. In Section II we review the discrete Wigner functions $`W`$ of GibbonsHW04 . In Section III we characterize the states with non-negative $`W`$, in Section IV we discuss positivity-preserving unitary dynamics in phase space and in Section V we summarize our results.
## II Discrete Wigner functions
In this section we review the class of discrete Wigner functions proposed in GibbonsHW04 and discuss some of their features.
Let us assume that we are describing a quantum state whose Hilbert space dimensionality $`d`$ is a power of a prime number $`p`$ ($`d=p^n`$). In such cases one can introduce a phase-space grid with $`d\times d`$ points and label the position and momentum coordinates $`(q,p)`$ with elements of the finite Galois field $`GF(p^n)`$ LidlN86 . At first the use of elements of $`GF(p^n)`$ for both phase-space coordinates could be seen as an unnecessary complication, but it turns out to be an essential step. The reason is that by doing this we can endow the phase-space grid with the same geometric properties as the ordinary plane. For example, in the finite $`d\times d`$ grid we can define lines as solutions to linear equations of the form $`aq+bp=c`$ \[where all elements and operations in this equation are in $`GF(p^n)`$\]. Each line will then consist of exactly $`d`$ points of the grid. The field structure of $`GF(p^n)`$ ensures the validity of properties such as: (i) there is only one line joining any given two points, (ii) two lines are either parallel (i.e. with no points in common) or they intersect at a single point. Moreover, it is possible to show that a set of $`d`$ parallel lines (which we will call a striationGibbonsHW04 ) is obtained by varying the parameter $`c`$ in the equation $`aq+bp=c`$. Finally, the number of different striations turns out to be $`(d+1)`$. The complete set of $`(d+1)`$ striations has been studied for a long time in discrete geometry, where it is called a finite affine plane Wootters04b ; LidlN86 . We will label the striations with an index $`\kappa =1,\mathrm{},d+1`$ and the lines within a striation with an index $`j=1,\mathrm{},d`$. In this way the $`j`$-th line belonging to the $`\kappa `$–th striation will be denoted as $`\lambda _j^{(\kappa )}`$.
A discrete phase space with the above properties was used by Gibbons, Hoffman and Wootters in GibbonsHW04 to define a class of discrete Wigner functions. As mentioned above, the crucial property of the continuous Wigner function is that its integral along any line $`\lambda `$ is equal to the expectation value of a projection operator $`\widehat{P}_\lambda `$, i.e. a probability. This essential feature is generalized to the discrete case in a straightforward way: every line in the $`d\times d`$ phase-space grid is associated to a rank one projection operator. As noted in GibbonsHW04 , this association cannot be arbitrary and must obey some simple geometric constraints. For example, we can define a set of $`d\times d`$ unitary operators $`\widehat{T}(q,p)`$ acting on the Hilbert space that faithfully represent discrete phase-space translations. For the association between lines and states to respect covariance under translations we must impose that the quantum state associated to a translated line should be identical to the state obtained by acting with the operator $`\widehat{T}(q,p)`$ on the original state. This covariance constraint can be used to show the validity of some very significant properties: a) the states associated to parallel lines must be orthogonal; b) the overlap between states associated to non-parallel lines must be equal to $`1/d`$. This is important and implies that the $`(d+1)`$ phase-space striations must be associated to an equal number of mutually unbiased bases (MUB), i.e. bases
$$MUB^{(\kappa )}=\{|\varphi _1^{(\kappa )},\mathrm{},|\varphi _d^{(\kappa )}\}$$
(1)
such that
$$|\varphi _j^{}^{(\kappa ^{})}|\varphi _j^{(\kappa )}|^2=\frac{1}{d}(1\delta _{\kappa ,\kappa ^{}})+\delta _{\kappa ,\kappa ^{}}\delta _{j,j^{}}.$$
(2)
As we see, mutually unbiased bases are orthonormal bases picked in such a way that any state in one basis is an equal–amplitude superposition of all the states of any other basis. A complete set of $`(d+1)`$ MUB is known to exist if the dimensionality of the space of states is a power of a prime number. In such case, many constructions of MUB have been proposed Ivanovic81 ; WoottersF89 ; LawrenceBZ02 ; BandyopadhyayBRV02 ; PittengerR04 . It has been shown that a complete set of $`(d+1)`$ MUB for $`d`$-dimensional systems can be chosen to consist solely of stabilizer states, i.e. simultaneous eigenstates of sets of (generalized) Pauli operators LawrenceBZ02 ; BandyopadhyayBRV02 ; PittengerR04 .
The defining feature of the discrete Wigner functions of GibbonsHW04 is the association between MUB and striations in the discrete phase space. As discussed in GibbonsHW04 this can be done in a variety of ways and each defines a different quantum net GibbonsHW04 , which will result in a different definition of the discrete Wigner function $`W`$. In this paper we propose a notion of classicality of quantum states which is based on non-negativity of $`W`$ for all quantum nets obtainable from a fixed complete set of MUB. It should be noted, however, that there have been proposals of criteria to narrow down the choice of quantum nets: in PazRS04b the criterion is covariance under the so-called discrete squeezing operator; and in PittengerR05 the net is chosen so as to enforce a natural relation between a separable state’s $`W`$ and the $`W`$ of its subsystems.
The quantum net is defined by associating each line $`\lambda _j^{(\kappa )}`$ in striation $`\kappa `$ to a projector $`\widehat{P}_j^{(\kappa )}=|\varphi _j^{(\kappa )}\varphi _j^{(\kappa )}|`$ onto a basis state of basis $`\kappa `$. Having fixed a quantum net, the discrete Wigner function is uniquely defined by imposing the condition that the sum of its values along any line should be equal to the expectation value of the projector corresponding to that line (see GibbonsHW04 for details). The resulting Wigner function at any phase-space point $`\alpha =(q,p)`$ can then be shown to be
$`W_\alpha `$ $`=`$ $`\mathrm{Tr}\left(\widehat{\rho }\widehat{A}(\alpha )\right),`$ (3)
$`\widehat{A}(\alpha )`$ $`=`$ $`{\displaystyle \frac{1}{d}}\left({\displaystyle \underset{\lambda _j^{(\kappa )}\alpha }{}}\widehat{P}_j^{(\kappa )}1𝐥\right),`$ (4)
where the sum is over projectors associated with all lines $`\lambda _j^{(\kappa )}`$ containing point $`\alpha `$. The construction of the striations guarantees that the sum above will contain exactly one projector from each basis.
The operators $`\widehat{A}(\alpha )`$ are known as phase-space point operators and form a complete basis for the space of operators, which is orthogonal in the Schmidt inner product (i.e. $`\mathrm{Tr}\left(\widehat{A}(\alpha )\widehat{A}(\beta )\right)=\delta _{\alpha ,\beta }/d`$). We can rewrite the expression for the Wigner function at phase-space point $`\alpha `$ using the probabilities associated with the projectors $`\widehat{P}_j^{(\kappa )}`$:
$$p_j^{(\kappa )}\mathrm{Tr}\left(\widehat{\rho }\widehat{P}_j^{(\kappa )}\right).$$
(5)
In terms of these probabilities, the Wigner function at the point $`\alpha `$ takes the form
$$W_\alpha =\frac{1}{d}\left(\underset{\lambda _j^{(\kappa )}\alpha }{}p_j^{(\kappa )}1\right).$$
(6)
The discrete Wigner function $`W`$ can be shown to have many of the features of the continuous Wigner function $`W(q,p)`$ GibbonsHW04 : it is real (but can be negative), normalized, and its values are obtained through eq. (6) from measurements onto MUB. Here the MUB projectors play the role that the quadratures $`a\widehat{Q}+b\widehat{P}`$ play in $`W(q,p)`$, forming a particularly symmetric set of observables whose measurement results completely characterize the state (in a process known as quantum tomography). For a discussion of further properties of $`W`$ see GibbonsHW04 ; PazRS04b ; PittengerR04 .
In the discussion that follows we will often be representing quantum states using the probabilities $`p_j^{(\kappa )}`$. As the projectors $`\widehat{P}_j^{(\kappa )}`$ form an over–complete basis for the space of density matrices, these probabilities completely characterize the state. Since for any striation $`_jp_j^{(\kappa )}=1`$, there are only $`(d1)`$ independent probabilities for each basis, resulting in a total of $`(d1)(d+1)=(d^21)`$ independent probabilities, exactly the number of real parameters necessary to describe a general normalized mixed quantum state in $`d`$-dimensional Hilbert space. Each quantum state is represented by a point $`\stackrel{}{p}`$ in this $`(d^21)`$-dimensional probability space.
As mentioned above, for power-of-prime $`d`$ it is possible to build a complete set of $`(d+1)`$ MUB using only stabilizer states, i.e. joint eigenstates of generalized Pauli operators. Let us discuss more explicitly such constructions for the case $`d=2^n`$, i.e. $`n`$ qubits (see LawrenceBZ02 for more details). In order to define a complete set of $`(2^n+1)`$ MUB we start by partitioning the $`(4^n1)`$ Pauli operators (excluding the identity) into $`(2^n+1)`$ sets $`S_i`$ of $`(2^n1)`$ Pauli operators each. We will require that the Pauli operators in each set $`S_i`$ be mutually commuting, but otherwise the partitioning can be completely arbitrary. If we add the identity and a $`\pm 1`$ phase to the Pauli operators in each set $`S_i`$, each will form a maximal Abelian subgroup of the Pauli group. The joint eigenstates of each such set $`S_i`$ form a basis for the Hilbert space, and due to properties of the Pauli operators the $`(2^n+1)`$ bases thus defined can be shown to be mutually unbiased LawrenceBZ02 .
The phase-space construction provides a natural procedure for partitioning the Pauli group into disjoint, mutually commuting sets. The idea, which is worth reviewing here, was described in GibbonsHW04 and further elaborated in PazRS04b . Pauli operators represent phase-space translations and can be labelled using binary $`n`$–tuples $`\stackrel{}{p}`$ and $`\stackrel{}{q}`$ ($`n`$–tuples $`\stackrel{}{q}`$ and $`\stackrel{}{p}`$ contain the coordinates of the field elements $`q`$ and $`p`$ in a given basis as described below). Each Pauli operator can be written as
$$\widehat{T}(\stackrel{}{q},\stackrel{}{p})=\underset{i=0}{\overset{n1}{}}\widehat{X}_i^{q_i}\widehat{Z}_i^{p_i}e^{i\frac{\pi }{2}q_ip_i},$$
(7)
where $`\widehat{X}_i`$ and $`\widehat{Z}_i`$ stand for the Pauli operators on qubit $`i`$, and the phase is chosen so as to make the operators Hermitian. The definition above will be written in shorthand as
$$\widehat{T}(\stackrel{}{q},\stackrel{}{p})=\widehat{X}^\stackrel{}{q}\widehat{Z}^\stackrel{}{p}e^{i\frac{\pi }{2}\stackrel{}{q}\stackrel{}{p}}.$$
(8)
The condition for two Pauli operators to commute turns out to be
$$[\widehat{T}(\stackrel{}{q},\stackrel{}{p}),\widehat{T}(\stackrel{}{q^{}},\stackrel{}{p^{}})]=0\text{ iff }\stackrel{}{q}\stackrel{}{p^{}}\stackrel{}{p}\stackrel{}{q^{}}=0(\mathrm{mod}2).$$
(9)
Let us consider a set of $`(d1)`$ Pauli operators
$$S_{(\stackrel{}{a},\stackrel{}{b})}=\{\widehat{T}(\stackrel{}{a}M^j,\stackrel{}{b}\stackrel{~}{M}^j)j=0,1,\mathrm{},d2\},$$
(10)
where $`M`$ is an arbitrary binary matrix, $`\stackrel{~}{M}`$ is its transpose and $`\stackrel{}{a}`$, $`\stackrel{}{b}`$ are binary $`n`$–tuples. Any two operators of this set commute. It is interesting to note that $`S_{(\stackrel{}{a},\stackrel{}{b})}`$ forms a maximal Abelian subgroup of the Pauli group if and only if $`M`$ is a generating element of the matrix representation of the field $`GF(2^n)`$. This can be seen as follows: the product of two elements of $`S_{(\stackrel{}{a},\stackrel{}{b})}`$ is itself an element of this set (up to a sign) iff the matrix $`M`$ is such that for every pair of integers $`j,j^{}`$, there is a third integer $`j^{\prime \prime }`$ such that $`M^j+M^j^{}=M^{j^{\prime \prime }}`$ \[where $`0j,j^{},j^{\prime \prime }(d2)`$ and $`jj^{}`$\]. Moreover, for the set to have exactly $`(d1)`$ different elements, the matrix $`M`$ should be such that all powers $`M^j`$ for $`j=0,\mathrm{},(d2)`$ are nonzero and different from each other. For $`M`$ satisfying these conditions, it can be seen that $`M^{d1}=1𝐥`$. Therefore, the elements of the set $`\{\mathrm{𝟎},1𝐥,M,M^2,\mathrm{},M^{d2}\}`$ form a finite field, and we see that the matrix $`M`$ and its powers form a matrix representation of $`GF(2^n)`$. A possible choice for $`M`$, used in PazRS04b , is the so–called “companion matrix” of the primitive polynomial which defines the product rule in the field. With such a matrix we can build $`(d+1)`$ disjoint sets of commuting Paulis of the form $`S_{(\stackrel{}{a},\stackrel{}{b})}`$ by choosing the binary $`n`$–tuples $`(\stackrel{}{a},\stackrel{}{b})`$ as explained below.
The association between each phase-space point $`(q,p)`$ and a Pauli operator $`\widehat{T}(\stackrel{}{q},\stackrel{}{p})`$ must respect the covariance of the construction under phase space translations. This is done as follows: the line formed by all phase-space points satisfying the equation $`bq+ap=c`$ is invariant under phase-space translations of the form $`q^{}=q+a\omega ^j`$, $`p^{}=p+b\omega ^j`$ (where $`\omega `$ is a generating element of the field). To this phase-space translation we must associate an operator acting in Hilbert space. The natural identification is to associate this with the operator $`\widehat{T}(\stackrel{}{a}M^j,\stackrel{}{b}\stackrel{~}{M}^j)`$. Here, the choice of $`n`$–tuples $`\stackrel{}{a}`$ and $`\stackrel{}{b}`$ is arbitrary. The important point is that once this choice is made \[i.e., once we arbitrarily assign two $`n`$–tuples to the point $`(a,b)`$\] we repeatedly apply the matrix $`M`$ ($`\stackrel{~}{M}`$) to the position (momentum) coordinates to obtain the $`n`$–tuples parametrizing the Pauli operators associated to the other phase-space points PazRS04b . In summary, this construction associates an operator $`\widehat{T}(\stackrel{}{q},\stackrel{}{p})`$ to every phase space point $`(q,p)`$ in such a way that the elements of the Abelian subgroup $`S_{(\stackrel{}{a},\stackrel{}{b})}`$ are associated to points in phase space that belong to the ray defined by the equation $`bq+ap=0`$ (a ray is defined as a line that contains the origin). In PazRS04b it was shown that by varying the $`n`$–tuples $`(\stackrel{}{a},\stackrel{}{b})`$ one can construct only $`(d+1)`$ different sets $`S_{(\stackrel{}{a},\stackrel{}{b})}`$. If we define the two $`n`$–tuples $`\stackrel{}{1}(1000\mathrm{}0)`$ and $`\stackrel{}{0}(0000\mathrm{}0)`$, these maximal mutually commuting sets of Pauli operators can be conveniently built by choosing $`(\stackrel{}{a},\stackrel{}{b})=(\stackrel{}{1},\stackrel{}{0)}`$ (which we will associate with the horizontal striation), $`(\stackrel{}{a},\stackrel{}{b})=(\stackrel{}{0},\stackrel{}{1)}`$ (the vertical striation) and $`(\stackrel{}{a},\stackrel{}{b})=(\stackrel{}{1},\stackrel{}{b)}`$ for $`\stackrel{}{b}\stackrel{}{0}`$ (the other striations).
Thus, the mapping between striations and MUB is naturally determined by the phase-space construction, as lines which are invariant under the transformations $`q^{}=q+a\omega ^j`$, $`p^{}=p+b\omega ^j`$ must be associated to states which are invariant under the corresponding transformations in Hilbert space, that is, the translation operators in $`S_{(\stackrel{}{a},\stackrel{}{b})}`$. Therefore, the lines of the form $`bq+ap=c`$ must be associated to common eigenstates of the set $`S_{(\stackrel{}{a},\stackrel{}{b})}`$. However, there is no criterion telling us how to associate each line in a striation with a projector in the corresponding basis. We can count the number of possible quantum nets as follows: for the ray of a given striation there are $`d`$ possible projectors to choose from; once this choice has been made the condition of covariance under translations determines which projector should be associated to each of the other lines in the same striation. As there are $`(d+1)`$ rays, the number of possible associations between lines and projectors is $`d^{d+1}`$, each of which defines a different quantum net, leading to a different definition of the Wigner function. The projectors associated to each of the lines in the vertical and horizontal striations can be chosen in such a way that the coordinates of each line correspond to the eigenvalues of the Paulis generating the set (the single qubit Paulis $`\widehat{Z}`$ and $`\widehat{X}`$, respectively). Then, there are still $`d^{d1}`$ possible quantum nets, each of them given by a particular choice of projectors to be associated to the rays of the remaining oblique striations.
There is a closely related methodology for constructing the Wigner functions for general dimension $`d=p^n`$ which emphasizes the link between the exponents of the $`\widehat{Z}`$ and $`\widehat{X}`$ operators and the finite geometry of $`V_2[GF(p^n)]`$, the two-dimensional vector space over the field $`GF(p^n)`$. One defines lines, rays and striations in this two dimensional space and then, using the properties of the algebraic field extension, defines an isomorphism with $`V_{2n}[GF(p)]`$. Vectors in this second space serve as exponents of the $`\widehat{Z}`$ and $`\widehat{X}`$ operators, and the commuting classes of generalized Pauli matrices correspond precisely to parallel lines in a striation in the first vector space. Details of this approach and a methodology for assigning projections to lines are given in PittengerR04 .
## III States with non-negative Wigner functions $`W`$
Following Galvao05 , let us now characterize the set of states having non-negative discrete Wigner functions $`W`$ simultaneously in all definitions proposed by Gibbons et al. GibbonsHW04 for power-of-prime dimension $`d`$.
The set $`C_d`$ is defined as the set of (pure or mixed) density matrices of systems in a $`d`$-dimensional Hilbert space having non-negative discrete Wigner function $`W`$ in all phase-space points and for all definitions of $`W`$ using a fixed set of mutually unbiased bases.
By definition, the set $`C_d`$ is specified as the intersection of a number of half-spaces in the $`(d^21)`$-dimensional $`\stackrel{}{p}`$-space. From (6) it can be seen that each half-space inequality is of the form
$$\underset{\lambda _j^{(\kappa )}\alpha }{}p_j^{(\kappa )}1,$$
(11)
where the probabilities appearing in the sum are associated with the lines containing phase-space point $`\alpha `$, and hence depend on $`\alpha `$ and on the quantum net chosen. The intersection of the half-spaces defined by these inequalities is a convex polytope in $`\stackrel{}{p}`$-space, given in an H-description (H standing for “Half-space”). Any convex polytope also admits an alternative V-description (V for “Vertices”), consisting of the list of vertices whose convex hull defines the polytope.
Galvão showed that for $`d5`$, the H-polytope $`C_d`$ has a V-description whose vertices are the MUB projectors Galvao05 , and conjectured this would also be true for general power-of-prime $`d`$. A geometrical argument showing the validity of this conjecture was given in Bengtsson04 ; BengtssonE04 . Let us now provide a constructive, analytical proof.
###### Theorem 1
For any power-of-prime Hilbert space dimension $`d`$, the H-polytope $`C_d`$ is equivalent to the V-polytope $`C_v`$ having the MUB projectors as vertices.
Proof:
Let us prove the theorem by first showing that the V-polytope $`C_v`$ is contained in $`C_d`$, and then the converse. From GibbonsHW04 we know that the Wigner function for any MUB projector $`\widehat{P}_j^{(\kappa )}=|\varphi _j^{(\kappa )}\varphi _j^{(\kappa )}|`$ is non-negative. Since the Wigner function depends linearly on the density matrices, $`W`$ is non-negative also for any state in the convex hull of the $`\widehat{P}_j^{(\kappa )}`$. This shows that any state in $`C_v`$ is also in $`C_d`$, as we wanted to prove.
Let us now prove the converse, i.e. that polytope $`C_v`$ is contained in polytope $`C_d`$. What is required now is to show that any state $`\widehat{\rho }C_d`$ can be written as a convex combination of the projectors $`\widehat{P}_j^{(\kappa )}`$:
$$\widehat{\rho }=\underset{\kappa =1}{\overset{d+1}{}}\underset{j=1}{\overset{d}{}}c_j^{(\kappa )}\widehat{P}_j^{(\kappa )},$$
(12)
with all $`c_j^{(\kappa )}0`$. Note that the decomposition is not unique.
Let us start by considering the general expression for Wigner function $`W`$ at phase-space point $`\alpha `$ \[eq. (6)\]. Given a state $`\widehat{\rho }`$ there is some Wigner function definition which, at some point $`\alpha `$, evaluates to a minimum value among all definitions and all points $`\alpha `$. This happens when the expression for $`W_\alpha `$ is such that the sum (6) includes only the smallest probability $`p_j^{(\kappa )}`$ from each MUB $`\kappa `$. Let us denote these $`W`$-minimizing probabilities $`p_{}^{(\kappa )}\mathrm{min}_j\{p_j^{(\kappa )}\}`$. States in $`C_d`$ are those for which all expressions of the form (6) are non-negative, and this happens if and only if the expression for $`W`$ involving only the $`p_{}^{(\kappa )}`$ is non-negative. In other words, a state has non-negative $`W`$ in all definitions if and only if:
$$\underset{\kappa =1}{\overset{d+1}{}}p_{}^{(\kappa )}1.$$
(13)
This is our hypothesis.
Any density matrix $`\widehat{\rho }`$ can be expanded in terms of the projectors $`\widehat{P}_j^{(\kappa )}=|\varphi _j^{(\kappa )}\varphi _j^{(\kappa )}|`$ as in eq. (12), with real (but possibly negative) coefficients $`c_j^{(\kappa )}`$. A first constraint on the coefficients $`c_j^{(\kappa )}`$ comes from the requirement that $`\mathrm{Tr}(\widehat{\rho })=1`$. Using property (2) of the MUB we can compute the trace, obtaining
$$\mathrm{Tr}(\widehat{\rho })=\underset{\kappa =1}{\overset{d+1}{}}\underset{j=1}{\overset{d}{}}c_j^{(\kappa )}=1.$$
(14)
Now let us use eq. (5) to calculate $`p_j^{(\kappa )}`$ explicitly from eq. (12), so as to obtain relations between the coefficients $`c_j^{(\kappa )}`$ and the probabilities $`p_j^{(\kappa )}`$:
$`p_j^{(\kappa )}=\varphi _j^{(\kappa )}|\widehat{\rho }|\varphi _j^{(\kappa )}={\displaystyle \underset{\mu =1}{\overset{d+1}{}}}{\displaystyle \underset{m=1}{\overset{d}{}}}c_m^{(\mu )}\left|\varphi _j^{(\kappa )}|\varphi _m^{(\mu )}\right|^2=`$
$`={\displaystyle \underset{\mu \kappa }{}}{\displaystyle \underset{m}{}}c_m^{(\mu )}\left|\varphi _j^{(\kappa )}|\varphi _m^{(\mu )}\right|^2+{\displaystyle \underset{m}{}}c_m^{(\kappa )}\left|\varphi _j^{(\kappa )}|\varphi _m^{(\kappa )}\right|^2=`$
$`={\displaystyle \underset{\mu \kappa }{}}{\displaystyle \underset{m}{}}c_m^{(\mu )}{\displaystyle \frac{1}{d}}+c_j^{(\kappa )},`$ (15)
where we have used the condition of mutual unbiasedness of the bases \[eq. (2)\]. Now we can use the trace condition (14) to rewrite this as
$$p_j^{(\kappa )}=c_j^{(\kappa )}+\frac{1}{d}\frac{1}{d}\underset{m}{}c_m^{(\kappa )}$$
(16)
or
$$c_j^{(\kappa )}=p_j^{(\kappa )}\frac{1}{d}+\frac{1}{d}\underset{m}{}c_m^{(\kappa )}.$$
(17)
Let us add $`0=p_{}^{(\kappa )}p_{}^{(\kappa )}`$ to the right-hand side of the equation above, to obtain
$$c_j^{(\kappa )}=\left(p_j^{(\kappa )}p_{}^{(\kappa )}\right)+x^{(\kappa )}$$
(18)
with
$$x^{(\kappa )}p_{}^{(\kappa )}\frac{1}{d}+\frac{1}{d}\underset{m}{}c_m^{(\kappa )}.$$
(19)
Eq. (18) tells us that each coefficient $`c_j^{(\kappa )}`$ can be written as the sum of a non-negative term $`\left(p_j^{(\kappa )}p_{}^{(\kappa )}\right)`$ plus a (possibly negative) constant $`x^{(\kappa )}`$. We can show, however, that the sum of those constants $`x^{(\kappa )}`$ has to be non-negative. We do that by using the normalization condition (14) on eq. (18):
$`{\displaystyle \underset{\kappa ,j}{}}c_j^{(\kappa )}=1{\displaystyle \underset{\kappa ,j}{}}p_j^{(\kappa )}{\displaystyle \underset{\kappa ,j}{}}p_{}^{(\kappa )}+{\displaystyle \underset{\kappa ,j}{}}x^{(\kappa )}=1`$
$`d+1d{\displaystyle \underset{\kappa }{}}p_{}^{(\kappa )}+d{\displaystyle \underset{\kappa }{}}x^{(\kappa )}=1`$
$`{\displaystyle \underset{\kappa }{}}x^{(\kappa )}={\displaystyle \underset{\kappa }{}}p_{}^{(\kappa )}1.`$ (20)
Now remember that our hypothesis is that $`_\kappa p_{}^{(\kappa )}1`$, which implies that $`_\kappa x^{(\kappa )}0`$. Let us now use this fact and expression (18) to obtain an expansion of the density matrix $`\widehat{\rho }`$ in terms of the projection operators $`\widehat{P}_j^{(\kappa )}`$, but now with non-negative coefficients only. Plugging eq. (18) into eq. (12) we obtain:
$$\widehat{\rho }=\underset{\kappa ,j}{}(p_j^{(\kappa )}p_{}^{(\kappa )}+x^{(\kappa )})\widehat{P}_j^{(\kappa )}.$$
(21)
Using the fact that $`_j\widehat{P}_j^{(\kappa )}=1𝐥`$, we can rewrite this as
$$\widehat{\rho }=\underset{\kappa ,j}{}\left(p_j^{(\kappa )}p_{}^{(\kappa )}+\frac{x}{d+1}\right)\widehat{P}_j^{(\kappa )}$$
(22)
where we defined $`x_\kappa x^{(\kappa )}`$. Note that $`p_j^{(\kappa )}p_{}^{(\kappa )}`$ by definition, and our hypothesis guarantees that $`x0`$. What we have now is then an expansion of $`\widehat{\rho }`$ in terms of the MUB projectors using only non-negative coefficients. QED
Our proof above is constructive – for any state in $`C_d`$ we can use equations (5) and (20) to obtain a convex decomposition of the state in terms of the MUB projectors and their associated probabilities, given by eq. (22). As noted in Galvao05 , some non-physical states (i.e. described by non-positive Hermitian matrices) can have non-negative Wigner functions in a single definition of $`W`$. Theorem 1 shows that imposing non-negativity of $`W`$ for all definitions of $`W`$ is sufficient to guarantee that the set $`C_d`$ contains only physical states.
In the light of Theorem 1 above, let us now discuss in which senses states with non-negative Wigner functions are classical. We have defined the set $`C_d`$ of states of a $`d`$-dimensional system with non-negative Wigner function $`W`$ in all definitions. Theorem 1 proves that the only pure states in $`C_d`$ are the MUB projectors, which can always be chosen to be stabilizer states, i.e. simultaneous eigenstates of Pauli operators LawrenceBZ02 ; BandyopadhyayBRV02 ; PittengerR04 . The stabilizer formalism then provides us with a way to represent pure states in $`C_d`$ using a number of bits which is polynomial in the number of qubits NielsenC00 . This contrasts with general quantum states whose classical description requires an exponential number of bits. We thus see that states whose discrete phase-space description avoids Feynman’s “negative probabilities” Feynman87 are classical also in the sense of having a classical-like short description.
Classicality witnesses are observables whose expectation values, if negative, indicate some non-classical property such as entanglement (see Terhal02 ; KorbiczCWL05 ). Our Theorem 1 gives such an interpretation to the phase-space point operators $`\widehat{A}(\alpha )`$: negative expectation values indicate (i.e. witness) non-classicality in the sense we discussed above.
The stabilizer formalism provides us with a framework in terms of which pure states in the set $`C_d`$ have an efficient classical description. Other choices of frameworks are possible, each choice resulting in a different set of quantum states with efficient classical descriptions. One example are the (mixed) separable density matrices of a collection of qubits, each of which has an efficient classical description in terms of single-qubit pure states. An efficient description, however, does not guarantee the existence of an efficient simulation scheme for the dynamics; the dynamics of separable mixed states in NMR quantum computation experiments provides us with an example of this problem MenicucciC02 . In the next section we discuss the issue of simulability of unitary dynamics within our set $`C_d`$.
Given these observations, it is not surprising that pure states in $`C_d`$ can behave non-classically in other ways, that is, with respect to other frameworks. For example, states in $`C_d`$ can be highly entangled, allowing for proofs of quantum non-locality and contextuality.
## IV Unitaries preserving non-negativity of $`W`$
In continuous phase-space we can define classical unitaries as the group of unitaries which preserve non-negativity of the Wigner function $`W(q,p)`$. It has been shown that this group is formed by all unitaries generated by Hamiltonians which are quadratic forms in phase space HilleryOSW84 . In this section we obtain an analogous result for our discrete Wigner functions $`W`$: using the ‘classical’ pure states in $`C_d`$, we define and characterize the group of unitaries $`\{U_c\}`$ that map pure states in $`C_d`$ to other pure states in $`C_d`$. In other words, we characterize the group of ‘classical’ unitaries $`\{U_c\}`$ that preserve non-negativity of $`W`$ for all quantum nets obtained from a fixed complete set of MUB.
The structure of the group $`\{U_c\}`$ may depend on the particular complete set of MUB we choose to define our discrete Wigner functions. We prove that for any MUB construction using Pauli operators, the group $`\{U_c\}`$ is a subgroup of the Clifford group (the group of unitaries mapping Pauli operators to Pauli operators under conjugation NielsenC00 ). For the particular construction in GibbonsHW04 ; PazRS04b we present some unitaries in $`\{U_c\}`$ and discuss their action in phase space.
### IV.1 $`\{U_c\}`$ is a subgroup of the Clifford group
Let us consider the $`(4^n1)`$ Pauli operators acting on the $`2^n`$-dimensional Hilbert space of $`n`$ qubits (excluding the identity). We can partition these $`(4^n1)`$ Pauli operators into $`(2^n+1)`$ sets $`S_i`$ of $`(2^n1)`$ commuting Paulis each. The joint eigenstates of the sets $`S_i`$ form a complete set of $`(2^n+1)`$ MUB, as discussed in section II. We can partition the Paulis in many different ways; each such partitioning defines a different complete set of MUB, which will be denoted as $`B_j`$. In this section we show that the ‘classical’ unitaries $`\{U_c\}`$ mapping MUB in a partition $`B_i`$ to MUB in the same partition form a subgroup of the Clifford group.
The strategy is as follows. We will consider a slightly more general problem, which is to characterize unitaries mapping MUB defined by an arbitrary partition $`B_1`$ of Pauli operators to MUB defined by a second partition $`B_2`$. A general unitary $`U`$ mapping $`B_1`$ to $`B_2`$ will in particular map two bases in $`B_1`$ to two other bases in $`B_2`$. Let us name them $`(S_1B_1)\stackrel{U}{}(S_2B_2)`$, $`(T_1B_1)\stackrel{U}{}(T_2B_2)`$. The first step is to prove there are two Clifford unitaries $`C_j(j=1,2)`$ that map basis $`S_j`$ to the computational ($`Z`$) basis, while mapping basis $`T_j`$ to the $`X`$-basis. These standard Clifford unitaries are the key to the proof. This is because $`U`$ is Clifford if and only if $`\stackrel{~}{U}C_2UC_1^{}`$ is Clifford. So it is enough to show $`\stackrel{~}{U}`$ is Clifford (done in Theorem 3), which is easier as by construction $`\stackrel{~}{U}`$ are unitaries that preserve both the $`Z`$ basis and the $`X`$ basis.
With this more general result in hand, we can consider the case when the two partitions are one and the same ($`B_1=B_2`$), and we will have what we wanted to prove, i.e. that our ‘classical’ unitaries $`\{U_c\}`$ are Clifford group operators.
We wish to show
###### Theorem 2
Let $`U`$ be a unitary transformation that maps a complete set of Pauli MUB $`B_1`$ to a second complete set of Pauli MUB $`B_2`$. Then $`U`$ is in the Clifford group, up to a global phase.
The first step involves proving the following Lemma:
###### Lemma 1
Let $`S`$ and $`T`$ be two maximal Abelian subgroups of the Pauli group, with $`ST=\left\{1𝐥\right\}`$. Then there exists a Clifford operation which maps $`SZ_1,\mathrm{},Z_n`$, $`TX_1,\mathrm{},X_n`$.
Proof:
Since $`S`$ and $`T`$ are maximal Abelian subgroups with trivial intersection, it follows that no (non-identity) element of $`T`$ commutes with every element of $`S`$ (or vice-versa). Let $`\{M_ii=1,\mathrm{},n\}`$ be a set of generators of $`S`$. For any particular element $`NT`$, we can define the syndrome $`\stackrel{}{\sigma }(N)`$ which is an $`n`$-tuple whose $`i`$-th component is given by $`\sigma _i=c(N,M_i)i=1,\mathrm{},n`$. Here $`c(N,M)=0`$ if $`N`$ and $`M`$ commute and $`c(N,M)=1`$ if $`N`$ and $`M`$ anticommute. Then it follows that if $`N,N^{}T`$, $`NN^{}`$, then $`\stackrel{}{\sigma }(N)\stackrel{}{\sigma }(N^{})`$ (since otherwise $`\stackrel{}{\sigma }(NN^{})=\stackrel{}{\sigma }(N)+\stackrel{}{\sigma }(N^{})=\stackrel{}{0}`$, and $`NN^{}`$ would commute with every element of $`S`$).
In particular, since there are $`2^n`$ elements of $`T`$ and $`2^n`$ different possible values of $`\stackrel{}{\sigma }`$, it follows that each value of $`\stackrel{}{\sigma }`$ is used exactly once. Thus, we can choose $`N_iT`$ such that $`\stackrel{}{\sigma }(N_i)=\stackrel{}{e_i}`$ (where $`\stackrel{}{e_i}`$ is the vector that is $`1`$ in the $`i`$-th position and $`0`$ elsewhere). That is, $`N_i`$ anticommutes with $`M_i`$ and commutes with $`M_j`$ ($`ij`$). The $`N_i`$’s are independent (because their $`\stackrel{}{\sigma }`$ vectors are independent) and they commute with each other (because $`T`$ is Abelian). Therefore, the set of $`M_i`$’s and $`N_i`$’s have the same commutation/anticommutation relationships as the $`Z_i`$’s and the $`X_i`$’s, so there exists a Clifford group operation that maps $`M_iZ_i`$ and $`N_iX_i`$. This provides the appropriate map on $`S`$ and $`T`$. QED
This Lemma can be adapted so it applies also to $`d`$-dimensional registers, the main difference being that the syndrome function $`\stackrel{}{\sigma }(N)`$ assumes values which are vectors modulo $`d`$ (see Gottesman99b ).
An immediate consequence of this lemma is that for any complete set of Pauli MUB $`B`$, we can choose any two of its bases, represented by stabilizers $`S`$ and $`T`$, and find a Clifford group operation that will map $`B`$ to another complete set of MUB containing the bases $`Z_1,\mathrm{},Z_n`$ and $`X_1,\mathrm{},X_n`$; and in particular, this Clifford group operation will map $`S`$ to $`Z_1,\mathrm{},Z_n`$ and $`T`$ to $`X_1,\mathrm{},X_n`$.
Therefore, if we have a general unitary $`U`$ that maps Pauli MUB $`B_1`$ to Pauli MUB $`B_2`$, we can choose bases $`S_1,T_1B_1`$ with $`S_2=U(S_1)`$, $`T_2=U(T_1)`$ (so $`S_2,T_2B_2`$), and then find Clifford operations $`C_1`$ and $`C_2`$ which map $`C_1:S_1Z_1,\mathrm{},Z_n`$, $`C_1:T_1X_1,\mathrm{},X_n`$, $`C_2:S_2Z_1,\mathrm{},Z_n`$, $`C_2:T_2X_1,\mathrm{},X_n`$. Then it follows that $`C_2UC_1^{}:Z_1,\mathrm{},Z_nZ_1,\mathrm{},Z_n`$, $`C_2UC_1^{}:X_1,\mathrm{},X_nX_1,\mathrm{},X_n`$. Denoting $`\stackrel{~}{U}=C_2UC_1^{}`$, it is easy to see that $`\stackrel{~}{U}`$ is a Clifford group operator iff $`U`$ is a Clifford group operation. Thus, to prove Theorem 2, it will be sufficient to prove
###### Theorem 3
If $`\stackrel{~}{U}`$ is a unitary operation which preserves both the $`Z`$ basis and the $`X`$ basis (i.e., maps eigenstates of $`Z_1,\mathrm{},Z_n`$ to other eigenstates of this set of operators, and the same is valid for eigenstates of $`X_1,\mathrm{},X_n`$), then $`\stackrel{~}{U}`$ is a Clifford group operation, up to a global phase.
Proof:
Since $`\stackrel{~}{U}`$ preserves the $`Z`$ basis, it has the form of a classical gate with possibly some phases changed:
$$\stackrel{~}{U}|\stackrel{}{z}_Z=e^{i\varphi (\stackrel{}{z})}|\stackrel{}{g}(\stackrel{}{z})_Z,$$
(23)
with $`\stackrel{}{g}(\stackrel{}{z})`$ a permutation of the $`2^n`$ possible values of $`\stackrel{}{z}`$.
In terms of the $`Z`$ basis, we can expand elements of the $`X`$ basis as follows:
$$|\stackrel{}{x}_X=\underset{\stackrel{}{z}}{}e^{i\pi (\stackrel{}{x}\stackrel{}{z})}|\stackrel{}{z}_Z.$$
(24)
Therefore,
$`\stackrel{~}{U}|\stackrel{}{x}_X`$ $`=`$ $`{\displaystyle \underset{\stackrel{}{z}}{}}e^{i(\pi \stackrel{}{x}\stackrel{}{z}+\varphi (\stackrel{}{z}))}|\stackrel{}{g}(\stackrel{}{z})_Z=`$ (25)
$`=`$ $`{\displaystyle \underset{\stackrel{}{z}}{}}e^{i(\pi \stackrel{}{x}\stackrel{}{g}^1(\stackrel{}{z})+\varphi (\stackrel{}{g}^1(\stackrel{}{z})))}|\stackrel{}{z}_Z.`$
In order to preserve the $`X`$ basis, we need
$$\stackrel{~}{U}|\stackrel{}{x}_X=e^{i\theta (\stackrel{}{x})}|\stackrel{}{h}(\stackrel{}{x})_X=\underset{\stackrel{}{z}}{}e^{i(\theta (\stackrel{}{x})+\pi \stackrel{}{h}(\stackrel{}{x})\stackrel{}{z})}|\stackrel{}{z}_Z,$$
(26)
where $`\stackrel{}{h}(\stackrel{}{x})`$ is a permutation of the values of $`\stackrel{}{x}`$. Equating (25) and (26), we find
$$\pi \stackrel{}{x}\stackrel{}{g}^1(\stackrel{}{z})+\varphi (\stackrel{}{g}^1(\stackrel{}{z}))=\theta (\stackrel{}{x})+\pi \stackrel{}{h}(\stackrel{}{x})\stackrel{}{z}.$$
(27)
This must be true for all $`\stackrel{}{x}`$ and $`\stackrel{}{z}`$. Plugging in $`\stackrel{}{x}=\stackrel{}{0}`$, we find
$$\varphi (\stackrel{}{g}^1(\stackrel{}{z}))=\theta _0+\pi \stackrel{}{h_0}\stackrel{}{z},$$
(28)
where $`\theta _0=\theta (\stackrel{}{0})`$ and $`\stackrel{}{h_0}=\stackrel{}{h}(\stackrel{}{0})`$. Therefore
$$\pi \stackrel{}{x}\stackrel{}{g}^1(\stackrel{}{z})=\pi \left[\stackrel{}{h}(\stackrel{}{x})\stackrel{}{h_0}\right]\stackrel{}{z}+\left[\theta (\stackrel{}{x})\theta _0\right].$$
(29)
Of course, eqs. 27-29 are understood to be modulo $`2\pi `$. It then follows that $`\stackrel{}{g}^1(\stackrel{}{z})`$ must be affine in $`\stackrel{}{z}`$:
$$\stackrel{}{g}^1(\stackrel{}{z})=A\stackrel{}{z}+\stackrel{}{b},$$
(30)
(where $`A`$ is an invertible $`n\times n`$ binary matrix) and thus, by (28), $`\varphi (\stackrel{}{z})`$ is also affine in $`\stackrel{}{z}`$:
$$\varphi (\stackrel{}{z})=\pi \stackrel{}{c}\stackrel{}{z}+d,$$
(31)
with $`A\stackrel{}{c}=\stackrel{}{h_0}`$ and $`\pi \stackrel{}{b}\stackrel{}{c}+d=\theta _0`$.
Thus we find
$$\stackrel{~}{U}|\stackrel{}{z}_Z=e^{i(\pi \stackrel{}{c}\stackrel{}{z}+d)}|A^1\stackrel{}{z}A^1\stackrel{}{b}.$$
(32)
We can easily identify this as a Clifford group operation, up to the global phase $`e^{id}`$: $`|\stackrel{}{z}|A^1\stackrel{}{z}`$ can be performed with CNOT gates, $`|\stackrel{}{z}|\stackrel{}{z}A^1\stackrel{}{b}`$ can be performed with $`X`$ operations, and $`|\stackrel{}{z}e^{i\pi (\stackrel{}{c}\stackrel{}{z})}|\stackrel{}{z}`$ can be performed with $`Z`$ operations. QED
The proof for $`d`$-dimensional registers is almost identical, except that we must replace $`\pi `$ everywhere with $`2\pi /d`$, and we need scalar multiplication gates as well as SUM gates to perform $`|\stackrel{}{z}|A^1\stackrel{}{z}`$ Gottesman99b .
For prime Hilbert space dimensions there is a unique Pauli MUB construction that uses all stabilizer states BandyopadhyayBRV02 . In that case our set of pure classical states coincides with the set of stabilizer states, and our group of classical unitaries $`\{U_c\}`$ coincides with the Clifford group. This is not the case for power-of-prime dimensions, where a complete set of Pauli MUB contains only a proper subset of the stabilizer states, resulting in classical unitaries $`\{U_c\}`$ which form a proper subgroup of the Clifford group.
The Gottesman-Knill theorem Gottesman97 states that Clifford group operations on stabilizer states can be simulated efficiently on a classical computer. Thus, our Theorems 1 and 2 guarantee that the group of classical unitaries $`\{U_c\}`$ applied on the set of classical pure states in $`C_d`$ can be efficiently simulated, i.e. with a number of time-steps that increases only polynomially in the number of qubits. This is to be contrasted with general quantum computation, which uses states and operations outside of our classical sets, and which is thought to provide exponential speedup for some problems. The necessity of negativity of $`W`$ for achieving universal quantum computation had been noted in Galvao05 for a particular computational model proposed recently by Bravyi and Kitaev BravyiK05 .
Our results point to an interesting convergence between two different notions of classicality. The first defines classical states as those whose description can be made in terms of non-negative quasi-probability distributions, in this case the discrete Wigner functions of GibbonsHW04 . The second is motivated by quantum computation: classical states and operations are those which can be efficiently simulated on a classical computer. For a related discussion of simulability in the context of continuous variables see BartlettSBN02 .
In this section and the previous one we only made claims of classicality for pure states in $`C_d`$ and their associated unitary dynamics. The problem seems to become much more involved when we consider mixed states in $`C_d`$ and their associated dynamics, which in this case will be (in general non-unitary) completely positive maps. It is not clear whether an efficient simulation scheme for this more general definition of classical dynamics can be devised.
### IV.2 Unitaries in $`\{U_c\}`$ and their action in phase space
We have just shown that when we build a complete set of MUB using Pauli operators, the ‘classical’ unitaries in $`\{U_c\}`$ turn out to form a subgroup of the Clifford group. The exact characterization of this subgroup will depend on which Pauli MUB construction we pick. In this section we restrict ourselves to the construction for $`N`$ qubits sketched in section II, and present some ‘classical’ unitaries in $`\{U_c\}`$ together with their associated action in phase space.
#### IV.2.1 Discrete phase-space translation operators
In PazRS04b it was shown that the Pauli operators act as discrete phase-space translations mapping phase-space lines into other lines. This means that the Pauli operators themselves are in our group $`\{U_c\}`$. Translation operators $`\widehat{T}(\stackrel{}{q},\stackrel{}{p})`$ operate on quantum states in such a way that their Wigner functions are transformed as flows: each phase-space point operator is mapped into another one because of the covariance condition imposed on the quantum net. Thus, the effect of a translation operator on a state’s Wigner function is a translation in phase space, and in this sense its action is “classical-like”. Since translation operators act as flows in phase space, they preserve positivity of $`W`$ for any single association between lines and MUB projectors.
#### IV.2.2 Discrete squeezing operator
The discrete squeezing operator PazRS04b maps the horizontal and the vertical striations into themselves, while cycling through all oblique striations. An explicit Clifford circuit for $`U_s`$ is given in PazRS04b . When acting on translation operators, $`U_s`$ maps them into other translations in a way that resembles a squeezing flow in phase space (see Fig. 1.b):
$$U_s\widehat{T}(\stackrel{}{q},\stackrel{}{p})U_s^{}=\pm \widehat{T}(\stackrel{}{q}M,\stackrel{}{p}\stackrel{~}{M}^1).$$
(33)
Besides covariance with respect to the discrete phase-space translations, we can impose on the quantum net also the constraint of covariance under $`U_s`$. In doing so, the freedom in picking the quantum net will be limited to the choice of which MUB projector to associate to a fixed oblique line, since the covariance requirement determines all other associations. In this way, the number of possible choices is greatly reduced from $`d^{d1}`$ to $`d`$.
If we choose a quantum net which is covariant under the squeezing operator, it can be shown that $`U_s`$ will map phase-space point operators into other point operators, i.e. $`U_s`$ acts like a phase-space flow. This may not be the case if the quantum net is not chosen to be covariant under $`U_s`$.
By definition, the group $`\{U_c\}`$ consists of unitaries that preserve non-negativity of the Wigner function $`W`$ for all possible quantum nets. This does not imply that operators in $`\{U_c\}`$ will preserve positivity for any single definition of $`W`$. This is because some states may have positive $`W`$ for a single definition, but negative $`W`$ for other definitions (and hence lie outside the set $`C_d`$). In the case of $`U_s`$, preservation of positivity for each single definition of $`W`$ is only guaranteed when $`U_s`$ acts as a flow in phase space, and this only happens when the quantum net is chosen to be covariant with respect to $`U_s`$.
#### IV.2.3 Finite Fourier transform
The finite Fourier transform $`F`$ KlimovS-SdG05 maps the horizontal and the vertical striations into one another; oblique striations are interchanged in pairs, and one of them \[the “main diagonal”, which corresponds to the eigenstates of the set of translations obtained by setting $`\stackrel{}{a}=\stackrel{}{b}`$ in eq. (10)\] is mapped into itself. For the particular case in which the canonical basis of the Galois field is self-dual, $`F`$ is just the Hadamard transform. The effect of $`F`$ on the translation operators is – up to a sign – a reflection with respect to the main diagonal of phase space (see Fig. 1.c):
$$F\widehat{T}(\stackrel{}{1}M^j,\stackrel{}{1}\stackrel{~}{M}^k)F^{}=\pm \widehat{T}(\stackrel{}{1}M^k,\stackrel{}{1}\stackrel{~}{M}^j)$$
(34)
The fact that $`F`$ interchanges translation operators by a reflection might suggest that its action on the states could be analogous, that is, that $`F`$ could reflect a state’s Wigner function with respect to the main diagonal, perhaps for some particular quantum nets (as is the case with $`U_s`$).
For $`F`$ to act on lines as a reflection with respect to the main diagonal, there should be one MUB projector associated to the axis of reflection, and $`F`$ should map this projector into itself. This projector should, then, be a common eigenstate of $`F`$ and all the Paulis that define the basis to which the state belongs. Using the fact that $`F`$ anti-commutes with some of them, and that the eigenvalues of $`F`$ and the Paulis are different from zero, it can be seen that no state can fulfill this requirement. Thus, there is no association between lines and MUB projectors that makes the action of $`F`$ on the Wigner function be a reflection flow.
Moreover, it can be seen that there is no quantum net for which $`F`$ acts as a flow in any way, because for $`F`$ to be a flow it should map phase-space point operators into other point operators. For this to happen, the $`(d+1)`$ lines that intersect in any given point must be mapped by $`F`$ into other lines that intersect in only one point. The vertical and the horizontal rays (i.e. lines containing the origin) are interchanged by $`F`$, so the other rays must be mapped into rays too (so that all the resulting lines intersect at the origin). This requires the ray in the main diagonal to be mapped into itself, and, as pointed out in the previous paragraph, this cannot be achieved.
Therefore, $`F`$ provides an example of an operator in $`\{U_c\}`$ which cannot be interpreted in terms of a flow for any choice of associations between lines and states, and so has no obvious continuous phase-space analogue.
## V Conclusion
We have characterized the set $`C_d`$ of states whose discrete Wigner functions $`W`$ (as defined in GibbonsHW04 ) are non-negative. We showed that the only pure states in $`C_d`$ are the mutually unbiased bases projectors used to define $`W`$, as conjectured in Galvao05 . Since these projectors can always be chosen to be stabilizer states, they admit an efficient classical description using the stabilizer formalism. Moreover, we proved that the unitaries which preserve non-negativity of $`W`$ for all such functions $`W`$ form a subgroup of the Clifford group. It is known that Clifford operations on stabilizer states can be simulated efficiently on a classical computer. We have thus identified a relation between two different notions of classicality: states which are classical in the sense of having non-negative quasi-probability distributions (the discrete Wigner functions of GibbonsHW04 ) can also be simulated efficiently on classical computers. Since general quantum computation is thought to be hard to simulate classically, our results mean that negativity of $`W`$ is necessary for exponential computational speedup with pure states.
There are many open problems worth investigating. The complete characterization of non-negativity preserving unitaries for different constructions of complete sets of MUB is still unsolved. It would also be interesting if one could relate non-classicality to negativity of $`W`$ in a quantitative way. Another research direction is to investigate the relationship between $`W`$ and a notable open problem, that of the existence of complete sets of mutually unbiased bases for general Hilbert space dimensions (see Bengtsson04 ; BengtssonE04 ; Wootters04b ). The original idea behind continuous-variable Wigner functions was to help visualize quantum dynamics in the familiar framework of classical phase space. Some research has been done on the visualization of quantum information protocols in discrete phase-space MiquelPS02 ; PittengerR05 ; Wootters04 ; PazRS04 ; KoniorczykBJ01 ; Paz02 ; LopezP03 ; PazRS04b ; further work might bring insights into existing applications, or suggest new ones.
Acknowledgments. EFG was partly supported by Canada’s NSERC. JPP was partially funded by Fundación Antorchas and a grant from ARDA. AOP was partially supported by NSF grants EIA-0113137 and DMS-0309042. |
warning/0506/astro-ph0506639.html | ar5iv | text | # Making the corona and the fast solar wind: a self-consistent simulation for the low-frequency Alfvén waves from photosphere to 0.3AU
## 1. Introduction
A certain portion of the solar surface is covered by the coronal holes. Their characteristic features are that the open magnetic fields are emanating into the interplanetary space and the hot plasma is streaming out. The high speed solar winds (800km/s at $`1`$AU) are accelerated from the polar coronal holes which stably exists except during the solar maximum phase (Phillips et al., 1995) The source of the energy for heating and accelerating the plasma is believed to be in the surface convection; this energy is lifted up through the magnetic fields and its dissipation results in the plasma heating. In this sense, the problem of the coronal heating and the solar wind acceleration is to solve how the solar atmosphere reacts to the footpoint motions of the magnetic fields.
In the coronal holes the Alfvén waves can play an important role because it can travel a long distance to both heat the corona and accelerate the solar wind. The Alfvén waves are excited by steady transverse motions of the field lines at the photosphere (e.g. Cranmer & van Ballegooijen 2005), while they can also be produced by reconnection events above the photosphere (Axford & Mckenzie, 1997). Although the latter process might be responsible for the generation of the high-frequency (up to $`10^4`$Hz) ioncyclotron wave highlighted for the heating of the heavy ions (Kohl et al., 1998), it is difficult to heat the protons (Cranmer, 2000). Because we focus on the heating of the major component of the plasma, we study the effect of the low-frequency ($`0.1`$Hz) Alfvén waves generated by the steady footpoint fluctuations.
Although the Alfvén waves have been intensively investigated, previous studies have a couple of limitations. Most of the calculations consider the waves from the fixed ‘coronal base’ of which density is $`9`$ orders of magnitudes lower than that of the photosphere (e.g. Oughton et al.2001; Ofman 2004). However, the corona dynamically interact with the chromosphere by chromospheric evaporation (Hammer, 1982) and wave transmission (Kudoh & Shibata, 1999) so that the coronal base varies time-dependently. Some calculations including the chromosphere even adopt an ad hoc heating function (e.g. Lie-Svendsen 2002). In contrast, we self-consistently treat the transfer of the mass/momentum/energy by solving the wave propagation from the photosphere to the interplanetary region. Our aim is to answer the problem of the heating and acceleration in the coronal holes by the forward approach (Gudiksen & Nordlund, 2005).
## 2. Method
We consider one-dimensional (1-D) super-radially open flux tube, measured by heliocentric distance, $`r`$. The computation domain is from the photosphere ($`r=1R_\mathrm{S}`$) with density, $`\rho =10^7`$g cm<sup>-3</sup>, to $`65R_\mathrm{S}`$ (0.3AU), where $`R_\mathrm{S}`$ is solar radius. Radial field strength, $`B_r`$, is given by conservation of magnetic flux as
$$B_rr^2f(r)=\mathrm{const}.,$$
(1)
where we adopt the same function as in Kopp & Holzer (1976) for superradial expansion, $`f(r)`$, but consider two-step expansions: the flux tube initially expands by a factor of 30 at $`r1.01R_\mathrm{S}`$ corresponding to ‘funnel’ structure (Tu et al., 2005), and followed by 2.5 times expansion at $`r1.2R_\mathrm{S}`$ due to the large scale magnetic fields. $`B_r`$ at the photosphere is set to be 161G. These give $`B_r=5`$G at $`r=1.02R_\mathrm{S}`$ (low corona), and $`B_r=2.1\mathrm{G}(R_\mathrm{S}/r)^2`$ in $`r>1.5R_\mathrm{S}`$. We give transverse fluctuations of the field line by the granulations at the photosphere, which excite Alfvén waves. We consider the fluctuations with power spectrum, $`P(\nu )\nu ^1`$, in frequency between $`6\times 10^4\nu 0.05`$Hz (period of 20 seconds — 30 minutes), and root mean squared (rms) average amplitude $`dv_{}0.7`$km/s corresponding to observed velocity amplitude $`1`$km/s (Holweger, Gehlsen, & Ruland, 1978). At the outer boundary, outgoing condition is imposed for all the MHD waves (Suzuki & Inutsuka 2005) to avoid the unphysical wave reflection.
We treat the dynamical evolution of the waves and the plasma by solving ideal MHD equations with the relevant physical processes :
$$\frac{d\rho }{dt}+\frac{\rho }{r^2f}\frac{}{r}(r^2fv_r)=0,$$
(2)
$$\rho \frac{dv_r}{dt}=\frac{p}{r}\frac{1}{8\pi r^2f}\frac{}{r}(r^2fB_{}^2)$$
$$+\frac{\rho v_{}^2}{2r^2f}\frac{}{r}(r^2f)\rho \frac{GM_\mathrm{S}}{r^2},$$
(3)
$$\rho \frac{d}{dt}(r\sqrt{f}v_{})=\frac{B_r}{4\pi }\frac{}{r}(r\sqrt{f}B_{}).$$
(4)
$$\rho \frac{d}{dt}(e+\frac{v^2}{2}+\frac{B^2}{8\pi \rho }\frac{GM_{}}{r})+\frac{1}{r^2f}\frac{}{r}[r^2f\{(p+\frac{B^2}{8\pi })v_r$$
$$\frac{B_r}{4\pi }(𝑩\mathbf{}𝒗)\}]+\frac{1}{r^2f}\frac{}{r}(r^2fF_\mathrm{c})+q_\mathrm{R}=0,$$
(5)
$$\frac{B_{}}{t}=\frac{1}{r\sqrt{f}}\frac{}{r}[r\sqrt{f}(v_{}B_rv_rB_{})],$$
(6)
where $`\rho `$, $`𝒗`$, $`p`$, $`e`$, $`𝑩`$ are density, velocity, pressure, specific energy, and magnetic field strength, respectively, and subscript $`r`$ and $``$ denote radial and tangential components. $`\frac{d}{dt}`$ and $`\frac{}{t}`$ denote Lagrangian and Eulerian derivatives, respectively. $`G`$ and $`M_\mathrm{S}`$ are the gravitational constant and the solar mass. $`F_\mathrm{c}`$ is thermal conductive flux and $`q_\mathrm{R}`$ is radiative cooling (Landini & Monsignori-Fossi, 1990; Anderson & Athay, 1989; Moriyasu et al., 2004). Note that the curvature effects appear as $`r\sqrt{f}`$ terms, instead of $`r`$ for the usual spherical coordinate. We adopt 2nd-order MHD-Godunov-MOCCT scheme for updating the physical quantities (Sano & Inutsuka 2005), of which an advantage is that no artificial viscosity is required even for strong MHD shocks. We fix 14000 grid points with variable sizes in a way to resolve Alfvén waves with $`\nu =0.05`$Hz by at least $`10`$ grids per wavelength to reduce the numerical damping.
The advantages of our calculations are automatic treatments of the wave propagation/dissipation and the plasma heating/acceleration without an ad hoc heating function but with the minimal assumption over the broadest regions with the huge density contrast. We emphasize that this is one of the most self-consistent simulations for the solar wind acceleration at present. The disadvantage is the 1-D MHD approximation, which will be discussed later.
## 3. Results and Discussions
We start the calculation from the static and cool atmosphere with temperature, $`T=10^4`$K; we do not set up the corona and the solar wind at the beginning. Figure 1 (on-line movie available) shows how the coronal heating and the solar wind acceleration are realized by the low-frequency Alfvén waves. We plot $`v_r`$(km/s), $`T`$(K), $`\rho `$(g/cm<sup>3</sup>), and $`dv_{}`$(km/s) averaged from $`v_{}`$ as a function of $`(rR_\mathrm{S})/R_\mathrm{S}`$ at different time, $`t=0,20,340`$ & 2573 minutes. As time goes on, the atmosphere is heated and accelerated effectively by the dissipation of the Alfvén waves. Temperature rises rapidly in the inner region even at $`t=20`$ minutes, and the outer region is eventually heated up by both outward thermal conduction and wave dissipation. Once the plasma is heated up to the coronal temperature, mass is supplied to the corona mainly by the chromospheric evaporation due to the downward thermal conduction. This is seen in temperature structure as an inward shift of the transition region, which is finally located around $`rR_\mathrm{S}=6\times 10^3R_\mathrm{S}`$ ($``$4000 km). As a result, the coronal density increases by two orders of magnitude. While the wind velocity exceeds 1000km/s at $`t=340`$ minutes on account of the initial low density, it gradually settles down to $`<1000`$ km/s as the density increases. $`dv_{}`$ also settles down to the reasonable value at the final stage.
We have found that the plasma is steadily heated up to $`10^6`$K in the corona and flows out as a transonic wind with $`v_r800`$km/s at the outer boundary ($`=`$0.3AU) when the quasi steady-state behaviors are achieved after $`t1800`$ minutes. This is the first numerical simulation which directly shows that the heated plasma actually flows out as the transonic wind, from the initially static and cool atmosphere, by the effects of the Alfvén waves. The sonic point where $`v_r`$ exceeds the local sound speed, $`c_\mathrm{s}=\sqrt{5p/3\rho }`$, is located at $`r2.5R_\mathrm{S}`$ and the Alfvén point for the Alfvén speed, $`B_r/\sqrt{4\pi \rho }`$, is at $`r24R_\mathrm{S}`$. The obtained proton flux at 0.3AU is $`N_pv(2\pm 0.5)\times 10^9`$(cm<sup>-2</sup>s<sup>-1</sup>), corresponding to $`N_pv(1.8\pm 0.5)\times 10^8`$(cm<sup>-2</sup>s<sup>-1</sup>) at 1AU for $`N_pvr^2`$, which is consistent with the observed high-speed stream around the earth (Aschwanden, Poland, & Rabin, 2001).
In fig.2 we compare the result at $`t=2573`$ minutes with recent observation in the high-speed solar winds from the polar regions by Solar & Heliospheric Observatory (SoHO) (Zangrilli et al., 2002; Teriaca et al., 2003; Fludra, Del Zanna, & Bromage, 1999; Lamy et al., 1997; Wilhelm et al., 1998; Banerjee et al., 1998; Esser et al., 1999) and Interplanetary Scintillation (IPS) measurements (Grall et al., 1996; Kojima et al., 2004; Canals et al., 2002). The figure shows that our forward-approach simulation naturally form the corona and the high-speed solar wind which are observed by the dissipation of the low-frequency Alfvén waves. Although the observed wind speed around $`r10R_\mathrm{S}`$ (Grall et al., 1996) appears to exceed our result, it might reflect the wave phenomena in the solar wind rather than the outflow (Harmon & Coles, 2005).
Figure 3 presents the dissipation of the waves. The top panel plots the following quantities,
$$S_c=\rho \delta v^2\frac{(v_r+v_{\mathrm{ph}})^2}{v_{\mathrm{ph}}}\frac{r^2f(r)}{r_c^2f(r_c)},$$
(7)
for outgoing Alfvén , incoming Alfvén , and outgoing slow MHD (sound) waves, where $`\delta v`$ and $`v_{\mathrm{ph}}`$ are amplitude and phase speed of each wave mode, and $`S_c`$’s are normalized at $`r_c=1.02R_\mathrm{s}`$. $`S_c`$ is an adiabatic constant derived from wave action (Jacques, 1977), and corresponds to the energy flux in static media. $`S_c`$ is conserved in the expanding atmosphere if the wave does not dissipate, while it is not the case for the energy flux. For the incoming Alfvén wave in Figure 3, we plot the opposite sign of $`S_c`$ so that it becomes positive in the sub-Alfvénic region. The outgoing and incoming Alfvén waves are decomposed by correlation between $`v_{}`$ and $`B_{}`$. Extraction of the slow wave is also from fluctuating components of $`v_r`$ and $`\rho `$. The bottom panel plots $`\delta v_{\mathrm{A},+}/v_\mathrm{A}`$, $`\delta v_{\mathrm{A},+}/c_\mathrm{S}`$, and $`\delta v_{\mathrm{S},+}/c_\mathrm{S}`$, where $`\delta v_{\mathrm{A},+}`$ and $`\delta v_{\mathrm{S},+}`$ are amplitudes of the outgoing Alfvén and slow modes. All the quantities in Figure 3 are time-averaged over 30 minutes to smooth out the variation due to the phase.
The top panel shows that the outgoing Alfvén waves dissipate quite effectively; $`S_c`$ becomes only $`10^3`$ of the initial value at the outer boundary. First, a sizable amount is reflected back downward below the coronal base ($`rR_\mathrm{S}<0.01R_\mathrm{S}`$), which is clearly illustrated as the incoming Alfvén wave following the outgoing component with slightly smaller level. This is because the wave shape is considerably deformed owing to the steep density gradient; a typical variation scale ($`<10^5`$km) of the Alfvén speed becomes comparable to or even shorter than the wavelength ($`=10^410^6`$km). Although the energy flux, $`5\times 10^5`$erg cm<sup>-2</sup>s<sup>-1</sup>, of the outgoing Alfvén waves ($`S_c`$ in the static region is equivalent with the energy flux) which penetrates into the corona is only $`15`$% of the input value, it satisfies the requirement for the energy budget in the coronal holes (Withbroe & Noyes, 1977).
Second, slow MHD waves are generated in the corona (see Sakurai 2002 for observation) as shown in the figure. The amplitude of the Alfvén waves is amplified through the upward propagation. The bottom panel shows that $`\delta v_{\mathrm{A},+}/c_\mathrm{S}1`$ in $`r2R_\mathrm{S}`$, which means that the wave pressure ($`B_{}^2/8\pi `$) exceeds the gas pressure (eq.) in spite of the weak nonlinearity, $`\delta v_{\mathrm{A},+}/v_\mathrm{A}0.1`$. This excites longitudinal slow waves (Kudoh & Shibata, 1999), which eventually become nonlinear $`\delta v_{\mathrm{S},+}/c_\mathrm{S}1`$ at $`r3R_\mathrm{S}`$ to lead to the shock dissipation (Suzuki, 2002). The coronal heating and wind acceleration are thus achieved by transferring the energy and momentum flux of the outgoing Alfvén waves. Fast MHD shocks by the direct steepening of the linearly polarized components of the Alfvén waves also contribute to the heating (Hollweg, 1982; Suzuki, 2004), though they are less efficient because of the weak nonlinearity. The incoming Alfvén waves are generated in the corona by the reflection of the outgoing ones by the density fluctuations due to the slow waves. The reflected waves further play a role in the wave dissipation by nonlinear wave-wave interaction.
A key mechanism in the heating and acceleration of the plasma is the generation of the slow MHD waves. Thus, one of the predictions from our simulation is the existence of the longitudinal fluctuations in the solar wind plasma. This is directly testable by in situ measurements of future missions, Solar Orbiter and Solar Probe, which will approach to $``$45 and 4 $`R_\mathrm{S}`$, respectively, corresponding to the inside of our computation domain.
We have shown by the self-consistent simulation that the dissipation of the low-frequency Alfvén waves through the generation of the compressive waves and shocks is one of the solutions for the heating and acceleration of the plasma in the coronal holes. However, the validity of the 1-D MHD approximation we have adopted needs to be examined. Generally, the shock dissipation tends to be overestimated in the 1-D simulation because the shocks cannot be diluted by the geometrical expansion. On the other hand, there are other dissipation mechanisms due to the multidimensionality (Ofman, 2004), such as turbulent cascade into the transverse direction (Oughton et al., 2001) and phase mixing (Heyvaerts & Priest, 1983). Therefore, the waves might also be dissipated by paths different from the shocks in the real situations. The self-consistent simulations including these various processes remain to be done to give the final conclusion.
We thank the referee for helpful comments and Drs. K. Shibata and T. Sano for many fruitful discussions. This work is supported in part by a Grant-in-Aid for the 21st Century COE “Center for Diversity and Universality in Physics” from the MEXT of Japan. T.K.S. is supported by the JSPS Research Fellowship for Young Scientists, grant 4607. SI is supported by the Grant-in-Aid (15740118, 16077202) from the MEXT of Japan. |
warning/0506/cond-mat0506098.html | ar5iv | text | # Phase diagram of spin-1 bosons on one-dimensional lattices
## Abstract
Spinor Bose condensates loaded in optical lattices have a rich phase diagram characterized by different magnetic order. Here we apply the Density Matrix Renormalization Group to accurately determine the phase diagram for spin-1 bosons loaded on a one-dimensional lattice. The Mott lobes present an even or odd asymmetry associated to the boson filling. We show that for odd fillings the insulating phase is always in a dimerized state. The results obtained in this work are also relevant for the determination of the ground state phase diagram of the $`S=1`$ Heisenberg model with biquadratic interaction.
The experimental realization of optical lattices general\_experiments has paved the way to study strongly correlated many-particle systems with cold atomic gases (see e.g. jaksch05 ; minguzzi03 ). The main advantages with respect to condensed matter systems lie on the possibility of a precise knowledge of the underlying microscopic models and an accurate and relatively easy control of the various couplings. Probably one of the most spectacular experiments in this respect is the observation greiner02 of a Superfluid - Mott Insulator transition previously predicted in jaksch98 by a mapping onto the Bose-Hubbard model fisher89 .
More recently the use of far-off-resonance optical traps has opened the possibility to study spinor condensates stamperkurn00 . Spin effects are enhanced by the presence of strong interactions and small occupation number, thus resulting in a rich variety of phases with different magnetic ordering. For spin-1 bosons it was predicted that the Mott insulating phases have nematic singlet demler02 or dimerized yip03 ground state depending on the mean occupation and on the value of the spin exchange. Since the paper by Demler and Zhou demler02 several works have addressed the properties of the phase diagram of spinor condensates trapped in optical lattices imambekov03 ; zhou03 ; duan03 ; svidzinsky03 ; imambekov04 ; snoek04 . The increasing attention in spinor optical lattices has also revived the attention on open problems in the theory of quantum magnetism. The spinor Bose-Hubbard model, when the filling corresponds to one boson per site, can be mapped onto the $`S=1`$ Heisenberg model with biquadratic interactions which exhibits a rich phase diagram including a long debated nematic to dimer quantum phase transition chubukov91 ; xian93 ; fath95 ; kawashima02 ; tsuchiya04 ; porras05 ; lauchli03 ; garcia-ripoll04 .
Up to now the location of the phase boundary of the spinor Bose-Hubbard model has been determined by means of mean-field and strong coupling approaches. A quantitative calculation of the phase diagram is however still missing. This might be particularly important in one dimension where non-perturbative effects are more pronounced. This is the aim of this Letter. We determine the location of the Mott lobes showing the even/odd asymmetry in the spinor case discussed in demler02 . We then discuss the magnetic properties of the first lobe, concluding that it is always in a dimerized phase.
The effective Bose-Hubbard Hamiltonian, appropriate for $`S=1`$ bosons, is given by
$`\widehat{}`$ $`=`$ $`{\displaystyle \frac{U_0}{2}}{\displaystyle \underset{i}{}}\widehat{n}_i(\widehat{n}_i1)+{\displaystyle \frac{U_2}{2}}{\displaystyle \underset{i}{}}\left(\widehat{𝐒}_i^22\widehat{n}_i\right)\mu {\displaystyle \underset{i}{}}\widehat{n}_i`$ (1)
$``$ $`t{\displaystyle \underset{i,\sigma }{}}\left(\widehat{a}_{i,\sigma }^{}\widehat{a}_{i+1,\sigma }+\widehat{a}_{i+1,\sigma }^{}\widehat{a}_{i,\sigma }\right).`$
The operator $`\widehat{a}_{i,\sigma }^{}`$ creates a boson in the lowest Bloch band localized on site $`i`$ and with spin component $`\sigma `$ along the quantization axis: $`\widehat{n}_i=_\sigma \widehat{a}_{i,\sigma }^{}\widehat{a}_{i,\sigma }`$ and $`\widehat{𝐒}_i=_{\sigma ,\sigma ^{}}\widehat{a}_{i,\sigma }^{}𝐓_{\sigma ,\sigma ^{}}\widehat{a}_{i,\sigma ^{}}`$ are the total number of particles and the total spin on site $`i`$ ($`\widehat{𝐓}`$ are the spin-1 operators). Atoms residing on the same lattice site have identical orbital wave function and their spin function must be symmetric. This constraint imposes that $`S_i+n_i`$ must be even. The uniqueness of the completely symmetric state with fixed spin and number makes it possible to denote the single site states with $`|n_i,S_i,S_i^z`$. The coupling constants, which obey the constraint $`1<U_2/U_0<1/2`$, can be expressed in terms of the appropriate Wannier functions imambekov03 . $`U_0`$ is set as the energy scale unit: $`U_0=1`$. We discuss only the anti-ferromagnetic case ($`0<U_2<1/2`$).
In the absence of spin dependent coupling a qualitative picture of the phase diagram can be drawn starting from the case of zero hopping ($`t=0`$). The ground state is separated from any excited state by a finite energy gap. For finite hopping strength, the energy cost to add or remove a particle $`\mathrm{\Delta }E_\pm `$ (excitation gap) is reduced and at a critical value $`t_c^\pm (\mu )`$ vanishes. This phase is named the Mott insulator. For large hopping amplitudes the ground state is a globally coherent superfluid phase. When $`U_2`$ is different from zero, states with lowest spins, compatible with the constraint $`n_i+S_i=\text{even}`$, are favoured. This introduces an even/odd asymmetry of the lobes: the amplitude of lobes with odd filling is reduced as compared with the lobes corresponding to even fillings demler02 . In the first lobe the extra energy required to have two particles on a site (instead of one) is $`1+2U_2\mu `$, thus lowering the chemical potential value where the second lobe starts. On the other hand, having no particles on a site gives no gain due to spin terms, accounting for the nearly unvaried bottom boundary of the lobe.
In order to determine the phase diagram of Eq.(1) we use the finite-size Density Matrix Renormalization Group (DMRG) with open boundary conditions white92-95 . The strategy of the DMRG is to construct a portion of the system (called the system block) and then recursively enlarge it, until the desired system size is reached. At every step the basis of the corresponding Hamiltonian is truncated, so that the size of the Hilbert space is kept manageable as the physical system grows. The truncation of the Hilbert space is performed by retaining the eigenstates corresponding to the $`m`$ highest eigenvalues of the block’s reduced density matrix.
The DMRG has been employed, for the spinless case, in kuhner98 ; kuhner00 . The presence of the spin degree of freedom makes the analysis considerably more difficult. In the numerical calculations the Hilbert space for the on-site part of the Hamiltonian is fixed by imposing a maximum occupation number $`\overline{n}_{max}`$. As the first lobe is characterized by an insulating phase with $`n=1`$ particle per site we choose $`\overline{n}_{max}=3`$ in this case; the dimension of the Hilbert space per site becomes $`d=20`$. We have checked, by increasing the value of $`\overline{n}_{max}`$, that this truncation of the Hilbert space is sufficient to compute the first lobe. In each DMRG iteration we keep up to $`m=300`$ states in order to guarantee accurate results. The numerical calculations of the second lobe ($`n=2`$ particles per site) have been performed with $`\overline{n}_{max}=4`$ (which corresponds to $`d=35`$).
Phase Diagram \- In the insulating phase the first excited state is separated by the ground state by a Mott gap. In the limit of zero hopping the gap is determined by the extra energy $`\mathrm{\Delta }E_\pm `$ needed to place/remove a boson at a given site. The finite hopping renormalizes the gap which will vanish at a critical value. Then the system becomes superfluid. This method has been employed for the spinless case by Freericks and Monien freericks96 , and in kuhner98 ; kuhner00 where it was combined with the DMRG. Here we use it for the spinor case. Three iterations of the DMRG procedure are performed, with projections on different number sectors; the corresponding ground states give the desired energies $`E_0`$, $`E_\pm =E_0+\mathrm{\Delta }E_\pm `$. As target energies we used those obtained by the mapping of the Bose-Hubbard system into effective models as described in imambekov03 . We considered chains up to $`L=128`$ sites for the first lobe, and $`L=48`$ for the second lobe. The extrapolation procedure to extract the asymptotic values was obtained by means of linear fit in $`1/L`$, as discussed in kuhner00 . A comparison with a quadratic fit shows that $`O(1/L^2)`$ corrections are negligible on the scale of Fig.1.
The plot of the phase diagram in the $`(\mu ,t)`$ plane for different values of the spin coupling $`U_2`$ is shown in Fig.1. The first lobe tends to reduce its size on increasing the spin coupling; in particular the upper critical chemical potential at $`t=0`$ is $`\mu _c^+(0)=12U_2`$, while the $`t^{}`$ value of the hopping strength over which the system is always superfluid is suppressed as $`U_2`$ increases. On the other hand, the second lobe grows up when $`U_2`$ increases. This even/odd effect, predicted in demler02 , is quantified in Fig.1.
Magnetic properties of the first Mott lobe \- The first lobe of the spinor Bose lattice has a very interesting magnetic structure. In presence of small hopping $`t`$ boson tunneling processes induce effective pairwise magnetic interactions between the spins, described by Hamiltonian imambekov03 :
$$\widehat{}_{\mathrm{eff}}=\kappa \underset{ij}{}\left[\mathrm{cos}\theta (\widehat{𝐒}_i\widehat{𝐒}_j)+\mathrm{sin}\theta (\widehat{𝐒}_i\widehat{𝐒}_j)^2\right]$$
(2)
with
$$\mathrm{tan}\theta =\frac{1}{12U_2}\kappa =\frac{2t^2}{1+U_2}\sqrt{1+\mathrm{tan}^2\theta }.$$
(3)
The absence of higher order terms, such as $`(\widehat{𝐒}_i\widehat{𝐒}_j)^3`$, is due to the fact that the product of any three spin operators can be expressed via lower order terms. In the case of anti-ferromagnetic interaction in Eq.(1), the parameter $`\theta `$ varies in the interval $`\theta [3/4\pi ,\pi /2[`$. Because of the form of the magnetic Hamiltonian, each bond tends to form a singlet-spin configuration, but singlet states on neighboring bonds are not allowed. There are two possible ground states that may appear in this situation. A nematic state can be constructed by mixing states with total spin $`S=0`$ and $`S=2`$ on each bond. This construction can be repeated on neighboring bonds, thereby preserving translational invariance. This state breaks the spin-space rotational group $`O(3)`$, though time-reversal symmetry is preserved. The expectation value of any spin operator vanishes ($`\widehat{S}_i^\alpha =0,\alpha =x,y,z`$), while some of the quadrupole operators have finite expectation values. The tensor $`𝒬^{ab}=\widehat{S}^a\widehat{S}^b\frac{2}{3}\delta ^{ab}`$ is a traceless diagonal matrix, due to invariance under spin reflections. Since it has two identical eigenvalues ($`(\widehat{S}_i^x)^2=(\widehat{S}_i^y)^2(\widehat{S}_i^z)^2`$), it can be written as $`𝒬^{ab}=Q\left(d^ad^b\frac{1}{3}\delta ^{ab}\right)`$ using an order parameter $`\widehat{Q}(\widehat{S}_i^z)^2(\widehat{S}_i^x)^2=`$ $`\frac{3}{2}(\widehat{S}_i^z)^21`$ and a unit vector $`𝐝=\pm 𝐳`$. However, since $`[\widehat{Q},\widehat{}_{\mathrm{eff}}]=0`$, it is not possible to get $`Q0`$ in finite-size systems, analogously to what happens for the magnetization without external field. Therefore we characterized the range of nematic correlations in the ground state by coupling this operator to a fictitious “nematic field”: $`\widehat{}_\lambda =\widehat{}_{\mathrm{eff}}+\lambda \widehat{Q}`$, and by evaluating the nematic susceptibility $`\chi _{\mathrm{nem}}`$ as a function of $`L`$:
$$\chi _{\mathrm{nem}}\frac{\mathrm{d}^2E_0(\lambda )}{\mathrm{d}\lambda ^2}|_{\lambda =0}=\underset{\gamma }{}\frac{\left|Q_{0,\gamma }\right|^2}{E_\gamma E_0},$$
(4)
where $`E_0(\lambda )`$ is the ground energy of $`\widehat{}_\lambda `$, $`Q_{0,\gamma }`$ is the matrix element between the ground and an excited state of $`\widehat{}_{\mathrm{eff}}`$ (respectively with energy $`E_0`$ and $`E_\gamma `$).
On the other hand a possibility to have $`SO(3)`$ symmetric solution stems from breaking translational invariance. Indeed a dimerized solution with singlets on every second bond satisfy these requirements. Dimerization could be described looking at the differences in expectation values of pair Hamiltonian $`\widehat{}_{\mathrm{eff}}^{(ij)}`$ on adjacent links ($`\widehat{}_{\mathrm{eff}}=_{ij}\widehat{}_{\mathrm{eff}}^{(ij)}`$footnote . The order parameter $`D`$ reads
$$D\left|\widehat{}_{\mathrm{eff}}^{(i1,i)}\widehat{}_{\mathrm{eff}}^{(i,i+1)}\right|.$$
(5)
It has been proposed chubukov91 that a narrow nematic region exists between the ferromagnetic phase boundary ($`\theta _F=3\pi /4`$, i.e. $`U_2=0`$) and a critical angle $`\theta _C0.7\pi `$ (i.e. $`U_210^2`$), whereas a dimerized solution is favoured in the remaining anti-ferromagnetic region $`\theta _C\theta \pi /2`$. This implies that the dimerization order parameter D should scale to zero in the whole nematic region. This possibility has been analyzed in Ref. fath95 where it was suggested that $`D`$ might go to $`0`$ exponentially near the ferromagnetic boundary, making it difficult to detect the effective existence of the nematic phase. This interesting challenge has motivated numerical investigations with different methods fath95 ; kawashima02 ; porras05 ; lauchli03 . We present new DMRG results which clarify the magnetic properties of the first Mott lobe (for sufficiently small hopping) and, consequently, of the biquadratic Heisenberg chain.
According to our numerical calculation there is no intermediate nematic phase, indeed we found a power law decay of the dimerization order parameter near $`\theta _F=3\pi /4`$. The simulations of the bilinear-biquadratic model (2) are less time and memory consuming than Bose-Hubbard ones, since the local Hilbert space has a finite dimension $`d=3`$. The number of block states kept during the renormalization procedure was chosen step by step in order to avoid artificial symmetry breaking. This careful treatment insures that there are no spurious sources of asymmetry like partially taking into account a probability multiplet. Here we considered up to $`m300`$ states in order to obtain stable results. Raw numerical data are shown in Fig.2, where the finite-size dimerization parameter $`D(L)`$ is plotted as a function of the chain length $`L`$ (see Eq. 5, and footnote ). Finite-size scaling was used to extrapolate to the thermodynamic limit. After the extrapolation to the $`L\mathrm{}`$ limit, see Fig. 3, we fitted the dimer order parameter with a power law
$$D=\left(\frac{\theta \theta _F}{\theta _0}\right)^\gamma $$
(6)
where $`\gamma 6.1502`$ and $`\theta _00.09177\pi `$ (Fig. 3, solid line). We also tried to fit our data by an exponential law
$$D=D_0e^{a/\sqrt{\theta \theta _F}}$$
(7)
as suggested in fath95 , with $`a2.911`$, $`D_09.617`$; this fit seems to work for narrower regions (Fig. 3, dashed line), however from our numerics we cannot exclude an exponential behavior of $`D`$ in the critical region. The dimerized phase thus seems to survive up to the ferromagnetic phase boundary, independently from the chosen fitting form. This is also confirmed by the fact that the scaled gap between the ground state $`E_0`$ and the lowest excited state $`E_2`$ (which is found to have total spin $`S_T=2`$) seems not to vanish in the interesting region $`\theta >0.75\pi `$ (see inset of Fig. 3).
Moreover we analyzed the susceptibility of the chain to nematic ordering $`\chi _{\mathrm{nem}}`$. The numerical data, presented in Fig. 4, show a power law behavior $`\chi _{\mathrm{nem}}(L)L^\alpha `$ as a function of the system size. The exponent $`\alpha `$ (shown in the inset) approaches the value $`\alpha =3`$ as $`\theta \theta _F`$. This can also be confirmed by means of a perturbative calculation around the exact solution available at $`\theta _F`$; indeed one obtains $`\left|Q_{0,\gamma }\right|^2L^2`$ and $`(E_\gamma E_0)L^1`$ to be inserted in Eq. (4). The increase of the exponent for $`\theta \theta _F`$ indicates, as suggested in porras05 , that a tendency towards the nematic ordering is enhanced as the dimer order parameter goes to zero.
Conclusions \- In this Letter we analyzed, by means of a DMRG analysis, the phase diagram of the one-dimensional spinor boson condensate on an optical lattice. We determined quantitatively the shape of the first two Mott lobes, and the even/odd properties of the lobes. We furthermore discussed the magnetic properties of the first lobe. Our results indicate that the Mott insulator is always in a dimerized phase.
This work was supported by IBM (2005 Faculty award), and by the European Community through grants RTNNANO, SQUBIT2. |
warning/0506/gr-qc0506028.html | ar5iv | text | # False vacuum decay in a brane world cosmological model
## 1 Introduction
The total energy-density of our universe is very close to the critical value corresponding to the Friedmann-Robertson-Walker metrics of the flat type. Numerically this means that the total matter-density parameter has the value in the thin range around $`1`$: $`\mathrm{\Omega }=1.02\pm 0.04`$. There are two possible interpretations of this fact: the universe is exactly flat; or the early evolution forced the universe to evolve to the present state which is very close to the flat Friedmann-Robertson-Walker metric, however the geometry may be both open or closed. There are models within open inflationary universe scenario, e.g., , that can successfully lead to an acceptable present-day universe. These models are built on Einstein’s theory of relativity in four dimensional space-time. There are also models of ”creation of an infinite universe within a finite bubble” based on modifications of the Einstein’s relativity theory like the model presented in where the false vacuum decay via Coleman-de Luccia (CdL) instanton and subsequent second phase of inflation within the nucleated bubble is studied in the context of Jordan-Brans-Dicke theory.
The progress in the superstring theory during last years has forced the cosmologists to consider the extra dimensions in various models of the universe evolution e.g. and . Especially, the gravity-scalar instantons have been considered in . The creation of an open or closed universe within the brane world scenario has been considered in . The physical and geometrical discussion of the semiclassical instability of the Randall-Sundrum brane world resulting in the vacuum decay via instantons is done in , .
In this paper we consider the decay of the false vacuum of a scalar field (inflaton) confined to a four dimensional brane in a five dimensional brane world model. We analyze the Euclidean cosmological equations in the Randall-Sundrum type II scenario that are supposed to describe semiclassically the false vacuum decay. We are inspired by the work done by del Campo, Herrera and Saavedra in which the authors investigate the possibility of realization of the open inflation scenario in the brane world models including the existence of the Coleman-de Luccia instanton providing the false vacuum decay and subsequent inflation within created open universe. The authors of the cited paper are interested in a specially chosen theory (the self interaction $`V`$ of the scalar field). They investigate the model only for fixed parameters of the self-interaction and obtain both CdL instanton and plausible evolution after tunneling. Our aim is to study the CdL instantons for arbitrary potentials and compare the properties of CdL instantons in the standard Einstein’s gravity in four dimensional space-time with those from our brane world. A similar problem has been recently considered in , where the authors study the vacuum decay on the brane within the thin-wall approximation. Unlike the authors of the paper we are interested in another problem that can be solved analytically, namely the CdL instanton(s) of the first order close to a Hawking-Moss instanton.
The paper is organized as follows: in section 2 we briefly review the basic fact about the CdL instantons in four dimensional de Sitter space-time to be able to compare them with the results of this paper. In section three we present the formulation of the instanton equations in considered brane world model and some consequences of these equations are discussed. Perturbative computation of the first order CdL instanton in our brane world model is presented in section four, and finally the main quantity characterizing the instanton - its action - is computed in the fifth section.
## 2 False vacuum decay via CdL instanton in four dimensional space-time within Einstein’s relativity theory
CdL instanton introduced in describes false vacuum decay in a de Sitter universe within the semiclassical approximation. If $`V=V(\mathrm{\Phi })`$ is the effective potential for the scalar field, this CdL instanton can be introduced as the $`O(4)`$-symmetric and finite-action solution of the Euclidean version of the Einstein equations. The $`O(4)`$-symmetry means that the scalar field $`\mathrm{\Phi }`$ lives on a (squeezed) four-sphere with the metric
$$\mathrm{d}s^2=\mathrm{d}\tau ^2+a^2(\tau )\left[\mathrm{d}\chi ^2+\mathrm{sin}^2(\chi )\mathrm{d}\mathrm{\Omega }_2^2\right],$$
(1)
with $`\tau [0,\tau _f]`$. The action (Euclidean version of the Einstein-Hilbert action) reads
$`S=`$
$`2\pi ^2{\displaystyle _0^{\tau _f}}\left[\left({\displaystyle \frac{1}{2}}(\mathrm{\Phi }^{})^2+V\right)a^2{\displaystyle \frac{1}{C}}(a(a^{})^2+1)\right]ad\tau ,`$ (2)
where $`C=8\pi /3`$ and the prime denotes the derivative with respect to $`\tau `$. Varying the action (2) we obtain the (Euclidean) equations of motion for the scale parameter and the inflaton:
$$a^{\prime \prime }(\tau )=C\left((\mathrm{\Phi }^{})^2+V\right)a,\mathrm{\Phi }^{\prime \prime }+3\frac{a^{}}{a}\mathrm{\Phi }^{}_\mathrm{\Phi }V=0.$$
(3)
The local energy conservation law and the requirement of finiteness of the action impose the boundary conditions on the functions $`a`$ and $`\mathrm{\Phi }`$
$$a(0)=0,a^{}(0)=1,\mathrm{\Phi }^{}(0)=\mathrm{\Phi }^{}(\tau _f)=0,$$
(4)
where $`\tau _f>0`$ is to be determined from $`a(\tau _f)=0`$. The action (2) of a CdL instanton can be considerably simplified by using the equations of motion (3):
$$S=\frac{4\pi ^2}{3C}_0^{\tau _f}ad\tau .$$
After Coleman and de Luccia have proposed the idea of the description of the vacuum decay in , many authors have studied the system (3) and (4). Coleman and de Luccia themselves have found the solution of the instanton equations in the thin-wall limit. We suppose the effective potential $`V`$ is non-negative function with two non-degenerate minima. The top of the potential barrier is reached at the point we denote by $`\mathrm{\Phi }_M`$. Furthermore, we denote: $`V_MV(\mathrm{\Phi }_M)`$, $`V_M^{\prime \prime }_\mathrm{\Phi }^2V(\mathrm{\Phi }_M)`$, $`H_M^2=CV_M`$ and $`H^2(\mathrm{\Phi })=CV(\mathrm{\Phi })`$. The properties of $`V`$ in the neighborhood of $`\mathrm{\Phi }_M`$ are crucial for the existence of CdL instanton. Motivated by the earlier works and the authors of the papers , and have achieved information about the solution(s) of the instanton equations interesting for us:
* if $`V_M^{\prime \prime }/H_M^2>4`$ then the CdL instanton exists
* if a CdL instanton exist then $`V^{\prime \prime }(\mathrm{\Phi })/H^2(\mathrm{\Phi })>4`$ somewhere in the barrier
* if $`V_M^{\prime \prime }/H_M^2>l(l+3)`$, where $`l`$ is an arbitrary integer, then CdL instaton crossing in the $`\mathrm{\Phi }`$-direction $`l`$-times the top of the barrier ($`l`$th order CdL instnton) exists
* if $`V_M^{\prime \prime }/H_M^2l(l+3)`$ we have the explicit approximative formulas for the instanton and its action,
## 3 False vacuum decay on a brane - elementary discussion of the instanton equations
The dynamics of the scalar field confined to a four dimensional brane in a five dimensional bulk in our model is defined by the action
$$S=M_{(5)}^3\mathrm{d}^5x\sqrt{|G|}\left[R_{(5)}2\mathrm{\Lambda }_{(5)}\right]\mathrm{d}^4x\sqrt{|g|}_\mathrm{\Phi }$$
(5)
with the matter term
$$_\mathrm{\Phi }=\frac{1}{2}_\mu \mathrm{\Phi }^\mu \mathrm{\Phi }V(\mathrm{\Phi }),$$
where $`R_{(5)}`$ is the scalar curvature of the five-dimensional bulk metric $`G`$, $`M_{(5)}`$ and $`\mathrm{\Lambda }_{(5)}`$ stand for the five-dimensional Planck mass and cosmological constant, respectively. These quantities are related to the effective four-dimensional cosmological constant $`\mathrm{\Lambda }_{(4)}`$, the brane tension $`\sigma `$ and the four-dimensional Planck mass $`M_{(4)}`$ by the relations
$`\mathrm{\Lambda }_{(4)}={\displaystyle \frac{4\pi }{M_{(5)}^3}}\left(\mathrm{\Lambda }_{(5)}+{\displaystyle \frac{4\pi }{3M_{(5)}^3}}\sigma ^2\right),`$
$`M_{(4)}=\left({\displaystyle \frac{3}{4\pi }}\right)^{1/2}{\displaystyle \frac{M_{(5)}^3}{\sigma ^{1/2}}}.`$ (6)
We will use the units where $`M_{(4)}=1`$. Following the works and we come at the Euclidean equations of motion for the inflaton $`\mathrm{\Phi }`$ on the brane and the induced metric (under the assumption of $`O(4)`$-symmetry which involves the line element of the form $`\mathrm{d}s^2=\mathrm{d}\tau ^2+a^2(\tau )[\mathrm{d}\chi ^2+\mathrm{sin}^2(\chi )\mathrm{d}\mathrm{\Omega }_2^2]`$):
$`a^{\prime \prime }=C\left\{(\mathrm{\Phi }^{})^2+V+{\displaystyle \frac{1}{8\sigma }}\left[(5(\mathrm{\Phi }^{})^2+2V)((\mathrm{\Phi }^{})^2+2V)\right]\right\}a,`$
$`\mathrm{\Phi }^{\prime \prime }+3{\displaystyle \frac{a^{}}{a}}\mathrm{\Phi }^{}_\mathrm{\Phi }V=0.`$ (7)
The functions $`a`$ and $`\mathrm{\Phi }`$ must obey the boundary conditions (4). We see that in the $`\sigma +\mathrm{}`$ limit we recover the standard general-relativistic equation for the scale parameter $`a`$ , the equation for $`\mathrm{\Phi }^{\prime \prime }`$ remains unchanged with respect to the Einstein’s relativity theory. If we assume that $`(\stackrel{~}{a},\stackrel{~}{\mathrm{\Phi }})`$ is a CdL instanton with $`\mathrm{\Phi }(0)=\mathrm{\Phi }_i`$ and $`\stackrel{~}{\mathrm{\Phi }}(\tau _f)=\mathrm{\Phi }_f`$, then we can write for $`\tau 0^+`$:
$$\stackrel{~}{\mathrm{\Phi }}^{\prime \prime }+\frac{3}{\tau }\stackrel{~}{\mathrm{\Phi }}^{}_\mathrm{\Phi }V(\mathrm{\Phi }_i)=0\stackrel{~}{\mathrm{\Phi }}=\mathrm{\Phi }_i+\frac{_\mathrm{\Phi }V(\mathrm{\Phi }_i)}{8}\tau ^2,$$
(8)
and for $`\tau \tau _f^{}`$:
$$\stackrel{~}{\mathrm{\Phi }}^{\prime \prime }+\frac{3}{\tau }\stackrel{~}{\mathrm{\Phi }}^{}_\mathrm{\Phi }V(\mathrm{\Phi }_f)=0\stackrel{~}{\mathrm{\Phi }}=\mathrm{\Phi }_f+\frac{_\mathrm{\Phi }V(\mathrm{\Phi }_f)}{8}(\tau _f\tau )^2.$$
(9)
Under the assumption that $`a`$ is a concave function (surely guaranteed by the positivity of the term proportional to $`1/\sigma `$ in the equation for $`a^{\prime \prime }`$) we deduce from eqs. (8) and (9) that the CdL instanton (in $`\mathrm{\Phi }`$-direction) must cross the value $`\mathrm{\Phi }_M`$ once at least. The asymptotic of the non-instanton solution (we can say - solutions with ”randomly” chosen initial value of $`\mathrm{\Phi }`$) of the system of equation (3) can be found in the same way as in Einstein’s general relativity.
We are interested in solutions which are close to the so-called Hawking-Moss (HM) instanton that describes the false vacuum tunneling as a process in which the inflaton ”jumps” (within a horizon-size domain) at the top ($`\mathrm{\Phi }_M`$) of the potential $`V`$. The HM instanton is the $`O(5)`$-symmetric (constant $`\mathrm{\Phi }`$) solution of the system (3):
$$\mathrm{\Phi }=\mathrm{\Phi }_M,a=\widehat{H}_M^1\mathrm{sin}\left(\widehat{H}_M\tau \right),$$
(10)
where $`\widehat{H}_M`$ is a modification of the Hubble parameter $`H_M`$ introduced in the previous section. $`\widehat{H}_M`$ is determined inserting the proposed solution into the first equation of (3). One easily obtains
$$\widehat{H}_M^2=\frac{8\pi }{3}\left(V_M+\frac{V_M^2}{2\sigma }\right)=H_M^2\left(1+\frac{V_M}{2\sigma }\right).$$
We can study the CdL instantons close to this HM instanton in the following way. We insert the expression for $`a`$ from (10) into the equation for $`\mathrm{\Phi }^{\prime \prime }`$, linearize the term $`_\mathrm{\Phi }V`$ and using new variables: $`x=\widehat{H}_M\tau `$ and $`y=\mathrm{\Phi }\mathrm{\Phi }_M`$ we obtain
$$\frac{\mathrm{d}^2y}{\mathrm{d}x^2}+3\mathrm{cot}(x)\frac{\mathrm{d}y}{\mathrm{d}x}\frac{V_M^{\prime \prime }}{\widehat{H}_M^2}y=0,$$
or transforming the independently variable $`x`$ to: $`z=\mathrm{cos}(\widehat{H}_M\tau )`$ we get the standard hypergeometric equation:
$$(1z^2)\frac{\mathrm{d}^2y}{\mathrm{d}z^2}4z\frac{\mathrm{d}y}{\mathrm{d}z}\frac{V_M^{\prime \prime }}{\widehat{H}_M^2}y=0.$$
The boundary conditions (4) restrict possible values of the parameter $`\frac{V_M^{\prime \prime }}{\widehat{H}_M^2}`$ to the eigenvalues of the Laplace-Beltrami operator on $`S^4`$:
$$\frac{V_M^{\prime \prime }}{\widehat{H}_M^2}=l(l+3),l\{0,1,2,\mathrm{}\}$$
(11)
and the solutions $`y=y_l`$ read for odd $`l`$:
$$y_l=c_lz_2F_1(\frac{1l}{2},2+\frac{l}{2},\frac{3}{2},z^2)$$
(12)
and for even $`l`$:
$$y_l=c_l{}_{2}{}^{}F_{1}^{}(\frac{3+l}{2},\frac{l}{2},\frac{1}{2},z^2),$$
(13)
where $`c_l`$ are arbitrary constants and $`{}_{2}{}^{}F_{1}^{}`$ stands for nondegenerate hypergeometric function. (In fact, the hypergeometric functions with special arguments according (12),(13) reduce to the Gegenbauer polynomials in the variable $`z`$. However, we will not need this explicitly.) The function $`y_0`$ correspond to the HM instanton and the functions $`y_l`$ approximate the $`l`$th order CdL instanton in its $`\mathrm{\Phi }`$-direction. The restriction (11) is formally the same as in the case of four dimensional space-time with $`H_M`$ changed to $`\widehat{H}_M`$. This change means that we have a new parameter (except the old one $`V_M^{\prime \prime }/H_M^2`$) which value is crucial for the existence of the CdL instanton, namely $`V_M/\sigma `$. It is obvious that for $`V_M/\sigma 1`$ the theory of vacuum decay on our brane reduces to the theory of vacuum decay in four-dimensional space-time. The first-order CdL instanton plays the most important role in the Einstein’s theory of gravity , , therefore we write down explicitly:
$$y_1=kz=k\mathrm{cos}\left(x\right),$$
(14)
with $`k`$ \- the amplitude of the inflaton during its Euclidean evolution. In the next we will be interested in the first-order CdL instanton only.
## 4 The first order CdL instanton - perturbative approach
The idea of our analysis of the system of equations (3) is to expand all the relevant quantities entering these equations (and the boundary conditions (4)) into the powers of the $`\mathrm{\Phi }`$ amplitude $`k`$, see (14). This means explicitly that the following formulas are of our interest:
$`y(x)={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}k^nu_n(x),a(x)=\widehat{H}_M^1{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}k^nv_n(x),`$
$`{\displaystyle \frac{V_M^{\prime \prime }}{\widehat{H}_M^2}}=4+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}k^n\mathrm{\Delta }_n,`$
together with the Taylor expansion of the potential $`V`$ around its local maximum. It can be useful to write down the form of the system of linear equations by which we replace equations (3):
$`u_n^{\prime \prime }+3\mathrm{cot}(x)u_n^{}+4u_n=𝒰_n,v_n^{\prime \prime }+v_n=𝒱_n\mathrm{sin}(x),`$
where the source-terms $`𝒰_n`$ and $`𝒱_n`$ are to be computed order-by-order. We know, from the previous section, that
$$u_0=0,v_0(x)=\mathrm{sin}(x),u_1(x)=\mathrm{cos}(x).$$
The boundary conditions imposed on $`a`$ and $`\mathrm{\Phi }`$ require that for all $`n1`$ we have $`v(0)=v^{}(0)=0`$. The right-end point $`x_f`$ at which the derivative of $`y`$ has to vanish is determined by $`a(x_f)=0`$, therefore we have also the expansion of this quantity:
$$x_f=\pi +\underset{n=1}{\overset{\mathrm{}}{}}k^nx_f^{(n)}.$$
The fact that the potential $`V`$ is supposed to have the local maximum at $`\mathrm{\Phi }_M`$ implies that $`v_1=0`$ and $`x_f^{(1)}=0`$. The contribution of the second order in $`k`$ to the scale factor $`a`$ is nonzero, and reads explicitly
$`v_2={\displaystyle \frac{1}{32}}\{[C+{\displaystyle \frac{159C}{4\sigma }}V_M]4x\mathrm{cos}(x)+`$
$`\left[5C+{\displaystyle \frac{51C45}{4\sigma }}V_M\right]\mathrm{sin}(x)`$
$`[3C+5(1+C){\displaystyle \frac{V_M}{4\sigma }}]\mathrm{sin}(3x)\}.`$ (15)
Having this results we come at the shift of the right-end point $`x_f`$
$$x_f^{(2)}=\frac{\pi }{8}\left[C+\frac{3V_M}{4\sigma }(53C)\right].$$
(16)
In the limit that $`\sigma \mathrm{}`$ this quantity is negative but for a finite value $`V_M/\sigma `$ it can be both negative or positive and it vanishes at $`V_M/\sigma 0.555`$. Careful and a little bit tedious computation shows that $`\mathrm{\Delta }_1=0`$ and that $`u_2`$ obeys equation
$$u_2^{\prime \prime }+3\mathrm{cot}(x)u_2^{}+4u_2=\frac{1}{2}\frac{V_M^{\prime \prime \prime }}{\widehat{H}_M^2}u_1^2.$$
The right-end point is still $`\pi `$, i.e. we seek for the solution(s) for which $`u_2^{}(0)=u_2^{}(\pi )=0`$. This determines $`u_2`$ as follows
$$u_2=\frac{1}{24}\frac{V_M^{\prime \prime \prime }}{\widehat{H}_M^2}\left[12\mathrm{cos}^2(x)\right].$$
(17)
We continue with time-consuming computations without any extra-idea and derive the equation for $`v_3`$ we do not write down. However, the explicit formula for $`v_3`$ is
$`v_3={\displaystyle \frac{V_M^{\prime \prime \prime }}{288\widehat{H}_M^2}}\{2C[2\mathrm{sin}(2x)+\mathrm{sin}(4x)]+`$
$`{\displaystyle \frac{V_M}{\sigma }}[16(C1)\mathrm{sin}(x)+2(3C5)\mathrm{sin}(2x)+`$
$`(C+1)\mathrm{sin}(4x)]\}.`$ (18)
We have divided the expression for $`v_3`$ into two parts: the first one does not contain the brane tension $`\sigma `$ and represents the contribution coming from the Einstein’s general relativity and the second one is connected with the brane-tension containing terms of the action (5). Finally, it remains to find the function $`u_3`$ to determine the amplitude of $`\mathrm{\Phi }`$. First of all we should realize that the shift of the right-end point (16) as well as the term $`k^2\mathrm{\Delta }_2`$ from the expansion of $`V_M^{\prime \prime }/\widehat{H}_M^2`$ enter the equation for $`u_3`$. This fact allows for determining a relation between $`k`$ and $`V_M^{\prime \prime }/H_M^2`$ as it is shown bellow. To keep the range of the argument of function $`u_3`$ equal to $`[0,\pi ]`$ we pass from the independent variable $`x`$ to $`w=(1k^2x_f^{(2)}/\pi )x`$. For the simplicity we introduce the notation
$`E=C+{\displaystyle \frac{159C}{4\sigma }}V_M,F=5C+{\displaystyle \frac{51C45}{4\sigma }}V_M,`$
$`G=3C+5(1+C){\displaystyle \frac{V_M}{4\sigma }}.`$
Within this notation we can derive straightforwardly differential equation for $`u_3`$:
$$u_3^{\prime \prime }+3\mathrm{cot}(w)u_3^{}+4u_3=A\mathrm{cos}(w)+B\mathrm{cos}^3(w),$$
(19)
where
$`A={\displaystyle \frac{5}{8}}C+{\displaystyle \frac{3}{8}}E+{\displaystyle \frac{3}{4}}G+{\displaystyle \frac{1}{24}}\left({\displaystyle \frac{V_M^{\prime \prime \prime }}{\widehat{H}_M^2}}\right)^2\mathrm{\Delta }_2,`$
$`B={\displaystyle \frac{3}{4}}G{\displaystyle \frac{1}{12}}\left({\displaystyle \frac{V_M^{\prime \prime \prime }}{\widehat{H}_M^2}}\right)^2+{\displaystyle \frac{1}{6}}{\displaystyle \frac{V_M^{\prime \prime \prime \prime }}{\widehat{H}_M^2}}.`$
The solution is: $`u_3=\beta \mathrm{cos}^3(w)`$, where the coefficient $`\beta `$ has to satisfy two conditions:
$$6\beta =A,14\beta =B.$$
However, the fixation of $`\beta `$ is, at the moment, only supplementary for us because the coefficient $`A`$ contains also $`\mathrm{\Delta }_2`$. Eliminating $`\beta `$ from the previous system of equations we obtain the expression for $`\mathrm{\Delta }_2`$ (we do not write down) and subsequently we get $`k^2`$ as the function of $`4V_M^{\prime \prime }/\widehat{H}_M^2`$:
$`k^2=7(4+{\displaystyle \frac{V_M^{\prime \prime }}{\widehat{H}_M^2}})\{16C+{\displaystyle \frac{1}{24}}\left({\displaystyle \frac{V_M^{\prime \prime \prime }}{\widehat{H}_M^2}}\right)^2+`$
$`{\displaystyle \frac{1}{2}}{\displaystyle \frac{V_M^{\prime \prime \prime \prime }}{\widehat{H}_M^2}}+{\displaystyle \frac{1}{32}}{\displaystyle \frac{V_M}{\sigma }}(43569C)\}^1.`$ (20)
Assuming $`4+V_M^{\prime \prime }/\widehat{H}_M^2<0`$ we need positive sign of the denominator in eq. (4). This sign can be changed, for given values of $`V_M`$ and $`\sigma `$, only due to $`V_M^{\prime \prime \prime \prime }/\widehat{H}_M^2`$. We introduce the critical value of $`V_M^{\prime \prime \prime \prime }/\widehat{H}_M^2`$
$$\zeta _c=\frac{1}{16}\frac{V_M}{\sigma }(69C435)32C\frac{1}{24}\left(\frac{V_M^{\prime \prime \prime }}{\widehat{H}_M^2}\right)^2$$
(21)
at which the mentioned denominator vanishes, and we have the first-order CdL instanton with inflaton amplitude given by the formula (4) in the two cases:
* for $`V_M^{\prime \prime \prime \prime }/\widehat{H}_M^2>\zeta _c`$ as $`4+V_M^{\prime \prime }/\widehat{H}_M^20^{}`$ ,
* for $`V_M^{\prime \prime \prime \prime }/\widehat{H}_M^2<\zeta _c`$ as $`4+V_M^{\prime \prime }/\widehat{H}_M^20^+`$ .
The first term in (21) is positive (if $`V_M`$ and $`\sigma `$ are positive). This cause the difference with respect to the situation when vacuum decays in four dimensional space-time within Einstein’s general relativity because now $`V_M^{\prime \prime \prime \prime }/\widehat{H}_M^2`$ can be both positive and less than $`\zeta _c`$ (In the first paper of refs. it is argued that the negative value of $`V_M^{\prime \prime \prime \prime }`$ is not like). Let us mention that the situation when $`V_M^{\prime \prime }/\widehat{H}_M^2>4`$ and $`V_M^{\prime \prime \prime \prime }/\widehat{H}_M^2<\zeta _c`$ does not mean automatically a kind of stabilization of the false vacuum because we still can have some CdL instanton with large amplitude in $`\mathrm{\Phi }`$ that is not included in our previous analysis and moreover the false vacuum can decay also via HM instanton. However, a consideration of such an instanton would involve some class of non-perturbative analysis of the system (3) or the numerical analysis, if we have a concrete potential $`V`$ and a brane tension $`\sigma `$.
## 5 The action of the first-order CdL instanton
The crucial quantity that tells us how probable is the vacuum decay via the CdL (or HM) instanton is its action. Our task is to find the approximative formula for the action of the first-order CdL instanton we investigated in previous section. In the authors have also considered the action of the CdL instanton. They were interested in the action within the thin-wall approximation that correspond with their example of the instanton. Following we can write the action
$`S=2\pi ^2{\displaystyle _0^{\tau _f}}\mathrm{d}\tau [a^3({\displaystyle \frac{1}{2}}(\mathrm{\Phi }^{})^2+V)+{\displaystyle \frac{a^3}{2\sigma }}({\displaystyle \frac{1}{2}}(\mathrm{\Phi }^{})^2+V)^2+`$
$`{\displaystyle \frac{1}{C}}(a^2a^{\prime \prime }+a(a^{})^2a)].`$
Using the equation of motion (3) we rewrite the action into a much simpler form
$`S={\displaystyle \frac{4\pi ^2}{3C}}{\displaystyle _0^{\tau _f}}ad\tau {\displaystyle \frac{\pi ^2}{\sigma }}{\displaystyle _0^{\tau _f}}a^3(\mathrm{\Phi }^{})^4d\tau =`$ (22)
$`{\displaystyle \frac{4\pi ^2}{3C}}S^{(I)}{\displaystyle \frac{\pi ^2}{\sigma }}S^{(II)}.`$ (23)
The structure of the $`S^{(I)}`$ term is the same as in the case of Einstein’s relativity theory. To avoid a confusion we must mention that brane-tension is included in this term throughout $`a`$. We can easily find the action of the HM instanton
$$S_{HM}=\frac{\pi }{\widehat{H}_M^2}=\frac{3}{8V_M}\frac{1}{1+\frac{V_M}{2\sigma }}.$$
(24)
We see that the action of the HM instanton in our brane world model is for fixed $`V_M`$ and any given (positive) $`\sigma `$ less than the action of corresponding HM instanton in Einstein’s relativity theory. Now we can compute the action of the first-order CdL instanton. Up to the second order in $`k`$ one has
$`S^{(I)}=\widehat{H}_M^2\left\{{\displaystyle _0^{\pi +k^2x_f^{(2)}}}v_0dx+k^2{\displaystyle _0^\pi }v_2dx\right\}=`$
$`\widehat{H}_M^2\left[2+k^2{\displaystyle \frac{4C5}{3}}{\displaystyle \frac{V_M}{\sigma }}\right]`$
and $`S^{(II)}`$ does not contribute because it is of the order $`k^4`$ at most. Putting these results together we obtain the difference between the actions of our first-order CdL instanton and related HM instanton
$`S_{CdL}S_{HM}={\displaystyle \frac{4\pi ^2}{3C}}{\displaystyle \frac{4C5}{3}}{\displaystyle \frac{1}{\widehat{H}_M^2}}{\displaystyle \frac{V_M}{\sigma }}k^2.`$ (25)
First of all: this difference vanishes as $`\sigma \mathrm{}`$ as it must be because in Einstein’s relativity theory the difference between the actions of our CdL instanton and related HM instanton is of the fourth-order in $`k`$ as it is shown in . For positive $`k^2`$ the difference (25) is negative (this coincides with the general-relativistic result ), this means that the vacuum decay via our CdL instanton is more probable than the decay via HM instanton. The fact that the difference of the actions (25) is of the $`k^2`$ order unlike the general-relativistic case where it is of the order $`k^4`$ means that the false vacuum decay rate in our brane world model is higher than for the conventional gravity. The same is shown for the false vacuum decay rate in another extremal case when the thin-wall approximation can be used .
## 6 Summary
We have compared, from the point of view of the semiclassical description of the false vacuum decay, the systems of differential equations (3) and (3) defining the CdL instantons in the standard Einstein’s relativity theory and in the brane world model. We have been interested in the situation when the effective curvature of the potential at its top is close to the critical value $`4`$. In this situation we were able to obtain explicit perturbative formulas for the instanton, and mainly for its actions that is the relevant quantity describing the instanton as ”mediator” of the vacuum decay. This allows for comparing our CdL instanton with always existing Hawking - Moss instanton. We have concluded that it is the first order CdL instanton that is preferred (to the corresponding Hawking - Moss instanton) in the considered region of parameters of the potential. A kind of lower symmetry of the instanton equations with respect to the situation in general relativity causes that the difference between actions of mentioned instantons is proportional to the square of inflaton amplitude in the instanton rather than to the fourth power of that amplitude. Physically, this difference of the actions is inversely proportional to the string tension ($`V_M/\sigma `$ is the parameter controlling the string tension influence) and therefore in the limit $`V_M/\sigma 0`$ we recover the results of the general relativity.
### Acknowledgement
This work was supported by the Slovak Scientific and Educational Grant Agency, project no. 1/0250/03. |
warning/0506/quant-ph0506019.html | ar5iv | text | # Enhanced algorithms for Local Search
## 1 Introduction
The Local Search problem consists in finding a local minimum of the function $`f`$ on $`G`$, that is a vertex $`v`$ such that $`f(v)`$ is not larger than the value of $`f`$ on the neighbors of $`v`$ in $`G`$. Obviously, such a vertex always exists, as a global minimum satisfies this constraint. Another easy argument shows how to find such a vertex : make a walk over vertices such that at each step the next vertex is the neighbor of the current vertex which has the smallest value; the walk will stop in a local minimum. Such a walk is called a steepest descent. Steepest descents are the basis of several approaches to efficiently find a local minimum.
The Local Search problem has been previously studied and there is already a large literature on its complexity. Its structural complexity, where the function and the graph are given as an input to a Turing machine, was studied in , and its query complexity, where the graph is known but the values of $`f`$ are accessed through an oracle, was investigated in . We focus on the query complexity, which complexity is obviously at most the size of the graph. Our query model is the standard one; see Section 3.2 for precise definitions.
The deterministic query complexity of Local Search on a graph $`G`$ of size $`n`$, maximum degree $`d`$ and genus $`g`$ (for a definition of the genus, see for instance ), is intimately connected to the size of separators of $`G`$:
###### Definition 1.
A separator for $`G`$ is a subset of $`V`$ whose removal leaves no connected component with more than $`2n/3`$ vertices.
In , a deterministic query algorithm was exhibited, which works using a sub-linear number of queries for large classes of graphs: a local minimum can be found on the graph $`G`$ using $`O(\mathrm{log}n)d+O(g)\sqrt{n}`$ queries. It is based on the recursive use of separators for smaller and smaller subgraphs of $`G`$, and their complexity analysis relies on the following result:
###### Theorem 1 (Gilbert, Hutchinson, Tarjan ).
The graph $`G`$ has a separator of size at most $`6\sqrt{gn}+2\sqrt{2n}+1`$.
In the randomized and quantum query models, the situation is quite different, as the size of separators of $`G`$. Also, the only known sub-linear query algorithm for general graphs is a randomized algorithm, that we call Randomized Steepest Descent, was exhibited by Aldous , and has a query complexity $`\mathrm{\Theta }(\sqrt{nd})`$. The idea of this algorithm is to choose $`\sqrt{nd}`$ vertices at random, query their values, start a steepest descent from the vertex with smallest value for at most $`\sqrt{nd}`$ steps, and to return the last visited vertex. This idea was later refined by Aaronson to give a sub-linear quantum query algorithm for general graphs, that we call Quantum Steepest Descent, using $`\mathrm{\Theta }(n^{1/3}d^{1/6})`$ queries.
On the side of lower-bounds, it follows from that the size of a smallest separator is a lower-bound on the deterministic complexity of Local Search. Also, from , we know that $`d`$ is a lower-bound on the deterministic query complexity, $`\mathrm{\Omega }(d)`$ a lower-bound for the randomized query complexity, and $`\mathrm{\Omega }(\sqrt{d})`$ a lower-bound for the quantum query complexity.
## 2 Results
In this note, we first improve Theorem 1 in Section 4.1, to obtain the following slightly stronger separation theorem:
###### Theorem 2 (strong separation for graphs of genus $`g`$).
Assume $`n3`$. There exists a separator $`C`$ for $`G`$ such that $`C`$ contains no more than $`(6+2\sqrt{2}12/n+6\sqrt{g}+4g/n+1/\sqrt{n})\sqrt{n}`$ vertices, and the subgraph induced on $`G`$ by $`VC`$ has maximal degree at most $`\sqrt{n}`$.
As a result it allows us to enhance, in Section 5, the deterministic algorithm of Llewellyn, Tovey and Trick whose complexity is of $`O(\mathrm{log}n)d+O(\sqrt{g})\sqrt{n}`$. We also derive a quantum algorithm from it. More precisely, we obtain the following result:
###### Theorem 3.
There exists a deterministic and a quantum query algorithms that find a local minimum of $`f`$ on $`G`$ using respectively $`d+O(\sqrt{g})\sqrt{n}`$ and $`O(\sqrt{d})+O(\sqrt[4]{g})\sqrt[4]{n}\mathrm{log}\mathrm{log}n`$ queries.
Our deterministic and quantum algorithms have smaller query complexities than the respective algorithms Randomized Steepest Descent of Aldous and of Quantum Steepest Descent of Aaronson for large classes of graphs, including graphs of bounded genus and planar graphs. We analyze this in detail in Section 6.
Independently from this work, Zhang has recently given an algorithm which finds a local minimum on the planar grid over $`\{1,\mathrm{},\sqrt{n}\}^2`$ using $`O(\sqrt[4]{n}(\mathrm{log}\mathrm{log}n)^2)`$ queries. Our quantum algorithm can be viewed as a strongly generalized, and slightly enhanced version of this algorithm.
## 3 Preliminaries
### 3.1 Notations
We denote by $`\mathrm{log}n`$ the natural logarithm of $`n`$, and for every positive real number $`b`$ we denote by $`\mathrm{log}_bn`$ the logarithm of $`n`$ in base $`b`$. If $`G`$ is a graph and $`v`$ is any vertex of $`G`$, we denote by $`_G(v)`$ the set of neighbors of $`v`$ in $`G`$.
### 3.2 Query complexity
In the query model of computation we count only queries made by the algorithm, but all other computations are free. The state of the computation is represented by three registers, the query register $`i\{1,\mathrm{},n\}`$, the answer register $`a\mathrm{\Sigma }`$, and the work register $`zW`$, where $`\mathrm{\Sigma }`$ and $`W`$ are finite sets. The computation takes place in the vector space spanned by all basis states $`|i|a|z`$. In the quantum query model introduced by Beals, Buhrman, Cleve, Mosca and de Wolf the state of the computation is a complex combination of all basis states which has unit length for the norm $`\mathrm{}_2`$, and the allowed operations on the state of the computation are all isometric operators for the $`\mathrm{}_2`$ norm acting over the computation space. In the randomized model, the state of the computation is a non-negative real combination of all basis states of unit length for the norm $`\mathrm{}_1`$, and the allowed operations on the state of the computation are all isometric operators for the norm $`\mathrm{}_1`$ acting over the computation space. In the deterministic model, the state of the computation is always one of the basis states, and the allowed operations are all operators mapping a basis state to another basis state.
Assume that $`x\mathrm{\Sigma }^n`$ is the input of the problem which can be accessed only through the oracle. The query operation $`𝒪_x`$ is the permutation which maps the basis state $`|i|a|z`$ into the state $`|i|(a+x_i)mod|\mathrm{\Sigma }||z`$ (here we identify $`\mathrm{\Sigma }`$ with the residue classes $`mod|\mathrm{\Sigma }|`$). Non-query operations are independent of $`x`$. A $`k`$-query algorithm is a sequence of $`(k+1)`$ operations $`(U_0,U_1,\mathrm{},U_k)`$ where $`U_i`$ is an allowed operation in the chosen model of computation. Initially the state of the computation is set to some fixed value $`|0|0|0`$, and then the sequence of operations $`U_0,𝒪_x,U_1,𝒪_x,\mathrm{},U_{k1},𝒪_x,U_k`$ is applied. The final state is denoted by $`\mathrm{\Phi }`$.
The output in the quantum model is an element $`zW`$ that appears with probability equal to the square of the $`\mathrm{}_2`$ norm of the orthogonal projection of $`\mathrm{\Phi }`$ over the vector space $`V`$ spanned by $`\{|i|a|z|i\{1,\mathrm{},n\},a\mathrm{\Sigma }\}`$. The output in the randomized model is an element $`zW`$ that appears with probability equal to the $`\mathrm{}_1`$ norm of the orthogonal projection of $`\mathrm{\Phi }`$ over the vector space $`V`$. The output in the deterministic model is the element $`zW`$ such that there exist $`i`$ and $`a`$ with $`\mathrm{\Phi }=|i|a|z`$.
Assume that $`R\mathrm{\Sigma }^n\times W`$ is a total relation (i.e. for every $`x\mathrm{\Sigma }^n`$ there exists $`zW`$ such that $`(x,z)R`$) that we want to compute. A quantum or randomized algorithm computes (with two-sided error) $`R`$ if its output yield some $`zW`$ such that $`(x,z)R`$ with probability at least $`2/3`$. A deterministic algorithm computes $`R`$ if its output yield some $`zW`$ such that $`(x,z)R`$.
Then the query complexity of a relation $`R`$ in a model of computation (deterministic, randomized or quantum) is the smallest $`k`$ for which there exists a $`k`$-query algorithm, in that model of computation, which computes $`R`$.
## 4 Tools
In this section, we recall and prove the results that we need in order to design the algorithms of Section 5.
### 4.1 Separation in graphs of higher genus
We first recall the following well-known theorem for graphs of higher genus (see for instance ) :
###### Theorem 4.
Any $`n`$-vertex graph of genus $`g`$ with $`n3`$ contains no more than $`3n6+2g`$ edges.
This result, together with Theorem 1, allows us to prove Theorem 2.
###### Proof of Theorem 2.
Theorem 1 shows that it is possible to find a separator $`C^{}`$ for $`G`$ such that $`C^{}`$ has size at most $`6\sqrt{gn}+2\sqrt{2n}+1`$. Let $`B`$ be the set of all vertices of degree greater than $`\sqrt{n}`$. Using Theorem 4, we have
$$|B|\sqrt{n}\underset{vV}{}\mathrm{d}(v)=2|E|6n12+4g.$$
The set $`C=C^{}B`$ being a superset of a separator of $`G`$ is also a separator of $`G`$. From the definition of $`B`$, the subgraph induced on $`G`$ by $`VC`$ obviously has maximal degree at most $`\sqrt{n}`$. ∎
The genus of a graph being in $`O(n^2)`$, the asymptotic inequality $`g/n=O(\sqrt{g})`$ holds and therefore Theorem 2 can be interpreted as stating the existence of a particular $`O(\sqrt{g})\sqrt{n}`$ separator $`G`$.
### 4.2 Minimum-finding algorithms
In this paragraph, we recall results about the query complexity of finding the minimum value of a function on a set.
Let $`n`$ be a positive integer, $`S`$ be a set of cardinality $`n`$, $`g:S`$ be a function and $`𝒪_g`$ be an oracle for $`g`$.
###### Definition 2.
Let $`\epsilon <1`$ be any positive real number. If $`𝒜`$ is a randomized algorithm that outputs the minimum value of the function $`g`$ on $`S`$ with probability at least $`1\epsilon >0`$, then we denote by $`\mathrm{argmin}_𝒜\{g(s)|sS\}`$ the random variable equal to its output.
It is obvious that, for every deterministic algorithm $`𝒜`$, computing $`\mathrm{argmin}_𝒜\{g(s)|sS\}`$ requires querying all $`n`$ values of $`g`$ to $`𝒪_g`$. It is natural that, for every randomized algorithm $`𝒜`$, computing $`\mathrm{argmin}_𝒜\{g(s)|sS\}`$ requires querying $`\mathrm{\Omega }(n)`$ values of $`g`$ to $`𝒪_g`$. It is more surprising that much less queries are needed when quantum queries are allowed:
###### Theorem 5 (Dürr, Høyer ).
There exists a quantum algorithm which finds the minimum value of $`g`$ with probability at least $`1/2`$, using $`O(\sqrt{n})`$ quantum queries to $`𝒪_g`$.
Amplification of the probability of success of the algorithm of Dürr and Høyer can be obtained by running the algorithm several times, and then taking the minimum value of all the values that have been returned by each repetition of the algorithm. After $`k`$ repetitions, the probability of having found the minimum is at least $`12^k`$. In particular, for every positive real number $`\epsilon <1`$, there exists a quantum algorithm $`𝒜`$ computing $`\mathrm{argmin}_𝒜\{g(s)|sS\}`$ with probability at least $`1\epsilon `$ using $`O(\sqrt{n}\mathrm{log}(1/\epsilon ))`$ quantum queries.
## 5 Enhanced algorithms for Local Search
In this paragraph, we prove Theorem 3. The proof will be in two steps: in Theorem 6 we will prove the correction of our algorithms, and in Theorem 7 we will prove their complexity.
The basic procedure of our algorithms follows the lines of the algorithms of Llewellyn, Tovey and Trick of and Santha and Szegedy . It is given in Algorithm 1. The main idea is to adopt a divide-and-conquer approach: the graph is split into connected components of small size by removing a separator; then, querying the values of the vertices in and close to that separator make it possible to find one of these connected components in which there is a local minimum of $`f`$ on $`G`$.
Notice that neither the way the separators are chosen, or how the minimum-finding algorithms $`𝒜_i`$ work for integers $`i1`$, are specified in Algorithm 1. Our algorithms consist in using the procedure described in Algorithm 1 with the following specific choices:
* a separator $`C`$ of a graph $`G^{}`$ will be chosen according to Theorem 2 if $`G^{}`$ has more than two vertices, and $`C`$ contains all vertices of $`G^{}`$ otherwise,
* the minimum-finding algorithm $`𝒜_i`$ will behave as follows when requested to minimize the function $`f`$ over a set $`SV`$ of vertices: in the deterministic case, the local minimum of $`f`$ is always found by exhaustive search<sup>1</sup><sup>1</sup>1One should observe that it is important at this point that $`f`$ takes distinct values on distinct vertices. This can be assumed, as one could for instance put a total order $``$ on $`V`$, and minimize the function $`g:v(f(v),v)`$ according to the lexicographic order induced on $`\times V`$ by $`<`$ and $``$, instead of minimizing $`f`$. The function $`g`$ takes distinct values on distinct vertices.. In the quantum cases, if the set $`S`$ has size at most $`3`$, then the minimum value of $`f`$ on $`S`$ is also found using exhaustive search. Otherwise, the output is the one found by the last measurement at the end of the quantum procedure described in paragraph 4.2; moreover, we request that the minimum-finding algorithm $`𝒜_1`$ has error probability $`1/12`$ using $`O(\sqrt{|S|})`$ queries, and $`𝒜_i`$ has error probability $`1/12\mathrm{log}_{3/2}n`$ using $`O(\sqrt{|S|}\mathrm{log}\mathrm{log}n)`$ queries, for $`i1`$.
###### Theorem 6.
With our choice of minimum-finding algorithm, Algorithm 1 always returns a local minimum in the deterministic case, and returns a local minimum with probability at least $`2/3`$ in the quantum case.
###### Proof.
Let $`j`$ be the largest value of the variable $`i`$ for a run of the algorithm. First, an easy inductions shows that for every iteration $`ij`$ of the main loop, and every $`vV^{(i)}`$ we have
$$_G(v)_{G^{(i)}}(v)C^{(1)}C^{(2)}\mathrm{}C^{(i)}.$$
So, to prove that $`f`$ is minimized on $`v^{(j)}`$, one must only prove that $`f(v^{(j)})`$ is not larger than $`\mathrm{min}\{f(v)|v_{G^{(j)}}(v)C^{(1)}C^{(2)}\mathrm{}C^{(j)}\}`$.
If, during the run of the algorithm, the calls to the algorithms $`𝒜_i`$, for $`1ij`$, have always successfully returned elements minimizing $`f`$, then for every positive integer $`ij`$ we have $`f(v^{(i)})\mathrm{min}\{f(v^{(i1)}),f(m^{(i)}),f(z^{(i)})\}`$. Therefore, an easy induction shows that $`f(v^{(i)})\mathrm{min}\{f(v)|vC^{(k)}\}`$, for every positive integers $`kij`$. Moreover, the equality $`v^{(j)}=m^{(j)}`$ implies $`f(v^{(j)})\mathrm{min}\{f(v)|v_{G^{(j)}}(v^{(j)})\}`$. So, if $`𝒜_i`$ never failed to find a minimizing element, then the criterion given in the previous paragraph shows that $`v^{(j)}`$ is a local minimum.
In the deterministic case, the algorithms $`𝒜_i`$, for $`1ij`$, always return an element minimizing $`f`$, and therefore Algorithm 1 always returns a local minimum.
In the quantum case, a call to $`𝒜_i`$ returns an element minimizing $`f`$ with error probability at most $`1/12`$ for $`i=1`$, and at most $`11/4\mathrm{log}_{3/2}n`$ for $`1<ij`$. The set $`C^{(i)}`$ being a separator of $`G^{(i1)}`$ for every positive integer $`ij`$, we have $`|V^{(i)}|2|V^{(i1)}|/3`$. This implies that $`j\mathrm{log}_{3/2}n`$, and the probability that $`𝒜_i`$ did not return an element minimizing $`f`$ at some point is at most $`21/12+2\mathrm{log}_{3/2}n/(12\mathrm{log}_{3/2}n)=1/3`$. So, in the quantum case, Algorithm 1 returns a local minimum with probability at least $`2/3`$. ∎
###### Theorem 7.
With our choices of separators, Algorithm 1 has a deterministic query complexity at most $`d+O(\sqrt{g})\sqrt{n}`$, and a quantum query complexity at most $`O(\sqrt{d})+O(\sqrt[4]{g})\sqrt[4]{n}\mathrm{log}\mathrm{log}n`$.
###### Proof.
Again, let $`j`$ be the largest value of the variable $`i`$ for a run of the algorithm. Let us denote by $`𝒞_{𝒜_i}(s)`$ the number of queries made by the minimum-finding algorithm $`𝒜_i`$ on a set of size $`s`$, and by $`L^i(n,d)`$ the number of queries that are made in the $`i`$-th iteration of the main loop of our algorithm on a graph $`G^{}`$ that has $`n`$ vertices and is of maximum degree $`d`$. We denote also by $`d_i`$ the maximum degree of $`|G^{(i)}|`$, for a non-negative integer $`ij`$. Analysis of the main loop of Algorithm 1 gives, for every positive integer $`ij`$,
$$L^i(n,d)𝒞_{𝒜_i}(|C^{(i)}|)+𝒞_{𝒜_i}(d_{i1})+3.$$
Let us denote by $`T_\gamma ^i(\alpha ,\beta )`$ the number of queries made by our algorithm in the main loop between its $`i`$-th iteration and the end of the algorithm if $`i<j`$ and $`0`$ if $`ij`$, on an input graph which has $`\alpha `$ vertices, is of maximum degree $`\beta `$ and has genus at most $`\gamma `$. Theorem 2 ensures that for every positive integer $`ij`$ we have $`|V^{(i)}|2|V^{(i1)}|/3`$, and $`d_i\sqrt{|V^{(i1)}|}`$. Moreover, the genus of be $`G^{(i)}`$ is not larger than the genus of $`G^{(i1)}`$. So, by induction we have $`|V^{(i)}|(2/3)^in`$, and the genus of $`|G^{(i)}|`$ is at most $`g`$. Therefore, for every integer $`1ij`$ we have $`|C^{(i)}|O(\sqrt{g})\sqrt{(2/3)^{i1}n}`$, $`d_0=d`$ and $`d_i\sqrt{(2/3)^{i1}n}`$. This leads to the following equations:
$`T_g^1(n,d)`$
$`L^1(n,d)+T_g^2(n,d)`$
$`𝒞_{𝒜_1}(O(\sqrt{g})\sqrt{n})+𝒞_{𝒜_1}(d)+3+T_g^2(n,d),`$
and for every $`i\{2,\mathrm{},\mathrm{log}_{3/2}n1\}`$,
$`T_g^i(n,d)`$
$`L^i(n,d)+T_g^{i+1}(n,d)`$
$`𝒞_{𝒜_i}\left(O(\sqrt{g})\sqrt{(2/3)^{i1}n}\right)+`$
$`𝒞_{𝒜_i}\left(\sqrt{(2/3)^{i2}n}\right)+3+T_g^{i+1}(n,d).`$
In the deterministic case we have $`𝒞_{𝒜_i}(k)=k`$ for all positive integers $`k`$ and $`i`$, and in the quantum case we have, for all positive integer $`k`$, $`𝒞_{𝒜_i}(k)=O(\sqrt{k})`$ when $`i=1`$, and $`𝒞_{𝒜_i}(k)=O(\sqrt{k}\mathrm{log}\mathrm{log}n)`$ when $`i1`$. So, in the deterministic case, summing all the previous inequalities gives
$$\begin{array}{c}T_g^1(n,d)O(\sqrt{g})\sqrt{n}\underset{i=0}{\overset{\mathrm{}}{}}\sqrt{2/3}^i+d+\hfill \\ \hfill \sqrt{n}\underset{i=0}{\overset{\mathrm{}}{}}\sqrt{2/3}^i+3\mathrm{log}_{3/2}n+\\ \hfill T_g^{\mathrm{log}_{3/2}n}(n,d),\end{array}$$
which shows $`T_g^1(n,d)=d+O(\sqrt{g})\sqrt{n}`$, as $`T_g^{\mathrm{log}_{3/2}n}(n,d)=O(1)`$, and the query complexity of our deterministic algorithm is $`T_g^1(n,d)=d+O(\sqrt{g})\sqrt{n}`$. In the quantum case, it gives
$$\begin{array}{c}T_g^1(n,d)O(\sqrt[4]{g})\sqrt[4]{n}\mathrm{log}\mathrm{log}n\underset{i=0}{\overset{\mathrm{}}{}}\sqrt[4]{2/3}^i+\hfill \\ \hfill O(\sqrt{d})+O(\sqrt[4]{n}\mathrm{log}\mathrm{log}n)\underset{i=0}{\overset{\mathrm{}}{}}\sqrt[4]{2/3}^i+\\ \hfill 3\mathrm{log}_{3/2}n+T_g^{\mathrm{log}_{3/2}n}(n,d),\end{array}$$
leading to a quantum query complexity $`T_g^1(n,d)=O(\sqrt{d})+O(\sqrt[4]{g})\sqrt[4]{n}\mathrm{log}\mathrm{log}n`$. ∎
## 6 Comparison with generic algorithms
Let us first compare the query complexity of our deterministic algorithm with the query complexity of the algorithm Randomized Steepest Descent of Aldous . The complexity of our algorithm is $`d+O(\sqrt{g})\sqrt{n}`$, and the complexity of Randomized Steepest Descent is $`\mathrm{\Theta }(\sqrt{nd})`$. As $`dn`$, our algorithm performs as well as Randomized Steepest Descent (up to a constant speedup factor) as soon as $`g=O(d)`$, and performs asymptotically better when $`g=o(d)`$. In particular, our deterministic algorithm has lower query complexity than Randomized Steepest Descent on classes of graphs with bounded genus, which includes the class of planar graphs.
Let us now compare the query complexity of our quantum algorithm with the query complexity of the algorithm Quantum Steepest Descent of Aaronson . The complexity of our algorithm is $`O(\sqrt{d})+O(\sqrt[4]{g})\sqrt[4]{n}\mathrm{log}\mathrm{log}n`$, and the complexity of Quantum Steepest Descent is $`\mathrm{\Theta }(n^{1/3}d^{1/6})`$. As $`dn`$, we have $`\sqrt{d}n^{1/3}d^{1/6}`$, and our algorithm performs as well as Quantum Steepest Descent (up to a constant speedup factor) as soon as $`g^{1/2}n^{1/4}\mathrm{log}\mathrm{log}n=O(n^{1/3}d^{1/6})`$, that is to say $`g=O(n^{1/6}d^{1/3}/(\mathrm{log}\mathrm{log}n)^2)`$. This holds if $`g=O(\sqrt{d}/(\mathrm{log}\mathrm{log}n)^2)`$. Also, our quantum algorithm performs asymptotically better when $`g=o(n^{1/6}d^{1/3}/(\mathrm{log}\mathrm{log}n)^2)`$, and therefore when $`g=o(\sqrt{d}/(\mathrm{log}\mathrm{log}n)^2)`$. In particular, our quantum algorithm has lower query complexity than Quantum Steepest Descent on classes of graphs with bounded genus, which includes the class of planar graphs.
In conclusion, the algorithms we have designed perform better than the known generic algorithms for some classes of graphs, in particular planar graphs and graphs of constant genus, both for classical (deterministic and randomized) computation, and for quantum computation. |
warning/0506/cs0506097.html | ar5iv | text | # A Flexible Thread Scheduler for Hierarchical Multiprocessor Machines
## 1 Introduction
“Disable HyperThreading!” That is unfortunately the most common pragmatic answer to performance losses noticed on HyperThreading-capable processors such as the Intel Xeon. This is of particular concern since hierarchy depth has increased over the past few years, making current computer architectures more and more complex (Sun WildFire , Sgi Origin , Bull NovaScale for instance).
Those machines look like Russian dolls: nested technologies allow them to execute several threads at the same time on the same core of one processor (SMT: Simultaneous Multi-Threading), to share cache memory between several cores (multicore chips), and finally to interconnect several multiprocessor boards (SMP) thanks to crossbar networks. The resulting machine is a NUMA (Non-Uniform Memory Architecture) computer, on which the memory access delay depends on the relative positions of processors and memory banks (this is called the “NUMA factor”).
The recent integration of SMT and multicore technologies make the structure of NUMA machines even more complex, yet operating systems still have not exploited previous NUMA machines efficiently. Hennessy and Patterson underlined that fact about systems proposed for SGI Origin and Sun Wildfire: *“There is a long history of software lagging behind on massively parallel processors, possibly because the software problems are much harder.”* The introduction of new hardware technologies emphasizes the need for software development. Our goal is to provide a *portable* solution to enhance the efficiency of high-performance multi-threaded applications on modern computers.
Obtaining optimal performance on such machines is a significant challenge. Indeed, without any information on tasks’ affinity, it is difficult to make good decisions about how to group tasks working on a common data set on NUMA nodes. Detecting such affinity is hard, unless the application itself somehow expresses it.
To relieve programmers from the burden of redesigning the whole task scheduling mechanism for each target machine, we propose to establish a communication between the execution environment and the application so as to automatically get an optimized schedule. The application describes the organization of its tasks by grouping those that work on the same data (memory affinity) for instance. The system scheduler can then exploit this information by adapting the task distribution to the hierarchical levels of the machine.
Of course, a universal scheduler that would get good results by using only such a small amount of information remains to be written. In the meantime, we provide facilities for applications to query the system about the topology of the underlying architecture and “drive” the scheduler. As a result, the programmer can easily try and evaluate different gathering strategies. More than a mere scheduling model, we propose a scheduling experimentation platform.
In this article, we first present the main existing approaches that exploit hierarchical machines, then we propose two new models describing application tasks and hierarchical levels of the machine, as well as a scheduler that takes advantage of them. Some implementation details and evaluation results are given before concluding.
## 2 Exploiting hierarchical <br>machines
Nowadays, multiprocessor machines like NUMAs with multi-threaded multicores are increasingly difficult to exploit. Several approaches have been considered.
### 2.1 Predetermined distribution and <br>scheduling
For very regular problems, it is possible to determine a task schedule and a data distribution that are suited to the target machine and its hierarchical levels. The application just needs to get the system to apply that schedule and that distribution, and excellent (if not optimal) performance can be obtained. The PaStiX large sparse linear systems solver is a good example of this approach. It first launches a simulation of the computation based on models of BLAS operators and communications on the target architecture. Then it can compute a static schedule of block-computations and communications.
So as to enforce these scheduling strategies, many systems (Aix, Linux, Solaris, Windows, …) allow process threads to be bound to processor sets, and memory allocations to be bound to memory nodes. Provided that the machine is dedicated to the application, the thread scheduling can be fully controlled by binding exactly one thread to each processor. To perform task switching, mere explicit context switches may be used: threads are only used as execution flow holders.
### 2.2 Opportunist distribution and scheduling
Greedy algorithms (called Self-Scheduling (SS) ) are dynamic, flexible and portable solutions for loop parallelization. Whatever the target machine, a Self-Scheduling algorithm takes care of both thread scheduling and data distribution. Operating systems schedulers are based on these algorithms.
They basically use a single list of ready tasks from which the scheduler just picks up the next thread to be scheduled. Hence the workload is automatically distributed between processors. For each task, the last processor on which it was scheduled is recorded, so as to try to reschedule it on the same processor as much as possible to avoid cache misses. These techniques are used in the Linux 2.4 and Windows 2000 operating systems. However, a unique thread list for the whole machine is a bottleneck, particularly when the machine has many processors.
To avoid such contention, Guided Self-Scheduling (GSS) and Trapezoid Self-Scheduling (TSS) algorithms make each processor take a whole part of the total work when they are idle, raising the risk of imbalances. AFfinity Scheduling (AFS) and Locality-based Dynamic Scheduling (LDS) algorithms use a per-processor task list. Whenever idle, a processor will steal work from the least loaded list, for instance. These latter algorithms are used by current operating systems (Linux 2.6 , FreeBSD 5.0 , Cellular Irix ). They also add a few rebalance policies: new processes are charged to the least loaded processor, for instance.
However, contention appears quickly with an increased number of processors, particularly on NUMA machines. Wang *et al*. propose a Clustered AFfinity Scheduling (CAFS) algorithm which groups $`p`$ processors in groups of $`\sqrt{p}`$. Whenever idle, rather than looking around the whole machine, processors steal work from the least loaded processor of their group, hence getting better localization of list accesses. Moreover, by aligning groups to NUMA nodes, data distribution is also localized. Finally, the Hierarchical AFfinity Scheduling (HAFS) (Wang *et al*. ) algorithm lets any idle group steal work from the most loaded group. This latter approach is being considered for latest NUMA-aware developments of operating systems such as Linux 2.6 and FreeBSD.
### 2.3 Negotiated distribution and scheduling
There are intermediate solutions between predetermined and opportunist scheduling. Some language extensions such as OpenMP , HPF (*High Performance Fortran*) or UPC (*Unified Parallel C*) let one achieve parallel programming by simply annotating the source code. For instance, a for loop may be annotated to be automatically parallelized. An HPF matrix may be annotated to be automatically split into rather independent domains that will be processed in parallel.
The distribution and scheduling decisions then belong to the compiler. To do this, it adds code to query the execution environment (the number of processors for instance) and compiles the program in a way generic enough to adapt to the different parallel architectures. In particular, it will have to handle threads for parallelized loops or distributed computing, and even handle data exchange between processors (in the case of distributed matrices of HPF). To date, expressiveness is limited mostly to “Fork-Join” parallelism, which means, for instance, that the programmer can not express imbalanced parallelism.
Programmers may also directly write applications that are able to adapt themselves to the target machine at runtime. Modern operating systems provide full information about the architecture of the machine (user-level libraries are available: lgroup for Solaris or numa for Linux). The application can then not only get the number of processors, but also get the NUMA nodes hierarchy, their respective number of processors and their memory sizes. Those systems also let the application choose the memory allocation policy (specific memory node, *first touch* or *round robin*) and bind threads to CPU sets. Thus, the application controls threads and memory distribution, but it is then in charge of balancing threads between processors.
### 2.4 Discussion
We chose to classify existing approaches into three categories. The *predetermined* category gives excellent performance. But it is portable only if the problem is regular, *i.e.,* its solving time depends on the data structure and not on the data itself. The *opportunist* approach scales well, but does not take task affinities into account, and thus, on average, does not get excellent performance. The *negotiated* approach lets the application adapt itself to the underlying machine, but requires rewriting of some parts of the scheduler in order to be flexible.
Our proposal is a mix between negotiated and opportunist approaches. We will give the programmers means to dynamically describe how their applications behave, and use this information to guide a generic opportunist scheduler.
## 3 Proposal: an application-guided scheduler
Our proposal is based on a collaboration between the application and its execution environment.
### 3.1 Bubbles modeling the application <br>structure
The application is asked to model the general layout of its threads in terms of nested sets called bubbles<sup>1</sup><sup>1</sup>1In a way relatively similar to some communication libraries such as MPI, which ask the application to specify *communicators*: groups of machines which will communicate..
Figure 1 shows such a model: the application groups threads into pairs, along with a communication thread (priorities will be discussed later). The concept of bubbles can be understood as a *coset with respect to a specific affinity relation*, and bubble nesting expresses refinement of a relation by another one. Indeed, several affinity relations can be considered, for instance:
It is a good idea to group threads that work on the same data so as to benefit from cache effects, or at least to avoid spreading the data throughout the NUMA nodes thereby incurring the NUMA factor penalty.
It can be beneficial to optimize the scheduling of threads which are to perform collective operations such as a synchronization barrier, which ensures that all involved threads have reached the barrier before they can continue executing.
Many attempts were made to address thread scheduling on Simultaneous Multi-Threading (SMT) processors, mostly by detecting affinities between threads at runtime . Indeed, in some cases, pairs of threads may be able to efficiently exploit the SMT technology: they can run in parallel on the logical processors of the same physical processor without interfering. If the programmer knows that some pairs of threads can work in such *symbiosis*, he can express this relation.
Other relations may be possible to express parallelism, sequentiality, preemption, etc. Yet, blindly expressing these relations may also be detrimental: Bulpin and Pratt show performance loss on SMT processors due to frequent cache misses for instance; Antonopoulos *et al.* also show performance loss when not taking the SMP bus bandwidth limit into account. But the programmer may try and test different refinements of the relations and thus experimentally reveal how the threads of an application should be related.
In order to cope with the emerging multiprocessor networks of the 1980’s, Ousterhout proposed to group data and threads by affinity into *gangs*. These gangs hold a fixed number of threads which are to be launched at the same time on the same machine of the network: this is called *Gang Scheduling*. However, processors may be left idle because a single machine can only run one gang at a time, even if it is “small”. Feitelson *et al*. propose a hierarchical control of the processors so as to execute several gangs on the same machine. Our approach is actually a generalization of this approach.
### 3.2 Task lists modeling the computing power structure
According to Dandamudi and Cheng , a hierarchy of task lists generally brings better performance than simple per-processor lists. This is why two-level list schedulers have been developed . Moreover, it makes task binding to processor sets easier. In a manner similar to Nikolopoulos et al.’s Nano-Threads list hierarchy , we have taken up and generalized this point of view.
Indeed, we model hierarchical machines by a hierarchy of task lists. Each component of each level of the hierarchy of the machine has one and only one task list. Figure 2 shows a hierarchical machine and its model. The whole machine, each NUMA node, each core, each physical SMT processor and each logical SMT processor has a task list.
For a given task, the list on which it is inserted expresses the scheduling area: if the task is on a list associated with a physical chip, it will be allowed to be run by any processor on this chip; if it is placed on the global list, it will be allowed to be run by any processor of the machine.
### 3.3 Putting both models together: a bubble scheduler
Once the application has created bubbles, threads and bubbles are just “tasks” that the execution environment distributes on the machine.
#### 3.3.1 Bubble evolution
As Figure 3 shows, the goal of a bubble is to hold tasks and bring them to the level where their scheduling will be most efficient. For this, the bubble goes down through lists to the *wanted* hierarchical level. It then “bursts”, *i.e.* held threads and bubbles are released and can be executed (or go deeper). The list of held tasks is recorded, for a potential later regeneration (see Section 3.3.3). The main issue is how to specify the right bursting level of a bubble.
In the long run, once we get good heuristics for a bubble scheduler, specifying such a parameter will no longer be necessary. For now, the goal is to provide an experimental platform for developing schedulers, and hence allow this parameter to be tuned by the scheduler developers. They can favor task affinity with the risk of making the load balance difficult (by setting deep bursting levels) or on the contrary favor processor use (by setting high bursting levels).
#### 3.3.2 Priorities
We choose to let the application attach integer priorities to tasks. When a processor looks for a task to be scheduled, it searches through the lists that “cover” this processor, from the most *local* one (*i.e.* on low levels) to the most *global* one, looking for a task with highest priority. It will then schedule that task, even if less prioritized tasks remain on more local lists.
Figure 1 shows an example using priorities. In this example, bubbles holding computing threads are *less* prioritized than the threads. Consequently, a bubble will burst only if every thread of the previously burst bubbles has terminated, or if there are not enough of them to occupy all the processors. This results in some *Gang scheduling* which automatically occupies all the processors.
#### 3.3.3 Bubble regeneration
Bubbles are automatically distributed by the scheduler over the different levels of task lists of the machine, hence distributing threads on the whole machine while taking affinity into account. However, it is possible that a whole thread group has far less work than others and terminates before them, leaving idle the whole part of the machine that was running it.
To *correct* such imbalance, some bubbles may be regenerated and moved up. Idle processors would then move some of them down on their side and have them re-burst there, getting a new distribution suited to the new workload while keeping affinity intact.
To *prevent* such imbalances, bubbles may periodically be regenerated<sup>2</sup><sup>2</sup>2In a way similar to Unix system thread preemption.: each bubble has its own *time slice* after which its threads are preempted and the bubble regenerated.
In the case of Figure 1, the preemption mechanism is extended to *Gang Scheduling*: whenever a bubble is regenerated (because its time slice expired), it is put back at the end of the task list while another bubble is burst to occupy the resulting idle processors.
### 3.4 Discussion
Bubbles give programmers the opportunity to express the structure of their application and to guide the scheduling of their threads in a *simple*, *portable* and *structured* way. Since the roles of processors and other hierarchical levels are not predetermined, the scheduler still has some degrees of freedom and can hence use an opportunist strategy to distribute tasks over the whole machine. By taking into account any irregularity in the application, this scheduler significantly enhances the underlying machine exploitation. Such preventive rebalancing techniques may still have side effects and lead to pathological situations (ping-ponging between tasks, useless bubble migration just before termination, etc.).
## 4 Implementation details
Marcel is a two-level thread library: in a way similar to manual scheduling (see section 2.1), it binds one kernel-level thread on each processor and then performs fast user-level context switches between user-level threads, hence getting complete control on threads scheduling<sup>3</sup><sup>3</sup>3We suppose that no other application is running, and neglect system daemons wake-ups. in userland without any further help from the kernel. Our proposal was implemented within Marcel’s user threads scheduler.
Figure 4 shows an example of using the interface to build and launch a bubble containing two Marcel threads.
The Marcel scheduler already had per-processor thread lists, so that integrating bubbles within the library did not need a thorough rewriting of the data structures. The scheduler code was modified to implement list hierarchy, bubble evolution and to take priorities (described in Section 3.3.2) into account.
So as to avoid contention, there is no global scheduling: processors just call the scheduler code themselves whenever they preempt (or terminate) a thread. The scheduler finds some thread that is ready to be executed by the processor. We added bubble management there: while looking for threads to execute, the scheduler code now also tries to “pull down” bubbles from high list levels and make them burst on a more local level. Getting an efficient implementation is complex, as explained below.
Given a processor, two passes are done to look for the task (thread or bubble) with maximum priority among all the tasks of the lists “covering” that processor. The first pass quickly finds the list containing the task with the highest priority, without the need of a lock. That list and the list holding the currently running task are locked<sup>4</sup><sup>4</sup>4By convention, locking lists is done by locking high-level lists first, and for a given level, according to the level elements identifiers.. A second pass is then used to check that the selected list still has a task of this priority, in case some other processor took it in the meantime. If the selected task is a thread, it is scheduled; otherwise it is a bubble that the processor deals with appropriately (going down / bursting). The implementation time-complexity is linear with respect to the number of hierarchical levels of the machine.
Regenerating a bubble is also a difficult operation. Replacing threads in a given bubble requires removing all of them from the task lists, except threads being executed. Those threads go back in the bubble by themselves when the processors executing them call the scheduler. Eventually, the last thread closes the bubble and moves it up to the list where it was initially released by the bubble holding it.
## 5 Performance evaluation
Our algorithm has some cost, but increases performance thanks to the resulting localization.
### 5.1 Bubble scheduler cost
We measured the performance impact of our implementation on the Marcel library running on a 2.66 GHz Pentium IV Xeon. Searching through lists has a reasonable cost, and our scheduler execution times are good compared to the Linux thread libraries LinuxThread (2.4 kernel) and NPTL (2.6 kernel), see Table 1.
Creation and destruction of a bubble holding a thread does not cost much more than creation and destruction of a simple thread: the cost increases from 3.3$`\mu `$s to 3.7$`\mu `$s.
Test-case examples of recursive creation of threads, such as divide-and-conquer *Fibonacci* show that the cost of systematically adding bubbles that express the natural recursion of threads creations is quickly balanced by the localization that they bring: Figure 5 shows that performance is affected when only a few threads are created, while on a HyperThreaded Bi-Pentium IV Xeon, the performance gain stabilizes at around 30 to 40% with 16 threads; on a NUMA $`4\times 4`$ Itanium II, the gain is 40% with 32 threads and gets up to 80% with 512 threads.
### 5.2 A real application
Marc Pérache used our scheduler in a comparison of the efficiency of various scheduling strategies for *heat conduction and advection* simulations. Results may be seen in Table 2. The target machine is a ccNUMA Bull NovaScale with 16 Itanium II processors and 64 GB of memory, distributed among 4 NUMA nodes. For a given processor, accessing the memory of its own node is about 3 times faster than accessing the memory of another node. The applications perform cycles of fully parallel computing followed by global hierarchical communication barrier.
In the *simple* version, the mesh is split into as many stripes as the number of processors, and an opportunist schedule is used. The *bound* version binds them to processors in a non-portable way. This gets far better performance: each thread remains on the same node, along with its data. Our proposal lets the application query Marcel about the number of NUMA nodes and processors and then automatically build bubbles according to the hierarchy of the machine (hence 4 bubbles of 4 threads in this example). It gets performance very similar to those of the *bound* version.
As can be seen, the use of bubbles attained performance close to that which may be achieved with a “handmade” thread distribution, but in a portable way.
These applications are a simple example in which the workload is balanced between stripes. The use of bubbles simply allowed it to automatically fit the architecture of the machine. However, in the future these applications will be modified to benefit from Adaptive Mesh Refinement (AMR) which increases computing precision on interesting areas. This will entail large workload imbalances in the mesh both *at runtime* and *according to the computation results*. It will hence be interesting to compare both development time and execution time of handmade-, opportunist-, and bubble-scheduled versions.
## 6 Conclusion
Multiprocessor machines are getting increasingly hierarchical. This makes task scheduling extremely complex. Moreover, the challenge is to get a scheduler that will perform “good” task scheduling on any multiprocessor machine with an arbitrary hierarchy, only guided by *portable* scheduling hints.
In this paper, we presented a new mechanism making significant progress in that direction: the bubble model lets applications express affinity relations of varying degrees between tasks in a portable way. The scheduler can then use these hints to distribute threads.
Ideally, the scheduler would need no other information to perform this. But practically speaking, writing such a scheduler is difficult and will need many experiments to be tuned. In the meantime, the programmer can use stricter guiding hints (indicating bubble bursting levels, for instance) so as to experiment with several strategies.
Performance observations on several test-cases are promising, far better than what opportunist schedulers can achieve, and close to what predetermined schedulers get. These observations were obtained on several architectures (Intel PC SMP, Itanium II NUMA).
This work opens numerous future prospects. In the short term, our proposal will be included within test-cases of real applications of CEA that run on highly hierarchical machines, hence stressing the bubble mechanism power. It will then be useful to develop analysis tools based on tracing the scheduler at runtime, so as to check and refine scheduling strategies. It will also be useful to let the programmer set other attributes than just priorities, and thus influence the scheduler: “strength” of the bubble (which expresses the amount of affinity that the bubble represents), preemptibility, some notion of amount of work, …
In the longer term, the goal is to provide a means of expression powerful and portable enough for the application to obtain an automatic schedule that gets close to the “optimal” whatever the underlying architecture. It could also be useful to provide more powerful memory allocation functions, specifying which scope of tasks (a bubble for instance) will use the allocated area. |
warning/0506/math0506177.html | ar5iv | text | # Max-plus definite matrix closures and their eigenspaces
## 1. Introduction
This paper is a contribution to geometrical understanding of some algebraic results on max-plus eigenspaces that were obtained by P. Butkovič in (see also ). The sources of geometrical inspiration for our work are and , as well as and . Our approach is closer to that of and , since we think of max-plus semiring with its simplifying total order rather than of generalized algebraic structures of idempotent analysis. However, our viewpoint differs from that of and in that we use basic tools of max-algebra instead of the more sophisticated machinery of convex geometry.
The paper is organized as follows. In Sect. 2 we recall the basic tools of max-algebra that we need. In Sect. 3 we introduce definite forms for max-plus matrices with nonzero (finite) permanent and prove that closures of all definite forms of a given matrix coincide. Thus we can introduce the ‘definite closure’ of any max-plus matrix with nonzero permanent. We also introduce definite eigenspaces, make some observations on systems of inequalities that define them, and consider an application to the cellular decomposition introduced in . In Sect. 4 we use a representation, due to , of the definite eigenspace that reveals some inner structures. Then we establish some facts about Hilbert distances between these inner structures and the boundary of the eigenspace.
The author wishes to thank P. Butkovič, A. Churkin, A. Kurnosov, A. Sobolevskiĭ, and anonymous referees for helpful comments on the paper.
## 2. Some tools of max-algebra
In its ordinary setting, max-algebra is linear algebra over the semiring $`_{\mathrm{max}}`$. This semiring is the set $`\{\mathrm{}\}`$ equipped with the operations of ‘addition’ $`=\mathrm{max}`$ and ‘multiplication’ $`=+`$. Its ‘zero’ 0 is equal to $`\mathrm{}`$ and its ‘unity’ 1 is equal to $`0`$.
The semiring $`_{\mathrm{max}}`$ resembles $`^+`$, the positive part of the field of real numbers: in both structures multiplication admits inverses and enjoys distributivity over addition, but subtraction is not allowed. In $`_{\mathrm{max}}`$, however, the addition $``$ is *idempotent*, i.e. for any $`\alpha _{\mathrm{max}}`$ we have $`\alpha \alpha =\alpha `$. In $`^+`$ this is certainly not the case. The semiring $`_{\mathrm{max}}`$ can be seen as an idempotent ‘dequantization’ of $`^+`$ (see, e.g., and ).
In max-algebra, it is also possible to exponentiate and to take roots. These operations are nothing but conventional multiplication and conventional division, respectively. Indeed, for any $`\alpha \text{0}`$ one has $`\alpha ^n=\alpha \times n`$ and $`\alpha ^{\frac{1}{n}}=\alpha /n`$ (for the remaining case $`\text{0}^n=\text{0}^{\frac{1}{n}}=\text{0}`$ since we assume that $`\alpha \text{0}=\text{0}`$ for any $`\alpha `$)
One of the principal objects of max-algebra is $`_{\mathrm{max}}^n`$, the set of $`n`$-component vectors with components in $`_{\mathrm{max}}`$. This set is equipped with ‘addition’ $`(xy)_i=x_iy_i`$ and with ‘multiplication’ by any *max-plus scalar* (i.e. by any element of $`_{\mathrm{max}}`$) $`(\alpha y)_i=\alpha y_i`$. The set $`_{\mathrm{max}}^n`$ equipped with these operations, as well as subsets of $`_{\mathrm{max}}^n`$ closed under these operations, will be called *max-plus spaces*. Below the notation $``$ will be frequently omitted.
These max-plus spaces resemble linear spaces in that the laws of associativity and distributivity hold, but again there is no subtraction and there is idempotency of addition. Structures of this kind are called *idempotent semimodules* and are a central object of the study in the *idempotent analysis*, see .
A max-plus space $`S_{\mathrm{max}}^n`$ is said to be *finitely generated* if there is a set of vectors $`\{v_1,\mathrm{},v_s\}`$ such that for any $`yS`$ one can find scalars $`\alpha _1,\mathrm{},\alpha _s`$ such that $`y=_{i=1}^s\alpha _iv_i`$. The set $`\{v_1,\mathrm{},v_s\}`$ is the *generating set* of $`S`$. It is *minimal* if no $`v_i`$ can be expressed as a linear combination of the other generators, i.e. if there are no equalities of the form
(1)
$$v_i=\underset{ji}{}\alpha _jv_j.$$
The minimal generating sets will be also called *bases*.
The following crucial result is due to Moller and to Wagneur (it is also contained in and ).
###### Proposition 1.
If $`\{u_1,\mathrm{},u_s\}`$ and $`\{v_1,\mathrm{},v_t\}`$ are two bases of a max-plus space, then $`s=t`$ and there is a permutation $`\sigma `$ and a set of nonzero scalars $`\{\alpha _1,\mathrm{},\alpha _s\}`$ such that $`u_i=\alpha _iv_{\sigma (i)}`$ for all $`i=1,\mathrm{},s`$.
Prop. 1 means that if we have found a finite base for a max-plus space, then it is in some sense unique: we can only multiply the vectors of the base by nonzero scalars. For more information on max-plus bases we refer the reader to .
The max-plus matrix algebra is formally analogous to the conventional matrix algebra (minus subtraction and plus idempotency): $`(AB)_{ij}=A_{ij}B_{ij}`$ and $`(AB)_{ij}=_kA_{ik}B_{kj}`$.
Let us introduce two important characteristics of max-plus matrices. The first characteristic deals with the cyclic permutations. Let $`A`$ be an $`n\times n`$ max-plus matrix. Here and below $`N`$ will stand for the $`n`$-element set $`\{1,\mathrm{},n\}`$. Denote by $`C_n`$ the set of all cyclic permutations $`\tau `$ that act on the subsets of the set $`N`$. For $`\tau C_n`$ denote by $`K(\tau )`$ the subset on which $`\tau `$ acts and by $`K(\tau )`$ the number of elements in this subset. Then
(2)
$$\lambda (A)=\underset{\tau C_n}{}(\underset{iK(\tau )}{}A_{i\tau (i)})^{\frac{1}{K(\tau )}}$$
is the *maximal cycle mean* of the matrix $`A`$ (the notation $`\lambda (A)`$ for the maximal cycle mean of $`A`$ will be used throughout the paper). The summand $`(_{iK(\tau )}A_{i\tau (i)})^{\frac{1}{K(\tau )}}`$ is called the *cycle mean* of the cyclic permutation $`\tau `$.
The cyclic permutation whose cycle mean is maximal will be called *critical*.
The second characteristic deals with the permutations of $`N`$. For any square $`n\times n`$ max-plus matrix $`A`$, its *permanent* is defined as
(3)
$$\text{per}(A)=\underset{\sigma S_n}{}\underset{i=1}{\overset{n}{}}A_{i\sigma (i)}.$$
Here $`S_n`$ is the group of all permutations of $`N`$.
The summand $`_{i=1}^nA_{i\sigma (i)}`$ is called the *weight of the permutation* $`\sigma `$, so the max-plus permanent is the maximal weight of all permutations of $`N`$. The permanent of $`A`$ is said to be *strong* (following ), if $`A`$ has only one *maximal permutation*, i.e., only one permutation with maximal weight.
In the papers and the max-plus permanent is called the ‘tropical determinant’. To some extent the max-plus permanent can overtake the role that the usual determinant plays, see and the just mentioned papers for details. We also refer the reader to and , for a symmetrized version of max-algebra, which admits subtraction and determinants.
An $`n\times n`$ max-plus matrix $`A`$ is called *invertible*, if there is another $`n\times n`$ max-plus matrix $`B`$ such that the products $`AB`$ and $`BA`$ are both equal to the max-plus identity matrix $`I`$. Prop. 1 implies that $`A`$ is invertible iff there is a permutation $`\sigma `$ and a set of nonzero scalars $`\alpha _1,\mathrm{},\alpha _n`$ such that
$$A_{ij}=\{\begin{array}{cc}\alpha _i,\hfill & \text{if }j=\sigma (i)\text{;}\hfill \\ \text{0},\hfill & \text{otherwise.}\hfill \end{array}$$
So the class of invertible matrices in max-algebra is very small. But it makes sense to calculate the series
(4)
$$A^{}=IAA^2\mathrm{},$$
where $`I`$ is the max-plus identity matrix. If the sum of this series exists, then it is called the *max-plus algebraic closure* of $`A`$. It is an obvious analogue of $`(IA)^1`$.
Powers of max-plus matrices and max-plus closures play important role in optimization on graphs. Indeed, if we associate an $`n`$-node graph with an $`n\times n`$ matrix $`A`$ and let $`A_{ij}`$ be the weight of the edge $`(i,j)`$ of this graph, then the entry $`(A^m)_{ij}`$ of the matrix $`A^m`$ will represent the maximal weight of paths of length $`m`$ running from $`i`$ to $`j`$. The entry $`(A^m)_{ii}`$ is the maximal weight of all cyclic paths, i.e. *cycles*, of length $`m`$ that traverse $`i`$. In other words, it is the maximal weight of the cyclic permutations $`\tau `$ such that $`iK(\tau )`$ and $`K(\tau )=m`$. Analogously, $`(A^{})_{ij}`$, for $`ij`$, is the maximal weight of all paths of any length running from $`i`$ to $`j`$. The path from $`i`$ to $`j`$ whose weight equals $`(A^{})_{ij}`$ is called *optimal*.
The following proposition, due to Carré , solves the problem of existence of the closure.
###### Proposition 2.
The closure of the matrix $`A`$ exists if and only if $`\lambda (A)\text{1}`$.
Max-plus closures enjoy the property
(5)
$$(A^{})^2=A^{}.$$
Otherwise stated, the inequality $`A_{ij}^{}A_{jk}^{}A_{ik}^{}`$ holds for all $`i`$,$`j`$, and $`k`$.
As $`A^{}`$ is an equivalent of $`(IA)^1`$, some algorithms of linear algebra can be adapted to calculate closures (that is, all optimal paths on a graph) in max-algebra, see , , and .
Another important tool is the max-plus spectral theory. For $`A`$, an $`n\times n`$ max-plus square matrix, the *max-plus spectral problem* consists in finding an $`n`$-element vector $`x`$, such that not all of its components are 0, and a scalar $`\lambda `$ such that
(6)
$$Ax=\lambda x.$$
The scalar $`\lambda `$ is a *max-plus eigenvalue*, and the vector $`x`$ is a *max-plus eigenvector*. The set of max-plus eigenvectors associated with $`\lambda `$ is closed under addition and multiplication by any nonzero scalar. Therefore it is called the *max-plus eigenspace* (associated with $`\lambda `$). The eigenspace associated with the maximal eigenvalue of $`A`$ will be denoted by $`\text{eig}(A)`$, and the space generated by the columns of $`A`$ will be denoted by $`\text{span}(A)`$.
One of the main results on max-plus spectral theory is the following.
###### Proposition 3.
The maximal eigenvalue of any max-plus square matrix is equal to its maximal cycle mean.
The eigenspace associated with the maximal eigenvalue is easy to describe.
###### Proposition 4.
Let $`A`$ be an $`n\times n`$ square max-plus matrix such that $`\lambda (A)=\text{1}`$. Then
* $`\text{eig}(A)`$ is generated by the columns $`A_j^{}`$ of $`A^{}`$ such that $`AA_j^{}=A_j^{}`$;
* $`AA_j^{}=A_j^{}`$ iff there is a critical cyclic permutation $`\tau `$ such that $`jK(\tau )`$.
Columns of $`A^{}`$ corresponding to vertices of the same cycle are proportional (in the max-plus sense) to each other.
###### Proposition 5.
Let $`A`$ be an $`n\times n`$ max-plus matrix such that $`\lambda (A)=\text{1}`$ and let $`\tau `$ be a critical cyclic permutation. Then for any $`lN`$ and $`iK(\tau )`$ we have
(7)
$$\begin{array}{c}A_{il}^{}=A_{i\tau (i)}A_{\tau (i)l}^{},\hfill \\ A_{li}^{}=A_{l\tau ^1(i)}^{}A_{\tau ^1(i)i}.\hfill \end{array}$$
If the graph associated with $`A`$ is strongly connected (i.e. for any $`i`$ and $`j`$ there is a path with nonzero weight running from $`i`$ to $`j`$), then $`A`$ is said to be *irreducible*. In this case the maximal cycle mean of $`A`$ is known to be its only eigenvalue. In the reducible case this eigenvalue need not be unique. If another eigenvalue exists, then it is the (only) eigenvalue of some maximal irreducible submatrix of $`A`$. However, not all maximal irreducible submatrices of $`A`$ yield eigenvalues for $`A`$.
For more details on max-plus spectral theory, as well as for the proofs of the above propositions, we refer the reader to , , , and .
## 3. Definite eigenspaces and definite closures
We begin this section with the following important proposition. The proof makes use of the uniqueness of the base (Prop. 1).
###### Proposition 6.
Let $`A`$ and $`B`$ be square max-plus matrices such that $`\lambda (A)\text{1}`$, and $`\lambda (B)\text{1}`$. Then $`\text{span}(A^{})=\text{span}(B^{})`$ if and only if $`A^{}=B^{}`$.
Proof. First let us prove that whenever
(8)
$$A_i^{}=\underset{ji}{}\alpha _jA_j^{},$$
the column $`A_i^{}`$ is proportional to $`A_j^{}`$ for some $`ji`$.
Suppose (8) holds. Then there is an $`l`$ such that $`A_{ii}^{}=\text{1}=\alpha _lA_{il}^{}`$ and that $`A_{li}^{}\alpha _l`$. Combining this we obtain that $`A_{il}^{}A_{li}^{}\text{1}`$, hence $`A_{il}^{}A_{li}^{}=\text{1}`$ (otherwise there is a cycle whose weight exceeds 1). This means that there is a critical cycle with weight 1 traversing $`i`$ and $`l`$, and due to Prop. 5 $`A_i^{}=\alpha _lA_l^{}`$. We conclude that no column of $`A^{}`$ can be expressed as linear combination (8) without being proportional to some of the columns involved in this combination.
Further let $`\{A_{r_1}^{},\mathrm{},A_{r_k}^{}\}`$ be the base of $`\text{span}(A^{})=\text{span}(B^{})`$. If we use columns of $`B^{}`$ to form the base, then, due to Prop. 1, it must be of the form $`\{B_{s_1}^{},\mathrm{},B_{s_k}^{}\}`$, so that $`B_{s_i}^{}=\alpha _iA_{r_{\sigma (i)}}^{}`$ for $`i=1,\mathrm{},k`$ and some nonzero $`\alpha _i`$. All remaining columns of $`A^{}`$ (of $`B^{}`$) are proportional to base columns of $`A^{}`$ (of $`B^{}`$). So every column of $`A^{}`$ is proportional to some column of $`B^{}`$, and vice versa. This implies that the rows $`A_i^{}`$ and $`A_j^{}`$ are proportional iff the rows $`B_i^{}`$ and $`B_j^{}`$ are proportional.
Let the columns $`A_i^{}`$ and $`A_j^{}`$ be proportional. Then $`A_{ij}^{}A_{ji}^{}=\text{1}`$, hence there is a critical cycle containing $`i`$ and $`j`$. Due to Prop. 5 the rows $`A_i^{}`$ and $`A_j^{}`$ are proportional, so are the rows $`B_i^{}`$ and $`B_j^{}`$, and, again due to Prop. 5, so are the columns $`B_i^{}`$ and $`B_j^{}`$. We conclude that the columns $`A_i^{}`$ and $`A_j^{}`$ are proportional iff so are the columns $`B_i^{}`$ and $`B_j^{}`$.
Now it is clear that $`\{A_{r_1}^{},\mathrm{},A_{r_k}^{}\}`$ is the base of $`\text{span}(A^{})`$ iff $`\{B_{r_1}^{},\mathrm{},B_{r_k}^{}\}`$ is also the base, so that $`B_{r_i}^{}=\alpha _iA_{r_{\sigma (i)}}^{}`$. We can assume w.l.o.g. that $`r_i=i`$. Consider the decomposition of $`\sigma `$ into cyclic permutations and let $`\tau `$ be one of them. Then $`B_{\tau (i)i}^{}=\alpha _i`$ and $`A_{i\tau (i)}^{}=\alpha _i^1`$ for all $`iK(\tau )`$. If $`_{iK(\tau )}\alpha _i>\text{1}`$, then $`B^{}`$ has a cycle with weight greater than 1, and if $`_{iK(\tau )}\alpha _i<\text{1}`$, then so does $`A^{}`$. The only remaining possibility $`_{iK(\tau )}\alpha _i=\text{1}`$ implies that all columns of $`A^{}`$ and $`B^{}`$ with indices belonging to $`K(\tau )`$ are proportional. Then $`K(\tau )`$ must be a singleton for any $`\tau `$, otherwise the minimality of the bases is violated. This implies that $`\sigma `$ is the identity permutation, and $`A_i^{}=B_i^{}`$ for any $`i=1,\mathrm{},k`$. Taking into account that there is a one-to-one correspondence between the sets of proportional columns of $`A^{}`$ and $`B^{}`$, and that all columns of $`A^{}`$ and $`B^{}`$ are proportional to some base columns, we conclude that $`A^{}=B^{}`$. $`\mathrm{}`$
Now consider matrices with maximal cycle mean equal to 1 and with all diagonal entries equal to 1. Following , we call such matrices *definite*. The following proposition contains some simple facts on eigenspaces of such matrices. The third statement is an easy consequence of , Ch. IV, Th. 2.2.4, and its proof is recalled here for convenience of the reader. We consider the general reducible case, in which the eigenvectors may have zero entries. For any $`y_{\mathrm{max}}^n`$, the index set $`K`$ such that $`y_i\text{0}`$ iff $`iK`$ is called the *support* of $`y`$ and is denoted by $`\text{supp}(y)`$.
###### Proposition 7.
If $`A`$ is a definite matrix, then
* it has a unique eigenvalue equal to 1;
* $`\text{eig}(A)=\text{span}(A^{})`$;
* an eigenvector with the support $`KN`$ exists iff $`A_{ij}=\text{0}`$ for all $`iN\backslash K`$ and $`jK`$.
Proof. 1) Any eigenvalue of $`A`$ is the maximal cycle mean of some of its submatrices, and the maximal cycle mean of any submatrix of $`A`$ is equal to 1.
2) follows from Prop. 4.
3) If an eigenvector $`x`$ such that $`\text{supp}(x)=K`$ exists, then $`A_{ij}x_j=\text{0}`$ for all $`iN\backslash K`$ and all $`jK`$. Hence $`A_{ij}=\text{0}`$ for all such $`i`$ and $`j`$.
Conversely, if the set $`K`$ satisfies the condition, we look for eigenvectors $`x`$ such that $`\text{supp}(x)=K`$. We may reduce the system $`Ax=x`$ to the system $`A_{KK}x_K=x_K`$, where $`A_{KK}`$ is a submatrix of $`A`$ standing on the rows and columns with indices belonging to $`K`$, and $`x_K`$ is a vector with $`K`$ nonzero components. The space $`\text{eig}(A_{KK})`$ is generated by the columns $`(A_{KK})_j^{}`$, where $`jK`$. Taking any combination of all these generators with all coefficients not equal to 0 we obtain an eigenvector with the support $`K`$. $`\mathrm{}`$
In the case of definite matrices we have one more implication of the uniqueness of the base. The proof is similar to that of Prop. 6.
###### Proposition 8.
If $`A`$ is definite and $`\text{span}(A)=\text{eig}(A)`$, then $`A=A^{}`$.
Proof. Let $`\{A_{s_1},\mathrm{}A_{s_k}\}`$ be the base of $`\text{span}(A)=\text{eig}(A)=\text{span}(A^{})`$. Due to Prop. 1, if we use columns of $`A^{}`$ to form the base, then it must be of the form $`\{A_{t_1}^{},\mathrm{}A_{t_k}^{}\}`$ so that $`A_{s_i}=\alpha _iA_{t_i}^{}`$ for $`i=1,\mathrm{},k`$ and some nonzero $`\alpha _i`$. More precisely, $`\alpha _i=A_{t_is_i}`$, and this implies $`A_{t_is_i}A_{s_it_i}^{}=\text{1}`$. Then there is a critical cycle traversing $`s_i`$ and $`t_i`$, so $`A_{s_i}^{}=A_{t_is_i}A_{t_i}^{}`$ according to Prop. 5. So $`A_{s_i}=A_{s_i}^{}`$ for all $`i=1,\mathrm{},k`$.
Now assume that there are columns $`A_j`$ and scalars $`\alpha _i`$ such that
(9)
$$A_j=\underset{i=1}{\overset{k}{}}\alpha _iA_{s_i},$$
and that no column of the base is proportional to $`A_j`$. It follows from (9) that there is an index $`m`$ such that $`A_{jj}=\text{1}=\alpha _mA_{js_m}`$. The columns $`A_j`$ and $`A_{s_m}`$ are not proportional, hence there is an $`l`$ such that $`A_{lj}>\alpha _mA_{ls_m}`$. This implies $`A_{lj}A_{js_m}>A_{ls_m}`$ and $`A_{ls_m}^{}>A_{ls_m}`$. This is a contradiction, since we have proved that $`A_{s_i}^{}=A_{s_i}`$ for any $`i=1,\mathrm{},k`$. So any column of $`A`$ is proportional to some column of the base. But if $`A_i`$ and $`A_j`$ are proportional, then $`A_{ij}A_{ji}=\text{1}`$, hence $`A_i^{}`$ and $`A_j^{}`$ are also proportional with the same coefficient. This implies $`A=A^{}`$. $`\mathrm{}`$
Let us now introduce a *definite form* of a matrix. Consider an $`n\times n`$ max-plus matrix $`A`$ that has nonzero permanent, i.e., at least one permutation whose weight is not equal to 0. Let $`\sigma `$ be one of the maximal permutations of $`A`$. If this permutation is not the identity permutation, then we turn this permutation into the identity permutation by rearranging the columns of $`A`$. Then we divide (in the max-plus sense) all columns by the corresponding diagonal entries, thus obtaining a matrix $`A^{}`$ with entries
(10)
$$A_{ij}^{}=A_{i\sigma (j)}A_{j\sigma (j)}^1.$$
Obviously, passing to a definite form of a matrix does not alter its span: $`\text{span}(A^{})=\text{span}(A)`$.
It is clear that $`A^{}`$ is definite, hence we call it the *definite form* of $`A`$ corresponding to the permutation $`\sigma `$. For example, if
$$A=\left(\begin{array}{ccc}1& 5& \mathrm{}\\ 2& 1& 7\\ 6& \mathrm{}& 2\end{array}\right),$$
then
$$A^{}=\left(\begin{array}{ccc}0& \mathrm{}& 5\\ 4& 0& 4\\ \mathrm{}& 5& 0\end{array}\right)$$
is the definite form of $`A`$ corresponding to $`\sigma =(231)`$. (Here and in the sequel we prefer not to use the symbols 0 and 1 in numerical examples.)
However, in general there are many maximal permutations and many definite forms corresponding to them. The fact that we want to prove is that closures of all definite forms coincide. It is convenient to pose the problem as follows. Let $`A`$ be definite, let it have maximal permutations different from the identity permutation, and let $`\sigma `$ be one of them. Let $`A^{}`$ be the definite form of $`A`$ corresponding to $`\sigma `$, then its entries are defined according to (10). $`A^{}`$ also has maximal permutations different from the identity permutation, and $`\sigma ^1`$ is one of them, since
(11)
$$A_{i\sigma ^1(i)}^{}=A_{\sigma ^1(i)i}^1.$$
First we prove the following proposition.
###### Proposition 9.
Let $`A`$ be a definite matrix and let $`\sigma `$ be one of its maximal permutations. For $`A^{}`$ the definite form of $`A`$ corresponding to $`\sigma `$, $`\text{eig}(A^{})=\text{eig}(A)`$.
Proof. Consider the decomposition of $`\sigma `$ into cyclic permutations. Let $`r`$ be the number of these permutations and denote these permutations by $`\tau _l`$ (where $`l=1,\mathrm{},r`$). Denote by $`K(\tau _l)`$ the index set on which $`\tau _l`$ acts. The sets $`K(\tau _l)`$ are pairwise disjoint and $`_{l=1}^rK(\tau _l)=N`$.
We prove the inclusion $`\text{eig}(A)\text{eig}(A^{})`$. Let $`y`$ be an eigenvector of $`A`$ with the support $`MN`$. Due to the third statement of Prop. 7 the set $`M`$ must be of the form $`_{lL}K(\tau _l)`$ for some $`L\{1,\mathrm{},r\}`$, and we have $`A_{ij}=\text{0}`$ for all $`iN\backslash M`$ and $`jM`$.
Now note that every cyclic permutation $`\tau _l`$ is critical, hence for any $`\tau _l`$ such that $`K_lM`$ we have
(12)
$$\underset{iK(\tau _l)}{}A_{i\tau _l(i)}y_{\tau _l(i)}y_i^1=\text{1}.$$
On the other hand we have $`A_{i\tau _l(i)}y_{\tau _l(i)}y_i`$. This inequality is an equality, otherwise the violation of (12) occurs. Thus we obtain
(13)
$$A_{i\sigma (i)}y_{\sigma (i)}=y_i$$
for all $`iM`$.
Now note that, since $`\sigma (j)M`$ for all $`jM`$, we have that $`A_{ij}^{}=\text{0}`$ for all $`iN\backslash M`$ and $`jM`$, the same as for $`A_{ij}`$. Therefore it suffices to prove that if $`A_{MM}y_M=y_M`$, then $`A_{MM}^{}y_M=y_M`$. In other words, we want to show that $`A_{ij}^{}y_jy_i`$ for all $`i,jM`$. But due to (10) and (13)
$$\begin{array}{cc}A_{ij}^{}y_j=\hfill & A_{i\sigma (j)}A_{j\sigma (j)}^1y_j=\hfill \\ \hfill =& A_{i\sigma (j)}y_{\sigma (j)}y_{\sigma (j)}^1A_{j\sigma (j)}^1y_j=\hfill \\ \hfill =& A_{i\sigma (j)}y_{\sigma (j)}y_i,\hfill \end{array}$$
and the proof of the inclusion $`\text{eig}(A)\text{eig}(A^{})`$ is complete. The opposite inclusion is proved analogously, after passing to the definite form associated with $`\sigma `$ (the matrices $`A`$ and $`A^{}`$ will interchange). $`\mathrm{}`$
The fact that we want to prove is now an immediate consequence of Propositions 6 and 9.
###### Proposition 10.
Closures of all definite forms of any matrix with nonzero permanent coincide.
We have just proved that closures of all definite forms coincide. Now we can define the *definite closure* of any $`n\times n`$ max-plus matrix with nonzero permanent to be the closure of any of its definite forms. Due to the second statement of Prop. 7, for any definite form $`A^{}`$ of $`A`$ we have $`\text{eig}(A^{})=\text{eig}((A^{})^{})=\text{span}((A^{})^{})`$. So the eigenspace of the definite closure of $`A`$ coincides with the eigenspace of any of its definite forms (they are all the same), and is generated by columns of the definite closure. The eigenspace of definite closure will be called the *definite eigenspace*.
The third statement of Prop. 7 suggests that, if we want to work with the space of eigenvectors of definite closure of non-full support, then we must confine ourselves to the corresponding submatrix.
Further we always assume that the eigenvectors considered have full support, i.e., that we study the eigenvectors with certain support and have passed to the corresponding submatrix.
Let us show that the definite eigenspace can be described by some system of inequalities.
###### Proposition 11.
Let $`A`$ be a definite matrix. Then its eigenspace (and the eigenspace of its closure) is the set $`X=\{xA_{ij}x_ix_j^1,ij,A_{ij}=A_{ij}^{}\}`$
Proof. Let $`x`$ be an eigenvector of $`A`$ corresponding to its maximal eigenvalue 1. It satisfies the system $`Ax=x`$, therefore it satisfies all inequalities of the form $`A_{ij}x_ix_j^1`$. Hence $`xX`$.
Conversely, if all inequalities $`A_{ij}x_ix_j^1`$, for $`ij`$ and $`A_{ij}=A_{ij}^{}`$, are satisfied, then absolutely all inequalities $`A_{ij}x_ix_j^1`$ are satisfied. Indeed, if $`A_{ij}<A_{ij}^{}`$, then the optimal path from $`i`$ to $`j`$ is not the edge $`(i,j)`$, but it traverses other nodes, say, $`i_1,\mathrm{},i_k`$. Then $`A_{ij}^{}=A_{ii_1}\mathrm{}A_{i_kj}`$ where $`A_{ii_1}=A_{ii_1}^{},\mathrm{},A_{i_kj}=A_{i_kj}^{}`$ (all edges $`(i,i_1),\mathrm{},(i_k,j)`$ are optimal paths). The inequality $`A_{ij}x_ix_j^1`$ is now an easy consequence of the inequalities $`A_{ii_1}x_ix_{i_1}^1,\mathrm{},A_{i_kj}x_{i_k}x_j^1`$ that are satisfied. So all inequalities $`A_{ij}x_ix_j^1`$ are satisfied and this implies $`Ax=x`$. $`\mathrm{}`$
An inverse problem can also be posed. Suppose that there is a system of inequalities $`\{a_{ij}x_ix_j^1\}`$, with at most one inequality per each pair $`(i,j)`$, and the set of vectors with full support defined by this system is not empty. Is it a full-support subspace of an eigenspace of a definite matrix?
To answer this question, consider the matrix $`A`$ whose entries $`A_{ij}`$ are equal either to 1, if $`i=j`$, or to $`a_{ij}`$, if there is an inequality of the form $`a_{ij}x_ix_j^1`$, or to 0, if there is no such inequality. Then we have the following proposition.
###### Proposition 12.
The set $`X=\{xA_{ij}x_ix_j^1\}`$ is nonempty if and only if $`A`$ is definite.
Proof. If $`A`$ is definite then the set of its eigenvectors associated with the eigenvalue 1 is nonempty and, due to Prop. 11, it is precisely the set $`X`$ (some of inequalities being redundant).
Conversely, let $`X`$ be nonempty, then there exists $`xX`$. Take an arbitrary cyclic permutation $`\tau `$. As a consequence of all inequalities $`A_{i\tau (i)}x_ix_{\tau (i)}^1`$, where $`iK(\tau )`$, we obtain that $`_{iK(\tau )}A_{i\tau (i)}\text{1}`$. Hence all cycle means of $`A`$ are not greater than 1, and $`A`$ is definite. $`\mathrm{}`$
Now consider the following application.
Let $`V`$ be an $`m\times n`$ max-plus matrix with at least one nonzero entry in each column, and let $`y_{\mathrm{max}}^m`$ have full support. Then, following , we can define the *combinatorial type* of $`y`$ with respect to $`V`$. It is an $`m`$-tuple $`S`$ of subsets $`S_1,\mathrm{},S_m`$ of $`\{1,\mathrm{},n\}`$, such that $`iS_j`$ whenever $`_{k=1}^mV_{ki}y_k^1=V_{ji}y_j^1`$. It is proved in that the collection of the sets
(14)
$$X_S=\{yV_{ki}V_{ji}^1y_ky_j^1\text{for }j,k=1,\mathrm{},m\text{ and }iS_j\}$$
defines a cellular decomposition of the full-support subspace of $`_{\mathrm{max}}^m`$.
Consider an $`m\times m`$ matrix $`V^S`$ whose columns are defined by
(15)
$$V_j^S=\underset{iS_j}{}V_{ji}^1V_i,$$
if $`S_j`$ is not empty, and by
(16)
$$V_{ij}^S=\{\begin{array}{cc}\text{1},\hfill & \text{if }i=j\text{;}\hfill \\ \text{0},\hfill & \text{if }ij\text{.}\hfill \end{array}$$
if $`S_j`$ is empty. Then it is clear that
(17)
$$X_S=\{yV_{ij}^Sy_iy_j^1\text{for }i,j=1,\mathrm{},m\},$$
and we immediately have the following proposition.
###### Proposition 13.
* The cell $`X_S`$ is nonempty if and only if $`\lambda (V^S)\text{1}`$;
* If $`\lambda (V^S)\text{1}`$, then $`X_S=\{y\text{span}((V^S)^{})\text{supp}(y)=N\}`$;
* If $`\lambda (V^S)\text{1}`$, then $`X_S=\{xV_{ij}^Sx_ix_j^1,ij,V_{ij}^S=(V^S)_{ij}^{}\}`$.
We close this section with three examples.
*Example 1* Consider the matrix
$$A=\left(\begin{array}{ccc}1& 3& 0\\ 2& 0& 0\\ 0& 1& 5\end{array}\right).$$
The only maximal permutation of $`A`$ is $`(231)`$, and the only definite form is
$$A^{}=\left(\begin{array}{ccc}0& 0& 1\\ 3& 0& 2\\ 4& 5& 0\end{array}\right).$$
The definite closure of $`A`$ is
$$(A^{})^{}=\left(\begin{array}{ccc}0& 0& 2\\ 2& 0& 2\\ 4& 4& 0\end{array}\right).$$
Fig. 1 displays the cross section by $`z=0`$ of $`\text{span}(A)`$ (left) and $`\text{span}((A^{})^{})=\text{eig}(A^{})`$ (right).
Note that $`\text{eig}(A^{})`$ is the set
$$\{(x,y,z)xy^10,yz^12,zx^14\},$$
in accordance with Prop. 11.
*Example 2* Consider the matrix
$$B=\left(\begin{array}{ccc}2& 0& 2\\ 1& 1& 3\\ 0& 3& 2\end{array}\right)$$
It has two maximal permutations: $`(13)(2)`$ and $`(231)`$. Therefore it has two definite forms, namely
$$B^{}=\left(\begin{array}{ccc}0& 1& 2\\ 1& 0& 1\\ 4& 4& 0\end{array}\right)$$
and
$$B^{\prime \prime }=\left(\begin{array}{ccc}0& 1& 2\\ 1& 0& 1\\ 3& 5& 0\end{array}\right).$$
But, in accordance with Prop. 10, the definite closure
$$(B^{})^{}=(B^{\prime \prime })^{}=\left(\begin{array}{ccc}0& 1& 2\\ 1& 0& 3\\ 3& 4& 0\end{array}\right)$$
is unique.
Fig. 2 displays the cross section by $`z=0`$ of $`\text{span}(B)`$ (left) and $`\text{span}((B^{})^{})=\text{eig}(B^{})`$ (right).
In accordance with Prop. 11, the space $`\text{eig}(B^{})`$ is the set
$$\{(x,y,z)yx^11,xy^11,xz^12,zy^14\},$$
or, equivalently, the set
$$\{(x,y,z)yx^11,xy^11,xz^12,zx^13\}.$$
The first system of inequalities corresponds to the definite form $`B^{}`$ and the second one corresponds to $`B^{\prime \prime }`$.
Comparing Fig. 1 with Fig. 2 we see that $`\text{eig}(A^{})`$ has ‘interior’ whereas $`\text{eig}(B^{})`$ does not have ‘interior’ (for the exact meaning of the term ‘interior’ see Sect. 4 below). As a consequence of , Th. 4.2, or , Th. 4.2, one can obtain that $`\text{eig}(A^{})`$, for $`A^{}`$ a definite form of $`A`$, has ‘interior’ if and only if $`A`$ (or equivalently $`A^{}`$) has strong permanent. This fact will be revisited in Prop. 14 of this paper.
*Example 3* Consider the matrix
$$V=\left(\begin{array}{cccc}1& 4& 6& 7\\ 4& 1& 5& 8\\ 0& 0& 0& 0\end{array}\right).$$
Let $`S=(\{2,3\},\{4\},\{1\})`$, $`P=(\mathrm{},\{4\},\{1,2,3\})`$, $`U=(\{3\},\{1,3,4\},\{2\})`$, and $`W=(\{1,4\},\{2\},\{3\})`$ be four combinatorial types. Do they exist in the cellular decomposition? If they do, what vectors generate the respective cells $`X_S`$, $`X_P`$, $`X_U`$, and $`X_W`$?
For $`S`$:
$$V^S=\left(\begin{array}{ccc}0& 1& 1\\ 1& 0& 4\\ 4& 8& 0\end{array}\right),(V^S)^{}=\left(\begin{array}{ccc}0& 1& 3\\ 0& 0& 4\\ 4& 5& 0\end{array}\right).$$
Hence $`X_S`$ exists and is generated by $`[\mathrm{4\hspace{0.33em}4\hspace{0.33em}0}]^T`$, $`[\mathrm{4\hspace{0.33em}5\hspace{0.33em}0}]^T`$, and $`[\mathrm{3\hspace{0.33em}4\hspace{0.33em}0}]^T`$.
For $`P`$:
$$V^P=\left(\begin{array}{ccc}0& 1& 6\\ \mathrm{}& 0& 5\\ \mathrm{}& 8& 0\end{array}\right)=(V^P)^{}.$$
Hence $`X_P`$ is generated by $`[0\mathrm{}\mathrm{}]^T`$, $`[\mathrm{7\hspace{0.33em}8\hspace{0.33em}0}]^T`$ and $`[\mathrm{6\hspace{0.33em}5\hspace{0.33em}0}]^T`$. However, $`[0\mathrm{}\mathrm{}]^T`$ does not have full support and does not belong to $`X_P`$.
For $`U`$:
$$V^U=\left(\begin{array}{ccc}0& 1& 4\\ 1& 0& 1\\ 6& 4& 0\end{array}\right),(V^U)^{}=\left(\begin{array}{ccc}0& 1& 4\\ 1& 0& 3\\ 5& 4& 0\end{array}\right).$$
Hence $`X^U`$ exists and is generated by $`[\mathrm{5\hspace{0.33em}4\hspace{0.33em}0}]^T`$ and $`[\mathrm{4\hspace{0.33em}3\hspace{0.33em}0}]^T`$.
For $`W`$:
$$V^W=\left(\begin{array}{ccc}0& 3& 6\\ 3& 0& 5\\ 1& 1& 0\end{array}\right).$$
We have $`\lambda (V^W)=3>\text{1}`$, hence $`W`$ does not exist in the cellular decomposition.
Fig. 3 displays $`\text{span}(V)`$ (blue), $`X_S`$ (red), $`X_P`$ (light grey), and $`X_U`$(dark green), projected onto $`z=0`$. The generators of $`\text{span}(V)`$ (the larger circles) and the generators of $`X_S`$ and $`X_U`$(the smaller squares) are also shown.
## 4. Inner structures of definite eigenspaces
We have introduced the definite closure operation and have given an external description of a definite eigenspace in terms of a system of inequalities. In this section we give an internal description of definite eigenspace and measure Hilbert distances between the structures involved in this description and the boundary.
For further considerations we need the following notions and notation. Let $`A`$ be a definite matrix. The sets $`X_{ij}=\{xA_{ij}=x_ix_j^1\}`$ will be called the *supporting planes* of $`\text{eig}(A)`$, and the sets $`\mathrm{\Gamma }_{ij}=X_{ij}\text{eig}(A)`$ will be called the *faces* of $`\text{eig}(A)`$. The *boundary*, i.e., the union of all faces of $`\text{eig}(A)`$ will be denoted by $`\mathrm{\Gamma }(\text{eig}(A))`$. The set of eigenvectors not belonging to the boundary will be called the *interior* of $`\text{eig}(A)`$ and denoted by $`\text{int}(\text{eig}(A))`$.
Also, denote by $`\stackrel{~}{A}`$ the matrix obtained from $`A`$ by replacing the diagonal 1’s by 0’s, and denote by $`A_\mu `$ the matrix obtained from $`\mu ^1\stackrel{~}{A}`$ by replacing the diagonal 0’s by 1’s. For example, if
$$A=\left(\begin{array}{ccc}0& 4& 1\\ 1& 0& 1\\ 5& 7& 0\end{array}\right)$$
then
$$\stackrel{~}{A}=\left(\begin{array}{ccc}\mathrm{}& 4& 1\\ 1& \mathrm{}& 1\\ 5& 7& \mathrm{}\end{array}\right)$$
and, e.g.,
$$A_1=\left(\begin{array}{ccc}0& 3& 2\\ 2& 0& 2\\ 4& 6& 0\end{array}\right)$$
Also, the maximal cycle mean of $`\stackrel{~}{A}`$ is $`\lambda (\stackrel{~}{A})=1.5`$, and
$$A_\lambda =\left(\begin{array}{ccc}0& 2.5& 2.5\\ 2.5& 0& 2.5\\ 3.5& 5.5& 0\end{array}\right)$$
(we write $`A_\lambda `$ instead of $`A_{\lambda (\stackrel{~}{A})}`$ for the sake of simplicity)
The following crucial proposition can be derived from .
###### Proposition 14.
Let $`A`$ be a definite matrix such that $`\lambda (\stackrel{~}{A})\text{0}`$.
* If $`A`$ does not have strong permanent then $`\text{eig}(A)`$ does not have interior;
* If $`A`$ has strong permanent then $`\text{eig}(A)`$ has interior, and
(18)
$$\text{int}(\text{eig}(A))=\underset{\lambda (\stackrel{~}{A})\mu <\text{1}}{}\text{eig}(A_\mu ).$$
Proof. 1) If $`A`$ does not have strong permanent, then it has a maximal permutation that differs from the identity permutation. Let $`\sigma `$ be such permutation, and assume that there is a $`y`$ belonging to $`\text{int}(\text{eig}(A))`$. Then $`A_{i\sigma (i)}y_{\sigma (i)}<y_i`$ for all $`i`$. After multiplying all these inequalities and then cancelling the product $`_iy_i`$ one obtains $`_iA_{i\sigma (i)}<\text{1}`$. This implies that $`\sigma `$ is not maximal, a contradiction.
2) $`y`$ belongs to $`\text{int}(\text{eig}(A))`$ if and only if $`_{ji}A_{ij}y_j<y_i`$ for all $`i`$. This takes place if and only if there is a $`\mu <\text{1}`$ such that $`_{ji}A_{ij}y_j\mu y_i`$, or, equivalently, $`A_\mu y=y`$. If we take $`\mu \lambda (\stackrel{~}{A})`$, then 1 is the maximal cycle mean of $`A_\mu `$, hence $`\text{eig}(A_\mu )`$ exists. The proof is complete.
$`\mathrm{}`$
If $`\lambda (\stackrel{~}{A})=\text{0}`$, i.e. if the graph associated with $`\stackrel{~}{A}`$ is acyclic, then representation (18) is replaced by the representation
(19)
$$\text{int}(\text{eig}(A))=\underset{\text{0}<\mu <\text{1}}{}\text{eig}(A_\mu ).$$
In the sequel, we always assume that $`A`$ is definite and has at least one off-diagonal entry not equal to 0.
According to Prop. 11, the eigenspace $`\text{eig}(A)`$ is the set
$$X=\{xA_{ij}x_ix_j^1,ij,A_{ij}=A_{ij}^{}\}.$$
Analogously, the eigenspace $`\text{eig}(A_\mu )`$, for any $`\mu `$ involved in (18) or (19), is the set
$$X_\mu =\{x\mu ^1A_{ij}x_ix_j^1,ij,\mu ^1A_{ij}=(A_\mu ^{})_{ij}\}.$$
We need the following proposition mainly for the proof of Prop. 18.
###### Proposition 15.
Let $`\mu `$ be a scalar such that $`\lambda (\stackrel{~}{A})\mu <\text{1}`$, if $`\lambda (\stackrel{~}{A})>\text{0}`$, or such that $`\text{0}<\mu <\text{1}`$, if $`\lambda (\stackrel{~}{A})=\text{0}`$. If $`(A_\mu ^{})_{ij}=(A_\mu )_{ij}`$, then $`A_{ij}^{}=A_{ij}`$.
Proof. In both cases considered the maximal cycle mean of $`A_\mu `$ is equal to 1, hence $`A_\mu ^{}`$ exists. Let $`A_{ij}^{}>A_{ij}`$, then $`A_{ij}^{}=A_{ii_1}\mathrm{}A_{i_kj}`$ for some $`i_1,\mathrm{},i_k`$ not equal to $`i`$ or $`j`$. Since $`\mu <\text{1}`$, we have $`\mu ^1A_{ii_1}\mathrm{}\mu ^1A_{i_kj}>\mu ^1A_{ij}`$, hence $`(A_\mu ^{})_{ij}>(A_\mu )_{ij}`$. $`\mathrm{}`$
The eigenspaces $`\text{eig}(A_\mu )`$ are the inner structures mentioned above. Now we are going to measure the Hilbert distances between these inner structures and the boundary $`\mathrm{\Gamma }(\text{eig}(A))`$.
The *Hilbert distance* between the two vectors $`x`$ and $`y`$ both having support $`K`$ is defined to be
(20)
$$d_H(x,y)=\underset{i,jK}{}x_ix_j^1y_i^1y_j.$$
Note that in the Hilbert distance is defined as an inverse of the quantity $`d_H(x,y)`$. If the supports of $`x`$ and $`y`$ differ, then we assume the Hilbert distance between $`x`$ and $`y`$ to be infinite.
It can be easily verified (see also , Th. 17) that the following properties hold:
* $`d_H(x,y)\text{1}`$, and $`d_H(x,y)=\text{1}`$ iff $`x=\lambda y`$, where $`\lambda `$ is a nonzero scalar;
* $`d_H(x,y)=d_H(y,x)`$;
* $`d_H(x,y)d_H(y,z)d_H(x,z)`$.
In fact these properties show that $`d_H`$ is a semidistance (recall that $`\text{1}=0`$ and $`=+`$). Indeed, $`d_H(x,y)=\text{1}=0`$ whenever $`x`$ is equivalent to $`y`$ modulo
$$xy\lambda \text{0}:x=\lambda y.$$
This semidistance is induced by the *range seminorm*
(21)
$$x=\underset{i,jK}{}x_ix_j^1,$$
introduced in , see also . However, by a slight abuse of language we will refer to $`d_H`$ as to a distance.
Now we measure the distance between an arbitrary $`y\text{eig}(A)`$ and the supporting plane $`X_{ij}=\{xx_ix_j^1=A_{ij}\}`$, i.e., the minimal distance between $`y`$ and $`xX_{ij}`$. From , Th. 18 it follows that this minimum is attained at the maximal vector of $`X_{ij}`$ not greater than $`y`$. Denote this vector by $`y^{ij}`$. Its coordinates are very easy to find:
(22)
$$\begin{array}{c}y_l^{ij}=\{x_lx_ly_l\}=y_l\text{for }li,j;\hfill \\ y_i^{ij}=\{x_ix_iy_i,A_{ij}^1x_iy_j\}=A_{ij}y_j;\hfill \\ y_j^{ij}=\{x_jx_jy_j,A_{ij}x_jy_i\}=y_j.\hfill \end{array}$$
The distance (20) between $`y`$ and $`X_{ij}`$ is then equal to
(23)
$$d_H(y,X_{ij})=A_{ij}^1y_j^1y_i.$$
However, what we need is the distance between $`y`$ and $`\mathrm{\Gamma }(\text{eig}(A))`$, i.e, the minimal distance between $`y`$ and $`\mathrm{\Gamma }_{ij}`$. The following proposition makes our life simpler.
###### Proposition 16.
The distance between $`y\text{eig}(A)`$ and the boundary $`\mathrm{\Gamma }(\text{eig}(A))`$ is equal to the minimal distance between $`y`$ and supporting planes.
Proof. Clearly the minimal distance between $`y`$ and supporting planes is not greater than the distance between $`y`$ and $`\mathrm{\Gamma }(\text{eig}(A))`$. Suppose $`X_{ij}`$ is the supporting plane such that the distance between $`y`$ and $`X_{ij}`$ is minimal. This distance is equal to the distance between $`y`$ and $`y^{ij}`$. If $`y^{ij}`$ belongs to $`\text{eig}(A)`$ and hence to $`\mathrm{\Gamma }_{ij}`$ then we are done. Suppose not; then the system of equalities $`_lA_{kl}y_l^{ij}=y_k^{ij}`$ must be violated for some $`kN`$. Note that, if $`ki`$, then $`y_k^{ij}=y_k`$ (see (22)), and there is no violation. So the violation must take place for $`k=i`$. There must be an $`l`$ such that $`A_{il}y_l^{ij}>y_i^{ij}`$, i.e. such that $`A_{il}y_l>A_{ij}y_j`$. Now consider $`z`$ such that $`z_i=A_{il}y_l`$ and $`z_k=y_k`$ for any $`ki`$. Then $`z`$ belongs to $`X_{il}`$ and $`d_H(y,z)=A_{il}^1y_l^1y_i`$. This distance is strictly less than the distance between $`y`$ and $`y^{ij}`$, a contradiction. $`\mathrm{}`$
Consequently,
(24)
$$d_H(y,\mathrm{\Gamma }(\text{eig}(A)))=\underset{ij,A_{ij}\text{0}}{}A_{ij}^1y_j^1y_i.$$
The key idea of Prop. 17 below is that $`\lambda (\stackrel{~}{A})^1`$, if $`\lambda (\stackrel{~}{A})`$ is invertible, is the largest radius of Hilbert balls contained in $`\text{eig}(A)`$. It can be said that $`\lambda (\stackrel{~}{A})^1`$ is the radius of *inscribed* Hilbert balls, as depicted on Fig. 5.
Let $`\tau `$ be any critical cyclic permutation of $`\stackrel{~}{A}`$.
###### Proposition 17.
* In the case $`\lambda (\stackrel{~}{A})>\text{0}`$ for any $`y\text{eig}(A)`$ the distance between $`y`$ and $`\mathrm{\Gamma }(\text{eig}(A))`$ is not greater than $`\lambda (\stackrel{~}{A})^1`$;
* Let $`\mu `$ be such that $`\lambda (\stackrel{~}{A})\mu <\text{1}`$, if $`\lambda (\stackrel{~}{A})>\text{0}`$, or such that $`\text{0}<\mu <\text{1}`$, if $`\lambda (\stackrel{~}{A})=\text{0}`$. Then for any $`y\text{eig}(A_\mu )`$ the distance between $`y`$ and $`\mathrm{\Gamma }(\text{eig}(A))`$ is not less than $`\mu ^1`$;
* In the case $`\lambda (\stackrel{~}{A})>\text{0}`$, for any $`i,jK(\tau )`$ such that $`j=\tau (i)`$, and any $`y\text{eig}(A_\lambda )`$, the distance between $`y`$ and the face $`\mathrm{\Gamma }_{ij}`$ is equal to $`\lambda (\stackrel{~}{A})^1`$.
Proof. 1) The distance between $`y`$ and $`\mathrm{\Gamma }(\text{eig}(A))`$ does not exceed the minimal distance between $`y`$ and supporting planes $`X_{ij}`$ that correspond to the edges $`(i,j)`$ of the cyclic path determined by $`\tau `$, and this minimal distance is not greater than $`\lambda (\stackrel{~}{A})^1`$:
$$\begin{array}{cc}d_H(y,\mathrm{\Gamma }(\text{eig}(A)))\hfill & _{iK(\tau )}A_{i\tau (i)}^1y_{\tau (i)}^1y_i\hfill \\ \hfill & (_{iK(\tau )}A_{i\tau (i)}^1)^{\frac{1}{K(\tau )}}=\lambda (\stackrel{~}{A})^1.\hfill \end{array}$$
2) If $`y\text{eig}(A_\mu )`$ then, since $`\text{eig}(A_\mu )=\text{span}(A_\mu ^{})`$, we have
(25)
$$y=\underset{kM}{}\alpha _k(A_\mu ^{})_k,$$
where $`M=\text{supp}(\alpha )`$. Substituting (25) into (24), we get
(26)
$$d_H(y,\mathrm{\Gamma }(\text{eig}(A)))=\underset{ij,A_{ij}\text{0}}{}A_{ij}^1\underset{kM_j}{}\alpha _k^1(A_\mu ^{})_{jk}^1\underset{l}{}\alpha _l(A_\mu ^{})_{il}.$$
Here by $`M_j`$ we denote the set $`M\text{supp}((A_\mu ^{})_j)`$. Now we estimate (26) from below and use the inequalities $`A_{ij}\mu (A_\mu ^{})_{ij}`$ and $`(A_\mu ^{})_{ij}(A_\mu ^{})_{jk}(A_\mu ^{})_{ik}`$ (see (5)):
$$\begin{array}{cc}d_H(y,\mathrm{\Gamma }(\text{eig}(A)))\hfill & _{ij,A_{ij}\text{0}}_{kM_j}A_{ij}^1(A_\mu ^{})_{jk}^1(A_\mu ^{})_{ik}\hfill \\ \hfill & _{ij,A_{ij}\text{0}}_{kM_j}\mu ^1(A_\mu ^{})_{ij}^1(A_\mu ^{})_{jk}^1(A_\mu ^{})_{ik}\mu ^1.\hfill \end{array}$$
3) The distance between $`y\text{eig}(A_\lambda )=\text{span}(A_\lambda ^{})`$ and the supporting plane $`X_{ij}`$ is equal to
(27)
$$d_H(y,X_{ij})=A_{ij}^1\underset{kM_j}{}\alpha _k^1(A_\lambda ^{})_{jk}^1\underset{lM_i}{}\alpha _l(A_\lambda ^{})_{il}.$$
The cyclic permutation $`\tau `$ of $`A_\lambda `$ has the weight 1. Hence for all $`i,jK(\tau )`$ such that $`j=\tau (i)`$ we have, according to Prop. 5, that $`A_{ij}=\lambda (\stackrel{~}{A})(A_\lambda ^{})_{ij}`$ and $`(A_\lambda ^{})_{ij}(A_\lambda ^{})_{jl}=(A_\lambda ^{})_{il}`$. Note that $`(A_\lambda ^{})_{ij}\text{0}`$ for all $`i,jK`$ and therefore $`(A_\lambda ^{})_{jl}=\text{0}`$ if and only if $`(A_\lambda ^{})_{il}=\text{0}`$, i.e. $`M_i`$ and $`M_j`$ coincide. Making use of all this we write the upper estimate for $`d_H(y,X_{ij})`$:
$$\begin{array}{c}d_H(y,X_{ij})_{lM_i}A_{ij}^1(A_\lambda ^{})_{jl}^1(A_\lambda ^{})_{il}=\hfill \\ =_{lM_i}\lambda (\stackrel{~}{A})^1(A_\lambda ^{})_{ij}^1(A_\lambda ^{})_{jl}^1(A_\lambda ^{})_{il}=\lambda (\stackrel{~}{A})^1.\hfill \end{array}$$
We also have $`d_H(y,\mathrm{\Gamma }(\text{eig}(A)))\lambda (\stackrel{~}{A})^1`$ and therefore (see Prop. 16) $`d_H(y,X_{ij})=d_H(y,\mathrm{\Gamma }_{ij})=\lambda (\stackrel{~}{A})^1`$. $`\mathrm{}`$
The sets $`\mathrm{\Gamma }(\text{eig}(A_\mu ))`$, for $`\mu <\text{1}`$, are the subsets of $`\text{eig}(A)`$ equidistant from $`\mathrm{\Gamma }(\text{eig}(A))`$, as Prop. 18 suggests.
###### Proposition 18.
For all $`\mu `$ such that $`\lambda (\stackrel{~}{A})\mu <\text{1}`$, if $`\lambda (\stackrel{~}{A})>\text{0}`$, or such that $`\text{0}<\mu <\text{1}`$, if $`\lambda (\stackrel{~}{A})=\text{0}`$, the distance $`d_H(y,\mathrm{\Gamma }(\text{eig}(A)))`$ is equal to $`\mu ^1`$ if and only if $`y\mathrm{\Gamma }(\text{eig}(A_\mu )))`$.
Proof. If $`\lambda (\stackrel{~}{A})>\text{0}`$ and $`\mu =\lambda (\stackrel{~}{A})`$ then the statement readily follows from the observation that $`A_\lambda `$ does not have strong permanent and therefore (see Prop. 14) $`\text{eig}(A_\lambda )`$ does not have interior.
Let us consider $`\mu >\lambda (\stackrel{~}{A})`$. First, the equality
$$\underset{ij,A_{ij}\text{0}}{}A_{ij}y_jy_i^1=\mu $$
implies $`A_\mu y=y`$. So, if $`d_H(y,\mathrm{\Gamma }(\text{eig}(A)))=\mu ^1`$ then $`y\text{eig}(A_\mu )`$. Assume that $`y`$ belongs to the interior of $`\text{eig}(A_\mu )`$. Since $`A_\mu `$ is definite and has strong permanent, we can use representation (18) or (19) and obtain $`\kappa <\text{1}`$ such that $`y\text{eig}(A_{\mu \kappa })`$. Now statement 2) of Prop. 17 implies that $`d_H(y,\mathrm{\Gamma }(\text{eig}(A)))(\mu \kappa )^1>\mu `$. This is a contradiction, so $`y\mathrm{\Gamma }(\text{eig}(A_\mu ))`$.
Suppose now that $`y\mathrm{\Gamma }(\text{eig}(A_\mu ))`$. It means that there are $`ij`$ such that $`y_iy_j^1=(A_\mu )_{ij}`$, where $`(A_\mu ^{})_{ij}=(A_\mu )_{ij}`$. According to Prop. 15, this face corresponds to the face of $`\text{eig}(A)`$ determined by the entry $`A_{ij}`$, and the distance between these two faces is clearly $`\mu ^1`$. $`\mathrm{}`$
Throughout this section we dealt with eigenvectors having full support. But let us recall the third statement of Prop. 7. It says that there might be eigenvectors with nontrivial support $`K`$. The distance between these eigenvectors and part of any face with full support would be infinite. Also, these eigenvectors are eigenvectors of the submatrix $`A_{KK}`$. Therefore it is presumable, in this case, to pose the problem of finding $`d_H(y,\mathrm{\Gamma }(\text{eig}(A_{KK})))`$.
We conclude this section with two examples.
*Example 1* In the beginning of this section we considered the definite matrix
$$A=\left(\begin{array}{ccc}0& 4& 1\\ 1& 0& 1\\ 5& 7& 0\end{array}\right),$$
with the maximal cycle mean of $`\stackrel{~}{A}`$ equal to $`\lambda =1.5`$. Now we pick the following three members of the $`\{A_\mu \}`$ family:
$$A_{0.5}=\left(\begin{array}{ccc}0& 3.5& 1.5\\ 1.5& 0& 1.5\\ 4.5& 6.5& 0\end{array}\right),A_1=\left(\begin{array}{ccc}0& 3& 2\\ 2& 0& 2\\ 4& 6& 0\end{array}\right),$$
and
$$A_\lambda =\left(\begin{array}{ccc}0& 2.5& 2.5\\ 2.5& 0& 2.5\\ 3.5& 5.5& 0\end{array}\right).$$
The left-hand side of Fig. 4 displays $`\text{span}(A)`$, and the right-hand side displays the sets $`\mathrm{\Gamma }(\text{eig}(A_\mu ))`$, for $`\mu =0`$ (dark blue) $`\mu =0.5`$ (green), $`\mu =1`$ (brown), and $`\mu =\lambda =1.5`$ (red). The lines corresponding to larger values of $`\mu `$ are given smaller weight (to help distinguishing between different values of $`\mu `$ in black and white printing). We see that there is an injection of systems of inequalities describing $`\text{eig}(A_\mu )`$ into the system of inequalities describing $`\text{eig}(A)`$ in accordance with Prop. 15.
Fig. 5 displays, together with $`\text{span}(A)`$, two Hilbert balls with the maximal radius $`\lambda ^1=1.5`$ inscribed in $`\mathrm{\Gamma }(\text{eig}(A))`$. The balls touch the ‘cycle faces’ $`\{(x,y,z)\text{eig}(A)yx^1=1\}`$ and $`\{(x,y,z)\text{eig}(A)xy^1=4\}`$, in accordance with the third statement of Prop. 17.
*Example 2* Consider a Hilbert ball with radius $`d`$ centered at $`\{\lambda x\}`$ ($`\lambda `$ is any nonzero scalar). It is the set
$$Y=\{y\underset{i,j}{}x_iy_i^1y_jx_j^1d\},$$
or, equivalently,
$$Y=\{yy_iy_j^1d^1x_ix_j^1\}.$$
Denote by $`D`$ the matrix with entries $`d_{ij}=d^1x_ix_j^1`$. It is easily verified that $`D=D^{}`$. Then it follows from Prop. 11 that the Hilbert ball is the eigenspace of $`D`$ and the columns of this matrix are its generators. The maximal cycle mean of $`\stackrel{~}{D}`$ is clearly $`d^1`$. The eigenspaces $`\text{eig}(D_\mu )`$ where $`d^1<\mu \text{1}`$ are Hilbert balls with radii $`(\mu d)^1`$ centered at $`\{\lambda x\}`$, and $`\text{eig}(D_{d^1})`$ is precisely $`\{\lambda x\}`$.
For a three-dimensional example, set $`x=[\mathrm{5\hspace{0.33em}4\hspace{0.33em}0}]^T`$ and $`d=3`$. Then
$$D=\left(\begin{array}{ccc}0& 2& 2\\ 4& 0& 1\\ 8& 7& 0\end{array}\right).$$
Fig. 6 displays the sets $`\mathrm{\Gamma }(\text{eig}(D_\mu )`$ for $`\mu =0`$ (dark blue), $`\mu =1`$ (green), and $`\mu =2`$ (brown) with the same convention about the weight of lines as in Fig. 4. These sets are concentric spheres centered at $`\text{eig}(D_3)=\{\lambda x\}`$ (the large red circle in the center of Fig. 6). |
warning/0506/physics0506159.html | ar5iv | text | # First principles electron transport: finite-element implementation for nanostructures
## I Introduction
Using small nano-scale lithographic structures, atomic aggregates and even single molecules, it is possible to fabricate new kind of electronic devices Datta (1995). The function and scale of these devices is based on quantum-mechanical phenomena, and cannot be described within the classical regime. Of particular relevance are the electron transport properties of these nanoscale devices, as this will determine their effectiveness in, for example, a new generation of transistors. As the experimental work on these devices grows, increasing emphasis is placed on developing a matching theoretical description Nitzan (2001); Nitzan and Ratner (2003). Although some efforts have included a full description of an electronic circuit Luo et al. (2002); Kim et al. (2005), current research is mainly focused to study single electronic components.
Density-functional theory (DFT) is widely used in atomistic modeling of materials properties and recently also properties and phenomena in nanostructures. The power of DFT is in its capacity to treat accurately systems with a hundreds of atoms, yet retain a full quantum-mechanical treatment. Although the full justification of use of the DFT in electron transport calculations is debated Burke et al. (2005); Thygesen and Jacobsen (2004) we adopt it as a practical scheme to describe the real systems and devices.
In the Kohn-Sham scheme of DFT the electron density is calculated using single-particle wave functions. The explicit use of the wave functions in constructing the density suffices well in two kinds of systems. Either the system has a repeating structure so that it can be modeled with periodic boundary conditions or the system is so small that it can be calculated as a whole. In nano electronics, however, a system consists usually of a small finite part, the nanostructure, which is connected to the surrounding infinite leads. If one enforces periodic boundary conditions even a large repeating super cell or calculation volume can cause finite-size effects with spurious results for electron transport.
A commonly used solution to this problem, which we have also employed, is to combine DFT with the Green’s function formalism Datta (1995). The Green’s functions are first constructed for the semi-infinite leads by using the analytically known or easily calculated wave functions. Once the Green’s function for the combined nanostructure and leads is constructed, the wave functions are no longer needed explicitly. This makes it possible to use open boundary conditions between the nanostructure and the lead. In this way we have an effectively infinite system without periodicity, making the finite-size effects small. It is also possible to calculate the electric current through the system for a finite bias voltage between the leads in a self-consistent manner with the electron density. The ensuing model for the current is analogous to the Landauer-Bütteker model Nitzan (2001). We have used non-local pseudopotentials for modeling atoms, and the ideal metal “jellium” model for the leads. The charge density in the leads can be varied according to the conducting properties of the leads we wish to model.
The use of Green’s functions instead of the explicit use of wave functions is computationally demanding. This is why special care has to be taken in choosing the numerical methods. The first implementations used tight-binding methods Sautet and Joachim (1988); Chico et al. (1996), but a more typical solution is to expand the Green’s functions in a special basis tailored for the system. Common examples are localized atomic orbitals Taylor et al. (2001); Brandbyge et al. (2002), an $`O(N)`$ optimized basis Nardelli et al. (2001), a wavelet basis Thygesen et al. (2003), full-potential linearized augmented plane-waves Wortmann et al. (2002), maximally localized Wannier functions Calzolari et al. (2004), a finite-difference method Khomyakov and Brocks (2004), and linear a finite-element method Polizzi and Ben (2005). Our solution is to use the finite-element method (FEM). It allows a systematic error control which is especially important in transport problems as there are many different properties which must be monitored. For example, the pole of the Green’s function can cause numerical problems. According to our experience electronic tunneling in particular is sensitive to numerical accuracy.
Besides systematic error control, the FEM has also other good properties which makes it a natural method for transport problems. It is a flexible method which allows one to take into account the geometry of the nano device exactly. Special boundary conditions are easy to derive without mixing the model with the numerical method and their implementation is straightforward. Moreover, the local nature of the basis produces sparse matrices for which efficient solving methods exist. Varying the size of the elements can be used to reduce the number of the basis functions and, consequently, the size of the system as compared to the finite-difference method. This is especially true for the high-order $`p`$-method. Finally, there exists a lot of theoretical work together with tested and reliable tools, such as mesh generators and optimized linear solvers. These are used as standard building blocks for any FEM implementation granting easy access to state-of-the-art algorithms. Using the FEM new theoretical or numerical ideas are easy to implement and test.
The structure of the paper is as follows: in Sec. II we describe the model itself in detail, including a discussion of the formalism of our implementation; in Sec. III we apply the model to two example systems, a Na atom chain and HfO<sub>2</sub>-Si interface between two leads. In Sec. IV we summarize the work. In this paper we use atomic units in all equations.
## II Model
The schematic picture of our model is shown in Fig. 1. Actually, the figure present our second test case, the HfO<sub>2</sub>-Si interface between two leads. We have an atomistic nanostructure between two semi-infinite leads. The system is divided into three parts, $`\mathrm{\Omega }`$ being the calculation volume, and $`\mathrm{\Omega }_L`$ and $`\mathrm{\Omega }_R`$ are left and right leads, respectively. The boundaries $`\mathrm{\Omega }_{L/R}`$ are open so that electrons can penetrate through them without any reflection or refraction. We use the DFT to model electron interactions. The basic quantity, the electron density, is calculated from single-particle Green’s functions. Then we use the density to calculate the effective potential as
$$V_{\mathrm{eff}}=V_{\mathrm{ext}}+V_\mathrm{c}+V_{\mathrm{xc}}+V_{\mathrm{bias}}+\widehat{V}_{\mathrm{nl}},$$
(1)
where $`V_{\mathrm{ext}}`$ is the external potential caused by positive background charges, local parts of the pseudopotential operators and the potential outside potential barriers. $`V_\mathrm{c}`$ is the Coulomb Hartree interaction part, and $`V_{\mathrm{xc}}`$ is the exchange-correlation part which we calculate using the local-density approximation parametrized by Perdew and Zunger Ceperley and Alder (1980); Perdew and Zunger (1981). $`V_{\mathrm{bias}}`$ sets the boundary conditions if a bias voltage is applied. $`\widehat{V}_{\mathrm{nl}}`$ is the nonlocal part of the pseudopotential operators.
The Hartree potential is calculated from the modified Poisson equation
$$^2V_c^ik_\mathrm{P}^2V_c^i=4\pi (\rho _+\rho )k_\mathrm{P}^2V_c^{i1},$$
(2)
where $`k_\mathrm{P}`$ is an adjustable parameter. $`k_\mathrm{P}`$ does not affect the final self-consistent result, but the stability and convergence of iterations are improved Arponen et al. (1973), because the Coulomb potential due to charge redistribution between adjacent iterations is screened. The non-local pseudopotential is an operator given by
$$\widehat{V}_{\mathrm{nl}}v(𝐫)=\underset{l,m}{}e_{l,m}\zeta _{l,m}(𝐫)_\mathrm{\Omega }\zeta _{l,m}(𝐫^{})v(𝐫^{})𝑑𝐫^{},$$
(3)
where $`e_{l,m}`$ and $`\zeta _{l,m}(𝐫)`$ are defined using the Troullier-Martins pseudopotentials Fuchs and Scheffler (1999); Troullier and Martins (1991). Eq. (3) uses the projection of the function $`v(𝐫)`$ (arbitrary function, which in practical calculations is a basis function) on the atomic-specific function $`\zeta _{l,m}`$ depending on the quantum numbers $`l`$ and $`m`$ corresponding to the angular momentum.
We have implemented the guaranteed-reduction-Pulay Bowler and Gillan (2000) method for the mixing of the self-consistent iterations. It uses potentials from the five previous iterations for computing a new potential in such a way that the predicted norm of the potential residue is minimized. The simplest mixing scheme in which potentials are mixed with a linear feed-back coefficient does not work well in open systems. The calculations are rather unstable so that quite a small feed-back coefficient has to be used. This is because the net charge in the calculation volume $`\mathrm{\Omega }`$ varies during the calculations.
### II.1 Green’s function model
The details of the Green’s function method for electron transport in nanostructures are explained, for example, in Ref. Xue et al. (2002). Here we give only a short introduction to the equations to be solved. The retarded Green’s function $`G^r`$ is solved from the equation
$$\left(\omega \widehat{H}(𝐫)\right)G^r(𝐫,𝐫^{};\omega )=\delta (𝐫𝐫^{}),$$
(4)
where $`\omega `$ is the electron energy and $`\widehat{H}`$ is the DFT Hamiltonian of the system,
$$\widehat{H}(𝐫)=\frac{1}{2}^2+V_{\mathrm{eff}}(𝐫).$$
(5)
When we know $`G^r`$ we can calculate the so-called lesser Green’s function $`G^<`$. In the equilibrium when no bias voltage is applied over the nanostructure it is obtained from
$$G^<(𝐫,𝐫^{};\omega )=2f_{L/R}(\omega )G^r(𝐫,𝐫^{};\omega ),$$
(6)
where $`f_{L/R}`$ are the Fermi functions of the leads. In the equilibrium, $`f_L=f_R`$. For a finite bias voltage $`f_{L/R}`$ we take into account the bias and a more complicated equation for $`G^<`$ has to be used. To obtain it we write Eq. (4) in the form
$$\left(\omega \widehat{H}_0\mathrm{\Sigma }_L^r(\omega )\mathrm{\Sigma }_R^r(\omega )\right)G^r(𝐫,𝐫^{};\omega )=\delta (𝐫𝐫^{}),$$
(7)
where $`\widehat{H}_0`$ is the Hamiltonian of the isolated volume $`\mathrm{\Omega }`$ and $`\mathrm{\Sigma }_{L/R}^r`$ are the so-called self-energies of the leads. We also define
$`i\mathrm{\Gamma }_{L/R}`$ $`=\mathrm{\Sigma }_{L/R}^r\mathrm{\Sigma }_{L/R}^a=2i\mathrm{Im}(\mathrm{\Sigma }_{L/R}^r),`$ (8)
and can solve $`G^<`$ for a finite bias voltage as
$`G^<`$ $`(𝐫,𝐫^{};\omega )=`$ (9)
$`if_R(\omega ){\displaystyle _{\mathrm{\Omega }_R}}{\displaystyle _{\mathrm{\Omega }_R}}G^r(𝐫,𝐫_R;\omega )\mathrm{\Gamma }_R(𝐫_R,𝐫_R^{};\omega )`$
$`\times G^a(𝐫_R^{},𝐫^{};\omega )d𝐫_Rd𝐫_R^{}`$
$`if_L(\omega ){\displaystyle _{\mathrm{\Omega }_L}}{\displaystyle _{\mathrm{\Omega }_L}}G^r(𝐫,𝐫_L;\omega )\mathrm{\Gamma }_L(𝐫_L,𝐫_L^{};\omega )`$
$`\times G^a(𝐫_L^{},𝐫^{};\omega )d𝐫_Ld𝐫_L^{}.`$
The first and second terms correspond to electrons originating from the right- and left leads, respectively. The electron density is calculated from
$$\rho (𝐫)=\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}\mathrm{Im}(G^<(𝐫,𝐫;\omega ))𝑑\omega $$
(10)
and the tunneling probability from
$`T(\omega )=`$ $`{\displaystyle _{\mathrm{\Omega }_L}}{\displaystyle _{\mathrm{\Omega }_L}}{\displaystyle _{\mathrm{\Omega }_R}}{\displaystyle _{\mathrm{\Omega }_R}}\mathrm{\Gamma }_L(𝐫_L,𝐫_L^{};\omega )G^r(𝐫_L^{},𝐫_R;\omega )`$ (11)
$`\times `$ $`\mathrm{\Gamma }_R(𝐫_R,𝐫_R^{};\omega )G^a(𝐫_R^{},𝐫_L;\omega )d𝐫_Ld𝐫_L^{}d𝐫_Rd𝐫_R^{},`$
from the values of the functions at the boundaries $`\mathrm{\Omega }_{L/R}`$. Finally the current is determined as
$$I=\frac{1}{\pi }_{\mathrm{}}^{\mathrm{}}T(\omega )\left(f_L(\omega )f_R(\omega )\right)𝑑\omega .$$
(12)
We use the FEM in the numerical implementation. Therefore we first cast Eq. (4) in the variational form with open boundary conditions (for the derivation, see Ref. Havu et al. (2004)) as
$`{\displaystyle _\mathrm{\Omega }}`$ $`\{v(𝐫){\displaystyle \frac{1}{2}}G^r(𝐫,𝐫^{};\omega )`$ (13)
$`+v(𝐫)[\omega V_{\mathrm{eff}}(𝐫)]G^r(𝐫,𝐫^{};\omega )\}d𝐫`$
$`<\widehat{\mathrm{\Sigma }}_LG^r,v><\widehat{\mathrm{\Sigma }}_RG^r,v>`$
$`=`$ $`v(𝐫^{}),`$
where the self energy-operators
$`<\widehat{\mathrm{\Sigma }}_LG^r,v>={\displaystyle _{\mathrm{\Omega }_L}}{\displaystyle _{\mathrm{\Omega }_L}}{\displaystyle \frac{1}{4}}G^r(𝐫_L^{},𝐫^{};\omega )`$ (14)
$`\times {\displaystyle \frac{^2g_e(𝐫_L^{},𝐫_L;\omega )}{𝐧_L𝐧_L^{}}}v(𝐫_L)d𝐫_L^{}d𝐫_L.`$
Above, $`g_e`$ is the Green’s function of the semi-infinite lead in the domain $`\mathrm{\Omega }_{L/R}`$ with the zero-value condition on the boundary $`\mathrm{\Omega }_{L/R}`$. In our implementation the leads are described by a uniform positive background charge and therefore $`g_e`$ can be calculated partly analytically. Thus our model means that the leads are of some kind of ideal generic metals. The important interface between the nanostructure, e.g. a molecule, and the actual metallic lead can be described accurately by including some lead metal atoms in the computational domain $`\mathrm{\Omega }`$. It is also possible to use fully atomistic leads by calculating numerically $`g_e`$ for them.
Note that the Eqs. (13) and (14) are analogous to those derivations of the open boundary conditions in which truncated matrices Datta (1995) are used. In the continuum limit these two forms give the same results. However, the weak form is more natural in the FEM formulation and more suitable for theoretical purposes when analyzing nonlinear partial differential equations. It is also straightforward to use and the error control is systematic. Note that this formulation can be used with any continuous basis set, not only with the FEM. In the context of basis set methods, the weak formulation case is known as the Galerkin method. In practice the Green’s functions are approximated with respect to this basis so that
$$G^r(𝐫,𝐫^{};\omega )\underset{i,j=1}{\overset{N}{}}g_{ij}(\omega )\varphi _i(𝐫)\varphi _j(𝐫^{}).$$
(15)
The coefficients $`g_{ij}(\omega )`$ can be solved from (13) by choosing $`v=\varphi _k`$ and evaluating the equations.
### II.2 Finite-element p-basis
In the FEM we partition the calculation volume to (in our case, tetrahedral) sub-domains called elements and the basis functions $`\varphi _i`$ are constructed using globally continuous (but not necessarily continuously differentiable) piecewise polynomials with respect to the finite element mesh. This gives both unique flexibility of the approximating functions as well as completeness of the basis with respect to almost any norm. Each basis function has a support that is concentrated to only a few neighboring elements. This makes the basis local and results in sparse system matrices.
There are several options of how to choose the finite-element basis and one has to be careful in achieving acceptable accuracy. The simplest basis is the linear one. It is easy to implement and works well, especially in systems with rapidly varying functions. A typical improvement to this basis is to use node-based higher-order elements. These elements converge faster to a smooth solution than the linear ones. However, practically only relatively low orders, two and three, can be used because of numerical stability problems.
In this work we have used so-called hierarchical $`p`$-elements. They also span higher-order polynomials, but the choice of the local basis ensures that stability problems do not appear. This is because the basis functions are chosen so that their derivatives are close to orthogonal in the $`L_2`$-norm. The hierarchical nature also makes it easy to change the order of the basis from element to element within the same mesh.
The actual FEM implementation consists of a reference element and reference basis that are mapped separately to each of the elements of the mesh. Our reference element is a tetrahedron with nodes at the coordinates $`1:(1,0,0)`$, $`2:(1,0,0)`$, $`3:(0,\sqrt{3},0)`$ and $`4:(0,\frac{1}{\sqrt{3}},2\sqrt{\frac{2}{3}})`$. One can easily show that there exists an affine map taking the reference element to any of the tetrahedron. The order of our basis is $`p`$ meaning that in each element polynomials of the order of $`p`$ are employed. The basis is constructed hierarchically. First, there are four linear node basis functions inside the elements sharing a common node. In the reference element they are
$`L_1`$ $`={\displaystyle \frac{1}{2}}\left(1\xi {\displaystyle \frac{\mu }{\sqrt{3}}}{\displaystyle \frac{\zeta }{\sqrt{6}}}\right)`$ (16)
$`L_2`$ $`={\displaystyle \frac{1}{2}}\left(1+\xi {\displaystyle \frac{\mu }{\sqrt{3}}}{\displaystyle \frac{\zeta }{\sqrt{6}}}\right)`$
$`L_3`$ $`={\displaystyle \frac{1}{\sqrt{3}}}\left(\mu {\displaystyle \frac{\zeta }{\sqrt{8}}}\right)`$
$`L_4`$ $`=\sqrt{{\displaystyle \frac{3}{8}}}\zeta `$
where $`\xi `$ $`\mu `$ and $`\zeta `$ are the cartesian coordinates of the reference element. Secondly, for $`p>`$1 we have $`6(p1)`$ edge functions. E.g. for the edge between the nodes 1 and 2
$`N_{i1}^{(1,2)}=L_1L_2\phi _i(L_2L_1),i=2,\mathrm{}p.`$ (17)
Here one usually sets
$`\phi _i(\xi )`$ $`={\displaystyle \frac{4\varphi _i(\xi )}{1\xi ^2}}`$ (18)
$`\varphi _i(\xi )`$ $`=\sqrt{{\displaystyle \frac{1}{2(2i1)}}}(P_i(\xi )P_{i2}(\xi ))`$
Above $`P_i`$ is the Legendre polynomial of the order of $`i`$. Third, we have $`2(p1)(p2)`$ face functions. For example for the face between the nodes 1,2, and 3 they are
$`N_{i,j}^{(1,2,3)}=L_1L_2L_3P_i(L_2L_1)P_j(2L_31),`$ (19)
$`i,j=0,\mathrm{},p3,i+j=0,\mathrm{},p3.`$
Fourth, we have $`\frac{1}{6}(p1)(p2)(p3)`$ bubble functions, which are supported only in a single element each. These are
$`N_{i,j,k}=L_1L_2L_3L_4P_i(L_2L_1)\times `$ (20)
$`P_j(2L_31)P_k(2L_41)`$
$`i,j,k=0,\mathrm{},p4,i+j+k=0,\mathrm{},p4.`$
When $`p`$-elements are used one must take care of the continuity of the basis. This is because, for example, the local basis function $`N_3^{(1,2)}`$ has an orientation on the boundary. The basis includes the function $`\phi (L_2L_1)`$, not $`\phi (L_1L_2)`$, which would be another possibility. This means that all the edges in the mesh have to have information about the direction. Otherwise there is very likely a continuity problem on some boundaries. In practice, for tetrahedral elements the orientation problem can be handled for arbitrary finite-element meshes using only two reference elements Ainsworth and Coyle (2003).
The benefits of selecting the basis described above are rather clear. The polynomial basis is very easy to realize and has good approximating properties. For smooth solutions the $`p`$-basis is known to give exponential convergence rates with respect to the number of basis functions used. In the DFT methods the theory is typically developed to the direction that the solutions are as smooth as possible. For example, pseudopotential operators are designed so that they produce as smooth an electron potential as smooth as possible. This is because the plane-wave basis set needs smooth solutions in order to work efficiently. On the other hand, in the case of non-smooth solutions one can benefit from the piecewise nature of the FEM basis allowing one to approximate even singular solutions to some extent. Moreover, the finite element mesh can be refined in regions where solution changes rapidly. When modeling molecules there is also a lot of empty space in the calculation domain. It is then practical to use large elements in the empty space and smaller ones near atoms.
### II.3 Linear algebra methods
The use of the Green’s function method is computationally demanding in comparison to explicit wave-function methods. Since the main computational burden of our method is to find a subset of the coefficients of the Green’s function in question, a special consideration must be taken when choosing the methods of linear algebra to be used. The eigenvalue problems in explicit wave-function methods are typically solved by iterative methods. In our case it is better to use direct solvers, because a set of linear equations needs to be solved. We have opted for the frontal method widely used in the solution of sparse linear systems Reid and Duff (1983); Duff (1996) and extremely suitable for finite-element matrices. The actual implementation is ME47 of the Harwell Subroutine Library (HSL) hsl (see Refs. Gupta (2000); Davis (2004); Demmel et al. (1999) for other similar approaches). In the frontal method, one first finds a permutation of the sparse matrix aiming to minimize the fill-in resulting from the factorization process. Next, a LU-decomposition (or Cholesky-decomposition for symmetric problems) of the matrix $`A`$ is found, and finally, two systems with triangular coefficient matrices, $`Lz=b`$ and $`Ux=z`$ (where $`U=L^T`$ for symmetric problems) are solved. To find all the required coefficients of the solution we must vary the right-hand side $`b`$ of the system.
For a three-dimensional problem the size of the linear system can grow so large that the CPU-time and memory requirements of different systems have to be addressed. The main question is how large systems can be calculated using these methods so that the calculation time for a single self-consistent iteration is not too large. Currently a system of several tens of thousands of unknowns can be solved in a commodity-CPU cluster environment.
In detail, the Green’s function method includes a computation of the elements for the inverse of a sparse matrix, so that the calculation time requirements increase relatively fast with the system size. A classical complexity result for the solution (and inversion) of a general $`N\times N`$ system with a direct method is $`O(N^3)`$. However, for sparse systems and modern frontal methods this bound is too pessimistic George and Liu (1981). The CPU time requirement depends on the fill-in of the inversion problem. For very simple cases one can show that the key statistic of the problem, the number of non-zeros (nnz) present in the factors $`L`$ and $`U`$, satisfies $`nnz(L)nnz(U)=O(N\mathrm{log}(N))`$ George and Liu (1981). Then the solution of each of the systems requires $`O(nnz(L)+nnz(R))`$ floating-point operations, and in the worst case we must solve these with $`N`$ different right-hand sides effectively giving us the inverse of the matrix $`A`$, so that the total cost is $`O(N(nnz(L)+nnz(R)))`$. However, in modern computer systems the complexity is not the only relevant measure since the performance may be highly nonlinear (see, e.g. atl ; Goto and van de Geijn (2002) for an example on BLAS-tuning).
Another topic related directly to the performance of modern computer systems is the relation between processor power and memory bandwidth. This is especially true for the computation of the Green’s function where the actual bottleneck is the lack of available memory bandwidth in commodity-based cluster systems used in calculations, not the floating-point performance of the processor itself.
It is likely that a better performance can be achieved by upgrading several parts of the algorithms. First, the current parallel solver is implemented using the Message Passing Interface (MPI) mpi . However, in Symmetric multiprocessor (SMP) systems it is likely that well-designed OpenMP ope (or similar) parallelism would reduce the need for data transfer and thus increase performance. It would also decrease the memory requirements of the problems. Second, at the moment the solution of the Green’s function is computed varying one vector on the right-hand side at a time. A better performance could be obtained if the equations could be solved for multiple right-hand sides at a time allowing the use of BLAS3 routines. Finally, it is likely that computations would benefit from a computer system having a larger memory bandwidth than our present commodity-based one.
## III Example systems
### III.1 Atomic wire
Using the atomic force microscope or the mechanically controlled break junction technique, a chain of atoms can be made of certain metals Agrait et al. (2003). It has been observed that the conductances of atom chains vary as a function of the number of atoms in the chain Smit et al. (2003). The conductances of these systems have been studied also theoretically in several works. In order to benchmark our results against other calculations, we use Na-atom chains as test systems. They have been simulated in several previous studies Sim et al. (2001); Lang (1997); Tsukamoto and Hirose (2002); Lee et al. (2004) using different models. According to these calculations the conductances of the wires show even-odd oscillation as a function of the number of atoms in the wire.
In our setup, the Na-atom chain is located between two leads, with the lead shape defined by a 70 cone angle (see Fig. 2). We consider two different connections of the atom chain to the electrodes. In model A we have just three Na atoms between the jellium leads. This resembles closely the system used in Ref. Lang (1997). In model B, there are four Na atoms at the tips of the leads in a square form. This makes the connection between the atom chain and the leads more realistic. This kind of structure is modeled also in Ref. Tsukamoto and Hirose (2002).
The conductances as a function of the number of chain atoms for systems A and B are shown in Fig. 3. In the Na-atom chain, electrons have only one conducting mode so that the conductance can be one conductance quantum $`2e^2/h`$ at maximum. Both systems A and B exhibit conductance oscillations as a function of the number of atoms. These oscillations arise from resonance states in the atom chain. Depending on the position of the resonances relative to Fermi-level the conductance has either a maximum or minimum value, so that the maxima and minima correspond to approximately half and fully occupied resonance states, respectively. The oscillation is within the range of 0.9 - 1.0 $`\times 2e^2/h`$ for system A and 0.6 - 1.0 $`\times 2e^2/h`$ for system B. The difference between the oscillation amplitudes is due to different strengths of the connection of the chains to the leads. System B has a weaker coupling to the leads than system A. Weak connections make the resonances also sharper, as is seen in the tunneling probability in Fig. 4. In contrast to Ref. Havu et al. (2002), we do not see a strong lead-shape dependence in the conductance. The widening of the cone angle lowers the conductance as the edges of the wire become sharper.
The electron tunneling probabilities through chains of three- and four- atom systems A and B are shown in Fig. 4. The probability function $`T(\omega )`$ is defined in Eq. (11). The conductance of the system in the zero-bias limit can be read at the Fermi-level. Here, as well as in Fig 3, we see that the conductance oscillations for systems A and B are in a different phase. This is because in system B the atom chain is effectively shorter than in A, as the first and the last chain atom are partly inside the square of the four Na atoms.
When we compare the conductance oscillations of system A (see Fig. 3) to those in Ref. Lang (1997) obtained by using semi-infinite jellium leads with planar surfaces ($`\alpha =180^{}`$), we see that the even-odd oscillations in the conductance are in the same phase. In the case of system B we can directly compare the tunneling probability of Fig. 4 with those in Ref. Tsukamoto and Hirose (2002) where the atom chain is connected also through a square of four Na atoms to jellium. The phase and the amplitude of the conductance oscillations of these results are in good agreement with our values in Fig. 3. Now that we have satisfied ourselves that the method provides a good model for electron transport we can consider a more interesting and demanding example.
### III.2 Thin insulating layer
The general increase in the performance of microelectronic devices in the past few decades has been made possible by continuous transistor scaling - based on a reduction in the thickness of the gate dielectric in typical metal-oxide-semiconductor field-effect transistors (MOSFET). At present the process has reached a bottleneck, as further reduction leads to a large increase in leakage current due to direct tunneling across the thin silicon dioxide (SiO<sub>2</sub>) layer. Several possible approaches to resolve this are being considered Kingon et al. (2000), but retaining conventional MOSFET design remains an economically attractive choice, and a leading option is just to replace SiO<sub>2</sub> with another oxide of higher dielectric constant (high-$`k`$). A high-$`k`$ oxide would provide higher effective capacitance to a comparable SiO<sub>2</sub> layer, hence allowing thicker layers to be used to reduce losses due to tunneling. The specific choice of oxide is determined by a set of requirements Huff et al. (2003) based on both the intrinsic properties of the grown oxide and its integration into the fabrication process, and at present hafnium oxide (HfO<sub>2</sub>) remains a leading candidate.
In order to study to transport properties of thin HfO<sub>2</sub> films we have simulated the growth of the oxide on a silicon surface via first principles molecular dynamics Hakala and Foster (2005). Here we consider three model interfaces: (i) a nonstoichiometric oxide interface (C), which is basically metallic due to Hf-Hf and Hf-Si bonds across the interface; (ii) a stoichiometric oxide interface ( D), which has a localized state in the band gap due to a few Hf-Hf bonds; (iii) a more idealistic interface (E), which remains insulating if no defects are present. The last model is based on the interface used in Ref. Fonseca et al. (2003), but slightly reduced in size to make it computationally manageable Gavartin et al. (2005). These models were calculated with periodic boundary conditions with k-points on the boundaries $`\mathrm{\Omega }_{P1/P2/P3/P4}`$. The effective potentials have been calculated for systems C and D using the gamma point, and for system E, four $`k`$-points. All the tunneling probabilities $`T(\omega )`$ are calculated using four $`k`$-points, which was enough to converge the probabilities to a good accuracy.
As shown in Fig. 1, the interface models are positioned between two leads. The charge density in the leads is chosen so that in the right lead $`r_s=2`$ (electron density $`n_e=3/(4\pi r_s^3)`$), representing a metal, and in the left one $`r_s=3.1`$, representing doped-silicon - as in a standard MOSFET design.
The tunneling profiles of the systems are shown in Fig. 5. Here it is seen that systems C and D show clearly metallic behavior, with a large tunneling probability at the Fermi energy. Although in principle, the stoichiometric interface (D) has a much lower density of metallic bonds, it is clear that both in interfaces C and D around two channels dominate the transport. The localized defect state in the band gap of system D plays an equivalent role in transport to the metallic bonds in interface C.
As expected, the tunneling probability for the more ideal interface E is an order of magnitude smaller at the Fermi energy than those for interfaces C and D. Yet we also see that it remains significant - this is largely due to the structure of the interface itself Fonseca et al. (2003). Although bulk HfO<sub>2</sub> is a wide-bandgap insulator, at the interface it exists as almost tetragonal HfSiO<sub>4</sub>, and the effective band gap is actually smaller than that of bulk silicon below the interface. This means that there is a negative conduction band offset between silicon and HfO<sub>2</sub>, and no real barrier for leakage. Although some of this is caused by the underestimation of the band gap in the DFT, this also reduces the silicon band gap (although the effect is not systematic).
The poor performance of interface D can also be seen in its capacity for dropping the potential. Fig. 6 shows the potential change for 0.25 V applied bias voltage. The potential drop across HfO<sub>2</sub> is less than 0.05 eV, demonstrating that the oxide hardly perturbs the electron flow from the right lead. The potential drops fastest at the right hand side of HfO<sub>2</sub>-layer where pure HfO<sub>2</sub> exists, and much more slowly in the thin layer of SiO<sub>2</sub> formed due to diffusion of oxygen. The large drop at the lead and silicon atoms is just an artifact of the boundary conditions of the Coulomb part of the effective potential.
In the rigid band approximation (used for example in Ref. Fonseca et al. (2003)) it is assumed that the shape of the tunneling probability stays constant and is only shifted in energy so that $`T(\omega ,V_{\mathrm{bias}})=T(\omega +\eta V_{\mathrm{bias}})`$, where $`\eta `$ is the ratio of potential drop at the other end of the nanostructure to the total drop over the nanostructure. In Fig. 7 we have studied how well this approximation works for interface D. The curves are plotted so that the zero-bias Fermi level is in the middle of the left and right Fermi levels of the biased interface. This corresponds to the symmetric case with $`\eta =0.5`$. We see that the tunneling probability curves roughly coincide. This indicates that potential drops symmetrically over the nanostructure and the rigid-band approximation gives a rather reasonable result.
The above results show, in agreement with previous calculations Fonseca et al. (2003) that tunneling through a more ideal, insulating interface is still significant due to a negative band offset with silicon. Since the only HfO<sub>2</sub> interfaces providing significant band offsets to silicon were built very idealistically (i.e. assuming no significant atom migration nor interfacial SiO<sub>2</sub> growth) Fiorentini and Gulleri (2002); Peacock and Robertson (2004), this indicates that fabricating a *good* interface directly between silicon and HfO<sub>2</sub> is very difficult. A more viable alternative maybe to sacrifice somewhat in dielectric constant, and grow HfO<sub>2</sub> onto a pre-existing SiO<sub>2</sub> layer. These possibilities will be explored in more detail in a further work Hakala and Foster (2005).
## IV Conclusions
In this paper we present a finite-element implementation of the non-equilibrium Green’s function method which is combined to the density-functional theory. Although the Green’s function method is computationally demanding, we demonstrate that by using hierarchical $`p`$-elements, large, physically relevant systems become tractable. More importantly, our method offers a much more rigorous control of accuracy than is usually possible in transport calculations.
We demonstrate the functionality of our implementation with two kinds of systems, the sodium atom chain wire and the silicon-HfO<sub>2</sub> interface. For the atom chain, we show that the method reproduces the previous results of other Green’s function transport methods. This gives us confidence to apply it to the more complex system: a thin layer of hafnium oxide on a silicon substrate. Here we show that the transport properties are an even more sensitive indicator of the role of defects than the electronic structure. Comparison of stoichiometric and non-stoichiometric HfO<sub>2</sub> oxide layers demonstrates that even one or two defects in a stoichiometric interface can result in tunneling comparable to that of a fully metallic non-stoichiometric interface.
###### Acknowledgements.
We are grateful to J. L. Gavartin and L. R. C. Fonseca for providing us with access to their interface structures, and for helpful discussions. We acknowledge the generous computer resources from the Center for Scientific Computing, Espoo, Finland. This research has been supported by the Academy of Finland through its Centers of Excellence Program (2000-2005). We have used the Harwell Subroutine Library in our calculations. |
warning/0506/math0506176.html | ar5iv | text | # Hamiltonian diffeomorphisms of toric manifolds
## 1. Introduction
A loop $`\psi `$ in the group $`\text{Ham}(M,\omega )`$ of Hamiltonian symplectomorphisms of the symplectic manifold $`(M^{2n},\omega )`$ determines a Hamiltonian fibration $`E\stackrel{\pi }{}S^2`$ with standard fibre $`M`$. On the total space $`E`$ we can consider the first Chern class $`c_1(VTE)`$ of the vertical tangent bundle of $`E`$. Moreover on $`E`$ is also defined the coupling class $`cH^2(E,)`$ . This class is determined by the following properties:
i) $`i_p^{}(c)`$ is the cohomology class of the symplectic structure on the fibre $`\pi ^1(p)`$, where $`i_p`$ is the inclusion of $`\pi ^1(p)`$ in $`E`$ and $`p`$ is an arbitrary point of $`S^2`$.
ii) $`c^{n+1}=0`$.
These canonical cohomology classes determine the characteristic number
(1.1)
$$I_\psi =_Ec_1(VTE)c^n.$$
$`I_\psi `$ depends only on the homotopy class of $`\psi `$. Moreover $`I`$ is an $``$-valued group homomorphism on $`\pi _1(\text{Ham}(M,\omega ))`$, so the non vanishing of $`I`$ implies that the group $`\pi _1(\text{Ham}(M,\omega ))`$ is infinite. That is, $`I`$ can be used to detect the infinitude of the corresponding homotopy group. Furthermore $`I`$ calibrates the Hofer’s norm $`\nu `$ on $`\pi _1(\text{Ham}(M,\omega ))`$ in the sense that $`\nu (\psi )C|I_\psi |`$, for all $`\psi `$, where $`C`$ is a positive constant .
In we gave an explicit expression for the value of the characteristic number $`I_\psi `$. This value can be calculated if one has a family of local symplectic trivializations of $`TM`$ at one disposal, whose domains cover $`M`$ and are fixed by the $`\psi _t`$’s. In fact we proved the following Theorem
###### Theorem 1.
Let $`\psi :S^1\text{Ham}(M,\omega )`$ be a closed Hamiltonian isotopy. If $`\{B_1,\mathrm{},B_m\}`$ is a set of symplectic trivializations for $`TM`$ which covers $`M`$ and such that $`\psi _t(B_j)=B_j`$, for all $`t`$ and all $`j`$, then
$$I_\psi =\underset{i=1}{\overset{m}{}}J_i_{B_i_{j<i}B_j}\omega ^n+\underset{i<k}{}N_{ik},$$
where
$$N_{ik}=n\frac{i}{2\pi }_{S^1}𝑑t_{A_{ik}}(f_t\psi _t)(d\mathrm{log}r_{ik})\omega ^{n1},$$
$`A_{ik}=(B_i_{r<k}B_r)B_k`$, $`J_i`$ is the Maslov index of $`(\psi _t)_{}`$ in the trivialization $`B_i`$ and $`r_{ik}`$ the corresponding transition function of $`\text{det}(TM)`$.
The homotopy type of $`\text{Ham}(M,\omega )`$ is completely known in a few particular cases only. When $`M`$ is a surface, $`\text{Diff}_0(M)`$ (the arc component of the identity map in the diffeomorphism group of $`M`$) is homotopy equivalent to the symplectomorphism group of $`M`$, so the topology of the groups $`\text{Ham}(M)`$ in dimension $`2`$ can be deduced from the description of the diffeomorphism groups of surfaces given in (see ). On the other hand, positivity of the intersections of $`J`$-holomorphic spheres in $`4`$-manifolds have been used in to prove results about the homotopy type of $`\text{Ham}(M)`$, when $`M`$ is a ruled surface. But these arguments which work in dimension $`2`$ or dimension $`4`$ cannot be generalized to higher dimensions.
Let $`𝒪`$ be a coadjoint orbit of a Lie group $`G`$. If $`G`$ is semisimple and acts effectively on $`𝒪`$, McDuff and Tolman have proved that the inclusion $`G\text{Ham}(𝒪)`$ induces an injection from $`\pi _1(G)`$ to $`\pi _1(\text{Ham}(𝒪))`$ . This result answers a question posed in . In we gave a lower bound for $`\mathrm{}\pi _1(\text{Ham}(𝒪))`$, when $`𝒪`$ is a quantizable coadjoint orbit of a compact Lie group. In particular we proved that $`\mathrm{}\pi _1(\text{Ham}(P^n))n+1`$.
In this note we use Theorem 1 to prove that $`\pi _1(\text{Ham}(M))`$ contains an infinite cyclic subgroup, when $`M`$ is a particular toric manifold. More precisely, when $`M`$ is the 6-manifold associated to the polytope obtained truncating the tetrahedron of $`^3`$ with vertices $`(0,0,0),(\tau ,0,0),(0,\tau ,0),(0,0,\tau )`$ by a horizontal plane ; that is, when $`M`$ is the one point blow up of $`P^3`$. Moveover we will give a sufficient condition for $`\pi _1(\text{Ham}(M))`$ to contain an infinite cyclic subgroup, when $`M`$ is a general toric manifold.
The paper is organized as follows. Section 2 is concerned with the determination of $`I_\psi `$ for a natural circle action on the one point blow up of $`P^3`$. In Section 3 we generalize the arguments developed in Section 2 to toric manifolds. From this generalization it follows the aforesaid sufficient condition for the existence of an infinite subgroup in $`\pi _1(\text{Ham}(M))`$, when $`M`$ is a toric manifold. Finally we check that this sufficient condition does not hold for $`P^n`$ with $`n=1,2`$. This is consistent with the fact that $`\pi _1(\text{Ham}(P^n))`$ is finite for $`n=1,2`$.
I thank Dusa McDuff for her enlightening comments.
## 2. Hamiltonian group of the one point blow up of $`P^3`$
Given $`\tau ,\mu _{>0}`$, with $`\mu <\tau `$, let $`M`$ be the following manifold
(2.1)
$$M=\{z^5:|z_1|^2+|z_2|^2+|z_3|^2+|z_5|^2=\tau /\pi ,|z_3|^2+|z_4|^2=\mu /\pi \}/𝕋^2,$$
where the action of $`𝕋^2`$ is defined by
(2.2)
$$(a,b)(z_1,z_2,z_3,z_4,z_5)=(az_1,az_2,abz_3,bz_4,az_5),$$
for $`a,bS^1`$.
$`M`$ is a toric $`6`$-manifold; more precisely, it is the toric manifold associated to the polytope obtained truncating the tetrahedron of $`^3`$ with vertices
$$(0,0,0),(\tau ,0,0),(0,\tau ,0),(0,0,\tau )$$
by a horizontal plane through the point $`(0,0,\lambda )`$, with $`\lambda :=\tau \mu `$ .
For $`0z_j`$ we put $`z_j=\rho _je^{i\theta _j}`$, with $`|z_j|=\rho _j`$. On the set of points $`[z]M`$ with $`z_i0`$ for all $`i`$ one can consider the coordinates
(2.3)
$$(\frac{\rho _1^2}{2},\phi _1,\frac{\rho _2^2}{2},\phi _2,\frac{\rho _3^3}{2},\phi _3),$$
where the angle coordinates are defined by
(2.4)
$$\phi _1=\theta _1\theta _5,\phi _2=\theta _2\theta _5,\phi _3=\theta _3\theta _4\theta _5.$$
Then the standard symplectic structure on $`^5`$ induces the following form $`\omega `$ on this part of $`M`$
$$\omega =\underset{j=1}{\overset{3}{}}d\left(\frac{\rho _j^2}{2}\right)d\phi _j.$$
Let $`0<ϵ<<1`$, we write
$$B_0=\{[z]M:|z_j|>ϵ,\text{for all}j\}.$$
For a given $`j\{1,2,3,4,5\}`$ we set
$$B_j=\{[z]M:|z_j|<2ϵ\text{and}|z_i|>ϵ,\text{for all}ij\}$$
The family $`B_0,\mathrm{},B_5`$ is not a covering of $`M`$, but if $`[z]B_k`$, then there are $`i,j`$, with $`ij`$ and $`|z_i|<ϵ>|z_j|.`$
On $`B_0`$ we will consider the well-defined Darboux coordinates (2.3). On $`B_1`$, $`\rho _j0`$ for $`j1`$; so the angle coordinates $`\phi _2`$ and $`\phi _3`$ of (2.4) are well-defined. We put $`x_1+iy_1:=\rho _1e^{i\phi _1}`$. In this way we take as symplectic coordinates on $`B_1`$
$$(x_1,y_1,\frac{\rho _2^2}{2},\phi _2,\frac{\rho _3^3}{2},\phi _3).$$
We will also consider the following Darboux coordinates: On $`B_2`$
$$(\frac{\rho _1^2}{2},\phi _1,x_2,y_2,\frac{\rho _3^2}{2},\phi _3),\text{with}x_2+iy_2:=\rho _2e^{i\phi _2}.$$
On $`B_3`$
$$(\frac{\rho _1^2}{2},\phi _1,\frac{\rho _2^2}{2},\phi _2,x_3,y_3),\text{where}x_3+iy_3:=\rho _3e^{i\phi _3}.$$
On $`B_4`$
$$(\frac{\rho _1^2}{2},\phi _1,\frac{\rho _2^2}{2},\phi _2,x_4,y_4),\text{with}x_4+iy_4:=\rho _4e^{i\phi _4}\text{and}\phi _4=\theta _4\theta _3+\theta _5.$$
On $`B_5`$
$$(x_5,y_5,\frac{\rho _2^2}{2},\chi _2,\frac{\rho _3^2}{2},\chi _3),$$
where
$$x_5+iy_5:=\rho _5e^{i\phi _5},\chi _2=\theta _2\theta _1,\chi _3=\theta _3\theta _1\theta _4,\chi _5=\theta _5\theta _1.$$
If $`[z_1,\mathrm{},z_5]`$ is a point of
$$M\underset{i=0}{\overset{5}{}}B_i,$$
then there are $`ab\{1,\mathrm{},5\}`$ such that $`|z_a|,|z_b|<ϵ`$. We can cover the set $`MB_i`$ by Darboux charts denoted $`B_6,\mathrm{},B_q`$ similar to the preceding $`B_i`$’s satisfying the following condition: The image of each $`B_a`$, with $`a=6,\mathrm{},q`$, is contained in a prism of $`^6`$ of the form
$$\underset{i=1}{\overset{6}{}}[c_i,d_i],$$
where at least four intervals $`[c_i,d_i]`$ have length of order $`ϵ`$.
By the infinitesimal “size” of the $`B_j`$, for $`j1`$, it turns out
(2.5)
$$_{B_j}\omega ^3=O(ϵ),\text{for}j1.$$
Let $`\psi _t`$ be the symplectomorphism of $`M`$ defined by
(2.6)
$$\psi _t[z]=[z_1e^{2\pi it},z_2,z_3,z_4,z_5].$$
Then $`\{\psi _t\}_t`$ is a loop in the group $`\text{Ham}(M)`$ of Hamiltonian symplectomorphisms of $`M`$. By $`f`$ is denoted the corresponding normalized Hamiltonian function. Hence $`f=\pi \rho _1^2\kappa `$ with $`\kappa `$ such that $`_Mf\omega ^3=0`$.
In the coordinates (2.3) of $`B_0`$, $`\psi _t`$ is the map $`\phi _1\phi _1+2\pi t`$. So the Maslov index $`J_{B_0}=0`$. It follows from (2.5) and Theorem 1
(2.7)
$$I_\psi =\underset{i<k}{}N_{ik}+O(ϵ),$$
with
$$N_{ik}=\frac{3i}{2\pi }_{A_{ik}}fd\mathrm{log}r_{ik}\omega ^2.$$
If $`[z]A_{ik}B_iB_k`$ , with $`1i<k`$, then at least the modules $`|z_a|`$ and $`|z_b|`$ of two components of $`[z]`$ are of order $`ϵ`$; so $`N_{ik}`$ is of order $`ϵ`$ when $`1i<k`$. Analogously $`N_{0k}`$ is of order $`ϵ`$, for $`k=6,\mathrm{},q`$. Hence (2.7) reduces to
(2.8)
$$I_\psi =\underset{k=1}{\overset{5}{}}N_{0k}+O(ϵ).$$
If we put
(2.9)
$$N_{0k}^{}=\frac{3i}{2\pi }_{A_{0k}^{}}fd\mathrm{log}r_{ik}\omega ^2,$$
with
$$A_{0k}^{}=\{[z]M:|z_k|=ϵ,|z_r|>ϵ\text{for all}rk\}$$
then
(2.10)
$$N_{0k}=N_{0k}^{}+O(ϵ).$$
Next we determine the value of $`N_{01}^{}`$. To know the transition function $`r_{01}`$ one needs the Jacobian matrix $`R`$ of the transformation
$$(x_1,y_1,\frac{\rho _2^2}{2},\phi _2,\frac{\rho _3^2}{2},\phi _3)(\frac{\rho _1^2}{2},\phi _1,\frac{\rho _2^2}{2},\phi _2,,\frac{\rho _3^2}{2},\phi _3)$$
in the points of $`A_{01}^{}`$; where $`\rho _1^2=x_1^2+y_1^2`$, $`\phi _1=\mathrm{tan}^1(y_1/x_1)`$. The function $`r_{01}=\rho (R)`$, where $`\rho :Sp(6,)U(1)`$ is the map which restricts to the determinant on $`U(3)`$ . The non trivial block of $`R`$ is the diagonal one
$$\left(\begin{array}{cc}x_1& y_1\\ r& s\end{array}\right),$$
with $`r=y_1(x_1^2+y_1^2)^1`$ and $`s=x_1(x_1^2+y_1^2)^1`$. The non real eigenvalues of $`R`$ are
$$\lambda _\pm =\frac{x_1+s}{2}\pm \frac{i\sqrt{4(s+x_1)^2}}{2}.$$
On $`A_{01}^{}`$ these non real eigenvalues occur when $`(s+x_1)^2<2`$, that is, if $`|\mathrm{cos}\phi _1|<2ϵ(ϵ^2+1)^1=:\delta `$. If $`y_1>0`$ then $`\lambda _{}`$ of the first kind (see ) and $`\lambda _+`$ is of the first kind, if $`y_1<0`$.
Hence, on $`A_{01}^{}`$,
$$\rho (R)=\{\begin{array}{cc}\lambda _+|\lambda _+|^1=x+iy,\hfill & \text{if }|\mathrm{cos}\phi _1|<\delta \text{ and }y_1<0\text{;}\hfill \\ \lambda _{}|\lambda _{}|^1=xiy,\hfill & \text{if }|\mathrm{cos}\phi _1|<\delta \text{ and }y_1>0\text{;}\hfill \\ \pm 1,\hfill & \text{otherwise.}\hfill \end{array}$$
where $`x=\delta ^1\mathrm{cos}\phi _1`$, and $`y=\sqrt{1x^2}`$.
If we put $`\rho (R)=e^{i\gamma }`$, then $`\mathrm{cos}\gamma =\delta ^1\mathrm{cos}\phi _1`$ (when $`|\mathrm{cos}\phi _1|<\delta `$), and
$$\mathrm{sin}\gamma =\{\begin{array}{cc}\sqrt{1\mathrm{cos}^2\gamma },\hfill & \text{if }\mathrm{sin}\phi _1>0\text{;}\hfill \\ \sqrt{1\mathrm{cos}^2\gamma },\hfill & \text{if }\mathrm{sin}\phi _1<0\text{.}\hfill \end{array}$$
So when $`\phi _1`$ runs anticlockwise from $`0`$ to $`2\pi `$, $`\gamma `$ goes round clockwise the circumference; that is, $`\gamma =h(\phi _1)`$, where $`h`$ is a function such that
(2.11)
$$h(0)=2\pi ,\text{and}h(2\pi )=0.$$
As $`r_{01}=\rho (R)`$, then $`dlogr_{01}=idh`$.
On $`A_{01}^{}`$ the form $`\omega `$ reduces to $`(1/2)d\rho _2^2d\phi _2+d\rho _3^2d\phi _3`$. From (2.9) one deduces
(2.12)
$$N_{01}^{}=\frac{3i}{4\pi }_{A_{01}^{}}if\frac{h}{\phi _1}𝑑\phi _1d\rho _2^2d\phi _2d\rho _3^2d\phi _3.$$
The submanifold $`A_{01}^{}`$ is oriented as a subset of $`B_0`$ and the orientation of $`B_0`$ is the one defined by $`\omega ^3`$, that is, by
$$d\rho _1^2d\phi _1d\rho _2^2d\phi _2d\rho _3^2d\phi _3.$$
Since $`\rho _1>ϵ`$ for the points of $`B_0`$, then $`A_{01}^{}`$ is oriented by $`d\phi _1d\phi _2^2d\phi _2d\rho _3^2d\phi _3.`$ On the other hand, the Hamiltonian function $`f=\kappa +O(ϵ)`$ on $`A_{01}^{}`$. Then it follows from (2.12) together with (2.11)
$$N_{01}^{}=6\pi ^2\kappa _0^{\mu /\pi }𝑑\rho _3^2_0^{\tau /\pi \rho _3^2}𝑑\rho _2^2+O(ϵ).$$
that is,
(2.13)
$$N_{01}^{}=3\kappa (\tau ^2\lambda ^2)+O(ϵ).$$
The contributions $`N_{02}^{},N_{03}^{},N_{04}^{},N_{05}^{}`$ to $`I_\psi `$ can be calculated in a similar way. One obtains the following results up to addends of order $`ϵ`$
(2.14)
$$N_{02}^{}=N_{05}=(\tau ^3\lambda ^3)+3\kappa (\tau ^2\lambda ^2),N_{03}^{}=\tau ^2(3\kappa \tau ),N_{04}^{}=\lambda ^2(3\kappa \lambda ).$$
As $`I_\psi `$ is independent of $`ϵ`$, it follows from (2.8), (2.10), (2.13) and (2.14)
(2.15)
$$I_\psi =6\kappa (2\tau ^2\lambda ^2)+\lambda ^33\tau ^3.$$
On the other hand, straightforward calculations give
$$_M\omega ^3=(\tau ^3\lambda ^3),\text{and}_M\pi \rho _1^2\omega ^3=\frac{1}{4}(\tau ^4\lambda ^4).$$
So
(2.16)
$$\kappa =\frac{1}{4}\left(\frac{\tau ^4\lambda ^4}{\tau ^3\lambda ^3}\right).$$
It follows from (2.15) and (2.16)
(2.17)
$$I_\psi =\frac{\lambda ^2(3\tau ^4+8\tau ^3\lambda 6\tau ^2\lambda ^2+\lambda ^4)}{2(\tau ^3\lambda ^3)}.$$
Hence $`I_\psi `$ is a rational function of $`\tau `$ and $`\lambda `$. It is easy to check that its numerator does not vanish for $`0<\lambda <\tau .`$ So we have
###### Proposition 2.
If $`\psi `$ is the closed Hamiltonian isotopy defined in (2.6), then the characteristic number $`I_\psi 0`$.
Next we consider the loop $`\stackrel{~}{\psi }`$ defined by
(2.18)
$$\stackrel{~}{\psi }_t[z]=[z_1,z_2,z_3e^{2\pi it},z_4,z_5].$$
The corresponding normalized Hamiltonian function is $`\stackrel{~}{f}=\pi \rho _3^2\stackrel{~}{\kappa },`$ where
(2.19)
$$\stackrel{~}{\kappa }=\frac{1}{4}\left(\frac{\tau ^44\tau \lambda ^3+3\lambda ^4}{\tau ^3\lambda ^3}\right).$$
As in the preceding case
(2.20)
$$I_{\stackrel{~}{\psi }}=\underset{j=1}{\overset{5}{}}\stackrel{~}{N}_{0j}^{}+O(ϵ),$$
where
$$\stackrel{~}{N}_{0j}^{}=\frac{3i}{2\pi }_{A_{0j}^{}}\stackrel{~}{f}d\mathrm{log}r_{0j}\omega ^2.$$
The expression for $`\stackrel{~}{N}_{01}^{}`$ can be obtained from (2.12) substituting $`f`$ for $`\stackrel{~}{f}`$; so
(2.21)
$$\stackrel{~}{N}_{01}^{}=3(\tau \stackrel{~}{\kappa })(\tau ^2\lambda ^2)+2(\tau ^3\lambda ^3)+O(ϵ).$$
Similar calculations give the following values for the $`\stackrel{~}{N}_{0j}^{}`$’s, up to summands of order $`ϵ`$,
(2.22)
$$\stackrel{~}{N}_{02}^{}=\stackrel{~}{N}_{05}^{}=3(\tau \stackrel{~}{\kappa })(\tau ^2\lambda ^2)+2(\tau ^3\lambda ^3),\stackrel{~}{N}_{03}^{}=3\stackrel{~}{\kappa }\tau ^2,\stackrel{~}{N}_{04}^{}=3\lambda ^2(\stackrel{~}{\kappa }\mu ).$$
It follows from (2.22), (2.21) and (2.20)
(2.23)
$$I_{\stackrel{~}{\psi }}=6\stackrel{~}{\kappa }(2\tau ^2\lambda ^2)3(\tau ^32\tau \lambda ^2+\lambda ^3).$$
After (2.19) we obtain
$$I_{\stackrel{~}{\psi }}=3I_\psi .$$
In the definition of $`M`$ the variables $`z_1,z_2,z_5`$ play the same role. However we can consider the following $`S^1`$ action on $`M`$
(2.24)
$$\widehat{\psi }_t[z]=[z_1,z_2,z_3,e^{2\pi it}z_4,z_5].$$
Its Hamiltonian is $`\widehat{f}=\pi \rho _4^2\widehat{\kappa },`$ with
(2.25)
$$\widehat{\kappa }=\frac{1}{4}\left(\frac{\lambda ^44\lambda \tau ^3+3\tau ^4}{\tau ^3\lambda ^3}\right).$$
The corresponding $`\widehat{N}_{0j}^{}`$ have the following values up summand of order $`ϵ`$
(2.26)
$$\widehat{N}_{01}^{}=\widehat{N}_{02}^{}=\widehat{N}_{05}^{}=3(\lambda +\widehat{\kappa })(\tau ^2\lambda ^2)2(\tau ^3\lambda ^3),\widehat{N}_{03}^{}=3\tau ^2(\widehat{\kappa }\mu ),\widehat{N}_{04}^{}=3\widehat{\kappa }\lambda ^2.$$
From the preceding formulae one deduces
$$I_{\widehat{\psi }}=I_{\stackrel{~}{\psi }}=3I_\psi .$$
###### Theorem 3.
Let $`M`$ be the toric manifold defined by (2.1) and (2.2). If $`\psi `$, $`\stackrel{~}{\psi }`$ and $`\widehat{\psi }`$ are the Hamiltonian loops in $`M`$ defined by (2.6), (2.18) and (2.24) respectively, then
$$I_{\widehat{\psi }}=I_{\stackrel{~}{\psi }}=3I_\psi ,$$
with
$$I_\psi =\frac{\lambda ^2(3\tau ^4+8\tau ^3\lambda 6\tau ^2\lambda ^2+\lambda ^4)}{2(\tau ^3\lambda ^3)},$$
$`\lambda `$ being $`\lambda :=\tau \mu `$.
###### Corollary 4.
Let $`(M,\omega )`$ be the toric manifold one point blow up of $`P^2`$, then $`\pi _1(\text{Ham}(M,\omega ))`$ contains an infinite cyclic subgroup.
###### Proof.
By Proposition 2 $`I_\psi 0`$. As $`I`$ is a group homomorphism then the class $`[\psi ^l]\pi _1(\text{Ham}(M,\omega ))`$ does not vanish, for all $`l\{0\}`$. ∎
## 3. Hamiltonian group of toric manifolds
In this Section we generalize the calculations carried out in Section 2 for the $`6`$-manifold one point blow up of $`P^3`$ to a general toric manifold.
Let $`𝕋`$ be the torus $`(S^1)^r`$, and $`𝔱=\mathrm{}`$ its Lie algebra. Given $`w_j^r`$, with $`j=1,\mathrm{},m`$ and $`\tau ^r`$ we put
(3.1)
$$M=\{z^m:\pi \underset{j=1}{\overset{m}{}}|z_j|^2w_j=\tau \}/𝕋,$$
where the relation defined by $`𝕋`$ is
(3.2)
$$(z_j)(z_j^{})\text{iff there is}\xi 𝔱\text{such that}z_j^{}=z_je^{2\pi iw_j,\xi }\text{for}j=1,\mathrm{},m.$$
We will assume that there is an open half space in $`^r`$ which contains all the vectors $`w_j`$ and that $`\{w_j\}_j`$ span $`^r`$. We also assume that $`\tau `$ ia a regular value of the map
$$z^m\pi \underset{j=1}{\overset{m}{}}|z_j|^2w_j^r.$$
Then $`M`$ is a closed toric manifold of dimension $`n:=2(mr)`$ .
When $`0z_a`$, we write $`z_a=\rho _ae^{i\theta _a}`$. The standard symplectic form on $`^m`$ gives rise to the symplectic structure $`\omega `$ on $`M`$. On
$$\{[z]M:z_j0\text{for all}j\}$$
$`\omega `$ can be written as
$$\omega =\underset{i=1}{\overset{n}{}}d\left(\frac{\rho _{ai}^2}{2}\right)d\phi _{ai},$$
with $`\phi _{ai}`$ a linear combination of the $`\theta _c`$’s.
Given $`0<ϵ<<1`$, we set
$$B_0=\{[z]M:|z_j|>ϵ\text{for all}j\}$$
$$B_k=\{[z]M:|z_k|<ϵ,|z_j|>ϵ\text{for all}jk\},$$
as in Section 2. On $`B_0`$ we will consider the Darboux coordinates
$$\{\frac{\rho _{ai}^2}{2},\phi _{ai}\}_{i=1,\mathrm{},n}.$$
Given $`k\{1,\mathrm{},m\}`$ we write $`\omega `$ in the form
$$\omega =d\left(\frac{\rho _k^2}{2}\right)d\phi _k+\underset{i=1}{\overset{n1}{}}d\left(\frac{\rho _{ki}^2}{2}\right)d\phi _{ki},$$
where $`\phi _k`$ and $`\phi _{ki}`$ are linear combinations of the $`\theta _c`$’s. Then we consider on $`B_k`$ the following Darboux coordinates
$$\{x_k,y_k,\frac{\rho _{ki}^2}{2},\phi _{ki}\}_{i=1,\mathrm{},n1},$$
with $`x_k+iy_k:=\rho _ke^{i\phi _k}.`$
We denote by $`\psi _t`$ the map
$$\psi _t:[z]M[e^{2\pi it}z_1,z_2,\mathrm{},z_m]M.$$
$`\{\psi _t:t[0,1]\}`$ is a loop in $`\text{Ham}(M)`$. By repeating the arguments of Section 2 one obtains
$$I_\psi =\underset{k=1}{\overset{m}{}}N_{0k}^{}+O(ϵ),$$
where
$$N_{0k}^{}=\frac{ni}{2\pi }_{A_{0k}^{}}fd\mathrm{log}r_{0k}\omega ^{n1},$$
$$A_{0k}^{}=\{[z]M:|z_k|=ϵ,|z_j|>ϵ\text{for all}jk\},$$
and $`f=\pi \rho _1^2\kappa `$, with
$$_M\pi \rho _1^2\omega ^n=\kappa _M\omega ^n.$$
As in Section 2, on $`A_{0k}^{}`$ the exterior derivative $`d\mathrm{log}r_{0k}=ih^{}(\phi _k)d\phi _k,`$ where $`h=h(\phi _k)`$ is a function such that $`h(0)=2\pi ,h(2\pi )=0`$. Then
$$N_{0k}^{}=n_{\{[z]:z_k=0\}}f\omega ^{n1}+O(ϵ).$$
Since $`I_\psi `$ is independent of $`ϵ`$, we obtain
(3.3)
$$I_\psi =n\underset{k=1}{\overset{m}{}}\left(_{\{[z]:z_k=0\}}(\pi \rho _1^2\kappa )\omega ^{n1}\right).$$
This formula together with the fact that $`I`$ is a group homomorphism give the following Theorem
###### Theorem 5.
Let $`(M,\omega )`$ be the toric manifold defined by (3.1) and (3.2). If
$$\underset{k=1}{\overset{m}{}}\left(_{\{[z]:z_k=0\}}(\pi \rho _1^2\kappa )\omega ^{n1}\right)0,$$
then $`\pi _1(\text{Ham}(M,\omega ))`$ contains an infinite cyclic subgroup.
Examples. We will check the above result calculating $`I_\psi `$ by (3.3) in two particular cases: When the manifold is $`P^1`$ and when is $`P^2`$.
For
$$P^1=\{[z_1,z_2]:|z_1|^2+|z_2|^2=\tau /\pi \}/S^1$$
and $`\psi _t[z_1,z_2]=[e^{2\pi it}z_1,z_2]`$, the normalized Hamiltonian is $`f=\pi \rho _1^2\tau /2`$, that is, $`\kappa =\tau /2.`$ In this case (3.3) reduces to $`I_\psi =\tau +2\kappa =0.`$ This is compatible with the fact that $`\pi _1(\text{Ham}(P^1))=/2.`$
For
$$P^2=\{[z_1,z_2,z_3]:|z_1|^2+|z_2|^2+|z_3|^2=\tau /\pi \}/S^1,$$
the Hamiltonian is $`f=\pi \rho _1^2\tau /3.`$ Moreover for for $`k\{1,2,3\}`$
$$_{\{[z]:z_k=0\}}\omega =\tau .$$
On the other hand, for $`k=2,3`$
$$_{\{[z]:z_k=0\}}\pi \rho _1^2\omega =\tau ^2/2.$$
After (3.3) $`I_\psi =2(\tau ^23\kappa \tau )=0`$. This result is consistent with the finiteness of $`\pi _1(\text{Ham}(P^2))`$, since $`\text{Ham}(P^2)`$ has the homotopy type of $`PU(3)`$ . |
warning/0506/hep-ph0506087.html | ar5iv | text | # Factorization in the Production and Decay of the 𝑋(3872)
## I Introduction
The $`X(3872)`$ is a narrow resonance near $`3872`$ MeV discovered by the Belle collaboration in 2003 Choi:2003ue . It has been observed through the exclusive decay $`B^\pm XK^\pm `$ Choi:2003ue ; Aubert:2004ns and through its inclusive production in proton-antiproton collisions Acosta:2003zx ; Abazov:2004kp . Its mass $`m_X`$ is extremely close to the threshold for the charm mesons $`D^0`$ and $`\overline{D}^0`$ Olsen:2004fp :
$$m_X(m_{D^0}+m_{D^0})=+0.6\pm 1.1\mathrm{MeV}.$$
(1)
The upper bound on its decay width $`\mathrm{\Gamma }_X`$ is Choi:2003ue
$`\mathrm{\Gamma }_X<2.3\mathrm{MeV}(90\%\mathrm{C}.\mathrm{L}.).`$ (2)
The $`X(3872)`$ was discovered through its decay into $`J/\psi \pi ^+\pi ^{}`$ Choi:2003ue . The Belle collaboration has recently presented evidence for the decays $`XJ/\psi \pi ^+\pi ^{}\pi ^0`$ and $`XJ/\psi \gamma `$ Abe:2005ix . The decay $`XJ/\psi \gamma `$ implies that the $`X`$ has positive charge conjugation. Upper bounds have been set on several other decay modes Abe:2003zv ; Yuan:2003yz ; Aubert:2004fc ; Abe:2004sd ; Metreveli:2004px .
The nature of the $`X(3872)`$ has not yet been definitively established. The presence of the $`J/\psi `$ among its decay products motivates its interpretation as a charmonium state with constituents $`c\overline{c}`$ Barnes:2003vb ; Eichten:2004uh ; Quigg:2004nv . Two interpretations that are motivated by the proximity of $`m_X`$ to the $`D^0\overline{D}^0`$ threshold are a hadronic molecule with constituents $`DD^{}`$ Tornqvist:2003na ; Tornqvist:2004qy ; Voloshin:2003nt ; Wong:2003xk ; Braaten:2003he ; Swanson:2003tb ; Swanson:2004pp and a “cusp” associated with the $`D^0\overline{D}^0`$ threshold Bugg:2004rk ; Bugg:2004sh . Other proposed interpretations include a tetraquark with constituents $`c\overline{c}q\overline{q}`$ Vijande:2004vt , a “hybrid charmonium” state with constituents $`c\overline{c}g`$ Close:2003mb ; Li:2004st , a glueball with constitutents $`ggg`$ Seth:2004zb , and a diquark-antidiquark bound state with constituents $`cu+\overline{c}\overline{u}`$ Maiani:2004vq . Measurements of the decays of the $`X`$ can be used to determine its quantum numbers and narrow down the options Close:2003sg ; Pakvasa:2003ea ; Rosner:2004ac ; Kim:2004cz ; Abe:2005iy . The most predictive of the proposed interpretations are charmonium and $`DD^{}`$ molecules. The $`C=`$ charmonium options are ruled out by the decay $`XJ/\psi \gamma `$. Evidence ruling out or disfavoring each of the $`C=+`$ charmonium options has been accumulating Olsen:2004fp ; Quigg:2004vf ; Abe:2005iy . The most difficult charmonium state to rule out is the $`\chi _{c1}(2P)`$, partly because it can have resonant S-wave interactions with $`D^0\overline{D}^0`$ and $`D^0\overline{D}^0`$ that transform it into a $`DD^{}`$ molecule Braaten:2003he .
The possibility that charm mesons might form molecular states was considered shortly after the discovery of charm Bander:1975fb ; Voloshin:ap ; DeRujula:1976qd ; Nussinov:1976fg . The first quantitative study of the possibility of molecular states of charm mesons was carried out by Tornqvist in 1993 using a one-pion-exchange potential model. He found that the isospin-0 combinations of $`D\overline{D}^{}`$ and $`D^{}\overline{D}`$ could form weakly-bound states in the S-wave $`1^{++}`$ channel and in the P-wave $`0^+`$ channel Tornqvist:1993ng . Since the binding energy is small compared to the 8.4 MeV splitting between the $`D^0\overline{D}^0`$ threshold and the $`D^+D^{}`$ threshold, there are large isospin breaking effects Tornqvist:2003na ; Tornqvist:2004qy . After the discovery of the $`X(3872)`$, Swanson considered a potential model that includes both one-pion-exchange and quark exchange, and found that the $`C=+`$ superposition of $`D^0\overline{D}^0`$ and $`D^0\overline{D}^0`$ could form a weakly-bound state in the S-wave $`1^{++}`$ channel Swanson:2003tb . Another mechanism for generating a $`DD^{}`$ molecule is the accidental fine-tuning of the mass of the $`\chi _{c1}(2P)`$ to the $`D^0\overline{D}^0`$/$`D^0\overline{D}^0`$ threshold which creates a $`DD^{}`$ molecule with quantum numbers $`1^{++}`$ Braaten:2003he .
The assumption that the $`X(3872)`$ is a weakly-bound $`DD^{}`$ molecule is very predictive Braaten:2003he . This assumption has been used by Voloshin to predict the rates and momentum distributions for the decays of $`X`$ into $`D^0\overline{D}^0\pi ^0`$ and $`D^0\overline{D}^0\gamma `$ Voloshin:2003nt . It has been used to calculate the rate for the exclusive decay of $`\mathrm{{\rm Y}}(4S)`$ into the $`X`$ and two light hadrons Braaten:2004rw , to estimate the decay rate for the discovery mode $`B^+XK^+`$ Braaten:2004fk , and to predict the suppression of the decay rate for $`B^0XK^0`$ Braaten:2004ai . The assumption that the $`X`$ is a weakly-bound $`DD^{}`$ molecule also has implications for inclusive production Braaten:2004jg ; Voloshin:2004mh .
In this paper, we point out that if $`X`$ is a loosely-bound molecule, its short-distance decay rates and its exclusive production rates satisfy simple factorization formulas. In Section II, we illustrate the factorization formulas using a two-channel scattering model. In Section III, we show how the rate for a short-distance decay mode of the $`X`$, such as $`J/\psi \pi ^+\pi ^{}\pi ^0`$, can be factorized into long-distance factor that involves the large scattering length and a short-distance factor that is insenstitive to $`a`$. In Section IV, we show how a production rate for the $`X`$, such as the decay rate for $`B^+XK^+`$, can be factorized into a long-distance factor and a short-distance factor. In Section V, we use factorization to calculate the shapes of the invariant mass distributions of $`D^0\overline{D}^0`$ and $`D^0\overline{D}^0`$ near threshold. In Section VI, we use factorization to calculate the line shape of the $`X`$ in any of its short-distance decay modes. The line shape depends on the real and imaginary parts of the scattering length and can differ substantially from a conventional Breit-Wigner resonance. A summary of our results is given in Section VII.
## II Factorization in a Simple Scattering Model
### II.1 Universality for large scattering length
Nonrelativistic few-body systems with short-range interactions and a large scattering length $`a`$ have universal properties that depend on the scattering length but are otherwise insensitive to details at distances small compared to $`a`$ Braaten:2004rn . In any specific system, there is a natural momentum scale $`\mathrm{\Lambda }`$ that sets the scale of most low-energy scattering parameters. The scattering length is large if it satisfies $`|a|\mathrm{\Lambda }^1`$. Universality predicts that the T-matrix element for 2-body elastic scattering with relative momentum $`p\mathrm{\Lambda }`$ is
$$𝒯(p)=\frac{2\pi /\mu }{1/aip},$$
(3)
where $`\mu `$ is the reduced mass of the two particles. If $`a`$ is real and positive, universality predicts that there is a weakly-bound state with binding energy
$$E_X=\frac{1}{2\mu a^2}.$$
(4)
The universal momentum-space wavefunction of this bound state is
$`\psi (p)={\displaystyle \frac{(8\pi /a)^{1/2}}{p^2+1/a^2}}.`$ (5)
The universal amplitude for transitions from the bound state to a scattering state consisting of two particles with relative momentum $`p\mathrm{\Lambda }`$ is
$$𝒜_X=\frac{\sqrt{2\pi }}{\mu }a^{1/2}.$$
(6)
These results are all encoded in the universal expression for the truncated connected transition amplitude:
$$𝒜(E)=\frac{2\pi /\mu }{1/a+\sqrt{2\mu E}}.$$
(7)
The universal amplitude in Eq. (7) can be obtained from a local effective field theory for the two particles. The particles interact through an S-wave contact interaction with Feynman rule $`i(2\pi /\mu )a_0(\mathrm{\Lambda })`$, where the parameter $`a_0(\mathrm{\Lambda })`$ is a bare scattering length that depends on the ultraviolet momentum cutoff $`\mathrm{\Lambda }`$. The amplitude $`𝒜(E)`$ can be obtained by summing the geometric series of Feynman diagrams in Fig. 1:
$$𝒜(E)=\frac{(2\pi /\mu )a_0(\mathrm{\Lambda })}{1+(2\pi /\mu )a_0(\mathrm{\Lambda })L(\mathrm{\Lambda },E)},$$
(8)
where $`L(\mathrm{\Lambda },E)`$ is the amplitude for the propagation of the particles between successive contact interactions:
$$L(\mathrm{\Lambda },E)=\frac{\mu }{\pi ^2}\left(\mathrm{\Lambda }\frac{\pi }{2}\sqrt{2\mu E}\right).$$
(9)
Renormalization is accomplished by eliminating $`a_0(\mathrm{\Lambda })`$ in favor of the physical scattering length:
$`a={\displaystyle \frac{a_0(\mathrm{\Lambda })}{1+(2/\pi )\mathrm{\Lambda }a_0(\mathrm{\Lambda })}}.`$ (10)
With this substitution, the expression for $`𝒜(E)`$ in Eq. (8) reduces without approximation to the universal result in Eq. (7). Note that the scattering length $`a`$ can be tuned to $`\pm \mathrm{}`$ by tuning the bare scattering length to a critical value of order $`\mathrm{\Lambda }^1`$:
$`a_0(\mathrm{\Lambda }){\displaystyle \frac{\pi }{2}}\mathrm{\Lambda }^1.`$ (11)
If the 2-body system has inelastic scattering channels, the large scattering length $`a`$ will have a negative imaginary part. It is convenient to express the complex scattering length in the form
$`{\displaystyle \frac{1}{a}}=\gamma _{\mathrm{re}}+i\gamma _{\mathrm{im}},`$ (12)
where $`\gamma _{\mathrm{re}}`$ and $`\gamma _{\mathrm{im}}`$ are real and $`\gamma _{\mathrm{im}}0`$. In the case $`\gamma _{\mathrm{re}}>0`$ where there is a weakly-bound state, it can decay into the inelastic channel. The expression for the binding energy on the right side of Eq. (4) is complex-valued. Its real part $`E_{X,\mathrm{pole}}`$ and its imaginary part $`\mathrm{\Gamma }_X/2`$ are given by
$`E_{X,\mathrm{pole}}`$ $`=`$ $`(\gamma _{\mathrm{re}}^2\gamma _{\mathrm{im}}^2)/(2\mu ),`$ (13a)
$`\mathrm{\Gamma }_X`$ $`=`$ $`2\gamma _{\mathrm{re}}\gamma _{\mathrm{im}}/\mu .`$ (13b)
These quantities specify that there is a pole in the S-matrix at the energy $`E=E_{X,\mathrm{pole}}i\mathrm{\Gamma }_X/2`$. As we shall see in Section VI, $`\mathrm{\Gamma }_X`$ can be interpreted as the full width at half-maximum of a resonance in the inelastic channel provided $`\gamma _{\mathrm{im}}<\gamma _{\mathrm{re}}`$. The peak of the resonance is below the threshold for the two particles by
$`E_X`$ $`=`$ $`\gamma _{\mathrm{re}}^2/(2\mu ).`$ (14)
We therefore interpret $`E_X`$ as the binding energy rather than $`E_{X,\mathrm{pole}}`$.
### II.2 Two-channel model
Cohen, Gelman, and van Kolck have constructed a renormalizable effective field theory that describes two scattering channels with S-wave contact interactions Cohen:2004kf . We will refer to this model as the two-channel scattering model. An essentially equivalent model has been used to describe the effects of $`\mathrm{\Delta }\mathrm{\Delta }`$ states on the two-nucleon system Savage:1996tb . The parameters of this model can be tuned to produce a large scattering length in the lower energy channel. It can therefore be used as a simple model for the effects on the $`D^0\overline{D}^0`$/$`D^0\overline{D}^0`$ system of other hadronic channels with nearby thresholds, such as $`J/\psi \rho `$, $`J/\psi \omega `$, and $`D^\pm D^{}`$.
The two-channel model of Ref. Cohen:2004kf describes two scattering channels with S-wave contact interactions only. We label the particles in the first channel $`1a`$ and $`1b`$ and those in the second channel $`2a`$ and $`2b`$. We denote the reduced masses in the two channels by $`\mu _1`$ and $`\mu _2`$. Renormalized observables in the 2-body sector are expressed in terms of 4 parameters: three interaction parameters $`a_{11}`$, $`a_{12}`$, and $`a_{22}`$ with dimensions of length and the energy gap $`\mathrm{\Delta }`$ between the two scattering channels, which is determined by the masses of the particles:
$`\mathrm{\Delta }`$ $`=`$ $`m_{2a}+m_{2b}(m_{1a}+m_{1b}).`$ (15)
The scattering parameters in Ref. Cohen:2004kf were defined in such a way that $`a_{11}`$ and $`a_{22}`$ reduce in the limit $`a_{12}\pm \mathrm{}`$ to the scattering lengths for the two channels. The truncated connected transition amplitude $`𝒜(E)`$ for this coupled-channel system is a $`2\times 2`$ matrix that depends on the energy $`E`$ in the center-of-mass frame. If that energy is measured relative to the threshold $`m_{1a}+m_{1b}`$ for the first scattering channel, the inverse of the matrix $`𝒜(E)`$ is<sup>1</sup><sup>1</sup>1The expression for the matrix $`T_s^1`$ in Eq. (2.18) of Ref. Cohen:2004kf should be equal to $`𝒜(E)^1`$ evaluated at $`E=p^2/(2\mu _1)`$. There is an error in the 22 component of $`T_s^1`$: the square root $`\sqrt{p^22\mu _2\mathrm{\Delta }}`$ should be $`\sqrt{(\mu _2/\mu _1)p^22\mu _2\mathrm{\Delta }}`$.
$`𝒜(E)^1={\displaystyle \frac{1}{2\pi }}\left(\begin{array}{cc}\mu _1\left[1/a_{11}+\sqrt{2\mu _1E}\right]& \sqrt{\mu _1\mu _2}/a_{12}\\ \sqrt{\mu _1\mu _2}/a_{12}& \mu _2\left[1/a_{22}+\sqrt{2\mu _2(\mathrm{\Delta }E)}\right]\end{array}\right).`$ (16)
The square roots are defined for negative real arguments by the prescription $`EE+iϵ`$ with $`ϵ0^+`$. The amplitudes defined by Eq. (16) are for transitions between states with the standard nonrelativistic normalizations. The transitions between states with the standard relativistic normalizations are obtained by multiplying by a factor $`\sqrt{2m_i}`$ for every particle in the initial and final state. We will need explicit expression for the $`11`$ and $`12`$ entries of this matrix:
$`𝒜_{11}(E)`$ $`=`$ $`{\displaystyle \frac{2\pi }{\mu _1}}\left({\displaystyle \frac{1}{a_{11}}}+\sqrt{2\mu _1E}{\displaystyle \frac{1}{a_{12}^2}}\left[1/a_{22}+\sqrt{2\mu _2(\mathrm{\Delta }E)}\right]^1\right)^1,`$ (17a)
$`𝒜_{12}(E)`$ $`=`$ $`{\displaystyle \frac{2\pi }{\sqrt{\mu _1\mu _2}}}\left({\displaystyle \frac{1}{a_{12}}}a_{12}\left[{\displaystyle \frac{1}{a_{11}}}+\sqrt{2\mu _1E}\right]\left[{\displaystyle \frac{1}{a_{22}}}+\sqrt{2\mu _2(\mathrm{\Delta }E)}\right]\right)^1.`$ (17b)
The T-matrix element for the elastic scattering of particles in the first channel with relative momentum $`p`$ is obtained by evaluating $`𝒜_{11}(E)`$ in Eq. (17a) at the energy $`E=p^2/(2\mu _1)`$:
$`𝒯_{11}(p)={\displaystyle \frac{2\pi }{\mu _1}}\left({\displaystyle \frac{1}{a_{11}}}ip{\displaystyle \frac{1}{a_{12}^2}}\left[1/a_{22}+\sqrt{2\mu _2\mathrm{\Delta }(\mu _2/\mu _1)p^2}\right]^1\right)^1.`$ (18)
Setting $`𝒯_{11}(0)=2\pi /(\mu _1a)`$, we can read off the inverse scattering length $`1/a`$:
$`{\displaystyle \frac{1}{a}}`$ $`=`$ $`{\displaystyle \frac{1}{a_{11}}}+{\displaystyle \frac{1}{a_{12}^2}}\left[\sqrt{2\mu _2\mathrm{\Delta }}1/a_{22}\right]^1.`$ (19)
If there is a bound state with energy $`\kappa ^2/(2\mu _1)`$ below the scattering threshold for the first channel, the matrix $`𝒜(E)`$ given by Eq. (16) has a pole at $`E=\kappa ^2/(2\mu _1)`$. The binding momentum $`\kappa `$ satisfies
$`\kappa ={\displaystyle \frac{1}{a_{11}}}+{\displaystyle \frac{1}{a_{12}^2}}\left[1/a_{22}+\sqrt{2\mu _2\mathrm{\Delta }+(\mu _2/\mu _1)\kappa ^2}\right]^1.`$ (20)
The behavior of the matrix $`𝒜(E)`$ as the energy $`E`$ approaches the pole associated with the bound state is
$`𝒜(E){\displaystyle \frac{1}{E+\kappa ^2/(2\mu _1)}}\left(\begin{array}{c}𝒜_{X1}\\ 𝒜_{X2}\end{array}\right)\left(\begin{array}{cc}𝒜_{X1}& 𝒜_{X2}\end{array}\right).`$ (21)
The components $`𝒜_{X1}`$ and $`𝒜_{X2}`$ of the column vector are the amplitudes for transitions from the bound state to particles in the first and second channels, respectively. The column vector is an eigenvector of the matrix $`𝒜(E)^1`$ in Eq. (16) with eigenvalue zero, so its components must satisfy
$`\mu _1[1/a_{11}+\kappa ]𝒜_{X1}+[\sqrt{\mu _1\mu _2}/a_{12}]𝒜_{X2}=0.`$ (22)
### II.3 Two-channel model with large scattering length
The two-channel model of Ref. Cohen:2004kf can be used as a phenomenological model for a system with a large scattering length $`a`$ in the first channel. The large scattering length requires a fine-tuning of the parameters $`a_{11}`$, $`a_{22}`$, $`a_{12}`$, and $`\mathrm{\Delta }`$. There are various ways to tune the parameters so that $`a\pm \mathrm{}`$. For example, if $`a_{11}<a_{12}^2/a_{22}`$, the energy gap $`\mathrm{\Delta }`$ can be tuned to the critical value where the right side of Eq. (19) vanishes. Alternatively, the scattering parameter $`a_{11}`$ can be tuned to the critical value $`a_{12}^2[\sqrt{2\mu _2\mathrm{\Delta }}1/a_{22}]`$. The coefficients in the low-energy expansion of $`𝒯_{11}(p)^1`$ are proportional to various powers of $`1/a_{11}`$, $`1/a_{12}`$, and $`\sqrt{2\mu _2\mathrm{\Delta }}`$. We assume that these momentum scales are comparable in magnitude. We refer to that common momentum scale as the natural low-energy scale and we denote it by $`\mathrm{\Lambda }`$.
For $`|a|\mathrm{\Lambda }^1`$ and $`|E|\mathrm{\Lambda }^2/(2\mu _1)`$, the amplitude $`𝒜_{11}(E)`$ in Eq. (17a) approaches the universal amplitude $`𝒜(E)`$ in Eq. (7) with $`\mu =\mu _1`$. It follows that for $`p\mathrm{\Lambda }`$ the T-matrix element $`𝒯_{11}(p)`$ in Eq. (18) approaches the universal T-matrix element $`𝒯(p)`$ in Eq. (3). For $`a\mathrm{\Lambda }^1`$, the solution to Eq. (20) for the binding momentum $`\kappa `$ approaches $`1/a`$, so $`\kappa ^2/(2\mu )`$ approaches the universal binding energy $`E_X`$ in Eq. (4). Finally the amplitude $`𝒜_{X1}`$ for transitions from the bound state to particles in the first channel, which is defined in Eq. (21), approaches the universal amplitude $`𝒜_X`$ in Eq. (6).
There are also universal features associated with transitions to the second channel. If $`|a|\mathrm{\Lambda }^1`$ and $`|E|\mathrm{\Lambda }^2/\mu _1`$, the leading term in the transition amplitude $`𝒜_{12}(E)`$ in Eq. (17b) reduces to
$`𝒜_{12}(E)={\displaystyle \frac{\sqrt{\mu _1/\mu _2}}{a_{12}}}\left[\sqrt{2\mu _2\mathrm{\Delta }}1/a_{22}\right]^1𝒜(E),`$ (23)
where $`𝒜(E)`$ is the universal amplitude in Eq. (7) with $`\mu `$ replaced by $`\mu _1`$. For $`a\mathrm{\Lambda }^1`$, the leading term in the amplitude $`𝒜_{X2}`$ for transitions of the weakly-bound state $`X`$ to particles in the second channel, which is defined in Eq. (21), reduces to
$`𝒜_{X2}={\displaystyle \frac{\sqrt{\mu _1/\mu _2}}{a_{12}}}\left[\sqrt{2\mu _2\mathrm{\Delta }}1/a_{22}\right]^1𝒜_X,`$ (24)
where $`𝒜_X`$ is the universal amplitude in Eq. (6). Note that the ratio $`𝒜_{12}(E)/𝒜_{X2}`$ of the amplitudes in Eqs. (23) and (24) is a universal function of $`a`$ and $`E`$ only.
The expressions for $`𝒜_{12}(E)`$ and $`𝒜_{X2}`$ in Eqs. (23) and (24) are examples of factorization formulas. They express the leading terms in the amplitudes as products of the same short-distance factor and different long-distance factors $`𝒜(E)`$ and $`𝒜_X`$. The long-distance factors involve the large scattering length $`a`$. The limit $`|a|\mathrm{}`$ has been taken in the short-distance factors. The conditions $`|E|\mathrm{\Lambda }^2/(2\mu )`$ or $`E=E_X`$ require the particles in the second channel to be off the energy shell by approximately $`\mathrm{\Delta }`$. In the short time $`1/\mathrm{\Delta }`$ allowed by the uncertainty principle, those particles can propagate only over short distances of order $`(2\mu _2\mathrm{\Delta })^{1/2}`$. This is small compared to the distance scales $`(2\mu |E|)^{1/2}`$ or $`|a|`$ associated with the particles in the first channel. Thus as far as they are concerned, the particles in the second channel act only as a point source for particles in the first channel. The amplitudes for particles from such a point source to evolve into particles of energy $`E`$ and into the weakly-bound state are $`L(\mathrm{\Lambda },E)𝒜(E)`$ and $`L(\mathrm{\Lambda },E_X)𝒜_X`$, respectively. By using the conditions $`|E|,E_X\mathrm{\Lambda }^2/(2\mu )`$, these amplitudes reduce to $`(\mu \mathrm{\Lambda }/\pi ^2)𝒜(E)`$ and $`(\mu \mathrm{\Lambda }/\pi ^2)𝒜_X`$, respectively. In these expressions, the short-distance factors are identical and the long-distance factors are the same as those in Eqs. (23) and (24).
### II.4 Unstable particle in the second channel
Now let us suppose one of the scattering particles in the second channel has a nonzero width. We take that particle to be $`2b`$. We assume that its width $`\mathrm{\Gamma }_{2b}`$ arises from its decay into particles with relativistic momenta that are much greater than the ultraviolet cutoff $`\mathrm{\Lambda }_{\mathrm{UV}}`$ that defines the domain of validity of the two-channel model. The momenta of the decay products are therefore also much greater than $`\sqrt{2\mu _2\mathrm{\Delta }}`$. We assume that $`\mathrm{\Gamma }_{2b}`$ is small compared to the mass $`m_{2b}`$, but not necessarily small compared to the energy gap $`\mathrm{\Delta }`$ between the two channels. This makes it necessary to take into account the contribution to the self-energy of particle $`2b`$ from the coupling to its decay products.
Taking into account the self-energy of particle $`2b`$ would modify the term $`\sqrt{2\mu _2(\mathrm{\Delta }E)}`$ in the inverse of the matrix of transition amplitudes given in Eq. (16). That term arises from the amplitude for the propagation of particles in the second channel between contact interactions, which is given by the integral
$`{\displaystyle \frac{d^3p}{(2\pi )^3}\frac{1}{E\mathrm{\Delta }p^2/(2\mu _2)+iϵ}}={\displaystyle \frac{\mu _2}{\pi ^2}}\left(\mathrm{\Lambda }_{\mathrm{UV}}{\displaystyle \frac{\pi }{2}}\sqrt{2\mu _2(\mathrm{\Delta }Eiϵ)}\right).`$ (25)
The cutoff constrains the momentum to the region in which the nonrelativistic approximation for the energy of the particle $`2b`$ is valid. In this region, the self-energy $`\mathrm{\Pi }`$ can be expressed as a function of $`E^{}=E\mathrm{\Delta }p^2/(2\mu _2)`$. It can be taken into account by replacing $`\mathrm{\Delta }`$ in the integral in Eq. (25) by $`\mathrm{\Delta }+\mathrm{\Pi }(E^{})`$. The assumption that the decay products of particle $`2b`$ have relativistic momenta comparable to $`m_{2b}`$ implies that their contributions to $`\mathrm{\Pi }(E^{})`$ have significant dependence on $`E^{}`$ only for variations in $`E^{}`$ that are comparable to $`m_{2b}`$. For energies satisfying $`|E|\mathrm{\Lambda }_{\mathrm{UV}}^2/(2\mu _2)`$ and loop momenta $`p<\mathrm{\Lambda }_{\mathrm{UV}}`$, the dependence on $`E^{}`$ can be neglected and the argument of $`\mathrm{\Pi }(E^{})`$ can be set to a constant, such as $`\mathrm{\Delta }`$. The prescription for taking into account the self-energy then reduces to replacing $`\mathrm{\Delta }`$ in the integral in Eq. (25) by $`\mathrm{\Delta }+\mathrm{\Pi }(\mathrm{\Delta })`$. The real part of $`\mathrm{\Pi }(\mathrm{\Delta })`$ can be absorbed into $`\mathrm{\Delta }`$ so that it becomes the physical threshold. The imaginary part of $`\mathrm{\Pi }(\mathrm{\Delta })`$ is related to the width of particle $`2b`$: $`\mathrm{Im}\mathrm{\Pi }(\mathrm{\Delta })=\mathrm{\Gamma }_{2b}/2`$. Thus the leading effect of the self-energy can be taken into account by replacing $`\mathrm{\Delta }`$ in Eq. (16) by the complex-valued energy gap
$`\mathrm{\Delta }`$ $`=`$ $`m_{2a}+m_{2b}(m_{1a}+m_{1b})i\mathrm{\Gamma }_{2b}/2.`$ (26)
If $`\mathrm{\Delta }`$ is complex, the solution to Eq. (20) for the binding momentum $`\kappa `$ is complex. It determines the pole mass $`m_{X,\mathrm{pole}}`$ and the width $`\mathrm{\Gamma }_X`$ of the weakly-bound state according to
$`m_{1a}+m_{1b}\kappa ^2/(2\mu _1)=m_{X,\mathrm{pole}}i\mathrm{\Gamma }_X/2.`$ (27)
The imaginary part reflects the fact that the bound state can decay into particle 2a and decay products of particle 2b. The quantities $`m_{X,\mathrm{pole}}`$ and $`\mathrm{\Gamma }_X`$ in Eq. (27) give the location of a pole in the S-matrix. They need not have the standard interpretations as the location of the peak and the full width at half maximum of a Breit-Wigner resonance.
### II.5 The $`𝑫^\mathrm{𝟎}\overline{𝑫}^\mathbf{}\mathrm{𝟎}`$/$`𝑫^\mathbf{}\mathrm{𝟎}\overline{𝑫}^\mathrm{𝟎}`$ System
The energy difference between the mass of $`X`$ and the $`D^0\overline{D}^0`$ threshold, which is given in Eq. (1), is small compared to the natural energy scale for binding by the pion exchange interaction: $`m_\pi ^2/2\mu 10`$ MeV, where $`\mu `$ is the reduced mass of $`D^0`$ and $`\overline{D}^0`$. The unnaturally small value of the energy difference implies that if the $`X`$ couples to $`D^0\overline{D}^0`$ and $`D^0\overline{D}^0`$, the S-wave scattering lengths for those channels must be large compared to the natural length scale $`1/m_\pi `$ associated with the pion exchange interaction. We assume that there is a large scattering length $`a`$ in the channel $`(DD^{})_+^0`$ with even charge conjugation defined by
$`|(DD^{})_+^0={\displaystyle \frac{1}{\sqrt{2}}}\left(|D^0\overline{D}^0+|D^0\overline{D}^0\right).`$ (28)
If the scattering length in the $`C=`$ channel is negligible in comparison, the scattering lengths for elastic $`D^0\overline{D}^0`$ scattering and elastic $`D^0\overline{D}^0`$ scattering are both $`a/2`$. We identify the $`X`$ as a bound state in the $`(DD^{})_+^0`$ channel.
The decays of the $`X`$ imply that the scattering length $`a`$ is complex-valued. It can be parameterized in terms of the real and imaginary parts of $`1/a`$ as in Eq. (12). Our interpretation of $`X`$ as a bound state requires $`\gamma _{\mathrm{re}}>0`$. The energy difference in Eq. (1) puts an upper bound on $`\gamma _{\mathrm{re}}`$:
$`\gamma _{\mathrm{re}}<40\mathrm{MeV}(90\%\mathrm{C}.\mathrm{L}.).`$ (29)
The upper bound on the width in Eq. (2) puts an upper bound on the product of $`\gamma _{\mathrm{re}}`$ and $`\gamma _{\mathrm{im}}`$:
$`\gamma _{\mathrm{re}}\gamma _{\mathrm{im}}<(33\mathrm{MeV})^2(90\%\mathrm{C}.\mathrm{L}.).`$ (30)
There is also a lower bound on the width of the $`X`$ from its decays into $`D^0\overline{D}^0\pi ^0`$ and $`D^0\overline{D}^0\gamma `$, which both proceed through the decay of a constituent $`D^{}`$. These decays involve interesting interference effects, but the decay rates have smooth limits as the binding energy is tuned to 0 Voloshin:2003nt . In this limit, the constructive interference increases the decay rate by a factor of 2: the partial width of $`X`$ reduces to $`2\mathrm{\Gamma }[D^0]`$. The width of $`D^0`$ has not been measured, but it can be deduced from other information about the decays of $`D^0`$ and $`D^+`$. Using the total width of the $`D^+`$, its branching fraction into $`D^+\pi ^0`$, and isospin symmetry, we can deduce the partial width of $`D^0`$ into $`D^0\pi ^0`$: $`\mathrm{\Gamma }[D^0D^0\pi ^0]=42\pm 10`$ keV. The total width of the $`D^0`$ can then be obtained by dividing by its branching fraction into $`D^0\pi ^0`$: $`\mathrm{\Gamma }[D^0]=68\pm 16`$ keV. The sum of the partial widths of $`X`$ into $`D^0\overline{D}^0\pi ^0`$ and $`D^0\overline{D}^0\gamma `$ is therefore $`136\pm 32`$ keV. The resulting lower bound on the product of $`\gamma _{\mathrm{re}}`$ and $`\gamma _{\mathrm{im}}`$ is
$`\gamma _{\mathrm{re}}\gamma _{\mathrm{im}}>(7\mathrm{MeV})^2(90\%\mathrm{C}.\mathrm{L}.).`$ (31)
By combining this with the upper bound on $`\gamma _{\mathrm{re}}`$ in Eq. (29), we can infer that $`\gamma _{\mathrm{im}}>1`$ MeV.
## III Short-distance decays of $`X`$
The decay modes of the $`X(3872)`$ can be classified into long-distance decays and short-distance decays. The long-distance decay modes are $`D^0\overline{D}^0\pi ^0`$ and $`D^0\overline{D}^0\gamma `$, which proceed through the decay of a constituent $`D^0`$ or $`\overline{D}^0`$. These decays are dominated by a component of the wavefunction of the $`X`$ in which the separation of the $`D`$ and $`D^{}`$ is of order $`1/|a|`$. These long-distance decays involve interesting interference effects between the $`D^0\overline{D}^0`$ and $`D^0\overline{D}^0`$ components of the wavefunction Voloshin:2003nt . The short-distance decays involve a component of the wavefunction in which the separation of the $`D`$ and $`D^{}`$ is of order $`1/m_\pi `$ or smaller. Examples are the observed decay modes $`J/\psi \pi ^+\pi ^{}`$, $`J/\psi \pi ^+\pi ^{}\pi ^0`$, and $`J/\psi \gamma `$.
Short-distance decays of the $`X`$ into a hadronic final state $`H`$ involve well-separated momentum scales. The $`DD^{}`$ wavefunction of the $`X`$ involves the momentum scale $`1/|a|`$ set by the large scattering length. The transition of the $`DD^{}`$ to $`H`$ involves momentum scales $`m_\pi `$ or larger. We will refer to momentum scales of order $`1/|a|`$ and smaller as long-distance scales and momentum scales of order $`m_\pi `$ and larger as short-distance scales. We denote the arbitrary boundary between these two momentum regions by $`\mathrm{\Lambda }`$.
The separation of scales $`|a|1/m_\pi `$ in the decay process $`XH`$ can be exploited through a factorization formula for the T-matrix element:
$`𝒯[XH]=\sqrt{2m_X}𝒜_X\times 𝒜_{\mathrm{short}}[(DD^{})_+^0H].`$ (32)
In the long-distance factor, $`𝒜_X`$ is the universal amplitude given in Eq. (6) and the factor of $`\sqrt{2m_X}`$ takes into account the difference between the standard nonrelativistic and relativistic normalizations of states. If the complex scattering length is parameterized as in Eq. (12), this factor is
$`𝒜_X={\displaystyle \frac{\sqrt{2\pi }}{\mu }}\left(\gamma _{\mathrm{re}}+i\gamma _{\mathrm{im}}\right)^{1/2}.`$ (33)
The short-distance factor $`𝒜_{\mathrm{short}}`$ in Eq. (32) is insensitive to $`a`$, and one can therefore take the limit $`|a|\mathrm{}`$ in this factor. The factorization formula in Eq. (32) can serve as a definition of the short-distance factor. The content of the factorization statement then resides in the fact that, up to corrections suppressed by powers of $`1/(am_\pi )`$, the same short-distance factor appears in the factorization formula for the T-matrix element for the scattering process $`D^0\overline{D}^0H`$ at energies $`E`$ near the $`D^0\overline{D}^0`$ threshold:
$`𝒯[D^0\overline{D}^0H]={\displaystyle \frac{1}{\sqrt{2}}}\sqrt{4m_{D^0}m_{D^0}}𝒜(E)\times 𝒜_{\mathrm{short}}[(DD^{})_+^0H].`$ (34)
In the long-distance factor, the $`1/\sqrt{2}`$ is the amplitude for $`D^0\overline{D}^0`$ to be in the channel $`(DD^{})_+^0`$ with the large scattering length, the factor of $`\sqrt{4m_{D^0}m_{D^0}}`$ takes into account the difference between the standard nonrelativistic and relativistic normalizations of states, and $`𝒜(E)`$ is the universal amplitude given in Eq. (7). If the complex scattering length is parameterized as in Eq. (12), this factor is
$`𝒜(E)={\displaystyle \frac{2\pi /\mu }{\gamma _{\mathrm{re}}i(\gamma _{\mathrm{im}}+\sqrt{2\mu E})}}.`$ (35)
The factorization formulas in Eqs. (34) and (32) are analogous to those in Eqs. (23) and (24) for the two-channel model with a large scattering length in the first channel.
The factorization formulas in Eqs. (32) and (34) can be motivated diagrammatically by separating virtual particles into soft particles and hard particles according to whether they are off their energy shells by less than or by more than $`\mathrm{\Lambda }^2/(2\mu )`$, where $`\mathrm{\Lambda }`$ is the arbitrary momentum separating the long-distance scale $`1/|a|`$ and the short-distance scale $`m_\pi `$. Any contribution from soft particles inside a subdiagram all of whose external legs are hard can be Taylor-expanded in the momentum of the soft particles, leading to suppression factors of $`1/(a\mathrm{\Lambda })`$. The diagrams with the fewest suppression factors will be ones that can be separated into a part for which all the internal lines are hard particles and a part that involves only soft particles. This separation leads to the factorization formula.
The leading terms in the T-matrix element for $`XH`$ in the limit $`|a|m_\pi 1`$ can be represented by the Feynman diagrams in Fig. 2 and can be expressed as
$`𝒯[XH]=\sqrt{2m_X}{\displaystyle \frac{d^3p}{(2\pi )^3}\psi (p)𝒜^{(\mathrm{\Lambda })}[(DD^{})_+^0H]},`$ (36)
where $`\psi (p)`$ is the universal wavefunction in Eq. (5). The factor $`𝒜^{(\mathrm{\Lambda })}`$, which is represented by a dot in Fig. 2, is an amplitude for the transition $`(DD^{})_+^0H`$ in which all virtual particles are off their energy shells by more than $`\mathrm{\Lambda }^2/(2\mu )`$. It is therefore insensitive to the relative momentum $`𝒑`$ of the $`D`$ and $`D^{}`$. If that momentum dependence is neglected and if the integral in Eq. (36) is regularized by a momentum cutoff $`|𝒑|<\mathrm{\Lambda }`$, the wavefunction factor reduces in the limit $`|a|\mathrm{\Lambda }^1`$ to
$`{\displaystyle \frac{d^3p}{(2\pi )^3}\psi (p)}=\sqrt{{\displaystyle \frac{2}{\pi ^3}}}\mathrm{\Lambda }a^{1/2}.`$ (37)
The factorization formula in Eq. (32) is then obtained by absorbing a factor of $`(\mu /\pi ^2)\mathrm{\Lambda }`$ into $`𝒜^{(\mathrm{\Lambda })}`$ to obtain the short-distance factor:
$`𝒜_{\mathrm{short}}[(DD^{})_+^0H]=\left({\displaystyle \frac{\mu }{\pi ^2}}\mathrm{\Lambda }\right)𝒜^{(\mathrm{\Lambda })}[(DD^{})_+^0H].`$ (38)
Since the T-matrix element in Eq. (32) is independent of the arbitrary separation scale, the dependence on $`\mathrm{\Lambda }`$ must cancel between the two factors on the right side of Eq. (38).
The leading term in the T-matrix element for $`D^0\overline{D}^0H`$ in the limits $`|a|\mathrm{\Lambda }^1`$ and $`E\mathrm{\Lambda }^2/(2\mu )`$ can be represented by the Feynman diagram in Fig. 3 and can be expressed as
$`𝒯[D^0\overline{D}^0H]={\displaystyle \frac{1}{\sqrt{2}}}\sqrt{4m_{D^0}m_{D^0}}𝒜(E)L(\mathrm{\Lambda },E)𝒜^{(\mathrm{\Lambda })}[(DD^{})_+^0H].`$ (39)
The factor $`L(\mathrm{\Lambda },E)`$ is the amplitude for the propagation of the $`D`$ and $`D^{}`$ between successive contact interactions, which is given in Eq. (9). The approximation $`E\mathrm{\Lambda }^2/(2\mu )`$ justifies neglecting the $`\sqrt{2\mu E}`$ term in $`L(\mathrm{\Lambda },E)`$. The factorization formula in Eq. (34) is then obtained by absorbing the remaining term $`(\mu /\pi ^2)\mathrm{\Lambda }`$ into $`𝒜^{(\mathrm{\Lambda })}`$ to obtain the short-distance factor in Eq. (38).
The factorization formula for the T-matrix element in Eq. (32) implies a factorization formula for the decay rate for $`XH`$. The decay rate $`\mathrm{\Gamma }[XH]`$ can be expressed as the product of a short-distance factor and the long-distance factor
$`|𝒜_X|^2={\displaystyle \frac{2\pi }{\mu ^2}}\sqrt{\gamma _{\mathrm{re}}^2+\gamma _{\mathrm{im}}^2}.`$ (40)
Using the expressions for the binding energy and the total width of the $`X`$ in Eqs. (14) and (13b), the long-distance factor in Eq. (40) can be expressed as
$`|𝒜_X|^2=\sqrt{{\displaystyle \frac{8\pi ^2}{\mu ^3}}}[E_X+\mathrm{\Gamma }_X^2/(16E_X)]^{1/2}.`$ (41)
If the partial width for a short-distance decay mode of the $`X`$ has been calculated using a model with a specific binding energy, the factorization formula for the decay rate can be used to extrapolate the prediction to other values of the binding energy and to take into account the effect of the width of the $`X`$. This is useful because numerical calculations in models often become increasingly unstable as the binding energy is tuned to zero. Swanson has estimated the partial widths for various short-distance decays of $`X`$ using a potential model, but only for binding energies down to about 1 MeV and without taking into account the effect of the width of the $`X`$ Swanson:2003tb ; Swanson:2004pp . His predictions can be extrapolated to other values of the binding energy and the width of the $`X`$ can be taken into account by using the long-distance factor in Eq. (41).
## IV Production of $`𝑿`$
The production of $`X`$ necessarily involves the long-distance momentum scale $`1/|a|`$ through the $`(DD^{})_+^0`$ wavefunction of the $`X`$. The production also involves much larger momentum scales. Unless there are already hadrons in the initial state containing a $`c`$ and $`\overline{c}`$, the production process involves the scale $`m_c`$ associated with the creation of a $`c\overline{c}`$ pair. Even if the initial state includes hadrons that contain $`c`$ and $`\overline{c}`$, such as $`J/\psi `$ or $`D^+`$ and $`D^{}`$, the production process involves the scale $`m_\pi `$ associated with the formation of the $`D^0`$ and $`\overline{D}^0`$ that bind to form the $`X`$. We will define a short-distance production process to be one for which the initial state either does not include any of the charm mesons $`D^0`$, $`\overline{D}^0`$, $`D^0`$, or $`\overline{D}^0`$, or if it does, the momentum of the charm meson in the rest frame of the $`X`$ is of order $`m_\pi `$ or larger. All practical production mechanisms for $`X`$ in high energy physics are short-distance processes. Long-distance production mechanisms could arise in a hadronic medium that includes charm mesons, such as that produced by relativistic heavy-ion collisions.
In a short-distance production process, the separation between the long-distance scale $`1/|a|`$ and all the shorter-distance momentum scales can be exploited through a factorization formula that expresses the leading term in the production rate as the product of a long-distance factor that involves $`a`$ and a short-distance factor that is insensitive to $`a`$. To be definite, we will consider the specific production process $`BXK`$. The factorization for any other short-distance production process will have the same long-distance factor but a different short-distance factor.
There are many momentum scales that play an important role in the decay $`BXK`$, ranging from the extremely short-distance scales $`m_W`$ and $`m_b`$ associated with the quark decay process $`bc\overline{c}s`$ to the smaller short-distance scales $`\mathrm{\Lambda }_{\mathrm{QCD}}`$ and $`m_\pi `$ involved in formation of the final-state hadrons to the long-distance scale $`1/|a|`$ associated with the $`(DD^{})_+^0`$ wavefunction of the $`X`$. We denote the arbitrary boundary between the long-distance momentum region and the short-distance momentum region by $`\mathrm{\Lambda }`$.
The separation between the long-distance scale $`1/|a|`$ and all the short-distance momentum scales in the decay $`BXK`$ can be exploited through a factorization formula for the T-matrix element:
$`𝒯[BXK]=𝒜_{\mathrm{short}}[B(DD^{})_+^0K]\times 𝒜_X\sqrt{2m_X}.`$ (42)
In the long-distance factor, $`𝒜_X`$ is the universal amplitude in Eq. (33). The short-distance factor in Eq. (42) is insensitive to $`a`$ and one can therefore take the limit $`|a|\mathrm{}`$ in that factor. The factorization formula in (42) can serve as the definition of the short-distance factor. The content of the factorization statement then resides in the fact that the same short-distance factor appears in the factorization formula for the T-matrix element for the decay $`BD^0\overline{D}^0K`$ when the $`DD^{}`$ invariant mass is near the $`D^0\overline{D}^0`$ threshold. This factorization formula is discussed in Section V.
The factorization formula in Eq. (42) can be motivated diagrammatically by separating the loop integrals in the decay amplitude according to whether the virtual particles are off their energy shells by less than or by more than $`\mathrm{\Lambda }^2/(2\mu )`$. The leading terms in the decay amplitude for $`BXK`$ are suppressed only by a factor of $`(am_\pi )^{1/2}`$. These terms can be represented by the Feynman diagram in Fig. 4 and can be expressed in the form
$`𝒯[BXK]=\sqrt{2m_X}{\displaystyle \frac{d^3p}{(2\pi )^3}\psi (p)𝒜^{(\mathrm{\Lambda })}[B(DD^{})_+^0K]},`$ (43)
where $`\psi (p)`$ is the universal wavefunction in Eq. (5). The factor $`𝒜^{(\mathrm{\Lambda })}`$, which is represented by a dot in Fig. 4, is an amplitude for the decay $`B(DD^{})_+^0K`$ in which all virtual particles are off their energy shells by more than $`\mathrm{\Lambda }^2/(2\mu )`$. It is therefore insensitive to the relative momentum $`𝒑`$ of the $`D`$ and $`D^{}`$. If that momentum dependence is neglected, the wavefunction factor in Eq. (43) reduces to Eq. (37). The factorization formula in Eq. (42) then requires the short-distance factor to be
$`𝒜_{\mathrm{short}}[B(DD^{})_+^0K]=𝒜^{(\mathrm{\Lambda })}[B(DD^{})_+^0K]\left({\displaystyle \frac{\mu }{\pi ^2}}\mathrm{\Lambda }\right).`$ (44)
Since the T-matrix element in Eq. (42) is independent of the arbitrary separation scale, the dependence on $`\mathrm{\Lambda }`$ must cancel between the two factors on the right side of Eq. (44).
We proceed to use the factored expression in Eq. (42) to evaluate the decay rate for $`B^+XK^+`$. Lorentz invariance constrains the short-distance amplitude $`𝒜_{\mathrm{short}}`$ at the $`DD^{}`$ threshold to have the form
$`𝒜_{\mathrm{short}}[B^+(DD^{})_+^0K^+]=c_+Pϵ_D^{},`$ (45)
where $`P`$ is the 4-momentum of the $`B^+`$ and $`ϵ_D^{}`$ is the polarization 4-vector of the $`D^{}`$. Heavy quark spin symmetry guarantees that the polarization vector $`ϵ_D^{}`$ can be identified with the polarization vector $`ϵ_X`$ of the $`X`$. The decay rate is obtained by squaring the amplitude in Eq. (42), summing over the spin of the $`X`$, and integrating over phase space. The resulting expression for the decay rate is
$`\mathrm{\Gamma }[B^+XK^+]=|c_+|^2{\displaystyle \frac{\lambda ^{3/2}(m_B,m_X,m_K)}{32\pi m_B^3m_X}}|𝒜_X|^2,`$ (46)
where $`\lambda (x,y,z)`$ is the triangle function:
$`\lambda (x,y,z)=x^4+y^4+z^42(x^2y^2+y^2z^2+z^2x^2).`$ (47)
The long-distance factor $`|𝒜_X|^2`$ is given in Eq. (40). The result in Eq. (46) was obtained in Ref. Braaten:2004fk for the special case $`\gamma _{\mathrm{im}}=0`$ and used to estimate the order of magnitude of the decay rate for $`B^+XK^+`$. The estimate is consistent with the measurement of the product of the branching fractions for $`B^+XK^+`$ and $`XJ/\psi \pi ^+\pi ^{}`$ Choi:2003ue provided $`J/\psi \pi ^+\pi ^{}`$ is one of the major decay modes of $`X`$.
The factorization formula for the decay rate for $`B^0XK^0`$ has the same form as in Eq. (46) except that the coefficient $`c_+`$ in the short-distance decay amplitude in Eq. (45) has a different value. In Ref. Braaten:2004ai , it was pointed out that the decay rate for $`B^0XK^0`$ should be suppressed compared to $`B^+XK^+`$. That suppression can be understood by considering the short-distance amplitude for $`BDD^{}K`$. The dominant contributions to most decay amplitudes of the $`B`$ meson are believed to be factorizable into the product of matrix elements of currents. The factorizable contributions to the decay amplitude for $`B^+(DD^{})_+^0K^+`$ have three terms: the product of $`B^+\overline{D}^0`$ and $`\mathrm{}D^0K^+`$ matrix elements, where $`\mathrm{}`$ is the QCD vacuum, the product of $`B^+\overline{D}^0`$ and $`\mathrm{}D^0K^+`$ matrix elements, and the product of $`B^+K^+`$ and $`\mathrm{}(DD^{})_+^0`$ matrix elements. The factorizable contributions to the decay amplitude for $`B^0(DD^{})_+^0K^0`$ have only one term: the product of $`B^0K^0`$ and $`\mathrm{}(DD^{})_+^0`$ matrix elements. Heavy quark symmetry implies that the $`\mathrm{}(DD^{})_+^0`$ matrix element vanishes at the $`D^0\overline{D}^0`$ threshold. The decay $`B^0(DD^{})_+^0K^0`$ near the $`D^0\overline{D}^0`$ threshold must therefore proceed through nonfactorizable terms in the decay amplitude. The resulting suppression of the coefficient $`c_+`$ in the short-distance factor for $`B^0(DD^{})_+^0K^0`$ results in a suppression of the rate for $`B^0XK^0`$ relative to the rate for $`B^+XK^+`$. In Ref. Braaten:2004ai , a quantitative analysis of Babar data on the branching fractions for $`BD^{()}D^{()}K`$ Aubert:2003jq was used to estimate the suppression factor to be an order of magnitude or more.
## V Production of $`𝑫^\mathrm{𝟎}\overline{𝑫}^\mathbf{}\mathrm{𝟎}`$ Near Threshold
It was pointed out in Ref. Braaten:2004fk that the identification of $`X`$ as a $`DD^{}`$ molecule could be confirmed by observing a peak in the invariant mass distribution for $`D^0\overline{D}^0`$ (or $`D^0\overline{D}^0`$) near the $`DD^{}`$ threshold in the decay $`BDD^{}K`$. The shape of that invariant mass distribution was given for a real scattering length $`a`$. The shape would be the same for any other short-distance production process. In this section, we consider the effect of an imaginary part of the scattering length on the $`DD^{}`$ invariant mass distribution for a short-distance production process. To be specific, we consider the short-distance production process $`BD^0\overline{D}^0K`$.
The separation between the long-distance scale $`1/|a|`$ and all the short-distance momentum scales in the decay $`BD^0\overline{D}^0K`$ can be exploited through a factorization formula for the T-matrix element:
$`𝒯[BD^0\overline{D}^0K]=𝒜_{\mathrm{short}}[B(DD^{})_+^0K]\times 𝒜(E)\sqrt{4m_{D^0}m_{D^0}}{\displaystyle \frac{1}{\sqrt{2}}}.`$ (48)
In the long-distance factor, $`𝒜(E)`$ is the universal amplitude in Eq. (35) and the factor $`1/\sqrt{2}`$ is the amplitude for $`D^0\overline{D}^0`$ to be in the channel $`(DD^{})_+^0`$ with the large scattering length. The short-distance factor $`𝒜_{\mathrm{short}}`$ is the same as in the factorization formula for $`BXK`$ in Eq. (42).
The factorization formula in Eq. (48) can be motivated diagrammatically by separating the loop integrals in the decay amplitude according to whether the virtual particles are off their energy shells by less than or by more than $`\mathrm{\Lambda }^2/(2\mu )`$. There are terms in the T-matrix element for the decay $`BD^0\overline{D}^0K`$ that are enhanced near the $`DD^{}`$ threshold by a factor of $`am_\pi `$. These terms can be represented by the Feynman diagrams in Fig. 5 and can be expressed in the form
$`𝒯[BD^0\overline{D}^0K]=𝒜^{(\mathrm{\Lambda })}[B(DD^{})_+^0K]L(\mathrm{\Lambda },E)𝒜(E)\sqrt{4m_{D^0}m_{D^0}}{\displaystyle \frac{1}{\sqrt{2}}}.`$ (49)
The first factor $`𝒜^{(\mathrm{\Lambda })}`$, which is represented by the dot in Fig. 5, is an amplitude for the decay into $`(DD^{})_+^0K`$ in which all virtual particles are off their energy shells by more than $`\mathrm{\Lambda }^2/(2\mu )`$. It is therefore insensitive to the relative momentum $`𝒑`$. The second factor $`L(\mathrm{\Lambda },E)`$ is the amplitude for the propagation of the $`D`$ and $`D^{}`$ between contact interactions, which is given in Eq. (9). The condition $`|E|\mathrm{\Lambda }^2/(2\mu )`$ implies that the $`\sqrt{2\mu E}`$ term in $`L(\mathrm{\Lambda },E)`$ can be neglected. The resulting expression for the T-matrix element is the factorization formula in Eq. (48), with the short-distance factor $`𝒜_{\mathrm{short}}`$ given in Eq. (44).
We proceed to use the factorized expression for the decay amplitude in Eq. (48) to calculate the $`DD^{}`$ invariant mass distribution near the $`D^0\overline{D}^0`$ threshold in the decay $`B^+D^0\overline{D}^0K^+`$. Lorentz invariance constrains the short-distance amplitude $`𝒜_{\mathrm{short}}`$ at the $`DD^{}`$ threshold to have the form in Eq. (45). The decay rate is obtained by squaring the amplitude in Eq. (48), summing over the spins of the $`\overline{D}^0`$, and integrating over phase space. The resulting expression for the differential decay rate with respect to the $`DD^{}`$ invariant mass $`M`$ near the $`D^0\overline{D}^0`$ threshold is
$`{\displaystyle \frac{d\mathrm{\Gamma }}{dM}}[B^+D^0\overline{D}^0K^+]=|c_+|^2{\displaystyle \frac{\mu \lambda ^{3/2}(m_B,M,m_K)}{256\pi ^3m_B^3M^2}}\lambda ^{1/2}(M,m_{D^0},m_{D^0})|𝒜(E)|^2,`$ (50)
where $`E`$ is the energy of the $`D^0\overline{D}^0`$ in its rest frame relative to the $`D^0\overline{D}^0`$ threshold:
$`E=M(m_{D^0}+m_{D^0}).`$ (51)
We have used the fact that $`M`$ is near the $`D^0\overline{D}^0`$ threshold to replace a factor of $`m_{D^0}m_{D^0}`$ in Eq. (50) by $`\mu M`$. The result in Eq. (50) was obtained previously in Ref. Braaten:2004fk for the special case $`\gamma _{\mathrm{im}}=0`$. For $`M`$ near the $`D^0\overline{D}^0`$ threshold, the only significant variation with $`M`$ is through the long-distance factor $`|𝒜(E)|^2`$ and the threshold factor
$`\lambda ^{1/2}(M,m_{D^0},m_{D^0})2M\sqrt{2\mu E}.`$ (52)
If the complex scattering length is parameterized as in Eq. (12), the long-distance factor is
$`|𝒜(E)|^2={\displaystyle \frac{4\pi ^2/\mu ^2}{((2\mu E)^{1/2}+\gamma _{\mathrm{im}})^2+\gamma _{\mathrm{re}}^2}}.`$ (53)
The shape of the $`D^0\overline{D}^0`$ invariant mass distribution in Eq. (50) is given by the factor $`\sqrt{2\mu E}|𝒜(E)|^2`$. Note that it depends on $`\gamma _{\mathrm{re}}`$ and $`\gamma _{\mathrm{im}}`$ but not on the sign of $`\gamma _{\mathrm{re}}`$. The invariant mass distribution is shown in Fig. 6 for $`\gamma _{\mathrm{im}}=10`$ MeV and three values of $`|\gamma _{\mathrm{re}}|`$: 10, 20, and 30 MeV. The peak in the invariant mass distribution occurs at $`E=|\gamma |^2/(2\mu )`$, where $`|\gamma |=\sqrt{\gamma _{\mathrm{re}}^2+\gamma _{\mathrm{im}}^2}`$. The value at the peak is proportional to $`(|\gamma |+\gamma _{\mathrm{im}})^1`$. The full width at half maximum is $`2(2|\gamma |+\gamma _{\mathrm{im}})[(|\gamma |+\gamma _{\mathrm{im}})(3|\gamma |+\gamma _{\mathrm{im}})]^{1/2}/\mu `$.
## VI The $`𝑿`$ Line Shape
The $`X`$ is observed as a peak in the invariant mass distribution of its decay products, such as $`J/\psi \pi ^+\pi ^{}`$. Its mass and width are extracted from that invariant mass distribution. For instance, the Belle collaboration obtained their value for the mass and the upper bound on the width by fitting the $`J/\psi \pi ^+\pi ^{}`$ invariant mass distribution in $`B^+J/\psi \pi ^+\pi ^{}K^+`$ near the $`D^0\overline{D}^0`$ threshold to a resolution-broadened Breit-Wigner function on top of a polynomial background. The shape of the invariant mass distribution of the decay products of the $`X`$ is called the line shape. The resonant interactions in the $`D^0\overline{D}^0/D^0\overline{D}^0`$ system can significantly modify the line shape, so it need not have the conventional Breit-Wigner form. In this section, we compute the line shape of the $`X`$ in short-distance decays of the $`X`$. To be definite, we consider the production process $`BHK`$, where $`H`$ is the hadronic system consisting of $`J/\psi \pi ^+\pi ^{}`$ with invariant mass near the $`D^0\overline{D}^0`$ threshold. However our results on the line shape will apply more generally to any short-distance production process for $`X`$ and any short-distance decay mode of $`X`$.
The separation between the long-distance scale $`1/|a|`$ and all the short-distance momentum scales in the decay $`BHK`$ can be exploited through a factorization formula for the T-matrix element:
$`𝒯[BHK]=𝒜_{\mathrm{short}}[B(DD^{})_+^0K]\times 𝒜(E)\times 𝒜_{\mathrm{short}}[(DD^{})_+^0H].`$ (54)
There is an implied sum over the spin states of the $`D^{}`$. The long-distance factor $`𝒜(E)`$ depends on the complex-valued scattering length $`a`$ and is given in Eq. (35). Its argument $`E`$ is the difference between the invariant mass $`M`$ of the hadronic system $`H`$ and the $`D^0\overline{D}^0`$ threshold, as given in Eq. (51). The short-distance factor $`𝒜_{\mathrm{short}}`$ associated with the initial state is the same one that appears in the factorization formulas for $`BXK`$ in Eq. (42) and for $`BD^0\overline{D}^0K`$ in Eq. (48). The short-distance factor $`𝒜_{\mathrm{short}}`$ associated with the final state is the same one that appears in the factorization formulas for $`XH`$ in Eq. (32) and for $`D^0\overline{D}^0H`$ in Eq. (34).
The factorization formula in Eq. (54) can be motivated diagrammatically by separating the loop integrals in the decay amplitude according to whether the virtual particles are off their energy shells by less than or by more than $`\mathrm{\Lambda }^2/(2\mu )`$. There are terms in the decay amplitude for $`BHK`$ that are enhanced by a factor of $`am_\pi `$ when the invariant mass of $`H`$ is near the $`D^0\overline{D}^0`$ threshold. These terms can be represented by the Feynman diagrams in Fig. 7 and can be expressed in the form
$`𝒯[BHK]`$ $`=`$ $`𝒜^{(\mathrm{\Lambda })}[B(DD^{})_+^0K]L(\mathrm{\Lambda },E)𝒜(E)L(\mathrm{\Lambda },E)𝒜^{(\mathrm{\Lambda })}[(DD^{})_+^0H].`$ (55)
The factors $`𝒜^{(\mathrm{\Lambda })}`$, which are represented by dots in Fig. 7, are amplitudes in which all virtual particles are off their energy shells by more than $`\mathrm{\Lambda }^2/(2\mu )`$. The factors of $`L(\mathrm{\Lambda },E)`$, which is given in Eq. (9), are the amplitudes for the propagation of the $`D`$ and $`D^{}`$ between contact interactions. The condition $`|E|\mathrm{\Lambda }^2/(2\mu )`$ implies that the $`\sqrt{2\mu E}`$ term in $`L(\mathrm{\Lambda },E)`$ can be neglected. The resulting expression for the T-matrix element is the factorization formula in Eq. (54), with the short-distance factors $`𝒜_{\mathrm{short}}`$ given in Eqs. (44) and (38).
The factorization formula for the T-matrix element in (54) implies a factorization formula for the invariant mass distribution for the hadronic system $`H`$ near the $`D^0\overline{D}^0`$ threshold. If the hadronic system consists of particles with momenta $`p_i`$ and invariant mass $`M`$, the factorization formula is
$`{\displaystyle \frac{d\mathrm{\Gamma }}{dM}}[B^+K^+H]`$ (56)
$`=`$ $`{\displaystyle \frac{M}{2\pi m_B}}{\displaystyle \left|𝒜_{\mathrm{short}}[B(DD^{})_+^0K]\right|^2(2\pi )^4\delta ^4(P_BP_HP_K)\frac{d^3P_H}{(2\pi )^32E_H}\frac{d^3P_K}{(2\pi )^32E_K}}`$
$`\times |𝒜(E)|^2\times {\displaystyle \left|𝒜_{\mathrm{short}}[(DD^{})_+^0H]\right|^2(2\pi )^4\delta ^4(P_H_ip_i)\underset{i}{}\frac{d^3p_i}{(2\pi )^32E_i}}.`$
For $`M`$ near the $`D^0\overline{D}^0`$ threshold, the only significant variation with $`M`$ is through the long-distance factor $`|𝒜(E)|^2`$, where $`E`$ is the energy defined in Eq. (51). If the complex scattering length is parameterized as in Eq. (12), the long-distance factor is
$`\left|𝒜(E)\right|^2`$ $`=`$ $`{\displaystyle \frac{4\pi ^2/\mu ^2}{(|2\mu E|^{1/2}\gamma _{\mathrm{re}})^2+\gamma _{\mathrm{im}}^2}}E0,`$ (57a)
$`=`$ $`{\displaystyle \frac{4\pi ^2/\mu ^2}{((2\mu E)^{1/2}+\gamma _{\mathrm{im}})^2+\gamma _{\mathrm{re}}^2}}E0.`$ (57b)
This factor gives the line shape of the $`X`$. Note that for $`E>0`$, the line shape does not depend on the sign of $`\gamma _{\mathrm{re}}`$. However for $`E<0`$, the line shape is completely different for $`\gamma _{\mathrm{re}}>0`$ and $`\gamma _{\mathrm{re}}<0`$.
In the case $`\gamma _{\mathrm{re}}>0`$, the peak in the invariant mass distribution occurs below the $`D^0\overline{D}^0`$ threshold by the amount $`\gamma _{\mathrm{re}}^2/(2\mu )`$. The $`X`$ line shape is illustrated in the upper panel of Fig. 8 for $`\gamma _{\mathrm{im}}=10`$ MeV and for three positive values of $`\gamma _{\mathrm{re}}`$: 10, 20, and 30 MeV. If $`\gamma _{\mathrm{im}}<\gamma _{\mathrm{re}}`$, the full width of the peak at half maximum is $`2\gamma _{\mathrm{re}}\gamma _{\mathrm{im}}/\mu `$. The line shape for $`E<0`$ is symmetric about the peak as a function of $`|E|^{1/2}`$ but not as a function of $`E`$. If $`\gamma _{\mathrm{im}}\gamma _{\mathrm{re}}`$, the line shape in Eq. (57) is sharply peaked at $`E=\gamma _{\mathrm{re}}^2/(2\mu )`$ and it can be approximated by a delta function:
$`|𝒜(E)|^2{\displaystyle \frac{4\pi ^3\gamma _{\mathrm{re}}}{\mu ^3\gamma _{\mathrm{im}}}}\delta \left(E+\gamma _{\mathrm{re}}^2/(2\mu )\right).`$ (58)
Note that the condition $`\gamma _{\mathrm{im}}\gamma _{\mathrm{re}}`$ is equivalent to $`\mathrm{\Gamma }_XE_X/4`$.
In the case $`\gamma _{\mathrm{re}}<0`$, the peak in the invariant mass distribution occurs at the $`D^0\overline{D}^0`$ threshold. The $`X`$ line shape is illustrated in the lower panel of Fig. 8 for $`\gamma _{\mathrm{im}}=10`$ MeV and for three negative values of $`\gamma _{\mathrm{re}}`$: $`10`$, $`20`$, and $`30`$ MeV. The line shape has a cusp at $`E=0`$. Bugg has proposed that the $`X`$ can be identified with this cusp at the $`D^0\overline{D}^0`$ threshold Bugg:2004rk . The normalization is the same in the upper and lower panels of Fig. 8. Note that the area under the cusp in the lower panel of Fig. 8 is much smaller than the area under the resonance in the upper panel for the same values of $`\gamma _{\mathrm{im}}`$ and $`|\gamma _{\mathrm{re}}|`$. Thus a cusp seems less likely as an interpretation for the $`X(3872)`$ than a resonance, although a quantitative analysis would be required to rule out that possibility.
The integral over all energies of the line shape of a conventional Breit-Wigner resonance is convergent. In contrast, the integral of the line shape in Eq. (57) diverges logarithmically as the endpoints $`E_{\mathrm{min}}`$ and $`E_{\mathrm{max}}`$ of the integral increase in magnitude. This follows from the fact that the line shape in Eq. (57) decreases as $`1/|E|`$ for $`(2\mu |E|)^{1/2}|\gamma _{\mathrm{re}}|,\gamma _{\mathrm{im}}`$. That expression for the line shape is of course only accurate for $`|E|`$ lower than $`\mathrm{\Lambda }^2/(2\mu )10`$ MeV, where $`\mathrm{\Lambda }m_\pi `$ is the natural momentum scale for low-energy $`DD^{}`$ scattering. Thus the logarithmic dependence on $`E_{\mathrm{min}}`$ and $`E_{\mathrm{max}}`$ holds only for $`|E_{\mathrm{min}}|,E_{\mathrm{max}}<\mathrm{\Lambda }^2/(2\mu )`$. It is convenient to define $`p_{\mathrm{min}}`$ and $`p_{\mathrm{max}}`$ by $`E_{\mathrm{min}}=p_{\mathrm{min}}^2/(2\mu )`$ and $`E_{\mathrm{max}}=+p_{\mathrm{max}}^2/(2\mu )`$. The integral of the factor in Eq. (57) reduces in the limit $`p_{\mathrm{min}},p_{\mathrm{max}}|\gamma _{\mathrm{re}}|,\gamma _{\mathrm{im}}`$ to
$`{\displaystyle _{E_{\mathrm{min}}}^{E_{\mathrm{max}}}}|𝒜(E)|^2𝑑E{\displaystyle \frac{4\pi ^2}{\mu ^3}}\left(\mathrm{log}{\displaystyle \frac{p_{\mathrm{min}}p_{\mathrm{max}}}{\gamma _{\mathrm{re}}^2+\gamma _{\mathrm{im}}^2}}+{\displaystyle \frac{\pi \gamma _{\mathrm{re}}}{\gamma _{\mathrm{im}}}}\theta (\gamma _{\mathrm{re}})f(\gamma _{\mathrm{re}}/\gamma _{\mathrm{im}})\right),`$ (59)
where $`f(x)`$ is the function
$`f(x)=x\mathrm{arctan}(1/x)+(1/x)\mathrm{arctan}(x).`$ (60)
This function has a limited range, varying from 1 at $`x=0`$ and $`x=\pm \mathrm{}`$ to $`\pi /2`$ at $`x=\pm 1`$.
The factorization formula for the invariant mass distribution of $`H`$ in the case $`\gamma _{\mathrm{re}}>0`$ has important implications for measurements of the branching fractions of $`X`$. Since the two short-distance factors in Eq. (56) are insensitive to $`E`$, we can set $`M`$ to $`m_{D^0}+m_{D^0}`$ or to $`m_X`$ in those factors. The short-distance factor associated with the decay of the $`B^+`$ reduces to $`\mathrm{\Gamma }[BXK]/(2\pi |𝒜_X|^2)`$. If $`\gamma _{\mathrm{re}}>0`$, the short-distance factor associated with the formation of $`H`$ reduces to $`\mathrm{\Gamma }[XH]/|𝒜_X|^2`$. Thus the differential decay rate in Eq. (56) reduces to
$`{\displaystyle \frac{d\mathrm{\Gamma }}{dM}}[BHK]=\mathrm{\Gamma }[BXK]\mathrm{Br}[XH]{\displaystyle \frac{\mathrm{\Gamma }_X|𝒜(E)|^2}{2\pi |𝒜_X|^4}}.`$ (61)
If the product of $`\mathrm{\Gamma }[BXK]`$ and $`\mathrm{Br}[XH]`$ is measured by integrating $`d\mathrm{\Gamma }/dM`$ over the energy interval from $`E_{\mathrm{min}}`$ to $`E_{\mathrm{max}}`$ with $`|E_{\mathrm{min}}|,E_{\mathrm{max}}E_X,\mathrm{\Gamma }_X`$, it will be in error by the factor
$`{\displaystyle \frac{\mathrm{\Gamma }_X}{2\pi |𝒜_X|^4}}{\displaystyle _{E_{\mathrm{min}}}^{E_{\mathrm{max}}}}|𝒜(E)|^2𝑑E`$ $``$ $`\left[1+{\displaystyle \frac{\mathrm{\Gamma }_X^2}{16E_X^2}}\right]^1`$ (62)
$`\times \left[1+\left(\mathrm{log}{\displaystyle \frac{|E_{\mathrm{min}}|^{1/2}E_{\mathrm{max}}^{1/2}}{E_X+\mathrm{\Gamma }_X^2/(16E_X)}}f(4E_X/\mathrm{\Gamma }_X)\right){\displaystyle \frac{\mathrm{\Gamma }_X}{4\pi E_X}}\right].`$
The error would cancel in the ratio of the branching fractions for any two short-distance decay modes of $`X`$. The error would not cancel in the ratio of the branching fractions for a short-distance decay mode of $`X`$ and one of the long-distance decay modes $`D^0\overline{D}^0\pi ^0`$ and $`D^0\overline{D}^0\gamma `$. This effect should be taken into account in analyzing the decays of the $`X(3872)`$.
## VII Summary
If the $`X(3872)`$ is a loosely-bound S-wave molecule corresponding to a $`C=+`$ superposition of $`D^0\overline{D}^0`$ and $`D^0\overline{D}^0`$, the scattering length $`a`$ in the $`(DD^{})_+^0`$ channel is large compared to all other length scales of QCD. The decays of the $`X`$ implies that the large scattering length has an imaginary part. It can be conveniently parameterized in terms of the real and imaginary parts of $`1/a`$ as in Eq. (12). The binding energy $`E_X`$ and the width $`\mathrm{\Gamma }_X`$ are expressed in terms of those parameters in Eqs. (14) and (13b).
The large scattering length can be exploited through factorization formulas for decay rates of $`X`$. For short-distance decay modes that do not proceed through the decay of a constituent $`D^{}`$ of the $`X`$, the long-distance factor in the factorization formula is proportional to $`1/|a|`$ and is given in Eq. (41). If a partial width of the $`X`$ is calculated using some model with a specific binding energy for the $`X`$, the factorization formulas can be used to extrapolate the prediction to other values of the binding energy and to take into account the width of the $`X`$.
The large scattering length can also be exploited through factorization formulas for production rates of $`X`$, $`D^0\overline{D}^0`$ near threshold, $`D^0\overline{D}^0`$ near threshold, and decay products of $`X`$ with invariant mass near the $`D^0\overline{D}^0`$ threshold. The long-distance factor in the factorization formula for production rates of $`X`$ is proportional to $`1/|a|`$. For production of $`D^0\overline{D}^0`$ and $`D^0\overline{D}^0`$ near threshold, the factorization formula implies that the dependence on the invariant mass is through the factor $`\sqrt{2\mu E}|𝒜(E)|^2`$, where $`E`$ is the invariant mass relative to the $`D^0\overline{D}^0`$ threshold and $`|𝒜(E)|^2`$ is given in Eq. (53). The peak in the invariant mass distribution is above the threshold by the amount $`E_X+\mathrm{\Gamma }_X^2/(16E_X)`$. The line shape of the $`X`$ can be measured through the invariant mass distribution of its decay products. In the case of short-distance decay modes, the factorization formulas imply that near the $`D^0\overline{D}^0`$ threshold, the shape of the invariant mass distribution is given by the factor $`|𝒜(E)|^2`$ in Eq. (57). If Re$`(a)<0`$, the distribution has a cusp at $`E=0`$, as shown in the lower panel of Fig. 8. If Re$`(a)>0`$, the distribution has a peak at $`E=E_X`$, as shown in the upper panel of Fig. 8. In contrast to a Breit-Wigner resonance, the integral over the line shape is logarithmically sensitive to the endpoints of the integration region. This effect should be taken into account in analyzing the production and decay of the $`X`$.
###### Acknowledgements.
This research was supported in part by the Department of Energy under grant DE-FG02-91-ER4069. |
warning/0506/hep-ph0506206.html | ar5iv | text | # Little Higgs Models and Electroweak Measurements
## I Introduction
Little Higgs models Arkani-Hamed:2001nc ; Arkani-Hamed:2002qx ; littlest ; su6 ; Schmaltz:2002wx ; Kaplan:2003uc ; Chang:2003un ; Skiba:2003yf ; Chang:2003zn ; Cheng:2003ju ; Cheng:2004yc ; Kaplan:2004cr ; simplest ; Kong:2004cv ; Low:2004xc ; Thaler:2005en have been proposed to stabilize the electroweak scale in the Standard Model (SM), see Ref. Schmaltz:2005ky for a review. In these models, the one-loop quadratically-divergent corrections to the Higgs mass from the SM particles are canceled by the corrections from new particles with masses at TeV scale. Eliminating the one-loop divergences allows the cutoff of the theory to be pushed to about $`4\pi `$ TeV, however this could be an optimistic estimate, see Ref. Chang:2003vs . The predicted particles are likely to be produced and observed in future colliders, especially the LHC. The presence of new particles also creates tension with the electroweak precision tests (EWPTs). To avoid fine-tuning of more than 10%, heavy fermions and gauge bosons with masses less than about 2 TeV and 5 TeV, respectively, should be introduced to cancel the top-loop and gauge-boson-loop divergences. However, the EWPTs do not indicate presence of new particles in a few TeV range if their couplings are generic. A successful model has to reconcile the tension between naturalness and the EWPTs.
Constraints on little Higgs models from the EWPTs have been considered for different models Csaki:2002qg ; Hewett:2002px ; Han:2003wu ; Csaki:2003si ; Gregoire:2003kr ; Chen:2003fm ; Casalbuoni:2003ft ; Kilian:2003xt ; Marandella:2005wd . In this article, we provide a more up-to-date and extensive analysis employing the effective theory approach we described in Ref. Han:2004az . In Ref. Han:2004az , we analyzed all flavor-independent and CP-conserving dimension-six operators written in terms of the SM fields that are tightly constrained by EWPTs. We calculated the corrections from these operators to the electroweak precision observables (EWPOs). The result is the $`\chi ^2`$ distribution in terms of the coefficients of these operators. In an extension of the SM one can integrate out the heavy fields and obtain the coefficients of the effective operators in terms of parameters in the model. Substituting the coefficients in the $`\chi ^2`$ distribution, one immediately obtains global constraints from all EWPOs. As we will discuss, this procedure fits most little Higgs models. We will only consider the tree-level diagrams when we integrate out the heavy fields. Loop diagrams involving heavy fields are usually suppressed by both the masses of the heavy fields and the loop factor. Thus, loop corrections from particles with TeV-scale masses are usually small and do not significantly affect the constraints we obtain. For the discussion of the fine-tuning problem, we will focus on the largest one-loop corrections to the Higgs mass arising from the top quark and the gauge bosons. Alternative estimates of fine-tuning associated with the sensitivity of the Higgs mass to all underlying parameters of little Higgs models are presented in Ref. Casas:2005ev .
In this article, we focus on the following little Higgs models: the $`SU(5)/SO(5)`$ or the littlest Higgs model littlest , the $`SU(6)/SP(6)`$ su6 model, and the models with the $`SU(3)\times U(1)`$ gauge group Kaplan:2003uc ; Skiba:2003yf ; simplest , as well as their variations. For simplicity, from now on we will refer to them as $`SU(5)`$, $`SU(6)`$ and $`SU(3)`$ models respectively, although the first two refer to their global symmetries and the last one refers to its gauge symmetry. The $`SU(3)`$ little Higgs models can have different global symmetries. In Sec. II, we discuss in general what kind of operators we expect from these models and how to constrain them. Secs. III, IV, V are devoted to detailed discussion of each of the three models. We summarize our results in Sec. VI.
## II Integrating out heavy fields
A complete set of independent dimension-six operators in the SM is given in Ref. Buchmuller:1985jz . Assuming flavor and CP conservation, in Ref. Han:2004az we narrowed this set down to 21 operators that are relevant to EWPTs. In a compact notation, the operators are:
$`O_h`$ $`=`$ $`|h^{}D_\mu h|^2,`$
$`O_{hf}^s`$ $`=`$ $`i(h^{}D^\mu h)(\overline{f}\gamma _\mu f)+\mathrm{h}.\mathrm{c}.,O_{hf}^t=i(h^{}\sigma ^aD^\mu h)(\overline{f}\gamma _\mu \sigma ^af)+\mathrm{h}.\mathrm{c}.,`$
$`O_{ff^{}}^s`$ $`=`$ $`{\displaystyle \frac{1}{1+\delta _{ff^{}}}}(\overline{f}\gamma ^\mu f)(\overline{f^{}}\gamma _\mu f^{}),O_{ff^{}}^t={\displaystyle \frac{1}{1+\delta _{ff^{}}}}(\overline{f}\gamma ^\mu \sigma ^af)(\overline{f^{}}\gamma _\mu \sigma ^af^{}),`$ (1)
where $`h`$ is the SM Higgs doublet, $`f,f^{}=q,l,u,d,e`$, are the left and right handed fermions. The operators are understood to be summed over flavor indices. The superscripts $`s`$ and $`t`$ stand for singlet and triplet $`SU(2)`$ contractions. For triplet couplings, $`f`$ has to be a SM doublet ($`q`$ or $`l`$). Note that $`O_h`$ corresponds to the oblique $`T`$ parameter Peskin:1991sw that breaks the custodial symmetry. We have omitted two operators from our list. The omitted operators are denoted $`O_{WB}`$ and $`O_W`$ in Ref. Han:2004az . The former corresponds to the oblique $`S`$ parameter and the latter modifies triple gauge-boson couplings. These two operators are not induced at tree level in the models analyzed here. Four-fermion operators involving only quark fields were not included in our list because they are not constrained as tightly as operators involving some leptons.
Including operators $`O_i`$, we can write the effective Lagrangian as
$$=_{SM}+\underset{i}{}a_iO_i.$$
(2)
In Ref. Han:2004az we calculated the corrections to EWPOs from the operators $`O_i`$ assuming arbitrary coefficients $`a_i`$. The corrections were combined with known experimental values and the SM predictions to obtain the total $`\chi ^2`$ distribution
$$\chi ^2=\chi ^2(a_i)=\chi _{min}^2+(a_i\widehat{a}_i)_{ij}(a_j\widehat{a}_j),$$
(3)
where $`\widehat{a_i}`$ are values of $`a_i`$ that minimize $`\chi ^2`$. In Ref. Han:2004az we calculated the corrections to EWPOs to linear order in $`a_i`$, so $``$ in Eq. (3) is a constant and positive-definite matrix. In a given model, we integrate out the heavy fields and obtain the coefficients $`a_i`$ as functions of the parameters in the model. Expressing $`a_i`$ in Eq. (3) in terms of parameters of a model allows us to immediately obtain constraints on the model without having to compute EWPOs.
Before we list and analyze in detail the coefficients of effective operators in the little Higgs models, we briefly discuss how the effective operators are generated and what interesting features these models have.
There are three kinds of heavy fields in these models: gauge bosons, scalars and fermions. We first discuss the effects of the heavy gauge bosons. The $`SU(5)`$ and $`SU(6)`$ models share the same gauge structure: $`[SU(2)\times U(1)]^2`$. The gauge group is broken to the diagonal $`SU(2)\times U(1)`$ which is identified with the SM gauge group. Thus half of the gauge bosons get masses and the other half remain massless until electroweak symmetry breaking takes place. The heavy gauge bosons include a triplet $`W^{}`$ and a singlet $`Z^{}`$, which couple to both the SM Higgs and fermions. The exchange of $`W^{}`$ generates the triplet coupling operators $`O_{hf}^t`$, $`O_{ff^{}}^t`$, while the exchange of $`Z^{}`$ generates $`O_h`$, $`O_{hf}^s`$, $`O_{ff^{}}^s`$. The gauge sector in the $`SU(3)`$ model is quite different: a gauged $`SU(3)\times U(1)`$ is broken to the SM $`SU(2)\times U(1)`$, leaving 5 heavy gauge bosons. One of them behaves like the $`Z^{}`$ and also induces $`O_h`$, $`O_{hf}^s`$ and $`O_{ff^{}}^s`$. The others decouple from the light fields. There are no $`SU(2)`$ triplet operators generated in this case.
Now we turn to discussing the scalars in the models. In the low-energy spectra, the $`SU(5)`$ and $`SU(3)`$ models contain one Higgs doublet and the $`SU(6)`$ model contains two Higgs doublets. The number of light Higgs doublets is irrelevant in our analysis since only the Higgs vev matters. Besides the doublet, the $`SU(5)`$ model also contains a heavy triplet scalar. When integrated out, the triplet generates the $`O_h`$ operator. In addition, there are heavy singlets in all the three models, but integrating them out does not generate dimension-six operators relevant to EWPTs.
Turning to fermions, heavy fermions are needed to cancel the quadratic divergence from the top loop. However, there are often multiple choices that insure the cancelation. In the $`SU(5)`$ little Higgs model littlest , a pair of vector-like heavy fermions is added to the SM fields. The right-handed heavy fermion mixes with the right-handed top quark so that the couplings between the top quark and the SM gauge bosons are modified. The loop divergences from the light two generations do not introduce fine-tuning for a cutoff as low as $`4\pi `$ TeV because the corresponding Yukawa couplings are small. Thus one does not introduce extra fermions to cancel the divergences for the light two generations. Therefore, fermion couplings in this model are flavor-dependent. However, since only the top quark mixes with the heavy fermion and no EWPO involves the top quark in the final or initial states, this flavor-dependent effect is not relevant. Of course, if we added fermions that mix with the first two generations as well, we could generate flavor-dependent operators that do affect EWPTs. Such operators can introduce FCNCs and would be severely constrained. A detailed analysis for this case is beyond the scope of this paper. We will assume approximate flavor-independence in our analysis.
The fermion sector in the $`SU(6)`$ model is similar to the $`SU(5)`$ model<sup>3</sup><sup>3</sup>3In Ref. su6 , the Yukawa structure also introduces mixing for the bottom quark, but we can make a similar choice as in the SU(5) model. Sec. IV contains an example., while the $`SU(3)`$ model merits a few more comments. Because the gauge group contains an $`SU(3)`$, every SM fermion doublet must be combined with an extra fermion to complete an $`SU(3)`$ triplet. Thus unlike the $`SU(5)`$ or the $`SU(6)`$ model, heavy fermions have to be added to all generations. If we assign fermions the same quantum numbers for the three generations and impose the constraint of small FCNCs, we will obtain flavor-independent operators $`O_{hq}^{s,t}`$ and $`O_{hl}^{s,t}`$ by integrating out the heavy fermions. This is a result of mixing between the SM and the heavy fermions. This mixing modifies the gauge couplings of the SM fermions.
Besides the operators obtained by integrating out the heavy fields, there exist dimension-six operators arising from expanding the kinetic term of the nonlinear sigma field to higher orders. It turns out that the dimension-six operators obtained this way do not affect EWPTs in the $`SU(5)`$ and $`SU(3)`$ models, while in the $`SU(6)`$ model, there are contributions to the $`O_h`$ operator when $`\mathrm{tan}\beta 1`$, where $`\mathrm{tan}\beta `$ is the ratio of the two Higgs vevs.
## III The $`SU(5)/SO(5)`$ model littlest
Detailed description of the $`SU(5)`$ little Higgs model can be found elsewhere. We will only specify the necessary conventions and notation. Throughout this paper, we use $`h=(h^+\text{ }h^0)^T`$ to denote the Higgs doublet, $`v`$ the vev of the Higgs, and $`g`$ and $`g^{}`$ the gauge coupling constants of the SM. The littlest Higgs model is based on a nonlinear sigma model with an $`SU(5)`$ global symmetry spontaneously broken to its $`SO(5)`$ subgroup. The $`SU(5)`$ breaking direction is given by the vev of the $`\mathrm{\Sigma }`$ field
$$\mathrm{\Sigma }_0\mathrm{\Sigma }=\left(\begin{array}{ccc}& & \text{1}_2\\ & 1& \\ \text{1}_2& & \end{array}\right).$$
(4)
$`\mathrm{\Sigma }`$ can be parameterized around its vev as
$$\mathrm{\Sigma }=e^{2i\mathrm{\Pi }/F}\mathrm{\Sigma }_0,$$
(5)
where the Goldstone boson matrix $`\mathrm{\Pi }`$ is defined as
$$\mathrm{\Pi }=\left(\begin{array}{ccc}0& \frac{\stackrel{~}{h}}{\sqrt{2}}& \varphi ^{}\\ \frac{\stackrel{~}{h}^{}}{\sqrt{2}}& 0& \frac{\stackrel{~}{h}^T}{\sqrt{2}}\\ \varphi & \frac{\stackrel{~}{h}^{}}{\sqrt{2}}& 0\end{array}\right).$$
(6)
In the equation above, $`\stackrel{~}{h}=i\sigma ^2h^{}`$, and $`\varphi `$ is a two by two symmetric matrix that represents a scalar triplet with mass $`M_\varphi `$ of order $`F`$. We have omitted the fields that are eaten when the gauge group is broken to the SM gauge group. The kinetic term of the nonlinear sigma model is given by
$$\frac{F^2}{8}\text{Tr}[D_\mu \mathrm{\Sigma }D^\mu \mathrm{\Sigma }^{}],$$
(7)
where the covariant derivative is defined as
$$D_\mu \mathrm{\Sigma }=_\mu \mathrm{\Sigma }i\underset{j=1}{\overset{2}{}}[g_jW_j^a(Q_j^a\mathrm{\Sigma }+\mathrm{\Sigma }Q_j^{aT})+g_j^{}B_j(Y_j\mathrm{\Sigma }+\mathrm{\Sigma }Y_j^T)].$$
(8)
$`Q_j^a`$ and $`Y_j`$ are the generators of the $`[SU(2)\times U(1)]^2`$ gauge group, defined in the same way as in Ref. littlest . $`W_j^a`$ and $`B_j`$ are the corresponding gauge bosons. The mass eigenstates of the gauge bosons are
$`W`$ $`=`$ $`sW_1+cW_2,W^{}=cW_1+sW_2,`$
$`B`$ $`=`$ $`s^{}B_1+c^{}B_2,Z^{}=c^{}B_1+s^{}B_2,`$ (9)
where
$$s,c=\frac{g_2,g_1}{\sqrt{g_1^2+g_2^2}},s^{},c^{}=\frac{g_2^{},g_1^{}}{\sqrt{g_1^2+g_2^2}}.$$
(10)
In Eq. (9), $`W`$ and $`B`$ are the SM gauge bosons, and $`W^{}`$ and $`Z^{}`$ are heavy gauge bosons with masses
$$M_W^{}=\frac{gF}{2sc},M_Z^{}=\frac{g^{}F}{\sqrt{20}s^{}c^{}},$$
(11)
where $`g=g_1s=g_2c`$ and $`g^{}=g_1^{}s^{}=g_2^{}c^{}`$
Most corrections to EWPOs come from the exchanges of the $`W^{}`$ or $`Z^{}`$ bosons. Another correction comes from integrating out the triplet field $`\varphi `$ which couples to $`h`$ as
$$i\lambda (h^T\varphi ^{}hh^{}\varphi h^{}),$$
(12)
where $`\lambda `$ is a dimensionful coupling of order $`F`$. The coefficient $`\lambda `$ is not determined in the low-energy theory since it arises from quadratically-divergent contributions. When integrated out, the triplet generates both dimension-four and dimension-six operators. The dimension-four operator is the $`|h|^4`$ term contributing to the Higgs potential. In order to insure that the Higgs potential is bounded from below and generates the correct vev for EWSB, the following relation has to be satisfied Han:2003wu
$$\frac{\lambda ^2F^2}{M_\varphi ^4}<\frac{1}{4}.$$
(13)
Besides the dimension-four operator, the triplet also generates the dimension-six $`O_h`$ with the coefficient $`2\lambda ^2/M_\varphi ^4`$. We will treat $`\lambda ^2/M_\varphi ^4`$ as a free parameter that is subject to the constraint (13). Note that if we did not integrate out the triplet, after EWSB the triplet would obtain a vev that breaks the custodial symmetry:
$$v^{}=\varphi =\frac{\lambda v^2}{2M_\varphi ^2}.$$
(14)
This is another way to understand the contribution to $`O_h`$.
We assume that fermions are all singlets under the second $`SU(2)`$ as in Ref. littlest , but allow them to be charged under both $`U(1)`$’s. The model includes a pair of vector-like heavy fermions to cancel the top loop contribution to the Higgs mass. As explained in Sec. II, the heavy fermions do not affect EWPOs.
Integrating out the heavy fields $`W^{}`$, $`Z^{}`$ and $`\varphi `$, we obtain:
$`a_h`$ $`=`$ $`{\displaystyle \frac{5(c^2s^2)^2}{2F^2}}+{\displaystyle \frac{2\lambda ^2}{M_\varphi ^4}},`$
$`a_{hq}^t`$ $`=`$ $`a_{hl}^t={\displaystyle \frac{(c^2s^2)c^2}{2F^2}},`$
$`a_{hf}^s`$ $`=`$ $`{\displaystyle \frac{5s^{}c^{}(c^2s^2)}{F^2}}\left(Y_2^f{\displaystyle \frac{s^{}}{c^{}}}Y_1^f{\displaystyle \frac{c^{}}{s^{}}}\right),`$
$`a_{lq}^t`$ $`=`$ $`a_{ll}^t={\displaystyle \frac{c^4}{F^2}},`$
$`a_{ff^{}}^s`$ $`=`$ $`{\displaystyle \frac{20s^2c^2}{F^2}}\left(Y_2^f{\displaystyle \frac{s^{}}{c^{}}}Y_1^f{\displaystyle \frac{c^{}}{s^{}}}\right)\left(Y_2^f^{}{\displaystyle \frac{s^{}}{c^{}}}Y_1^f^{}{\displaystyle \frac{c^{}}{s^{}}}\right),`$ (15)
where $`Y_1^f`$ and $`Y_2^f`$ are the charges of fermion $`f`$ under the two $`U(1)`$’s. $`Y_1^f`$ and $`Y_2^f`$ are assumed to be generation independent. The SM hypercharge is $`Y=Y_1+Y_2`$. If the $`U(1)`$ charge assignment is given, we can substitute the coefficients $`a_i`$ into Eq. (3) and obtain $`\chi ^2`$ as a function of $`f`$, $`c`$, $`c^{}`$ and $`\lambda ^2/M_\varphi ^4`$.
As we will see shortly, the coefficients $`a_i`$ can put tight constraints on $`F`$. If $`F`$ is too large, we will reintroduce fine-tuning to the theory. Therefore it is interesting to consider how to choose the parameters in Eqs. (15) to suppress $`a_i`$. It is easy to see that if $`c1`$ and $`s^{}c^{}`$, the coefficients $`a_h`$, $`a_{hf}^s`$, $`a_{hf}^t`$ and $`a_{ff^{}}^t`$ are all suppressed. If we further assume that the fermions are charged equally under the two $`U(1)`$’s ($`Y_1^f=Y_2^f`$), the coefficients $`a_{ff^{}}^s`$ also vanish. However, it is impossible to render the two $`U(1)`$ charges equal unless one allows the $`U(1)`$’s to be outside of the global $`SU(5)`$. Changing $`Y_1Y_1+bI`$, and $`Y_2Y_2bI`$ and setting $`b=1/80`$ Csaki:2003si , we can obtain equal charges<sup>4</sup><sup>4</sup>4In this case, the coefficients listed in Eqs. (15) are modified because of the change in the $`Z^{}`$ mass. The coefficients $`a_{hf}^t`$ and $`a_{ff^{}}^t`$ do not change while $`a_{hf}^s`$, $`a_{ff^{}}^s`$ and the term multiplying $`(c^2s^2)^2`$ in $`a_h`$ are rescaled by a factor of $`1/(1+100b^2)`$.. These parameter choices are very much like the “near-oblique” limit for the $`SU(6)`$ model discussed in Ref. Gregoire:2003kr . In this limit, the only significant corrections come from the $`\lambda ^2/M_\varphi ^4`$ term in $`a_h`$. In addition, $`c`$ can not be arbitrarily small and has to be greater than $`g/4\pi `$, otherwise the second $`SU(2)`$ will be strongly coupled.
To make the constraints and thus the associated fine-tuning more transparent, we trade $`F`$ for the heavy top mass $`M_t^{}`$ and the heavy triplet gauge boson mass $`M_W^{}`$. For the Higgs mass of order 200 GeV, 10% fine-tuning corresponds to $`M_t^{}2`$ TeV or $`M_W^{}6`$ TeV. $`M_W^{}`$ is given in Eq. (11). As discussed in Ref. littlest , $`M_t^{}\sqrt{2}\lambda _tF`$, where $`\lambda _t1`$ is the top Yukawa coupling. For simplicity, we set $`M_t^{}=\sqrt{2}F`$. Fig. 1 shows the bounds on $`M_t^{}`$ and $`M_W^{}`$ around the “near-oblique” limit as functions of $`c`$ and $`\lambda ^2F^2/M_\varphi ^4`$. We have set $`s^{}=c^{}`$ and $`Y_1^f=Y_2^f`$ in the two plots. As expected, the bounds on $`M_t^{}`$ are loose near $`c=g/4\pi `$ and $`\lambda ^2f^2/M_\varphi ^4=0`$. On the other hand, because $`M_W^{}1/sc`$, the bounds on $`M_W^{}`$ are quite tight near $`c=g/4\pi `$, introducing more fine-tuning than that from the top sector. Nevertheless, there clearly exists a region with less than 10% fine-tuning. It would be desirable if the charge assignment and the parameter space limit could come naturally from a UV extension of the model.
Another way to suppress the coefficients in Eqs. (15) is by taking $`c,c^{}1`$ and assuming the fermions are charged under only the first $`U(1)`$ ($`Y_1^f=Y^f`$, $`Y_2^f=0`$). In this limit, all coefficients except $`a_h`$ are suppressed. It turns out that $`a_h`$ alone can still put tight constraints on $`F`$, as can be seen from Fig. 2. In Fig. 2, we plot the 95% confidence level (CL) bounds on $`M_t^{}`$ for two fermion charge assignments, as a function of $`c,c^{}`$. The left (middle) plot corresponds to all fermions charged under only the first (second) $`U(1)`$. In the plots, we allow both $`M_t^{}`$ and $`\lambda ^2f^2/M_\varphi ^4`$ to vary. The shown bounds on $`M_t^{}`$ are the minimal values for an arbitrary value of $`\lambda ^2f^2/M_\varphi ^4`$ within the \[0,1/4\] interval. For a two-parameter fit, $`\mathrm{\Delta }\chi ^2=5.99`$. The limit mentioned above corresponds to the region near the origin in the left plot.
The bounds shown in Fig. 2 are quite stringent. For all of the parameter space, the bounds on $`M_t^{}`$ exceed $`6`$ TeV, introducing fine-tuning of more than 1%. This is also true for the parameter limit discussed in the previous paragraph, which makes this limit seem uninteresting. However, as proposed in Ref. Chang:2003zn , it is possible to enlarge one of the $`U(1)`$’s to $`SU(2)`$ and make its coupling relatively strong so that there exists an approximate custodial symmetry that suppresses the coefficient $`a_h`$.
It is worth mentioning that LEP2 data included in our fit contributes significantly to the constraints. For comparison, the plot on the right in Fig. 2 shows the bounds from data excluding LEP2 measurements for the $`Y^f=Y_1^f`$ case. The bounds are significantly relaxed compared with the bounds obtained using all data. This is because LEP2 experiments are very sensitive to the 4-fermion operators $`O_{ff^{}}^{s,t}`$ generated by $`Z^{}`$ and $`W^{}`$ exchanges, while LEP1 and other measurements are only sensitive to a few of them.
Other modifications of the model include gauging only $`U(1)_Y`$ Csaki:2003si and applying a $`T`$ parity Low:2004xc . The former is similar to the first limit discussed above, where the coefficients $`a_h`$, $`a_{hf}^s`$, $`a_{ff^{}}^s`$ are suppressed due to the lack of the $`Z^{}`$ boson contribution. The latter avoids generating operators at tree level and thus constraints on the model from EWPTs are less stringent.
## IV The $`SU(6)/SP(6)`$ model su6
The $`SU(6)`$ model has the same gauge structure as the $`SU(5)`$ model but different global symmetry. The nonlinear sigma model terms have the same form as in Eqs. (7) and (8). The $`\mathrm{\Sigma }_0`$ and $`\mathrm{\Pi }`$ are six by six matrices:
$`\mathrm{\Sigma }_0=\left(\begin{array}{cccc}& \mathrm{𝟏}_\mathrm{𝟐}& & \\ \mathrm{𝟏}_\mathrm{𝟐}& & & \\ & & & 1\\ & & 1& \end{array}\right),\mathrm{\Pi }={\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{cc}\mathrm{𝟎}_\mathrm{𝟒}& \begin{array}{cc}h_1& h_2\\ h_2^{}& h_1^{}\end{array}\\ \begin{array}{cc}h_1^{}& h_2^T\\ h_2^{}& h_1^T\end{array}& \mathrm{𝟎}_\mathrm{𝟐}\end{array}\right),`$ (26)
where the two Higgs doublets $`h_1`$, $`h_2`$ have the $`SU(2)\times U(1)`$ SM quantum numbers $`(\mathrm{𝟐},+1/2)`$ and $`(\mathrm{𝟐},1/2)`$, respectively. We have omitted the eaten fields from the $`\mathrm{\Pi }`$ matrix, as well as omitted singlet pseudo-Goldstone bosons that do not affect EWPTs. We have adopted the basis used in Ref. Gregoire:2003kr here, but changed the definition of $`h_2`$ for convenience. The $`[SU(2)\times U(1)]^2`$ generators are also given in Ref. Gregoire:2003kr .
The gauge boson mixings are described by Eqs. (9) and (10) and their masses are
$$M_W^{}=\frac{gF}{2sc},M_Z^{}=\frac{g^{}F}{\sqrt{8}s^{}c^{}}.$$
(27)
Integrating out the heavy gauge bosons, we obtain
$`a_h`$ $`=`$ $`{\displaystyle \frac{1}{F^2}}[(c^2s^2)^2+{\displaystyle \frac{1}{2}}\mathrm{cos}^2(2\beta )],`$
$`a_{hq}^t`$ $`=`$ $`a_{hl}^t={\displaystyle \frac{1}{2F^2}}(c^2s^2)c^2,`$
$`a_{hf}^s`$ $`=`$ $`{\displaystyle \frac{2s^{}c^{}(c^2s^2)}{F^2}}\left(Y_2^f{\displaystyle \frac{s^{}}{c^{}}}Y_1^f{\displaystyle \frac{c^{}}{s^{}}}\right),`$
$`a_{lq}^t`$ $`=`$ $`a_{ll}^t={\displaystyle \frac{c^4}{F^2}},`$
$`a_{ff^{}}^s`$ $`=`$ $`{\displaystyle \frac{8s^2c^2}{F^2}}\left(Y_2^f{\displaystyle \frac{s^{}}{c^{}}}Y_1^f{\displaystyle \frac{c^{}}{s^{}}}\right)\left(Y_2^f^{}{\displaystyle \frac{s^{}}{c^{}}}Y_1^f^{}{\displaystyle \frac{c^{}}{s^{}}}\right).`$ (28)
Since only the vevs of $`h_1`$ and $`h_2`$ matter, we have combined their contributions to a single $`h`$. The term proportional to $`\mathrm{cos}^2(2\beta )`$ in $`a_h`$ comes from expanding the kinetic term in Eq. (7).
The Yukawa couplings can be constructed in several different ways. Here we choose the Yukawas in a way similar to the $`SU(5)`$ model to avoid extra corrections to EWPOs. For example, we add a pair of fermions ($`t_L^{},t_R^{}`$) and define $`\chi =(u_{3L},d_{3L},0,0,0,t_L^{})^T`$. Then the top Yukawa coupling comes from the Lagrangian
$$\lambda _1fϵ_{ijk}ϵ_{xy}\overline{\chi _i}\mathrm{\Sigma }_{jx}\mathrm{\Sigma }_{ky}u_{3R}+\lambda _2f\overline{t^{}}_Lt_R^{},$$
(29)
where $`i,j,k\{1,2,6\}`$ and $`x,y\{3,4\}`$. This Lagrangian only induces mixing between the top quark and the heavy fermion and thus does not affect EWPTs.
Comparing Eq. (15) and Eq. (28), we see that the two sets of coefficients have similar structure. Thus we expect that the constraints from EWPTs in the $`SU(6)`$ model to have similar features as in the $`SU(5)`$ model. First, if the fermions are charged under only one $`U(1)`$, the contributions from the $`Z^{}`$ exchange are large and one expects tight constraints. In this case, we have verified that the bounds on $`M_t^{}`$ are greater than $`6`$ TeV for all choices of $`c`$, $`c^{}`$ and $`\mathrm{tan}\beta `$. Second, the two limits that lead to suppression of coefficients are present here as well. Instead of the triplet contribution in the $`SU(5)`$ model, there is another custodial symmetry breaking term generated by the nonlinear structure of the $`\mathrm{\Sigma }`$ field in this model. For comparison with the $`SU(5)`$ model, when $`s^{}=c^{}`$ and $`Y_1^f=Y_2^f`$ we draw 95% CL bounds on $`M_t^{}`$ and $`M_W^{}`$ as functions of $`c`$ and $`\mathrm{tan}\beta `$ in Fig. 3. The “near-oblique” limit discussed in Ref. Gregoire:2003kr corresponds to the region where $`\mathrm{tan}\beta 1`$ and $`c0`$ in Fig. 3. In this limit, the new particles could be so light that loop corrections have to be considered to obtain accurate bounds, as has been done in Ref. Gregoire:2003kr . Like in the $`SU(5)`$ model, the other limit: $`Y_1^f=Y^f`$, $`Y_2^f=0`$ and $`c,c^{}1`$ suppresses all coefficients except $`a_h`$.
The two $`U(1)`$ generators can also be modified. Since the global symmetry is larger than in the $`SU(5)`$ model, we have more choices. For example, by redefining
$$Y_1Y_1+bI+b^{}\text{Diag}\{\mathrm{𝟏}_2,\mathrm{𝟏}_2,0,0\},Y_2Y_1bIb^{}\text{Diag}\{\mathrm{𝟏}_2,\mathrm{𝟏}_2,0,0\},$$
(30)
we can rescale the $`Z^{}`$ mass and the corresponding $`a_i`$. We can also change the $`U(1)`$ generators in a way that the Higgs bosons are charged differently under the two $`U(1)`$’s. For example, taking
$$Y_1=\left(\begin{array}{ccc}0_4& & \\ & \frac{5}{8}& \\ & & 0\end{array}\right),Y_2=\left(\begin{array}{ccc}0_4& & \\ & \frac{1}{8}& \\ & & \frac{1}{2}\end{array}\right)$$
(31)
yields the correct $`U(1)_Y`$ charges for $`h_1`$ and $`h_2`$, while the coupling between $`h_2`$ and $`Z^{}`$ becomes
$$\frac{1}{2}igZ^\mu (D_\mu h_2)^{}(\frac{5c^{}}{8s^{}}\frac{3s^{}}{8c^{}})h_2+h.c.$$
(32)
There are similar terms for the $`h_1`$ to $`Z^{}`$ coupling. Correspondingly, $`a_h`$ and $`a_{hf}^s`$ are no longer proportional to $`(c^2s^2)`$. Such changes affect the near-oblique condition, but they do not introduce essentially new features in our analysis. It would be interesting if such structure comes naturally from some underlying theories.
## V The $`SU(3)\times U(1)`$ models Kaplan:2003uc ; Skiba:2003yf ; simplest
As effective theories, the little Higgs models do not need to be anomaly free. The original $`SU(3)`$ model Kaplan:2003uc is anomalous, indicating that additional fermions must be present at the cutoff. Anomaly-free versions of the model have also been constructed simplest ; Kong:2004cv . However, they require assigning different charges for different generations. Our methods can not be applied to this generation-dependent model. Therefore, we concentrate on the original model only.
By integrating out the heavy fermions, we obtain the following coefficients. The notation follows Ref. simplest . For one generation,
$`a_{hl}^s`$ $`=`$ $`a_{hl}^t={\displaystyle \frac{1}{4}}{\displaystyle \frac{f_2^2}{F^2f_1^2}},`$ (33)
$`a_{hq}^s`$ $`=`$ $`a_{hq}^t={\displaystyle \frac{1}{4}}{\displaystyle \frac{(\lambda _1^{u2}\lambda _2^{u2})^2f_1^2f_2^2}{[(f_1\lambda _1^u)^2+(f_2\lambda _2^u)^2]^2F^2}},`$ (34)
where $`F^2=f_1^2+f_2^2`$. For three generations, $`\lambda _1^u`$ and $`\lambda _2^u`$ are flavor dependent and in general should be $`3\times 3`$ matrices. However, in order to avoid large FCNCs, one should set one of the matrices to be proportional to the identity matrix and the other one to be hierarchical Kaplan:2003uc . For example, if $`\lambda _2^u`$ is proportional to the identity matrix and $`\lambda _1^u`$ is hierarchical, for the first two generations, $`\lambda _1^u\lambda _2^u`$, so we can approximate $`\lambda _1^u0`$. For the third generation, since only the top quark mixes with the heavy quark the EWPTs are not affected. The relation $`a_{hq}^s=a_{hq}^t`$ is a consequence of the fact that the bottom quark does not mix with the heavy fermion. We can assign any values to $`a_{hq}^s`$ and $`a_{hq}^t`$ for the third generation as long as the relation holds. Therefore we can ignore all terms containing $`\lambda _1^u`$ in Eq. (34). Similarly, if $`\lambda _2^u`$ is hierachical instead, we can ignore all terms containing $`\lambda _2^u`$. In this case the coefficients $`a_{hq}^s`$ and $`a_{hq}^t`$ are identical to the coefficients of the corresponding lepton operators in Eq. (33) and the constraints on the model turn out to be even less stringent.<sup>5</sup><sup>5</sup>5We thank Heather Logan for drawing our attention to this possibility. Thus, we obtain the following flavor-independent operators
$`a_{hq}^s=a_{hq}^t={\displaystyle \frac{1}{4}}{\displaystyle \frac{f_1^2}{F^2f_2^2}}\text{ (}\lambda _1^u\text{ hierarchical)},`$
$`a_{hq}^s=a_{hq}^t={\displaystyle \frac{1}{4}}{\displaystyle \frac{f_2^2}{F^2f_1^2}}\text{ (}\lambda _2^u\text{ hierarchical)}.`$ (35)
By integrating out the heavy gauge bosons, we obtain
$`a_h`$ $`=`$ $`{\displaystyle \frac{9}{4F^2}}{\displaystyle \frac{(1\frac{2}{3}x^2)^2}{(3+x^2)^2}},`$
$`a_{hf}^s`$ $`=`$ $`{\displaystyle \frac{9}{4F^2}}{\displaystyle \frac{1\frac{2}{3}x^2}{(3+x^2)^2}}(\sqrt{3}T^{8f}+x^2Y_x^f),`$
$`a_{ff^{}}^s`$ $`=`$ $`{\displaystyle \frac{9}{2F^2}}{\displaystyle \frac{1}{(3+x^2)^2}}(\sqrt{3}T^{8f}+x^2Y_x^f)(\sqrt{3}T^{8f^{}}+x^2Y_x^f^{}),`$ (36)
where $`x=g_x/g`$, while $`T^{8f}=1/(2\sqrt{3}),1/(2\sqrt{3}),0,0,0`$ and $`Y_x^f=1/3,1/3,2/3,1/3,1`$ are fermion charges for $`f=q,l,u,d,e`$, respectively. Moreover, $`g_x`$ is related to $`g`$ and $`g^{}`$ by
$$g_x^2=\frac{3g^2g^2}{3g^2g^2}.$$
(37)
The complete list of coefficients is given by Eqs. (33), (35) and (36).
In this model, the mass of the heavy fermion and the top Yukawa coupling are given by simplest
$$M_t^{}^2=\lambda _1^2f_1^2+\lambda _2^2f_2^2,\lambda _t=\lambda _1\lambda _2\frac{F}{M_t^{}}.$$
(38)
Unlike the $`SU(5)`$ and $`SU(6)`$ models, for given $`f_1`$ and $`f_2`$, $`M_t^{}`$ is not uniquely determined because of the freedom in choosing the ratios $`\lambda _1/\lambda _2`$ and $`f_1/f_2`$. Thus, $`M_t^{}`$ is not tightly constrained. However, the mass of the heavy gauge doublet $`W^{}`$ is determined uniquely in terms of $`F`$: $`M_W^{}^2=g^2F^2/2`$.
Given the ratio $`f_1/f_2`$, we can obtain the bound on $`M_W^{}`$. Fig. 4 shows 95% CL bounds on $`M_W^{}`$ as a function of $`f_1/f_2`$. Since the contribution to the Higgs mass from $`W^{}`$ is not significant until $`M_W^{}>6`$ TeV the fine-tuning problem is not severe. Depending on the choice of the Yukawa matrices in the quark sector corresponding to the two possibilities in Eq. (35), the least severe constraints are either when $`f_1/f_21`$ for hierarchical $`\lambda _1^u`$, or when $`f_1/f_2>2`$ for hierarchical $`\lambda _2^u`$.
The $`SU(3)\times U(1)`$ gauge sector can be embedded in different global groups. The above discussion makes use of the setup of Ref. simplest , where the coset space is $`[SU(3)/SU(2)]^2`$. But the constraints also apply to the model constructed in Ref. Skiba:2003yf , where the coset space is $`SU(9)/SU(8)`$. The model is also extended to an $`[SU(4)/SU(3)]^4`$ model with an $`SU(4)`$ group gauged in Ref. Kaplan:2003uc . These variations mainly stem from theoretical considerations, rather than the need to avoid the experimental constraints.
## VI Summary and discussion
Using effective field theory approach, we have obtained constraints from EWPTs for three types of little Higgs models: $`SU(5)/SO(5)`$, $`SU(6)/SP(6)`$ and $`SU(3)\times U(1)`$. We have carried out the analysis by first integrating out the heavy fields to obtain the effective operators, and then calculating the constraints from the bounds on the linear combinations of the coefficients of these operators.
We gained two main benefits from following through this procedure. First, expressing the corrections in terms of effective operators allowed us to treat the corrections in a model-independent way, which we have done in Ref. Han:2004az . Thus we can simply use the results in Ref. Han:2004az once we acquire the coefficients of the operators. For the $`SU(5)`$ and $`SU(6)`$ models, we use the heavy top mass $`M_t^{}`$ and the heavy triplet gauge boson mass $`M_W^{}`$ to illustrate the bounds. We show that there exist regions in the parameter space where the bounds on $`M_t^{}`$ and $`M_W^{}`$ are less than 2 TeV and 6 TeV respectively. In these regions, no significant fine-tuning is required to yield light Higgs bosons. These regions are in the vicinity of the points where the two $`U(1)`$ gauge sectors in the models have the same coupling strength and fermions are charged identically under the two $`U(1)`$’s. Away from this limit, the bounds are usually much tighter. In particular, for the cases that fermions are charged under only one $`U(1)`$, we obtain the bounds $`M_t^{}>6`$ TeV at 95% CL for all of the parameter space, requiring fine-tuning of more than 1%. For the $`SU(3)`$ model, $`M_t^{}`$ remains a free parameter for a given $`F`$. Thus we only obtain the bounds for $`M_W^{}`$. At 95% CL, $`M_W^{}>1.8`$ TeV. The associated fine-tuning is not significant in this case.
It is interesting that LEP2 results are very useful. The models induce many 4-fermion operators which LEP2 observables are sensitive to. The LEP2 results were not included in most of previous papers considering EWPT constraints on littlest Higgs models. As a comparison, we have checked our results for the $`SU(5)`$ and $`SU(6)`$ models against Ref. Csaki:2002qg ; Csaki:2003si . We agree with the results in the two references when we use the same set of observables, which contain most precisely measured LEP1 and some low-energy observables. After we include the LEP2 observables, the bounds become significantly tighter.
In Ref. Marandella:2005wd the authors also obtain constraints on little Higgs models from EWPTs including LEP2 data. Technically, our approach is more general. The authors of Ref. Marandella:2005wd have to assume that the models are approximately universal and all corrections can be condensed in four oblique parameters. In particular, the heavy gauge bosons should couple to fermions “universally”, which means that the fermion currents that couple to the heavy gauge bosons are proportional to the SM current. This is not true for the $`SU(3)`$ model, so the authors of Ref. Marandella:2005wd have to neglect some of the EWPOs. In addition, they do not discuss the effects of the heavy fermions in the $`SU(3)`$ model, which are not universal either. In the $`SU(5)`$ and $`SU(6)`$ models we have analyzed, the fermion charge assignments do yield universal couplings between fermions and gauge bosons. Thus their method apply as well. However, one could obtain non-universal couplings from more general fermion assignments, which would make their method inapplicable, but would not introduce any difficulty in our approach. In Marandella:2005wd only the simplest fermion assignments made in the original papers were considered, while we are interested in the variations which relax the electroweak constraints.
The second advantage of the effective theory approach is that it makes it transparent how to modify the models to avoid tight constraints. From the compact but complete lists of operator coefficients in the $`SU(5)`$ and $`SU(6)`$ models, we easily deduce two regions of parameter space that lead to suppressions of the corrections. One of the regions for the $`SU(6)`$ model case is discussed and termed “near-oblique” limit in Ref. Gregoire:2003kr . In Ref. Chang:2003zn , the other limit has been utilized in a model with enlarged gauge sector and an approximate custodial symmetry. With their weaker constraints, such model variations are interesting since they address the fine-tuning problem better than the original models. From experimental point of view, new particles are light enough to be observable at the LHC only if there is no significant fine tuning.
## Acknowledgments
This research was supported in part by the US Department of Energy under grant DE-FG02-92ER-40704 and by Outstanding Junior Investigator grant. WS thanks the Aspen Center for Physics for its hospitality. |
warning/0506/hep-ph0506117.html | ar5iv | text | # 1 Quark line diagrams for 𝐷-decays. (a) is the typical spectator model that underlies the idea that the resulting 𝐾𝜋 system has just 𝐼=1/2. (b) involves internal 𝑊-boson conversion. For semi-leptonic decays only the analogue of (a) is possible. This is shown as (c).
DCPT/05/42
IPPP/05/21
Estimating the $`I=3/2`$ $`K\pi `$ interaction
in $`D`$ decay
L. Edera and M.R. Pennington
Institute for Particle Physics Phenomenology,
University of Durham, Durham, DH1 3LE, U.K.
Abstract
Heavy flavour decay to light hadrons is the key to understanding many aspects of the Standard Model from CP violation to strong dynamics. It is often presumed in line with the simple quark spectator model of $`D`$ decay to $`K\pi \pi `$ that the $`K\pi `$ system has only $`I=1/2`$. E791 have recently presented an analysis of their results on $`D^+\left(K^{}\pi ^+\right)\pi ^+`$ using a generalised isobar picture of two body interactions. While higher $`K\pi `$ waves are described by sums of known resonances, the $`S`$-wave amplitude and phase are determined bin-by-bin in $`K\pi `$ mass. The phase variation is found not to be that of $`K^{}\pi ^+`$ elastic scattering. This hints at a different mixture of $`I=1/2`$ and $`I=3/2`$ $`S`$-wave interactions than in elastic scattering. Applying Watson’s theorem to this generalised isobar model allows us to estimate the $`I=3/2`$ $`K\pi `$ $`S`$-wave component. We indeed find that this is larger than in hadronic scattering or semileptonic processes.
Drawing quark line diagrams for heavy flavour decays provides a guide to the weak interaction dynamics that takes place. Thus as shown in Fig. 1a, in $`D`$ decay to $`K\pi \pi `$, the $`c`$ quark is seen to change into an $`s`$ quark by emitting an off-shell $`W^+`$ that materialises as a $`\pi ^+`$. The other quark in the $`D`$, for instance a $`d`$, acts merely as a spectator and with the $`s`$ quark forms on the creation of a $`u\overline{u}`$ or $`d\overline{d}`$ a $`K\pi `$ system. Since the isospin of the $`c`$ and $`s`$ quarks is zero, it is then natural to assume that the $`K\pi `$ system has only the isospin of the $`u`$ or $`d`$ quark, namely $`I=1/2`$. This presumption is in keeping with a simple isobar picture of this decay in which well-known $`I=1/2`$ $`K^{}`$ resonances, like the $`K^{}(890)`$ or $`K_2^{}(1430)`$, are readily seen. However, there are, of course, $`\pi \pi `$ resonances too, like the $`\rho `$, in other charge channels, that clearly indicate that the spectator picture is too simplistic. The $`W`$ can be internally converted into hadrons as shown in Fig. 1b, and then the $`K\pi `$ system can have $`I=3/2`$ as well as $`I=1/2`$. From wholly hadronic reactions we know the $`I=3/2`$ interactions do not produce resonances, in keeping with the $`q\overline{q}`$ model of hadrons. Moreover, in low energy hadron scattering unitarity is a key constraint and restricts the magnitude of every partial wave amplitude to be bounded by unity. However, in the weak decay reaction we have no such constraint. Whilst in $`K`$ decays we know there is the well-known $`\mathrm{\Delta }I=1/2`$ rule, no such property is known for charmed decays. Consequently, we cannot presume that the $`I=1/2`$ component always dominate over $`I=3/2`$ $`K\pi `$ interactions, even though only the former is dominated by resonances. Here we estimate how big the fraction of $`I=3/2`$ amplitude is in $`D`$ decay.
This is made possible by a recent model-independent treatment of the E791 results on $`DK\pi \pi `$ decay presented at La Thuile by Meadows . There a Dalitz plot analysis is performed that isolates the phase and magnitude of the $`K\pi `$ $`S`$-wave interaction bin by bin in $`K\pi `$ mass. To explain how this allows us to estimate the $`I=1/2`$ and $`3/2`$ components, we begin with $`K\pi K\pi `$ scattering. For this the $`S`$-wave amplitude, in particular, is given by
$$𝒯(K\pi K\pi ;s)=C_{1/2}𝒯^{1/2}(s)+C_{3/2}𝒯^{3/2}(s),$$
(1)
where $`s`$ is the square of the $`K\pi `$ mass, and the $`C_I`$ are the appropriate Clebsch-Gordan coefficients depending on the charges of the $`K\pi `$ system. The hadronic amplitude $`𝒯^I`$ with definite isospin has a magnitude and phase related by elastic unitarity for $`\sqrt{s}<(m_K+m_\eta ^{})`$, the first strongly coupled threshold, and given by:
$$𝒯^I(s)=\frac{1}{\rho }\mathrm{sin}\delta ^I\mathrm{exp}(i\delta ^I),$$
(2)
where $`\rho =\mathrm{\hspace{0.17em}2}k/\sqrt{s}`$ with $`k`$ the $`K\pi `$ c.m. 3-momentum. For the famous LASS experiment the $`K^{}\pi ^+K^{}\pi ^+`$ amplitude is thus given by
$$𝒯(K^{}\pi ^+K^{}\pi ^+;s)=\frac{2}{3\rho }\left[\mathrm{sin}\delta ^{1/2}\mathrm{exp}(i\delta ^{1/2})+\frac{1}{2}\mathrm{sin}\delta ^{3/2}\mathrm{exp}(i\delta ^{3/2})\right].$$
(3)
These data, plus results on $`K^{}\pi ^{}K^{}\pi ^{}`$ , fix $`\delta ^I`$ for both $`I=1/2,3/2`$. The LASS group provide parametrisations of these data, which are displayed as the dashed lines in the lower part of Fig. 2.
Though these are valid in the experimental region LASS accessed, i.e. above 825 MeV, their extrapolations down to $`K\pi `$ threshold are known to be inconsistent with Chiral Perturbation Theory. Consequently, we consider a parametrisation of the elastic phases, shown as the solid lines in Fig. 2, that is consistent with the calculations of Chiral Perturbation Theory as shown in the dispersive analysis by Büttiker et al These phases are the key ingredients in our discussion.
We now turn to $`D`$ decay, firstly in its semi-leptonic mode to $`K\pi \mu \nu `$. The $`S`$-wave amplitude in the $`K\pi `$ channel is given by
$$(D(K\pi )\mu \nu ;s)=_{sl}^{1/2}(s)+_{sl}^{3/2}(s),$$
(4)
where the $`_{sl}^I(s)`$ are the semileptonic production amplitudes with $`I=1/2,3/2`$. In this case the emitted $`W^+`$ is the source of the dilepton system, which of course, has no strong interaction with the $`K\pi `$ pair, as illustrated in Fig. 1c. Now elastic unitarity requires that the phase of the $`K\pi `$ interaction with definite spin and isospin is the same as that in elastic $`K\pi `$ scattering, so that:
$$_{sl}^I(s)=_{sl}^I(s)\mathrm{exp}\left[i\delta ^I(s)\right].$$
(5)
This is the famous final state interaction theorem of Watson . Results from FOCUS confirm that this relationship holds and that the $`I=3/2`$ component is small or negligible in this process.
With this confirmation of Watson’s theorem in heavy flavour decays, we consider the wholly hadronic channel: $`D^+`$ decay to $`K^{}\pi ^+\pi ^+`$. This has been analysed by E791 in a “new approach to the analysis of 3 body decays” . The structure of the Dalitz plot has marked bands in both $`K^{}\pi ^+`$ mass combinations for the well known $`K^{}(892)`$ and $`K_2^{}(1430)`$ that feature so prominently in $`K^{}\pi ^+K^{}\pi ^+`$ scattering . This suggests the dominance of 2-body strong interactions. Consequently, E791 assume a generalised isobar picture in which genuine 3-body interactions are neglected. As is by now standard , the decay matrix element is represented in the way summarised, for instance, by the CLEO-c group in Section II of Ref. and so for the $`D`$-decay we consider here is then a sum of the amplitudes for the two possible combinations of $`K^{}\pi ^+`$ interactions with their appropriate spin and vertex factors. As in the usual isobar picture, all waves with $`J>0`$ are represented by a sum of resonance terms. This incorporates the known $`I=1/2`$ $`K\pi `$ resonances like the $`K^{}(892)`$ $`K_1^{}(1410)`$ and $`K_1^{}(1680)`$ for $`J=1`$, for instance. With these resonances being well separated, their amplitude can be reasonably described by a sum of Breit-Wigner forms. Each partial wave has a production phase from the $`D\pi K^{}`$ coupling, which is assumed to be a constant as is also standard . Since the parameters of each resonance are those found in $`K\pi `$ scattering, e.g. by LASS , the phase variation of each $`K\pi `$ partial wave in $`D`$-decay is exactly that of $`K\pi `$ scattering in the elastic region, in agreement with the application of Watson’s theorem to such an isobar picture . Since one overall phase is not determinable, the strong $`P`$-wave is taken as the reference wave with its production phase set to zero.
What is new in the E791 analysis is the description of the $`K\pi `$ $`S`$-wave. There the well-established wide $`K_0^{}(1430)`$ appears, with perhaps an even broader $`\kappa `$ at lower mass. These are not simply describable by a sum of isolated Breit-Wigners, as the $`K^{}\pi ^+K^{}\pi ^+`$ amplitude of Fig. 2 illustrates. Rather than enforce some ad hoc prescription of this key wave, E791 represent this by a magnitude and phase in each bin of $`K\pi `$ mass in each $`K^{}\pi ^+`$ combination. The beauty of this analysis is that each $`S`$-wave $`K\pi `$ mass band overlaps with a crossed $`K\pi `$ band in a $`P`$-wave and so the relative phase is determined, as well as the magnitude of the $`S`$-wave. This provides as close to a model independent determination of the $`K^{}\pi ^+`$ $`S`$-wave interaction in this $`D`$-decay as is presently possible. The results are shown in Fig. 3. From the application of Watson’s theorem to this generalised isobar picture , one would expect the phase variation of this $`S`$-wave amplitude to follow that of $`K\pi `$ scattering in the region of elastic unitarity. As already mentioned (just before Eq. (2)), elastic unitarity is found to be very nearly exact upto $`K\eta ^{}`$ threshold, despite the fact that this is above the opening of the $`K\pi \pi `$ and $`K\eta `$ channels, from which we infer these to be negligible below 1450 MeV. While the phase of the $`K\pi `$ $`S`$-wave in $`D`$-decay and that of $`K\pi `$ scattering (compare the lower graphs of Figs. 2 and 3) both have an upward trend in the “elastic” region, they do not match to the precision expected. The reasons for this can be manifold:
* the isobar assumption of only 2-body $`K\pi `$ interactions may not be true,
* even if it is, then perhaps Watson’s theorem does not apply and so the phase variation is not required to be the same,
* even if Watson’s theorem applies, should the phase variation of the $`K^{}\pi ^+`$ interactions be that of $`K^{}\pi ^+K^{}\pi ^+`$ scattering or that of just its $`I=1/2`$ component?
These are all questions that have been raised at the BaBar Dalitz Workshop in December 2004 .
Since the generalised isobar picture of pairwise $`AB`$ interactions defines a model for $`ABAB`$ scattering, we adopt the view that this must be in keeping with data on $`ABAB`$ scattering and the phase variation must agree in the region of elastic unitarity . This has to be so, for the model to be consistent. Thus the $`K\pi `$ $`S`$-wave phase variation found by E791 must agree with that for $`K\pi K\pi `$ scattering. However, unitarity is only diagonalised by partial wave amplitudes with definite quantum numbers such as isospin, and not in general for individual charged channels, like $`K^{}\pi ^+`$. Consequently, it is the phase variation with definite isospin that should match. There is no reason that the combination of $`I=1/2`$ and $`I=3/2`$ contributions determined by simple Clebsch-Gordan coefficients in $`K^{}\pi ^+`$ scattering (as in Eq. (3)) is that formed in $`D`$-decay, to which we already alluded in the introduction. The relationship between the phase variations then allows us to estimate the relative $`I=1/2`$ and $`3/2`$ components of the E791 $`S`$-wave, as we shall shortly describe. Of course, there could also be $`I=3/2`$ components in the higher partial waves too. These have been neglected in the E791 resonance-dominated description of these amplitudes. Nevertheless, it is known that the $`J1`$ waves with $`I=3/2`$ vary by less than $`3^o`$ between $`K\pi `$ threshold and 1.8 GeV. Consequently, any such components should have only a tiny effect on our estimate.
In keeping with the generalised isobar description that E791 adopt, the effect of the spectator pion is to produce an additional production phase, $`\beta _I`$, which we take to be a constant, as they do for all other waves. The $`K^{}\pi ^+`$ $`S`$-wave determined by E791 we call $``$, which is a function of $`s`$, the square of the $`K\pi `$ invariant mass, $`E`$. This is then given by
$$(s)=_{had}^{1/2}(s)+_{had}^{3/2}(s)𝒜(s)\mathrm{exp}[i\varphi (s)],$$
(6)
where the amplitudes with definite quantum numbers are given by:
$$^I(s)_{had}=_{had}^I(s)\mathrm{exp}\left[i\delta ^I(s)+i\beta _I\right],$$
(7)
in the region of elastic unitarity below 1450 MeV. The phases $`\beta _I`$ reflect the structure of the complete set of quark line graphs of Figs. 1a, b in the strong coupling limit before final state interactions are included. Since different graphs contribute to each $`K\pi `$ isospin, the production phase $`\beta _I`$ depends on $`I`$. It is the magnitude $`𝒜`$ and phase, $`\varphi `$, of the total $`S`$-wave $`K\pi `$ amplitude, displayed in Fig. 3, that E791 have determined . Given this and the $`K\pi `$ phases, $`\delta ^I`$ shown in Fig. 2, the aim is to deduce the magnitude of the $`I=1/2`$ and $`3/2`$ $`K\pi `$ amplitudes in $`D`$ decay.
Let us see how to do this. In vector terms it is like finding the vectors $`^{\mathrm{𝟏}/\mathrm{𝟐}}`$ and $`^{\mathrm{𝟑}/\mathrm{𝟐}}`$ knowing only their sum $`𝒜`$, i.e.
$$𝒜=^{\mathrm{𝟏}/\mathrm{𝟐}}+^{\mathrm{𝟑}/\mathrm{𝟐}}.$$
(8)
Clearly there are an infinite number of vectors $`^𝐈`$ that satisfies this. If the production phases, $`\beta _I`$, were zero (or otherwise determined), the solution would be simple, since we know from $`K\pi `$ elastic scattering what the phases $`\delta _I`$ are at each energy in the elastic region, Fig. 2. Knowing the directions of the vectors, their magnitudes can easily be found. However, here we do not know in advance what the production phases are, but we do know that these are the same phases at every $`K\pi `$ mass, $`E=\sqrt{s}`$. This sets the scene for determining the amplitudes.
Starting from
$$𝒜\mathrm{exp}(i\varphi )=^{1/2}+^{3/2},$$
(9)
where the $`^I`$ are the complex $`I=1/2,\mathrm{\hspace{0.17em}3}/2`$ amplitudes, it is straightforward to check that the solution is :
$`^{1/2}(E)`$ $`=`$ $`𝒜{\displaystyle \frac{\mathrm{sin}\left(\delta ^{3/2}(E)+\beta _{3/2}\varphi (E)\right)}{\mathrm{sin}\left(\delta ^{3/2}(E)\delta ^{1/2}(E)\beta _{1/2}+\beta _{3/2}\right)}}\mathrm{exp}\left[i(\delta ^{1/2}(E)+\beta _{1/2})\right],`$ (10)
$`^{3/2}(E)`$ $`=`$ $`𝒜{\displaystyle \frac{\mathrm{sin}\left(\delta ^{1/2}(E)+\beta _{1/2}\varphi (E)\right)}{\mathrm{sin}\left(\delta ^{1/2}(E)\delta ^{3/2}(E)+\beta _{1/2}\beta _{3/2}\right)}}\mathrm{exp}\left[i(\delta ^{3/2}(E)+\beta _{3/2})\right].`$ (11)
The phase difference, $`\delta ^{1/2}(E)\delta ^{3/2}(E)`$, varies from 0 to $`180^o`$ in the elastic region below $`K\eta ^{}`$ threshold, and so we see that the denominator in each expression, namely
$$\mathrm{sin}\left(\delta ^{1/2}(E)\delta ^{3/2}(E)+\beta _{1/2}\beta _{3/2}\right)$$
(12)
will inevitably vanish at at least one energy in this region. Let this reference energy be $`E_r`$. Since the amplitudes $`^I`$ are finite at real energies, the numerator must vanish at this same reference energy. This fixes the production angles. Thus
$`\beta _{1/2}`$ $`=`$ $`\varphi (E_r)\delta ^{1/2}(E_r)+m\pi `$ (13)
$`\beta _{3/2}`$ $`=`$ $`\varphi (E_r)\delta ^{3/2}(E_r)+n\pi ,`$ (14)
where $`m,n`$ are integers (including zero). Consequently, we can write the individual isospin amplitudes in terms of just one unknown parameter, the reference energy $`E_r`$, as:
$`^{1/2}(E)`$ $`=`$ $`𝒜{\displaystyle \frac{\mathrm{sin}\left(\varphi (E)\varphi (E_r)\delta ^{3/2}(E)+\delta ^{3/2}(E_r)\right)}{\mathrm{sin}\left(\delta ^{1/2}(E)\delta ^{3/2}(E)\delta ^{1/2}(E_r)+\delta ^{3/2}(E_r)\right)}}`$ (15)
$`\times \mathrm{exp}\left[i(\delta ^{1/2}(E)\delta ^{1/2}(E_r)+\varphi (E_r))\right]`$
$`^{3/2}(E)`$ $`=`$ $`𝒜{\displaystyle \frac{\mathrm{sin}\left(\delta ^{1/2}(E)\delta ^{1/2}(E_r)\varphi (E)+\varphi (E_r)\right)}{\mathrm{sin}\left(\delta ^{1/2}(E)\delta ^{3/2}(E)\delta ^{1/2}(E_r)+\delta ^{3/2}(E_r)\right)}}`$ (16)
$`\times \mathrm{exp}\left[i(\delta ^{3/2}(E)\delta ^{3/2}(E_r)+\varphi (E_r))\right].`$
For each value of $`E_r`$ we can then determine the amplitudes. Consequently, we have a continuous range of possible amplitudes. Within this set, there are only a small group that are physically (as opposed to mathematically) allowed. Watson’s theorem is ensured by the phase variation of each of $`^{1/2}`$ and $`^{3/2}`$ being given by the elastic phases of the corresponding elastic scattering amplitudes $`𝒯^I`$. However, the dynamics of the final state interactions requires that the amplitudes $`^I`$ should be smoothly connected to the scattering amplitudes $`𝒯^I`$ through coupling functions $`\overline{\alpha _I}(E)`$. A consequence of this is that any resonance poles in $`K\pi `$ scattering appear in the decay process with the same mass and width. The functions $`\overline{\alpha _I}`$ fix the coupling to the decay channel. Thus
$$^I(E)=\overline{\alpha _I}(E)𝒯^I(E)\mathrm{exp}(i\beta _I).$$
(17)
The dynamics of the pseudoscalar interactions imposes Adler zeros in each of the $`K\pi `$ scattering amplitudes with $`I=\mathrm{\hspace{0.17em}1}/2,\mathrm{\hspace{0.17em}3}/2`$ at $`s=s_I^{\mathrm{\hspace{0.17em}0}}`$. Such zeros may appear in the decay process but not necessarily at the same position as in elastic scattering. We take account of this by defining reduced elastic scattering amplitudes with the Adler zero divided out. We then specify new coupling functions, $`\alpha _I`$ which have to be smooth since they contain no explicit $`s`$channel dynamics. Thus we have
$$^I(E)=\alpha _I(E)\frac{𝒯^I(E)}{ss_I^{\mathrm{\hspace{0.17em}0}}}\mathrm{exp}(i\beta _I),$$
(18)
where the coupling functions $`\alpha _I`$ are representable by low order polynomials in $`s=E^2`$.
We therefore determine the coupling functions $`\alpha _I(E)`$ for each choice of $`E_r`$ and select those that are representable by low order polynomials, while still providing an accurate fit to the amplitude determined by the E791 analysis. With 26 values of magnitude and 26 phases in the elastic region below $`K\eta ^{}`$ threshold, this is a severe constraint.
The “best” fit is found with $`E_r\mathrm{\hspace{0.17em}800}`$ MeV, for which the production phases are $`\beta _{1/2}=\mathrm{\hspace{0.17em}72}^o`$ and $`\beta _{3/2}=73^o`$. The moduli of the $`I=1/2`$ and $`3/2`$ amplitudes, using Eqs. (15,16), are shown in Fig. 4. As can be seen these are accurately determined away from $`E_r`$, where the errors are very small. These show that the $`\alpha _I`$ of Eq. (18) are well represented by quadratic polynomials in the square of the $`K\pi `$ mass, $`s`$. The precise results above 1 GeV fix the continuation to lower $`K\pi `$ mass as illustrated in Fig. 3. The requirement that the $`D`$-decay amplitude of Eq. (6) should be well described is illustrated by the (red) lines, labelled 2 in Fig. 3. Allowing a one standard deviation change by varying $`E_r`$ alters the $`I=1/2`$ and $`3/2`$ amplitudes, while still fitting the E791 magnitudes and phase. This result is shown in Figs. 3, 5. The lines 1, 3 (in blue and green, respectively) in Fig. 3 show the corresponding small difference in the description of the E791 amplitude, while Fig. 5 shows the range of variation of the $`I=1/2`$ and $`3/2`$ then permitted. The bands reflect a change of $`E_r`$ from 750 to 950 MeV, or equivalently a range of
$$\beta _{1/2}=\mathrm{\hspace{0.17em}65}^o\mathrm{\hspace{0.17em}85}^o,\beta _{3/2}=86^o38^o,$$
(19)
which are strongly correlated.
We, of course, only know that the coupling functions $`\alpha _I`$ of Eq. (18) should be smooth functions of $`K\pi `$ mass since they do not contain direct $`K\pi `$ dynamics. A constant or linear function of $`s`$ does not give an acceptable fit in terms of $`\chi ^2`$ for any value of $`E_r`$. A quadratic is the lowest order polynomial to give fits of acceptable confidence, as seen from the curvature of the magnitudes in Fig. 4 determined by the data-points with small error bars above 1.1 GeV, as shown in Fig. 3. Adding higher order terms in the polynomial representation of the $`\alpha _I`$ does not significantly improve the confidence level.
We see by comparing the upper plot of Fig. 2 with Fig. 5 that the $`I=1/2`$ and $`3/2`$ $`K\pi `$ components of $`D`$ decay are quite different from those of elastic scattering. The $`I=3/2`$ component is more than 50% of the $`I=1/2`$ above 1.1 GeV. This is surprising because the $`I=1/2`$ amplitude contains the broad $`K_0^{}(1430)`$, while the $`I=3/2`$ is entirely non-resonant. Consequently, any analysis of the $`D`$-decay data which neglects the $`I=3/2`$ amplitude must place this component incorrectly in other contributions and so lead to false conclusions about resonant branching fractions. Neither isospin component appears to contain an Adler zero and hence they grow towards threshold. At low mass we see the $`I=1/2`$ amplitude is dominant, whether this is because there is a near threshold $`\kappa `$ resonance requires a far better determination of the amplitude and phases.
An important ingredient in fixing the isospin components of the $`D`$-decay amplitude are the phases of $`K\pi `$ scattering in the elastic region. Here these have been set by the LASS experiment above 825 MeV and their continuation down to threshold by the predictions of Chiral Perturbation Theory. However, increased statistics on $`DK\pi \pi `$ decay would reduce the error bars on the amplitude and phases in Fig. 3. These might well become sufficiently small that an acceptable fit cannot be obtained without changing the $`K\pi `$ phases $`\delta ^I`$ away from their LASS results. This would mean that the $`D`$-decay results could indeed increase the precision within which the $`K\pi `$ elastic phases are known. Such improvement in $`D`$-decay statistics would be particularly welcome below 1 GeV down to threshold, where the current errors are sizeable — Figs. 4, 3. There it holds out the prospect of revealing whether there is indeed a low mass $`\kappa `$ or not. Such a state can only be exposed by analytic continuation of the appropriate amplitudes into the complex $`s`$-plane to deduce whether a pole exists or not. This requires far greater precision than the E791 data and its analysis we have considered here. $`D`$-decay results from $`B`$-factories may make this possible. The present analysis points the way.
Acknowledgements
We are most grateful to Brian Meadows for discussions that prompted this analysis. We acknowledge the key support of the EU-RTN Programme, Contract No. HPRN-CT-2002-00311, “EURIDICE”. One of us (MRP) is grateful to the hospitality provided by the LNF Spring Institute Programme when this work was completed. |
warning/0506/cond-mat0506280.html | ar5iv | text | # Structure of Stochastic Dynamics near Fixed Points
## I Introduction
In equilibrium statistical mechanics an important role is played by the principle of detailed balance and by the related fluctuation–dissipation theorem. Einstein used the principle that the excess energy that is put into each mode of an equilibrium system in the course of thermal fluctuations is also removed from the same mode by dissipative forces. This is implicit in his work on Brownian movement einstein05 , and explicit in later works on the photoelectric effect einstein12 , and on the relation between spontaneous and induced emission of electromagnetic radiation einstein17 . This principle was formulated as the principle of detailed balance by Bridgman bridgman28 , and used to explain Johnson noise in electrical circuits by Nyquist nyquist28 . This is related to the fact that the same processes that drive fluctuations in the neighborhood of a typical equilibrium configuration also drive the configuration back towards a typical equilibrium or steady state configuration when it is displaced from equilibrium by an amount which is small, but large compared with the fluctuations in thermal equilibrium. In this situation the equilibrium distribution in phase space is just the Boltzmann distribution, proportional to $`e^{\beta E}`$, where $`\beta `$ is inversely proportional to temperature, and $`E`$ is the energy of the point in phase space. Configurations that differ significantly from those that contribute to the minimum of the free energy are driven back to the neighborhood of this minimum by dissipative effects such as thermal or electrical conduction, or viscosity, and the magnitude of these effects is related to the equilibrium fluctuations of related variables.
In many situations there is no thermodynamic equilibrium, but external steady and fluctuating forces drive the system into a steady or very slowly varying state for which the principle of detailed balance does not hold. A light bulb powered by an external battery, or a chemical reaction in which the reactants are introduced at a steady rate and the products of the reaction removed at a steady rate, would both be examples of such a situation. Even in a situation which is almost in equilibrium, such as a system which is started in equilibrium at a local minimum of the free energy, but which can can go over a saddle point to a deeper minimum, the behavior near the saddle point does not satisfy the principle of detailed balance, since there is a current over the saddle point.
For such systems without detailed balance there is no general method of obtaining the equilibrium distribution from a knowledge of the steady and stochastic forces, such as the Boltzmann distribution provides for a system with detailed balance. In recent work one of us ao03 has developed a method valid near a stable fixed point which, even when detailed balance does not hold, obtains a cost function analogous to the energy for the Boltzmann distribution. If this method can be extended away from the linear region in the neighborhood of a fixed point it may provide a new method for dealing with problems of this sort zhu03 .
Great efforts have been spent on finding such a cost function ever since the work of Onsager onsager . Results up to 1990 have been summarized, for example, by van Kampen vankampen . In general such efforts have been regarded as not very successful cross .
In spite of the difficulty, there have been continuous efforts on the construction of cost function and related topics. Elegant results have been obtained in several directions. Tanase-Nicola and Kurchan kurchan have considered explicitly the saddle points of gradient systems. They started from the existence of potential or cost function to avoid the most difficult problem of the irreversibility. The gain is that they can now obtain a powerful computational method to count the saddle points and to compute the escape rate. They also provide an extensive list of related literature.
The study on the mismatch of the fixed points of the drift force and the extremals of the steady state distribution has been reviewed by Lindner et al. lindner . Rich phenomena have been observed, but the mismatch has been treated as “experimental” result. There is no mathematical/theoretical explanation on why it should happen.
In another survey the useful and constructive role played by the noise has been demonstrated by examples wio . It is argued that the noise is essential to establish the functions of dynamical systems. Again, the mismatch problem is encountered and the constructed potential function is often regarded as approximation.
From a different perspective, there has been an effort to provide a solid foundation for non-equilibrium processes based on the chaotic hypothesis gbg . The chaotic hypothesis presumes that the system is sufficiently chaotic that variation of parameters of the system leads to a unique parameter-dependent steady state, even though the Gibbs entropy change is not path independentruelle03 . Under this hypothesis an interesting and important fluctuation theorem has been obtained, which further suggests the existence of the Boltzmann-like steady state distribution function. Hence, a cost function very likely exists under this situation. A difficulty with this approach is that extremely few practical physical systems have been shown to satisfy the chaotic hypothesis.
Because the metastability is such an important phenomenon and because of the difficulty encountered in the construction of cost function, efforts have been made to go around the cost function problem when computing the life time of a metastable state. The effort results in the so-called Machlup-Onsager functional method, summarized by Freidlin and Wentzell freidlin . This approach has been actively pursued recently maier ; beri .
In this paper we give a careful discussion of the basis of our method of constructing a cost function in the linear region close to any fixed point, whether stable or unstable. We show that this approach gives an unambiguous prescription under wide conditions. The only case we have found which does not give an unambiguous expression for the cost function has a subspace within which the noise does not act, and out of which the force does not carry the state.
In section II we give a general discussion of a linear system with noise. In section III we show how a cost function can be constructed by a decomposition of the force matrix into two factors, one of which is a symmetric cost function matrix, and that this gives a probability density of exponential form, with the exponent proportional to the cost function. The general proof of this result is obtained in the Appendix, where we exploit the Jordan transformation for matrices with an incomplete set of eigenvectors. In section IV we discuss the singularities of the transformation, and identify two types of singularities, one of which corresponds to a flat subspace of the cost function, while the other corresponds to the possibility that the dynamics separates the system into two or more disjoint subspaces. In section V we discuss solutions other than the Boltzmann-like solution of the equation for a stationary distribution, and argue that such solutions may be significant in any attempt to extend this solution to the nonlinear regime. In section VI we discuss the significance of this decomposition of the force matrix, and its relation to the principle of detailed balance.
## II Stating the Problem
Many processes in natural sciences can be modeled quantitatively. One particularly important class of such modeling is that described by first order differential equations kaplan , supplemented by stochastic terms vankampen . We start with the nonlinear dynamic equation
$$\dot{x}=f(x)+\zeta (x,t),$$
(1)
which gives the stochastic evolution of the state represented by the real $`d`$-dimensional vector, $`x^\tau =(x_1,\mathrm{},x_d)`$. Here the superscript $`\tau `$ denotes the transposed vector. The force vector is the $`d`$-component $`f(x)`$, and this gives the deterministic time evolution of the system. For simplicity we take the noise term $`\zeta (x,t)`$ to be Gaussian white noise, with zero mean, $`\zeta (t)=0`$, and variance $`\zeta (x,t)\zeta ^\tau (x,t^{})=2D(x)\delta (tt^{})`$. The angular brackets denote the average over noise distribution, and $`\delta (t)`$ is the Dirac delta function. In this work we assume $`D(x)`$ to be independent of $`x`$. The probability distribution function $`\rho (x,t)`$ then satisfies the Fokker-Planck equation
$$\frac{\rho (x,t)}{t}=\frac{}{x_i}\left[f_i(x)+D_{ij}\frac{}{x_j}\right]\rho (x,t).$$
(2)
In the neighborhood of a fixed point, which we take to be at the origin, the force can be replaced by its linear approximation
$$f_i(x)=F_{ij}x_j.$$
(3)
It was noticed by Ao ao03 that, in the linear region near a stable fixed point, Eqs. (1) and (3) can be decomposed in the form
$$(S+A)\dot{x}=Ux+\xi (t),$$
(4)
where the symmetric matrix $`S`$ is semi-positive definite and the matrix $`A`$ is antisymmetric. The noise function $`\xi `$ has variance given by $`\xi (x,t)\xi ^\tau (x,t^{})=2S\delta (tt^{})`$. In this paper we adopt an equivalent, but simpler, approach of factorizing the force matrix as
$$F=(D+Q)U=(S+A)^1U,$$
(5)
where $`D`$ is the symmetric diffusion matrix, $`Q`$ is an antisymmetric matrix which can be determined, and $`U`$ is the symmetric cost function matrix, which was called a potential matrix in ref. References. This breaks the force $`Fx`$ into two components, $`F^{(d)}x=DUx`$, which generates a motion towards the origin if $`U`$ is positive definite, and $`F^{(c)}x=QUx`$, which gives a motion on the manifold of constant $`𝒰(x)=(1/2)x^\tau Ux`$. The quadratic form $`𝒰(x)`$ is the cost function.
It is immediately obvious that if the vector $`f`$ is replaced by $`F^{(d)}x=DUx`$ in Eq. (2), a Boltzmann-like stationary distribution of the form:
$$\rho (x)\mathrm{exp}\{𝒰(x)\}$$
(6)
satisfies the equation, since the current density
$$j_i(x)=\left[f_i(x)D_{ij}/x_j\right]\rho $$
(7)
vanishes. Since $`F^{(c)}x`$ generates a current density
$$j^{(c)}(x)=QUx\rho (x),$$
(8)
which is divergence free, and conserves $`𝒰`$ and $`\rho `$, the combination $`f(x)=[F^{(d)}+F^{(c)}]x`$ also conserves the distribution given by Eq. (6), so that this is a stationary solution of Eq. (2).
This decomposition of the force matrix allows an explicit time-independent solution of the Fokker-Planck equation to be written down. In ref. References the solution of the equation for $`Q`$ was obtained by a power series expansion, without any discussion of the convergence of this series. In the next section we show that there is a unique solution for the equation for $`Q`$ under rather wide conditions. It is not even required that the fixed point of $`F`$ be stable, although if it is not the stationary solution given by Eq. (6) is unbounded, and could only, at best, give a useful solution in a neighborhood of the fixed points with boundary conditions that do not perturb this solution too strongly.
It is worth noting that, under the coordinate transformation $`xy=M^1x`$, the mapping $`F`$ that relates $`\dot{x}`$ to $`x`$ transforms as
$$FM^1FM,$$
(9)
while the symmetric matrices $`U`$ and $`D`$, which represent quadratic forms, transform as
$$UM^\tau UM,DM^1D\left(M^\tau \right)^1,$$
(10)
and $`Q`$ transforms in the same way as $`D`$. These transformations preserve the separation of $`D+Q`$ into symmetric and antisymmetric parts.
Starting from the work of Onsager onsager , there has been an extensive literature on dynamical behavior near a stable fixed point sampling . The new construction clearly works for this important situation, and has indeed offered a new angle. However, it is not sufficient to generalize the decomposition of the force in Eq. (5) to the nonlinear regime by an equation of the form $`f=(D+Q)\mathrm{grad}𝒰`$, since the Boltzmann form, Eq. (6) may not satisfy the Fokker-Planck equation (2) if the antisymmetric matrix $`Q`$ is space dependent. The generalization to nonlinear systems needs further study. In the next three sections we establish the decomposition firmly in the linear regime, and investigate its limitations and implications.
## III Decomposition of the Force Matrix
In this section we develop a general method for making the decomposition of the force matrix given in Eq. (5). Since this equation, together with the symmetry of $`U,D`$, and antisymmetry of $`Q`$, leads to
$$U=(D+Q)^1F=F^\tau (DQ)^1,$$
(11)
the equation to determine $`Q`$ is
$$FQ+QF^\tau =FDDF^\tau .$$
(12)
This is a system of $`d(d1)/2`$ linear equations to determine the same number of independent components of $`Q`$, so it has a unique solution unless the set of equations is singular. Inversion of the matrix $`D+Q`$ then gives the matrix $`S+A`$ of Eq. (4).
Our method of solution is best illustrated by considering the case that $`F`$ is real symmetric or has distinct eigenvalues, so that it can be diagonalized in terms of its left and right eigenvectors. Equation (12) then takes the form
$$(\lambda _\alpha +\lambda _\beta )\stackrel{~}{Q}_{\alpha \beta }=(\lambda _\alpha \lambda _\beta )\stackrel{~}{D}_{\alpha \beta },$$
(13)
where the $`\lambda _\alpha `$ are the eigenvalues of $`F`$, and the tilde denotes this representation in terms of eigenvectors. This gives an immediate solution for $`Q`$ provided that no pair of the eigenvalues of $`F`$ adds up to zero. The eigenvalues can only add to zero when the fixed point is unstable, and this is discussed in section IV.
For completeness we must consider the general case with degenerate eigenvalues for asymmetric $`F`$, in which case there may not be a complete set of eigenvectors. This case is dealt with in the Appendix, using the Jordan representation of a nonsymmetric matrix.
## IV Singularities of the Transformation
There are two places in our argument where the transformation from the force matrix $`F`$ to the symmetric cost function matrix $`U`$ might be singular. Equation (12) for the antisymmetric matrix $`Q`$ can be solved, and we have an explicit solution in section III unless the determinant of the coefficients in $`d(d1)/2`$ inhomogeneous equations is zero. The second possibility is that the matrix $`D+Q`$ whose inverse appears in Eq. (11) might have zero determinant.
In section III we showed that the conditions for the equation for $`Q`$ to be singular are that two of the eigenvalues of $`F`$ sum to zero, or, as can be seen from Eq. (33) in the Appendix, when the null space of $`F^2`$ has two or more dimensions. Neither of these cases arise for a stable fixed point. There are two distinct cases of $`\lambda _\alpha +\lambda _\beta =0`$, according to whether the two eigenvalues are real eigenvalues of opposite sign, or whether they form a complex conjugate pair. We study the behavior of the eigenvalues and eigenvectors of $`U`$ in these two cases, assuming that the two eigenvalues of $`F`$ are nondegenerate, and that $`F`$ has no zero eigenvalue.
Instead of studying the eigenvectors of $`U`$ directly, we study the eigenvalues and eigenvectors of
$$U^1=F^1(D+Q)=R\mathrm{\Lambda }^1(\stackrel{~}{D}+\stackrel{~}{Q})R^\tau ,$$
(14)
which has the same eigenvectors but reciprocal eigenvalues. Here $`R`$ is the matrix whose columns are the right eigenvectors of $`F`$, and $`\mathrm{\Lambda }`$ is the diagonal matrix with the eigenvalues of $`F`$ as its diagonal elements. The generalization of these definitions of $`R`$ and $`\mathrm{\Lambda }`$ to the case where the eigenvalues of $`F`$ are not complete is given in the Appendix in Eqs. (28), (30) and (35). For a pair of eigenvalues with $`\lambda _\alpha +\lambda _\beta 0`$, with no other sums of two eigenvalues small and no other small individual eigenvalues, the only large terms in the matrix $`\stackrel{~}{U}^1=\mathrm{\Lambda }^1(\stackrel{~}{D}+\stackrel{~}{Q})`$ are, according to Eq. (13),
$$(\stackrel{~}{U}^1)_{\alpha \beta }=(\stackrel{~}{U}^1)_{\beta \alpha }=\frac{2\stackrel{~}{D}_{\alpha \beta }}{\lambda _\alpha +\lambda _\beta }.$$
(15)
When all other matrix elements are neglected, this approximation, combined with Eq. (14), gives
$$\underset{j}{}\left(U^1\right)_{ij}L_{\gamma j}\frac{2\stackrel{~}{D}_{\alpha \beta }}{\lambda _\alpha +\lambda _\beta }\left(\delta _{\gamma \alpha }R_{i\beta }+\delta _{\gamma \beta }R_{i\alpha }\right).$$
(16)
This equation, combined with the relation $`RL=I`$, shows that the approximate eigenvector corresponding to a large eigenvalue $`w^1`$ can be written as
$$a_\alpha R_{i\alpha }+a_\beta R_{i\beta }=\underset{j}{}(a_\alpha R_{j\alpha }+a_\beta R_{j\beta })\underset{\gamma }{}R_{j\gamma }L_{\gamma i},$$
(17)
provided the amplitudes and eigenvalues satisfy the equation
$$\begin{array}{cc}w^1a_\alpha & =\frac{2\stackrel{~}{D}_{\alpha \beta }}{\lambda _\alpha +\lambda _\beta }\left(a_\alpha _jR_{j\alpha }R_{j\beta }+a_\beta _jR_{j\beta }^2\right),\\ w^1a_\beta & =\frac{2\stackrel{~}{D}_{\alpha \beta }}{\lambda _\alpha +\lambda _\beta }\left(a_\alpha _jR_{j\alpha }^2+a_\beta _jR_{j\alpha }R_{j\beta }\right).\end{array}$$
(18)
This gives the two small real eigenvalues of $`U`$ as
$$w\frac{\lambda _\alpha +\lambda _\beta }{2\stackrel{~}{D}_{\alpha \beta }}\left(\underset{j}{}R_{j\alpha }R_{j\beta }\pm \sqrt{\underset{i}{}R_{i\alpha }^2\underset{j}{}R_{j\beta }^2}\right)^1,$$
(19)
and the corresponding eigenvectors as
$$\frac{a_\alpha }{a_\beta }\pm \sqrt{\frac{\underset{j}{}R_{j\beta }^2}{_jR_{j\alpha }^2}}.$$
(20)
For the case of a pair of real eigenvalues of opposite signs we can see that, as the sign of $`\lambda _\alpha +\lambda _\beta `$ is changed by a change in the parameters of $`F,D,`$ the stable and unstable manifolds of $`U`$ change places with one another. The stable and unstable manifolds of $`U`$ bisect the stable and unstable manifolds of $`F`$ in the original representation, as is obvious if one normalizes the real eigenvectors of $`\stackrel{~}{U}`$ by $`\widehat{R}_{i\alpha }=R_{i\alpha }/\sqrt{_jR_{j\alpha }^2}`$ and $`\widehat{R}_{i\beta }`$ similarly, giving eigenvectors, $`\widehat{R}_{i\alpha }\pm \widehat{R}_{i\beta }`$. In the limit $`\lambda _\alpha =\lambda _\beta `$, $`U`$ is flat in this two-dimensional subspace.
For a complex conjugate pair of eigenvalues with $`\lambda _\alpha +\lambda _\alpha ^{}0`$ the behavior is a little different. Using the property $`R_{j\alpha }=R_{j\beta }^{}`$, we find the term inside the bracket in Eq. (19) is always positive. The two eigenvalues of $`U`$ then have the same sign, so $`U`$ is either stable or unstable, depending on the sign of $`\lambda _\alpha +\lambda _\beta `$, in this two dimensional subspace. As the real part of $`\lambda _\alpha `$ changes sign, a two-dimensional stable manifold becomes unstable, or vice versa.
It can be seen from Eq. (32) in the Appendix that where one of the eigenvalues satisfying $`\lambda _\alpha +\lambda _\beta =0`$ corresponds to a higher dimensional subspace there may be higher order zeros of the eigenvalues of $`U`$.
If the matrix $`D`$ is positive definite there is no possibility that $`D+Q`$ could be singular. If $`u`$ is a vector in the null space of $`D+Q`$, we have
$$0=u^\tau (D+Q)u=u^\tau Du,$$
(21)
since the antisymmetry of $`Q`$ makes its expectation value vanish. Therefore $`D\pm Q`$ cannot be singular when $`D`$ is positive definite.
However, we do not usually want to specify that the noise acts on all coordinates. Typically, when two of the coordinates are the position and momentum of a particle, people will take the noise to change the momentum but not the position of the particle. However, Eq. (21) shows that for non-negative definite $`D`$, vectors in the null space of $`D+Q`$ are in the intersection of the null spaces of $`D`$ and $`Q`$. Equation (12) then shows that, for such a vector $`u`$ in the null space of $`D`$ and $`Q`$,
$$0=F(DQ)u=(D+Q)F^\tau u,$$
(22)
and so $`u`$ is only in this null space if $`F^\tau u`$, or any power of $`F^\tau `$ acting on $`u`$, is still in the null space.
This condition is in agreement with what one should expect. The noise does not have to act directly on all coordinates, but, if there is a subspace in which there is no noise, and which is left invariant by the motion, there can be no equilibration within that subspace except collapse towards a stable fixed point.
## V Other Stationary Solutions
Although the Boltzmann-like form given in Eq. (6) gives a stationary solution of the Fokker-Planck equation, it is only the unique solution under certain rather restrictive boundary conditions. One can see clearly why this might be an issue by considering the one-dimensional form of the equation near a stable fixed point, which can be written as
$$\frac{d^2\rho }{dx^2}+\frac{d}{dx}(x\rho )=0.$$
(23)
In addition to the Boltzmann-like solution $`\rho ^{(0)}=(2\pi )^{1/2}\mathrm{exp}(x^2/2)`$, this has a current-carrying solution of the form
$$\rho ^{(1)}(x)e^{x^2/2}_0^xe^{x^2/2}𝑑x^{}.$$
(24)
This expression is proportional to $`1/x`$ for large values of $`x`$, so, if the linear approximation to the equation is valid up to fairly large values of $`x`$, the coefficient of such a term must be exponentially small to prevent the probability density given by $`\rho ^{(0)}+\rho ^{(1)}`$ from being negative.
In $`d`$-dimensional systems there are similar solutions falling off like $`1/|x|^d`$ for large $`|x|`$. For such solutions of Eq. (2) the current at large distances from the origin is primarily driven by the linear force $`Fx`$, and the diffusive motion is a small correction, so, while the density falls off like $`|x|^d`$, the conserved current falls off like $`|x|^{d+1}`$. Again, current conservation shows that this contribution to the density must be positive and negative in different parts of space, so that its coefficient at the origin must be exponentially small, with an exponent that depends on the size of the region in which the linear approximation is valid.
Near a maximum of the cost function, where $`\rho ^{(0)}\mathrm{exp}(x^2/2)`$, the one-dimensional current-carrying solution has the form
$$\rho ^{(1)}(x)e^{x^2/2}_0^xe^{x^2/2}𝑑x^{}.$$
(25)
This grows at large distances in the same way as $`\rho ^{(0)}`$, with a change of sign at the origin. In $`d`$ dimensions a saddle point can sustain a relatively large current across it, because there are current-carrying states for which $`\rho `$ is of the same order of magnitude as $`\rho ^{(0)}`$.
To get such a current across a stationary point in a linear system, it is necessary to impose some external current sources and sinks. However, if we want to describe a nonlinear force field in terms of its approximately linear behavior in the neighborhood of its fixed points, adjacent neighborhoods can generate external current sources and sinks for one another, so we should not be surprised to find such currents if we linearize in a local region. These “external” currents will not only produce flows at the boundaries, but will shift the flow lines away from the surfaces of constant cost function $`U`$ shown in Eq. (8).
In the neighborhood of the minimum of the cost function, such current-carrying solutions resembling the one-dimensional Eq. (24) shift the maximum of the density away from fixed point, since the gradient of $`\rho ^{(1)}`$ is nonzero. Since, as we remarked in connection with this equation, the amplitude of such a term must fall off exponentially with the size of the region of linearization, in order to prevent negative densities, it should not be possible to obtain such a term by a conventional perturbation theory in the neighborhood of the stable fixed point. Our numerical exploration of nonlinear systems of this sort suggests that these current-carrying solutions are significant, since the maximum of the density is displaced from the zero of the force. One possibility is to introduce such current-carrying states, in addition to a $`\rho ^{(0)}`$ determined by the cost function, in order to make this method applicable to nonlinear systems.
## VI Discussion
The main result of this work is to show that, for a system with a deterministic motion controlled by a linear force, and a diffusive motion driven by constant white noise, the force matrix $`F`$ can be decomposed into two parts, $`DU`$ and $`QU`$. Provided the fixed point of $`F`$ is stable, the first of these components leads to a steady state distribution of the Boltzmann form, $`\mathrm{exp}(𝒰(x))`$, with no probability current, the usual form of an equilibrium distribution when detailed balance holds. The second part gives a flow on the surface of constant $`𝒰`$.
The cost function matrix $`U`$ can be diagonalized by an orthogonal transformation, and if it is positive definite, it can be transformed to the identity matrix by choosing a new scale for the variables. In this representation the dissipative part of the force matrix is $`F^{(d)}=D`$, which is suggestive of Einstein’s relation between diffusion and dissipation einstein05 , or of the fluctuation–dissipation theorem callen .
When the cyclic motion induced by $`F^{(c)}`$ is included the relation between the eigenvalues of $`F`$ and $`D`$ becomes more complicated. If the motion in a two-dimensional subspace is dominated by a fast cyclic motion within the subspace, there will be a complex conjugate pair of eigenvalues of $`F`$, so that they have a common relaxation rate given by the real part of the eigenvalues. These possibilities require much more detailed work than we have yet given them.
This is actually not so different from the situation usually encountered in statistical mechanics, at least if a classical rather than a quantum description is used. For a damped harmonic oscillator there is a cyclic motion in phase space, as well as the thermal noise and viscous damping acting on the momentum coordinate, while for black-body radiation in a cavity there is cyclic motion between electric and magnetic fields, in addition to the resistive damping and noise from the walls of the cavity.
What is remarkable is not that the steady state density can be written as the exponential of a cost function, since if there is a steady state we could always define the cost function as minus the logarithm of the steady state density. We find it remarkable that for a linear stochastic system of this sort it is generally true that the force can be decomposed into two parts, one of which gives detailed balance in its strictest sense, while the other gives a cyclic motion on the surfaces of constant cost function.
The question of whether this technique can be extended to nonlinear systems is an important one, but requires careful investigation. We have done preliminary work based on perturbative inclusion of nonlinear terms, and made comparison with numerical calculations. It is clear that Eqs. (4) and (5) need some modification for nonlinear effects, and also that if we try to solve the problem by matching limited regions in which linearity holds approximately, the matching may introduce solutions of the linearized Fokker-Planck equation other than the Boltzmann-like $`\rho ^{(0)}`$. Different regions will have to serve as sources and sinks for their adjacent regions.
###### Acknowledgements.
Discussions with M. Cross, S.W. Rhee, L. Yin, and X.-M. Zhu are highly appreciated. We are particularly grateful to E. Siggia, whose criticisms of this paper have led us to make some important modifications. This work was supported in part by Institute for Systems Biology (P.A.), by the National Institutes of Health through Grant number HG002894, and by the National Science Foundation through Grant number DMR-0201948.
## Appendix A Appendixes
For cases in which the force matrix $`F`$ has degenerate eigenvalues and is not symmetric, so that it may not have a complete set of eigenvectors, we use the Jordan transformation hirsch of a general real matrix. The Jordan transformation uses a complete set of independent column vectors $`v^\alpha `$ with the property
$$Fv^\alpha =\lambda _\alpha v^\alpha +\mu _{\alpha 1}v^{\alpha 1},$$
(26)
where $`\mu _\alpha `$ is zero if $`\lambda _\alpha \lambda _{\alpha +1}`$ and is either unity or zero for $`\lambda _\alpha =\lambda _{\alpha +1}`$. The set of row vectors $`u^\alpha `$ with the orthonormality property $`u^\alpha v^\beta =\delta _{\alpha \beta }`$ then satisfies the equation
$$u^\alpha F=\lambda _\alpha u^\alpha +\mu _\alpha u^{\alpha +1}.$$
(27)
Let us define matrices $`R`$ and $`L`$ as $`R_{i\alpha }=v_i^\alpha `$ and $`L_{\alpha i}=u_i^\alpha `$. Since the vectors $`u^\alpha `$ are orthonormal to the set $`v^\alpha `$, we have $`LR=I`$, so that these matrices are inverse to one another. $`R_{i\alpha }`$, $`L_{\alpha i}`$ are real for real $`\lambda _\alpha `$. For complex $`\lambda _\alpha `$, its complex conjugate is also an eigenvalue, $`\lambda _\beta =\lambda _\alpha ^{}`$, since $`F`$ is a real matrix. In this case $`R_{\alpha i}=R_{\beta i}^{}`$, $`L_{i\alpha }=L_{i\beta }^{}`$. The Jordan transformation is then given by
$$LFR=\mathrm{\Lambda }.$$
(28)
The matrix $`\mathrm{\Lambda }`$ is now block diagonal, where each nonzero block has identical diagonal elements, which are degenerate eigenvalues, and unity in each place immediately above the diagonal.
With this result, Eq. (12) can be rewritten as
$$\mathrm{\Lambda }\stackrel{~}{Q}+\stackrel{~}{Q}\mathrm{\Lambda }^\tau =\mathrm{\Lambda }\stackrel{~}{D}\stackrel{~}{D}\mathrm{\Lambda }^\tau ,$$
(29)
where
$$\stackrel{~}{Q}=LQL^\tau ,\stackrel{~}{D}=LDL^\tau .$$
(30)
The matrix $`\stackrel{~}{Q}`$ remains antisymmetric and $`\stackrel{~}{D}`$ remains symmetric. In this representation Eq. (29) takes on the form
$`(\lambda _\alpha +\lambda _\beta )\stackrel{~}{Q}_{\alpha \beta }+\mu _\alpha \stackrel{~}{Q}_{\alpha +1,\beta }+\mu _\beta \stackrel{~}{Q}_{\alpha ,\beta +1}`$
$`=(\lambda _\alpha \lambda _\beta )\stackrel{~}{D}_{\alpha \beta }+\mu _\alpha \stackrel{~}{D}_{\alpha +1,\beta }\mu _\beta \stackrel{~}{D}_{\alpha ,\beta +1}.`$ (31)
It has the solution, for $`\alpha <\beta `$,
$`\stackrel{~}{Q}_{\alpha \beta }`$ $`=`$ $`{\displaystyle \frac{\lambda _\alpha \lambda _\beta }{\lambda _\alpha +\lambda _\beta }}\stackrel{~}{D}_{\alpha \beta }`$ (32)
$`+2\lambda _\beta {\displaystyle \underset{\mu 1}{}}{\displaystyle \frac{(1)^{\mu 1}}{(\lambda _\alpha +\lambda _\beta )^{\mu +1}}}\stackrel{~}{D}_{\alpha +\mu ,\beta }`$
$`2\lambda _\alpha {\displaystyle \underset{\nu 1}{}}{\displaystyle \frac{(1)^{\nu 1}}{(\lambda _\alpha +\lambda _\beta )^{\nu +1}}}\stackrel{~}{D}_{\alpha ,\beta +\nu }`$
$`+2{\displaystyle \underset{\mu ,\nu 1}{}}{\displaystyle \frac{(\mu +\nu 1)!}{\mu !\nu !}}{\displaystyle \frac{(1)^{\mu +\nu +1}}{(\lambda _\alpha +\lambda _\beta )^{\mu +\nu +1}}}`$
$`(\mu \lambda _\beta \nu \lambda _\alpha )\stackrel{~}{D}_{\alpha +\mu ,\beta +\nu },`$
where the sums go over all values of $`\mu ,\nu `$ for which the indices lie within the same block of the block-diagonal matrix $`\mathrm{\Lambda }`$. For the case $`\lambda _\alpha =\lambda _\beta `$ this reduces to
$`\stackrel{~}{Q}_{\alpha \beta }`$ $`=`$ $`{\displaystyle \underset{\mu 1}{}}{\displaystyle \frac{(1)^\mu }{(2\lambda _\alpha )^\mu }}\stackrel{~}{D}_{\alpha +\mu ,\beta }+{\displaystyle \underset{\nu 1}{}}{\displaystyle \frac{(1)^\nu }{(2\lambda _\alpha )^\nu }}\stackrel{~}{D}_{\alpha ,\beta +\nu }`$ (33)
$`{\displaystyle \underset{\mu ,\nu 1}{}}{\displaystyle \frac{(\mu +\nu 1)!}{\mu !\nu !}}{\displaystyle \frac{(1)^{\mu +\nu }}{(2\lambda _\alpha )^{\mu +\nu }}}(\mu \nu )\stackrel{~}{D}_{\alpha +\mu ,\beta +\nu }.`$
For the case in which all the $`\mu _\alpha `$ are zero, Eq. (32) is equivalent to Eq. (13).
In this Jordan representation of force matrix $`F`$, the state variable is transformed by
$$y=Lx,$$
(34)
and the corresponding transformation of $`U`$ is given by
$$\stackrel{~}{U}=R^\tau UR.$$
(35)
Equation (11) then becomes
$$(\stackrel{~}{D}+\stackrel{~}{Q})\stackrel{~}{U}=\mathrm{\Lambda },$$
(36)
which is form-invariant with Eq. (11) under the Jordan transformation. |
warning/0506/hep-th0506232.html | ar5iv | text | # 1 Introduction
## 1 Introduction
D3 branes living at conical Calabi-Yau singularities are a good laboratory for the AdS/CFT correspondence since its early days. The world-volume theory on the branes is dual to a type IIB background of the form $`AdS_5\times H`$, where $`H`$ is the horizon manifold . Supersymmetry requires that $`H`$ is a Sasaki-Einstein manifold. Until few months ago, the only known Sasaki-Einstein metrics were the round sphere $`S^5`$ and $`T^{1,1}`$, the horizon of the conifold. Recently, various infinite classes of new regular Sasaki-Einstein metrics were constructed and named $`Y^{p,q}`$ and $`L^{p,q,r}`$. For infinite values of the integers $`p,q,r`$ one obtains smooth Sasaki-Einstein manifolds. With the determination of the corresponding dual gauge theory (see for the $`Y^{p,q}`$ manifolds and for the $`L^{p,q,r}`$), new checks of the AdS/CFT correspondence were possible . As well known, the central charge of the CFT and the dimension of some operators can be compared with the volumes of $`H`$ and of some of its submanifolds. In particular, the a-maximization technique now allows for a detailed computation of the relevant quantum field theory quantities. Needless to say, the agreement of the two computations is perfect.
The number of explicit metrics for Sasaki-Einstein horizons than can be used in the AdS/CFT correspondence is rapidly increasing. However, to demystify a little bit the importance of having an explicit metric, we should note that all relevant volumes are computed for calibrated divisors. This means that these volumes can be computed without actually knowing the metric. There exist moreover a beautiful geometrical counterpart of the a-maximization : this is the volume minimization proposed in for determining the Reeb vector for toric cones. This procedure only relies on the vectors defining the toric fan. This suggests that with a correspondence between toric diagrams and gauge theories, many checks of the AdS/CFT correspondence can be done without an explicit knowledge of the metric. It is the purpose of this paper, indeed, to show that the knowledge of the toric data is sufficient to determine many properties of the dual gauge theory and to perform all the mentioned checks, for every singularity.
The precise correspondence between conical Calabi-Yau singularities and superconformal gauge theories is still unknown. However, a remarkable progress has been recently made for the class of Gorenstein toric singularities. The brane tiling (dimers) construction , an ingenious generalization of the Brane Boxes , introduces a direct relation between an Hanany-Witten realization for gauge theory and the toric diagram. In particular, from the quiver associated with a non-chiral superconformal gauge theory one can determine the dual brane tiling configuration, a dimer lattice. It is then possible to associate a toric diagram with each of these lattices, identifying the dual Calabi-Yau. The inverse process (to associate a gauge theory with a given singularity) is more difficult. However, for the mentioned checks of the AdS/CFT correspondence, we don’t really need the full quiver description of the gauge theory. We just need to know the R-charges and the multiplicities of chiral fields. In this paper, elaborating on existing results in the literature , we propose a general assignment of charges and multiplicities for the gauge theory dual to a generic Gorenstein singularity. This assignment is made using only the toric data of the singularity. We then compare the result of a-maximization with that of volume minimization showing that the two procedures are completely equivalent. This agreement is remarkable. We have two different algebraic procedures for computing the R-symmetry charges of the fields and the volumes. The first is based on the maximization of the central charge . The second one can be efficiently encoded in a geometrical minimization procedure for determining the Reeb vector . The two procedures deal with different test quantities (the R-charges on one side and the components of the Reeb vector on the other) and with different functions to be extremized. However, we will show that, with a suitable parametrization, the two functions ($`a`$ and the inverse volume) are equal, even before extremization.
The agreement of results in the gauge theory and the supergravity side can be regarded as a general non-trivial check of the AdS/CFT correspondence, valid for all the theories living on branes at toric singularities.
The paper is organized as follows. In Section 2 we briefly review the general features of the gauge theories dual to conical singularities. In Section 3 we propose the assignment of R-charges and multiplicities for the gauge theory in terms of geometrical data. In Section 4 we show the equivalence of the a-maximization and the volume minimization. Section 5 contains several examples based on known gauge theories and various observations. In particular, as a by product of our analysis, we discuss in detail the case of the manifolds $`X^{p,q}`$ introduced in whose general analysis was missing in the literature. We also make some observations on the identification of fields using the brane tiling technology. Finally, the Appendix contains the proofs of various results that are too long and boring for the main text.
## 2 Generalities about the gauge theory
We consider $`N`$ D3-branes living at a conical Gorenstein singularity. The internal manifold is a six-dimensional symplectic toric cone; its base, or horizon, is a five-dimensional Sasaki-Einstein manifold $`H`$ . As well known, the $`𝒩=1`$ gauge theory living on the branes is superconformal and dual in the AdS/CFT correspondence to the type IIB background $`AdS_5\times H`$. The gauge theory on the world-volume of the D3 branes is not chiral and represents a toric phase , where all gauge groups have the same number of colors $`N`$ and the only matter fields are bi-fundamentals. By applying a Seiberg duality we can obtain a different theory that flows in the IR to the same CFT. If we dualize a gauge group with number of flavors equal to $`2N`$ we remain in a toric phase where all gauge groups have number of colors $`N`$. In this process the number of gauge groups remains constant but the number of matter fields changes. In a toric phase the following relation between the number of gauge groups $`F`$, the number of chiral fields $`E`$ and the number of terms in the superpotential $`V`$
$$VE+F=0$$
(2.1)
is valid . Indeed for a gauge theory living on branes placed at the tip of toric CY cone, one can extend the quiver diagram, drawing it on a torus $`T^2`$. The dual graph, known as the brane tiling associated with the gauge theory , has $`F`$ faces, $`E`$ edges and $`V`$ vertices and it is still defined on a torus. The previous formula then follows from the Euler formula for a torus .
We can assign an R-charge to all the chiral fields. The most general non-anomalous R-symmetry is determined by the cancellation of anomalies for each gauge group and by the requirement that each term in the superpotential has R-charge $`2`$. This would seem to imply $`F+V=E`$ linear conditions for $`E`$ unknowns with an unique solution. However, in the cases we are interested in, not all the conditions are linearly independent. This is reflected by the fact that the R-symmetry can mix with all the non anomalous $`U\left(1\right)`$ global symmetries. We can count the number of global non-anomalous $`U\left(1\right)`$ symmetries from the number of massless vectors in the $`AdS`$ dual. Since the manifold is toric, the metric has three $`U\left(1\right)`$ isometries. One of these (the Reeb one) corresponds to the R-symmetry while the other two are related to non-anomalous global $`U\left(1\right)`$s. Other gauge fields in $`AdS`$ come from the reduction of the RR four form on the non-trivial three-spheres in the horizon manifold $`H`$. The number of three-cycles depends on the topology of the horizon, and, as we will review soon, can be computed using the toric data of the singularity. In the supergravity literature the vector multiplets obtained from RR four form are known as the Betti multiplets. On the gauge theory side, these gauge fields correspond to baryonic symmetries.
At the fixed point, only one of the possible non-anomalous R-symmetry enters in the superconformal algebra. It is the one in the same multiplet as the stress-energy tensor. The actual value of the R-charges at the fixed point can be found by using the a-maximization technique . As shown in , we have to maximize the a-charge
$$a\left(R\right)=\frac{3}{32}\left(3\mathrm{T}\mathrm{r}R^3\mathrm{Tr}R\right)$$
(2.2)
It is not difficult to show that the absence of anomalies implies $`\mathrm{Tr}R=0`$ so that we can equivalently maximize $`\mathrm{Tr}R^3`$.
The results of the maximization give a complete information about the values of the central charge and the dimensions of chiral operators at the conformal fixed point. These can be compared with the prediction of the AdS/CFT correspondence . The first important point is that the central charge is related to the volume of the internal manifold
$$a=\frac{\pi ^3}{4\mathrm{V}\mathrm{o}\mathrm{l}\left(H\right)}$$
(2.3)
Moreover, recall that in the AdS/CFT correspondence a special role is played by baryons. The gravity dual describes a theory with $`SU\left(N\right)`$ gauge groups. The fact that the groups are $`SU\left(N\right)`$ and not $`U\left(N\right)`$ allows the existence of dibaryons. Each bi-fundamental field $`\mathrm{\Phi }_\alpha ^\beta `$ gives rise to a gauge invariant baryonic operator
$$ϵ^{\alpha _1\mathrm{}\alpha _N}\mathrm{\Phi }_{\alpha _1}^{\beta _1}\mathrm{}\mathrm{\Phi }_{\alpha _N}^{\beta _N}ϵ_{\beta _1\mathrm{}\beta _N}$$
It is sometime convenient to think about the baryonic symmetries as non-anomalous combinations of $`U\left(1\right)`$ factors in the enlarged $`U\left(N\right)`$ theory. In the AdS dual the baryonic symmetries correspond to the reduction of the RR four form and the dibaryons are described by a D3-brane wrapped on a non-trivial three cycle. The R-charge of the $`i`$-th field can be computed in terms of the volume of the corresponding cycle $`\mathrm{\Sigma }_i`$ using the formula
$$R_i=\frac{\pi \mathrm{Vol}\left(\mathrm{\Sigma }_i\right)}{3\mathrm{V}\mathrm{o}\mathrm{l}\left(H\right)}$$
(2.4)
## 3 Geometrical formulae for the R-charges
In this Section we propose a general formula for the R-charges and the multiplicity of chiral fields based only on the toric data <sup>1</sup><sup>1</sup>1For the necessary elements of toric geometry see and the review part of .. This proposal is the natural combination of existing results and it is substantially implicit in previous papers on the subject. It is based on a formula for multiplicities first derived using mirror symmetry . The same proposal appeared for the case of $`L^{p,q:r,s}`$ manifolds in , under the name of “folded quiver”.
The fan associated with a six-dimensional symplectic toric cone is generated by $`d`$ integers primitive vectors in $`^3`$, which we call $`V_i`$, $`i=1,2\mathrm{}d`$. When the cone is a Calabi-Yau manifold, we can perform an $`SL(3,)`$ transformation to put the first coordinates of the $`V_i`$’s equal to 1. The intersection of the fan with the plane of points having the first coordinate $`x=1`$ is thus a convex polygon $`P`$, called toric diagram, and we shall call the vectors associated to its sides $`v_i`$, $`i=1,2\mathrm{}d`$, as in Figure 2. In Figure 2 we draw the corresponding $`(p,q)`$ web: the vectors $`v_i`$ have the same length than the edges of the polygon $`P`$ <sup>2</sup><sup>2</sup>2With a little abuse of notation we call $`v_i`$ both the sides of $`P`$ and the vectors of the $`(p,q)`$ web. In fact they differ only by a rotation of $`90^o`$. When some of the sides of the polygon $`P`$ pass through integer points, that is for singular horizons, we should consider more complicated $`(p,q)`$ webs; here we are ignoring such subtleties. We claim that this does not affect the process of a-maximization, since it is equivalent to setting to zero the charges $`b_i`$ associated with integers points on the sides of $`P`$ (see subsection 5.6).. Let us also define the symbol:
$$w_i,w_jdet(w_i,w_j)$$
(3.1)
that is the determinant of the matrix with $`w_i`$ and $`w_j`$ as first and second line respectively, where $`w_i`$ and $`w_j`$ are two vectors in the plane of $`P`$. This is the oriented area of the parallelogram generated by $`w_i`$ and $`w_j`$.
Some of the data of the gauge theory can be extracted directly from the geometry of the cone. In particular, there exist simple formulae for the number of gauge groups $`F`$ and the total number of chiral bi-fundamental fields $`E`$
$`F`$ $`=`$ $`2\mathrm{A}\mathrm{r}\mathrm{e}\mathrm{a}\left(P\right)`$
$`E`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i,j}{}}\left|v_i,v_j\right|`$ (3.2)
Notice that the expression for $`E`$ refers to a particular toric phase of the gauge theory. The number of toric phases of a theory can be large; hopefully, the value of $`E`$ in formula (3.2) refers to the phase with the minimal number of fields.
The R-charges and the multiplicity of fields with given R-charge are more difficult to determine. Here we make a proposal based on the following general observations. Each chiral field is associated with a dibaryon and, consequently, with a supersymmetric three cycle in the horizon $`H`$. The cone over this cycle is a divisor in the symplectic cone $`C\left(H\right)`$. Each edge $`V_i`$ in the fan determines a divisor $`D_i`$ and the collection of the $`D_i`$, subject to some relations, is a complete basis of divisors for $`C\left(H\right)`$ . We can therefore associate a type of chiral field to each vector $`V_i`$ and assign it a trial R-charge $`a_i`$. It is important to stress that more than one chiral field is associated with a single divisor: as pointed out in , a D3 brane wrapped on the cycle $`\mathrm{\Sigma }`$ may have more than one supersymmetric vacuum and each of these corresponds to a different bi-fundamental but with the same R-charge. As shown in the volumes of the base three cycles $`\mathrm{\Sigma }_i`$ of the divisors $`D_i`$ satisfy the relation
$$\underset{i=1}{\overset{d}{}}\mathrm{Vol}\left(\mathrm{\Sigma }_i\right)=\frac{6}{\pi }\mathrm{Vol}\left(H\right)$$
(3.3)
which implies, using formula (2.4), $`_{i=1}^da_i=2`$. In general, these $`d`$ fields will not exhaust all the different types of chiral fields. We expect the existence of other dibaryons obtained from divisors which are linear combinations of the $`D_i`$s. The R-charges of the corresponding fields will not be independent but they will be determined as a linear combination of the $`a_i`$s. Indeed, we claim that the $`a_i`$s parametrize the most general R-symmetry <sup>3</sup><sup>3</sup>3For a generalization of this sentence to singular horizons see subsection 5.6.. The number of independent parameters in the trial R charge is equal to the number of global $`U\left(1\right)`$ symmetries. We always have two global symmetries from the toric action and a number of baryonic symmetries equal to the number of three cycles. As shown in , the latter is equal to $`d3`$; each baryonic symmetry $`B_a`$ is indeed associated with a linear relation among the edges $`V_i`$
$$\underset{i=1}{\overset{d}{}}B_i^aV_i=0$$
(3.4)
and there are exactly $`d3`$ such relations. In conclusion, we have a total number of $`d1`$ global $`U\left(1\right)`$ symmetries which matches the number of independent parameters $`a_i`$.
Collecting all these pieces of information, we can propose the following assignments of R-charges and multiplicities for the chiral fields in the gauge theory:
* Associate with each edge vector $`V_i`$ a chiral field with trial R-charge $`a_i`$, with the constraint,
$$\underset{i=1}{\overset{d}{}}a_i=2$$
(3.5)
* Call $`C`$ the set of all the unordered pairs of vectors in the $`(p,q)`$ web; we label an element of $`C`$ with the ordered indexes $`(i,j)`$, with the convention that the vector $`v_i`$ can be rotated to $`v_j`$ in the counter-clockwise direction with an angle $`180^o`$. With our conventions $`\left|v_i,v_j\right|=v_i,v_j`$. Associate with any element of $`C`$ the divisor
$$D_{i+1}+D_{i+2}+\mathrm{}D_j$$
(3.6)
and a type of chiral field in the field theory with multiplicity $`v_i,v_j`$ and R-charge equal to $`a_{i+1}+a_{i+2}+\mathrm{}a_j`$. The indexes $`i`$, $`j`$ are always understood to be defined modulo $`d`$. For example in Figure 2 the field associated to the pair $`(d,3)`$ has R-charge $`a_1+a_2+a_3`$ and multiplicity $`v_d,v_3`$. The total number of fields is the sum of all the multiplicities:
$$E\underset{(i,j)C}{}\left|v_i,v_j\right|$$
(3.7)
and thus reproduces formula (3.2).
More generally, we can assign global symmetry charges to all the fields. The algorithm is very similar to that for R-charges:
* Assign global charges $`a_i`$ to the fields corresponding to vertices $`V_i`$. The only difference is that now $`a_i`$ satisfy the relation:
$$\underset{i=1}{\overset{d}{}}a_i=0$$
(3.8)
* The global charges of composite chiral fields are then: $`a_{i+1}+a_{i+2}\mathrm{}+a_j`$ for the fields corresponding to $`(i,j)`$ in $`C`$.
With a small abuse of notation, we will use the same letter $`a_i`$ for R and global symmetries; in the first case they satisfy $`_{i=1}^da_i=2`$, while in the latter $`_{i=1}^da_i=0`$.
Note that with the assignment (3.8) we parametrize all the possible $`d1`$ global symmetries, the $`d3`$ baryonic ones and the two flavor ones. We can explicitly identify the baryonic symmetries as follows. As shown in , the chiral fields associated with the edges $`V_i`$ have a charge under the baryonic symmetry $`B_a`$ equal to the coefficient $`B_i^a`$ in the linear relations (3.4). Notice that the baryonic charges of the fields associated with the edges $`V_i`$ sum up to zero
$$\underset{i=1}{\overset{d}{}}B_i^a=0$$
(3.9)
and therefore satisfy eq. (3.8) as a consequence of the Calabi-Yau condition; the latter requires that all the vectors $`V_i`$ lie on a plane which, in our conventions, means that the first coordinate of all $`V_i`$ is $`1`$. In conclusion, among the global symmetries, those satisfying also the constraint (3.4) (with $`a_i=B_i^a`$) are the baryonic ones, the remaining two (for which there is not a natural basis, being mixed with baryonic symmetries ) are the flavor ones.
We conjecture that for every Gorenstein toric singularity there exists a toric phase of the dual gauge theory where the R-charges and the multiplicities of all chiral fields can be computed with the algorithm above. This toric phase has generally the minimal number of chiral fields (3.2), as we have checked in many known cases. To be concrete look at Figure 4 corresponding to $`L^{p,q;r,s}`$. There are six kinds of fields: the four with charge $`a_i`$, fields with charge $`a_3+a_4`$ and others with charge $`a_2+a_3`$, but there are not for instance fields with charge $`a_1+a_2`$, since the region formed by $`v_4`$ and $`v_2`$ which includes $`a_1`$ and $`a_2`$ in the $`(p,q)`$ web has always an angle greater than $`180^o`$. Note that in general the number of different kinds of fields is $`d\left(d1\right)/2`$, the number of elements of $`C`$. Note also that the R-charges of composite chiral fields can be written as sum of consecutive $`a_i`$s; since $`P`$ is convex the ordering of vectors $`v_i`$ in the $`(p,q)`$ web is always equal to the ordering of $`v_i`$ in $`P`$.
With this assignment, we have a trial central charge $`a`$ given by:
$$a=\frac{9}{32}\mathrm{tr}R^3=\frac{9}{32}\left(F+\underset{(i,j)C}{}\left|v_i,v_j\right|\left(a_{i+1}+a_{i+2}+\mathrm{}a_j1\right)^3\right)$$
(3.10)
Recall that $`F`$ is the double area of the polygon $`P`$ (3.2). The values of the R-charges $`a_i`$ can be found by (locally) maximizing this formula. Note that this formula and the algorithm proposed above are obviously invariant under translations and $`SL(2,)`$ transformations in the plane of $`P`$, since $`v_i,v_j`$ are conserved, also in sign <sup>4</sup><sup>4</sup>4If the determinant is $`1`$ all the signs are reversed and so relative orientations do not change..
We can make several checks of this proposal. First of all, it is easy to compare the proposal to the case where the quiver gauge theory is explicitly known. Several examples are discussed in Section 5. In some cases, the fields and their multiplicity can be determined by using mirror symmetry; this was done for the toric delPezzo in where the formula for the multiplicities based on the $`(p,q)`$ web first appeared. The multiplicity of the fields associated with the edges $`V_i`$ was computed in and agrees with our proposal:
$$\text{multiplicity of fields }(i1,i)C=v_{i1},v_i=\text{det}(V_{i1},V_i,V_{i+1})$$
(3.11)
since the fields corresponding to the pair $`(i1,i)`$ in $`C`$ have R-charge $`a_i`$; we have used that the first coordinates of $`V_i`$ are equal to 1. A proposal identical to ours was used in to determine the gauge theory for the $`L^{p,q,r}`$ manifolds.
We can next study the consistency of our proposal with the general properties of the $`U\left(1\right)`$ symmetries in our theories. First of all, we must have
$$\mathrm{Tr}G=0.$$
(3.12)
where $`G`$ is a general R-charge or global symmetry charge. In particular $`\mathrm{Tr}R=\mathrm{Tr}B^a=0`$. The proof of this formula is relatively easy and is reported in the Appendix. Another non trivial check of our proposal is the proof, reported in the Appendix, that, for baryonic symmetries,
$$\mathrm{Tr}B_a^3=0$$
(3.13)
This condition, which is true also for mixed baryonic symmetries, is a consequence of the vanishing of the cubic t’Hooft anomaly for a baryonic symmetry. This follows from the fact that on the stack of D3 branes in type IIB the baryonic symmetries are actually gauged. The counterpart of this statement in the AdS dual is that cubic anomalies are computed from the Chern-Simons terms in the five dimensional supergravity and no such term can contain three vector fields coming from reduction of the RR four-form .
The best check of the proposal is however the computation of the R-charges at the fixed point using a-maximization and the comparison with volumes of three cycles in $`H`$. Now that we have an algorithm to extract the field content of the gauge theory from the toric diagram, it is not difficult to write down an algorithm on a computer and check the agreement of a-maximization with Z-minimization on arbitrary large polytopes. The complete agreement of the a-maximization with the volume minimization of will be discussed in details in the next Section, where a general analytic proof will be given.
We finish this Section by making some remarks about other toric phases of the same CFT with more chiral fields than the minimal phase presented above. In practical examples we often meet toric phases with the same trial central charge $`a`$ than the minimal phase; these phases generally contain all the kinds of fields of the minimal phase, but with greater multiplicities. In fact there are other possible assignments of R-charges and multiplicities leading to the same $`a`$ charge. For example, to each element in $`C`$ we may assign two different types of chiral fields, one associated with the divisor
$$C_{i,j}=D_{i+1}+D_{i+2}+\mathrm{}+D_j$$
(3.14)
with R-charge $`a_{i,j}=a_{i+1}+\mathrm{}a_j`$, and a second one associated with the divisor <sup>5</sup><sup>5</sup>5Recall that, in our conventions, the indexes $`i,j`$ are always defined modulo $`d`$.
$$\underset{i=1}{\overset{d}{}}D_iC_{i,j}=D_{j+1}+D_{j+2}\mathrm{}+D_i$$
(3.15)
with R-charge $`a_{j+1}+\mathrm{}a_i=2a_{i,j}`$. If we assign multiplicities $`n_{i,j}`$ and $`\stackrel{~}{n}_{i,j}`$ to the two types of fields with the constraint
$$n_{i,j}\stackrel{~}{n}_{i,j}=\left|v_i,v_j\right|$$
(3.16)
it is easy to see that the equations $`\mathrm{Tr}R=\mathrm{Tr}B_a=\mathrm{Tr}B_a^3=0`$ are still satisfied. Moreover the expression for the trial central charge $`a`$ is unchanged. Indeed the contribution of the integers $`n_{i,j}`$ to the central charge cancels:
$$n_{i,j}\left(a_{i,j}1\right)^3+\stackrel{~}{n}_{i,j}\left(1a_{i,j}\right)^3\left|v_i,v_j\right|\left(a_{i,j}1\right)^3.$$
(3.17)
The formula (3.2) for the number of chiral fields is obviously no more satisfied. Each time a field is split and a new arbitrary integer $`n_{i,j}`$ is introduced, the total numbers of fields increase. Formula (3.2) is strictly valid for the minimal presentation. We do not expect that for all arbitrary choices of $`n_{i,j}`$ and pairs $`(i,j)`$ there exists a non minimal toric phase with multiplicities of chiral fields described by this splitting mechanism, even though many known toric phases are characterized by multiplicities determined in this simple way. In fact all the examples of (non minimal) toric phases considered in this paper are described by this splitting mechanism, and it would be interesting to know whether this is true in general.
## 4 a-maximization is volume minimization
For the purposes of the AdS/CFT correspondence, the R-charges of the chiral fields have to be matched with the volumes of the three-cycles bases $`\mathrm{\Sigma }_i`$ of the corresponding divisors. In the previous Section we proposed a formula for computing the R-charges and the trial central charge $`a`$ directly from the toric diagram. Moreover in it was shown that all the geometric information on the volumes can be extracted from the toric data, through the process known as volume minimization (or Z-minimization), without any explicit knowledge of the metric. The reason for that is the following: supersymmetric cycles are calibrated and the volumes can be extracted only from the Kahler form on the cone. Therefore now it is possible to compare directly R-charges in the gauge theory and volumes in the geometry, checking the correctness of the AdS/CFT predictions for every toric CY cone. In this Section we discuss the equivalence of a-maximization and Z-minimization.
We start by reviewing the work of and reducing their formulas in the plane containing the convex polygon $`P`$. The Reeb vector $`K`$ of a symplectic toric cone can be expanded in a basis $`e_i`$ for the $`T^3`$ effective action on the fiber:
$$K=\underset{i=1}{\overset{3}{}}b_ie_i$$
(4.1)
where the vector of coordinates $`b=(b_1,b_2,b_3)`$ lives inside the toric fan of the cone. The Reeb vector is associated with an R-symmetry in the dual gauge theory; by varying the vector we change the R-symmetry by mixing it with the global symmetries. From the geometrical point of view, the variation of the Reeb vector changes the metric and the volumes. For only one choice of vector $`\overline{b}`$ there exists a Calabi-Yau metric for the cone. The vector $`\overline{b}`$ has $`\overline{b}_1=3`$ and can be determined through the minimization of a certain function Z of the variables $`b_2`$ and $`b_3`$ . We rephrase this process in the plane containing $`P`$ by writing $`b=3(1,x,y)`$ and allowing the point $`B(x,y)`$ to vary inside the convex polygon $`P`$: note in fact that $`b`$ is inside the fan. Define the functions:
$$\mathrm{Vol}_{\mathrm{\Sigma }_i}(x,y)=\frac{2\pi ^2}{9}\frac{v_{i1},v_i}{r_{i1},v_{i1}r_i,v_i}\frac{2\pi ^2}{9}l_i(x,y)$$
(4.2)
where $`r_i`$ is the plane vector going from $`B`$ to the vertex $`V_i`$ (see Figure 3).
As shown in , these are the volumes of the base three-cycles associated with the divisors $`D_i,i=1,\mathrm{},d`$. Define also the function:
$$\mathrm{Vol}_H(x,y)=\frac{\pi }{6}\underset{i=1}{\overset{d}{}}\mathrm{Vol}_{\mathrm{\Sigma }_i}(x,y)$$
(4.3)
which determines the total volume of the horizon $`H`$. The two previous equations are just equations (3.25) and (3.26) of . The function to minimize is just $`\mathrm{Vol}_H(x,y)`$ <sup>6</sup><sup>6</sup>6This is the function $`Z`$ in up to a constant multiplicative factor. and the position of the minimum $`(\overline{x},\overline{y})`$ gives the Reeb vector $`\overline{b}=3(1,\overline{x},\overline{y})`$ for the CY cone. It was proved in that such minimum exists and is unique.
The values of $`\mathrm{Vol}_H(\overline{x},\overline{y})`$ and $`\mathrm{Vol}_{\mathrm{\Sigma }_i}(\overline{x},\overline{y})`$ at the minimum are the total volume of $`H`$ and the volumes of $`\mathrm{\Sigma }_i`$ to be compared with the central charge $`a`$ and the R-charges $`a_i`$ of the field theory through the AdS/CFT relations (2.3) and (2.4). To facilitate this comparison we define the geometrical function:
$$a^{MSY}(x,y)=\frac{\pi ^3}{4\mathrm{V}\mathrm{o}\mathrm{l}_H(x,y)}$$
(4.4)
and the functions:
$$f_i(x,y)=\frac{2l_i(x,y)}{_{j=1}^dl_j(x,y)}$$
(4.5)
corresponding to the R-charges $`R_i`$ through equation (2.4). The process of Z-minimization can be restated as a maximization of $`a^{MSY}(x,y)`$ with $`(x,y)`$ varying in the interior of $`P`$.
On the other side of the correspondence we have the gauge theory with trial central charge $`a`$ which is a function of the $`d`$ variables $`a_i`$:
$$a(a_1,a_2,\mathrm{}a_d)=\frac{9}{32}\left(F+\underset{(i,j)C}{}\left|v_i,v_j\right|\left(a_{i+1}+a_{i+2}+\mathrm{}a_j1\right)^3\right)$$
(4.6)
We are considering a formal extension of the trial central charge to $`^d`$ defined by equation (3.10). This function has to be locally maximized with the constraint (3.5) (and $`a_i>0`$). To impose this constraint it is enough to introduce a Lagrange multiplier $`\lambda `$ and to extremize the function:
$$a(a_1,a_2,\mathrm{}a_i)\lambda \left(a_1+a_2\mathrm{}+a_d2\right)$$
(4.7)
By deriving with respect to $`a_i`$ we get the conditions <sup>7</sup><sup>7</sup>7This means that the gradient of the extended function $`a`$ in the local maximum is parallel to the vector $`(1,\mathrm{}1)`$. So to extremize $`a`$ it is enough to impose that the variations of $`a`$ along the $`d1`$ vectors $`S^a`$ orthogonal to $`(1,\mathrm{}1)`$ vanish:
$$\underset{i=1}{\overset{d}{}}S_i^a\frac{a}{a_i}=0,\text{if}\underset{i=1}{\overset{d}{}}S_i^a=0$$
(4.8) But note that, in the language of Section 3, the space of $`S^a`$ is just the space of the $`d1`$ global symmetries (compare with (3.8)).:
$$\frac{a}{a_i}=\lambda i=1,\mathrm{}d$$
(4.9)
If we call $`\overline{a}_i`$ the values of $`a_i`$ at the local maximum, we have to prove that:
$$\begin{array}{c}a^{MSY}(\overline{x},\overline{y})=a(\overline{a}_1,\overline{a}_2,\mathrm{}\overline{a}_d)\hfill \\ f_i(\overline{x},\overline{y})=\overline{a}_ii=1,\mathrm{}d\hfill \end{array}$$
(4.10)
This is a highly non trivial check to perform: a-maximization and Z-minimization use different functions and different trial charges; it is not at all obvious why the result should be the same. First of all a-maximization is done on a total of $`d1`$ independent trial parameters while the volume minimization is done only on two parameters $`(x,y)`$. The trial central charge $`a`$ is a cubic polynomial in $`a_i`$, whereas $`a^{MSY}`$ is a rational function of $`(x,y)`$. These parameters, in both cases, are somehow related to the possible global symmetries: the Reeb vector in the geometry is connected to R-symmetries of the gauge theory and changing the position of $`B`$ in the directions $`x`$ and $`y`$ means adding to the R-symmetry the two flavor global symmetries <sup>8</sup><sup>8</sup>8Recall that flavor symmetries are mixed with baryonic ones, so actually we are moving also in the space of baryonic symmetries.. In any case, the volume minimization is done by moving only in a two dimensional subspace of the set of global symmetries, while a-maximization is done on the entire space. Fortunately, as often claimed in the literature, a-maximization can be always performed on a two dimensional space of parameters related to flavor symmetries. Indeed, on general grounds, one can parametrize the trial R-symmetry as a contribution $`R(X,Y)`$ from the flavor symmetries plus a baryonic part
$$R=R(X,Y)+\underset{a=1}{\overset{d3}{}}h_aB_a$$
(4.11)
and the elimination of the variables $`h_a`$ is simple: imposing that the derivatives of $`\mathrm{tr}R^3`$ with respect to $`h_a`$ vanish, one gets the equations
$$\mathrm{Tr}R^2B_a=0.$$
(4.12)
These conditions read
$$\mathrm{Tr}\left(R(X,Y)+\underset{b=1}{\overset{d3}{}}h_bB_b\right)^2B_a=0$$
(4.13)
which is a linear system of $`d3`$ equations in the $`d3`$ variables $`h_a`$. Linearity in $`h_a`$ follows from the fact that the cubic mixed t’Hooft anomaly for baryonic symmetries is zero: $`\mathrm{Tr}B_a^3=0`$. So one can solve for $`h_a`$ in function of $`X`$, $`Y`$ and substitute into the trial charge (4.11); the central charge $`a`$ is now a rational function only of $`X`$ and $`Y`$. So we have reduced a-maximization to a maximization over a set of two parameters.
In the previous argument, the choice of a basis of flavor and baryonic symmetries in (4.11) was quite arbitrary. In our specific case we can choose a more natural parameterization for the two dimensional space over which to reduce a-maximization. This space is just the space of coordinates $`(x,y)`$ of the plane where $`P`$ lies: consider the map from $`^2`$ to $`^d`$ given by
$$\begin{array}{c}f:(x,y)(a_1,a_2,\mathrm{}a_d)\hfill \\ (x,y)a_i=\frac{2l_i(x,y)}{_{j=1}^dl_j(x,y)}=f_i(x,y)\hfill \end{array}$$
(4.14)
We claim that the local maximum $`(\overline{a}_1,\mathrm{}\overline{a}_d)`$ of the a-maximization is found on the image $`f\left(P\right)`$ of the interior of $`P`$ under this map. In fact it is not difficult to prove that the gradient of the trial central charge along the $`d3`$ baryonic directions evaluated on $`f\left(P\right)`$ is always zero:
$$\underset{i=1}{\overset{d}{}}B_i^a\frac{a}{a_i}_{|a_i=f_i(x,y)}=0$$
(4.15)
where $`B^a`$ is a baryon charge and where the equality holds for every $`(x,y)`$ in the interior of $`P`$. We give the general proof of (4.15) in the Appendix. Note that equation (4.15) is completely equivalent to the condition (4.12) when the trial R-charge is evaluated with $`a_i=f_i(x,y)`$. Therefore we have clarified in which sense the baryonic symmetries decouple from the process of a-maximization.
At this point we have to compare the two functions $`a^{MSY}(x,y)`$ and the field theory trial central charge evaluated on the surface $`f\left(P\right)`$, which are two functions only of $`(x,y)`$. Remarkably one discovers that they are equal even before maximization:
$$a(a_1,\mathrm{}a_d)_{|a_i=f_i(x,y)}=a^{MSY}(x,y)$$
(4.16)
for every $`(x,y)`$ inside the interior of $`P`$. We give a general (still long) analytic proof of (4.16) in the Appendix.
The proof of the equivalence of a-maximization and Z-minimization is now almost finished: we know that $`a^{MSY}(x,y)`$ has a unique maximum $`(\overline{x},\overline{y})`$ inside the polygon $`P`$. In this point we have for the field theory $`a`$:
$$\frac{f_i}{x_h}(\overline{x},\overline{y})\left(\frac{a}{a_i}\right)_{|a_i=f_i(\overline{x},\overline{y})}=\frac{d}{dx_h}a(f_1(x,y),\mathrm{}f_d(x,y))_{|\overline{x},\overline{y}}=\frac{d}{dx_h}a^{MSY}(\overline{x},\overline{y})=0$$
(4.17)
where $`x_h`$, $`h=1,2`$ is $`x`$ or $`y`$. So we see that, in the point $`(\overline{x},\overline{y})`$, also the two vectors:
$$\frac{f_i}{x}(\overline{x},\overline{y}),\frac{f_i}{y}(\overline{x},\overline{y})$$
(4.18)
belonging to the space of global symmetries (since $`_if_i=2`$) are orthogonal to the gradient of $`a(a_1,\mathrm{},a_d)`$. Together with the $`d3`$ baryon symmetries they span the whole $`d1`$ space of global symmetries, thus proving (4.8) in the point $`(\overline{x},\overline{y})`$. Therefore the extremum point for the trial central charge lies on the surface $`f\left(P\right)`$. One should also check that the Hessian matrix is negative definite to prove that this is a local maximum. The agreement between the volumes of $`\mathrm{\Sigma }_i`$ and the total volume with the R-charges $`\overline{a}_i`$ and the central charge in $`(\overline{x},\overline{y})`$ follows immediately from the parametrization (4.14) and from (4.16).
## 5 Examples
In this Section we provide various examples of our proposal using manifolds $`H`$ where it is possible to determine the dual gauge theory explicitly. Needless to say, we find a remarkable agreement.
### 5.1 The $`Y^{p,q}`$ manifolds
The superconformal theory dual to $`AdS_5\times Y^{p,q}`$ has been determined in . The cone $`C\left(Y^{p,q}\right)`$ determines a polytope $`P`$ with vertices
$$(0,0),(1,0)(0,p)(1,p+q)$$
(5.1)
with a $`(p,q)`$ web given by the vectors
$$(0,1),(p,1)(q,1)(pq,1)$$
(5.2)
With a toric diagram with four sides, we expect six different types of fields corresponding to the number of pairs $`(i,j)`$. However, due to the non-abelian isometry of the manifolds, there is an accidental degeneration. Our proposal and the comparison with the known results is reported in Table 1 using the notations of . Recall that as usual $`a_4=2a_1a_2a_3`$. With this assignment we would perform the a-maximization on a three dimensional space of parameters. The enhanced global symmetry allows to reduce the parameter space to a two-dimensional one, as done in . Indeed, in the a-maximization, $`R`$ can mix only with abelian symmetries ; we still have $`d3=1`$ baryonic symmetries, but only one $`U\left(1\right)`$ flavor symmetry since the other is enhanced to $`SU\left(2\right)`$. In any event, without knowing about the $`SU\left(2\right)`$ symmetry we can perform a-maximization on three parameters and discover at the end that $`a_2=a_4`$. In Table 1, four fields are associated with the four edges of the fan. For $`Y^{p,q}`$ we obtain the fields $`Y`$,$`Z`$ and two copies of the fields $`U`$ with the same multiplicity $`p`$: they combine to give the $`SU\left(2\right)`$ doublet $`U_\alpha `$. The remaining two types of fields are associated with the divisors $`D_2+D_3`$ and $`D_3+D_4`$, they have multiplicity $`q`$ and combine to give the doublets $`V_\alpha `$.
| $`(i,j)C`$ | multiplicity | $`U(1)_R`$ | fields |
| --- | --- | --- | --- |
| $`(4,1)`$ | p+q | $`a_1`$ | $`Y`$ |
| $`(1,2)`$ | p | $`a_2`$ | $`U`$ |
| $`(2,3)`$ | p-q | $`a_3`$ | $`Z`$ |
| $`(3,4)`$ | p | $`a_4`$ | $`U`$ |
| $`(1,3)`$ | q | $`a_2+a_3`$ | $`V`$ |
| $`(2,4)`$ | q | $`a_3+a_4`$ | $`V`$ |
Table 1: Charges and multiplicities for $`Y^{p,q}`$.
In the previous assignment $`D_2+D_3`$ has been chosen instead of $`D_4+D_1`$ because $`v_2,v_3>0`$. A similar argument applies to $`D_3+D_4`$. It is also easy to check that all the toric phases of $`Y^{p,q}`$ described in can be obtained in the way discussed at the end of Section 3 (cfr Table 1 in ).
### 5.2 The $`L^{p,q;r,s}`$ manifolds
The superconformal theory dual to $`AdS_5\times L^{p,q;r,s}`$ has been determined in . The cone $`C\left(L^{p,q;r,s}\right)`$ determines a polytope $`P`$ with vertices
$$(0,0),(1,0)(P,s)(k,q)$$
(5.3)
where $`k`$ and $`P`$ are determined through the Diophantine equation
$$rksPq=0$$
(5.4)
Recall that $`p+q=r+s`$. As explained in , we can always choose $`prsq`$ without any loss of generality. The $`(p,q)`$ web is given by the vectors
$$(0,1),(s,1P)(qs,k+P)(q,k)$$
(5.5)
The toric diagram and $`(p,q)`$ web for $`L^{p,q;r,s}`$ are reported in Figure 4.
| $`(i,j)C`$ | multiplicity | $`U(1)_R`$ | fields |
| --- | --- | --- | --- |
| $`(4,1)`$ | q | $`a_1`$ | $`Y`$ |
| $`(1,2)`$ | s | $`a_2`$ | $`\stackrel{~}{W}`$ |
| $`(2,3)`$ | p | $`a_3`$ | $`Z`$ |
| $`(3,4)`$ | r | $`a_4`$ | $`X`$ |
| $`(1,3)`$ | q-s | $`a_2+a_3`$ | $`W`$ |
| $`(2,4)`$ | q-r | $`a_3+a_4`$ | $`\stackrel{~}{X}`$ |
Table 2: Charges and multiplicities for $`L^{p,q;r,s}`$.
The toric diagram has four sides, and we expect six different types of fields corresponding to the number of pairs $`(i,j)`$. In this case the isometry is $`U\left(1\right)^3`$ and we don’t expect any degeneration. Our proposal and the comparison with the known results is reported in Table 2 using the notations of . Recall that as usual $`a_4=2a_1a_2a_3`$. The same analysis was performed in .
### 5.3 The $`X^{p,q}`$ manifolds
It is interesting to check the case of the manifolds $`X^{p,q}`$ discussed in . These correspond to toric cones with five facets which can be blown down to the cones over $`Y^{p,q}`$. The corresponding gauge theory can be determined by an inverse Higgs mechanism . The general assignment of R-charges and the a-maximization has not been performed in the literature except for particular $`p`$ and $`q`$; therefore this model is an interesting laboratory.
The toric diagram is given by (see Figure 6):
$$(1,p),(0,pq+1)(0,pq)(1,0)(2,0)$$
(5.6)
and the $`(p,q)`$ web is given by the vectors $`v_i`$ (see Figure 6):
$$(q+1,1),(1,0)(p+q,1)(0,1)(p,1)$$
(5.7)
With a toric diagram with five sides, we expect ten different types of fields corresponding to the number of pairs $`(i,j)`$. Our proposal is reported in Table 3. Recall that as usual $`_ia_i=2`$ for an R-symmetry. We have four independent parameters because there are now two baryonic symmetries.
| $`(i,j)C`$ | multiplicity | $`U(1)_R`$ |
| --- | --- | --- |
| $`(5,1)`$ | p+q-1 | $`a_1`$ |
| $`(1,2)`$ | 1 | $`a_2`$ |
| $`(2,3)`$ | 1 | $`a_3`$ |
| $`(3,4)`$ | p-q | $`a_4`$ |
| $`(4,5)`$ | p | $`a_5`$ |
| $`(1,3)`$ | p-1 | $`a_2+a_3`$ |
| $`(2,4)`$ | 1 | $`a_3+a_4`$ |
| $`(1,4)`$ | q-1 | $`a_2+a_3+a_4`$ |
| $`(5,2)`$ | 1 | $`a_1+a_2`$ |
| $`(3,5)`$ | q | $`a_4+a_5`$ |
Table 3: Charges and multiplicities for $`X^{p,q}.`$
We can explicitly determine the gauge theory and assign the R-charges to bi-fundamental fields. This can be done more efficiently using the brane tiling description of the $`X^{p,q}`$ theory. We refer to for a detailed discussion of the brane tiling. Here we use the method we employed for the $`L^{p,q;r,s}`$ manifolds in . The tiling for $`X^{p,q}`$ is pictured in Figure 7. Similarly to $`Y^{p,q}`$ theories, the dimer configuration of $`X^{p,q}`$ can be obtained using only one column of $`n`$ hexagons, and $`m+1`$ consecutive cut hexagons. The horizontal identification has a shift $`k=1`$, as for $`Y^{p,q}`$. The main difference is that for $`X^{p,q}`$ the last cut hexagon has a cut in the opposite direction than the other $`m`$ cuts.
To fit the number of fields, gauge groups and superpotential terms for $`X^{p,q}`$ we must choose: $`n=2q1`$, $`m=pq`$. We report also the general form of the Kasteleyn matrix, with vertices numbered in the same way as in (see Figure 7). The determinant of $`K`$ is then:
$$\mathrm{det}K=1+z+z^1w^{n+m+1}+z^1w^{n+m}+w^{\frac{n+1}{2}+m}+\mathrm{}$$
(5.8)
where we have not been careful about signs and the omitted terms are powers of $`w`$ with lower exponent. In the plane $`(z,w)`$ one gets the toric diagram:
$$(0,0),(1,0)(0,p)(1,p+q)(1,p+q1)$$
(5.9)
Translating by $`(0,p)`$, applying the $`SL(2,)`$ transformation $`((1,0),(p,1))`$ and translating by $`(1,0)`$, one recovers the equivalent diagram (5.6). This shows that the dimer configuration reproduces the geometry. By comparing (5.9) with (5.1), it is manifest that the cone $`C\left(X^{p,q}\right)`$ can be obtained by blowing up $`C\left(Y^{p,q}\right)`$. Using the tiling in Figure 7 and the algorithm described in <sup>9</sup><sup>9</sup>9the four independent symmetries are determined by the assignments $`v^1`$, $`v^2`$, $`v^3`$ as in Appendix A.2 of plus an assignment built as a second “cycle” starting from the cut hexagon at position $`m+1`$. An alternative and more general method for determining the distribution of R-charges on the dimer is described in subsection 5.4. we can find the different types of fields and their distribution on the tiling in the general case. The agreement with our proposal given in Table 3 is complete.
### 5.4 Assigning R-charges on the dimer: a general conjecture
In this subsection we propose, and check also on specific examples, a general conjecture to assign R-charges to chiral fields.
In it was suggested a natural one to one correspondence between auxiliary fields in the Witten sigma model associated with a quiver theory and perfect matchings of the dimer configuration. Recall that a perfect matching of a bipartite graph is a choice of links such that every white and black vertex is taken exactly once. We will concentrate on theories for which the multiplicities of the auxiliary fields in the associated Witten sigma model <sup>10</sup><sup>10</sup>10not to be confused with the multiplicity of the “real” fields in the gauge theory associated to the vertex of the toric diagram. corresponding to vertices of the toric diagram are all equal to one. That is we consider dimer configurations with only one perfect matching corresponding to each vertex of the toric diagram. This is always true in all known theories we considered and we think this may be true also in general.
In fact not only there exist many equivalent descriptions (dimer configurations) of the same physical theory, generally connected by Seiberg dualities, but there are also dimer configurations that do not have any AdS/CFT dual. As an example consider the dimers that can be built using only one column of $`n`$ hexagons and $`m`$ (consecutive) cut hexagons as in . In that paper it was pointed out that, using an horizontal identification with shift $`k=1`$, one can obtain the whole family of $`Y^{p,q}`$ theories with the choice $`n=2q`$ and $`m=pq`$. Note that $`n`$ is always even, and the toric diagram of $`Y^{p,q}`$ is:
$$(0,0),(1,0),(0,n/2+m),(1,n+m)\text{for }n\text{ even}$$
(5.10)
For these configurations the number of perfect matchings associated to any vertex $`V_i`$ is always one, as it is easy to prove from the general expression of the Kasteleyn matrix reported in . Moreover these configurations survive the test of the equivalence between a-maximization and Z-minimization.
Instead if we build tilings with an odd number $`N`$ of normal hexagons and $`M`$ cut hexagons (shift again k=1) we get surprising results. The toric diagram is now given by:
$$(0,0),(1,0),(0,\frac{N1}{2}+M),(1,N+M)\text{for }N\text{ odd}$$
(5.11)
as one can see from the Kasteleyn matrix. Note that, up to auxiliary fields multiplicities, we get the same toric configuration if we choose:
$$N=n1M=m+1$$
(5.12)
but with N odd there is a vertex of the toric diagram (precisely $`(0,\left(N1\right)/2+M)`$) having more than one corresponding perfect matching, as one can see again from the Kasteleyn matrix. Moreover it is easy to see that the theories corresponding to configurations $`\left(N,M,k=1\right)`$ with $`N`$ odd do not match the Z-minimization results of the corresponding toric diagrams. These quiver theories do not have a conformal fixed point satisfying the unitary bounds. In fact it is easy to convince oneself that they have only two <sup>11</sup><sup>11</sup>11With $`N`$ odd the cycle described in Appendix A.2 of to build the third charge do not cover all the cut hexagons. global symmetries instead of three $`U\left(1\right)`$ symmetries of $`Y^{p,q}`$ (one of these $`U\left(1\right)`$ is however enlarged to $`SU\left(2\right)`$ for $`Y^{p,q}`$). The trial R-charges associated with some fields (those corresponding to the cuts of the hexagons) are zero and this violates the unitary bound since the corresponding gauge invariant dibaryon operators would have dimension zero. In this way, we have built an infinite family of quiver gauge theories, that can be represented with dimer configurations, but cannot have any geometric AdS/CFT dual. We analyzed some other cases of theories without a geometric dual by varying also $`k`$, and always found that such theories have at least a vertex of the toric diagram with number of perfect matchings associated greater than one. We conjecture that the request of having only a perfect matching corresponding to each vertex of the toric diagram is necessary for the existence of an AdS/CFT dual, but this statement should be further studied. In the following we only consider theories that satisfy such request.
Our conjecture is that it is possible to assign R charges (or global charges) once the perfect matchings corresponding to the vertices of the toric diagram are known <sup>12</sup><sup>12</sup>12We consider here smooth horizons for which the edges of the toric diagram do not pass through integer points, for an extension of this conjecture also to non smooth horizons see subsection 5.6.. The method is simple: assign R-charge (or global charge) $`a_i`$ to the perfect matching corresponding to the vertex $`V_i`$ of the toric diagram. The charges $`a_i`$ satisfy (3.5) if they are R-charges or (3.8) if they are global charges. The (R-)charge of a link in the dimer configuration is then the sum of all (R-)charges of the perfect matchings (corresponding to vertices of the toric diagram) to which the link belongs.
We have checked in many known cases that this method works, also in different toric phases of the same theory. For phases with the minimal number of fields it reproduces our formula for the multiplicities of the different kinds of fields. For example it is not difficult to extract the perfect matchings associated to vertices of the $`X^{p,q}`$ theories from the Kasteleyn matrix reported in the previous subsection. And then one can check that the distribution of R-charges in the dimer obtained with the method proposed is a good distribution, that is one verifies that at every vertex the sum of R-charges is 2 (invariance of the superpotential) and for every face the sum of R-charges is equal to the number of edges minus 2 (beta functions equal to zero). We give other explicit examples of this method in the following subsection.
It would be interesting to check whether this method works in general. Obviously the invariance of the superpotential is guaranteed, since every perfect matching takes every vertex once and the sum of the $`a_i`$ is 2. It would be necessary also to prove the condition for faces (zero beta functions).
In the toric phases with minimal number of fields the method for computing multiplicities of fields described in Section 3 should hold. Every perfect matching is made up with $`V/2`$ links, where $`V`$, the number of vertices in the dimer configuration, is computed in minimal phases from the toric diagram as $`V=EF`$. The method proposed in this subsection implies that there are exactly $`V/2`$ fields containing the charge $`a_1`$, and the same is true for every $`a_i`$, $`i=1,\mathrm{}d`$. Consistence with our formulas for computing multiplicities from the $`(p,q)`$ web requires that the sum of multiplicities of all fields containing $`a_i`$ is equal to $`V/2`$ independently from $`i`$. This is true and is proved in Appendix A.1.
As a final remark, let us remind that in it was discovered that a chiral field (a link in the dimer) in the gauge theory can be written as the product of all auxiliary fields associated to perfect matchings to which the field belongs, and not only to perfect matchings corresponding to vertices of the toric diagram. Hence we have claimed that only the perfect matchings associated to vertices are charged under R or global symmetries, whereas other perfect matchings have charges equal to zero.
### 5.5 The toric del Pezzo 3
In this subsection we consider the example of the theories associated with the complex cone over $`dP_3`$. This toric manifold is interesting since its toric diagram has six edges and four different toric phases are known. All the corresponding quivers are given in . We draw in Figures 9 and 9 the toric diagram and $`(p,q)`$ web for $`dP_3`$; we also show the assignment of charges $`a_i`$ in our conventions. Remember that for R-charges the sum of all $`a_i`$ is equal to 2.
The area of the toric diagram is 3, and therefore the number of gauge groups is $`F=6`$. Model I of $`dP_3`$ has 12 fields $`E=12`$ and hence $`V=6`$ terms in the superpotential. This model has the least number of fields among the toric phases of $`dP_3`$, in agreement with equations (3.2). We draw the dimer configuration for Model I in Figure 10; we label the chiral fields $`X_i`$ with numbers $`i=1,\mathrm{}12`$ typed in blue and vertices with letters $`A,B,\mathrm{}F`$. The identification of faces is as in .
To compute the R-charges of the theory we can use the method suggested in the previous subsection; first we have to know the perfect matchings associated to the vertices. A fast way to compute them is by writing the determinant of a modified Kasteleyn matrix:
$$K=\begin{array}{cccc}& B& D& F\\ & & & \\ A& X_2+wX_5& X_3wzX_6& X_4+zX_1\\ C& X_7& X_9& w^1X_8\\ E& z^1X_{12}& X_{11}& X_{10}\end{array}$$
(5.13)
where we have written for every field not only the usual weight in function of $`w`$ and $`z`$ , but also the name of the field itself. Note that it is not necessary to be careful about signs. The coefficient of $`w^iz^j`$ in the expression of $`\mathrm{det}K`$ gives the perfect matching(s) associated with the point at position $`(i,j)`$ in the plane of the toric diagram. So we find that the perfect matchings associated with the vertices are:
$`a_1`$ $``$ $`X_3X_8X_{12}`$
$`a_2`$ $``$ $`X_4X_9X_{12}`$
$`a_3`$ $``$ $`X_5X_9X_{10}`$
$`a_4`$ $``$ $`X_6X_7X_{10}`$
$`a_5`$ $``$ $`X_1X_7X_{11}`$
$`a_6`$ $``$ $`X_2X_8X_{11}`$
where on the left we have written the R-charge associated with the vertex/perfect matching. We can then compute the R-charges of the fields $`X_i`$ as described in the previous subsection by summing all the charges of the vertex perfect matchings to which a field $`X_i`$ belongs. We thus get the following table for R-charges:
$$\begin{array}{cccccccccccc}X_1& X_2& X_3& X_4& X_5& X_6& X_7& X_8& X_9& X_{10}& X_{11}& X_{12}\\ a_5& a_6& a_1& a_2& a_3& a_4& a_4+a_5& a_1+a_6& a_2+a_3& a_3+a_4& a_5+a_6& a_1+a_2\end{array}$$
We have found five independent trial R-charges (there is relation (3.5) among the $`a_i`$), and indeed it is not difficult to show that they are the correct ones, for example by writing the matrix $`C_{ij}`$ as in Appendix A.2 of .
Note that the multiplicities (equal to 1 for $`dP_3`$) and the kinds of different fields just found are in agreement with the general formula we propose in this paper. So we recognize in Model I the minimal toric phase of $`dP_3`$ for which the formulae proposed in this paper strictly hold.
There are three other phases of $`dP_3`$ with more than 12 fields. We have performed a similar analysis also for these phases, taking the dimer diagrams from . We do not report here all the calculations, but make some useful comments. First of all we have checked that one can use the algorithm described in the previous subsection to determine the R-charges; this is an efficient and fast algorithm.
Model II, III and IV fit in the general analysis at the end of Section 3. Model II and III of $`dP_3`$ have $`F=6`$, $`E=14`$, $`V=8`$. They both have the same distribution of fields: there are all the fields that appeared in Model I with the same R-charges (5.5) plus two other fields: one has R-charge $`a_3+a_4+a_5`$ and the other $`a_1+a_2+a_6`$. Their contribution to the trial central charge $`a`$ cancels:
$$\left(a_3+a_4+a_51\right)^3+\left(a_1+a_2+a_61\right)^3=0$$
(5.15)
because of (3.5). So the trial a charge to maximize is the same as in Model I.<sup>13</sup><sup>13</sup>13Note that there may exist other parametrizations of trial R-charges. For example in Model III, one can find an equivalent distribution interchanging $`a_1`$ and $`a_4`$: this still satisfies the linear constraints from the vanishing of beta functions and conservation of superpotential. The expression of the trial central charge to maximize is different, but the results at the maximum, where $`a_1=a_4`$, are the same. This is due to the high degree of symmetry of the toric diagram of $`dP_3`$. Model IV of $`dP_3`$ has $`F=6`$, $`E=18`$, $`V=12`$. There are all the fields appearing in Model I plus the six fields with R-charge: $`a_3+a_4+a_5`$, $`a_1+a_2+a_6`$, $`a_5+a_6+a_1`$, $`a_2+a_3+a_4`$, $`a_1+a_2+a_3`$, $`a_4+a_5+a_6`$. Again their contribution to the trial central charge cancels.
### 5.6 Orbifolds and singular horizons
In this subsection we deal with the problem of toric cones over non smooth five dimensional horizons; their toric diagram is characterized by the fact that some of its sides pass through integer points: let’s call $`p`$ the total number of such points on the sides. The global symmetries are now $`d+p1`$. So we have to add new variables to the $`a_i`$, $`i=1,\mathrm{}d`$ if we want to find all the global charges. Let’s call the new variables $`b_i`$, $`i=1,\mathrm{}p`$.
For simplicity we shall work on a specific example: a particular realization of $`L^{2,6;2,6}`$ whose toric diagram and $`(p,q)`$ web are drawn in Figures 14 and 14. This example has $`d=4`$ and $`p=4`$. The double area of the toric diagram is $`F=8`$
We have considered two toric phases of $`L^{2,6;2,6}`$. Their dimers are represented in Figures 14 and 14 respectively. In fact it is not difficult to get the gauge theory by partial resolution of $`C^3/\left(_3\times _3\right)`$, by resolving the point of coordinates $`(3,0)`$. The orbifold has 9 gauge groups and its gauge theory is described in . The only way to get a theory with 8 gauge groups is by eliminating (any) one of the links in the dimer of $`C^3/\left(_3\times _3\right)`$. Integrating out the massive fields one gets Model II, Figure 14, which has $`E=22`$, $`V=14`$. Performing a Seiberg duality with respect to the gauge group corresponding to face E in the dimer of Figure 14, one gets Model I for this theory, which has fewer fields: $`E=20`$, $`V=12`$.
We identify Model I with the toric phase with a minimal number of fields for which our formulae should work. In fact it is possible to extend the algorithm described in Section 3 to extract multiplicities from the toric diagram. Now one should assign charge $`a_i`$ to the $`d`$ vertices $`V_i`$ and $`b_j`$ to the $`p`$ integer points along the edges of $`P`$. Then the multiplicities are extracted using all the vectors of the $`(p,q)`$ web as in Figure 14. In our particular example one gets the fields:
$$\begin{array}{c}\begin{array}{ccccccc}a_1,& a_1+b_1,& b_1+a_2,& b_1+a_2+a_3,& b_1+a_2+a_3+b_4,& a_2,& a_2+a_3,\end{array}\hfill \\ \begin{array}{ccccccc}a_2+a_3+b_4,& a_3,& a_3+b_4,& b_4+a_4,& b_4+a_4+b_3,& b_4+a_4+b_3+b_2,& a_4,\end{array}\hfill \\ \begin{array}{cccccc}a_4+b_3,& a_4+b_3+b_2,& b_3+b_2+a_1,& b_3+b_2+a_1+b_1,& b_2+a_1,& b_2+a_1+b_1\end{array}\hfill \end{array}$$
(5.16)
all with multiplicity equal to one (the total number of fields is thus 20, as in Model I). Note that, differently from the case of $`a_i`$, there is no chiral field with charge, say, $`b_1`$, since the $`b_i`$ are always included between parallel vectors (forming a parallelogram with area zero). Indeed it is not difficult to find a distribution of R-charges in the dimer configuration of Model I with these kinds of fields. Remember that the constraints are:
$$\underset{i=1}{\overset{d}{}}a_i+\underset{j=1}{\overset{p}{}}b_j=2$$
(5.17)
if we are dealing with R-charges, and
$$\underset{i=1}{\overset{d}{}}a_i+\underset{j=1}{\overset{p}{}}b_j=0$$
(5.18)
if we are dealing with global charges. The trial R-charge depends both on $`a_i`$ and $`b_j`$, however we have verified in this case that the point that maximizes the central charge has all $`b_i`$ equal to zero. We conjecture that this may be true in general. In practice one could have started with the $`(p,q)`$ web drawn in Figure 15 for $`L^{2,6;2,6}`$; this is simply built ignoring the fact that there are points on the sides of $`P`$: the vectors are not the primitive ones, but they have the same length as the vectors of $`P`$.
Using the usual method for multiplicities as in Section 3 with the $`(p,q)`$ web in Figure 15, we get this table of multiplicities for the 20 fields:
$$\begin{array}{cccccc}\mathrm{R}\mathrm{charge}:\hfill & a_1& a_2& a_3& a_4& a_2+a_3\\ \mathrm{multiplicity}:\hfill & 6& 2& 2& 6& 4\end{array}$$
(5.19)
to which (5.16) obviously reduces after setting $`b_i=0`$. Then the a-maximization can be performed also keeping into account only the charges $`a_i`$ and it is easy to check in this example that it reproduces the volumes of Z-minimization.
Let make also some comments about the generalization of the method described in subsection 5.4 for assigning (R-)charges. The multiplicities of perfect matchings associated with vertices are again equal to one. Then we assign to the corresponding perfect matching (R-)charge $`a_i`$. But in general there is more than one perfect matching corresponding to a certain point along a side of $`P`$. In Model I of the example at hand the multiplicities of perfect matchings corresponding to points $`b_1`$, $`b_2`$, $`b_3`$, $`b_4`$ are respectively 2, 3, 3, 2. Therefore for every point along the sides we can choose a particular perfect matching and give it (R-)charge $`b_i`$ (and zero charge to all other perfect matchings). Then we can compute the charge of chiral fields as sums of charges of the perfect matchings to which they belong, as in subsection 5.4. In this way one always find R-charges (or global charges). However not all charges built in this way are linearly independent: this depends on the choice of perfect matchings. We verified in the case at hand that there are choices of perfect matchings for the $`b_i`$ that allow to find all the 7 independent (R-)charges, some of them also reproducing the fields content given in (5.16).
The same conclusions hold for Model II of $`L^{2,6:2,6}`$. The number of fields now is 22: again a-maximization can be performed by setting to zero the $`b_i`$. We have all the fields appearing in table 5.19 plus one field with charge $`a_1+a_2`$ and one with charge $`a_3+a_4`$, so that the trial R-charge is the same as in Model I for the mechanism described at the end of Section 3.
In conclusion in this subsection we have generalized our results to the case of non smooth horizon, checking in detail the algorithms on a particular example. This analysis deserves further study in order to verify whether it is true in the general case. In particular we guess that charges associated with points along the sides of the toric diagram are never relevant for a-maximization.
## 6 Conclusions
In this paper we computed the central charge and the R-charges of chiral fields for all the superconformal gauge theories living on branes at toric conical singularities. We also showed that the a-maximization technique is completely equivalent to the volume minimization technique proposed in . This, by itself, is an absolutely general check of the AdS/CFT correspondence, valid for all toric singularities.
In this general construction, something is obviously missing. We have now, using the tiling construction , a direct determination of the singularity associated with a given gauge theory. The inverse process is still incomplete: we can determine R-charges and multiplicities of fields but not the specific distribution of bi-fundamentals in the quiver theory. We are quite confident that, in the long period, the dimers technology will allow to define a one-to-one correspondence between CFTs and toric singularities.
It would be also interesting to derive the assignment of charges and multiplicities we propose here. A possible way of deriving it goes though mirror symmetry. It would be interesting to perform the analysis done in in the general case. This analysis would probably teach us also about the many toric phases that are associated with the same superconformal gauge theory.
Acknowledgments
This work is supported in part by by INFN and MURST under contract 2001-025492, and by the European Commission TMR program HPRN-CT-2000-00131.
Appendix
### A.1 A useful formula
Let us define the sets $`C_h`$, $`h=1,2\mathrm{}d`$ which are subsets of $`C`$: a couple $`(i,j)`$ is in $`C_h`$ iff the R-charge of the corresponding chiral field is a sum $`a_{i+1}+a_{i+2}+\mathrm{}a_j`$ containing $`a_h`$. In practice $`C_h`$ is made up of all the couples $`(i,j)`$ such that the region $`a_h`$ in Figure 2 is contained in the angle $`180^o`$ generated by $`v_i`$ and $`v_j`$.
In this Appendix we shall prove the useful formula
$$S_h\underset{(i,j)C_h}{}\left|v_i,v_j\right|=\frac{V}{2}$$
(A.1)
where $`V`$ is defined as:
$$VEF$$
(A.2)
$`V`$ is the number of vertices of the associated dimer configuration. Note that (when the convex polygon $`P`$ has integer coordinates) equation (A.1) proves also that $`V`$ is even. This agrees with the fact that there is an equal number of white and black vertices in the dimer configuration.
Given a vector $`v_j`$ in the $`(p,q)`$ web let us extend it (as in Figure 17 for the case $`j=1`$) and call $`v_{k_j}`$ the vector in the $`(p,q)`$ web just before this extension (moving in counter-clockwise direction). Note that:
$$\begin{array}{c}\left|v_j,v_{j+1}\right|+\left|v_j,v_{j+2}\right|\mathrm{}+\left|v_j,v_{k_j}\right|\left|v_j,v_{j1}\right|\left|v_j,v_{j2}\right|\mathrm{}\left|v_j,v_{k_j+1}\right|\hfill \\ =v_j,v_{j+1}+v_{j+2}+\mathrm{}v_{k_j}+v_{k_j+1}+\mathrm{}v_{j1}\hfill \\ =v_j,v_j=0\hfill \end{array}$$
(A.3)
where we have used that the sum of all $`v_i`$ in the $`(p,q)`$ web is zero. Remember that our indexes are always defined modulo $`d`$. Note that equation (A.3) is just the difference $`S_{j+1}S_j`$, so we have proved that all $`S_j`$ are equal.
To prove (A.1) we can choose $`h=1`$ by a relabeling of vertices and sides (see Figure 17). Let us consider the vector $`v_1`$ and write in the first line of a table all the multiplicities made up with $`v_1`$ (see below). We divide this line into two parts: on the left we write the pairs from $`\left|v_1,v_2\right|`$ to $`\left|v_1,v_{k_1}\right|`$ (those which do not contain $`a_1`$) and on the right the pairs from $`\left|v_1,v_d\right|`$ to $`\left|v_1,v_{k_1+1}\right|`$ (that contain $`a_1`$)<sup>14</sup><sup>14</sup>14If there is a vector $`v_j`$ lying just on the extension of $`v_1`$, the multiplicity $`|v_1,v_j|`$ is zero and so it can be ignored.. We repeat this procedure writing in the second line of the table all the pairs in $`C`$ that contain $`v_2`$, again dividing the line into two parts: on the left the pairs from $`\left|v_2,v_3\right|`$ to $`\left|v_2,v_{k_2}\right|`$ and on the right the pairs from $`\left|v_2,v_1\right|`$ to $`\left|v_2,v_{k_2+1}\right|`$. We continue to fill in the lines with this ordering up to line $`k_1`$; in the remaining lines from $`k_1+1`$ to $`d`$ we reverse the order in which we divide lines in a left and right part. For example line $`k_1+1`$ contains the multiplicities formed with $`v_{k_1+1}`$ and we write on the left the pairs from $`\left|v_{k_1+1},v_{k_1}\right|`$ to $`\left|v_{k_1+1},v_{k_{k_1+1}+1}\right|`$ and on the right the pairs from $`\left|v_{k_1+1},v_{k_1+2}\right|`$ to $`\left|v_{k_1+1},v_{k_{k_1+1}}\right|`$: the idea is that all the pairs on the left do not contain $`a_1`$ whereas the pairs on the right may or may not contain $`a_1`$.
$$\mathtt{}\begin{array}{cccccccc}\left|v_1,v_2\right|\hfill & \left|v_1,v_3\right|\hfill & \mathrm{}\hfill & \left|v_1,v_{k_1}\right|\hfill & \left|v_1,v_d\right|\hfill & \left|v_1,v_{d1}\right|\hfill & \mathrm{}\hfill & \left|v_1,v_{k_1+1}\right|\hfill \\ \left|v_2,v_3\right|\hfill & \left|v_2,v_4\right|\hfill & \mathrm{}\hfill & \left|v_2,v_{k_2}\right|\hfill & \left|v_2,v_1\right|\hfill & \left|v_2,v_d\right|\hfill & \mathrm{}\hfill & \left|v_2,v_{k_2+1}\right|\hfill \\ \text{ }\mathrm{}\hfill & & & & \text{ }\mathrm{}\hfill & & & \\ \left|v_{k_1},v_{k_1+1}\right|\hfill & \left|v_{k_1},v_{k_1+2}\right|\hfill & \mathrm{}\hfill & \left|v_{k_1},v_{k_{k_1}}\right|\hfill & \left|v_{k_1},v_{k_11}\right|\hfill & \left|v_{k_1},v_{k_12}\right|\hfill & \mathrm{}\hfill & \left|v_{k_1},v_{k_{k_1}+1}\right|\hfill \\ \left|v_{k_1+1},v_{k_1}\right|\hfill & \left|v_{k_1+1},v_{k_11}\right|\hfill & \mathrm{}\hfill & \left|v_{k_1+1},v_{\mathrm{}}\right|\hfill & \left|v_{k_1+1},v_{k_1+2}\right|\hfill & \left|v_{k_1+1},v_{k_1+3}\right|\hfill & \mathrm{}\hfill & \left|v_{k_1+1},v_{\mathrm{}}\right|\hfill \\ \text{ }\mathrm{}\hfill & & & & \text{ }\mathrm{}\hfill & & & \\ \left|v_d,v_{d1}\right|\hfill & \left|v_d,v_{d2}\right|\hfill & \mathrm{}\hfill & \left|v_d,v_{k_d+1}\right|\hfill & \left|v_d,v_1\right|\hfill & \left|v_d,v_2\right|\hfill & \mathrm{}\hfill & \left|v_d,v_{k_d}\right|\hfill \end{array}$$
Note that the multiplicity associated with every pair of vectors $`v_i`$, $`v_j`$ in $`C`$ appears twice in the above table: once in line $`i`$ and once in line $`j`$. Hence the total sum of multiplicities in the table is $`2E`$. But the sum of multiplicities in each line on the left equals the sum of multiplicities on the right in the same line because of (A.3). Hence the total sum of multiplicities on the right (or left) side of the table equals $`E`$. Moreover the pairs $`(i,j)`$ of vectors with multiplicity $`\left|v_i,v_j\right|`$ in the left side of this table do not belong to $`C_1`$. All the pairs of vectors in $`C_1`$ appear (twice) in the right side of the table; but on the right side there appear also pairs that do not belong to $`C_1`$: $`\left|v_2,v_1\right|`$ in the second line, $`\left|v_3,v_2\right|`$ and $`\left|v_3,v_1\right|`$ in the second line, and so on. The total sum of such multiplicities is:
$$\begin{array}{c}v_1,v_2+v_1+v_2,v_3\mathrm{}+v_1+v_2\mathrm{}+v_{k_11},v_{k_1}+\hfill \\ +v_{k_1+1},v_{k_1+2}+v_{k_1+3}\mathrm{}+v_d\mathrm{}+v_{d1},v_d=\hfill \\ =v_1,v_2+v_1+v_2,v_3\mathrm{}+v_1+v_2\mathrm{}+v_{k_11},v_{k_1}+\hfill \\ v_1+v_2\mathrm{}+v_{k_1},v_{k_1+1}+v_1+v_2\mathrm{}+v_{d2},v_{d1}=2\mathrm{A}\mathrm{r}\mathrm{e}\mathrm{a}\left(P\right)=F\hfill \end{array}$$
(A.4)
where in the first equality we have used that the sum of all $`v_i`$ is zero and the bilinearity and antisymmetry of the determinant. The sum of all multiplicities in the right side of the table above that do not belong to $`C_1`$ is thus equal to the double area of $`P`$, see Figure 17. The sum we had to compute is therefore:
$$S_1\underset{(i,j)C_1}{}\left|v_i,v_j\right|=\frac{EF}{2}=\frac{V}{2}$$
(A.5)
which is relation (A.1).
### A.2 Charges
We now show that our proposed formula for extracting multiplicities of chiral fields from the toric diagram correctly gives $`U\left(1\right)`$ baryon, flavor and R-charges with trace equal to zero. Let’s start with a charge commuting with supersymmetry; as explained in Section 3 it can be built by assigning charges $`a_i`$ to chiral fields associated with vectors $`V_i`$ of the fan with (3.8):
$$\underset{i=1}{\overset{d}{}}a_i=0$$
(A.6)
Therefore we have $`d1`$ global symmetries, 2 of which are flavor symmetries and the remaining $`d3`$ are baryonic symmetries (remember that for non smooth horizons we have to consider also charges associated to integer points lying along the sides of the convex polygon $`P`$; the total sum of all charges associated to “fundamental” fields (A.6) must still be zero). The charge of a generic “composite” chiral field associated with the pair $`(i,j)C`$ is simply the sum $`a_{i+1}+\mathrm{}a_j`$.
The trace of a generic $`U\left(1\right)`$ global symmetry is thus:
$$\begin{array}{cc}\mathrm{tr}U\left(1\right)\hfill & =\underset{(i,j)C}{}\left|v_i,v_j\right|\left(a_{i+1}+a_{i+2}\mathrm{}+a_j\right)\hfill \\ & =\underset{h=1}{\overset{d}{}}a_h\underset{(i,j)C_h}{}\left|v_i,v_j\right|=\frac{V}{2}\underset{h=1}{\overset{d}{}}a_h=0\hfill \end{array}$$
(A.7)
where we have used that $`S_h`$ in (A.1) does not depend on $`h`$.
Let us now turn to R-symmetry; to build the generic trial R-symmetry we have to associate a R-charge $`a_i`$ to the chiral fields corresponding to divisors $`V_i`$ (and also to fields corresponding to vertices along sides for non smooth horizons); the only difference with the global case is that now the sum must satisfy (3.5):
$$\underset{i=1}{\overset{d}{}}a_i=2$$
(A.8)
The trace of a generic $`U\left(1\right)_R`$ symmetry is now
$$\begin{array}{cc}\mathrm{tr}U\left(1\right)_R\hfill & =F+\underset{(i,j)C}{}\left|v_i,v_j\right|\left(a_{i+1}+a_{i+2}\mathrm{}+a_j1\right)\hfill \\ & =F+\underset{h=1}{\overset{d}{}}a_h\underset{(i,j)C_h}{}\left|v_i,v_j\right|\underset{(i,j)C}{}\left|v_i,v_j\right|\hfill \\ & =F+\frac{V}{2}\left(\underset{h=1}{\overset{d}{}}a_h\right)E=0\hfill \end{array}$$
(A.9)
where we have used equation (A.1). The term $`F=2\mathrm{Area}\left(P\right)`$ comes from gauginos, since we know that the double area gives the number of gauge groups. This also shows that for gauge theories dual to toric geometries the trial R-charge always reduces to $`a=9/32\mathrm{tr}R^3`$.
Let us now prove that the trace of cubic t’Hooft anomaly and mixed cubic anomaly for baryonic symmetries are always zero with the multiplicities and charges for chiral fields that we have conjectured in this paper. The vanishing of such anomalies is required by the AdS/CFT correspondence and is always true for the quiver gauge theories under consideration since the global baryonic symmetries are (the non anomalous) linear combinations of the $`U\left(1\right)`$ part of the original gauge groups $`U\left(N\right)`$ (after the AdS/CFT limit they generally become $`SU\left(N\right)`$ gauge groups). But since we have only conjectured the multiplicities of chiral fields and a full algorithm for extracting the whole gauge theory from toric geometry is still lacking, the proof of zero cubic anomaly for baryonic symmetries is a non trivial check of our conjecture.
First of all recall that, as discovered in , the $`d1`$ baryonic symmetries are simply the linear relations between the $`d`$ generators of the toric fan $`V_i`$: <sup>15</sup><sup>15</sup>15Again recall that for non smooth horizons one has to add to the set of $`V_i`$ all the vectors in the fan arriving at the integer points along the sides of $`P`$. if $`(a_1,a_2,\mathrm{}a_d)`$ are the charges of a baryonic symmetry associated with chiral fields corresponding to the vectors $`V_i`$ we have equation (3.4)
$$\underset{i=1}{\overset{d}{}}a_iV_i=0$$
(A.10)
Knowing that $`V_i`$ have first coordinate equal to 1, and that the other two components are the coordinates $`(x_i,y_i)`$ of the vertices of $`P`$ in the plane, the previous equation can also be restated by saying that $`(a_1,a_2,\mathrm{}a_d)`$ must satisfy (A.6), as all global symmetries, and moreover the constraint:
$$a_2v_1+a_3\left(v_1+v_2\right)\mathrm{}+a_d\left(v_1+v_2+\mathrm{}v_{d1}\right)=0$$
(A.11)
where we have started to compute the coordinates of the vertices of $`P`$ from the first vertex (see Figure 17), but one could have started from any other point in the plane of $`P`$ because of (A.6). Note also that a basis for the two flavor symmetries orthogonal to the baryonic ones is given by the $`x`$ and $`y`$ coordinates of the vertices of $`P`$ in the plane containing $`P`$ referred to the barycenter of $`P`$, so that (A.6) holds.
So take now three different (or equal) baryonic symmetries: $`(a_1,a_2,\mathrm{}a_d)`$, $`(a_1^{},a_2^{},`$ $`\mathrm{}a_d^{})`$ and $`(b_1,b_2,\mathrm{}b_d)`$ all satisfying (A.6) and (A.11). To avoid writing too long formulae we will consider first the case when two symmetries are equal, say $`a_i=a_i^{}`$, and then we will extend our results to the general case. The mixed cubic t’Hooft anomaly with our formula for multiplicities becomes:
$`\mathrm{tr}\left(U\left(1\right)_B^a\right)^2U\left(1\right)_B^b`$ $`=`$ $`{\displaystyle \underset{(i,j)C}{}}\left|v_i,v_j\right|\left(a_{i+1}+a_{i+2}\mathrm{}+a_j\right)^2\left(b_{i+1}+b_{i+2}\mathrm{}b_j\right)`$ (A.12)
$`=`$ $`{\displaystyle \underset{h=1}{\overset{d}{}}}b_h\left({\displaystyle \underset{(i,j)C_h}{}}\left|v_i,v_j\right|\left(a_{i+1}+a_{i+2}\mathrm{}+a_j\right)^2\right)`$
$``$ $`{\displaystyle \underset{h=1}{\overset{d}{}}}b_hc_h`$
where the coefficients $`c_h`$ are defined by the last equality. We have to prove that the vector formed by $`c_h`$ is orthogonal to a generic baryonic symmetry, that is that the vector of $`c_h`$ is a linear combination of $`x`$ and $`y`$ coordinates of vertices of $`P`$ up to some multiple of $`(1,\mathrm{},1)`$. So let’s compute the differences:
$`c_{j+1}c_j=`$
$`=\left|v_j,v_{j+1}\right|\left(a_{j+1}\right)^2+\left|v_j,v_{j+2}\right|\left(a_{j+1}+a_{j+2}\right)^2\mathrm{}+\left|v_j,v_{k_j}\right|\left(a_{j+1}+a_{j+2}\mathrm{}+a_{k_j}\right)^2`$
$`\left|v_j,v_{j1}\right|\left(a_j\right)^2\left|v_j,v_{j2}\right|\left(a_j+a_{j1}\right)^2\mathrm{}\left|v_j,v_{k_j+1}\right|\left(a_j+a_{j1}\mathrm{}+a_{k_j+2}\right)^2`$
where there survive only the sum over pairs that contain $`a_{j+1}`$ and do not contain $`a_j`$, minus the sum over pairs that contain $`a_j`$ and do not contain $`a_{j+1}`$, since all other pairs cancel. The symbols $`k_j`$ are defined as in Appendix A.1. The previous equation can be rewritten as
$$c_{j+1}c_j=v_j,T_j$$
(A.13)
where $`T_j`$ is the vector:
$`T_j`$ $`=`$ $`v_{j+1}\left(a_{j+1}\right)^2+v_{j+2}\left(a_{j+1}+a_{j+2}\right)^2\mathrm{}+v_{k_j}\left(a_{j+1}+a_{j+2}\mathrm{}+a_{k_j}\right)^2`$ (A.14)
$`+v_{j1}\left(a_j\right)^2+v_{j2}\left(a_j+a_{j1}\right)^2\mathrm{}+v_{k_j+1}\left(a_j+a_{j1}\mathrm{}+a_{k_j+2}\right)^2`$
$`=`$ $`v_{j+1}\left(a_{j+1}\right)^2+v_{j+2}\left(a_{j+1}+a_{j+2}\right)^2\mathrm{}+v_{k_j}\left(a_{j+1}+a_{j+2}\mathrm{}+a_{k_j}\right)^2`$
$`+v_{k_j+1}\left(a_{j+1}+a_{j+2}\mathrm{}a_{k_j+1}\right)^2\mathrm{}+v_{j1}\left(a_{j+1}+a_{j+2}\mathrm{}+a_{j1}\right)^2`$
where in the last line we have reordered the sum and used that the sum of all $`a_i`$ is zero (A.6).
Now we want to show that all vectors $`T_j`$ are equal: $`T_1=T_2\mathrm{}=T_dT`$; it is enough to prove that consecutive vectors $`T_j`$ are equal and, by a relabeling of vectors and vertices, it is enough to prove this for, say $`T_1`$ and $`T_2`$. A straightforward computation then yields:
$`T_2T_1=`$ (A.15)
$`=`$ $`v_3\left(a_3\right)^2+v_4\left(a_3+a_4\right)^2\mathrm{}+v_d\left(a_3+a_4\mathrm{}+a_d\right)^2+v_1\left(a_3+a_4\mathrm{}+a_d+a_1\right)^2`$
$`v_2\left(a_2\right)^2v_3\left(a_2+a_3\right)^2v_4\left(a_2+a_3+a_4\right)^2\mathrm{}v_d\left(a_2+a_3\mathrm{}+a_d\right)^2`$
$`=`$ $`a_2^2\left(v_2+v_3\mathrm{}+v_d\right)2a_2\left[v_3a_3+v_4\left(a_3+a_4\right)\mathrm{}+v_d\left(a_3+a_4\mathrm{}a_d\right)\right]+v_1\left(a_2\right)^2`$
$`=`$ $`2a_2^2v_12a_2\left[v_3a_3+v_4\left(a_3+a_4\right)\mathrm{}+v_d\left(a_3+a_4\mathrm{}a_d\right)\right]`$
$`=`$ $`2a_2\left[v_3a_3+v_4\left(a_3+a_4\right)\mathrm{}+v_d\left(a_3+a_4\mathrm{}a_d\right)+v_1\left(a_3+a_4\mathrm{}+a_d+a_1\right)\right]`$
$`=`$ $`2a_2\left[a_1v_1+a_d\left(v_1+v_d\right)+a_{d1}\left(v_1+v_d+v_{d1}\right)\mathrm{}+a_3\left(v_1+v_d+v_{d1}\mathrm{}+v_3\right)\right]`$
$`=`$ $`0`$
where we have used (A.6) and that the sum of $`v_i`$ is zero. In the last step we have used that $`(a_1,\mathrm{}a_d)`$ is a baryonic symmetry, since the last sum is one of the kind of (A.11), centered in the second vertex of the polygon $`P`$.
Now we get for the differences:
$`c_2c_1`$ $`=`$ $`v_1,T`$
$`c_3c_1`$ $`=`$ $`\left(c_3c_2\right)+\left(c_2c_1\right)=v_1+v_2,T`$
$`\mathrm{}`$
$`c_dc_1`$ $`=`$ $`v_1+v_2\mathrm{}+v_{d1},T`$ (A.16)
and for the cubic t’Hooft anomaly of baryonic symmetries:
$`\mathrm{tr}\left(U\left(1\right)_B^a\right)^2U\left(1\right)_B^b={\displaystyle \underset{h=1}{\overset{d}{}}}b_hc_h`$
$`=c_1\left(b_1+b_2+b_3\mathrm{}+b_d\right)+b_2\left(c_2c_1\right)+b_3\left(c_3c_1\right)\mathrm{}+b_d\left(c_dc_1\right)`$
$`=b_2v_1+b_3\left(v_1+v_2\right)\mathrm{}+b_d\left(v_1+v_2\mathrm{}+v_{d1}\right),T`$
$`=0`$ (A.17)
where we have used that $`(b_1,\mathrm{}b_d)`$ is a baryonic symmetry thus satisfying (A.6) and (A.11). It is easy to generalize to the case $`a_ia_i^{}`$: the coefficients $`c_h`$ are given now by:
$$c_h=\underset{(i,j)C_h}{}\left|v_i,v_j\right|\left(a_{i+1}+a_{i+2}\mathrm{}+a_j\right)\left(a_{i+1}^{}+a_{i+2}^{}\mathrm{}+a_j^{}\right)$$
(A.18)
and one has to repeat all the steps leading to (A.15) keeping products of sums of $`a_i`$ and $`a_i^{}`$ instead of squares. It is to see that now (A.15) reads:
$`T_2T_1=`$ (A.19)
$`=`$ $`a_2\left[a_1^{}v_1+a_d^{}\left(v_1+v_d\right)+a_{d1}^{}\left(v_1+v_d+v_{d1}\right)\mathrm{}+a_3^{}\left(v_1+v_d+v_{d1}\mathrm{}+v_3\right)\right]`$
$`a_2^{}\left[a_1v_1+a_d\left(v_1+v_d\right)+a_{d1}\left(v_1+v_d+v_{d1}\right)\mathrm{}+a_3\left(v_1+v_d+v_{d1}\mathrm{}+v_3\right)\right]`$
$`=`$ $`0`$
so that one has to use that both $`a_i`$ and $`a_i^{}`$ are baryonic. The proof then proceeds as before (A.17). This concludes our proof for the cubic anomaly of baryonic symmetries:
$$\mathrm{tr}U\left(1\right)_B^aU\left(1\right)_B^a^{}U\left(1\right)_B^b=0.$$
(A.20)
### A.3 Decoupling baryon charges in a-maximization
In this Appendix we shall prove equation (4.15):
$$\underset{h=1}{\overset{d}{}}b_h\frac{a}{a_h}_{|a_i=f_i(x,y)}=0$$
(A.21)
for every baryonic symmetry with charges $`b_i`$ for the chiral fields associated to $`V_i`$. The functions $`f_i(x,y)`$ and $`l_i(x,y)`$ are defined as in (4.5) and (4.2):
$$f_i=\frac{2l_i}{S},l_i=\frac{v_{i1},v_i}{A_{i1}A_i}$$
(A.22)
where we have defined the sum
$$S\underset{i=1}{\overset{d}{}}l_i$$
(A.23)
and the double area of triangles in Figure 3:
$$A_ir_i,v_i=r_{i+1},v_i.$$
(A.24)
Remember that $`r_{i+1}r_i=v_i`$. Note that $`l_i`$ is positive inside the interior of $`P`$ and diverges on the edges $`v_i`$ and $`v_{i1}`$.
We will need some useful relations among these quantities. In particular we can prove the vectorial identity:
$$l_ir_i=\frac{v_{i1}}{A_{i1}}\frac{v_i}{A_i}.$$
(A.25)
In fact a straightforward computation gives
$$l_ir_i\left(\frac{v_{i1}}{A_{i1}}\frac{v_i}{A_i}\right)=\frac{v_{i1},v_ir_iv_{i1}r_i,v_i+v_ir_i,v_{i1}}{A_{i1}A_i}\frac{N}{A_{i1}A_i}$$
(A.26)
Then we get for the numerator $`N,v_i=N,v_{i1}=0`$. So $`N`$ has to be parallel both to $`v_i`$ and $`v_{i1}`$ which are two linearly independent vectors. Therefore $`N=0`$ and we have proved (A.25).
By summing up (A.25) we get another important property:
$$\underset{i=1}{\overset{d}{}}l_ir_i=0$$
(A.27)
which says that the $`l_i`$ are (proportional to) the weights that should be put on the vertices $`V_i`$ of $`P`$ to keep it in equilibrium if we want to suspend it by the internal point $`B`$.
Equation (A.21) then reads
$$\underset{h=1}{\overset{d}{}}b_h\frac{a}{a_h}_{|a_i=f_i(x,y)}=\frac{27}{32}\underset{h=1}{\overset{d}{}}b_hd_h$$
(A.28)
where we have defined
$$d_h=\left(\underset{(i,j)C_h}{}\left|v_i,v_j\right|\left(a_{i+1}+a_{i+2}\mathrm{}+a_j1\right)^2\right)_{|a_i=f_i(x,y)}$$
(A.29)
that comes from deriving (4.6) with respect to $`a_h`$. Note that (A.28) is, up to a constant factor, equal to $`\mathrm{tr}R^2b`$, with $`R`$ the trial R symmetry and $`b`$ the baryon charge. Note in fact the similarities with equation (A.12): the main difference here being that we are dealing with R-symmetry, so the constraint on $`a_i`$ is (3.5), automatically implemented by the substitution (A.22).
Again the idea is to compute the differences $`d_{j+1}d_j`$ and, similarly to (A.13), to rewrite them as
$$d_{j+1}d_j=v_j,\stackrel{~}{W}_j$$
(A.30)
where now the vector $`\stackrel{~}{W}_j`$ reads:
$$\begin{array}{c}\stackrel{~}{W}_j=[v_{j+1}(a_{j+1}1)^2+v_{j+2}(a_{j+1}+a_{j+2}1)^2\mathrm{}+v_{k_j}(a_{j+1}+a_{j+2}\mathrm{}+a_{k_j}1)^2\hfill \\ +v_{j1}(a_j1)^2+v_{j2}(a_j+a_{j1}1)^2\mathrm{}+v_{k_j+1}(a_j+a_{j1}\mathrm{}+a_{k_j+2}1)^2]_{|a_i=f_i}\hfill \end{array}$$
where the symbols $`k_j`$ are defined as in Appendix A.1. Performing the substitution (A.22) $`a_i=f_i`$ and taking the common denominator we get
$`S^2\stackrel{~}{W}_j=`$ (A.31)
$`=`$ $`v_{j+1}\left(l_{j+1}l_{j+2}l_{j+3}\mathrm{}l_j\right)^2+v_{j+2}\left(l_{j+1}+l_{j+2}l_{j+3}\mathrm{}l_j\right)^2+\mathrm{}`$
$`+v_{k_j}\left(l_{j+1}+l_{j+2}\mathrm{}+l_{k_j}l_{k_j+1}\mathrm{}l_j\right)^2+`$
$`+v_{j1}\left(l_jl_{j1}l_{j2}\mathrm{}l_{j+1}\right)^2+v_{j2}\left(l_j+l_{j1}l_{j2}\mathrm{}l_{j+1}\right)^2+\mathrm{}`$
$`+v_{k_j+1}\left(l_j+l_{j1}\mathrm{}+l_{k_j+2}l_{k_j+1}\mathrm{}l_{j+1}\right)^2`$
$`=`$ $`v_{j+1}\left(l_{j+1}l_{j+2}l_{j+3}\mathrm{}l_j\right)^2+v_{j+2}\left(l_{j+1}+l_{j+2}l_{j+3}\mathrm{}l_j\right)^2+\mathrm{}`$
$`+v_{k_j}\left(l_{j+1}+l_{j+2}\mathrm{}+l_{k_j}l_{k_j+1}\mathrm{}l_j\right)^2+`$
$`+v_{k_j+1}\left(l_{j+1}+l_{j+2}\mathrm{}+l_{k_j+1}l_{k_j+2}\mathrm{}l_j\right)^2+\mathrm{}`$
$`+v_{j1}\left(l_{j+1}+l_{j+2}\mathrm{}+l_{j1}l_j\right)^2`$
where in the last step we have reordered the sum. For later convenience, let us add to $`\stackrel{~}{W}_j`$ two terms proportional to $`v_j`$ defining the new vector $`W_j`$ as:
$`S^2W_j=`$ (A.33)
$`=`$ $`v_{j+1}\left(l_{j+1}l_{j+2}l_{j+3}\mathrm{}l_j\right)^2+v_{j+2}\left(l_{j+1}+l_{j+2}l_{j+3}\mathrm{}l_j\right)^2+\mathrm{}`$
$`+v_{j1}\left(l_{j+1}+l_{j+2}\mathrm{}+l_{j1}l_j\right)^2+`$
$`+v_j\left(l_{j+1}+l_{j+2}\mathrm{}+l_{j1}+l_j\right)^28S{\displaystyle \frac{v_j}{A_j}}`$
and because of antisymmetry of the determinant we still have:
$$d_{j+1}d_j=v_j,W_j$$
(A.34)
We want to prove that all $`W_j`$ are equal: $`W_1=W_2\mathrm{}=W_dW`$. As in the previous Appendix, it is enough to show the equality of consecutive $`W_j`$, $`W_{j+1}`$, and, up to a relabeling of indexes, it is enough to show that $`W_2=W_1`$. So let’s compute the difference:
$`S^2\left(W_2W_1\right)=`$ (A.35)
$`=`$ $`v_3\left(l_3l_4l_5\mathrm{}l_1l_2\right)^2+v_4\left(l_3+l_4l_5\mathrm{}l_1l_2\right)^2+\mathrm{}`$
$`+v_1\left(l_3+l_4+l_5\mathrm{}+l_1l_2\right)^2+v_2\left(l_3+l_4+l_5\mathrm{}+l_1+l_2\right)^2`$
$`v_2\left(l_2l_3l_4l_5\mathrm{}l_1\right)^2v_3\left(l_2+l_3l_4l_5\mathrm{}l_1\right)^2`$
$`v_4\left(l_2+l_3+l_4l_5\mathrm{}l_1\right)^2\mathrm{}v_1\left(l_2+l_3+\mathrm{}+l_d+l_1\right)^2`$
$`8S\left({\displaystyle \frac{v_2}{A_2}}{\displaystyle \frac{v_1}{A_1}}\right)`$
$`=`$ $`4l_2[v_2(l_3+l_4+l_5\mathrm{}+l_1)+v_3(l_3+l_4+l_5\mathrm{}+l_1)`$
$`+v_4(l_3l_4+l_5\mathrm{}+l_1)\mathrm{}+v_1(l_3l_4l_5\mathrm{}l_1)]`$
$`+8S\left({\displaystyle \frac{v_1}{A_1}}{\displaystyle \frac{v_2}{A_2}}\right)`$
where in the last step we have computed the differences between factors with the same $`v_i`$ keeping in consideration that each time only the term $`l_2`$ changes relative sign. Now we reorder the first term in the square bracket and we use equation (A.25) (with $`i=2`$) for the last term:
$`S^2\left(W_2W_1\right)=`$ (A.36)
$`=`$ $`4l_2[l_3(v_2v_3v_4\mathrm{}v_1)+l_4(v_2+v_3v_4\mathrm{}v_1)+\mathrm{}`$
$`+l_1(v_2+v_3+v_4\mathrm{}+v_dv_1)]+8Sl_2r_2`$
$`=`$ $`8l_2\left[l_3v_2+l_4\left(v_2+v_3\right)\mathrm{}+l_1\left(v_2+v_3\mathrm{}+v_d\right)\right]+8Sl_2r_2`$
$`=`$ $`8l_2\left[l_2r_2+l_3\left(r_2+v_2\right)+l_4\left(r_2+v_2+v_3\right)\mathrm{}+l_1\left(r_2+v_2+v_3\mathrm{}+v_d\right)\right]`$
$`8l_2r_2\left({\displaystyle \underset{j=1}{\overset{d}{}}}l_j\right)+8Sl_2r_2`$
where in the second equality we have used that $`_iv_i=0`$, and in the third equality we have added and subtracted the same term. Now the last two terms cancel and, noting that $`r_2+v_2+v_3\mathrm{}v_{i1}=r_i`$ (look at Figure 3) the sum in the square brackets becomes:
$$S^2\left(W_2W_1\right)=8l_2\left(\underset{j=1}{\overset{d}{}}l_jr_j\right)=0$$
(A.37)
where we have used (A.27). Hence we conclude that $`W_1=W_2\mathrm{}=W_dW`$. Now the conclusion of the proof of (A.21), that is $`_hb_hd_h=0`$, proceeds as in (A.17) (with the appropriate substitutions $`TW`$, $`c_hd_h`$). In this step we use that $`b_i`$ are baryonic. This concludes our proof.
### A.4 The equality of $`a`$ and $`a^{MSY}`$
In this Appendix we give a general proof of equation (4.16), that shows the agreement of the central charge $`a`$ and the total volume even before maximization, once the substitution $`a_i=f_i2l_i/S`$ has been performed.
Taking into consideration that $`a=9/32\mathrm{t}\mathrm{r}R^3`$, the definition of $`a^{MSY}`$ in (4.4) and equations (3.3), (4.2), what we have to prove is:
$$\mathrm{tr}R_{\mathtt{}|a_i=f_i}^3=\frac{24}{S}$$
(A.38)
where $`S`$ is the sum of $`l_i`$, as in the previous Appendix (A.23).
In this Appendix we will use the notation $`b=(x,y)`$ to indicate the point $`B`$ in the plane of $`P`$ (recall that the Reeb vector can be parametrized as $`3(1,x,y)`$). With a little abuse of notation, we will call $`V_i`$ the coordinates $`(x_i,y_i)`$ in the plane of $`P`$ of the vertices $`V_i`$. Hence we have $`v_i=V_{i+1}V_i`$ and $`r_i=V_ib`$.
To simplify the calculation of $`\mathrm{tr}R^3`$, choose a point $`(x^0,y^0)`$ in $`P`$, in general distinct from the “Reeb point” $`b=(x,y)`$. For every field in the quiver gauge theory (in the minimal toric phase described in Section 3) associated with the pair $`(i,j)C`$ consider its R-charge:
$$a_{i,j}a_{i+1}+a_{i+2}\mathrm{}+a_j$$
(A.39)
and perform the substitution $`a_i=f_i(x,y)`$; we get a rational function of $`(x,y)`$. Perform the Taylor expansion of this function around the point $`(x^0,y^0)`$ and denote with $`\stackrel{~}{a}_{i,j}`$ the truncation of this expansion up to linear terms in $`(x,y)`$:
$`a_{i,j}(x,y)`$ $`=`$ $`a_{i,j}(x_0,y_0)+\left(x_hx_h^0\right){\displaystyle \frac{}{x_h}}a_{i,j}(x_0,y_0)+O\left(\left(x_hx_h^0\right)^2\right)`$ (A.40)
$``$ $`\stackrel{~}{a}_{i,j}(x,y)+O\left(\left(x_hx_h^0\right)^2\right)`$
where $`x_h`$, $`h=1,2`$, is $`x`$ or $`y`$.
We will use the fact that
$$\mathrm{tr}R^3=\mathrm{tr}R^2\stackrel{~}{R}$$
(A.41)
where $`\stackrel{~}{R}`$ stands for the vector of truncated R-charges $`\stackrel{~}{a}_{i,j}`$. In this formula and in the following we always understand the substitutions $`a_i=f_i(x,y)`$. To prove (A.41) note that, by multiplying by $`2/S`$ equation (A.27), we get:
$$\underset{i=1}{\overset{d}{}}a_ir_i=0,\underset{i=1}{\overset{d}{}}a_iV_i=2b$$
(A.42)
since $`r_i=V_ib`$. Note that this is just equation (2.86) in . Deriving the last relation with respect to $`x`$ and/or $`y`$ we get:
$$\underset{i=1}{\overset{d}{}}\left(\frac{}{x_h}\right)^ka_iV_i=0,\text{if}k2$$
(A.43)
where the derivatives can be mixed in $`x`$, $`y`$ and have total degree $`k2`$. In fact $`b=(x,y)`$ is linear in $`(x,y)`$. Deriving instead the relation $`_ia_i=2`$ we get
$$\underset{i=1}{\overset{d}{}}\left(\frac{}{x_h}\right)^ka_i=0$$
(A.44)
The two previous relations tell us that the derivatives of $`a_i`$ with degree 2 or higher, calculated in any point $`(x^0,y^0)`$, are baryonic symmetries: see equations (3.4) and (3.9). In the previous Appendix we proved that for any baryonic symmetry $`\mathrm{tr}R^2B=0`$ for $`a_i=f_i(x,y)`$. Hence we have
$$\mathrm{tr}R^3=\mathrm{tr}R^2\left(\stackrel{~}{R}+\text{higher derivatives}\right)=\mathrm{tr}R^2\stackrel{~}{R}$$
(A.45)
since the other terms in the Taylor expansion are derivatives with degree $`k2`$.
In the following we will choose $`(x_0,y_0)`$ as the first vertex $`V_1`$ of $`P`$ and we will calculate $`\stackrel{~}{a}_{i,j}(x,y)`$ in the point $`b=(x,y)`$. So we need to get the explicit expressions for the charges
$$\stackrel{~}{a}_i(x,y)=a_i\left(V_1\right)r_1\stackrel{}{}a_i\left(V_1\right),\stackrel{}{}a_i=(\frac{a_i}{x},\frac{a_i}{y})$$
(A.46)
since $`r_1=V_1b=(x_0x,y_0y)`$ and in the second term we have written the scalar product of this vector with the gradient of $`a_i`$. The charges of composite fields are obviously given by $`\stackrel{~}{a}_{i,j}=\stackrel{~}{a}_{i+1}\mathrm{}+\stackrel{~}{a}_j`$.
Let us study first the behavior of $`a_i(x,y)=2l_i/S`$ when $`(x,y)=V_1+tv_1`$ approaches the point $`V_1`$ along the first side of $`P`$, $`0<t<1`$. Note that $`A_1`$ goes to zero, whereas the other areas $`A_i`$ are strictly positive. Hence $`l_1`$ and $`l_2`$ goes to $`+\mathrm{}`$ and the other $`l_i`$ remain finite. Hence all $`a_i(x,y)=2l_i/S`$ different from $`a_1`$ and $`a_2`$ goes to zero, since they have a finite numerator and are divided by $`S`$ which diverges. Performing the limit $`(x,y)V_1+tv_1`$ for $`a_1`$ and $`a_2`$ we get:
$`a_1\left(V_1+tv_1\right)`$ $`=`$ $`{\displaystyle \frac{2v_d,v_1A_2}{A_2v_d,v_1+A_dv_1,v_2}}_{|(x,y)=V_1+tv_1}=2\left(1t\right)`$
$`a_2\left(V_1+tv_1\right)`$ $`=`$ $`{\displaystyle \frac{2v_1,v_2A_d}{A_2v_d,v_1+A_dv_1,v_2}}_{|(x,y)=V_1+tv_1}=2t`$ (A.47)
where we used $`A_d=v_d,tv_1`$ and $`A_2=\left(1t\right)v_1,v_2`$ when $`(x,y)=V_1+tv_1`$. Note in particular that for $`t=0`$, we obtain for the vertex $`V_1`$: $`a_1\left(V_1\right)=2`$ and all other $`a_i`$ equal to zero. Repeating this analysis on the last side $`v_d`$ of $`P`$ we obtain:
$`a_1\left(V_1tv_d\right)`$ $`=`$ $`2\left(1t\right)`$
$`a_d\left(V_1tv_d\right)`$ $`=`$ $`2t`$ (A.48)
and all other charges $`a_i`$ equal to zero.
Deriving the previous relations with respect to $`t`$ we obtain the gradient of the $`a_i`$ along the sides $`v_1`$ and $`v_d`$ of $`P`$:
$$\begin{array}{c}v_1\stackrel{}{}a_1\left(V_1\right)=2\hfill \\ v_d\stackrel{}{}a_1\left(V_1\right)=2\hfill \end{array}\begin{array}{c}v_1\stackrel{}{}a_2\left(V_1\right)=2\hfill \\ v_d\stackrel{}{}a_2\left(V_1\right)=0\hfill \end{array}\begin{array}{c}v_1\stackrel{}{}a_d\left(V_1\right)=0\hfill \\ v_d\stackrel{}{}a_d\left(V_1\right)=2\hfill \end{array}$$
(A.49)
and zero for all other charges different from $`a_d`$, $`a_1`$, $`a_2`$. Finally relation (A.25)
$$r_1=\frac{v_d}{A_dl_1}\frac{v_1}{A_1l_1}$$
(A.50)
allows to compute $`\stackrel{~}{a}_i(x,y)`$ from (A.46):
$$\{\begin{array}{c}\stackrel{~}{a}_1=22\alpha 2\beta \hfill \\ \stackrel{~}{a}_d=2\alpha \hfill \\ \stackrel{~}{a}_2=2\beta \hfill \end{array}\{\begin{array}{c}\alpha \frac{A_1}{v_d,v_1}\hfill \\ \beta \frac{A_d}{v_d,v_1}\hfill \end{array}$$
(A.51)
All other $`\stackrel{~}{a}_i`$ different from $`\stackrel{~}{a}_d`$, $`\stackrel{~}{a}_1`$, $`\stackrel{~}{a}_2`$ are zero. This fact, together with (A.41), allows to disentangle the complex combinatorics and to perform a straightforward, but quite long, computation of $`\mathrm{tr}R^3`$.
So we obtain:
$`\mathrm{tr}R^3=\mathrm{tr}R^2\stackrel{~}{R}=F+{\displaystyle \underset{(i,j)C}{}}v_i,v_j\left(a_{i,j}1\right)^2\left(\stackrel{~}{a}_{i,j}1\right)`$ (A.52)
$`=`$ $`F{\displaystyle \underset{(i,j)C}{}}v_i,v_j\left(a_{i,j}1\right)\left(\stackrel{~}{a}_{i,j}1\right)+{\displaystyle \underset{h=1}{\overset{d}{}}}a_h\left({\displaystyle \underset{(i,j)C_h}{}}v_i,v_j\left(a_{i,j}1\right)\left(\stackrel{~}{a}_{i,j}1\right)\right)`$
$``$ $`F{\displaystyle \underset{(i,j)C}{}}v_i,v_j\left(a_{i,j}1\right)\left(\stackrel{~}{a}_{i,j}1\right)+{\displaystyle \underset{h=1}{\overset{d}{}}}a_hc_h`$
where $`c_h`$ are defined by the last equality. With similar tricks as in previous Appendices, we see that:
$$c_{j+1}c_j=v_j,T_j$$
(A.53)
where the vector $`T_j`$ is:
$`T_j`$ $`=`$ $`\overline{T}_j+v_j{\displaystyle \frac{4}{S}}{\displaystyle \frac{v_j}{A_j}}`$
$`\overline{T}_j`$ $`=`$ $`v_{j+1}\left(a_{j+1}1\right)\left(\stackrel{~}{a}_{j+1}1\right)\mathrm{}+v_{k_j}\left(a_{j+1}\mathrm{}+a_{k_j}1\right)\left(\stackrel{~}{a}_{j+1}\mathrm{}+\stackrel{~}{a}_{k_j}1\right)`$ (A.54)
$`+v_{j1}\left(a_j1\right)\left(\stackrel{~}{a}_j1\right)\mathrm{}+v_{k_j+1}\left(a_j\mathrm{}+a_{k_j+2}1\right)\left(\stackrel{~}{a}_j\mathrm{}+\stackrel{~}{a}_{k_j+2}1\right)`$
$`=`$ $`v_{j+1}\left(a_{j+1}1\right)\left(\stackrel{~}{a}_{j+1}1\right)+v_{j+2}\left(a_{j+1}+a_{j+2}1\right)\left(\stackrel{~}{a}_{j+1}+\stackrel{~}{a}_{j+2}1\right)\mathrm{}`$
$`+v_{j1}\left(a_{j+1}+a_{j+2}\mathrm{}+a_{j1}1\right)\left(\stackrel{~}{a}_{j+1}+\stackrel{~}{a}_{j+2}\mathrm{}+\stackrel{~}{a}_{j1}1\right)`$
and the pieces proportional to $`v_j`$ have been introduced for later convenience. For the difference of consecutive $`\overline{T}_j`$ we obtain:
$`\overline{T}_2\overline{T}_1=`$ (A.55)
$`=`$ $`v_3\left(a_31\right)\left(\stackrel{~}{a}_31\right)+v_4\left(a_3+a_41\right)\left(\stackrel{~}{a}_3+\stackrel{~}{a}_41\right)`$
$`\mathrm{}+v_1\left(a_3+a_4\mathrm{}+a_11\right)\left(\stackrel{~}{a}_3+\stackrel{~}{a}_4\mathrm{}+\stackrel{~}{a}_11\right)`$
$`v_2\left(a_21\right)\left(\stackrel{~}{a}_21\right)v_3\left(a_2+a_31\right)\left(\stackrel{~}{a}_2+\stackrel{~}{a}_31\right)`$
$`\mathrm{}v_d\left(a_2+a_3\mathrm{}+a_d1\right)\left(\stackrel{~}{a}_2+\stackrel{~}{a}_3\mathrm{}+\stackrel{~}{a}_d1\right)`$
$`=`$ $`a_2\stackrel{~}{a}_2\left[v_2+v_3\mathrm{}+v_d\right]+v_2a_2+v_2\stackrel{~}{a}_2v_2+v_1\left(a_21\right)\left(\stackrel{~}{a}_21\right)`$
$`a_2\left[v_3\left(\stackrel{~}{a}_31\right)+v_4\left(\stackrel{~}{a}_3+\stackrel{~}{a}_41\right)\mathrm{}+v_d\left(\stackrel{~}{a}_3+\stackrel{~}{a}_4\mathrm{}+\stackrel{~}{a}_d1\right)\right]`$
$`\stackrel{~}{a}_2\left[v_3\left(a_31\right)+v_4\left(a_3+a_41\right)\mathrm{}+v_d\left(a_3+a_4\mathrm{}+a_d1\right)\right]`$
$`=`$ $`a_2\left[v_3\left(\stackrel{~}{a}_31\right)+v_4\left(\stackrel{~}{a}_3+\stackrel{~}{a}_41\right)\mathrm{}+v_1\left(\stackrel{~}{a}_3+\stackrel{~}{a}_4\mathrm{}+\stackrel{~}{a}_11\right)\right]`$
$`\stackrel{~}{a}_2\left[v_3\left(a_31\right)+v_4\left(a_3+a_41\right)\mathrm{}+v_1\left(a_3+a_4\mathrm{}+a_11\right)\right]`$
$`+v_2a_2+v_2\stackrel{~}{a}_2+v_1v_2`$
$`=`$ $`a_2\left[\stackrel{~}{a}_1v_1+\stackrel{~}{a}_d\left(v_1+v_d\right)\mathrm{}+\stackrel{~}{a}_3\left(v_1+v_d\mathrm{}+v_3\right)\right]`$
$`\stackrel{~}{a}_2\left[a_1v_1+a_d\left(v_1+v_d\right)\mathrm{}+a_3\left(v_1+v_d\mathrm{}+v_3\right)\right]+v_1v_2`$
$`=`$ $`v_1v_2+a_2\left({\displaystyle \underset{i=1}{\overset{d}{}}}\stackrel{~}{a}_i\left(V_iV_2\right)\right)+\stackrel{~}{a}_2\left({\displaystyle \underset{i=1}{\overset{d}{}}}a_i\left(V_iV_2\right)\right)`$
$`=`$ $`v_1v_22a_2r_22\stackrel{~}{a}_2r_2=v_1v_2{\displaystyle \frac{4}{S}}{\displaystyle \frac{v_1}{A_1}}+{\displaystyle \frac{4}{S}}{\displaystyle \frac{v_2}{A_2}}2\stackrel{~}{a}_2r_2`$
where in the last step we used (A.42) (which is also true for $`\stackrel{~}{a}_i`$, as one deduces from its Taylor expansion up to linear terms), $`r_2=V_2b`$ and (A.25). By relabelling indices:
$$T_{j+1}T_j=2\stackrel{~}{a}_{j+1}r_{j+1}$$
(A.56)
Note that
$$T_2=T_3\mathrm{}=T_{d1}=T_12\stackrel{~}{a}_2r_2$$
(A.57)
since $`\stackrel{~}{a}_i`$ are zero for $`i=3,4,\mathrm{}d1`$. We obtain then
$`{\displaystyle \underset{h=1}{\overset{d}{}}}a_hc_h=c_1\left(a_1+a_2\mathrm{}+a_d\right)+a_2\left(c_2c_1\right)\mathrm{}+a_d\left(c_dc_1\right)`$ (A.58)
$`=`$ $`2c_1+a_2v_1,T_1+a_3\left(v_1,T_1+v_2,T_2\right)\mathrm{}+a_d\left(v_1,T_1\mathrm{}+v_{d1},T_{d1}\right)`$
$`=`$ $`2c_1+a_2v_1+a_3\left(v_1+v_2\right)\mathrm{}+a_d\left(v_1+v_2\mathrm{}v_d\right),T_1`$
$`2\stackrel{~}{a}_2a_3v_2+a_4\left(v_2+v_3\right)\mathrm{}+a_d\left(v_2+v_3\mathrm{}+v_{d1}\right),r_2`$
$`=`$ $`2c_1+{\displaystyle \underset{i=1}{\overset{d}{}}}a_i\left(V_iV_1\right),T_12\stackrel{~}{a}_2{\displaystyle \underset{i=1}{\overset{d}{}}}a_i\left(V_iV_2\right),r_2+2\stackrel{~}{a}_2\left(V_1V_2\right)a_1,r_2`$
$`=`$ $`2c_12r_1,T_1+4\stackrel{~}{a}_2r_2,r_22\stackrel{~}{a}_2a_1v_1,r_2`$
$`=`$ $`2c_12r_1,T_1+4a_1{\displaystyle \frac{A_dA_1}{v_d,v_1}}=2c_12r_1,T_1+{\displaystyle \frac{4a_1}{l_1}}`$
$`=`$ $`2c_12r_1,T_1+{\displaystyle \frac{8}{S}}`$
where we have used the explicit expression (A.51) for $`\stackrel{~}{a}_2`$, and performed the substitution $`a_1=2l_1/S`$.
From the definition we now compute:
$`\overline{T}_1=`$ (A.59)
$`=`$ $`v_2\left(a_21\right)\left(\stackrel{~}{a}_21\right)+v_3\left(a_2+a_31\right)\left(\stackrel{~}{a}_2+\stackrel{~}{a}_31\right)\mathrm{}`$
$`+v_d\left(a_2+a_3\mathrm{}+a_d1\right)\left(\stackrel{~}{a}_2+\stackrel{~}{a}_3\mathrm{}+\stackrel{~}{a}_d1\right)`$
$`=`$ $`\left(\stackrel{~}{a}_21\right)\left[v_2\left(a_21\right)+v_3\left(a_2+a_31\right)\mathrm{}+v_d\left(a_2\mathrm{}+a_d1\right)\right]+\stackrel{~}{a}_dv_d\left(1a_1\right)`$
$`=`$ $`\left(\stackrel{~}{a}_21\right)\left[\left(v_2+v_3\mathrm{}+v_d\right){\displaystyle \underset{i=1}{\overset{d}{}}}a_i\left(V_iV_1\right)\right]+\stackrel{~}{a}_dv_d\left(1a_1\right)`$
$`=`$ $`\left(\stackrel{~}{a}_21\right)\left(v_1+2r_1\right)+\stackrel{~}{a}_dv_d\left(1a_1\right)`$
and hence
$`r_1,T_1`$ $`=`$ $`r_1,\left(\stackrel{~}{a}_21\right)\left(v_1+2r_1\right)+\stackrel{~}{a}_dv_d\left(1a_1\right)+v_1{\displaystyle \frac{4}{S}}{\displaystyle \frac{v_1}{A_1}}`$ (A.60)
$`=`$ $`\stackrel{~}{a}_2r_1,v_1+\stackrel{~}{a}_d\left(1a_1\right)r_1,v_d{\displaystyle \frac{4}{SA_1}}r_1,v_1`$
$`=`$ $`{\displaystyle \frac{2A_dA_1}{v_d,v_1}}+{\displaystyle \frac{2A_dA_1}{v_d,v_1}}\left(1a_1\right){\displaystyle \frac{4}{S}}={\displaystyle \frac{2}{l_1}}\left(2{\displaystyle \frac{2l_1}{S}}\right){\displaystyle \frac{4}{S}}`$
$`=`$ $`{\displaystyle \frac{4}{l_1}}{\displaystyle \frac{8}{S}}`$
where again we have used the explicit expressions for $`\stackrel{~}{a}_i`$ in (A.51) and the substitution $`a_1=2l_1/S`$. Collecting pieces together we obtain for (A.58):
$$\underset{h=1}{\overset{d}{}}a_hc_h=2c_1+\frac{24}{S}\frac{8}{l_1}$$
(A.61)
Going back to (A.52) we obtain:
$`\mathrm{tr}R^3=`$ (A.62)
$`=`$ $`F{\displaystyle \underset{(i,j)C}{}}v_i,v_j\left(a_{i,j}1\right)\left(\stackrel{~}{a}_{i,j}1\right)+{\displaystyle \underset{h=1}{\overset{d}{}}}a_hc_h`$
$`=`$ $`F+{\displaystyle \underset{(i,j)C}{}}v_i,v_j\left(a_{i,j}1\right){\displaystyle \underset{(i,j)C}{}}v_i,v_j\left(a_{i,j}1\right)\stackrel{~}{a}_{i,j}+{\displaystyle \underset{h=1}{\overset{d}{}}}a_hc_h`$
$`=`$ $`{\displaystyle \underset{(i,j)C}{}}v_i,v_j\left(a_{i,j}1\right)\stackrel{~}{a}_{i,j}+{\displaystyle \underset{h=1}{\overset{d}{}}}a_hc_h`$
where in the last step we used $`\mathrm{tr}R=0`$ (A.9).
Let us expand the first term in the previous equality; we use the explicit form (A.51) of the $`\stackrel{~}{a}_i`$ putting in evidence the factors of $`2`$, $`\alpha `$ and $`\beta `$:
$`{\displaystyle \underset{(i,j)C}{}}v_i,v_j\left(a_{i,j}1\right)\stackrel{~}{a}_{i,j}=`$
$`=`$ $`2{\displaystyle \underset{(i,j)C_1}{}}v_i,v_j\left(a_{i,j}1\right)`$
$`+2\alpha \left({\displaystyle \underset{(i,j)\left(C_dC_1\right)}{}}v_i,v_j\left(a_{i,j}1\right)\right)2\alpha \left({\displaystyle \underset{(i,j)\left(C_1C_d\right)}{}}v_i,v_j\left(a_{i,j}1\right)\right)`$
$`+2\beta \left({\displaystyle \underset{(i,j)\left(C_2C_1\right)}{}}v_i,v_j\left(a_{i,j}1\right)\right)2\beta \left({\displaystyle \underset{(i,j)\left(C_1C_2\right)}{}}v_i,v_j\left(a_{i,j}1\right)\right)`$
Expanding the factor $`2c_1`$ in (A.61) we obtain:
$`2c_1=`$
$`=`$ $`2{\displaystyle \underset{(i,j)C_1}{}}v_i,v_j\left(a_{i,j}1\right)\left(\stackrel{~}{a}_{i,j}1\right)`$
$`=`$ $`2{\displaystyle \underset{(i,j)C_1}{}}v_i,v_j\left(a_{i,j}1\right)+2{\displaystyle \underset{(i,j)C_1}{}}v_i,v_j\left(a_{i,j}1\right)\stackrel{~}{a}_{i,j}`$
$`=`$ $`2{\displaystyle \underset{(i,j)C_1}{}}v_i,v_j\left(a_{i,j}1\right)+4{\displaystyle \underset{(i,j)C_1}{}}v_i,v_j\left(a_{i,j}1\right)`$
$`4\alpha \left({\displaystyle \underset{(i,j)\left(C_1C_d\right)}{}}v_i,v_j\left(a_{i,j}1\right)\right)4\beta \left({\displaystyle \underset{(i,j)\left(C_1C_2\right)}{}}v_i,v_j\left(a_{i,j}1\right)\right)`$
Then equation (A.62), using (A.61), (LABEL:term1), (LABEL:term2), becomes:
$`\mathrm{tr}R^3=`$ (A.65)
$`=`$ $`2\alpha \left({\displaystyle \underset{(i,j)\left(C_1C_d\right)}{}}+{\displaystyle \underset{(i,j)\left(C_dC_1\right)}{}}\right)v_i,v_j\left(a_{i,j}1\right)`$
$`2\beta \left({\displaystyle \underset{(i,j)\left(C_1C_2\right)}{}}+{\displaystyle \underset{(i,j)\left(C_2C_1\right)}{}}\right)v_i,v_j\left(a_{i,j}1\right)+{\displaystyle \frac{24}{S}}{\displaystyle \frac{8}{l_1}}`$
and it is easy to compute:
$`\left({\displaystyle \underset{(i,j)\left(C_1C_d\right)}{}}+{\displaystyle \underset{(i,j)\left(C_dC_1\right)}{}}\right)v_i,v_j\left(a_{i,j}1\right)=`$ (A.66)
$`=`$ $`v_d,v_1\left(a_11\right)+v_d,v_2\left(a_1+a_21\right)\mathrm{}+v_d,v_{k_d}\left(a_1+a_2\mathrm{}+a_{k_d}1\right)`$
$`v_d,v_{d1}\left(a_d1\right)\mathrm{}v_d,v_{k_d+1}\left(a_d+a_{d1}\mathrm{}+a_{k_d+2}1\right)`$
$`=`$ $`v_d,v_1\left(a_11\right)+v_2\left(a_1+a_21\right)\mathrm{}+v_{d1}\left(a_1+a_2\mathrm{}+a_{d1}1\right)`$
$`=`$ $`v_d,(v_1+v_2\mathrm{}+v_{d1})+a_1(v_1+v_2\mathrm{}+v_{d1})+a_2(v_2\mathrm{}+v_{d1})\mathrm{}+a_{d1}v_{d1}`$
$`=`$ $`v_d,{\displaystyle \underset{i=1}{\overset{d}{}}}a_i\left(V_iV_d\right)=2v_d,r_d=2A_d`$
and similarly, with the opportune changes in the indexes:
$$\left(\underset{(i,j)\left(C_2C_1\right)}{}+\underset{(i,j)\left(C_1C_2\right)}{}\right)v_i,v_j\left(a_{i,j}1\right)=2A_1$$
(A.67)
Finally equation (A.65) becomes:
$`\mathrm{tr}R^3`$ $`=`$ $`{\displaystyle \frac{24}{S}}{\displaystyle \frac{8}{l_1}}2{\displaystyle \frac{A_1}{v_d,v_1}}\left(2A_d\right)2{\displaystyle \frac{A_d}{v_d,v_1}}\left(2A_1\right)`$ (A.68)
$`=`$ $`{\displaystyle \frac{24}{S}}`$
and this concludes our proof. |
warning/0506/quant-ph0506223.html | ar5iv | text | # Quantum estimation of relative information
## I Introduction
Estimating the state of a given quantum system is a fundamental primitive of many quantum information tasks. This problem is usually translated to the estimation of the value of a physical parameter describing specific properties of the preparation procedure. In many instances mp ; gp ; scudo1 ; scudo2 ; bagan1 ; bagan2 ; lindner ; dariano ; holevo1 global parameters of the state space define a natural scheme for encoding quantum information. The global parameters describe collective degrees of freedom of a system with respect to the external environment and are often related to an overall symmetry transformation of the state.
However, encoding information into global degrees of freedom may be often problematic, due to lack of knowledge of the reference frame with respect to which they were prepared, or due to collective decoherence by which they are affected. Encoding information into relative degrees of freedom, possible whenever a quantum system is decomposable into parts, can overcome many of the difficulties encountered in these situations. Such an encoding scheme has been demonstrated experimentally banaszek ; kwiat and can be applied to quantum computation zanardi , communication bourennane ; commnorf and cryptography walton ; boileau .
The aim of this work is to develop efficient preparations and measurement schemes for the relative parameters describing symmetries between different components of a system. We note that such measurements can induce relative relations when previously absent, as in the case of the relative phase between two Fock states or the relative position between two momentum eigenstates hugo .
In this paper we specifically confront the task of efficient estimation of relative rotation angles between two representation vectors of the $`SU(2)`$ symmetry group. This problem was first addressed by Bartlett et. al. terry , who explicitly worked out the estimation of $`SU(2)`$ rotation angles between two spin coherent states. In this article we proceed along the lines of earlier work that Asher Peres started together with us, and propose an extension of the previous methods. Our approach is based on an optimization of the quantum states used in such protocols with respect to some average measure of success of the estimation task, which we shall refer to as the fidelity. The key to this problem lies in the decomposition of both the signal and the measurement elements in irreducible components, invariant under global rotation transformations.
In the following section we discuss the general mathematical structure of the problem. In section III.1 we derive the optimal measurement for the case in which one system is comprised of two spin-$`1/2`$, and the other of one spin-$`1/2`$. We find that preparing two spin-$`1/2`$ parallel to each other leads to a marginally higher fidelity than the antiparallel case; then we determine an optimal preparation procedure which gives a higher fidelity then the ones achieved by the above preparations. This is in contrast to the known results for transmitting a spatial direction gp . We then proceed in section III.2 by replacing the single spin state of the second system with a spin-$`j`$ coherent state and determine the optimal preparations for the cases of anti-parallel and parallel spins, which now yield nearly the same fidelity for any value of $`j`$. We then study the quantum/classical correspondence by considering the limit in which $`j`$ becomes very large. In section III.3 it is shown that in this limit the problem reduces to estimating the first state with respect to a classical reference direction. Our results establish the correspondence between relative degrees of freedom in quantized systems and collective degrees of freedom defined with respect to a classical reference frame.
## II Estimation of relative parameters
### II.1 Formulation of the general problem and basic notations
We begin this section by formulating the problem for a general symmetry group and introducing the basic notation that is useful for our general scheme. In the following, let $`G`$ denote a symmetry group, compact or finite, that describes the global properties of the system through its action on a set of parameters $`T`$. We shall consider $`G=SU(2)`$ acting on quantum spin states, parametrized by the set of rotation angles. The possible states of the system are pure states $`\mathrm{\Psi }(\mathrm{\Theta }),\mathrm{\Theta }T`$ in a $`d`$–dimensional Hilbert space $``$ carrying a unitary representation $`\{U(g),gG\}`$ of $`G`$. To introduce relative symmetry transformations, we assume that the representation space $``$ is a tensor product of two components $`=_1_2`$ of dimensions $`d_1,d_2`$ respectively. The representation of $`G`$ on $``$ is decomposed into the product $`\{U_1(g)U_2(g)\}`$ on $`_1_2`$. This product representation of $`G\times G`$, in which each component is transformed by the same element $`gG`$, is isomorphic to $`G`$ itself. We introduce another set, $`t`$, of parameters $`\theta `$, to describe a relative symmetry between the two components, represented by a group of transformations $`\stackrel{~}{G}`$, which can be the same $`G`$ as before or a subgroup of the latter. $`\stackrel{~}{G}`$ refers to a symmetry property of one of the subsystems (say 2) with respect to the other (say 1), such as, in our case, a relative rotation angle. We call $`U_2(h)`$ its representation operators on the space $`_2`$. We shall restrict ourselves to the case where each of the subsystems is prepared in a pure state, although the formalism can easily be extended to mixed states. The total state on $``$ can be written in terms of the two sets of parameters as $`\mathrm{\Psi }(\mathrm{\Theta },\theta )`$. Its transformation under a global operation is
$$U(g)\mathrm{\Psi }(\mathrm{\Theta },\theta )U_1(g)U_2(g)\mathrm{\Psi }(\mathrm{\Theta },\theta )=\mathrm{\Psi }(g\mathrm{\Theta },\theta ).$$
(1)
The objective of the construction is to define an efficient estimation procedure for the relative parameters $`\theta `$, overlooking the information carried by the global parameters $`\mathrm{\Theta }`$. Note that we might be interested only in estimating a subset of the relative parameters (which will be the case in the following sections). In order to quantify the efficiency of our estimation procedure, we choose a utility function $`f(\mu ,\theta )`$ which measures the deviation of the estimated parameters $`\mu `$ from their true values $`\theta `$. We consider only utility functions which are invariant under global rotations. The measurement apparatus, represented by the POVM $`\{E_\mu \}`$, should be constructed such that it maximizes the average fidelity, denoted by $`F`$, and given by
$$F\left\{E_\mu \right\}=\underset{\mu }{}𝑑\theta 𝑑\mathrm{\Theta }P(\mathrm{\Theta },\theta )\mathrm{Tr}\left[\rho (\mathrm{\Theta },\theta )E_\mu \right]f(\mu ,\theta ),$$
(2)
where $`P(\mathrm{\Theta },\theta )`$ is a prior probability distribution over the global and relative parameters, and $`\rho (\mathrm{\Theta },\theta )=|\mathrm{\Psi }(\mathrm{\Theta },\theta )\mathrm{\Psi }(\mathrm{\Theta },\theta )|`$. As noted in terry , we can assume that the global and relative parameters are independent random variables, and that the global parameter is uniformly distributed on its domain of definition, i.e.,
$$P(\mathrm{\Theta },\theta )d\theta d\mathrm{\Theta }=p(\theta )d\theta d\mathrm{\Theta }.$$
(3)
Now, using the definition for the global transformation (1) and the properties of the trace, we can write
$$F\left\{E_\mu \right\}=\underset{\mu }{}𝑑\theta p(\theta )\mathrm{Tr}\left[\overline{\rho }(\theta )E_\mu \right]f(\mu ,\theta ),$$
(4)
where
$$\overline{\rho }(\theta )=𝑑gU^{}(g)\rho (\mathrm{\Theta },\theta )U(g).$$
(5)
In the last equation, $`dg`$ is the invariant measure for the group $`G`$. As a consequence of equations (4) and (5), we need only to consider a reduced form of the input state which is manifestly invariant under global transformations. Schur’s lemma sternberg then assures that the input state is block diagonal in the irreducible representations of $`G`$,
$$\overline{\rho }(\theta )=\underset{J}{}\rho ^{(J)}(\theta ).$$
(6)
The above considerations also have implications on the form of the optimal measurement. In fact, using the global invariance of $`\overline{\rho }`$, we have
$`F\left\{E_\mu \right\}`$ $`=`$ $`F\left\{U(g)E_\mu U^{}(g)\right\}`$
$`=`$ $`F\left\{{\displaystyle 𝑑gU(g)E_\mu U^{}(g)}\right\}.`$
This relation implies that when searching for the optimal estimation procedure it suffices to consider POVMs which are invariant under a global transformation
$$U(g)E_\mu U^{}(g)=E_\mu ,gG.$$
(7)
Equation (7) is by itself a strong prerequisite on the structure of the POVM elements. Indeed, if combined with Schur’s lemma, it implies that whenever the total Hilbert space $``$ can be decomposed into a direct sum of irreducible representations under the global transformation, the optimal POVM elements labelling the different outcomes have the simple form
$$E_\mu =\underset{J}{}p_{\mu ,J}E_\mu ^J,$$
(8)
where each of the operators $`E_\mu ^J`$ has support only on the representation labelled by $`J`$ and $`_\mu p_{\mu ,J}E_\mu ^J=\text{1l}_J`$. Here the symbol $`\text{1l}_J`$ denotes the identity operator in the $`J`$ subspace. The search for optimal POVMs may be further restricted to the case in which all elements have support on only one representation space, due to the linearity of the fidelity functional (given a POVM of the form (8), the POVM with elements $`E_{\mu ,J}p_{\mu ,J}E_\mu ^J`$ yields the same fidelity as the original one by linearity of the trace and obeys the restriction that each element has support on only one representation).
In the following sections, the general considerations outlined above will be applied to the problem of transmitting relative rotation angles of the $`SU(2)`$ group.
### II.2 Estimation of a relative angle between spin coherent states
In estimating relative, as opposed to absolute, rotation angles, we assume no prior knowledge of the overall orientation of the classical frame in which the system is defined. Following the notation introduced above, the problem may be illustrated as follows.
Imagine that, with no prior knowledge on the absolute spatial orientation of two observers, Alice and Bob, we were requested to estimate the angle $`\beta `$ between two unit vectors $`\widehat{𝐧}_1`$ and $`\widehat{𝐧}_2`$, each chosen by one of them, by measuring a pair of $`SU(2)`$ spin states prepared in the corresponding reference frame of each observer. This task is in general possible owing to the fact that states belonging to an $`SU(2)`$ representation space can be used as intrinsic direction indicators mp ; gp ; scudo1 ; scudo2 ; bagan1 ; bagan2 ; lindner ; dariano ; holevo1 and therefore it makes sense to consider relative angles between them.
The simplest way to achieve the task would be to consider two $`SU(2)`$ coherent states, corresponding, say, to spins $`j_1,j_2`$,
$$\widehat{𝐧}_1𝐉_1|\psi _1=j_1|\psi _1,\widehat{𝐧}_2𝐉_2|\psi _2=j_2|\psi _2.$$
(9)
Without loss of generality, we can assume that Alice and Bob choose the $`z`$-axis of their reference frame. A state denoted by $`|\psi _2`$ in Bob’s frame is written in Alice’s reference frame as $`U_2^z(\alpha )U_2^y(\beta )U_2^z(\gamma )|\psi _2`$, where $`U_k^x(\alpha )=e^{i𝐉_k𝐱\alpha }`$, and similarly for the other directions. The angles $`\alpha ,\beta ,\gamma `$ are the three Euler angles relating Alice’s reference frame to the one of Bob. Note that the angle $`\beta [0,\pi ]`$ is also the angle between $`\widehat{𝐧}_1`$ and $`\widehat{𝐧}_2`$. We introduce a global reference frame which is rotated with respect to Alice’s frame by an angle $`\alpha `$ around the $`z`$-axis. In this frame the composite state is given by
$$|\mathrm{\Psi }(\alpha ,\beta ,\gamma )=U_1^z(\alpha )|\psi _1U_2^z(\beta )U_2^z(\gamma )|\psi _2.$$
(10)
For spin coherent states, as in Eq. (9), the above equation is simplified to
$$|\mathrm{\Psi }(\beta )=|j_1,m_1=j_1U_2^y(\beta )|j_2,m_2=j_2,$$
(11)
up to an overall phase. Notice that in Eq. (10),(11) we have implicitly specified the global parameter $`\mathrm{\Theta }`$, which we shall omit from now on from our notation. So far we have a set of three relative parameters, $`\theta =\{\alpha ,\beta ,\gamma \}`$, with a joint probability distribution given by the Haar-measure
$$p(\alpha ,\beta ,\gamma )=\frac{1}{8\pi ^2}\mathrm{sin}\beta ,$$
(12)
which corresponds to a random orientation of Alice’s and Bob’s reference frames. Since the party making the measurement (Bob) is not interested in estimating $`\alpha ,\gamma [0,2\pi ]`$, these parameters are averaged out by integrating over their range. We are then left with the probability distribution $`p(\beta )=\mathrm{sin}\beta /2`$, which corresponds to the probability density of the angle between two random unit vectors in three dimensions.
Let us denote by $`_1,_2`$ the Hilbert spaces of the systems prepared by Alice and Bob respectively, carrying the $`SU(2)`$ representations $`j_1`$ and $`j_2`$. The composite Hilbert space $`=_1_2`$ carries a diagonal product representation of $`SU(2)SU(2)`$, $`U^{j_1}(g)U^{j_2}(g),gSU(2)`$, which corresponds to a global symmetry operation labelled by the parameter $`g`$, acting identically on the two subspaces. This representation may be reduced as $`_{J=|j_1j_2|}^{j_1+j_2}_J`$, where each component has multiplicity one. Invariance of the measurement operators under a global rotation (as explained in the previous section) reduces the signal state to the form
$$\overline{\rho }(\beta )=\underset{J}{}p_J(\beta )\mathrm{\Pi }_J,$$
(13)
where $`\mathrm{\Pi }_J`$ are projectors on the representations $`J`$. Since each representation of the global rotation (specified above by $`J`$) appears only once in the signal state, the measurement process amounts to estimating a probability distribution over the relative angle $`\beta `$. The scenario described above, where each of the parties prepares a spin coherent state, is the one examined in terry .
## III General Scheme
In the procedure outlined in the last section, the two parties, Alice and Bob, use spin coherent states in order to indicate their chosen direction. However, for the task considered here this is not the optimal preparation. It is known, in fact, that optimal $`SU(2)`$ direction indicators exploit entanglement between components belonging to different irreducible representations scudo1 ; bagan1 ; bagan2 ; scudo2 ; dariano . This suggests to consider the following general encoding procedure. Let
$$|\mathrm{\Phi }=\underset{j_1=0}{\overset{j_{max}}{}}\underset{m_1=j_1}{\overset{j_1}{}}a_{m_1}^{j_1}|j_1m_1$$
(14)
be a generic state in $`_1`$. By choosing a unit vector $`\widehat{𝐧}_1`$, Alice would prepare the state $`U(\widehat{𝐧}_1)|\psi _1`$,
$$|\psi _1=U(\widehat{𝐧}_1)|\mathrm{\Phi },$$
(15)
where $`U(\widehat{𝐧}_1)`$ is a unitary operator corresponding to the rotation which carries Alice’s $`\widehat{𝐳}`$-axis onto $`\widehat{𝐧}_1`$. Our goal is to find the optimal state $`|\mathrm{\Phi }`$ for the case in which Bob indicates his direction $`\widehat{𝐧}_2`$ with a coherent state $`|\psi _2`$ satisfying
$$\widehat{𝐧}_2𝐉_2|\psi _2=j_2|\psi _2.$$
(16)
Note that equations (15) and (16) are written in Alice’s and Bob’s reference frames, respectively. In a global reference frame, specified as in the last section, the total state is given by
$`|\mathrm{\Psi }(\alpha ,\beta ,\gamma )`$ $`=`$ $`U_1^z(\alpha )|\psi _1U_2^y(\beta )U_2^z(\gamma )|\psi _2`$ (17)
$`=`$ $`{\displaystyle \underset{j_1m_1}{}}a_{m_1}^{j_1}e^{im_1\alpha }|j_1m_1U_2^y(\beta )e^{ij_2\gamma }|j_2j_2.`$
As before, the state $`|\mathrm{\Psi }`$ is expressed in terms of an arbitrary orientation of the global reference frame, so that the parameter $`\mathrm{\Theta }`$ may be omitted. However, one can see that the angles $`\alpha ,\gamma `$ now induce relative phases between the different components of the state. Since our protocol does not deal with the estimation of $`\alpha `$ and $`\gamma `$ (only the angle $`\beta `$ is considered), we will average over them in the expression of the fidelity
$$F=\underset{\mu }{}d_{\alpha \beta \gamma }𝑑g\mathrm{Tr}\left[U(g)\rho (\alpha ,\beta ,\gamma )U^{}(g)E_\mu \right]f(\mu ,\beta ),$$
(18)
where we denoted $`d_{\alpha \beta \gamma }1/8\pi ^2\mathrm{sin}(\beta )d\alpha d\beta d\gamma `$. Integrating over $`\alpha `$ and $`\gamma `$, we have
$$\frac{1}{4\pi ^2}𝑑\alpha 𝑑\gamma \rho (\alpha ,\beta ,\gamma )=\underset{m_1}{}c_{m_1}\rho _{m_1}U_2^y(\beta )|j_2j_2j_2j_2|U_2^y(\beta ),$$
(19)
with
$$\rho _{m_1}=|\psi _{m_1}\psi _{m_1}|,|\psi _{m_1}=\underset{j_1=m_1}{\overset{j_{max}}{}}a_{m_1}^{j_1}|j_1,m_1,$$
(20)
and $`c_{m_1}`$ given by
$$c_{m_1}=\underset{j_1}{}|a_{m_1}^{j_1}|^2,\underset{m_1}{}c_{m_1}=1.$$
(21)
In the following we will search for the optimal generic state $`|\mathrm{\Phi }`$. Note that after integrating over $`\alpha `$, $`\gamma `$ we get a convex combination of states with different $`m_1`$. Thus the fidelity will contain a linear combination of contributions from the different $`m_1`$ sectors,
$$F=\underset{m_1}{}c_{m_1}F(m_1)\underset{m_1}{\mathrm{max}}F(m_1).$$
(22)
From the above discussion it is clear thanks that the optimal generic state can be taken with $`m_1`$ fixed, i.e, of the form
$$|\mathrm{\Phi }_{m_1}=\underset{j_1=m_1}{\overset{j_{max}}{}}a_{m_1}^{j_1}|j_1m_1.$$
(23)
In the following we shall restrict ourselves to generic states of this form. The rotation in Eq. (17) by the relative angle $`\beta [0,\pi ]`$ is expressed in the standard Euler angle notation as a rotation around the $`y`$-axis, given by the matrix
$$U_2^y(\beta )|j_2j_2\underset{m}{}d_{mj_2}^{j_2}(\beta )|j_2m,$$
(24)
where the $`d_{m^{}m}^{j_2}(\beta )`$ can be expressed using Jacobi polynomials (see for example edmonds ). The superscript in the above equation refers to the $`(2j_2+1)`$–dimensional irreducible representation of spin $`j_2`$. The signal state (17) is given explicitly (up to an overall phase) by
$$|\mathrm{\Psi }(\beta )=\underset{j_1}{}a_{m_1}^{j_1}\underset{m_2}{}d_{m_2j_2}^j(\beta )|j_1m_1,j_2m_2,$$
(25)
where $`|j_1m_1,j_2m_2|j_1m_1|j_2m_2`$. Note that unlike the example discussed in section II.2 and in terry , we now exploit repeated irreducible representations $`J`$, since each value of $`j_1`$ gives rise to a series of total angular momentum $`J=|j_1j_2|,\mathrm{},(j_1+j_2)`$. The equivalent repetitions of the representation $`J`$ are labelled by $`j_1`$, with $`j_2`$ being fixed. The state can be written in the basis $`(JM,j_1j_2)`$ using the unitary transformation
$$|j_1m_1,j_2m_2=\underset{JM}{}C_{j_1m_1j_2m_2}^{JM}|JM,j_1j_2,$$
(26)
with $`C`$ denoting the Clebsch–Gordan coefficients, and $`M`$ denoting the $`z`$ component of the total angular momentum.
We now need to compute the averaged state $`\overline{\rho }(\beta )`$, following Eq. (5). A corollary to Schur’s lemma sternberg then states that, for irreducible representations $`\sigma ,\tau `$ of a group $`G`$, a group–averaged operator satisfies
$$𝑑gU^{(\sigma )}(g)AU^{(\tau )}(g)^{}=\frac{\delta _{\sigma ,\tau }\mathrm{Tr}_\sigma (A)\text{1l}}{d_\sigma },$$
(27)
where $`\delta _{\sigma ,\tau }`$ is a Kronecker delta over the inequivalent representations and the trace is computed over the $`d_\sigma `$-dimensional space of the irreducible representation $`\sigma `$. The above corollary is applied to obtain the invariant reduced density operator $`\overline{\rho }(\beta )`$. Remembering that the global rotation operator $`U(\mathrm{\Omega })`$ is a direct sum of operators $`U(\mathrm{\Omega })=_JU^{(J)}(\mathrm{\Omega })`$, we see from Eq. (27) that $`\overline{\rho }(\beta )`$ is also block diagonal in the representations $`J`$,
$`\overline{\rho }(\beta )`$ $`=`$ $`{\displaystyle U(\mathrm{\Omega })\rho (\beta )U^{}(\mathrm{\Omega })𝑑\mathrm{\Omega }}`$ (28)
$`=`$ $`{\displaystyle \underset{J}{}}\overline{\rho }^{(J)}(\beta ),`$
where the operators $`\rho ^{(J)}(\beta )`$ bear the indices $`j_1,j_1^{}`$ and are given by
$$\overline{\rho }^{(J)}(\beta )_{j_1,j_1^{}}=\underset{M}{}J,M,j_1^{}|\rho (\beta )|J,M,j_1.$$
(29)
Note that by averaging over the global rotation one does not diagonalize the operator $`\overline{\rho }(\beta )`$ with respect to the additional quantum number $`j_1`$. The invariant signal state $`\overline{\rho }(\beta )`$ is block diagonal over the irreducible representations with off diagonal elements across the repeated ones
$$\overline{\rho }(\beta )=\underset{J}{}\underset{j_1j_1^{}}{}\frac{1}{(2J+1)}\underset{Mm_2}{}a_{m_1}^{j_1}a_{m_1}^{j_1^{}}\left(d_{j_2m_2}^{j_2}(\beta )\right)^2C_{j_1m_1j_2m_2}^{JM}C_{j_1^{}m_1j_2m_2}^{JM}|J,j_1J,j_1^{}|.$$
(30)
Remember that $`m_1`$ is fixed.
Note that the state $`\overline{\rho }(\beta )`$ does not, even implicitly, depend on the orientation of the reference frame of the measurement apparatus. We will now use $`\overline{\rho }(\beta )`$ to determine the state $`|\psi _1`$ which will enable optimal estimation of $`\beta `$. To this end, we fix a convenient figure of merit as measure of the discrepancy between the estimated and the given value of $`\beta [0,\pi ]`$, namely the quadratic utility function $`f(\mu ,\beta )=\mathrm{cos}^2((\mu \beta )/2)`$, where $`\mu `$ is the estimated value of the parameter. The choice of the utility function is not unique, and a different choice might lead to different optimal states and POVMs. However, the optimization procedure, as described below, is independent of this choice. The above utility function has the advantage of having been broadly used throughout earlier literature bagan2 ; bagan1 ; scudo1 ; scudo2 ; lindner ; dariano ; mp ; gp .
We denote the POVM elements by $`\{E_\mu \}`$, with $`\mu [0,\pi ]`$, and $`_\mu E_\mu =1`$. The average fidelity with respect to the given figure of merit, integrating over all possible transmitted angles $`\beta `$ and all possible inferred values $`\mu `$, is
$$F[\{\mu \},\{E_\mu \}]=\underset{\mu }{}\mathrm{Tr}[\overline{\rho }(\beta )E_\mu ]\mathrm{cos}^2\left(\frac{\mu \beta }{2}\right)\mathrm{sin}\beta d\beta /2.$$
(31)
Note that the fidelity is a functional both of the set of estimates $`\{\mu \}`$ and of the POVM used for the estimation procedure $`\{E_\mu \}`$, where to each estimate corresponds a
(single) POVM element. The probability of estimating $`\mu `$ for a true angle $`\beta `$ is $`\mathrm{Tr}[\overline{\rho }(\beta )E_\mu ].`$
The above expression can be rewritten by exchanging the order of the integral with the trace (due to the linearity of the integration and finiteness of the sum) and gives the fidelity in the form
$$F[\{\mu \},\{E_\mu \}]=\underset{\mu }{}\mathrm{Tr}\{A_\mu E_\mu \}.$$
(32)
with
$$A_\mu \overline{\rho }(\beta )\mathrm{cos}^2((\mu \beta )/2)\mathrm{sin}\beta d\beta /2.$$
(33)
Since $`\overline{\rho }(\beta )`$ is block diagonal, also $`A_\mu `$ can be written as a direct sum $`A_\mu =_JA_\mu ^J`$.
In terms of the basis representation states of angular momenta $`(J;j_1,j_2)`$, and using Eq. (30), the operator $`A_\mu `$ can be explicitly written as
$`A_\mu `$ $`=`$ $`{\displaystyle \underset{J}{}}{\displaystyle \underset{j_1j_1^{}}{}}{\displaystyle \frac{1}{2J+1}}{\displaystyle \underset{Mm_2}{}}\{a_{m_1}^{j_1}a_{m_1}^{j_1^{}}I_{m_2}^{j_2}(\mu )`$ (34)
$`C_{j_1m_1j_2m_2}^{JM}C_{j_1^{}m_1j_2m_2}^{JM}\}|J,j_1J,j_1^{}|,`$
where
$$I_{m_2}^{j_2}(\mu )\left(d_{j_2m_2}^{j_2}(\beta )\right)^2\mathrm{cos}^2((\mu \beta )/2)\mathrm{sin}\beta d\beta /2.$$
(35)
This expression may be evaluated using the properties of the Wigner functions $`d_{mm^{}}^j`$ and their representation in terms of Jacobi polynomials edmonds . The average fidelity can now be written as a sum of contributions from each subspace of given $`J`$
$$F[\{\mu \},\{E_\mu \}]=\underset{J}{}\underset{\mu }{}\mathrm{Tr}(A_\mu ^JE_\mu ^J),$$
(36)
where $`_\mu E_\mu ^J=\text{1l}_J`$. Given a set of estimates $`\{\mu \}`$, the task of maximizing the expression $`_\mu \mathrm{Tr}(A_\mu ^JE_\mu ^J)`$ is straightforward, at least numerically (for example, by using semidefinite programming yonina ). However, in our approach, in order to maximize the average fidelity (36), we need to maximize over the set $`\{\mu \}`$, so that the maximal fidelity $`F_{\mathrm{max}}`$ will actually be given by
$$F_{\mathrm{max}}=\underset{\{\mu \}}{\mathrm{max}}\underset{\{E_\mu \}}{\mathrm{max}}F[\{\mu \},\{E_\mu \}].$$
(37)
### III.1 The case of $`j_2=1/2`$, $`j_1\{0,1\}`$
We will now solve the optimization problem of Eq. (37) for the following case: the state $`|\psi _2`$ is a spin $`1/2`$ coherent state, while $`|\psi _1`$ is composed of two spin $`1/2`$ systems, so that $`j_1\{0,1\}`$. According to the discussion in the previous section, we can restrict ourselves to two classes of generic states $`|\mathrm{\Phi }`$
$$|\mathrm{\Phi }_0=a|j_1=0m_1=0+\sqrt{1a^2}|j_1=1m_1=0,$$
(38)
and
$$|\mathrm{\Phi }_1=|j_1=1m_1=1$$
(39)
Let us first discuss the case where $`|\mathrm{\Phi }=|\mathrm{\Phi }_1`$. In this simple case, the state $`|\psi _1`$ is just a spin-$`1`$ coherent state or, viewed as composed of two spins, it is a polarized state with parallel spins along the vector $`\widehat{𝐧}_1`$. Coupling the representations of $`\psi _1`$ and $`\psi _2`$ gives
$$\text{1}\text{1/2}=\text{1/2}\text{3/2},$$
and the operators $`A_\mu ^J`$ are one-dimensional, making the optimization trivial. In this case, the optimal measurement simply consists in the projections onto the $`J=1/2`$ and $`J=3/2`$ subspaces, as there are no repeated representations. The fidelity achieved with this state is $`F_{\mathrm{max}}\left[\text{parallel}\right]=0.90983`$.
Next, we consider $`|\mathrm{\Phi }=|\mathrm{\Phi }_0`$. Coupling the representations of $`\psi _1`$ and $`\psi _2`$ gives in this case
$$(\text{1}\text{0})\text{1/2}=\text{1/2}\text{1/2}\text{3/2},$$
and therefore the density matrix $`\overline{\rho }(\beta )`$ contains two blocks of dimensions $`2`$ and $`1`$ corresponding to $`J=1/2`$ and $`J=3/2`$, respectively. The operators $`A_\mu ^J`$ are given by
$`A_\mu ^{1/2}=\left(\begin{array}{cc}\frac{a^2(4+\pi \mathrm{sin}\mu )}{8}& \frac{a\sqrt{1a^2}\mathrm{cos}\mu }{6\sqrt{3}}\\ \frac{a\sqrt{1a^2}\mathrm{cos}\mu }{6\sqrt{3}}& \frac{6(2a^2)+3\pi \left(1a^2\right)\mathrm{sin}\mu }{72}\end{array}\right)`$ (42)
and
$$A_\mu ^{3/2}=\frac{12\left(1a^2\right)+3\left(1a^2\right)\mathrm{sin}\mu }{36}.$$
(43)
We are seeking the set $`\mu `$ and $`E_\mu ^J`$ which maximizes the mean fidelity (36). Let us start with the $`J=3/2`$ subspace. Since this subspace is one-dimensional, the restriction to operators which are invariant under a global rotation leaves us with one operator only, $`E_{\mu _{3/2}}^{3/2}`$, which is the projection operator on the $`J=3/2`$ subspace. The estimate $`\mu _{3/2}`$ which maximizes the corresponding expression for $`A_\mu ^{3/2}`$ in reference to Eq. (43), is obviously given by $`\mu _{3/2}=\pi /2`$.
Next, we consider the 2-dimensional subspace of $`J=1/2`$. Following helstrom , we define an operator $`\mathrm{{\rm Y}}`$ as
$$\mathrm{{\rm Y}}=\underset{\mu }{}A_\mu E_\mu .$$
(44)
For a set $`\{\mu \}`$, a POVM $`\{E_\mu \}`$ is optimal if and only if it satisfies the following set of conditions
$$\mathrm{{\rm Y}}A_\mu 0$$
(45)
for each $`\mu `$ in the set of estimates $`\{\mu \}`$, with the additional requirement that $`\mathrm{{\rm Y}}`$ be hermitian. The inequality sign in Eq. (45) means that the operator $`\mathrm{{\rm Y}}A_\mu `$ must be positive semi-definite. The maximal fidelity will then be given by
$$F_{\mathrm{max}}=\mathrm{Tr}\mathrm{{\rm Y}}.$$
(46)
In order to see that equations (45) and (44) indeed lead to the maximization of the mean fidelity, consider a different POVM $`\{E_\mu ^{}\}`$, such that
$$\underset{\mu }{}E_\mu ^{}=\text{1l}.$$
(47)
The difference between the fidelity achieved with this POVM and the one achieved with the optimal one is
$$F_{\mathrm{max}}F^{}=\mathrm{Tr}\underset{\mu }{}(\mathrm{{\rm Y}}A_\mu )E_\mu ^{},$$
(48)
thanks to Eqs. (46) and (47). Now, if $`C`$ and $`D`$ are positive semi-definite hermitian operators, then they satisfy
$$\mathrm{Tr}(CD)0.$$
(49)
Setting $`C=\mathrm{{\rm Y}}A_\mu `$ and $`D=E_\mu ^{}`$, we obtain
$$F_{\mathrm{max}}F^{}0,$$
(50)
as desired.
Let us first maximize the average fidelity for only two estimates $`\mu _1`$ and $`\mu _2`$, in correspondence to which we have the POVM elements $`E_{\mu _1}^{1/2}+E_{\mu _1}^{1/2}=\text{1l}_{J=1/2}`$. Then
$`\mathrm{{\rm Y}}A_{\mu _1}^{1/2}`$ $`=`$ $`A_{\mu _1}^{1/2}E_{\mu _1}^{1/2}+A_{\mu _2}^{1/2}E_{\mu _2}^{1/2}A_{\mu _1}^{1/2}`$ (51)
$`=`$ $`(A_{\mu _2}^{1/2}A_{\mu _1}^{1/2})E_{\mu _2}^{1/2}0,`$
since $`\mathrm{{\rm Y}}A_{\mu _1}^{1/2}`$ is non-negative if the POVM is optimal. Let us denote by $`\eta _i`$ and $`|\eta _i`$ the eigenvalues and corresponding eigenvectors of the operator $`\mathrm{\Delta }A_{\mu _2}^{1/2}A_{\mu _1}^{1/2}`$. For each $`|\eta _i`$ we can write
$$\eta _i|(A_{\mu _2}^{1/2}A_{\mu _1}^{1/2})E_{\mu _2}^{1/2}|\eta _i=\eta _i\eta _i|E_{\mu _2}^{1/2}|\eta _i0,$$
(52)
using (51). If we assume that $`\eta _i`$ is negative, then (52) gives
$$\eta _i|E_{\mu _2}^{1/2}|\eta _i0;$$
on the other hand, since $`E_{\mu _2}^{1/2}`$ is positive semi-definite, we must have
$$\eta _i|E_{\mu _2}^{1/2}|\eta _i=0,\mathrm{if}\eta _i<0,$$
and similarly
$$\eta _i|E_{\mu _1}^{1/2}|\eta _i=0,\mathrm{if}\eta _i>0.$$
Thus $`E_{\mu _2}^{1/2}`$ projects onto the subspace spanned by $`|\eta _i`$ with $`\eta _i0`$ and $`E_{\mu _1}^{1/2}`$ projects onto the subspace of positive eigenvalues of the operator $`\mathrm{\Delta }`$. The subspace with $`\eta _i=0`$ does not contribute to the fidelity, so that the maximal fidelity is given by
$`F_{\mathrm{max}}=\mathrm{Tr}\mathrm{{\rm Y}}=\mathrm{Tr}A_{\mu _1}^{1/2}+{\displaystyle \underset{\eta _i0}{}}\eta _i.`$ (53)
Let us now assume that $`\mu _1=\mu `$ and $`\mu _2=\pi \mu `$. We shall see that this choice will lead to the optimal measurement for the class of states under consideration. Indeed, we have
$$\mathrm{\Delta }=\left(\begin{array}{cc}0& \frac{a\sqrt{1a^2}\mathrm{cos}\mu }{3\sqrt{3}}\\ \frac{a\sqrt{1a^2}\mathrm{cos}\mu }{3\sqrt{3}}& 0\end{array}\right),$$
(54)
with eigenvalues $`\pm \frac{a\sqrt{1a^2}\mathrm{cos}\mu }{3\sqrt{3}}`$ and corresponding eigenvectors $`|+=(1,1)^T`$ and $`|=(1,1)^T`$. The contribution to the fidelity from the $`J=1/2`$ subspace $`F^{1/2}`$ is now given by
$$F^{1/2}(\mu )=\frac{a\sqrt{1a^2}\mathrm{cos}\mu }{3\sqrt{3}}+\mathrm{Tr}A_\mu ^{1/2}.$$
(55)
At this point, in order to find the maximal mean fidelity under our assumptions, it suffices to maximize the function $`F^{1/2}(\mu )`$. A simple calculation shows that the maximum is attained for
$$\nu =\mathrm{tan}^1\left[3\sqrt{3}(1+2a^2)\pi /(8a\sqrt{1a^2})\right].$$
(56)
It remains to check that indeed the choice
$$\mu _1=\nu ,\mu _2=\pi \nu $$
(57)
leads to the maximal fidelity for the $`J=1/2`$ subspace. A proof of this fact is provided by the following argument. Consider a general set of estimates $`\{𝝁\}=\{\mu _1,\mu _2,\mu _3,\mathrm{},\mu _n\}`$, with a corresponding set of POVM elements $`E_\mu ^{1/2}`$, and let
$$F^{1/2}\left[\{𝝁\}\right]=\underset{\{E_\mu ^{1/2}\}}{\mathrm{max}}F^{1/2}[\{𝝁\},\{E_\mu ^{1/2}\}]$$
(58)
be the maximal fidelity achieved by this set. We would like to show that by adding $`\nu `$ and $`\pi \nu `$ to the set $`\{𝝁\}`$ the mean fidelity can never decrease with respect to the optimal bound and, at the same time, the bound is attained by these two values alone. Let $`\{\stackrel{~}{𝝁}\}`$ denote the new set obtained by adding $`\nu `$ and $`\pi \nu `$, defined by Eq. (56), to the set $`\{𝝁\}`$. The first property,
$$F^{1/2}\left[\{\stackrel{~}{𝝁}\}\right]F^{1/2}\left[\{𝝁\}\right],$$
(59)
simply follows from the fact that adding estimates to a given set can only increase the mean fidelity. To complete the argument we still have to show that
$$F^{1/2}\left[\{\stackrel{~}{𝝁}\}\right]=F^{1/2}\left[\{\nu ,\pi \nu \}\right].$$
(60)
The optimal measurement for $`\{\nu ,\pi \nu \}`$, following the earlier discussion, is defined by the projectors on the negative and positive eigenvectors of the operator $`A_\nu ^{1/2}A_{\pi \nu }^{1/2}`$, i.e.,
$$E_\nu ^{1/2}=||,E_{\pi \nu }^{1/2}=|++|.$$
(61)
Consider a POVM for the set $`\{\stackrel{~}{𝝁}\}`$ which consists of the two operators in (61), and of
$$E_\mu ^{1/2}=0\mathrm{if}\mu \nu ,\pi \nu .$$
This POVM is optimal also for the set $`\{\stackrel{~}{𝝁}\}`$. To see this, we need to check whether the condition
$$\mathrm{{\rm Y}}A_\mu ^{1/2}0$$
holds for all $`\mu \{\stackrel{~}{𝝁}\}`$, with
$$\mathrm{{\rm Y}}=A_\nu ^{1/2}||+A_{\pi \nu }^{1/2}|++|.$$
(62)
Let us evaluate the entries of the operators $`\mathrm{{\rm Y}}A_\mu ^{1/2}`$ in the basis $`|+,|`$. These are given by
$`|\mathrm{{\rm Y}}A_\mu ^{1/2}|=|A_\nu ^{1/2}||A_\mu ^{1/2}|`$
$`+|\mathrm{{\rm Y}}A_\mu ^{1/2}|+=|A_\nu ^{1/2}|+|A_\mu ^{1/2}|+`$
$`+|\mathrm{{\rm Y}}A_\mu ^{1/2}|=+|A_\nu ^{1/2}|+|A_\mu ^{1/2}|,`$ (63)
where we have used $`+|A_{\pi \nu }^{1/2}|+=|A_\nu ^{1/2}|`$. The eigenvalues $`\lambda _1(\mu ),\lambda _2(\mu )`$ of the operator $`\mathrm{{\rm Y}}A_\mu ^{1/2}`$ can now be calculated from (63), and their positivity can be verified (at least numerically). The positivity of these eigenvalues (of which we do not report here the explicit expression) implies then (60). For all input states $`|\psi _1`$ discussed in this paper we have verified that the POVM given in (61), with $`\nu `$ given by (56), is indeed optimal.
From the above optimization procedure we see that the maximal fidelity, as a function of the parameter $`a`$, is given by
$$F_{\mathrm{max}}\left[m_1=0\right]=\frac{a\sqrt{1a^2}\mathrm{cos}\nu }{3\sqrt{3}}+\mathrm{Tr}A_\nu ^{1/2}+\mathrm{Tr}A_{\pi /2}^{3/2},$$
(64)
with $`\nu `$ given by (56). The fidelity as a function of the state parameter $`a`$ is plotted in Fig. 1.
The maximal fidelity is achieved in correspondence of the state $`|\psi _{\mathrm{opt}}`$, by setting $`a=0.609`$, and is $`F_{\mathrm{max}}\left[\psi _{\mathrm{opt}}\right]=0.91092`$. For comparison, the anti-parallel spin state
$$|\mathrm{\Phi }_{\mathrm{anti}}=|=\frac{1}{\sqrt{2}}|00+\frac{1}{\sqrt{2}}|10$$
leads to a fidelity of $`F_{\mathrm{max}}[|]=0.90982`$, lower than the one obtained using the parallel spin state by only a factor $`10^5`$. Similar results for parallel and anti-parallel spin states were obtained by N. Gisin and S. Iblisdir gi .
To conclude this part, we compare the above results to earlier results on quantum direction indicators, where a quantum system carrying a representation of the rotation group is used to transmit a spatial direction between two observes that do not share a common reference frame. As shown in gp ; scudo1 , if the state of the quantum system is constructed from two spin $`1/2`$, (i.e., constrained to have maximal spin 1), encoding the directional information into anti-parallel spins proves to be the optimal strategy. Here we see instead that if the receiver is interested only in the relative orientation of this state with respect to another state, the anti-parallel spin state gives nearly the same fidelity as the parallel one, which is well below the optimum.
### III.2 Higher values of $`j_2`$
The estimation of the relative orientation of two states can be seen as a process in which the first party (Alice) encodes a direction into a quantum state while the receiving party (Bob) attempts to estimate the signal without having a classical reference frame relative to which he can measure it. Therefore Bob resorts to finding the relative orientation of the signal state with respect to the orientation of some given state (say, a coherent state of spin $`j_2`$), which serves as a quantum reference frame. So far the value of $`j_2`$ has been kept equal to $`1/2`$. We move on to consider what happens as we increase the value of $`j_2`$. The limit $`j_2\mathrm{}`$ could be regarded as the limit in which the quantum reference direction becomes a classical one.
As before, we need to consider states with $`m_1=0`$ and $`m_1=1`$. Different blocks of $`\overline{\rho }(\beta )`$ are found by coupling the representations of $`|\psi _1`$ and $`|\psi _2`$. For the $`m_1=1`$ sector, i.e. the case of parallel spins, $`\overline{\rho }(\beta )`$ has three one-dimensional blocks, since
$$\mathrm{𝟏}𝒋_\mathrm{𝟐}=(𝒋_\mathrm{𝟐}\mathbf{}\mathrm{𝟏})𝒋_\mathrm{𝟐}(𝒋_\mathrm{𝟐}\mathbf{+}\mathrm{𝟏}).$$
The optimal measurement is then given by projections on subspaces of total angular momentum $`J`$.
For the $`m_1=0`$ sector we have
$$(\mathrm{𝟎}\mathrm{𝟏})𝒋_\mathrm{𝟐}=𝒋_\mathrm{𝟐}(𝒋_\mathrm{𝟐}\mathbf{}\mathrm{𝟏})𝒋_\mathrm{𝟐}(𝒋_\mathrm{𝟐}\mathbf{+}\mathrm{𝟏}),$$
thus $`\overline{\rho }(\beta )`$ and the operators $`A_\mu `$ have two $`1`$-dimensional blocks and one $`2`$-dimensional block. The POVM elements acting on the $`J=j_21`$ and $`J=j_2+1`$ subspaces are rank one projectors onto these subspaces. In the $`J=j_2`$ subspace, an optimization procedure similar to the one in Sec. III.1 needs to be done.
Carrying out the optimization, we find that the optimal state $`|\psi _{\mathrm{opt}}`$ belongs to the $`m_1=0`$ sector for all values of $`j_2`$. The optimal state is not fixed but rather depends on the value of $`j_2`$. The limit $`j_2\mathrm{}`$ yields the optimal asymptotic state
$$\underset{j_2\mathrm{}}{lim}|\psi _{\mathrm{opt}}=a_{\mathrm{}}|00+\sqrt{1a_{\mathrm{}}^2}|10,$$
(65)
with $`a_{\mathrm{}}=0.595`$. The dependence of the state $`|\psi _{\mathrm{opt}}`$ on $`j_2`$ is plotted in Fig 2.
Comparing the fidelity achieved with the optimal state, as a function of $`j_2`$, to the fidelity obtained using both the anti-parallel and the parallel spin states, we get, quite remarkably, almost the same fidelity for any value of $`j_2`$. The parallel spins give a slightly higher fidelity, with the difference (already very small for $`j_2=1/2`$) rapidly decreasing for increasing values of $`j_2`$. This comparison is plotted in Fig 3, which shows that the plot for the parallel and anti parallel spin states coincide.
Notice that although we expect the limit $`j_2\mathrm{}`$ to be equivalent to measuring the state $`|\psi _1`$ against a classical reference direction, the anti-parallel spin states do not become optimal in this limit. This seems to be in contradiction with the result of gp ; scudo1 , who showed that if only two spins are available, the optimal direction indicator is provided by anti-parallel spins along that direction. The resolution to this apparent contradiction relies on the fact that the $`j_2\mathrm{}`$ limit considered here is not equivalent to an estimation of a direction with respect to a classical reference frame, as bagan2 ; bagan1 ; scudo2 ; lindner ; dariano ; scudo1 ; mp ; gp , but rather to that of an angle between the same vector, and, say, the $`z`$-axis of such a frame. In the next section we will consider the latter estimation task and show that it coincides with the limit $`j_2\mathrm{}`$ discussed here, demonstrating that a macroscopic spin can be treated as a classical reference direction.
It is conceivable to consider a different estimation task in which one observer, say Alice, indicates a direction in space using the state $`|\psi _1`$, while Bob encodes his reference frame into the state $`|\psi _2`$, and finally a third observer is interested in the orientation of Alice’s direction in Bob’s frame. This task, however, cannot be performed using a spin coherent state $`|\psi _2=|j_2,j_2`$. It would be interesting to see what state would encode Bob’s frame in an optimal manner, and whether taking the appropriate limit would reproduce previous results for direction alignment.
### III.3 Quantum-classical correspondence
Let us consider a scenario in which Alice prepares a state, $`|\psi _1`$, in order to indicate a chosen direction, and Bob is requested to estimate the angle between this direction and the $`z`$-axis of his classical reference frame, which replaces the quantum reference direction ($`|\psi _2`$) in the previous study. We assume no knowledge of Alice’s reference frame, and without loss of generality we can assume that Alice chooses to indicate her $`z`$-axis. If the transformation relating Alice’s frame to Bob’s is parameterized by the Euler angles $`\chi ,\beta `$ and $`\varphi `$, then, in the latter reference frame, Alice’s state is given by $`U(\chi ,\beta ,\varphi )|\psi _1`$. Bob’s task is now to estimate the angle $`\beta `$ between Alice’s $`z`$-axis and his $`z`$-axis. As before, given $`|\psi _1`$, we are seeking a POVM that maximizes the fidelity
$$F\left\{E_\mu \right\}=\underset{\mu }{}d_{\chi \beta \varphi }\mathrm{Tr}\left[\sigma _1(\chi ,\beta ,\varphi )E_\mu \right]f(\mu ,\beta ),$$
(66)
where $`d_{\chi \beta \varphi }=\mathrm{sin}\beta d\beta d\varphi d\chi /8\pi ^2`$ is the invariant measure of the rotation group and
$$\sigma _1(\chi ,\beta ,\varphi )=U(\chi ,\beta ,\varphi )|\psi _1\psi _1|U(\chi ,\beta ,\varphi )^{}.$$
Since we are not interested in estimating the angles $`\varphi `$ and $`\chi `$, we can integrate over them and define an averaged density matrix function of $`\beta `$ only,
$$\overline{\sigma }(\beta )=\frac{1}{4\pi ^2}𝑑\varphi 𝑑\chi \sigma _1(\chi ,\beta ,\varphi ).$$
(67)
Let us analyze the form of the density matrix $`\overline{\sigma }(\beta )`$. Writing $`|\psi _1`$ as in Eq. (14), and using the definition of the rotation operator matrix elements
$$jm^{}|U(\chi ,\beta ,\varphi )|jm=e^{im^{}\chi }d_{m^{}m}^j(\beta )e^{im\varphi },$$
(68)
we derive the matrix elements of $`\overline{\sigma }(\beta )`$ as
$$j^{}m^{}|\overline{\sigma }(\beta )|jm=\underset{r}{}\delta _{m^{}m}a_r^j^{}a_r^jd_{m^{}r}^j^{}(\beta )d_{mr}^j(\beta ).$$
(69)
The matrix $`\overline{\sigma }(\beta )`$ is thus diagonal in the indices $`m`$, but has off-diagonal elements with different values of $`j`$. In this respect it is similar to the matrix $`\overline{\rho }(\beta )`$ defined in Eq. (30), which was diagonal in the representation $`J`$ with off-diagonal elements between different values of $`j_1`$. Indeed, by taking the limit $`j_2\mathrm{}`$ in Eq. (30) and interchanging the indices $`Jm`$, we get the asymptotic equivalence between the matrices $`\overline{\rho }(\beta )`$ and $`\overline{\sigma }(\beta )`$, i.e.,
$$\underset{j_2\mathrm{}}{lim}Jj_1|\overline{\rho }(\beta )|Jj_1^{}=j_1m|\overline{\sigma }(\beta )|j_1^{}m.$$
(70)
This result shows that the fidelity achieved by relating any state $`|\psi _1`$ to a classical reference direction is identical to the one achieved in the limit $`j_2\mathrm{}`$ discussed in the previous section (with the same state $`|\psi _1`$). Consequently, also the optimal state will be identical in the two cases.
## IV Concluding remarks
In this work we studied the problem of estimating relative rotation angles of quantum signals. By considerations of global symmetry invariance, we derived the general form of the signal state and of the appropriate set of measurements and estimation strategies. For special low-dimensional cases, explicit optimization is carried out.
With these tools we have studied the state preparations which maximize the average fidelity of the estimation procedure, and compared these with the results found in the estimation of absolute rotations. We have also discussed the asymptotic limit in which one of the quantum states becomes a macroscopic spin. In this limit, the resulting estimation task is identical to an estimation of the orientation of a quantum state with respect to a classical reference direction.
Many important questions remain open for investigation. A broader extension of the problem would lead to the estimation of the orientation of a quantum state relative to another one that encodes a full reference frame (three axes), rather then a single direction. In such a scenario one has to estimate two angles (polar and azimutal). An even more elaborate framework is one in which each quantum state encodes a full reference frame, and one is interested in estimating the transformation that would align these two reference frames. It would be interesting to compare the optimal states found in all these cases with those found in prior studies scudo2 ; scudo1 ; bagan1 ; bagan2 ; dariano , when one of the reference frames is classical. We would like, finally, to emphasize that we did not include in our communication scheme the possibility of the two parties sharing a prior entangled state. In this case, the encoding and the detection of relative information could proceed via a covariant dense coding scheme, which could increase the efficiency of the estimation.
## Acknowledgements
The Authors acknowledge interesting discussions with Daniel Terno, Sofyan Iblisdir, Terry Rudolph and Giulio Chiribella. Special thanks are due to Rob Spekkens for early contributions to this work. Part of the research was carried out during a visit at the Perimeter Institute, which N.H.L. and P.F.S. gratefully acknowledge for the hospitality. During the completion of this work, we learned about similar results obtained by Nicolas Gisin and Sofyan Iblisdir gi , whom we kindly acknowledge for their correspondence. We also wish to thank Emilio Bagan, Sofyan Iblisdir and Ramon Muñoz-Tapia for constructive criticism of an earlier version of this paper.
Work by N.H.L was supported by a grant from the Technion Graduate School. P.F.S is grateful to the EU (Grant HPRN-CT-2002-002777) and to Prof. Joseph Avron for supporting this work. |
warning/0506/math0506385.html | ar5iv | text | # Ramanujan’s Inverse Elliptic Arc Approximation
## Acknowledgment
Support from the Vicerrectoría de Investigación of the University of Costa Rica is acknowledged. |
warning/0506/cond-mat0506575.html | ar5iv | text | # Unconventional Integer Quantum Hall effect in graphene
## Abstract
Monolayer graphite films, or graphene, have quasiparticle excitations that can be described by $`2+1`$ dimensional Dirac theory. We demonstrate that this produces an unconventional form of the quantized Hall conductivity $`\sigma _{xy}=(2e^2/h)(2n+1)`$ with $`n=0,1,\mathrm{}`$, that notably distinguishes graphene from other materials where the integer quantum Hall effect was observed. This unconventional quantization is caused by the quantum anomaly of the $`n=0`$ Landau level and was discovered in recent experiments on ultrathin graphite films.
The quantum Hall effect (QHE) is one of the most remarkable phenomena in condensed matter discovered in the second half of the 20th century. The basic experimental fact characterizing QHE is that the diagonal electric conductivity of a two-dimensional electron system in a strong magnetic field is vanishingly small $`\sigma _{xx}0`$, while the non-diagonal conductivity is quantized in multiples of $`e^2/h`$: $`\sigma _{xy}=\nu e^2/h`$, where $`\nu `$ is an integer (the integer quantum Hall effect (IQHE)) or a fractional number (the fractional QHE). In a recent paper Novoselov2004Science the fabrication of free-standing monocrystalline graphite films with thickness down to a single atomic layer was reported. This new material, called graphene, possesses truly remarkable properties such as excellent mechanical characteristics, scalability to the nanometer sizes, and the ability to sustain huge ($`>10^8A/cm^2`$) electric currents. By using the electric field effect Novoselov2004Science , it is possible to change the carrier concentration in samples by tens times and even to change the carrier type from electron to hole when the sign of applied gate voltage is altered. All this make graphene a promising candidate for applications in future micro- and nanoelectronics.
On the theoretical side, the linear, Dirac-like, spectrum of quasiparticle excitations (up to energies of the order of $`1000`$ K) and the pseudospin degeneracy make graphene a unique truly two-dimensional ”relativistic” electronic system. The thinnest graphite films can be described by a low-energy (2+1) dimensional effective massless Dirac theory Semenoff1984PRL . Of special interest are the properties of graphene in a magnetic field. The important differences between the Dirac and Schrödinger theories may be observed in thermodynamic and magnetotransport measurements Kopelevich2003PRL ; Novoselov2004Science ; Morozov2005 ; Berger2004JPCB ; Zhang2005PRL . For instance, the phase of de Haas van Alphen and Shubnikov de Haas oscillations for Dirac quasiparticles is shifted Sharapov2004PRB ; Gusynin2005PRB ; Luk'yanchuk2004PRL ; Geim2005 compared to the phase of non-relativistic quasiparticles. Moreover, the Dingle and temperature factors in the amplitude of oscillations explicitly depend on the carrier density in the case of a Dirac-like spectrum Sharapov2004PRB ; Gusynin2005PRB .
Because of the large value of the cyclotron gap, it is expected that the QHE in this material can be observed for much higher temperatures and lower magnetic fields than in conventional semiconductors. Therefore it is naturally to ask whether the fundamental difference between the properties of Landau levels (LL) (see Eqs. (5) and (10) below) in the Dirac and Schrödinger theories can be observed experimentally in the Hall conductivity? The purpose of this letter is to show that the Dirac-like dynamics of graphene results in an unconventional form of the Hall quantization
$$\sigma _{xy}=\frac{2e^2}{h}(2n+1),n=0,1,\mathrm{}$$
(1)
We argue that the quantization rule (1) is caused by the quantum anomaly of the $`n=0`$ LL , i.e. by the fact that it has a twice smaller degeneracy than the levels with $`n>0`$ and its energy does not depend on the magnetic field Gusynin1995PRD . Remarkably this quantization is observed experimentally Geim2005 for ultrathin graphite films which exhibit the behavior expected for ideal 2D graphene.
We begin with the Lagrangian density of noninteracting quasiparticles in a single graphene sheet that in the continuum limit reads Semenoff1984PRL
$$=\underset{\sigma =\pm 1}{}\overline{\mathrm{\Psi }}_\sigma \left[i\gamma ^0(\mathrm{}_ti\mu _\sigma )+iv_F\gamma ^i(\mathrm{}_ii\frac{e}{c}A_i)\right]\mathrm{\Psi }_\sigma ,$$
(2)
where $`\mathrm{\Psi }_\sigma =(\psi _{1\sigma }(t,𝐫),\psi _{2\sigma }(t,𝐫))`$ is the four-component Dirac spinor combined from two spinors $`\psi _{1\sigma },\psi _{2\sigma }`$ \[corresponding to $`𝐊`$ and $`𝐊^{}`$ points of the Fermi surface, respectively\] that describe the Bloch states residing on the two different sublattices of the biparticle hexagonal lattice of the graphene sheet, and $`\sigma =\pm 1`$ is the spin. In Eq. (2) $`\gamma ^\mu `$ with $`\mu =0,1,2`$ are $`4\times 4`$ $`\gamma `$ matrices belonging to a reducible representation in $`2+1`$, $`\overline{\mathrm{\Psi }}_\sigma =\mathrm{\Psi }_\sigma ^{}\gamma _0`$ is the Dirac conjugated spinor, $`e<0`$ is the electron charge, $`v_F`$ is the Fermi velocity. We set $`k_B=1`$, but kept Planck constant $`\mathrm{}=h/2\pi `$.
The external magnetic field $`𝐁`$ is applied perpendicularly to the plane along the positive z axis and the vector potential is taken in the symmetric gauge $`𝐀=(B/2y,B/2x)`$. In contrast to the truly relativistic $`(3+1)`$ case foot1 , the Zeeman interaction term still has to be explicitly added to the Lagrangian (2), because it originates from nonrelativistic many-body theory. This can be done by considering spin splitting $`\mu _\sigma =\mu \sigma \mu _BB`$ of the chemical potential $`\mu `$, where $`\mu _B=e\mathrm{}/(2mc)`$ is the Bohr magneton. However, for realistic values of $`v_F10^5\text{m/s}`$ in graphene the distance between LL is very large compared to the Zeeman splitting Gusynin2005PRB , so that in what follows we will not consider this term and simply multiply all relevant expressions by 2 to count the spin degeneracy. While simple tight-binding calculations made for the hexagonal lattice of a single graphene sheet predict that $`\mu =0`$, the real picture is more complicated and the actual value of $`\mu `$ can be nonzero due to finite doping and/or disorder. Moreover, nonzero and even tunable value of $`\mu `$ \[including the change of the character of carriers, either electrons or holes\] is possible in electric-field doping experiments Novoselov2004Science ; Morozov2005 ; Geim2005 . In our notations $`\mu >0`$ corresponds to electrons and accordingly to the positive gate voltage.
Using the Kubo formalism and modeling the LL by Lorentzians with a constant width $`\mathrm{\Gamma }`$ the following expression for the diagonal conductivity was obtained in Refs. Gorbar2002PRB ; Gusynin2005PRB
$$\sigma _{xx}(B,\mu ,\mathrm{\Gamma })=\frac{2e^2}{\mathrm{}}_{\mathrm{}}^{\mathrm{}}𝑑\omega [n_F^{}(\omega \mu )]𝒜_{xx}(\omega ,B,\mathrm{\Gamma }),$$
(3)
where $`n_F(\omega )=1/[\mathrm{exp}(\omega /T)+1]`$ is the Fermi distribution and the function $`𝒜_{xx}`$ that incorporates the effect of all LL is given by Eq. (11) of Ref. Gusynin2005PRB . Now this result is extended for the Hall conductivity and we derive a general analytical expression for $`\sigma _{xy}(B,\mu ,\mathrm{\Gamma })`$. Gusynin2005 The resulting dependence $`\sigma _{xy}(\mu )`$ is shown in Fig. 1, where one sees that the plateaux of $`\sigma _{xy}`$ follow Eq. (1). This agrees with the latest experimental results Geim2005 and resemble earlier theoretical predictions Zheng2002PRB .
However, to demonstrate result (1) in the most transparent way it is useful to write down a simpler conventional representation Schakel1991PRD ; Gorbar2002PRB for $`\sigma _{xy}`$ obtained in the clean limit $`\mathrm{\Gamma }0`$:
$$\sigma _{xy}=\frac{ec\rho }{B}\frac{e^2\text{sgn}(eB)\text{sgn}\mu }{\pi \mathrm{}}\nu _B.$$
(4)
Here we introduced the filling factor of LL, $`\nu _B=\pi \mathrm{}c|\rho |/|eB|`$ with $`\rho `$ being the carrier imbalance ($`\rho n_en_h`$, where $`n_e`$ and $`n_h`$ are the densities of “electrons” and “holes”, respectively). This filling factor can be represented as a sum over LL
$$M_n=\sqrt{\mathrm{\Delta }^2+2n\mathrm{}v_F^2|eB|/c},n=0,1,\mathrm{}$$
(5)
of the Dirac theory:
$$\begin{array}{cc}\hfill \text{sgn}\mu \nu _B& =\frac{1}{2}[\mathrm{tanh}\frac{\mu +\mathrm{\Delta }}{2T}+\mathrm{tanh}\frac{\mu \mathrm{\Delta }}{2T}\hfill \\ & +2\underset{n=1}{\overset{\mathrm{}}{}}(\mathrm{tanh}\frac{\mu +M_n}{2T}+\mathrm{tanh}\frac{\mu M_n}{2T})],\hfill \end{array}$$
(6)
where we separated out the level with $`n=0`$ because its degeneracy is only half of the degeneracy of the levels with $`n>0`$. To illustrate this rather peculiar property of the Dirac theory in a perspicuous way, we included in $`M_n`$ and $`\nu _B`$ the mass (excitonic gap) $`\mathrm{\Delta }`$ which was discussed recently to explain some experiments Khveshchenko2001PRL ; Gorbar2002PRB . Our consideration of $`\sigma _{xy}`$ is in fact independent of the presence of $`\mathrm{\Delta }`$, so in what follows we set $`\mathrm{\Delta }=0`$. A zero value of $`\mathrm{\Delta }`$ is expected for noninteracting quasiparticles on the hexagonal lattice of graphene.
The first equality in Eq. (4) corresponds to a classical straight line $`\sigma _{12}\nu _B`$. As discussed, for example, in Ref. Hajdu.book , this line emerges from two step function dependences, viz. $`\mu (n)`$ and $`\sigma _{12}(\mu )`$. Indeed, using $`\mathrm{tanh}(\omega /2T)=\mathrm{sgn}(\omega )`$ for $`T0`$, we obtain from Eqs. (4) and (6) (compare with Ref. Schakel1991PRD )
$$\sigma _{xy}=\frac{2e^2\text{sgn}(eB)\text{sgn}\mu }{h}\left(1+2\left[\frac{\mu ^2c}{2\mathrm{}|eB|v_F^2}\right]\right),$$
(7)
where $`[x]`$ denotes the integer part of $`x`$. The usual argumentation (see e.g. Ref. Hajdu.book ) for the occurrence of the IQHE states that in the presence of disorder the dependence of $`\sigma _{12}(\mu )`$ remains the same, while $`\mu (n)`$ becomes a smooth function. The classical (4) and quantum (7) Hall conductivities coincide only for the fillings, $`\nu _B=2n+1`$ (see Fig. 2).
The odd integer rather than integer fillings that produces the quantization rule (1) appears due to the above-mentioned halved degeneracy of the $`n=0`$ LL. Another interesting feature of Eq. (7) (see also Figs. 1 and 2) is that $`\sigma _{xy}=\pm 2e^2/h`$ for the fillings $`\nu _B<1`$ and it crosses 0 only when $`\mu `$ changes sign. On the contrary, in a conventional IQHE $`\sigma _{xy}=0`$ for $`\nu _B<1`$.
Although Eqs. (4) - (7) are obtained in the clean limit and using a simple bare bubble expression for conductivity, our main result (1) is model independent and is only based on the $`n=0`$ level anomaly.
Now we rewrite Eqs. (4)-(6) in terms of the Fermi distribution
$$\begin{array}{cc}\hfill \sigma _{xy}=& \frac{2e^2\mathrm{sgn}(eB)}{h}\underset{n=0}{\overset{\mathrm{}}{}}(2n+1)\hfill \\ \hfill \times & [n_F(M_n\mu )+n_F(M_n\mu )\hfill \\ & n_F(M_{n+1}\mu )n_F(M_{n+1}\mu )]\hfill \end{array}$$
(8)
to compare it with Eq. (18) of Ref. Jonson1984PRB that was obtained for an ideal two-dimensional electron gas
$$\sigma _{xy}=\frac{e^2}{h}\underset{n=0}{\overset{\mathrm{}}{}}(n+1)[n_F(\omega _n^{\mathrm{nonrel}})n_F(\omega _{n+1}^{\mathrm{nonrel}})]$$
(9)
with nonrelativistic spectrum
$$\omega _n^{\mathrm{nonrel}}=\frac{e\mathrm{}B}{mc}\left(n+\frac{1}{2}\right).$$
(10)
There is a commonsense reasoning Dresselhaus2002AP ; Novoselov2004Science that graphene is a two-band \[the first band would corresponds to the electrons with $`\omega _n=M_n\mu `$ and the second band, to the holes with $`\omega _n=M_n\mu `$\], two-valley \[corresponding to $`𝐊`$ and $`𝐊^{}`$ points of graphen’s Fermi surface\] semiconductor with zero gap $`\mathrm{\Delta }`$ between the bands. Accordingly its Hall conductivity can be directly obtained from (9) by summing over all these bands and valleys
$$\begin{array}{cc}\hfill \sigma _{xy}^{\mathrm{semicond}}& =\frac{2e^2\mathrm{sgn}(eB)}{h}\underset{n=0}{\overset{\mathrm{}}{}}2(n+1)\hfill \\ \hfill \times & [n_F(M_n\mu )+n_F(M_n\mu )\hfill \\ & n_F(M_{n+1}\mu )n_F(M_{n+1}\mu )],\hfill \end{array}$$
(11)
where we also counted spin degeneracy. It is easy to see that Eqs. (8) and (11) correspond to two completely different Hall conductivity quantization rules, viz. Eq. (8) which correctly counts the degeneracy of the $`n=0`$ level produces Eq. (1), while the semiconducting analogy (11) leads to
$$\sigma _{xy}^{\mathrm{semicond}}=\frac{4e^2}{h}n,n=0,1,\mathrm{}.$$
(12)
Here we assumed that $`e,B,\mu >0`$. Although previous experimental observations supported the picture based on Eq. (12), the latest experiments made on thin films Geim2005 are in accord with the unconventional Hall quantization (1). This shows that in an applied magnetic field the semiconducting interpretation of graphene’s band structure that led us to Eq. (11) becomes invalid (see also Ref. Ando2005JPSJ ). The drastic difference between Eqs. (1) and (12) is caused by the above-mentioned fact that the lowest LL in Dirac theory is special and has twice smaller degeneracy than the levels with $`n>0`$, because depending on the sign $`eB`$ it is occupied either by electrons or holes, while higher levels contain both electrons and holes Gusynin1995PRD ; Johnson1949PR . In the nonrelativistic theory when the Lande factor $`g2`$ all Landau levels have the same degeneracy foot1 . It turns out that graphene for which the valence and conduction bands intersect in discrete points foot2 , is reasonably well described by the Dirac formalism which naturally embodies the $`n=0`$ level anomaly.
We now consider the phenomenon of quantum magnetic oscillations in graphene which is closely related to the quantization of $`\sigma _{xy}`$ and discuss the specific of the $`n=0`$ level. The de Haas van Alphen and Shubnikov de Haas effects in graphene were studied in Refs. Sharapov2004PRB ; Luk'yanchuk2004PRL ; Gusynin2005PRB . In particular, in Ref. Gusynin2005PRB it was shown that the oscillatory part of the diagonal conductivity (3) is given by
$$\sigma _{xx}\underset{k=1}{\overset{\mathrm{}}{}}\mathrm{cos}\left[\frac{\pi k\mu ^2}{\mathrm{}v_F^2|eB|/c}\right]R_T(k)R_D(k)R_s(k),$$
(13)
where $`R_T`$, $`R_D`$ and $`R_s`$ are respectively the temperature, Dingle and spin factors. Using the relationship $`\mu ^2=\pi \mathrm{}^2v_F^2|\rho |`$ valid for $`T=\mathrm{\Gamma }=B=0`$ Gorbar2002PRB one can check that the minima of the diagonal conductivity (3) occur at the fillings $`\nu _B=2n+1`$ giving an indication of the possible positions of the plateaux in the IQHE Geim2005 . \[Note that in thick films the minima of $`\sigma _{xx}`$ occur at integer fillings Morozov2005 .\] Obviously for $`\mu =0`$ there is no oscillations of $`\sigma _{xx}`$, the conductivity $`\sigma _{xx}(\mu =0)=2e^2/(\pi ^2\mathrm{})`$ becomes a field independent universal foot3 quantity that is another distinctive feature of the $`n=0`$ level anomaly.
Although the quantization (1) can be understood by considering noninteracting Dirac quasiparticles placed in an external magnetic field, even this simple model reveals other unusual properties Gusynin1995PRD intimately related to nontrivial dynamics of quasiparticles from the $`n=0`$ level. For example, the $`U(4)`$ symmetry of the Lagrangian (2) is spontaneously broken down to $`U(2)\times U(2)`$ at $`\mu =0`$ in non-zero magnetic field even in the absence of additional interaction between fermions Gusynin1995PRD thus leading to the emergence of the chiral condensate $`\overline{\mathrm{\Psi }}\mathrm{\Psi }`$. Including many body effects such as an attractive interaction between quasiparticles, could further generate a gap for quasiparticles like the above mentioned gap $`\mathrm{\Delta }`$ (see e.g. Refs. Gorbar2002PRB ; Khveshchenko2001PRL ). Fortunately in the case of the IQHE the presence of the condensate does not affect our consideration. On the other hand, a possible gap generation for the fermions from the lowest LL might become important for the fractional quantum Hall effect and this issue certainly deserves further experimental and theoretical study.
To conclude, we have shown that the integer numbers associated with quantized Hall conductivity in graphene have an unusual pattern $`\sigma _{xy}h/e^2=2,6,10,14\mathrm{}`$. We argued that it is related to the fact that a theoretical description of graphene is based on $`2+1`$ dimensional Dirac theory, where the lowest Landau level has half of the higher Landau levels degeneracy.
We are indebted to A. Geim for showing us his experimental results prior to publication and for stimulating discussions. We also thank J.P. Carbotte, V.M. Loktev and V.A. Miransky for useful discussions and W.A. de Heer, P. Kim for informing us about their latest results. S.G.Sh. was supported by the Natural Science and Engineering Council of Canada (NSERC) and by the Canadian Institute for Advanced Research (CIAR). |
warning/0506/astro-ph0506621.html | ar5iv | text | # Discovery of a widely separated binary system of very low mass stars Based on observations made at the European Southern Observatory, La Silla, Chile
## 1 Introduction
Binary systems have been studied for decades to measure accurate stellar masses, and to test evolutionary models and star formation theories. Considerable attention has recently been paid to very low-mass (VLM) companions to low-mass stars (Duquennoy & Mayor duquennoy91 (1991); Fischer & Marcy fischer (1992); Delfosse et al. delfosse99 (1999); Reid et al. reid01 (2001); Beuzit et al. beuzit04 (2004); Forveille et al. forveille04 (2004)), as well as to binaries among ultracool dwarfs (spectral types later than M6) in the solar neighborhood (Martín et al. 1999a ; Close et al. close02 (2002), close03 (2003); Bouy et al. bouy (2003); Burgasser et al. burgasser03 (2003); Gizis et al. gizis03 (2003); Siegler et al. siegler (2005); Forveille et al. forveille05 (2005)) and in nearby young open clusters and associations (Martín et al. martin98 (1998), martin00 (2000), martin03 (2003); Chauvin et al. chauvin (2004)). The main results of the high spatial resolution imaging surveys of VLM stars and brown dwarfs (BDs) can be summarized as follows: (1) the binary frequency in the separation range 1-15 AU is about 15% , (2) the frequency of wide binary systems (semi-major axis $`>`$ 15 AU) is very low, $`<`$1% , (3) the mass ratios are strongly biased towards nearly-equal mass binaries, beyond the expected observation selection effects. The companions to the more massive low mass stars, by contrast, span a wider range of both separations and mass ratios.
The properties of VLM binaries are an important constraint for models of star-formation and evolution. It has been debated in the literature whether the properties of VLM binaries and stellar binaries differ, implying different formation mechanisms (Kroupa et al. kroupa03 (2003)), or whether the binary properties instead show continuous trends with decreasing primary mass, implying that VLM binaries form through the same processes as stellar binaries (Luhman 2004b ). Clearly there is a need for a larger sample of observed VLM binaries, particularly at wide separations where few of them are known. One leading model of BD formation is that they form and are ejected in unstable multiple systems within small clusters (Bate et al. bate02 (2002), bate03 (2003)). Numerical simulations of decaying N-body clusters indicate that the typical separations of binaries composed of stars with masses ranging between 0.1 $`M_{}`$ and 0.2 $`M_{}`$ would be below 10 AU (Sterzik & Durisen sterzik (1998)), in rough agreement with the observations. The ejection models (Reipurth & Clarke reipurth (2001); Bate et al. bate02 (2002)) suggest that the binary BD systems that do exist must be close (separations $``$ 10 AU). The detection of wide VLM binary systems has thus become an important test of the ejection models. The first wide binary BDs have been found in young ($`<`$10 Myr) associations or clusters: 2MASS J1101-7732 (240 AU separation, in Chamaeleon I, Luhman 2004a ); 2MASS J1207-3932 (55 AU, in TW Hydrae, Chauvin et al. chauvin (2004)). Their existence is at first sight inconsistent with the ejection models, though the statistics of the numerical simulations is currently limited.
In this Letter, we present a new wide binary consisting of two VLM stars in the field. Sec. 2 presents the observations and data reduction, while Sec. 3. discusses the results in the context of the binary properties of VLM stars and BDs.
## 2 Observations and data reduction
Phan-Bao et al. (phan-bao03 (2003)) identified LP 714-37 (DENIS-P J0410-1251) as a mid-M dwarf while cross-identifying the NLTT (Luyten luytenb (1980)) and DENIS (Epchtein epchtein97 (1997)) catalogues, and derived a photometric distance of 15.3 pc. Cruz et al. (cruz03 (2003)) independently determined a spectrophotometric distance of 15.4 pc and an M6.0 spectral type. Allowing for the binary nature of the system pushes its distance out to 18.1 pc (Sec. 3).
As part of our ongoing spectroscopic follow-up of new nearby M dwarfs detected with DENIS, one of us (Céline Reylé) observed LP 714-37 with the 3.6-meter New Technology Telescope (NTT) at La Silla Observatory (ESO, Chile) in November 2003; and the data was analyzed in January 2005. Two objects were in fact detected on the acquisition image (Fig. 1).
Optical low-resolution spectra were obtained for both objects in the Red Imaging and Low-dispersion spectroscopy (RILD) observing mode of the EMMI instrument. The detector was a 2048$`\times `$4096 MIT/LL CCD, used in normal readout mode with 2$`\times `$2 binning. Grism#2 gave a dipersion of 3.5 Å pixel<sup>-1</sup>, and the OG530#645 order sorting filter defined a wavelength range of 520 to 950 nm. The slit width of 1 arsec corresponds to a spectral resolution of 10.4 Å, and the seeing was 0.7 arcsec. The exposure time was 30 s. LTT 2415 and Feige 110 (selected from the ESO list, ftp.eso.org/pub/stecf/standards/ctiostan/) were observed as spectrophotometric standards. The reduction of spectra was performed with MIDAS packages. We normalized the spectra to 1 over the 754-758 nm interval, the denominator of the PC3 index (Martín et al. 1999b ) and a region with a good flat pseudo-continuum. Figure 2 shows the two resulting spectra. The presence of the NaI and KI doublets and the absence of the CaII triplet immediately demonstrate that both stars are M dwarfs rather than M giants, as does the general shape of the TiO bands. Analysis of the R band acquisition image with SExtractor (Bertin & Arnouts bertin (1996)) gives a separation of 5.53 pixels, or 1.83″ with the 0.33″/pixel focal scale of EMMI. Since the discovery is serendipitous we did not obtain images with other filters and have no photometric calibrations for the R band. We therefore only use the image for astrometry.
## 3 Discussion
The calibration of the PC3 index to spectral type (Martín et al. 1999b ) gives spectral types of M5.5 for component A and M7.5 for component B, with an uncertainty of $`\pm `$0.5 subclass. The observed H<sub>α</sub> emission in component B with spectral type of M7.5 is consistent with the observations by Gizis et al. (2000a ) pointed out that all M7-M8 dwarfs show the H<sub>α</sub> line in emission.
The proper motion of the system is $`\mu _\alpha `$ = $``$117 mas/yr and $`\mu _\delta `$ = $``$382 mas/yr (Phan-Bao et al. phan-bao03 (2003)), and it has therefore moved by 6.8 arcsec between the epoch of the SERC-I plate and our NTT observation. Figure 3 shows that there is no background star at the position of the system in either the SERC-I image ($`I`$ band) or the DENIS-I image, and therefore demonstrates that the system is a physical binary.
To estimate the distance, we assume that this system has no additional unresolved component(s). We can then calculate the absolute magnitudes of the system from those of the two components, computed from their PC3 indices with the PC3 index to magnitudes relation of Crifo et al. (crifo (2005)) (Table 1). Comparison with the DENIS apparent magnitudes ($`I=12.99`$, $`J=10.94`$, $`K=9.89`$) then gives the distance: $`d_\mathrm{I}=17.8`$$`\pm 2.1`$ pc; $`d_\mathrm{J}=18.9`$$`\pm 2.3`$ pc and $`d_\mathrm{K}=17.6`$$`\pm 2.1`$ pc. We adopt the mean distance of $`d=18.1`$$`\pm 2.2`$ pc.
We estimate the mass of each component, using $`I`$, $`J`$, and $`K`$-band mass-luminosity relations (Fig. 4 for the $`K`$ band) from the Baraffe et al. (baraffe98 (1998)) models. We adopt their 5 Gyr model, an intermediate age for the solar neighbourhood (Caloi et al. caloi (1999)), but the results are insensitive to this choice (Fig. 4). The masses of component A and B are respectively 0.11$`\pm `$0.01 M and 0.09$`\pm `$0.005 M, slightly above the hydrogen-burning limit (Chabrier & Baraffe chabrier (1997)). The total mass of LP 714-37AB is 0.2$`\pm 0.015`$ M, within one sigma of the $`M_{\mathrm{tot}}<0.185M_{}`$ convention adopted by Close et al. (close02 (2002)) for their sample of VLM binaries. Thus, this system can be considered a VLM binary. At the 18.1 pc of the system, its 1.83″ separation corresponds to 33.1 AU. Assuming a face-on circular orbit, the corresponding orbital period is approximately 400 years.
Recent surveys demonstrate that VLM binaries with large separations ($`>`$ 30 AU) are rare in the field (Fig. 5, and references therein), but can be found in young associations and clusters (Chauvin et al. chauvin (2004), Luhman 2004a ). Those authors reported the respective discoveries of 2MASSW J1207-3932 (8 Myr, TW Hydrae), and 2MASSW J1101-7732 (1 Myr, Chamaeleon I) with separations of 55 AU and 240 AU. We report here the discovery of a 33 AU VLM binary, and Martín et al. (martin00 (2000)) found that CFHT-Pl-18 is a 35 AU VLM binary (Fig. 5), both in the field.
Since ejection models predict a maximum separation of about 10 AU for VLM binaries, the existence of the wide binaries is at first sight inconsistent with these ejection models. One should note however that the numerical models to date suffer from small number statistics. A further caveat if that the relevant quantity is the total mass of the system, and that the apparent binaries could possibly be triple or higher order multiple systems, with a correspondingly higher total mass. This would make them analogs of the GJ 1245ABC triple system, which consists of two M5.5 and one $``$M8 dwarfs (Reid et al. reid95 (1995), Henry et al. henry99 (1999)) with separations of 32 and 5 AUs (McCarthy et al. mccarthy88 (1988)). We note that Gizis et al. (2000b ) discovered a 230 AU VLM binary in the field, LP 213-67 (M6.5) and LP 213-68 (M8.0), which in the recent adaptive optics survey of Close et al. (close03 (2003)) turned out to be triple when LP 213-67AB (or 2MASS J1047+4026) was resolved into two components (see Fig. 5). Bouy et al. (bouy05 (2005)) further reported that the DENIS-P J020529.0-115925 VLM binary is a probable triple system. Triple systems could thus potentially explain the apparent excess of wide VLM binaries, and adaptive optics imaging of LP 714-37 would be of obvious interest to clarify its true multiplicity.
Amongst the recently discovered VLM field binaries (Bouy et al. 2003; Close et al. close02 (2002), close03 (2003); Siegler et al. siegler (2005); Forveille et al. forveille05 (2005)), LP 714-37 has one of the widest separation, making it of great interest as a constraint for VLM binary star formation theories. We note that late-M dwarfs detected by Phan-Bao et al. (phan-bao01 (2001), phan-bao03 (2003)) over 5700 square degrees in the DENIS database provide a valuable well defined sample for studies of VLM field binaries.
###### Acknowledgements.
Ngoc Phan-Bao is grateful to the DENIS consortium for access to the DENIS data used by his very low mass stars search. Ngoc Phan-Bao also thanks Françoise Crifo for her comments on the manuscript. Celine Reylé acknowledges help during the observations by Olivier Hainaut and the NTT team at the European Southern Observatory. We thank the referee, Kevin Luhman, for a fast and useful report. Partial funding was provided by NSF grant AST 02-05862. |
warning/0506/cond-mat0506504.html | ar5iv | text | # Sequence of multipolar transitions: Scenarios for URu2Si2
## 1 Introduction
There is increasing interest in orbital ordering phenomena, and their relationship to magnetism. Many of these phenomena can be at least partially understood in terms of localized electron models. This is clearly justified for Mott-localized $`d`$-electrons. Some $`f`$-electron systems are semiconductors (like NpO<sub>2</sub>), but more often, we wish to describe multipole ordering in metals like rare-earth-filled skutterudites or URu<sub>2</sub>Si<sub>2</sub>. In the standard model of the majority of rare earth elements the $`f`$-shell is localized, and we may invoke a similar feature for the systems of interest to us. It can be argued that itineracy and multipolar ordering are complementary features, and that they may be manifest in different phases of the same $`f`$-electron system. In this case, the localized description is acceptable for the ordered phases, even though we know that it could not be extended to the whole phase diagram.
The common starting point is the existence of an $`N`$-dimensional local Hilbert space $`|n`$, which allows the definition of $`N^2`$ local operators $`|nm|`$ (of these, $`N^21`$ are non-trivial, and are also called local order parameters). $`N(N+1)/2`$ operators have symmetrical, and $`N(N1)/2`$ operators have antisymmetrical character. They can be chosen as
$$X_{mn}^s=\frac{1}{2}(|mn|+|nm|)$$
(1)
for symmetrical and
$$X_{mn}^a=\frac{1}{2i}(|mn||nm|)$$
(2)
for antisymmetrical operators where $`i`$ was inserted to ensure hermiticity.
There are two canonical cases: the basis may consist of
1. N/2 time-reversed pairs (the N-fold degeneracy arises as 2$`\times `$(N/2),
(Kramers degeneracy)$`\times `$(non-Kramers degeneracy)
2. real basis states each of which is time reversal invariant (the N-fold degeneracy is purely non-Kramers degeneracy)
Case I. The local Hilbert space consists of $`(N/2)`$ time-reversed pairs. There are $`N(N1)/2`$ time-reversal-even order parameters (this includes the trivial $`\widehat{1}`$), and $`N(N+1)/2`$ time-reversal-odd order parameters.
$`N=4`$ is realized by the $`\mathrm{\Gamma }_8`$ representation of the cubic double group, the irrep of the ground state level for either the $`4f^1`$ compound CeB<sub>6</sub> (Refs. ), or the $`5f^3`$ compound NpO<sub>2</sub> (Ref. ). Six operators: $`\widehat{1}`$ and the five quadrupoles are time-reversal-even, while ten order parameters: the three dipoles and seven magnetic octupoles, change sign under time reversal. A general form of intersite interaction is a sum over the quadratic invariants, for a cubic system with six independent coupling constants (one dipolar, two quadrupolar, three octupolar). It was, however, soon realized that the consideration of the SU(4) symmetrical model (with all couplings set equal) should be enlightening. Later research showed that CeB<sub>6</sub> can be regarded as a ”nearly SU(4)-symmetrical” system.
Case II. If all basis states can be chosen real, the $`X_{m,n}^s`$ are real operators and the $`X_{m,n}^a`$ are imaginary operators. Choosing one of the real operators as the projection onto the entire local Hilbert space $`\widehat{1}=_m|mm|`$, and orthogonalizing the remaining $`N1`$ diagonal operators, we are left with $`N(N+1)/21`$ time-reversal invariant local order parameters. $`N(N1)/2`$ local order parameters (the imaginary operators $`X_{m,n}^a`$) change sign under time reversal.
This case is realized for non-Kramers ions which have an even number of electrons. We are interested in the $`5f^2`$ (U<sup>4+</sup>), and $`4f^2`$ (Pr<sup>3+</sup>) configurations as they appear in URu<sub>2</sub>Si<sub>2</sub> and PrFe<sub>4</sub>P<sub>12</sub>, resp. It is always debatable which $`N`$ to choose. Crystal field doublets ($`N=2`$) have been suggested for both systems but it is unlikely for PrFe<sub>4</sub>P<sub>12</sub>, and at least not widely accepted for URu<sub>2</sub>Si<sub>2</sub>. In fact, for both systems we settle for a local Hilbert space which is composed of the bases of several irreps. For PrFe<sub>4</sub>P<sub>12</sub> the quasi-quartet $`\mathrm{\Gamma }_1+\mathrm{\Gamma }_4`$ seems to work. Under tetrahedral symmetry, the nine time-reversal-even order parameters are a singlet, five quadrupoles, and three hexadecapoles, while the six time-reversal-odd order parameters are three dipoles and three octupoles.
### 1.1 The nature of the higher multipoles
While dipoles and quadrupoles are well-known from decades of research experience, it is only recently that octupoles were seriously considered, and higher multipoles are virtually never discussed. It is worth forming a picturial idea of octopule-, hexadecapole-, and triakontadipole-carrying shells.
The $`J=4`$ Hund’s rule two-electron states of U$`{}_{}{}^{4+}=5f^2`$ ionic cores are rather complicated objects, superposed of a number of $`l_js_j`$ components by the Clebsch-Gordan coefficients. However, there is strong similarity between the multipoles composed of two $`f`$-orbitals and two spins via a number of projections, and simplified multipoles built of $`l=4`$ orbitals (atomic $`g`$-states). We derive the charge cloud shapes and current distribution patterns from acting with the analogous purely orbital (e.g., $`\overline{J_z(J_x^2J_y^2)}\overline{l_z(l_x^2l_y^2)}`$) operators on fictitious atomic states (for our purposes, fourth-order spherical harmonics $`Y_m^4(\vartheta ,\varphi )`$ for $`g`$-states). In this Section, we do not discuss crystal field effects.
We diagonalize multipole operators in the 9-dimensional $`l=4`$ space, and base our pictures on the state with the (in absolute value) largest eigenvalue. This can be interpreted as finding the ground state of an effective field which orders the multipole. We mention that free-ion multipole eigenstates are not necessarily the same that we find within the restricted Hilbert spaces of crystal field problems. This will be illustrated by comparing Fig. 2 and Fig. 5. Furthermore, smaller Hilbert spaces may not give a representation of certain multipoles at all.
For odd-rank (magnetic) multipoles, the eigenstates appear in time-reversed pairs, similar to the $`l_z=\pm m`$ dipole eigenstates, and the multipole spectrum is symmetrical about 0. For even-rank (electric) multipoles, each eigenstate can be chosen time-reversal invariant.
#### 1.1.1 Octupoles
First, let us consider the octupole
$$𝒯_{xyz}\frac{1}{6}\overline{l_xl_yl_z}=\frac{i}{4}\left(l_+l_zl_{}l_{}l_zl_{}\right)$$
(3)
The simplest octupolar states are realized in the $`l=2`$ subspace of $`e_g`$ functions
$$\mathrm{\Phi }_\pm (l=2)=\frac{1}{\sqrt{2}}\left(|3z^2r^2\pm i|x^2y^2\right)$$
(4)
$`\mathrm{\Phi }_\pm (l=2)`$ form a time-reversed pair.
The charge and the current distribution of $`\mathrm{\Phi }_+(l=2)`$ is shown in Fig. 1. The object got its name from the eight magnetic poles: of the eight current eddies we see in the figure<sup>1</sup><sup>1</sup>1Two general remarks about Figures 15: charge density angular dependences are shown, i.e., the lobe shapes do not include the radial fall-off of the atomic wave functions. Flow lines are to indicate the sense of circulation of the current but (in contrast to the textbook interpretation) the density of flow lines is not associated with higher current density, but is rather arbitrary. Flow lines were calculated by solving the differential equation for the tangential curves for the calculated vector field of currents, and initial values were randomly generated. (We found that the direct representation of the vectors would give unattractive figures)., four belong to magnetic field lines entering, and four to those leaving the surface. Note that the magnetic field pattern measured by $`\mu `$SR in NpO<sub>2</sub> bears an overall similarity to what we expect in the neighborhood of an octupole-moment bearing shell. For the time-reversed partner $`\mathrm{\Phi }_{}(l=2)`$ we would find the same charge cloud with reversed currents.
The concept of octupolar ordering was pioneered by Korovin and Kudinov who envisaged an antiferro-octupole pattern of the $`e_g`$ states (4) resulting from spin-orbital exchange in Mott insulators. A similar possibility arises in the $`E`$ doublets of trigonal compounds where, however, the octupole moment is mixed with the orbital moment $`l_z`$. In any case, antiferro-orbital order is likely to be combined with spin ferromagnetism. Metallic phases, including the itinerant octupolar phase, were investigated for the $`e_g`$ band Hubbard model by Takahashi and Shiba.
Currently known realizations of octupolar order appear in $`f`$-electron systems. Field-induced octupoles play a role in understanding the phase diagram of CeB<sub>6</sub> (Ref. ). Kusunose and Kuramoto called attention to the fact that octupole moments are ideally suited for the role of ”hidden” primary order parameters. A detailed study by Kubo and Kuramoto makes a convincing case that the ”Phase IV” of Ce<sub>1-x</sub>La<sub>x</sub>B<sub>6</sub> is an antiferro-octupolar phase. At about the same, it became accepted that the long-standing mystery of the nature of the 25K transition of NpO<sub>2</sub> is solved by identifying it with the triple-q ordering of $`\mathrm{\Gamma }_{5u}`$ octupoles. A symmetry analysis was successful in identifying the unique octupolar signature in NMR spectra on NpO<sub>2</sub>, and Ce<sub>1-x</sub>La<sub>x</sub>B<sub>6</sub> (Ref. . The phenomenology of the two systems shows certain similarities, as it is to be expected since, e.g., the relationship between the anomalies of the linear, and the third-order, susceptibilities follow from general thermodynamic reasoning. The same should be true of URu<sub>2</sub>Si<sub>2</sub> whos hidden order, we argued, is also of octupolar nature.
Fig. 1 illustrates the symmetry of the (uniform) octupolar ground state with $`𝒯_{xyz}`$ as the order parameter. We leave to Sec. 3.2 the construction of an octupolar symmetry group, which will be carried out for a combination of tetragonal crystal field and octupolar effective field, the case relevant for URu<sub>2</sub>Si<sub>2</sub>. Here we use Fig. 1 to visualize the hybrid nature of some of the symmetry elements. The charge distribution is highly symmetrical (octahedral). However, the true symmetry is that of the magnetic field pattern, so in purely geometrical terms it is tetrahedral. It is obvious, however, that $`𝒞_4`$ rotations are allowed if they are combined with time reversal $`𝒯`$ (which reverses the currents and the field). Including elements like $`𝒞_4𝒯`$, a higher non-unitary symmetry group could be derived.
Figure. 1 shows why octupolar order can be so well ”hidden”. The charge distribution (which influences directional bonding, i.e., the structure of the crystal) is cubic, hiding the fact that if current directions are considered, the pure $`𝒞_4`$ axes are not symmetry elements. Similar considerations arose when it was asked whether the pseudocubic phase of Sr-doped LaMnO<sub>3</sub> is not, in fact, octupolar ordered.
One might have thought that the character of the octupoles is the same whatever Hilbert space we represent them on, but there are interesting nuances here.
Within the $`l=4`$ shell (which mimicks<sup>2</sup><sup>2</sup>2However, a calculation of the currents for realistic two-electron states would be desirable. URu<sub>2</sub>Si<sub>2</sub>’s $`5f^2J=4`$ multiplet), $`𝒯_{xyz}`$ has the minimal eigenvalue $`6\sqrt{3}`$, corresponding to the eigenstate (expressed in the $`|l_z`$ basis)
$$\mathrm{\Phi }_0^{\mathrm{oct}}=\sqrt{\frac{7}{50}}(|+4+|4)2i\sqrt{\frac{3}{50}}(|+2+|2)\sqrt{\frac{1}{5}}|0$$
(5)
(the maximal eigenvalue $`6\sqrt{3}`$ corresponds to the time-reversed of $`\mathrm{\Phi }_0^{\mathrm{oct}}`$). The current distribution for $`\mathrm{\Phi }_0^{\mathrm{oct}}`$ is shown in Fig. 2. The overall pattern of eight alternating vortices is the same as in Fig. 1. However, within each of the major eddies four new sub-eddies appear: three (arranged like the petals of a flower) rotating in the sense of the eddy as a whole, and a little central eddy counter-rotating. The symmetry of the current distribution in Fig. 2 is the same as that in Fig. 1, but there are differences in the magnetic field pattern. While in Fig. 1 the magnetic field is maximum in the center of an octant, and has the same sign everywhere within the octant, in the $`l=4`$ solution (Fig. 2) the field changes sign in a small central region of the octant and appears smaller than in the ”petals”. On the whole, the magnetic field of the octupolar currents is weaker for $`l=4`$ states ($`f^2`$ configurations) than for the simplest $`l=2`$ solution. This may account for the weakness of the internal fields in URu<sub>2</sub>Si<sub>2</sub>. It should be noted, though, that the current distribution shown in Fig. 2 was obtained for a free ion subject to the octupolar effective field only. The situation changes if crystal field effects are included (Sec. 3.2).
It would be of obvious interest to calculate internal field distributions for octupoles, and maybe also for triakontadipoles, for these predictions could be checked against $`\mu `$SR and NMR observations. We note that Kubo and Kuramoto discussed the situation for the two well-established octupolar systems NpO<sub>2</sub> and Ce<sub>1-x</sub>La<sub>x</sub>B<sub>6</sub>. for which 500G and 40G, respectively, are measured; both are in excess of the theoretical estimate. The measured internal field in URu<sub>2</sub>Si<sub>2</sub> is much weaker: <sup>29</sup>Si-NMR linewidth gives $`12`$G (Ref.), while at the muon stopping sites, the internal field is at most 1–2G (Ref.).
#### 1.1.2 Hexadecapoles
The hexadecapole $`\overline{l_xl_y(l_x^2l_y^2)}`$ has the largest eigenvalue for the eigenstate
$$\mathrm{\Phi }_{\mathrm{hex}}=\frac{1}{\sqrt{2}}(|2i|2).$$
(6)
In spite of the appearance of $`i`$, $`\mathrm{\Phi }_{\mathrm{hex}}`$ is time reversal invariant. It can be chosen real and its positive and negative lobes can be identified (Fig. 3). There is a general similarity to quadrupolar eigenstates (which also have positive and negative lobes), only the number of lobes is higher. We do not construct the symmetry group of the hexadecapole effective field but we notice that, like the octupolar symmetry group, it would have composite symmetry elements: purely geometrical transformations combined with ”lobe sign reversal”.
Hexadecapoles as order parameters were discussed within the tetrahedral symmetry classification, valid for Pr-filled skutterudites. We are not aware of the existence of primary hexadecapole order in any system.
#### 1.1.3 Triakontadipoles
There are 11 triakontadipoles but we consider only $`\overline{l_xl_yl_z(l_x^2l_y^2)}`$ which appears as the lowest-rank $`A_{1u}`$ multipole in the tetragonal classification (Table 1). Its ground state is
$$\mathrm{\Phi }_{\mathrm{tria}}=\frac{1}{2}\left(|4+|4+i\sqrt{2}|0\right).$$
(7)
with eigenvalue $`6\sqrt{35}`$. As shown in Fig. 4, there are 32 cells of alternatingly flowing current, so the associated magnetic field is more short-ranged than in the case of an octupole (Fig. 1). As in the case of octupoles, the symmetry of the charge distribution is higher than that of the current distribution: the charge cloud has a $`𝒞_8`$ axis which is reduced to $`𝒞_4`$ for the currents. However, the combination of the $`\pi /4`$ rotation with time reversal $`𝒞_8𝒯`$ is a symmetry operation. It would be interesting to meet triakontadipolar order in nature.
## 2 Recent experimental developments on URu<sub>2</sub>Si<sub>2</sub>
The magnetic behavior and phase diagram of the intermetallic compound URu<sub>2</sub>Si<sub>2</sub> have been the focus of attention for over two decades. Specific heat measurements show that electronic entropy of $`O(\mathrm{ln}2)k_\mathrm{B}`$ is released by the time the temperature reaches 30K, and a sizeable fraction of it is associated with the $`\lambda `$-anomaly at $`T_017\mathrm{K}`$. Thus it is not unjustified to think of the phase transition as the full-scale ordering of a localized degree of freedom, but the nature of the order parameter remains hidden. It is obviously not the tiny ($`0.03\mu _\mathrm{B}`$) antiferromagnetic moment which is observed by neutron scattering .
Small (either static or slowly fluctuating) moments have long been held to be an attribute of heavy fermion systems on the borderline between localized and itinerant $`f`$-electron phases. The specific heat value would allow to classify URu<sub>2</sub>Si<sub>2</sub> as a ”light heavy fermion system”, raising the question whether its mysterious properties may be related to an exotic itinerant phase of strongly correlated $`f`$-electrons. Orbital magnetism due to plaquette currents, unconventional density waves, and Pomeranchuk instability leading to a nematic state are in this category.
The RKKY interaction mediates a variety of multipole-multipole interactions between the $`f`$-shells. However, the ordering may be foiled by the collective Kondo effect: the formation of a heavy Fermi sea with a large (Luttinger) Fermi surface may swallow up the localized moments. The Kondo-to-RKKY transition has been extensively studied for the case when the relevant local degrees of freedom are spin dipoles but work on the general multipolar problem has only just started. It is an intriguing possibility that the 17K phase transition in URu<sub>2</sub>Si<sub>2</sub> coincides with the itinerant-to-localized transition of the $`f`$-electrons. This would lend credence to describing the $`T<T_0`$ order in terms of a localized $`f`$-electron model even if a more satisfactory description will have to encompass itinerant aspects.
It has long been known that the symmetry classification of local order parameters in tetragonal crystal fields follows Table 1. For each of the irreps, only the lowest-rank multipole is listed. The result can be viewed as arising from the tetragonal splitting of the operator irreps of the cubic classification scheme.
Progress with uncovering the true nature of the hidden order of URu<sub>2</sub>Si<sub>2</sub> has been particularly slow because of a seemingly extrinsic property of almost all of the samples: they show a kind of micro-antiferromagnetism with moments $`m0.020.04\mu _\mathrm{B}\widehat{z}`$ with the simple alternating order $`𝐐=(0,0,1)`$. However, evidence from microscopic measurements shows that the nominally small moment should be understood as a relatively large moment at a minority of the sites.
The relationship between the hidden order parameter $`\psi `$ and the antiferromagnetic moment $`m`$ has been, and still remains, a matter of debate. There are two basic possibilities: A) hidden order and antiferromagnetism share the same symmetry, or B) they are of different symmetries. In Case A), the relative amplitude of $`m`$ and $`\psi `$ can be tuned continuously (e.g., by pressure), and there is no sharp distinction between the two orders. In Case B), the two orders are incompatible, and there must be a first order transition from the $`\psi 0`$, $`m=0`$ phase to the $`\psi =0`$, $`m0`$ phase. In the present discussion (as in Ref.) we take the conclusion drawn from recent $`\mu `$SR experiments as our starting point: there is a first-order transition, the symmetry of $`\psi `$ is different from that of $`m`$, and therefore the hidden order parameter must be sought from among other entries in Table 1.
As far as gross features like linear and non-linear susceptibility, specific heat, etc. are concerned, with a little adjustment of the crystal field level scheme, $`B_{1g}`$ and $`B_{2g}`$ quadrupolar , and $`B_{1u}`$ and $`B_{2u}`$ octupolar models do about equally well. Evidence beyond this simple range of experiments has to be invoked to choose between the quadrupolar and octupolar scenarios.
We cite two crucial (but as yet unpublished) experiments to argue that the hidden order is not quadrupolar. First, there is a remarkable mechanical–magnetic cross–effect (at $`T<17`$K): in the presence of uniaxial stress applied perpendicularly to the tetragonal main axis, large-moment antiferromagnetism with the simple pattern described above becomes visible for neutron magnetic Bragg scattering . This is unambiguous proof that the background order breaks time reversal invariance, and is thus certainly not quadrupolar; the simplest remaining choice is octupolar order.
The second evidence is coming from recent NQR measurements by Saitoh et al . The temperature dependence of the electric field gradients was followed carefully from relatively high temperatures (70K) to well below $`T_0=17`$K. $`V_{zz}`$ changes all the time, reflecting the $`T`$-dependence of the tetragonal crystal field component $`𝒪_2^0`$. Saitoh et al’s findings can be formulated in two statements: (a) Neither of the field gradients shows anomalies at $`T_0=17`$K, the onset temperature of hidden order, so the hidden order cannot be quadrupolar. (b) There is an anomaly in the NQR signal at $`T^{}13.5`$K, so something happens to the quadrupoles there. A possibility is that $`T^{}`$ is the ordering transition of quadrupoles; then it has to be a second HO transition following the first one at $`T_0`$. We emphasize that a lot more experimental evidence (especially specific heat, susceptibility, etc) is needed before the existence of a phase transition at $`T^{}`$ may become accepted. It is not our aim to describe the two transitions in any detail. We merely emphasize that, given the symmetry of URu<sub>2</sub>Si<sub>2</sub>, octupolar ordering can be followed by quadrupolar ordering; but then further induced magnetic multipoles should be observable.
While circumstantial evidence for octupolar ordering looks encouraging, attempts for its direct verification have as yet yielded negative results. An early, very specific neutron scattering investigation ruled out either $`B_{1u}`$ or $`B_{2u}`$ octupoles at the selected wavevectors, including the (0,0,1) periodicity of the weak-moment antiferromagnetism. It has been argued that resonant X-ray scattering would not see the octupoles. Last, but not least, though the octupolar scenario is as yet the only one to account for the observed fact that transverse uniaxial stress induces $`\widehat{z}`$ antiferromagnetism, in a simple mean field version it forces a choice between $`\sigma (100)`$ and $`\sigma (110)`$ directions of stress, while experiments seem to tell us that these directions are equivalent. An essential insight is missing here.
A microscopic theory will have to address the questions raised above. Our present investigation is of a limited scope: we use general arguments to classify the symmetry-allowed ordering transitions of URu<sub>2</sub>Si<sub>2</sub>. We are particularly interested in the possibility of a sequence of such continuous phase transitions. We have no statement to make if the lower-temperature transition is of first order.
We address the simplest questions for which no knowledge of microscopic details is required: Once we had an octupolar ordering transition, is there any compelling reason to expect a second symmetry-breaking transition? What may be the order parameters? Should we expect still more, as yet undiscovered, phase transitions?
In an abstract sense, our question is the following: assuming the symmetry of the high-temperature (”para”) symmetry group $`𝒢_{\mathrm{para}}`$ has been broken by a spontaneous ordering transition, introducing order $`𝒪_10`$ lowers the symmetry to $`𝒢(𝒪_1)𝒢_{\mathrm{para}}`$. What is the structure of $`𝒢(𝒪_1)`$? What is the new classification of the order parameters? Assuming that such order parameters are found, does this imply that further continuous phase transitions necessarily happen?
In Ref., we discussed in some detail two questions of this nature: the symmetry classification of the order parameters in the presence of an 1) external magnetic field $`\widehat{z}`$, and 2) a uniaxial stress. Formally the same question arises if instead of externally applied fields, we assume the presence of an effective field associated with either a dipole ordering transition (1) or quadrupolar order (2).
In Sec. 3.1, we rephrase our earlier results on the effect of an external magnetic field, and add some remarks about the field direction dependence. In Sec. 3.2, we analyze the symmetry in the presence of a $`𝒯_{xyz}`$ octupolar (effective) field, and the possibility of a sequence of phase transitions. Sec. 4 illustrates the general arguments with a simple mean field calculation.
## 3 Symmetry lowering in external, or effective, fields
### 3.1 Magnetic field
In the presence of an external magnetic field $`H_z(0,0,1)`$ or equivalently $`J_z>0`$, the remaining purely geometrical symmetry elements are $``$, $`2𝒞_4`$, and $`𝒞_2^2`$ (and, naturally, the inversion $``$). The geometrical symmetry (described by unitary operations) is lowered to $`𝒞_{4h}`$, but the full<sup>3</sup><sup>3</sup>3Magnetic field is invariant under space inversion. It is understood that all symmetry groups would get doubled if we included the inversion $``$. symmetry group $`𝒢_{\mathrm{tetr}}(J_z)`$ contains also non-unitary elements, namely $`2𝒞_2^{}𝒯`$, and $`2𝒞_2^{\prime \prime }𝒯`$, where $`𝒯`$ is time reversal. $`𝒢_{\mathrm{tetr}}(J_z)`$ contains eight elements, its character table (and multiplication table) is the same as that of the tetragonal point group $`𝒟_4`$. The character table is given in Table 2.
Reducing the symmetry from $`𝒢_{\mathrm{tetr}}`$ to $`𝒢_{\mathrm{tetr}}(J_z)`$, some of the originally inequivalent irreps of $`𝒢_{\mathrm{tetr}}`$ become equivalent. This parentage of the irreps is shown in Table 2. It also follows that the corresponding order parameters get mixed. This was discussed in Ref. .
$`\widehat{x}`$ is a lower-symmetry direction than $`\widehat{z}`$ and correspondingly a system subjected to a magnetic field $`𝐁=(B_x,0,0)`$ has a smaller symmetry group $`𝒢_{\mathrm{tetr}}(J_x)`$ (Table 3). In contrast to the case $`𝐁\widehat{z}`$ where irreps taken from the zero field case kept their dimensionality only pairwise merged, now the two-dimensional irreps split. The parentage of the $`𝐁\widehat{x}`$ irreps is as follows $`A_{1g},B_{1g}\mathrm{\Gamma }_1`$, $`A_{2g},B_{2g}\mathrm{\Gamma }_2`$, $`A_{1u},B_{1u}\mathrm{\Gamma }_3`$, $`A_{2u},B_{2u}\mathrm{\Gamma }_4`$, $`E_g\mathrm{\Gamma }_3+\mathrm{\Gamma }_4`$, $`E_u\mathrm{\Gamma }_1+\mathrm{\Gamma }_2`$.
The presence of non-identity irreps in Tables 2 and 3 shows that for fields $`𝐁\widehat{z}`$, and even for $`𝐁\widehat{x}`$, a variety of symmetry breaking transitions are possible. Since the field has broken time reversal invariance, and induced a polarization in its direction<sup>4</sup><sup>4</sup>4Assuming the local Hilbert space allows a Zeeman splitting for the particular field direction. E.g., the quasidoublet $`\{|t_1,|t_4\}`$ of URu<sub>2</sub>Si<sub>2</sub> would not split in a field $`𝐁\widehat{z}`$., all these transitions have the character of orbital ordering (lifting some residual non-Kramers degeneracy), and also lead to the appearance of time-reversal-odd transverse polarization components. We enumerate the possibilities ($`𝐁\widehat{z}`$)
* A field $`𝐁\widehat{z}`$ allows the ordering of $`E`$ type quadrupoles and this is accompanied by the appearance of dipole polarization in the $`xy`$ plane. This was suggested for the disjoint high-field phase observed in experiments on URu<sub>2</sub>Si<sub>2</sub>. Analogous phenomena are observed in high field experiments on the tetrahedral skutterudite PrOs<sub>4</sub>Sb<sub>12</sub>.
* The ordering is of $`B_1`$ type. In our suggestion for the low-field order of URu<sub>2</sub>Si<sub>2</sub>, the primary order was $`𝒯_z^\beta `$ octupolar (this survives when the field is switched off), and $`𝒪_2^2`$ quadrupoles were field-induced. Santini’s scenario was the opposite: primary quadrupolar order, and field-induced octupoles.
* $`B_2`$ order parameters: primary $`𝒯_{xyz}`$ octupolar order and field-induced $`𝒪_{xy}`$ quadrupoles. In the present paper, we assume this is the hidden order with onset temperature $`T_0=17`$K. Lacking a microscopic model, it is impossible to decide between this scenario and the previous $`B_1`$ scheme.
* The exotic possibility of $`A_2`$ order: $`_1`$ hexadecapoles and $`\overline{J_xJ_yJ_z(J_x^2J_y^2)}`$ triakontadipoles, has not been considered yet.
The discussion of the cases $`𝐁\widehat{y}`$ and $`𝐁(\widehat{x}\pm \widehat{y})`$ would be analogous to $`𝐁\widehat{x}`$: the remaining symmetry group has order 4 (or 8, if we include inversion). The symmetry group for fields lying in the $`x`$$`y`$ plane in non-special directions is generated by $`𝒞_4^2𝒯`$ and $``$. For fields in general out-of-plane directions, the only remaining symmetry is inversion.
We conclude that a spontaneous symmetry breaking transition in external fields remains possible if the field is either parallel or perpendicular to $`\widehat{z}`$. In particular, the octupolar transition can remain second order for a range of field strengths. However, $`𝒯_{xyz}`$ octupoles get mixed with $`𝒪_{xy}`$ ($`B_2`$) quadrupoles for $`𝐁\widehat{z}`$, and with $`E`$-derived quadrupoles (a suitable combination of $`𝒪_{yz}`$ and $`𝒪_{xz}`$) for $`𝐁\widehat{z}`$. For general field directions, a smearing of the octupolar transition should be observed. This effect may be helpful in deciding whether the hidden order is of octupolar nature.
The above reasoning can also be extended to discuss further symmetry breakings after a phase transition to a dipole-ordered phase has taken place. The effective field of the ordered moments acts the same way as an external field. E.g., we may conclude that following a transition to a $`J_z0`$ phase at the higher critical temperature $`T_c^>`$, tranverse dipoles (from the $`E`$ doublet $`\{J_x,J_y\}`$) may also order at a lower critical temperature $`T_c^<`$, and that $`J_x0`$ induces $`𝒪_{xz}0`$. Though symmetry allows this (or some other scenario identifiable from Table 2) to happen, whether this potentiality is realized depends on the nature of the microscopic model.
### 3.2 The symmetry group of the octupolar phase
We may ask whether following the onset of, say, $`𝒯_{xyz}`$ type octupolar order a further symmetry lowering transition is possible. This is the same question as to whether spontaneous symmetry breaking in an external octupolar field of $`B_{1u}`$ symmetry is possible. The question is partly academical, but it is also motivated by the NQR findings on URu<sub>2</sub>Si<sub>2</sub> by Saitoh et al. In particular, they claimed that either $`𝒪_{xy}`$ or $`𝒪_2^2`$ quadrupoles appear at $`T^{}13.5\mathrm{K}`$, well below $`T_017\mathrm{K}`$ which we associate with the onset of octupolar order. However, as we emphasized before, the connection of our arguments with experiments remains tenuous.
We have to identify which operations leave the $`𝒯_{xyz}`$ octupoles unchanged, i.e., we are looking for the octupolar symmetry group $`𝒢_{\mathrm{tetr}}(𝒯_{xyz})`$ as a subgroup of $`𝒢_{\mathrm{tetr}}`$. Schematically, $`𝒯_{xyz}`$ is represented in Fig. 5 where the current distribution in the ground state of $`𝒯_{xyz}`$ is shown for two- and three-dimensional subspaces selected from the 9-dimensional Hilbert space of the tetragonal crystal field eigenstates. It is obvious that the selection of the basis functions (which is ”done” by the crystal field potential) has a strong influence on the details of the current pattern. However, though the pattern is much more decorated for the $`\{|t_1,|t_3\}`$ doublet than for the triplet $`\{|t_1,|t_2,|t_4\}`$ (the crystal field model used in Ref.), the symmetries are the same. In either case, there is a current distribution with eight major eddies: four positive, and four negative (whether within these, there are sub-eddies, has no influence on the symmetry). Though there are local magnetic fields, the total magnetic moment is zero. A $`𝒞_4`$ rotation takes positive eddies into negative ones, and vice versa; however, reversing also the direction of currents, the original state is restored. From the $`𝒞_2^{}`$ and $`𝒞_2^{\prime \prime }`$ elements of $`𝒢_{\mathrm{tetr}}`$, the latter two have to be combined with $`𝒯`$. The character table of the symmetry group $`𝒢_{\mathrm{tetr}}(𝒯_{xyz})`$ is shown in Table 4 (as before, it is understood that the space inversion $``$ would generate the other half of the complete symmetry group). We have also given the resulting symmetry classification of some of the order parameters in the last column of Table 4.
It is interesting to note the similarities and dissimilarities to the symmetry group of $`J_z`$ (Table 2). The character table is the same only in $`𝒢_{\mathrm{tetr}}(𝒯_{xyz})`$, $`2𝒞_4`$ and $`2𝒞_2^{}`$, while in $`𝒢_{\mathrm{tetr}}(J_z)`$, $`2𝒞_2^{}`$ and $`2𝒞_2^{\prime \prime }`$ have to be combined with $`𝒯`$.
Analogous results would have been obtained if we had assumed a $`𝒯_z^\beta `$ octupolar field instead of $`𝒯_{xyz}`$. That was our starting assumption in Ref. .
$`𝒢_{\mathrm{tetr}}(𝒯_{xyz})`$ has two generators. There are several choices:
* $`𝒞_4𝒯`$ and $`𝒞_{2x}^{}`$ or $`𝒞_{2y}^{}`$
* $`𝒞_4𝒯`$ and $`𝒞_{2,x+y}^{\prime \prime }𝒯`$ or $`𝒞_{2,xy}^{\prime \prime }𝒯`$
* either of $`𝒞_2^{}`$ and either of $`𝒞_2^{\prime \prime }`$
At $`T^{}<T_0`$, one of the order parameters appearing in Table 4 acquires non-zero expectation value, and the symmetry of the system is lowered to one of the subgroups of $`𝒢_{\mathrm{tetr}}(𝒯_{xyz})`$. The list of the subgroups is
$`𝒢_1`$ $`=`$ $`\{,𝒞_4𝒯,𝒞_4^2,𝒞_4^3𝒯\}`$
$`𝒢_2`$ $`=`$ $`\{,𝒞_{2,x}^{}\}`$
$`𝒢_3`$ $`=`$ $`\{,𝒞_{2,y}^{}\}`$
$`𝒢_4`$ $`=`$ $`\{,𝒞_{2,x+y}^{\prime \prime }𝒯\}`$
$`𝒢_5`$ $`=`$ $`\{,𝒞_{2,xy}^{\prime \prime }𝒯\}`$
$`𝒢_6`$ $`=`$ $`\{,𝒞_{2,x}^{},𝒞_{2,y}^{},𝒞_4^2\}`$
$`𝒢_7`$ $`=`$ $`\{,𝒞_{2,x+y}^{\prime \prime }𝒯,𝒞_{2,xy}^{\prime \prime }𝒯,𝒞_4^2\}`$
$`𝒢_8`$ $`=`$ $`\{,𝒞_4^2\}`$
Each of the order parameters appearing in Table 4 breaks one, or several, of the symmetries in $`𝒢_{\mathrm{tetr}}(𝒯_{xyz})`$, and thus reduces the symmetry to one of the subgroups. All the possibilities are listed below
$`𝒯_z^\beta \text{ and }\overline{J_xJ_y(J_x^2J_y^2)}`$ $`𝒢_1`$
$`𝒪_{yz}\text{ and }J_x`$ $`𝒢_2`$
$`𝒪_{xz}\text{ and }J_y`$ $`𝒢_3`$ (8)
$`𝒪_2^2\text{ and }\overline{J_xJ_yJ_z(J_x^2J_y^2)}`$ $`𝒢_6`$
$`𝒪_{xy}\text{ and }J_z`$ $`𝒢_7`$ (9)
Time-reversal-even and time-reversal-odd order parameters appear in pairs. Since the octupolar background already breaks time reversal invariance, it cannot be ”unbroken” by the next transition, so the new phase has a new pattern of non-zero current densities. All the quadrupoles allowed by tetragonal symmetry can appear in a symmetry-breaking transition but they bring either dipoles or higher magnetic multipoles with them. We list the possibilities:
* Most straightforward is the case of $`𝒪_{xy}`$ quadrupolar ordering accompanied by magnetic moments in the $`z`$ direction. The coupling of these three multipoles: $`𝒯_{xyz}`$, $`𝒪_{xy}`$, and $`J_z`$ has been discussed from several points of view. In Ref. , we argued that on the background of $`𝒯_{xyz}`$ octupole ordering, applying uniaxial stress $`\sigma (1,1,0)`$ creates $`𝒪_{xy}`$ quadrupoles, and therefore magnetic moments $`J_z`$, offering an interpretation of the experimental results by Yokoyama et al. The existence of a third-order invariant $`𝒯_{xyz}𝒪_{xy}J_z`$ implies that if in an ordered phase $`𝒯_{xyz}0`$ then also the correlator $`𝒪_{xy}J_z0`$. The non-vanishing correlator does not automatically give non-zero values to $`𝒪_{xy}`$, and $`J_z`$, but at least it gives a hint that $`𝒪_{xy}`$, and $`J_z`$ are likely actors in a cooperative phenomenon. Here we discuss them as coupled order parameters of a low-temperature phase transition following a high-temperature octupolar transition.
* $`𝒪_{xz}`$ ($`𝒪_{yz}`$) quadrupoles would be accompanied by the transverse magnetization component $`J_x`$ ($`J_y`$), so the experimental verification (or refusal) of this scenario should be straightforward.
* More exotic is the possibility of $`𝒪_2^2`$ quadrupolar ordering (one of the likely candidates according to Ref. ) which should be accompanied by $`\overline{J_xJ_yJ_z(J_x^2J_y^2)}`$ triakontadipole ordering.
* Finally, the symmetry of the $`𝒯_{xyz}`$ octupolar field can be spontaneously broken by ordering the $`𝒯_z^\beta `$ octupoles; this should be accompanied by hexadecapole order of the $`_1=\overline{J_xJ_y(J_x^2J_y^2)}`$ type. We note that $`_1`$ hexadecapoles at the U sites can combine to an effective quadrupole at the Ru sites, so this scenario is not necessarily in conflict with the findings by Saitoh et al.. It is an added attraction that the simultaneous presence of $`𝒯_{xyz}0`$ and $`𝒯_z^\beta 0`$ would allow that uniaxial press applied either in the $`\sigma (100)`$ or the $`\sigma (110)`$ direction induce $`J_z0`$, as observed. We note that the measurements of Yokoyama et al. were carried out at 1.4K, well below either $`T_0`$ or $`T^{}`$, so both octupolar amplitudes would be near their saturation values. We have to admit, though, that in our scenario it would be difficult to get them equal.
### 3.3 Octupolar phase in external magnetic field
It is of some interest to combine the previous two cases to discuss the remaining possibilities of symmetry breaking if the established $`𝒯_{xyz}`$ octupolar order is subject to an external magnetic field.
A field $`𝐁\widehat{z}`$ reduces the symmetry to the four-element subgroup $`𝒢_7`$. There are four irreps: $`A_1^{}`$, $`A_2^{}`$,$`A_3^{}`$, and $`A_4^{}`$. In terms of the irreps shown in Table 4, their parentage is:
$$A_1,B_2A_1^{},$$
$$A_2,B_1A_2^{},$$
$$EA_3^{}+A_4^{}$$
Thus now $`𝒯_{xyz}`$, $`J_z`$, and $`𝒪_{xy}`$ are all present in the identity representation.
A spontaneous symmetry breaking transition is possible to a phase with $`𝒪_2^20`$; but then $`𝒯_z^\beta `$, $`_1`$, and $`\overline{J_xJ_yJ_z(J_x^2J_y^2)}`$ are induced order parameters.
There are two more possibilities of ordering, both involving a transverse dipole and a quadrupole (e.g., $`J_xJ_y`$ and $`𝒪_{xz}𝒪_{yz}`$).
A field $`𝐁\widehat{x}`$ reduces the symmetry to the two-element subgroup $`𝒢_2`$. Still, there remains one symmetry element to break: it can be done by $`J_z`$, and a number of associated multipoles.
To conclude this subsection: if the $`T<T_0`$ hidden order is octupolar, the possibilities of remaining symmetry breaking at $`T^{}`$ depend very much on the direction of the applied field. It seems that the only direction offering non-trivial possibilities is $`𝐁\widehat{z}`$, where the $`𝒪_2^2`$ quadrupolar symmetry breaking is accompanied by induced octupolar, hexadecapole, and triakontadipole moments.
## 4 Mean field calculations
From the fact that the symmetry group of a model contains non-trivial elements, it does not necessarily follow that symmetry breaking phase transitions occur until the only remaining symmetry element is the unit operator. It is possible that the local Hilbert space is not large enough, or it does not have the right structure, to support a sequence of ordering transitions. An ordering transition is an inevitability only if in its absence, the ground state is degenerate. In models of octupolar ordering, this is usually not the case. Santini’s and our crystal field model is built with three singlets, so the single ion ground state is in any case non-degenerate. Still, symmetry breaking transitions can occur by the induced moment mechanism if the level splittings are not too large to begin with. After the first ordering transition - which we assume is $`𝒯_{xyz}`$ octupolar ordering - has taken place at $`T_0`$, the symmetry is reduced to $`𝒢_{\mathrm{tetr}}(𝒯_{xyz})`$. The ionic level scheme is formed by the combined action of the crystal field potential and the octupolar effective field, and the ground state is again non-degenerate. Nevertheless, the induced moment mechanism can be effective again and we may ask whether another second order transition may take place at $`T^{}<T_0`$, corresponding to one of the options listed in (8). A mean field calculation shows that this is indeed possible for a wide range of parameters. Naturally, one has to assume non-zero couplings for the multipoles which appear in the definition of the new order (for instance, $`𝒪_2^2`$ and $`\overline{J_xJ_yJ_z(J_x^2J_y^2)}`$, or $`𝒪_{xy}`$ and $`J_z`$), but this is in any case more plausible than setting the said couplings to zero. The starting values of the crystal field splittings do not play any particular role, except that they have to be small enough to allow ordering.
Let us observe that $`𝒢_{\mathrm{tetr}}(𝒯_{xyz})`$ has a two-dimensional ($`E`$) representation. It shows that if the crystal field scheme includes a low-lying doublet, then this feature may be preserved also in the octupole-ordered phase. The two-fold degeneracy may then be lifted at the second ordering transition. In this sense, crystal field schemes with either a doublet ground state, or a low-lying doublet, offer a more direct route to a second symmetry breaking transition. We note that the crystal field level scheme is not undisputed: while models with low-lying singlets have been most widely discussed, doublet ground state was also considered . However, a sequence of two continuous transitions is possible either with, or without, a low-lying doublet.
### 4.1 Three singlets
Since $`f^2`$ is a non-Kramers configuration, it is possible to choose time reversal invariant basis functions. For the singlets $`A_1`$, $`A_2`$, and $`B_2`$, these are
$`|t_1`$ $`=`$ $`b\left(|4+|4\right)+\sqrt{12b^2}|0`$ (10)
$`|t_2`$ $`=`$ $`{\displaystyle \frac{i}{\sqrt{2}}}\left(|4|4\right)`$ (11)
$`|t_4`$ $`=`$ $`{\displaystyle \frac{i}{\sqrt{2}}}\left(|2|2\right)`$ (12)
In this basis, electric (magnetic) multipoles are real (imaginary) operators (see Appendix).
Choosing the subspace of the singlets (10-12) is the common starting point for a mean field calculation in Santini’s and our work. The presence of these three states is useful for getting good overall agreement with measured macroscopic quantities. Depending on the assumptions about the sequence and the splitting of the levels, and about the matrix elements connecting them, the three-singlet scheme may support $`J_z`$ dipolar, $`𝒪_2^2`$ or $`𝒪_{xy}`$ quadrupolar, $`𝒯_z^\beta `$ or $`𝒯_{xyz}`$ octupolar, or $`_1`$ hexadecapole order in the first ordered phase, and a combination of two more orders from the same list below a second critical temperature. Up to now, mainly scenarios with a single transition were considered. Santini chose quadrupolar order while we preferred octupolar order. There is no effective way to choose between $`𝒯_z^\beta `$ and $`𝒯_{xyz}`$. Our earlier discussion was based on postulating $`𝒯_z^\beta `$ ordering. Here we phrase our arguments on the alternative assumption that $`𝒯_{xyz}`$ order appears first.
The energy (free energy) density can be expressed as a sum of invariants. A number of third order invariants contain $`𝒯_{xyz}`$. The corresponding third-order contribution to the free energy is
$`\mathrm{\Delta }F_3`$ $`=`$ $`c_1(𝐐_1,𝐐_2,𝐐_3)𝒯_{xyz}(𝐐_1)𝒯_\beta ^z(𝐐_2)_1(𝐐_3)`$ (13)
$`+c_2(𝐐_1,𝐐_2,𝐐_3)𝒯_{xyz}(𝐐_1)𝒪_{xy}(𝐐_2)J_z(𝐐_3)`$
$`+c_3(𝐐_1,𝐐_2,𝐐_3)𝒯_{xyz}(𝐐_1)\left[J_x(𝐐_2)𝒪_{zx}(𝐐_3)+J_y(𝐐_2)𝒪_{yz}(𝐐_3)\right]`$
$`+c_4(𝐐_1,𝐐_2,𝐐_3)𝒯_{xyz}(𝐐_1)𝒪_2^2(𝐐_2)𝒦(𝐐_3)`$
where $`𝐐_1+𝐐_2+𝐐_3=\mathrm{𝟎}`$, and $`𝒦=\overline{J_xJ_yJ_z(J_x^2J_y^2)}`$.
If $`𝒯_{xyz}(𝐐_1)`$ is the hidden order, i.e., $`𝒯_{xyz}(𝐐_1)0`$ at $`T<T_0`$, then in the same temperature range the double correlators read off from above become non-zero, e.g., $`𝒪_{xy}(𝐐_2)J_z(𝐐_3)0`$. This is helpful, but in itself not sufficient, for the ordering of $`𝒪_{xy}`$ and $`J_z`$ separately. For that, dipole-dipole and quadrupole-quadrupole interactions have to be non-zero, and sufficiently strong to overcome the level splittings arising from the combined effect of the crystal field and the octupolar effective field. Playing with parameters introduces a lot of arbitrariness into a crystal field model, and at present it would be pointless to try to fit them to experiments, especially as there is no agreement about the relevant Q vectors. We merely wish to demonstrate the possibility of a second ordering transition at $`T^{}<T_0`$.
For the order parameters appearing at $`T^{}`$, any of the pairs in (9) could be chosen. For the moment, we arbitrarily pick $`𝒪_{xy}`$ and $`J_z`$.
Ordering is possible by the induced moment mechanism in spite of crystal field splittings, but its nature is basically the same as it would be for three degenerate singlets<sup>5</sup><sup>5</sup>5Crystal field splittings establish an asymmetry between the complementary variables $`𝒪_{xy}`$ and $`J_z`$ because $`𝒪_{xy}`$ has a matrix element between $`|t_1`$ and $`|t_2`$, while $`J_z`$ between $`|t_1`$ and $`|t_4`$. Asymmetry arises also from unequal multipole couplings.. We introduced three mean field couplings: $`\lambda (𝒯_{xyz})`$, $`\lambda (𝒪_{xy})`$, and $`\lambda (J_z)`$, and solved the self-consistency equations for the three order parameters.
Typical results are shown in Fig. 6. The onset of $`𝒯_{xyz}`$ octupolar order at $`T_0=0.645`$ is followed by the ordering of $`𝒪_{xy}`$ and $`J_z`$ at $`T^{}0.495`$. To bring out the contrast<sup>6</sup><sup>6</sup>6The critical temperature is the same, but the amplitudes unequal. between the complementary parameters $`𝒪_{xy}`$ and $`J_z`$, we used the parameter set $`\mathrm{\Delta }_{21}=1`$ as the energy unit, $`\mathrm{\Delta }_{41}/\mathrm{\Delta }_{21}=2`$, $`\lambda (𝒪_{xy})=\lambda (J_z)`$, $`\lambda (J_z)/\mathrm{\Delta }_{21}=10`$, and $`\lambda (𝒯_{xyz})/\mathrm{\Delta }_{21}=17`$.
The second step of ordering is assisted by the pre-existing octupolar order. This can be seen in the susceptibility plot Fig. 7 (left). We have chosen the quadrupolar susceptibility which belongs to the ordering degree of freedom $`𝒪_{xy}(𝐐_2)`$. The high-temperature susceptibility extrapolates to a lower quadrupolar ordering temperature than the behavior calculated within the octupolar phase. An alternative argument is based on the Landau expansion (we drop the Qs)
$``$ $`=`$ $`a_{\mathrm{oct}}(TT_0)𝒯^2+b_{\mathrm{oct}}𝒯^4+a_{\mathrm{quad}}(TT_1)(𝒪^2+J^2)`$ (14)
$`+b_{\mathrm{quad}}(𝒪^4+J^4)+e𝒯^2(𝒪^2+J^2)+f𝒯𝒪J`$
The last term corresponds to (13). Whatever the sign of $`f`$, the sign of $`𝒪J`$ can be chosen so that the term lowers the free energy, and enhances the lower critical temperature from $`T_1`$ to $`T^{}`$, as observed.
The octupolar order responds to the onset of quadrupolar-dipolar order by a change of the slope of the temperature dependence (Fig. 7, right).
We emphasize that we did not make any attempt to fine-tune our mean field parameters. If we plotted the specific heat, it would show two lambda-anomalies, in conflict with the experiments known to us. We think that until further experiments are done, it would be pointless to refine our calculation.
## 5 Conclusion
We studied possible sequences of symmetry breaking transitions in tetragonal $`f^2`$ systems. Assuming that octupolar order sets in first, the next transition may lead to mixed dipole–quadrupole, quadrupole–triakontadipole, or octupole–hexadecapole order. There is an interesting similarity to the recent NQR results by Saitoh et al but we feel that any detailed comparison would be premature.
## Acknowledgments
We are greatly indebted to S. Takagi for valuable advice, and for informing us about his results prior to publication. We thank H. Amitsuka, F. Bourdarot, W.J.L. Buyers, B. Fåk, G. Kriza, and G. Solt for enlightening discussions and/or correspondence. At all points of this work, we were helped by advice from, and discussions with, K. Penc. We were supported by the Hungarian National Grants T038162, T049607, and TS049881.
## Appendix A
The off-diagonal multipoles in basis (10-12) are
$$J_z=4\sqrt{2}b\left(\begin{array}{ccc}0& i& 0\\ i& 0& 0\\ 0& 0& 0\end{array}\right)\text{ }𝒪_2^2=2\sqrt{7}\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& 1\\ 0& 1& 0\end{array}\right)$$
(15)
$$𝒪_{xy}=(2\sqrt{14}b+6\sqrt{510b^2})\left(\begin{array}{ccc}0& 0& 1\\ 0& 0& 0\\ 1& 0& 0\end{array}\right)$$
(16)
$$𝒯_z^\beta =3\sqrt{3}(\sqrt{70}b+5\sqrt{12b^2})\left(\begin{array}{ccc}0& 0& i\\ 0& 0& 0\\ i& 0& 0\end{array}\right)$$
(17)
$$𝒯_{xyz}=3\sqrt{105}\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& i\\ 0& i& 0\end{array}\right)\text{ }_1=3\sqrt{3570b^2}\left(\begin{array}{ccc}0& 1& 0\\ 1& 0& 0\\ 0& 0& 0\end{array}\right)$$
(18)
All the matrix elements of the triakontadipole operator $`\overline{J_xJ_yJ_z(J_x^2J_y^2)}`$ vanish in this subspace. |
warning/0506/hep-ex0506005.html | ar5iv | text | # TOP QUARK PRODUCTION AND PROPERTIES AT THE TEVATRON
## 1 Introduction
The top quark is special among the fermions of the Standard Model because of its large mass. Currently, the top quark can only be studied at the two Tevatron experiments CDF and DØ, where measurements of top quark production and properties are one of the key physics goals of Run II.
The top quark mass is discussed in a separate article. This article focuses on measurements of the total $`t\overline{t}`$ production cross-section, searches for new physics in $`t\overline{t}`$ production and top quark decay, and on the search for electroweak (single) top quark production. While the CDF and DØ experiments have both collected more than $`500\mathrm{pb}^1`$ of data so far during Tevatron Run II, surpassing the Run I integrated luminosity by a factor $``$5, the measurements summarized in this article typically use about $`200\mathrm{pb}^1`$.
In section 2, general aspects of top production and event topologies at the Tevatron are briefly discussed. Section 3 discusses the measurements of the total $`t\overline{t}`$ production cross-section, while further measurements in $`t\overline{t}`$ events are presented in section 4. The search for single top quark production is presented in section 5.
## 2 Top Quark Production at the Tevatron
In the Standard Model, the production of top quarks at a hadron collider can in principle proceed via two mechanisms: $`t\overline{t}`$ pair production via the strong interaction, and single top (or antitop) production via the electroweak interaction. The leading order Feynman diagrams are shown in figure 1 together with the Standard Model cross-sections in $`p\overline{p}`$ collisions at a centre-of-mass energy of 1.96 TeV (corresponding to the Tevatron collider at Run II).
In the following, the main characteristics of $`t\overline{t}`$ and single top events at the Tevatron are discussed.
### 2.1 Classification of $`t\overline{t}`$ Event Topologies
In the Standard Model, the branching fraction of top quark decays to a b quark and an on-shell $`W`$ boson is close to $`100\%`$, other decay modes not being observable with Tevatron luminosities. The subsequent $`W`$ decays determine the event topology seen in the detector, and $`t\overline{t}`$ events are classified as follows:
* Dilepton events, where both $`W`$ bosons decay into an $`e\nu `$ or $`\mu \nu `$ final state, are characterized by two energetic, isolated leptons of opposite charge, two energetic b jets, and missing transverse energy. While the product branching ratio is only about $`5\%`$, pure event samples can be obtained requiring the two leptons in the event to be reconstructed.
* In lepton+jets events, one $`W`$ boson decays hadronically and the other into an $`e\nu `$ or $`\mu \nu `$ final state. This topology is characterized by an energetic, isolated electron or muon, four energetic jets (two b jets and two light-quark jets from the $`W`$ decay), and missing transverse energy. The product branching ratio of $`30\%`$ is larger than for dilepton events, and the main background is from $`W`$+jets events.
* In hadronic events, one expects 6 energetic jets (of which two are b jets) and no significant missing transverse energy. Because of large backgrounds from QCD jet production, identifying $`t\overline{t}`$ events in the hadronic channel is challenging, despite the large product branching ratio of $`44\%`$.
* In about $`21\%`$ of the $`t\overline{t}`$ events, at least one $`W`$ boson decays into a $`𝝉𝝂`$ final state. Depending on its decay, the $`\tau `$ lepton can be identified as a narrow jet, an isolated track, or an electron or muon. Two energetic b jets, missing transverse energy, and the decay products from the second $`W`$ boson complete the event topology.
In general, the reconstruction and selection of $`t\overline{t}`$ event candidates is based on reconstructing the directions and energies/momenta of isolated electrons or muons and jets, and on reconstructing the missing transverse energy $`E/_T`$ from the transverse momentum balance in the event. The purity of the event samples can be enhanced by identifying jets that originated from a b quark (b tagging), since in the Standard Model, every $`t\overline{t}`$ event contains two b jets. Both CDF and DØ use
* secondary vertex algorithms, based on explicit reconstruction of the decay vertex of the b hadron within the jet;
* impact parameter based algorithms that classify tracks inside a jet according to their distance of closest approach to the primary event vertex; and
* soft leptons from semileptonic bottom or charm hadron decay (only muons are used so far)
to identify b jets. The requirements on the jet multiplicity, the minimum jet transverse energies, b identification of the jets, and event kinematic information can be balanced to minimize the measurement error; depending on the selection, not all jets need to be explicitly reconstructed.
### 2.2 Single Top Quark Production
The total cross-section for single top quark production is only a factor $``$2 smaller than that for $`t\overline{t}`$ production; however, the relevant backgrounds are substantially larger ($`W`$+2jet instead of $`W`$+4jet events). To reduce the background, the selection of single top event candidates focuses on top decays with leptonic $`W`$ decays and on the identification of the b jet(s) in the event. Figure 2 shows the expected pseudorapidity distributions of the charged lepton and jets in single top events. For s-channel events, two b jets are expected in the center of the detector. In general only one b jet can be reconstructed in case of the t-channel, but here an additional light quark jet can be observed.
## 3 Measurements of the Total $`𝒕\overline{𝒕}`$ Production Cross-Section
The goal is to measure the $`t\overline{t}`$ cross-section in as many different modes as possible to check the predictions of the Standard Model. The measurements in the different channels are described in the following sub-sections.
### 3.1 Lepton+Jets Channel, Topological Analyses
Both CDF and DØ have measured the $`t\overline{t}`$ cross-section in the lepton+jets channel without relying on b jet identification. In the CDF analysis, events with one isolated electron with $`E_T>20\mathrm{GeV}`$ or muon with $`p_T>20\mathrm{GeV}`$, missing transverse energy $`E/_T>20\mathrm{GeV}`$, and at least 3 jets with $`E_T>15\mathrm{GeV}`$ within the pseudorapidity range $`|\eta |<2.0`$ are selected. For $`E/_T<30\mathrm{GeV}`$, there is an additional cut requiring the angle $`\mathrm{\Delta }\varphi `$ between the missing transverse energy and the highest $`E_T`$ jet in the transverse plane to satisfy $`0.5<\mathrm{\Delta }\varphi <2.5`$. After this selection, the $`t\overline{t}`$ signal can be seen in the $`H_T`$ distribution ($`H_T`$ is the scalar sum of transverse energies of the lepton, jets, and the missing transverse energy) as shown in figure 3, and the $`t\overline{t}`$ cross-section is measured to be
$$\sigma (t\overline{t})=(4.7\pm 1.6(\mathrm{stat}.)\pm 1.8(\mathrm{syst}.))\mathrm{pb}$$
(1)
using a $`195\mathrm{pb}^1`$ data sample. To optimise the measurement, 7 quantities have been chosen for training a neural network (NN) to separate the $`t\overline{t}`$ signal from the background. The NN output distribution is also shown in figure 3. From a fit to this distribution a result of
$$\sigma (t\overline{t})=(6.7\pm 1.1(\mathrm{stat}.)\pm 1.6(\mathrm{syst}.))\mathrm{pb}$$
(2)
is obtained from the same dataset. In both analyses, the main systematic error is from the uncertainty in the jet energy scale; it is however reduced from $`30\%`$ in the $`H_T`$ based measurement to $`16\%`$ in the optimised analysis.
In the DØ topological analysis, events with one isolated electron or muon with $`p_T>20\mathrm{GeV}`$, missing transverse energy ($`E/_T>20\mathrm{GeV}`$ in the e+jets case and $`E/_T>17\mathrm{GeV}`$ for $`\mu `$+jets events), and at least four jets with $`E_T>15\mathrm{GeV}`$ within $`|\eta |<2.5`$ are selected. So in contrast to CDF, four jets are required to be reconstructed, but a larger pseudorapidity region is allowed. To further separate $`t\overline{t}`$ events from background, a likelihood is constructed using angular variables and ratios of energy dependent variables, to avoid direct dependence on the jet energy scale. The resulting distributions are given in figure 4. The combined fit to the e+jets and $`\mu `$+jets distributions from $`141144\mathrm{pb}^1`$ of data yields
$$\sigma (t\overline{t})=(7.2_{2.4}^{+2.6}(\mathrm{stat}.)_{1.7}^{+1.6}(\mathrm{syst}.)\pm 0.5(\mathrm{lumi}.))\mathrm{pb}.$$
(3)
### 3.2 Lepton+Jets Channel, b Tagging Analyses
For the DØ measurements that make use of lifetime b tagging information, events are selected with the same criteria as above. The $`t\overline{t}`$ cross-section is then determined from a combined fit to the jet multiplicity distributions for events with exactly one b tagged jet and events with at least two b tagged jets. Using a data sample of $`158169\mathrm{pb}^1`$, DØ obtains
$$\sigma (t\overline{t})=(8.2\pm 1.3(\mathrm{stat}.)_{1.6}^{+1.9}(\mathrm{syst}.)\pm 0.5(\mathrm{lumi}.))\mathrm{pb}$$
(4)
using secondary vertex b tagging and
$$\sigma (t\overline{t})=(7.2_{1.2}^{+1.3}(\mathrm{stat}.)_{1.4}^{+1.9}(\mathrm{syst}.)\pm 0.5(\mathrm{lumi}.))\mathrm{pb}$$
(5)
with a track impact parameter based algorithm. The jet multiplicity distributions obtained with secondary vertex b tagging are shown in figure 5.
In a separate analysis, DØ analyzes events with a semimuonic bottom or charm decay, resulting in
$$\sigma (t\overline{t})=(11.2\pm 4.0(\mathrm{stat}.)\pm 1.3(\mathrm{syst}.)\pm 1.1(\mathrm{lumi}.))\mathrm{pb}$$
(6)
based on $`93\mathrm{pb}^1`$.
Several CDF analyses make use of b tagging information. The preselection of events requires one lepton, missing transverse energy, and three jets as in section 3.1. When at least one jet is required to be secondary vertex tagged, a measurement of
$$\sigma (t\overline{t})=\left(5.6_{1.1}^{+1.2}(\mathrm{stat}.)_{0.6}^{+0.9}(\mathrm{syst}.)\right)\mathrm{pb}$$
(7)
is obtained from the jet multiplicity distribution (where a value of $`H_T>200\mathrm{GeV}`$ is required for events with three or more jets) in $`162\mathrm{pb}^1`$ shown in figure 6. Alternatively, the fraction of $`t\overline{t}`$ events in events with at least 3 jets is obtained from a fit to the $`E_T`$ distribution of the leading jet, which is also shown in figure 6, yielding
$$\sigma (t\overline{t})=(6.0\pm 1.6(\mathrm{stat}.)\pm 1.2(\mathrm{syst}.))\mathrm{pb}.$$
(8)
For the CDF multiple tag analysis, a special version of the b tagging algorithm has been developed with looser criteria to increase the statistics. From the jet multiplicity distributions obtained with the regular and the loose b tag, measurements of
$`\sigma (t\overline{t})=\left(5.0_{1.9}^{+2.4}(\mathrm{stat}.)_{0.8}^{+1.1}(\mathrm{syst}.)\right)\mathrm{pb}`$ $`(\mathrm{regular}\mathrm{b}\mathrm{tag})\mathrm{and}`$ (9)
$`\sigma (t\overline{t})=\left(8.2_{2.1}^{+2.4}(\mathrm{stat}.)_{1.0}^{+1.8}(\mathrm{syst}.)\right)\mathrm{pb}`$ $`(\mathrm{loose}\mathrm{b}\mathrm{tag})`$ (10)
are obtained. Finally, CDF also uses events with b jets identified by an impact parameter based algorithm, yielding
$$\sigma (t\overline{t})=(5.8_{1.2}^{+1.3}(\mathrm{stat}.)\pm 1.3(\mathrm{syst}.))\mathrm{pb}$$
(11)
in $`162\mathrm{pb}^1`$, as well as events with a semimuonic bottom or charm hadron decay, resulting in
$$\sigma (t\overline{t})=\left(5.2_{1.9}^{+2.9}(\mathrm{stat}.)_{1.0}^{+1.3}(\mathrm{syst}.)\right)\mathrm{pb}$$
(12)
using $`200\mathrm{pb}^1`$.
### 3.3 Dilepton Channel
The CDF selection of $`t\overline{t}`$ events in the dilepton channel requires two isolated tracks ($`p_T>20\mathrm{GeV}`$) and missing transverse energy ($`E/_T>25\mathrm{GeV}`$). The $`t\overline{t}`$ production cross-section is determined from the jet multiplicity distribution of events where both tracks are identified as leptons, or events where only one identified lepton is required. The combined result is
$$\sigma (t\overline{t})=(7.0_{2.1}^{+2.4}(\mathrm{stat}.)_{1.1}^{+1.6}(\mathrm{syst}.)\pm 0.4(\mathrm{lumi}.))\mathrm{pb}$$
(13)
using $`200\mathrm{pb}^1`$. The jet multiplicity distribution for the analysis with at least one identified lepton is shown in figure 7.
One of the main backgrounds to $`t\overline{t}`$ production in the dilepton channel is diboson (mostly $`W`$$`W`$) production. In a separate analysis, CDF fits the two-dimensional jet multiplicity vs. $`E/_T`$ distribution to measure the $`t\overline{t}`$, $`WW`$, and $`Z\tau \tau `$ cross-sections simultaneously. This analysis yields
$$\sigma (t\overline{t})=(8.6_{2.4}^{+2.5}(\mathrm{stat}.)\pm 1.1(\mathrm{syst}.))\mathrm{pb}$$
(14)
For the DØ dilepton analysis, events with two isolated leptons with $`p_T>15\mathrm{GeV}`$ ($`p_T>20\mathrm{GeV}`$ in the dielectron channel), two jets with $`E_T>20\mathrm{GeV}`$, missing transverse energy $`E/_T>35\mathrm{GeV}`$ ($`E/_T>25\mathrm{GeV}`$ in the $`e\mu `$ channel), and $`H_T^{\mathrm{lead}.\mathrm{}}>120(140)\mathrm{GeV}`$ in the $`\mu \mu `$ ($`e\mu `$) channel are selected, where $`H_T^{\mathrm{lead}.\mathrm{}}`$ includes all jets and the leading lepton. Additional cuts reject events consistent with a $`Z\mathrm{}\mathrm{}`$ hypothesis. With no b tagging criteria applied, the analysis yields
$$\sigma (t\overline{t})=(14.3_{4.3}^{+5.1}(\mathrm{stat}.)_{1.9}^{+2.6}(\mathrm{syst}.)\pm 0.9(\mathrm{lumi}.))\mathrm{pb}$$
(15)
using $`140156\mathrm{pb}^1`$. When requiring at least one jet to be secondary vertex b tagged, the $`e\mu `$ channel alone yields
$$\sigma (t\overline{t})=(11.1_{4.3}^{+5.8}(\mathrm{stat}.)\pm 1.4(\mathrm{syst}.)\pm 0.7(\mathrm{lumi}.))\mathrm{pb}$$
(16)
with a very high purity sample, see figure 7.
### 3.4 Hadronic Channel
To separate $`t\overline{t}`$ events in the hadronic channel from the large multijet background, both tight kinematic cuts and b tagging information are applied. CDF selects events with 6 to 8 jets and no isolated leptons and applies kinematic cuts. In the distribution of the number of b tagged jets as a function of jet multiplicity (see figure 8) the $`t\overline{t}`$ cross-section is then measured to be
$$\sigma (t\overline{t})=(7.8\pm 2.5(\mathrm{stat}.)_{2.3}^{+4.7}(\mathrm{syst}.))\mathrm{pb}$$
(17)
in $`165\mathrm{pb}^1`$.
DØ selects events with 6 or more jets, of which exactly one is required to be b tagged. A chain of NNs feeding into each other is used, and the $`t\overline{t}`$ cross-section is determined from the excess of events after a cut on the last NN output over background. A data sample of $`162\mathrm{pb}^1`$ yields
$$\sigma (t\overline{t})=(7.7_{3.3}^{+3.4}(\mathrm{stat}.)_{3.8}^{+4.7}(\mathrm{syst}.)\pm 0.5(\mathrm{lumi}.))\mathrm{pb}.$$
(18)
### 3.5 Events with $`W\tau \nu `$ Decays
CDF searches for $`t\overline{t}`$ events where one $`W`$ decays electronically or muonically, while the other decays into a $`\tau \nu `$ final state with a subsequent $`\tau `$ decay into hadron(s) and a neutrino. Events with an electron or muon with $`E_T>20\mathrm{GeV}`$, a tau lepton with $`E>15\mathrm{GeV}`$ and opposite charge, missing transverse energy $`E/_T>20\mathrm{GeV}`$, at least two jets ($`E_T(1)>25\mathrm{GeV}`$, $`E_T(2)>15\mathrm{GeV}`$), and $`H_T>205\mathrm{GeV}`$ are selected. The two events observed in $`193.5\mathrm{pb}^1`$ are consistent with the Standard Model expectation, and a limit of
$$Br(tb\tau \nu )<5.0Br_{\mathrm{SM}}(tb\tau \nu )$$
(19)
is derived at 95% confidence level.
### 3.6 Summary of $`t\overline{t}`$ Cross-Section Measurements
The $`t\overline{t}`$ cross-section measurements at Tevatron Run II are summarized in figure 9. The measurements in all decay channels and by both CDF and DØ are mutually consistent and consistent with the prediction of the Standard Model.
## 4 Further $`𝒕\overline{𝒕}`$ Measurements
It is conceivable that physics beyond the Standard Model does not change the total $`t\overline{t}`$ cross-section, but either only affects differential cross-sections or top quark decays.
### 4.1 Searches for New Physics in $`t\overline{t}`$ Production
In a model independent analysis, CDF have searched for anomalous kinematic properties in their dilepton $`t\overline{t}`$ sample. Four kinematic distributions where new physics signatures are expected to be likely to be seen were chosen a priori. While one of them, the leading lepton $`p_T`$ spectrum shown in figure 10, shows an excess at low transverse momenta in $`193\mathrm{pb}^1`$, the other distributions agree with the expectation. The overall compatibility with the Standard Model prediction has been computed to be in the $`1.04.5\%`$ range.
CDF searches explicitly for production of fourth generation quarks. If these $`t^{}`$ quarks are heavier than the top quark, an excess of events at large $`H_T`$ is expected. From a fit to the measured $`H_T`$ distribution, which is consistent with the Standard Model expectation, upper limits on the cross-section of $`t^{}\overline{t^{}}`$ events can be placed, see figure 11.
### 4.2 Searches for New Physics in Top Quark Decays
The Standard Model predicts the fractions of longitudinal and left-handed $`W`$ bosons from top decay to be $`F_0=1/(1+2m_W^2/m_{top}^2)0.7`$ and $`F_{}=1F_0`$, while the fraction $`F_+`$ of right-handed W bosons is essentially zero because the bottom quarks from top quark decay are left-handed due to the large mass difference between top and bottom quarks. The predicted distribution of the decay angle $`\theta ^{}`$ in the $`W`$ rest frame is shown in figure 12. From measurements of this distribution (or quantities that depend on $`\mathrm{cos}\theta ^{}`$), one can either search for non-zero constributions from right-handed $`W`$ bosons ($`F_+>0`$) or, assuming $`F_+=0`$, for deviations from the predicted ratio $`F_0/F_{}`$.
CDF have measured the fraction $`F_0`$ from the charged lepton $`p_T`$ spectrum (using $`200\mathrm{pb}^1`$) to be
$$F_0=0.27_{0.24}^{+0.35},$$
(20)
and from explicit reconstruction of the value of $`\mathrm{cos}\theta ^{}`$ (using $`162\mathrm{pb}^1`$) to be
$$F_0=0.89_{0.34}^{+0.30}(\mathrm{stat}.)\pm 0.17(\mathrm{syst}.),$$
(21)
respectively – to be compared with the DØ Run I value of $`F_0=0.56\pm 0.31`$. All of these values are consistent with the Standard Model expectation.
The two DØ measurements both use explicit $`\mathrm{cos}\theta ^{}`$ reconstruction in an event sample obtained with a topological selection or using b tagging in $`159169\mathrm{pb}^1`$. The $`\mathrm{cos}\theta ^{}`$ distribution from the b tagging analysis is shown in figure 12. Both analyses each yield a limit of
$$F_+<0.24\mathrm{at}90\%\mathrm{confidence}\mathrm{level},$$
(22)
to be compared with the CDF Run I exclusion limit of $`F_+<0.18`$ at 95% confidence level.
In supersymmetric models with $`m_{H^\pm }<m_{top}`$, the top quark may decay into a charged Higgs and a bottom quark. Depending on the values of $`\mathrm{tan}\beta `$ and $`m_{H^\pm }`$, one expects the following changes in the observed $`t\overline{t}`$ event topologies:
* an excess of $`\tau `$ decays due to $`H^+\tau ^+\nu `$ decays for large $`\mathrm{tan}\beta `$,
* an excess of hadronic top decays due to $`H^+c\overline{s}`$ decays for small $`\mathrm{tan}\beta `$ and small $`m_{H^\pm }`$, or
* $`t\overline{t}`$ events with two extra b jets from $`H^+W^+b\overline{b}`$ decays for small $`\mathrm{tan}\beta `$ and large $`m_{H^\pm }`$.
The CDF collaboration has therefore taken their measurements of the $`t\overline{t}`$ cross-section in the dilepton and lepton+jets channels as well as their limit on $`t\overline{t}\mathrm{}+\tau `$ events to place limits on $`tH^+b`$ decays in the $`m_{H^\pm }`$ vs. $`\mathrm{tan}\beta `$ plane, as shown in figure 13.
Both CDF and DØ have compared their $`t\overline{t}`$ cross-section measurements obtained with different numbers of b tagged jets to determine the branching ratio $`Br(tWb)/Br(tWq)`$, where $`q`$ denotes any down-type quark. The results are
$$\begin{array}{ccc}1.11_{0.19}^{+0.21}\hfill & \hfill \mathrm{CDF},& 162\mathrm{pb}^1,\hfill \\ 0.65_{0.30}^{+0.34}(\mathrm{stat}.)_{0.12}^{+0.17}(\mathrm{syst}.)\hfill & \hfill \mathrm{D}\mathrm{Ø},& 158169\mathrm{pb}^1,\mathrm{and}\hfill \\ 0.70_{0.24}^{+0.27}(\mathrm{stat}.)_{0.10}^{+0.11}(\mathrm{syst}.)\hfill & \hfill \mathrm{D}\mathrm{Ø},& 158169\mathrm{pb}^1,\hfill \end{array}$$
(23)
where the first DØ result has been obtained with impact parameter b tagging and the second with secondary vertex b tagging. They show no sign of a deviation from the Standard Model expectation close to 1. It should be noted that this quantity does not constrain the value of $`|V_{tb}|^2`$ in models where top quark decays into quarks from more than three quark generations are allowed.
In summary, from measurements of $`t\overline{t}`$ production, there is currently no sign of physics beyond the Standard Model.
## 5 Search for Single Top Quark Production
The production cross-section for single top quarks is proportional to $`|V_{tb}|^2`$. Also, any differences to the Standard Model prediction could provide hints for new physics.
In their searches for single top quark production, the Tevatron experiments concentrate on s-channel and t-channel production with expected cross-sections of $`0.9\mathrm{pb}`$ and $`2.0\mathrm{pb}`$, respectively, see figure 1.
Both CDF and DØ select events with an energetic isolated charged lepton, missing transverse energy, and exactly 2 (CDF) or 2–4 (DØ) jets out of which at least one must be b tagged. CDF then selects events with a reconstructed top quark mass between $`140\mathrm{GeV}`$ and $`210\mathrm{GeV}`$, while DØ requires $`H_T>150\mathrm{GeV}`$. As shown in figure 14, single top events can be found at intermediate values of $`H_T`$, and s-channel and t-channel events can be disentangled using the lepton charge signed distribution of the pseudorapidity of the identified b jet. With the current data sets, sensitivity for Standard Model single top quark production has not yet been reached. No significant excess of events has been observed, and the following 95% confidence level limits have been placed on the single top quark cross-section:
$$\begin{array}{cccc}\mathrm{experiment}& \mathrm{s}\mathrm{channel}& \mathrm{t}\mathrm{channel}& \mathrm{s}+\mathrm{t}\mathrm{channel}\\ & & & \\ \mathrm{CDF}& 13.6\mathrm{pb}& 10.1\mathrm{pb}& 17.8\mathrm{pb}\\ \mathrm{D}\mathrm{Ø}& 19\mathrm{pb}& 25\mathrm{pb}& 23\mathrm{pb}\end{array}$$
(24)
With more data being taken and analysed and refined methods being developed, sensitivity for Standard Model single top quark production is within reach for Tevatron Run II.
## 6 Conclusions
The current status of top quark measurements at the Tevatron experiments CDF and DØ has been summarized, with the exception of the results for the top quark mass which are covered in a separate article.
A wealth of measurements of the total $`t\overline{t}`$ production cross-section are available from Tevatron Run II. Measurements have been performed for dilepton, lepton+jets, all-hadronic events, and events with top quark decays involving $`\tau `$ leptons. They all yield results that are both mutually consistent and in agreement with the Standard Model prediction.
The event samples have been further interpreted by looking for non-Standard Model $`t\overline{t}`$ production mechanisms and top quark decays. No signs for physics beyond the Standard Model have been found so far, supporting the interpretation of the signal as $`t\overline{t}`$ production via QCD and top quark decay to $`Wb`$ final states.
In the search for single (electroweak) production of top quarks, the sensitivity of the experiments has been improved over Run I. Even for Standard Model single top quark production, a significant cross-section measurement at the Tevatron is within reach in the near future.
## 7 Acknowledgements
The author would like to thank the organizers for a very interesting and enjoyable conference. |
warning/0506/nlin0506032.html | ar5iv | text | # A generalized model of active media with a set of interacting pacemakers: Application to the heart beat analysis
## 1 Introduction
Representation of an active distributed system by ensembles of coupled excitable or oscillatory elements is very useful method of the analysis because it allows to understand main dynamical processes inherent in the considered medium. As is known, this approach goes back to the model of Wiener and Rosenblueth \[Wiener & Rosenblueth, 1946\], according to which a medium consists of single elements being in one of three possible states: excited, refractory or rest. Later such models as coupled limit cycle oscillators and chaotic maps \[Kaneko, 1990; Shibata & Kaneko, 1998\] have played an important role not only in a quite realistic description of active media but also in the understanding of a possible behavior of systems far from equilibrium. Many useful concepts like phase–locked patterns, synchronization and spatio–temporal chaos have become popular due to detailed studies of similar nonlinear models \[Kuramoto, 1984; Kuramoto, 1995; Winfree, 2000\].
Investigations of such an example of an active medium as cardiac tissue are of significant scientific interest owing to vital importance of its rhythm stability. Real heart cells exhibit oscillatory properties (can be reset and entrained), they are excitable and have a refractory time, during which they do not respond to external stimulation. Hence, the heart can be considered as consisting of oscillatory (conductive cardiomyocytes, which have automaticity) and excitable (contractile heart cells, which do not initiate electrical activity under normal conditions) elements.
Due to extraordinary complexity of the heart, many qualitative discrete and continuous models have been tested. Majority of computational models of cardiac tissue of last generation takes into consideration the kinetics of excitable cells, how the excitation propagates from cell to cell and how contractile cardiomyocytes are arranged and connected in space \[Clayton, 2001\]. Such models mainly serve for studying sustained by re-entrant activity lethal arrhythmia – ventricular fibrillation, during which the spatio–temporal behavior is very complex \[Clayton et al., 2006\].
Other models treat the cardiac tissue as an active conductive system, taking into account oscillatory properties of heart cells. In this case the cardiac rhythms can be described on the basis of the dynamical system theory (see e.g. \[Courtemanche et al., 1989; Goldberger, 1990; Bub & Glass, 1994; Glass et al., 2002; Loskutov et al., 2004\] and refs. therein). Hereinafter we hold this approach.
Under normal conditions the electrical activity of the heart (action potentials) is spontaneously initiated in a region of the right atrium, sino–atrial (SA) node, so–called leading pacemaker. Automatic excitation is a distinctive feature not only of the cells of the SA node, but also of other conductive heart cells, so-called latent pacemakers. In addition, contractile cardiomyocytes can initiate a spontaneous action potential in pathology. Electrophysiological studies have suggested that the activity of cardiomyocytes with automaticity (e.g. P-cells of the SA node, of the atrioventricular (AV) junction, Purkinje cells) can be modulated by current pulses stimulating (super-threshold depolarizing) applied extracellularly \[Jalife & Moe, 1976; Sano et al., 1978; Jalife et al., 1980; Antzelevitch et al., 1982\].
Effects of external stimuli on biological oscillators are observed in a wide range of species. Experimentally obtained characteristics can be represented by a phase response curve (PRC) \[Jalife & Moe, 1976; Antzelevitch et al., 1982; Reiner & Antzelevitch, 1985\]. To establish the shape of the PRC experimentally, stimulation of an oscillator at various phases of its intrinsic cycle is applied \[Jalife et al., 1980; Guevara & Shrier, 1987\]. It has been found that in different pacemaker cells early stimuli delay the next pacemaker discharge and late pulses advance it. Therefore, the typical PRC shape is biphasic \[Jalife et al., 1980; Reiner & Antzelevitch, 1985\].
The rhythm of autonomous biological oscillators can also undergo an external periodic perturbation (e.g. activity of cells of the AV junction is subjected to sinus rhythm), depending on both the stimulus magnitude and its phase within the cycle. It is known that when the frequency and the amplitude of the external periodic stimulation are varied, a diversity of phase diagrams can be established between the stimulus and the self-sustained oscillator (see e.g. \[Loskutov et al., 2004\]). In some situations the rhythm of the biological oscillator is entrained (or phase-locked) to the external stimulation so that for each $`M`$ cycles of the stimulation there are $`N`$ cycles of the autonomous oscillator rhythm. This occurs at a fixed phase (or phases) of the stimulus and is called $`M:N`$ phase-locking or entrainment, which appears as a time–periodic sequence. In particular, entrainment of $`1:1`$, during which the rhythms of the oscillator and external stimulus are matched, is defined as synchronization phenomenon.
In the present paper we develop a general simplified model describing a network of pulse oscillators coupled by their response to the internal depolarization of mutual stimulations. Our primary aim is to keep the model as simple as possible and to introduce a minimal number of parameters. Therefore, in our model the pacemakers are fully characterized by their intrinsic cycle length and are represented as pulse oscillators. Their interaction is described by PRCs. At first, we will consider two bidirectionally interacting pacemakers to demonstrate the basic concepts of the model. Then we will apply this approach to construct a pacemaker network model with global coupling. As the following step, we will analyze two specific cases of this PRC based model of coupled pulse oscillators: two and then three interacting cardiac nodes. An additional pacemaker can also be expounded as an external stimulater. Our further intention is to go on to the next (microscopic) level and represent each pacemaker as an ensemble of interacting oscillatory elements. Extrapolation of our approach to one– and two–dimensional matrixes (or lattices) of pacemaker cells interacting via nearest neighbors concludes the present study.
## 2 Development of the General Model
In this part we construct a system with two pacemakers and then consider a quite general model of a set of interacting pacemakers coupled by their PRCs.
### 2.1 Two interacting pacemakers: outline of the approach
Consider two interacting pulse oscillators (or pacemakers) $`A`$ and $`B`$ with intrinsic periods of their autonomous beating $`T_a`$ and $`T_b`$ respectively. An interaction between oscillators is governed by so-called phase response curve (PRC). This means that a phase shift of one of the oscillators happens as a result of an impact of the another one. To construct an adequate mathematical model, it is necessary to accept some restrictions concerning the character of the interaction. We describe them briefly \[Ikeda, 1982; Glass et al., 1986; Glass & Zeng, 1990\].
1. The phase of the disturbed oscillator is shifted to a new value instantly after an impact.
2. The phase shift depends only on two main parameters: a) on the phase difference of oscillators and b) on the influence strength. In turn, this influence strength depends on its amplitude and the coupling coefficient of the oscillators. In a real system the coupling coefficient is the average factor that shows how the strength of the pulse decreases during its passing from one oscillator to another. Thus, the phase shift $`\mathrm{\Phi }`$ determining a new phase of the disturbed oscillator with period $`T`$ can be represented as follows:
$$\mathrm{\Phi }=\mathrm{\Delta }/T\mathrm{\Phi }(\phi ,\epsilon ),$$
(1)
where $`\mathrm{\Delta }`$ is the time shift of the disturbed oscillator, $`\phi `$ is the phase difference of the oscillators and $`\epsilon `$ is the influence strength.
Pacemakers can be represented as a set of separated firing peaks on a time scale. Assume that the instants of the last firings of the oscillators $`A`$ and $`B`$ are $`a`$ and $`b`$ respectively (Fig. 1). Note that $`a`$ and $`b`$ are the moments of the impacts after all previous phase shifts of the oscillators. In other words, one can observe the oscillators’ firings at these particular instants. Then it is necessary to analyze two cases:
1. $`b<a`$. This is the case 1 in Fig. 1, i.e. when the oscillator $`B`$ has fired before the oscillator $`A`$. Let us follow the dynamics of the system in real time. The nearest event, that affects on the further behavior of the entire system, is an appearance of the pulse $`A`$ at time $`a`$. Let us stop on at this moment and make the forecast. To this end we define a concept of the moments of expected firings of the oscillators, i.e. $`a^e`$ and $`b^e`$. Let us imagine that we have shifted back in time with respect to the moment $`a`$. Since the oscillator $`A`$ has not fired yet, one should expect the appearance of the next $`A`$ and $`B`$ pulses at the moments $`a^e=a`$ and $`b^e=b+T_b`$ respectively, where $`T_b`$ is the period of $`B`$. We call this situation as “$`A`$ fires and $`B`$ is at an expected state” and denote symbolically as $`(a,b^e)`$. Now we consider the instant $`a`$. Since $`A`$ fires, the next expected values can be transformed to
$$\begin{array}{ccccc}a_{next}^e\hfill & =\hfill & a+T_a\hfill & =\hfill & a^e+T_a,\hfill \\ b_{next}^e\hfill & =\hfill & b+T_b+\mathrm{\Delta }_b(\phi _b,\epsilon _b)\hfill & =\hfill & b^e+\mathrm{\Delta }_b(\phi _b,\epsilon _b),\hfill \end{array}$$
where $`\mathrm{\Delta }_b(\phi _b,\epsilon _b)`$ is the time shift of the oscillator $`B`$ due to the impact of $`A`$. It depends on the phase $`\phi _b`$ of the pacemaker $`A`$ with respect to $`B`$ and the influence strength $`\epsilon _b`$. The phase $`\phi _b`$ can be calculated as follows
$$\phi _b=\frac{ab}{T_b}(\text{mod}\mathrm{\hspace{0.33em}1})$$
or, in terms of the expected values,
$$\phi _b=\frac{a^eb^e}{T_b}(\text{mod}\mathrm{\hspace{0.33em}1}).$$
Note that the phase $`\phi _b`$ is a positive value, and it belongs to the $`[0,1]`$ segment (negative values in the two previous expressions are eliminated by the mod 1 operation).
To determine which oscillator fires next, one should compare $`a_{next}^e`$ and $`b_{next}^e`$. If $`a_{next}^e<b_{next}^e`$, then $`A`$ fires and $`B`$ remains at an expected state until $`b_{next}^e`$, i.e. the system moves to the state $`(a,b^e)_{next}`$. Otherwise, if $`b_{next}^e<a_{next}^e`$, then $`B`$ fires, and $`A`$ jumps to an expected state and the entire system’s state becomes $`(a^e,b)_{next}`$.
2. $`b>a`$. This is the case 2 in Fig. 1. This inverse situation is analogous to the previous one with the difference in speculations owing to the firing of the pacemaker $`A`$ prior to $`B`$. Then expected values $`a^e`$ and $`b^e`$ can be written accordingly: $`a^e=a+T_a`$ and $`b^e=b`$. One can call this case as “$`B`$ fires and $`A`$ is at an expected state” and denote by $`(a^e,b)`$. The next expected values are given by the following expression:
$$\begin{array}{ccccc}a_{next}^e\hfill & =\hfill & a^e+\mathrm{\Delta }_a(\phi _a,\epsilon _a),\hfill & & \\ b_{next}^e\hfill & =\hfill & b^e+T_b,\hfill & & \end{array}$$
where $`\mathrm{\Delta }_a(\phi _a,\epsilon _a)`$ is the time shift of the oscillator $`A`$. It depends on the phase of the oscillator $`B`$ with respect to $`A`$, i.e. $`\phi _a=(b^ea^e)/T_a(\text{mod}\mathrm{\hspace{0.33em}1})`$, and the influence strength $`\epsilon _a`$. Further analysis is also similar to the case 1. Namely, if $`b_{next}^e<a_{next}^e`$, then $`B`$ fires and $`A`$ jumps to the expected state $`a_{next}^e`$, i.e. the system moves to the state $`(a^e,b)_{next}`$. Otherwise, if $`a_{next}^e<b_{next}^e`$, then $`A`$ fires and $`B`$ remains in the expected state and the system state becomes $`(a,b^e)_{next}`$.
Summarizing the above calculations, the model can be represented by the following scheme:
$$\begin{array}{cccccccccccccccccccc}\begin{array}{cccccccccccccccccccc}\left(\begin{array}{cccccccccccccccccccc}a& & & & & & & & & & & & & & & & & & & \\ b^e& & & & & & & & & & & & & & & & & & & \end{array}\right)& & \{\begin{array}{cccccccccccccccccccc}a_{next}^e=a+T_a& & & & & & & & & & & & & & & & & & & \\ b_{next}^e=b^e+\mathrm{\Delta }_b(\phi _b,\epsilon _b)& & & & & & & & & & & & & & & & & & & \end{array}& & & & & & & & & & & & & & & & & \\ & & & & & & & & & & & & & & & & & & & \\ \left(\begin{array}{cccccccccccccccccccc}a^e& & & & & & & & & & & & & & & & & & & \\ b& & & & & & & & & & & & & & & & & & & \end{array}\right)& & \{\begin{array}{cccccccccccccccccccc}a_{next}^e=a^e+\mathrm{\Delta }_a(\phi _a,\epsilon _a)& & & & & & & & & & & & & & & & & & & \\ b_{next}^e=b+T_b& & & & & & & & & & & & & & & & & & & \end{array}& & & & & & & & & & & & & & & & & \end{array}\}& & \{\begin{array}{cccccccccccccccccccc}\left(\begin{array}{cccccccccccccccccccc}a& & & & & & & & & & & & & & & & & & & \\ b^e& & & & & & & & & & & & & & & & & & & \end{array}\right)_{next}\text{if}a_{next}^e<b_{next}^e& & & & & & & & & & & & & & & & & & & \\ & & & & & & & & & & & & & & & & & & & \\ \left(\begin{array}{cccccccccccccccccccc}a^e& & & & & & & & & & & & & & & & & & & \\ b& & & & & & & & & & & & & & & & & & & \end{array}\right)_{next}\text{if}a_{next}^e>b_{next}^e& & & & & & & & & & & & & & & & & & & \end{array}& & & & & & & & & & & & & & & & & \end{array}$$
(2)
In the notions of the expected values, which we denote for convenience as $`a^e\widehat{a}`$ $`b^e\widehat{b}`$, the dynamics can be described by the following difference equation:
$$\left(\begin{array}{cccccccccccccccccccc}\widehat{a}_{n+1}& & & & & & & & & & & & & & & & & & & \\ \widehat{b}_{n+1}& & & & & & & & & & & & & & & & & & & \end{array}\right)=\left(\begin{array}{cccccccccccccccccccc}\widehat{a}_n& & & & & & & & & & & & & & & & & & & \\ \widehat{b}_n& & & & & & & & & & & & & & & & & & & \end{array}\right)+\{\begin{array}{cccccccccccccccccccc}\left(\begin{array}{cccccccccccccccccccc}T_a& & & & & & & & & & & & & & & & & & & \\ \mathrm{\Delta }_b(\phi _n^b,\epsilon _b)& & & & & & & & & & & & & & & & & & & \end{array}\right),\widehat{a}_n<\widehat{b}_n,\text{and then }A\text{ fires at time}\widehat{a}_n& & & & & & & & & & & & & & & & & & & \\ & & & & & & & & & & & & & & & & & & & \\ \left(\begin{array}{cccccccccccccccccccc}\mathrm{\Delta }_a(\phi _n^a,\epsilon _a)& & & & & & & & & & & & & & & & & & & \\ T_b& & & & & & & & & & & & & & & & & & & \end{array}\right),\widehat{b}_n<\widehat{a}_n,\text{and then }B\text{ fires at time}\widehat{b}_n,& & & & & & & & & & & & & & & & & & & \end{array}$$
(3)
where:
$$\phi _n^a=\frac{\left(\widehat{b}_n\widehat{a}_n\right)}{T_a}\left(\text{mod}\mathrm{\hspace{0.33em}1}\right),\phi _n^b=\frac{\left(\widehat{a}_n\widehat{b}_n\right)}{T_b}\left(\text{mod}\mathrm{\hspace{0.33em}1}\right).$$
To simulate the dynamics of two interacting oscillators coupled by PRCs, it is necessary to carry out the iterative process (3) for the expected pulses and put sequentially in the time scale the firing moments of the oscillators $`A`$ and $`B`$ depending on the result of the comparison $`\widehat{a}_n`$ and $`\widehat{b}_n`$. An investigation of the case of two interacting pacemakers in more detail with the description of the possible modes of behavior is given in Section 3.1.
### 2.2 Derivation of the basic model equation
Assume that there are $`N`$ autonomous pulse oscillators (or pacemakers). Suppose that all the pacemakers are different. This means that each has its own intrinsic cycle length $`T_i`$, $`i=1,\mathrm{},N`$, and the beating amplitude. To define coupling between pacemakers, it is necessary to determine the topology of the system space. In other words, one should specify the nearest neighbors of each pacemaker in the space. Vice versa, it is obvious that the determination of coupling between the elements of such spatially discrete system sets its topology.
First of all we develop the general model of $`N`$ mutually coupled pulse oscillators. Suppose that all the pacemakers interact with each other, i.e. so-called global coupling is realized. The model derived in Section 2.1 can be easily generalized to the case of $`N`$ pacemakers. We operate with the expected values introduced in Section 2.1, the real firings of the pacemakers are found by the analysis of the expected impacts series.
Let the set of expected firings $`\{\widehat{a}_i\}_{1,\mathrm{},N}`$ be located in a time axis (see Fig. 2). This means that in the absence of coupling, pacemakers strike at these instants. Suppose now that some oscillator acts on another one by means of the PRC $`\{\mathrm{\Delta }_{ij}(\phi ^{ij},\epsilon _{ij})\}_{1,\mathrm{},N}`$, where $`\phi ^{ij}`$ is the phase of the $`j`$-th pacemaker with respect to the $`i`$-th one and $`\epsilon _{ij}`$ is a total parameter defining coupling between the $`j`$-th and $`i`$-th elements.
The next values of the expected firings can be calculated by the same manner as in Section 2.1. Because the $`j`$-th oscillator appears before all others (see. Fig. 2), it does not undergo any influence and fires as a real impact of the $`j`$-th pacemaker. Thus, the given event makes the shifts of all other oscillators according to the set of the PRCs $`\{\mathrm{\Delta }_{ij}\}`$. The $`j`$-th oscillator is shifted to a new expected moment as an unperturbed one, i.e. by adding its own cycle length $`T_j`$. To get the next sequential expected values, one should make the same procedure with the newly obtained expected pulses. The dynamics of the system can be easily represented as the following iterative relation:
$$\widehat{a}_{n+1}^i=\widehat{a}_n^i+\{\begin{array}{cccccccccccccccccccc}T_i,& i=j,& & & & & & & & & & & & & & & & & & \\ \mathrm{\Delta }_{ij}(\phi _n^{ij},\epsilon _{ij}),& ij,& & & & & & & & & & & & & & & & & & \end{array}j:a_n^j=\mathrm{min}\{a_n^i\}_{i=1,\mathrm{},N},$$
(4)
where
$$\phi _n^{ij}=\frac{\widehat{a}_n^j\widehat{a}_n^i}{T_i}\left(\text{mod}\mathrm{\hspace{0.33em}1}\right),$$
It is convenient to rewrite the PRCs by normalizing $`\{\mathrm{\Delta }_{ij}\}`$ on the intrinsic pacemaker cycle lengths $`T_i`$ and to define as follows:
$$\mathrm{\Delta }_{ij}(\phi _n^{ij},\epsilon _{ij})=\{f_{ij}(\phi _n^{ij},\epsilon _{ij})T_i\},\phi _n^{ij}[0,1],$$
where $`\{f_{ij}(\phi _n^{ij},\epsilon _{ij})\}`$ are the dimensionless functions, which are also called the PRCs. The real pacemakers can have identical nature but differ in the intrinsic cycle lengths. Then the form of dimensionless PRCs is identical for them, i.e. their $`f(\phi ,\epsilon )`$ coincide, while $`\mathrm{\Delta }(\phi ,\epsilon )`$ are different. Moreover, it is convenient to use the functions $`\{f_{ij}(\phi _n^{ij},\epsilon _{ij})\}`$ in the construction of the equations for the dimensionless phase differences between pacemaker pairs. In Section 3 we will demonstrate this approach for the systems of two and three coupled pulse oscillators.
## 3 Applications
### 3.1 Two interacting pulse oscillators
Let us investigate the system of two interacting pacemakers in detail. As well as in Section 2.1, suppose that the system consists of two oscillators $`A`$ and $`B`$ coupled by means of the phase response curves $`\mathrm{\Delta }_a(\phi _n^a,\epsilon _a)`$ and $`\mathrm{\Delta }_b(\phi _n^b,\epsilon _b)`$. It is clear that the map (3) specifying the dynamics of such a system is a particular case of the general model (4) restricted to $`N=2`$. Rewrite (3) using the expressions for the dimensionless PRCs:
$$\left(\begin{array}{cccccccccccccccccccc}\widehat{a}_{n+1}& & & & & & & & & & & & & & & & & & & \\ \widehat{b}_{n+1}& & & & & & & & & & & & & & & & & & & \end{array}\right)=\left(\begin{array}{cccccccccccccccccccc}\widehat{a}_n& & & & & & & & & & & & & & & & & & & \\ \widehat{b}_n& & & & & & & & & & & & & & & & & & & \end{array}\right)+\{\begin{array}{cccccccccccccccccccc}\left(\begin{array}{cccccccccccccccccccc}T_a& & & & & & & & & & & & & & & & & & & \\ f_b(\phi _n^b,\epsilon _b)T_b& & & & & & & & & & & & & & & & & & & \end{array}\right),\widehat{a}_n<\widehat{b}_n,\text{and then }A\text{ fires at time}\widehat{a}_n,& & & & & & & & & & & & & & & & & & & \\ & & & & & & & & & & & & & & & & & & & \\ \left(\begin{array}{cccccccccccccccccccc}f_a(\phi _n^a,\epsilon _a)T_a& & & & & & & & & & & & & & & & & & & \\ T_b& & & & & & & & & & & & & & & & & & & \end{array}\right),\widehat{b}_n<\widehat{a}_n,\text{and then }B\text{ fires at time}\widehat{b}_n,& & & & & & & & & & & & & & & & & & & \end{array}$$
(5)
where:
$$\phi _n^a=\frac{\left(\widehat{b}_n\widehat{a}_n\right)}{T_a}\left(\text{mod}\mathrm{\hspace{0.33em}1}\right),\phi _n^b=\frac{\left(\widehat{a}_n\widehat{b}_n\right)}{T_b}\left(\text{mod}\mathrm{\hspace{0.33em}1}\right).$$
Let us assume that the pacemakers have one the same nature. Hence, one can accept $`f_a(\phi ,\epsilon )f_b(\phi ,\epsilon )f(\phi ,\epsilon )`$, where $`f(\phi ,\epsilon )`$ is an one-parametric function. The parameter $`\epsilon `$ integrally determines a total influence of one oscillator on another one. In the symmetric case we have $`\epsilon _a=\epsilon _b\epsilon `$. In our further consideration we do not accept this symmetry.
From the mathematical point of view the system (5) is a map of the real plane into itself depending on four parameters: $`T_a`$, $`T_b`$, $`\epsilon _a`$ and $`\epsilon _b`$. It is easy to comprehend that the effective parameter that changes the behavior of the system (5) is the dimensionless ratio $`\delta =T_b/T_a`$. Specifying parameters $`\epsilon _a`$, $`\epsilon _b`$, $`\delta `$ and the function $`f(\phi ,\epsilon )`$, one can iterate directly the equations (5) marking on the time axis the values of the real firings of the pacemakers $`A`$ and $`B`$. Some examples of this simulation are presented below.
In spite of the fact that the expression (5) defines the two-dimensional map, the values $`\widehat{a}_n`$ and $`\widehat{b}_n`$ increase gradually because they represent the time series of the expected firings of the oscillators. Therefore, trajectories of the map (5) are infinite. Thus, this is not informative for us. It is more important to derive a map that determines the dynamics of the phase difference of the pacemakers.
Let us introduce the dimensionless phase difference of the pacemakers $`A`$ and $`B`$:
$$x_n=\frac{\widehat{a}_n\widehat{b}_n}{T_a}.$$
The choice of $`T_a`$ as a normalization factor is not essential. Having selected $`T_b`$, as a result we obtain the similar expressions. Subtracting the second equation of the system (5) from the first one and dividing the result by $`T_a`$ we get the following:
$$x_{n+1}=\{\begin{array}{cc}x_n+1\delta f(\frac{x_n}{\delta }(\text{mod}\mathrm{\hspace{0.17em}1}),\epsilon _b),\hfill & x_n<0,\hfill \\ x_n+f(x_n(\text{mod}\mathrm{\hspace{0.17em}1}),\epsilon _a)\delta ,\hfill & x_n>0,\hfill \end{array}$$
(6)
where $`\delta =T_b/T_a`$. Here we take into account that $`x_n<0`$ if $`\widehat{a}_n<\widehat{b}_n`$ and $`x_n>0`$ if $`\widehat{a}_n>\widehat{b}_n`$.
Let us make a brief analysis of the developed map. Because in general $`x(\mathrm{};\mathrm{})`$, the equation (6) represents a one-dimensional nonlinear map of the real axis into itself. Note that the map can not be reduced to a circle map by the restriction of $`x`$ to the range $`[0;1]`$ as it is usually done for two pacemakers interacting by PRCs. It is essentially asymmetric with respect to changing $`x`$ to $`x`$ (see Fig. 3). If $`f(x,\epsilon )`$ is a nonmonotonic function, then the map (6) is nonlinear, and it can exhibit a big variety of the behavior: from complex periodic motion to chaotic dynamics. Owing to $`x(\mathrm{};\mathrm{})`$, in a rigorous sense it is not a difference of phases of the pacemakers. In this context one can call $`x`$ as the generalized phase difference. Analyzing Eq. (6) one can find out which oscillator, $`A`$ or $`B`$, fires at the given discrete time $`n`$. This depends on the sign of $`x`$: $`A`$ fires if $`x_n<0`$ and $`B`$ fires when $`x_n>0`$. Thus, it makes possible to determine the phase-locking degree of the pacemakers. However, taking into consideration only the values $`x_n`$, we may not reconstruct the initial time series of firing events of the pacemakers $`A`$ and $`B`$. Using $`x_n`$, one can say only about their phase difference.
Let us show how the model equations (5), (6) can be applied to the investigation of the behavior of two interacting pacemakers.
Analysis of real systems \[Glass et al., 1987\] shows that the function $`f(x,\epsilon )`$ can take different forms. But, as a rule, it obeys a number of general properties. For example, $`f(0,\epsilon )=f(1,\epsilon )=0`$. Usually it has a maximum and a minimum. Sometimes instead of extrema it has breaks. Let the function $`f(x,\epsilon )`$ be in an elementary form that is often used (see, e.g. \[Glass & Zeng, 1990\]). Namely, let $`f(x,\epsilon )=\epsilon \mathrm{sin}2\pi x`$. This leads to an array of dynamical equations:
$`\left(\begin{array}{cccccccccccccccccccc}\widehat{a}_{n+1}& & & & & & & & & & & & & & & & & & & \\ \widehat{b}_{n+1}& & & & & & & & & & & & & & & & & & & \end{array}\right)=\left(\begin{array}{cccccccccccccccccccc}\widehat{a}_n& & & & & & & & & & & & & & & & & & & \\ \widehat{b}_n& & & & & & & & & & & & & & & & & & & \end{array}\right)+\{\begin{array}{cccccccccccccccccccc}\left(\begin{array}{cccccccccccccccccccc}T_a& & & & & & & & & & & & & & & & & & & \\ \epsilon _b\mathrm{sin}\left(2\pi \phi _n^b\right)T_b& & & & & & & & & & & & & & & & & & & \end{array}\right),\widehat{a}_n<\widehat{b}_n,& & & & & & & & & & & & & & & & & & & \\ \left(\begin{array}{cccccccccccccccccccc}\epsilon _a\mathrm{sin}\left(2\pi \phi _n^a\right)T_a& & & & & & & & & & & & & & & & & & & \\ T_b& & & & & & & & & & & & & & & & & & & \end{array}\right),\widehat{b}_n<\widehat{a}_n.& & & & & & & & & & & & & & & & & & & \end{array}`$ (17)
$`x_{n+1}=\{\begin{array}{cc}x_n+1\delta \epsilon _b\mathrm{sin}({\displaystyle \frac{x_n}{\delta }}(\text{mod}\mathrm{\hspace{0.17em}1})2\pi ),\hfill & x_n<0,\hfill \\ x_n+\epsilon _a\mathrm{sin}(x_n(\text{mod}\mathrm{\hspace{0.17em}1})2\pi )\delta ,\hfill & x_n>0.\hfill \end{array}`$ (20)
In Fig. 3 examples of direct simulation of the system (17), (20) are presented. In the left column some possible phase-lockings found on the basis of Eq. (17) are displayed. In the right column the corresponding map (20), its periodic orbit and values of Lyapunov exponents are shown. Fig. 3d represents the existence of the chaotic behavior for the system at $`\epsilon _a=0.1205`$, $`\epsilon _b=1.2845`$ and $`\delta =0.6465.`$
As a conclusion of this section it should be noted that the system of two bidirectionally acting pacemakers has been already intensively investigated (see, e.g., \[Ikeda et al., 2004\] and references herein) as a qualitative model of cardiac arrhythmia known as modulated parasystole. The interacting oscillators represented the pairs of cardiac pacemakers such as: the sinoatrial (SA) node and the ventricle contracting by different factors, e.g. ectopic pacemaker, etc. Hereby the authors used various kinds of PRCs $`f(x,\epsilon )`$ approximating the experimental data on stimulating the cardiac cells of animals by single electric current pulses. Recently in the paper \[Loskutov et al., 2004\] the system of two interacting pacemakers similar to (17), (20) has been analyzed. In this paper taking into account the refractory time the various types of the smooth PRCs were examined and the bifurcation diagrams of possible phase-lockings were constructed. However, in \[Loskutov et al., 2004\] the behavior regimes when the firings of the pacemakers are not alternated, i.e. cases of strong discrepancies in the intrinsic cycle length ($`T_aT_b`$ and vice versa), were not investigated. The system (17), (20) is more general and takes into account all possible variants.
### 3.2 Three pacemakers
Below we explain why the special case of three interacting oscillators is worth individual attention. First, as is known there are three drivers of the rhythm in the cardiac conductive system: the SA node is a leading pacemaker, the AV junction and Purkinje fibers are latent pacemakers, which under normal conditions are suppressed by the sinus rhythm. However, at violations of conductivity and pulse initiation, the cardiac pacemakers can influence on each other, i.e. so-called bidirectional coupling may be realized. Second, in pathology a group of contractile cardiomyocytes can also initiate action potentials: an ectopic (abnormal) cardiac pacemaker may emerge and start to compete with the SA node for leading the heart rhythm. Third, at stimulating the cardiac tissue by external current impulses (cardio-stimulation), the cardiac rate is changed. The external stimulaters can be naturally included in our general model of interacting pacemakers as additional leading pulse oscillators.
Thus, representation of the heart conductive system as at least three coupled autonomous oscillators (Fig. 4) is very useful for understanding which influence of an additional pacemaker exerts on a system of two bidirectionally interacting drivers of the rhythm (i.e. the SA and AV nodes) considered above. Investigation of the possible behavior modes of a larger amount of interacting pacemakers turns out to be a very complicated both analytical (mathematical) and numerical problem. For example for a system of five coupled pulse oscillators we have 25 various functions and 29 different parameters (25 values of $`\epsilon _{ij}`$ and 4 independent ratios $`T_i/T_j`$). This is due to the fact that, in general, the cardiac pacemakers have various frequencies (or intrinsic cycle lengths $`T_i`$ ) and the different nature. This means that the PRCs $`\{f_{ij}(x,\epsilon _{ij})\}`$ determining coupling between a pair of pacemakers have different forms. For some heart pacemakers PRCs have been measured by the direct experiments on cardiomyocytes of animals (see \[Glass et al., 1986\]). Other PRCs can be chosen using general principles based on the nature of nodes or on the basis of the collateral measurements \[Ikeda et al., 1988\].
Applying the general model equations (4) for three interacting pacemakers, one can get the following system:
$$\left(\begin{array}{cccccccccccccccccccc}\widehat{a}_{n+1}& & & & & & & & & & & & & & & & & & & \\ \widehat{b}_{n+1}& & & & & & & & & & & & & & & & & & & \\ \widehat{c}_{n+1}& & & & & & & & & & & & & & & & & & & \end{array}\right)=\left(\begin{array}{cccccccccccccccccccc}\widehat{a}_n& & & & & & & & & & & & & & & & & & & \\ \widehat{b}_n& & & & & & & & & & & & & & & & & & & \\ \widehat{c}_n& & & & & & & & & & & & & & & & & & & \end{array}\right)+\{\begin{array}{cccccccccccccccccccc}\left(\begin{array}{cccccccccccccccccccc}T_a& & & & & & & & & & & & & & & & & & & \\ \mathrm{\Delta }_{ba}(\phi _n^{ba},\epsilon _{ba})& & & & & & & & & & & & & & & & & & & \\ \mathrm{\Delta }_{ca}(\phi _n^{ca},\epsilon _{ca})& & & & & & & & & & & & & & & & & & & \end{array}\right),\text{if}\widehat{a}_n<\widehat{b}_n,\widehat{c}_n;A\text{fires at time}\widehat{a}_n& & & & & & & & & & & & & & & & & & & \\ & & & & & & & & & & & & & & & & & & & \\ \left(\begin{array}{cccccccccccccccccccc}\mathrm{\Delta }_{ab}(\phi _n^{ab},\epsilon _{ab})& & & & & & & & & & & & & & & & & & & \\ T_b& & & & & & & & & & & & & & & & & & & \\ \mathrm{\Delta }_{cb}(\phi _n^{cb},\epsilon _{cb})& & & & & & & & & & & & & & & & & & & \end{array}\right),\text{if}\widehat{b}_n<\widehat{a}_n,\widehat{c}_n;B\text{fires at time}\widehat{b}_n& & & & & & & & & & & & & & & & & & & \\ & & & & & & & & & & & & & & & & & & & \\ \left(\begin{array}{cccccccccccccccccccc}\mathrm{\Delta }_{ac}(\phi _n^{ac},\epsilon _{ac})& & & & & & & & & & & & & & & & & & & \\ \mathrm{\Delta }_{bc}(\phi _n^{bc},\epsilon _{bc})& & & & & & & & & & & & & & & & & & & \\ T_c& & & & & & & & & & & & & & & & & & & \end{array}\right),\text{if}\widehat{c}_n<\widehat{a}_n,\widehat{b}_n;C\text{fires at time}\widehat{c}_n& & & & & & & & & & & & & & & & & & & \end{array}$$
(21)
where
$$\begin{array}{c}\phi _n^{ba}=\frac{\widehat{a}_n\widehat{b}_n}{T_b}\left(\text{mod}\mathrm{\hspace{0.33em}1}\right),\phi _n^{ca}=\frac{\widehat{a}_n\widehat{c}_n}{T_c}\left(\text{mod}\mathrm{\hspace{0.33em}1}\right)\hfill \\ \phi _n^{ab}=\frac{\widehat{b}_n\widehat{a}_n}{T_a}\left(\text{mod}\mathrm{\hspace{0.33em}1}\right),\phi _n^{cb}=\frac{\widehat{b}_n\widehat{c}_n}{T_c}\left(\text{mod}\mathrm{\hspace{0.33em}1}\right)\hfill \\ \phi _n^{ac}=\frac{\widehat{c}_n\widehat{a}_n}{T_a}\left(\text{mod}\mathrm{\hspace{0.33em}1}\right),\phi _n^{bc}=\frac{\widehat{c}_n\widehat{b}_n}{T_b}\left(\text{mod}\mathrm{\hspace{0.33em}1}\right)\hfill \end{array}$$
The response functions $`\mathrm{\Delta }_{ij}`$ are supposed to have a form $`\mathrm{\Delta }_{ij}(\phi _n^{ij},\epsilon _{ij})=f_{ij}(\phi _n^{ij},\epsilon _{ij})T_i,\phi _n^{ij}[0;1],i,j=a,b,c`$. The example of a system of three bidirectionally interacting cardiac pacemakers: the SA node, the AV junction and ectopic pacemaker, is shown in Fig. 4a.
It is convenient to study Eqs. (21) introducing phase differences of the pacemakers. Let us define,
$$x_n=\frac{\widehat{a}_n\widehat{b}_n}{T_b},y_n=\frac{\widehat{c}_n\widehat{b}_n}{T_b},\alpha =\frac{T_a}{T_b},\beta =\frac{T_c}{T_b}.$$
Then for the phase differences we obtain the map of the real plane into itself:
$$\begin{array}{cccccccccccccccccccc}\left(\begin{array}{cccccccccccccccccccc}x_{n+1}& & & & & & & & & & & & & & & & & & & \\ y_{n+1}& & & & & & & & & & & & & & & & & & & \end{array}\right)=\left(\begin{array}{cccccccccccccccccccc}x_n& & & & & & & & & & & & & & & & & & & \\ y_n& & & & & & & & & & & & & & & & & & & \end{array}\right)+\hfill & & & & & & & & & & & & & & & & & & & \\ & & & & & & & & & & & & & & & & & & & \\ +\{\begin{array}{cccccccccccccccccccc}\left(\begin{array}{cccccccccccccccccccc}\alpha f_{ba}(\left\{x_n\right\},\epsilon _{ba})& & & & & & & & & & & & & & & & & & & \\ \beta f_{ca}(\left\{\frac{x_ny_n}{\beta }\right\},\epsilon _{ca})f_{ba}(\left\{x_n\right\},\epsilon _{ba})& & & & & & & & & & & & & & & & & & & \end{array}\right),x_n<0,x_n<y_n& & & & & & & & & & & & & & & & & & & \\ & & & & & & & & & & & & & & & & & & & \\ \left(\begin{array}{cccccccccccccccccccc}\alpha f_{ab}(\left\{\frac{x_n}{\alpha }\right\},\epsilon _{ab})1& & & & & & & & & & & & & & & & & & & \\ \beta f_{cb}(\left\{\frac{y_n}{\beta }\right\},\epsilon _{cb})1& & & & & & & & & & & & & & & & & & & \end{array}\right),x_n>0,y_n>0& & & & & & & & & & & & & & & & & & & \\ & & & & & & & & & & & & & & & & & & & \\ \left(\begin{array}{cccccccccccccccccccc}\alpha f_{ac}(\left\{\frac{y_nx_n}{\alpha }\right\},\epsilon _{ac})f_{bc}(\left\{y_n\right\},\epsilon _{bc})& & & & & & & & & & & & & & & & & & & \\ \beta f_{bc}(\left\{y_n\right\},\epsilon _{bc})& & & & & & & & & & & & & & & & & & & \end{array}\right),y_n<x_n,y_n<0& & & & & & & & & & & & & & & & & & & \end{array}\hfill & & & & & & & & & & & & & & & & & & & \end{array}$$
(22)
Here we denote $`\{x\}`$ as the fractional part of $`x`$. The conditions in the right hand side of (22) divide the plane into three areas corresponding to the firing pacemaker (see Fig. 5). The map (22) has breaks on the boundaries of these areas.
Suppose that the pacemaker $`B`$ is an external stimulus (see Fig. 4b). Then the expressions (21), (22) become simpler. The stimulus $`B`$ acts on the pacemakers $`A`$ and $`C`$ but does not experience any influence. In this situation $`f_{ba}(x,\epsilon _{ba})f_{bc}(x,\epsilon _{bc})0`$ and the model (21), (22) can be rewritten as follows:
$$\left(\begin{array}{cccccccccccccccccccc}\widehat{a}_{n+1}& & & & & & & & & & & & & & & & & & & \\ \widehat{b}_{n+1}& & & & & & & & & & & & & & & & & & & \\ \widehat{c}_{n+1}& & & & & & & & & & & & & & & & & & & \end{array}\right)=\left(\begin{array}{cccccccccccccccccccc}\widehat{a}_n& & & & & & & & & & & & & & & & & & & \\ \widehat{b}_n& & & & & & & & & & & & & & & & & & & \\ \widehat{c}_n& & & & & & & & & & & & & & & & & & & \end{array}\right)+\{\begin{array}{cccccccccccccccccccc}\left(\begin{array}{cccccccccccccccccccc}T_a& & & & & & & & & & & & & & & & & & & \\ 0& & & & & & & & & & & & & & & & & & & \\ \mathrm{\Delta }_{ca}(\phi _n^{ca},\epsilon _{ca})& & & & & & & & & & & & & & & & & & & \end{array}\right),\text{if}\widehat{a}_n<\widehat{b}_n,\widehat{c}_n;A\text{fires at time}\widehat{a}_n& & & & & & & & & & & & & & & & & & & \\ & & & & & & & & & & & & & & & & & & & \\ \left(\begin{array}{cccccccccccccccccccc}\mathrm{\Delta }_{ab}(\phi _n^{ab},\epsilon _{ab})& & & & & & & & & & & & & & & & & & & \\ T_b& & & & & & & & & & & & & & & & & & & \\ \mathrm{\Delta }_{cb}(\phi _n^{cb},\epsilon _{cb})& & & & & & & & & & & & & & & & & & & \end{array}\right),\text{if}\widehat{b}_n<\widehat{a}_n,\widehat{c}_n;B\text{fires at time}\widehat{b}_n& & & & & & & & & & & & & & & & & & & \\ & & & & & & & & & & & & & & & & & & & \\ \left(\begin{array}{cccccccccccccccccccc}\mathrm{\Delta }_{ac}(\phi _n^{ac},\epsilon _{ac})& & & & & & & & & & & & & & & & & & & \\ 0& & & & & & & & & & & & & & & & & & & \\ T_c& & & & & & & & & & & & & & & & & & & \end{array}\right),\text{if}\widehat{c}_n<\widehat{a}_n,\widehat{b}_n;C\text{fires at time}\widehat{c}_n& & & & & & & & & & & & & & & & & & & \end{array}$$
(23)
The corresponding expression for the phase differences takes the form:
$$\begin{array}{cccccccccccccccccccc}\left(\begin{array}{cccccccccccccccccccc}x_{n+1}& & & & & & & & & & & & & & & & & & & \\ y_{n+1}& & & & & & & & & & & & & & & & & & & \end{array}\right)=\left(\begin{array}{cccccccccccccccccccc}x_n& & & & & & & & & & & & & & & & & & & \\ y_n& & & & & & & & & & & & & & & & & & & \end{array}\right)+\hfill & & & & & & & & & & & & & & & & & & & \\ & & & & & & & & & & & & & & & & & & & \\ +\{\begin{array}{cccccccccccccccccccc}\left(\begin{array}{cccccccccccccccccccc}\alpha & & & & & & & & & & & & & & & & & & & \\ \beta f_{ca}(\left\{\frac{x_ny_n}{\beta }\right\},\epsilon _{ca})& & & & & & & & & & & & & & & & & & & \end{array}\right),x_n<0,x_n<y_n& & & & & & & & & & & & & & & & & & & \\ & & & & & & & & & & & & & & & & & & & \\ \left(\begin{array}{cccccccccccccccccccc}\alpha f_{ab}(\left\{\frac{x_n}{\alpha }\right\},\epsilon _{ab})1& & & & & & & & & & & & & & & & & & & \\ \beta f_{cb}(\left\{\frac{y_n}{\beta }\right\},\epsilon _{cb})1& & & & & & & & & & & & & & & & & & & \end{array}\right),x_n>0,y_n>0& & & & & & & & & & & & & & & & & & & \\ & & & & & & & & & & & & & & & & & & & \\ \left(\begin{array}{cccccccccccccccccccc}\alpha f_{ac}(\left\{\frac{y_nx_n}{\alpha }\right\},\epsilon _{ac})& & & & & & & & & & & & & & & & & & & \\ \beta & & & & & & & & & & & & & & & & & & & \end{array}\right),y_n<x_n,y_n<0.& & & & & & & & & & & & & & & & & & & \end{array}\hfill & & & & & & & & & & & & & & & & & & & \end{array}$$
(24)
In Fig. 5 some examples of phase-lockings and complex dynamics for the system (21), (22) are shown. The PRCs are chosen in the sinus form, i.e. $`f_{ij}(x,\epsilon _{ij})=\epsilon _{ij}\mathrm{sin}(2\pi x)`$. $`T_i`$ and $`\epsilon _{ij}`$ are indicated in the legends. It is obvious that the model of three interacting pacemakers requires further investigation in detail. In particular, we believe that appropriate choice of influence strength and period of external stimulation in the equations (24) will remove the system from undesirable complex behavior to the normal heart rhythm. It will be represented within forthcoming works.
## 4 Lattices of Coupled Pulse Oscillators
Let us demonstrate briefly the way how one can approximate discrete distributed media on the basis of the general model of coupled oscillators (4). Considering a cardiac pacemaker on the microscopic level, one can represent it as a large group of cells, which generate the heart rhythm and synchronize their action potentials to initiate the cardiac contraction. Thus, instead of examining a single pacemaker we should construct a lattice of coupled pulse oscillators. In this paper we restrict ourselves to the one- (chain) and two-dimensional (lattice) cases.
Let us suppose that the autonomous pacemakers are located in sites of the two-dimensional square lattice of the size $`N\times M`$. An element of the lattice with coordinates $`(i,j)`$ is designated as $`A^{i,j}`$, where $`i=1,\mathrm{},N`$ and $`j=1,\mathrm{},M`$. We restrict consideration to the homogeneous medium and accept a number of limitations on anisotropy. This means that the pacemakers of the lattice are identical, i.e. they have the same cycle length $`T_{i,j}T`$, $`i=1,\mathrm{},N,j=1,\mathrm{},M`$ (however, in fact, cells from the periphery of the sinus pacemaker show the shortest cycle length, although the centre acts as the leading pacemaker site). This limitation decreases a quantity of parameters of the system and hence simplifies the study of the model. Now we should define coupling between elements. In works devoted to the coupled map lattices (CML) two main types of coupling are usually assumed: the nearest neighbors and global couplings (see, e.g., \[Kaneko, 1989\]). Since in the previous sections we supposed that pacemakers interact with each other, this time as an example we consider lattices with the coupling to nearest neighbors: first, a two-dimensional lattice, and then a chain of coupled pulse oscillators.
Within the square lattice each pacemaker $`A^{i,j}`$ interacts with four of the neighbors according to the schematic picture in Fig. 6a. Taking into account the limitation of homogeneous medium, one can assert that all couplings of the lattice are identical, i.e. each two adjacent elements interact with each other by a general law defined by the identical PRCs $`f(x,\epsilon )`$. Moreover, suppose that the coupling between a pair of elements is isotropic in such a sense that $`\epsilon _{(i,j)(i^{},j^{})}=\epsilon _{(i^{},j^{})(i,j)}`$, and is equal to one of the values $`\epsilon _1`$ or $`\epsilon _2`$ depending on the relative position of elements. This means that there is an anisotropy of the influence strength in vertical and horizontal directions. In other words, if two pacemakers are the neighbors in the vertical direction, they interact by $`f(x,\epsilon _1)`$, and if they are horizontal neighbors, then they are coupled by $`f(x,\epsilon _2)`$.
Note that all these restrictions are made for the simplification of the analytical form of the resulting model. In general, it is possible to write the expressions for a two-dimensional lattice of coupled pulse oscillators without any limitations. It is a subject of a standalone investigation and lies beyond the framework of the given study.
Let us write equations that determine the iterative dynamics of the expected firings of the pacemakers $`\{\widehat{a}^{i,j}\}_{_{j=1,\mathrm{},M}^{i=1,\mathrm{},N}}`$ on the basis of the approach represented in Sec.2.2. To obtain the $`(n+1)`$-th value of the individual element $`\widehat{a}^{i,j}`$, it is necessary to analyze all elements of the lattice since they are coupled with each other by means of a local coupling. In other words, the considered element cannot be affected by others at the $`n`$-th time step because the latter is suppressed by the influence of other elements and remains at the expected state until the $`(n+1)`$-th step. Then dynamics of the model can be described by the following expression:
$$\begin{array}{cc}\widehat{a}_{n+1}^{i,j}=\widehat{a}_n^{i,j}+\hfill & \{\begin{array}{cc}T,\hfill & \text{if}\widehat{a}_n^{i,j}=\widehat{a}^{\mathrm{min}},\hfill \\ f(\phi _n^{(i,j)(0,+1)},\epsilon _1)T,\hfill & \text{if}\widehat{a}_n^{i,j+1}=\widehat{a}^{\mathrm{min}},\hfill \\ f(\phi _n^{(i,j)(0,1)},\epsilon _1)T,\hfill & \text{if}\widehat{a}_n^{i,j1}=\widehat{a}^{\mathrm{min}},\hfill \\ f(\phi _n^{(i,j)(+1,0)},\epsilon _2)T,\hfill & \text{if}\widehat{a}_n^{i+1,j}=\widehat{a}^{\mathrm{min}},\hfill \\ f(\phi _n^{(i,j)(1,0)},\epsilon _2)T,\hfill & \text{if}\widehat{a}_n^{i1,j}=\widehat{a}^{\mathrm{min}},\hfill \\ 0,\hfill & \text{otherwise}\hfill \end{array}\widehat{a}^{\mathrm{min}}=\mathrm{min}\{\widehat{a}_n^{i,j}\}_{_{j=1\mathrm{}M}^{i=1\mathrm{}N}}\hfill \end{array}$$
(25)
where the phases are
$$\begin{array}{c}\phi _n^{(i,j)(0,+1)}=\left\{\frac{\widehat{a}_n^{i,j+1}\widehat{a}_n^{i,j}}{T}\right\}\hfill \\ \\ \phi _n^{(i,j)(0,1)}=\left\{\frac{\widehat{a}_n^{i,j1}\widehat{a}_n^{i,j}}{T}\right\}\hfill \\ \\ \phi _n^{(i,j)(+1,0)}=\left\{\frac{\widehat{a}_n^{i+1,j}\widehat{a}_n^{i,j}}{T}\right\}\hfill \\ \\ \phi _n^{(i,j)(1,0)}=\left\{\frac{\widehat{a}_n^{i1,j}\widehat{a}_n^{i,j}}{T}\right\}\hfill \end{array}$$
The constructed model demands detailed investigation on the basis of the approach developed for the CML (see, e.g., \[Kaneko, 1989\]). It will be represented in the succeeding works.
As the second example let us consider a chain of the identical pulse oscillators coupled by the nearest neighbor principle. We restrict ourselves to a homogeneous case with anisotropy of the right and left direction in the influence strength between the nearest neighbors. The schematic picture of the chain is shown in Fig. 6b. Similarly to the above consideration one can get:
$$\begin{array}{cc}\widehat{a}_{n+1}^i=\widehat{a}_n^i+\hfill & \{\begin{array}{cc}T,\hfill & \text{if}\widehat{a}_n^i=\widehat{a}^{\mathrm{min}},\hfill \\ f(\phi _n^{i,+1},\epsilon _2)T,\hfill & \text{if}\widehat{a}_n^{i+1}=\widehat{a}^{\mathrm{min}},\hfill \\ f(\phi _n^{i,1},\epsilon _1)T,\hfill & \text{if}\widehat{a}_n^{i1}=\widehat{a}^{\mathrm{min}},\hfill \\ 0,\hfill & \text{otherwise}\hfill \end{array}\widehat{a}^{\mathrm{min}}=\mathrm{min}\{\widehat{a}_n^i\}_{i=1\mathrm{}N}\hfill \end{array}$$
(26)
where the phases are the following
$$\begin{array}{c}\phi _n^{i,+1}=\left\{\frac{\widehat{a}_n^{i+1}\widehat{a}_n^i}{T}\right\}\phi _n^{i,1}=\left\{\frac{\widehat{a}_n^{i1}\widehat{a}_n^i}{T}\right\}.\hfill \end{array}$$
If $`\epsilon _1=0`$ (or $`\epsilon _2=0`$), then Eqs (26) define the so-called open-flow model \[Willeboordse & Kaneko, 1994\].
Because in this work we present a general approach of developing models without the detailed analysis of their behavior, the type of boundary conditions for both lattices has not been indicated. Hence, to investigate such systems analytically or numerically, one should set the boundary conditions along with the PRCs $`f(x,\epsilon )`$. Usually the boundary conditions are chosen as periodic, i.e. $`\widehat{a}^{i,j+M}\widehat{a}^{i,j};\widehat{a}^{i+N,j}\widehat{a}^{i,j}`$ for the two-dimensional lattice and $`\widehat{a}^{i+N}\widehat{a}^i`$ for the one-dimensional one. For the open-flow model a condition of the fixed left boundary, $`\widehat{a}_n^1const`$, is frequently accepted.
The described models (4), (25) and (26) admit generalization to a natural inhomogeneous case by placing different intrinsic cycle lengths of the pacemakers, PRCs and influence strengths for various groups of elements. However, consideration of inhomogeneous anisotropic lattices is extremely difficult problem even for numerical analysis. The first attempts of investigating inhomogeneous lattices of coupled maps (ICML) were described in \[Vasil’ev et al., 2000; Loskutov et al., 2002; Rybalko & Loskutov, 2004\].
## 5 Summary and Limitations
In the present study we propose a quite general discrete model of active media by introducing a simple phase response curve interaction between leading centers. We have shown that the PRC can be a useful “tool” for representation of the interaction between pacemakers in cardiac tissue both on a large and small scales. This PRC based model together with demonstrating complex (chaotic) behavior, can describe the entrainment and synchronization phenomena of interacting pulse oscillators. It can also aid to understand their response to an external stimulus with variable intensity and duration (see Fig. 4b), as previously observed in experimental studies \[Jalife et al., 1976; Jalife et al., 1980\].
Starting with consideration of two interacting pulse oscillators and introducing new concepts of expected values, we have extrapolated our PRC based approach to investigate the mutual influence among an arbitrary large ensemble of pacemakers. The specific cases of the proposed model show that it can be very useful for investigating the dynamical interaction of cardiac nodes.
The last part of our study suggests that the derived general model can be easily applied to construct one– and two–dimensional lattices of active elements interacting by the nearest neighbors type. Extension of the model to a three–dimensional case is straightforward.
Finally, some limitations of our approach should be mentioned. First, the proposed model is not complete, there is no a time delay in pulse propagation among pacemakers, which can be very important for describing cardiac arrhythmias. Second, we represented cardiac tissue as a discrete one and used iterative approach to investigate its behavior. However, a large amount of realistic examples of active media is treated as continuous. Nevertheless, cardiac tissue is not a continuum, but is built up by discrete cardiomyocytes (or nodes with approximate dimensions 0.15 mm $`\times `$ 0.02 mm $`\times `$ 0.01 mm) \[Kuramoto, 1984\].
Third and most important, to analyze the essential features governing dynamics of network of active elements, we have not included many important properties of the real conductive cells. These include the relaxation after stimulating, the prolonged (non-peak) form of pulses profiles, realistic topological structure, etc. Further investigations are required to incorporate these features to the general combined model.
## 6 Acknowledgements
This paper was partially supported by INTAS fellowship No 03-55-1920, granted to Ekaterina Zhuchkova. Also, we would like to thank Prof.Alexander Loskutov for critical reading of the manuscript.
References
Antzelevitch, C., Jalife, J. & Moe, G. K. “Electrotonic modulation of pacemaker activity - further biological and mathematical observations on the behavior of modulated parasystole,” Circulation 66(6), 1225–1232.
Bub, G. & Glass, L. “Bifurcations in a continuous circle map: A theory for chaotic cardiac arrhythmia,” Int. J. Bifurcation and Chaos 5(2), 359–371.
Clayton, R. H. “Computational models of normal and abnormal action potential propagation in cardiac tissue: linking experimental and clinical cardiology,” Physiol. Meas. 22, R15–R34.
Clayton, R. H., Zhuchkova, E. A. & Panfilov, A. V. “Phase singularities and filaments: Simplifying complexity in computational models of ventricular fibrillation,” Prog. Biophys. Mol. Biol. 90, 378–398.
Courtemanche, M., Glass, L., Belair, J., Scagliotti, D. & Gordon, D. “A circle map in a human heart,” Physica D 49, 299–310.
Glass, L., Goldberger, A. L. & Belair, J. “Dynamics of pure parasystole,” Am. J. Physiol. 251(4), H841–H847.
Glass, L., Goldberger, A. L., Courtemanche, M. & Shrier, A. “Nonlinear dynamics, chaos and complex cardiac arrhythmias,” Proc. R. Soc. London Ser. A-Math. Phys. Eng. Sci. 413(1844), 9–26.
Glass, L. & Zeng, W. Z. “Complex bifurcations and chaos in simple theoretical models of cardiac oscillations,” Ann. N.Y. Acad. Sci. 591, 316–327.
Glass, L., Nagai, Yo., Hall, K., Talajie, M. & Nattel, S. “Predicting the entrainment of reentrant cardiac waves using phase resetting curves,” Phys. Rev. E 65, 021908-1–021908-10.
Goldberger, A. L. “Nonlinear dynamics, fractals and chaos: Applications to cardiac electrophysiology,” Ann. Biomed. Eng. 18(2), 195–198.
Guevara, M. R. & Shrier, A. “Phase resetting in a model of cardiac Purkinje fiber,” Biophys. J. 52(2), 165–175.
Ikeda, N. “Model of bidirectional interaction between myocardial pacemakers based on the phase response curve,” Biol. Cybern. 43(3), 157–167.
Ikeda, N., DeLand, E., Miyahara, H., Takeuchi, A., Yamamoto, H. & and Sato, T. “A personal computer-based arrhythmia generator based on mathematical models of cardiac arrhythmia,” J. Electrocardiol. 21(Suppl).
Ikeda, N., Takeuchi, A., Hamada, A., Goto, H., Mamorita, N. & Takayanagi, K. “Model of bidirectional modulated parasystole as a mechanism for cyclic bursts of ventricular premature contractions,” Biol. Cybern. 91(1), 37–47.
Jalife, J. & Moe, G. K. “Effects of electronic potentials on pacemaker activity of canine Purkinje fibers in relation to parasystole,” Circ. Res. 39(6), 801–808.
Jalife, J., Hamilton, A. J., Lamanna, V. R. & and Moe, G. K. “Effects of current flow on pacemaker activity of the isolated kitten sinoatrial node,” Am. J. Physiol. 238(3), H307–H316.
Kaneko, K. “Spatiotemporal chaos in one-dimensional and two-dimensional coupled map lattices,” Physica D 37(1-3), 60–82.
Kaneko, K. “Clustering, coding, switching, hierarchical ordering, and control in a network of chaotic element,” Physica D 41(2), 137–172.
Kuramoto, Y. Chemical Oscillations, Waves, and Turbulence (Springer-Verlag, Berlin).
Kuramoto, Y. “Scaling behavior of turbulent oscillators with nonlocal interaction,” Prog. Theor. Phys. 94(3), 321–330.
Loskutov, A., Prokhorov, A. K. & Rybalko, S. D. “Dynamics of inhomogeneous chains of coupled quadratic maps,” Theor. Math. Phys. 132(1), 983–999.
Loskutov, A., Rybalko, S. & Zhuchkova, E. “Model of cardiac tissue as a conductive system with interacting pacemakers and refractory time,” Int. J. Bifurcation and Chaos 14(7), 2457–2466.
Reiner, V. S. & Antzelevich, C. “Phase resetting and annihilation in a mathematical model of sinus node,” Am. J. Physiol. 249, H1143–H1153.
Rybalko, S. & Loskutov, A. “Dynamics of inhomogeneous one-dimensional coupled map lattices,” http://arxiv.org/abs/nlin.CD/0409014.
Sano, T., Sawanobori, T. & Adaniya, H. “Mechanism of rhythm determination among pacemaker cells of the mammalian sinus node,” Am. J. Physiol. 235, H379–H384.
Shibata, T. & Kaneko, K. “Collective chaos,” Phys. Rev. Lett. 81(19), 4116–4119.
Vasil’ev, K. A., Loskutov, A., Rybalko, S. D. & Udin, D. N. “Model of a spatially inhomogeneous one-dimensional active medium,” Theor. Math. Phys. 124(3), 1286–1297.
Wiener, N. & Rosenblueth, A. “Conduction of impulses in cardiac muscle,” Arch. Inst. Cardiol. Mex. 16, 205–265.
Willeboordse, F. H. & Kaneko, K. “Bifurcations and spatial chaos in an open flow model,” Phys. Rev. Lett. 73(4), 533–536.
Winfree, A. T. The Geometry of Biological Time (Springer-Verlag, New York), 2nd ed. |
warning/0506/cond-mat0506006.html | ar5iv | text | # From Electrons to Finite Elements: A Concurrent Multiscale Approach for Metals
## Abstract
We present a multiscale modeling approach that concurrently couples quantum mechanical, classical atomistic and continuum mechanics simulations in a unified fashion for metals. This approach is particular useful for systems where chemical interactions in a small region can affect the macroscopic properties of a material. We discuss how the coupling across different scales can be accomplished efficiently, and we apply the method to multiscale simulations of an edge dislocation in aluminum in the absence and presence of H impurities.
Some of the most fascinating problems in all fields of science involve multiple spatial and/or temporal scales: processes that occur at a certain scale govern the behavior of the system across several (usually larger) scales. In the context of materials science, the ultimate microscopic constituents of materials are ions and valence electrons; interactions among them at the atomic level determine the behavior of the material at the macroscopic scale, the latter being the scale of interest for technological applications. Conceptually, two categories of multiscale simulations can be envisioned, sequential, consisting of passing information across scales, and concurrent, consisting of seamless coupling of scales review . The majority of multiscale simulations that are currently in use are sequential ones, which are effective in systems where the different scales are weakly coupled. For systems whose behavior at each scale depends strongly on what happens at the other scales, concurrent approaches are usually required. In contrast to sequential approaches, concurrent simulations are still relatively new and only a few models have been developed to date review ; qc ; qm ; maad ; noam ; lid .
A successful concurrent multiscale method is the Quasicontinuum (QC) method originally proposed by Tadmor et al. qc . The idea underlying this method is that atomistic processes of interest often occur in very small spatial domains while the vast majority of atoms in the material behave according to well-established continuum theories. To exploit this fact, the QC method retains atomic resolution only where necessary and grades out to a continuum finite element description elsewhere. The original formulation of QC was limited to classical potentials for describing interactions between atoms. Since many materials properties depend explicitly on the behavior of electrons, such as bond breaking/forming at crack tips or defect cores, chemical reactions with impurities, and surface reactions and reconstructions, it is desirable to incorporate appropriate quantum mechanical descriptions into the QC formalism. In this Letter, we extend the original QC approach so that it can be directly coupled with quantum mechanical calculations based on density functional theory (DFT) for metallic systems. We refer to the new approach as QCDFT.
The goal of the QC method is to model an atomistic system without explicitly treating every atom in the problem qc ; qc1 . This is achieved by replacing the full set of $`N`$ atoms with a small subset of $`N_r`$ “representative atoms” or repatoms ($`N_rN`$) that approximate the total energy through appropriate weighting. The energies of individual repatoms are computed in two different ways depending on the deformation in their immediate vicinity. Atoms experiencing large deformation gradients on an atomic-scale are computed in the same way as in a standard fully-atomistic method. In QC these atoms are called nonlocal atoms to reflect the fact that their energy depends on the positions of their neighbors in addition to their own position. In contrast, the energies of atoms experiencing a smooth deformation field on the atomic scale are computed based on the deformation gradient in their vicinity as befitting a continuum model. These atoms are called local atoms because their energy is based only on the deformation gradient at the point where it is computed. The total energy $`E_{\mathrm{tot}}`$ (which for a classical system can be written as $`E_{\mathrm{tot}}=_{i=1}^NE_i`$, with $`E_i`$ the energy of atom $`i`$) is approximated as
$$E_{\mathrm{tot}}^{\mathrm{QC}}=\underset{i=1}{\overset{N^{\mathrm{nl}}}{}}E_i(\{𝐪\})+\underset{j=1}{\overset{N^{\mathrm{loc}}}{}}n_jE_j^{\mathrm{loc}}(\{𝐅\}).$$
(1)
The total energy has been divided into two parts: an atomistic region of $`N^{\mathrm{nl}}`$ nonlocal atoms and a continuum region of $`N^{\mathrm{loc}}`$ local atoms ($`N^{\mathrm{nl}}+N^{\mathrm{loc}}=N^r`$). The calculation in the atomistic region is identical to that in fully atomistic methods with the energy of the atom depending on the coordinates $`\{𝐪\}`$ of the surrounding repatoms. However, in the coarse-grained continuum region each repatom can represent a large region of $`n_i`$ atoms on the atomic scale. Rather than depending on the positions of neighboring atoms, the energy of a local repatom depends on the deformation gradients $`\{𝐅\}`$ characterizing the finite strain around its position. The basic assumption employed is the Cauchy-Born rule which relates the continuum deformation at a point to the motion of the atoms in the underlying lattice represented by this point. To obtain the necessary deformation gradients, a finite element mesh is defined with the representative atoms as its nodes. It is important to note that the calculations of $`E_j^{\mathrm{loc}}(\{𝐅\})`$ in the continuum regions make use of the same interatomic potential used in the nonlocal atomistic region. This makes the passage from the atomistic to continuum regions seamless since the same material description is used in both. This seamless description enables the model to adapt automatically to changing circumstances, for example the nucleation of new defects or the migration of existing defects. The adaptability of QC is one of its main strengths, which is missing in many other multiscale methods. A consequence of the partitioning into local and nonlocal regions and the existence of a well-defined total energy for the entire system is the presence of non-physical ghost forces at the interface. These can be eliminated by self-consistent application of dead load corrections qc1 .
The original QC formulation assumes that the total energy can be written as a sum over individual atom energies. This condition is not satisfied by quantum mechanical models. To address this limitation, in the present QCDFT approach the material of interest is partitioned into three distinct types of domains (see Fig. 1): (1) a nonlocal quantum mechanical DFT region (region I); (2) a nonlocal classical region where Embedded-Atom Method (EAM) EAM potentials are used (region II); and (3) a local region that employs the same EAM potentials as region II (region III). The total energy of the QCDFT system is then
$$E_{\mathrm{tot}}^{\mathrm{QCDFT}}=E[\mathrm{I}+\mathrm{II}]+\underset{j=1}{\overset{N^{\mathrm{loc}}}{}}n_jE_j^{\mathrm{loc}}(\{𝐅\}),$$
(2)
where $`E[\mathrm{I}+\mathrm{II}]`$ is the total energy of regions I and II together (the assumption here is that region I is embedded within region II). The coupling between regions II and III is achieved seamlessly via the QC formulation, while the coupling between regions I and II is accomplished by a scheme recently proposed by Choly et al. choly . Based on this coupling strategy, $`E[\mathrm{I}+\mathrm{II}]`$ can be written as
$$E[\mathrm{I}+\mathrm{II}]=E_{\mathrm{DFT}}[\mathrm{I}]+E_{\mathrm{EAM}}[\mathrm{II}]+E^{\mathrm{int}}[\mathrm{I},\mathrm{II}],$$
(3)
where $`E_{\mathrm{DFT}}`$\[I\] is the energy of region I in the absence of region II computed using the DFT model, $`E_{\mathrm{EAM}}`$\[II\] is the energy of region II in the absence of region I computed using the EAM model, and $`E^{\mathrm{int}}`$\[I,II\] represents a formal interaction energy added to give the correct total energy. The interaction energy between the two subsystems can be rewritten as:
$`E^{\mathrm{int}}[\mathrm{I},\mathrm{II}]`$ $``$ $`E[\mathrm{I}+\mathrm{II}]E[\mathrm{I}]E[\mathrm{II}],`$
$`=`$ $`E_{\mathrm{EAM}}[\mathrm{I}+\mathrm{II}]E_{\mathrm{EAM}}[\mathrm{I}]E_{\mathrm{EAM}}[\mathrm{II}].`$
The first equation serves as a general definition of the interaction energy whereas the second equation represents one particular implementation of $`E^{\mathrm{int}}`$, which is used in this work. Eq. From Electrons to Finite Elements: A Concurrent Multiscale Approach for Metals is not contradictory to Eq. 3 because EAM has its root in DFT and the EAM energy can be viewed as an approximation to the DFT energy. Different combinations of quantum mechanical and classical atomistic methods other than DFT/EAM may also be implemented choly . The great advantage of the present implementation is its simplicity. It demands nothing beyond what is required for a DFT calculation and an EAM QC calculation. Furthermore, by substituting Eq. From Electrons to Finite Elements: A Concurrent Multiscale Approach for Metals into Eq. 3, we arrive at
$$E[\mathrm{I}+\mathrm{II}]=E_{\mathrm{DFT}}[\mathrm{I}]E_{\mathrm{EAM}}[\mathrm{I}]+E_{\mathrm{EAM}}[\mathrm{I}+\mathrm{II}].$$
(5)
The forces on the EAM atoms in region II are then
$$𝐅_i^{\mathrm{II}}=\frac{E_{\mathrm{tot}}^{\mathrm{QCDFT}}}{𝐪_i^{\mathrm{II}}}=\frac{E_{\mathrm{EAM}}[\mathrm{I}+\mathrm{II}]}{𝐪_i^{\mathrm{II}}}+\frac{_{j=1}^{N^{\mathrm{loc}}}n_jE_j^{\mathrm{loc}}(\{𝐅\})}{𝐪_i^{\mathrm{II}}},$$
(6)
where $`𝐪_i^{\mathrm{II}}`$ are the Cartesian coordinates of atom $`i`$ in region II. It is clear from this equation that the forces on the atoms in region II are identical to those that would be obtained from a fully-classical QC calculation. The same applies to the region III atoms, that is, as far as forces are concerned, regions II and III behave as though the entire model were classical. This is a very desirable property in terms of achieving a seamless coupling between region I and the rest of the model. At the same time, the forces on the DFT atoms in region I will have contributions from both DFT atoms and the nearby EAM atoms in region II. The error in forces on the DFT atoms due to the coupling is thus given by the difference between calculated forces with DFT and EAM on these atoms. To minimize this error, we propose to use a class of interatomic potentials which are generated by matching the forces obtained from the EAM method to those from DFT calculations E-A ; Ta . Another important practical advantage of the present QCDFT method is that, if region I contains many different atomic species while region II contains only one atom type, there is no need to develop reliable EAM potentials that can describe each species and their interactions. This is because if the various species of atoms are well within region I, then the energy contributions of these atoms are canceled out in the total energy calculation (the last two terms in Eq. 5). This advantage renders the method particularly useful in dealing with impurities, which is an exceedingly difficult task for empirical potential simulations.
The equilibrium structure of the system is obtained by minimizing the total energy in Eq. 2 with respect to all degrees of freedom. Because the time required to evaluate $`E_{\mathrm{DFT}}`$\[I\] is considerably more than that required for computation of the other EAM terms in $`E_{\mathrm{tot}}^{\mathrm{QCDFT}}`$, an alternate relaxation scheme turns out to be rather efficient. The total system can be relaxed by using the conjugate gradient approach on the DFT atoms alone, while fully relaxing the EAM atoms in region II and the displacement field in region III at each step. Similar to Choly et al. choly , an auxiliary energy function can be defined as
$$E^{}[\{𝐪^\mathrm{I}\}]\underset{\{𝐪^{\mathrm{II}}\},\{𝐪^{\mathrm{III}}\}}{\mathrm{min}}E_{\mathrm{tot}}^{\mathrm{QCDFT}}[\{𝐪\}],$$
(7)
which allows for the following relaxation scheme: (i) Minimize $`E_{\mathrm{tot}}^{\mathrm{QCDFT}}`$ with respect to the atoms in regions II ($`\{𝐪^{\mathrm{II}}\}`$) and the atoms in region III ($`\{𝐪^{\mathrm{III}}\}`$), while holding the atoms in region I fixed; (ii) Calculate $`E_{\mathrm{tot}}^{\mathrm{QCDFT}}[\{𝐪\}]`$, and the forces on the region I atoms; (iii) Perform a step of conjugate gradient minimization of $`E^{}`$; (iv) Repeat until the system is relaxed. In this manner, the number of DFT calculations performed is greatly reduced, albeit at the expense of more EAM and local QC calculations. A number of tests have shown that the total number of DFT energy calculations for the relaxation of an entire system is about the same as that required for DFT relaxation of region I alone. Further computational speed-up can be achieved for the DFT calculations by using converged electronic charge density and wave functions from the previous step, so that the charge (potential) self-consistency can be reached faster for the next DFT calculation because the atomic relaxation is usually very small between two consecutive DFT moves.
In the remainder of the paper, we apply the present QCDFT approach to study the core structure of an edge dislocation in Al in the absence and presence of H impurities. We chose this system as an example because results from both experiments and simulations are available for comparison. The QCDFT model for an edge dislocation with a Burgers vector $`\frac{a}{2}`$ ($`a=3.97`$ Å) is presented in Fig. 1. Convergence tests on the size of region I indicate that a DFT box of 30 Å $`\times `$ 9 Å $`\times `$ 4.86 Å (84 DFT atoms) is sufficient to capture the dislocation splitting behavior accurately; hence the following calculations are all based on this DFT box. A force-matching potential for Al E-A was used for EAM calculations. The DFT calculations were performed by using the plane-wave pseudopotential VASP code vasp for a cluster with 8 Å vacuum in both the x and y directions. The energy cutoff for pure Al and Al+H is 129 eV and 200 eV, respectively. We find that 10 $`k`$ points along the one-dimensional Brillouin zone are adequate for good convergence. Fig. 2 presents the simulation results for both a standard EAM-based QC calculation and the QCDFT method, showing the dissociation of the edge dislocation into two equivalent 60 Shockley partials. The splitting distance (obtained from an analysis of the displacement jump across the slip plane) in the standard QC calculation is 15.4 Å, whereas the splitting distance obtained with QCDFT method is 5.6 Å, a value very close to the experimentally observed value of 5.5 Å mills . This result demonstrates that even for a simple metal like Al which should be the best candidate for use of an EAM potential, a quantum mechanical calculation is necessary to obtain correct results.
The most important advantage of QCDFT approach, however, is that it allows the study of impurity effects on mechanical response, an impossible task for simpler empirical potentials. Fig. 3 shows the effect of adding one column of H impurities at the dislocation core. The presence of the H atoms results in a spreading of the core (the splitting distance is now increased to about 13 Å). This finding is consistent with the fact, confirmed by earlier DFT calculations lu\_h , that H can lower the stacking fault energy. The fact that the dislocation becomes wider may explain the H-enhanced dislocation mobility that is believed to lead to H embrittlement phenomena via the enhanced local plasticity theory. A similar core structure is also found for two columns of H atoms placed at the dislocation core.
In order to understand the underlying origin of the H-enhanced dislocation mobility, we calculate the electron density distribution at the dislocation core in the absence and presence of H impurities, as shown in Fig. 4 . In the absence of the H impurity, the electron bonding is stronger and with a distinct covalent character. The bonding is more directional above the slip plane, and it becomes more spherical below the slip plane where there are two extra atomic planes, corresponding to the two partial dislocations. In the presence of H atoms, charge accumulation develops at these H atoms as the H impurities attract the valence electrons from the Al atoms and become negatively charged. The covalent bonding across the slip plane between Al atoms is disrupted by the H atoms, and at the same time, ionic bonding between the oppositely charged H and Al ions is developed. The fact that the directional covalent bonds are replaced by more homogeneous ionic bonds near the core leads to the wider dislocation core seen in Fig. 2.
In summary, we have introduced a multiscale modeling approach which concurrently couples quantum mechanical, classical atomistic and continuum mechanics simulations, in a unified fashion for metals. Our QCDFT method provides a useful framework for multiscale modeling of metallic materials because it does not require the existence of localized covalent bonds for computing the coupling energy as all other multiscale methods do maad ; qm ; noam ; lid . Furthermore, this approach is completely general and versatile: it can be applied to diverse materials problems, such as dislocations, cracks, surfaces, and grain boundaries. Finally, the automatic adaption feature of the QCDFT method allows the DFT and/or EAM region to move and change in response to the current deformation state, when for example, defects are being nucleated in an otherwise perfect region. To demonstrate the unique strength of this method in dealing with impurities, we have applied it to study H-dislocation interactions in Al.
This research was partly supported by an award from Research Corporation (GL). ET and GL thank the Institute for Mathematics and its Applications (IMA) for hosting them in the fall of 2004 during which time part of this work was done. |
warning/0506/quant-ph0506261.html | ar5iv | text | # A unified approach to realize universal quantum gates in a coupled two-qubit system with fixed always-on coupling
(June 29, 2005)
## Abstract
We demonstrate that in a coupled two-qubit system any single-qubit gate can be decomposed into two conditional two-qubit gates and that any conditional two-qubit gate can be implemented by a manipulation analogous to that used for a controlled two-qubit gate. Based on this we present a unified approach to implement universal single-qubit and two-qubit gates in a coupled two-qubit system with fixed always-on coupling. This approach requires neither supplementary circuit or additional physical qubits to control the coupling nor extra hardware to adjust the energy level structure. The feasibility of this approach is demonstrated by numerical simulation of single-qubit gates and creation of two-qubit Bell states in rf-driven inductively coupled two SQUID flux qubits with realistic device parameters and constant always-on coupling.
During the past decade, a variety of physical qubits have been explored for possible implementation of quantum gates. Of those solid-state qubits based on superconducting devices have attracted much attention because of their advantages of large-scale integration and easy connection to conventional electronic circuits Leggett (2002); Clarke (2003); Blatter (2003); Nakamura et al. (1999); Vion et al. (2002); Yu et al. (2002); Han et al. (2001); Martinis et al. (2002); Friedman et al. (2000); van der Wal et al. (2000); Chiorescu et al. (2003); Pashkin et al. (2003); Yamamoto et al. (2003); Berkley et al. (2003); McDermott et al. (2005). Superconducting single-qubit gates Nakamura et al. (1999); Vion et al. (2002); Yu et al. (2002); Han et al. (2001); Martinis et al. (2002); Friedman et al. (2000); van der Wal et al. (2000); Chiorescu et al. (2003) and two-qubit gates Pashkin et al. (2003); Yamamoto et al. (2003); Berkley et al. (2003); McDermott et al. (2005) have been demonstrated recently.
However, building a practical quantum computer requires to operate a large number of multi-qubit gates simultaneously in a coupled multi-qubit system. It has been demonstrated that any type of multi-qubit gate can be decomposed into a set of universal single-qubit gates and a two-qubit gate, such as the controlled-NOT (CNOT) gate Nielsen and Chuang (2000); Barenco et al. (1995). Thus it is imperative to implement the universal single-qubit and two-qubit gates in a multi-qubit system with the minimum resource and maximum efficiency Plourde et al. (2004).
Implementing universal single-qubit gates and two-qubit gates in coupled multi-qubit systems can be achieved by turning off and on the coupling between qubits Plourde et al. (2004); Averin and Bruder (2003); Blais et al. (2003). In these schemes, supplementary circuits were required to control inter-qubit coupling. However, rapid switching of the coupling results in two serious problems. The first one is gate errors caused by population propagation between qubits. Because the computational states of the single-qubit gates are not a subset of the eigenstates of the two-qubit gates the populations of the computational states propagate from one qubit to another when the coupling is changed, resulting in additional gate errors. The second one is additional decoherence introduced by the supplementary circuits Zhou et al. (2002); Benjamin (2002). This is one of the biggest obstacles for quantum computing with solid-state qubits, particularly in coupled multi-qubit systems Clarke (2003). In addition, the use of supplementary circuits also significantly increase the complexity of fabrication and manipulation of the coupled qubits.
To circumvent these problems, a couple of alternative schemes, such as those with untunable coupling Zhou et al. (2002) and always-on interaction Benjamin and Bose (2003), have been proposed. In the first scheme, each logic qubit is encoded by extra physical qubits and coupling between the encoded qubits is constant but can be turned off and on effectively by putting the qubits in and driving them out of the interaction free subspace. In the second scheme, the coupling is always on but the transition energies of the qubits are tuned individually or collectively. These schemes can overcome the problem of undesired population propagation but still suffer from those caused by the supplementary circuits needed to move the encoded qubits and tune the transition energies. Moreover, the use of encoded qubits also requires a significant number of additional physical qubits.
In this Letter, we present a unified approach to implement universal single-qubit and two-qubit gates in a coupled two-qubit system with fixed always-on coupling. In this approach, each single-qubit gate is realized via two conditional two-qubit gates and each conditional two-qubit gate is implemented with a manipulation analogous to that used for a controlled two-qubit gate (e.g., CNOT) in the same subspace of the coupled two-qubit system without additional circuits or physical qubits. Since the computational states of the single-qubit gates are a subset of those of the two-qubit gates the gate errors due to population propagation are completely eliminated. The effectiveness of the approach is demonstrated by numerically simulating the single-qubit gates and creating the Bell states in a unit of inductively coupled two superconducting quantum interference device (SQUID) flux qubits with realistic device parameters and constant always on coupling.
Consider a basic unit consisting of two coupled qubits which we call a control qubit ($`𝒞q`$) and a target qubit ($`𝒯q`$) for convenience. An eigenstate of the coupled qubits is denoted by $`|n)=|ij`$, which can be well approximated by the product of an eigenstate of $`𝒞q`$, $`|i`$, and that of $`𝒯q`$, $`|j`$, $`|n)|ij=|i|j`$, for weak inter-qubit coupling. The computational states of the coupled qubits are $`|1)=|00`$, $`|2)=|01`$, $`|3)=|10`$, and $`|4)=|11`$, which correspond to the eigenstates for $`i,`$ $`j=0,`$ $`1`$ and $`n=1,`$ $`2,`$ $`3,`$ $`4`$, respectively. In general, the result of an operation on $`𝒞q`$ depends on the state of $`𝒯q`$ since $`𝒞q`$ is coupled to and hence influenced by $`𝒯q`$. Suppose $`U`$ is a unitary operator acting on a single qubit in the four-dimensional (4D) Hilbert space spanned by $`|n)`$. In order to perform $`U`$ on $`𝒞q`$ one wants the state of $`𝒞q`$ evolves according to $`U`$ independent of the state of $`𝒯q`$ which is left unchanged. This operation is denoted by $`U^{\left(2\right)}=U_C^{\left(1\right)}I_T^{\left(1\right)}`$, where, the 4$`\times `$4 unitary matrix $`U^{\left(2\right)}`$ and the 2$`\times `$2 unitary matrix $`U_C^{(1)}`$ are the representations of $`U`$ in the Hilbert space of the coupled qubits and in the subspace of $`𝒞q`$ while the 2$`\times `$2 unitary matrix $`I_T^{\left(1\right)}`$ is the representation of the identity operation in the subspace of $`𝒯q`$. If the matrix elements of $`U_C^{\left(1\right)}`$ are $`u_{ij}`$ ($`i,j=1,2`$) the operation can be decomposed into two operations as
$$U^{\left(2\right)}=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& u_{11}& 0& u_{12}\\ 0& 0& 1& 0\\ 0& u_{21}& 0& u_{22}\end{array}\right)\left(\begin{array}{cccc}u_{11}& 0& u_{12}& 0\\ 0& 1& 0& 0\\ u_{21}& 0& u_{22}& 0\\ 0& 0& 0& 1\end{array}\right)=U_1^{\left(2\right)}U_0^{\left(2\right)},$$
(1)
where, the 4$`\times `$4 unitary matrices $`U_0^{\left(2\right)}`$ and $`U_1^{\left(2\right)}`$ are also in the 4D Hilbert space of the coupled qubits. Since $`U_0^{\left(2\right)}`$ only involves the states $`|00`$ and $`|10`$ and $`U_1^{\left(2\right)}`$ only involves $`|01`$ and $`|11`$ they represent operations on $`𝒞q`$ when $`𝒯q`$ is in the state $`|0`$ and $`|1`$, respectively. If the state of $`𝒯q`$ is $`|0`$ ($`|1`$), $`U_0^{\left(2\right)}`$ ($`U_1^{\left(2\right)}`$) implements the desired operation $`U`$ while $`U_1^{\left(2\right)}`$ ($`U_0^{\left(2\right)}`$) does nothing. Thus $`U_0^{\left(2\right)}`$ and $`U_1^{\left(2\right)}`$ represent two conditional two-qubit gates of the coupled qubits. Eq.(1) shows that for the coupled qubits any single-qubit gate can be realized via two conditional two-qubit gates. This property is universal and is shown schematically in Fig. 1, where the single-qubit operation $`U^{\left(2\right)}`$ and its equivalent operation via two conditional two-qubit gates $`U_0^{\left(2\right)}`$ and $`U_1^{\left(2\right)}`$ are illustrated. The conditional gates are denoted by the open and closed circles for the state of $`𝒯q`$ being $`|0`$ and $`|1`$, respectively Nielsen and Chuang (2000).
The most obvious advantage of implementing single-qubit gates this way is that the conditional two-qubit gates can be realized using essentially the same set of operations used for controlled two-qubit gates. To show this, let us consider a rf-driven coupled two-qubit system with sufficiently different energy-level spacings $`\mathrm{\Delta }E_{13}`$, $`\mathrm{\Delta }E_{24}`$, $`\mathrm{\Delta }E_{12}`$, and $`\mathrm{\Delta }E_{34}`$, where $`\mathrm{\Delta }E_{nn^{}}`$ is the level spacing between the states $`|n)`$ and $`|n^{})`$. To implement, for example, a single-qubit NOT gate, we apply two $`\pi `$ pulses to the $`𝒞q`$: the first one is resonant with $`\mathrm{\Delta }E_{13}`$ and the second with $`\mathrm{\Delta }E_{24}`$. Because both pulses are largely detuned from $`\mathrm{\Delta }E_{12}`$ and $`\mathrm{\Delta }E_{34}`$, they do not induce unintended population transfer when the fields are sufficiently weak. If the state of $`𝒯q`$ is $`|0`$ $`\left(|1\right)`$, the first (second) pulse accomplishes the desired transformation. Therefore, when the two microwave pulses are applied to $`𝒞q`$, either sequentially or simultaneously, the NOT gate is accomplished and the gate time is essentially the same as that of a stand-alone $`𝒞q`$ if the both pulses are applied simultaneously.
To implement a two-qubit gate such as CNOT in coupled qubits, we apply a $`\pi `$ pulse resonant with $`\mathrm{\Delta }E_{34}`$ to $`𝒯q`$ Zhou et al. (2005, in press). In this case, the state of $`𝒯q`$ flips if and only if the state of $`𝒞q`$ is $`|1`$. Since the energy level structure for the conditional two-qubit gates and the CNOT gate could be the same the universal single-qubit gates and CNOT gate can be implemented in the same fashion without adjusting inter-qubit coupling as long as $`𝒞q`$ and $`𝒯q`$ can be addressed individually by microwave pulses, which is rather straightforward to realize with solid-state qubits in general and with rf SQUID flux qubits in particular.
For concreteness, we demonstrate how to realize the single-qubit and two-qubit gates in rf-driven two SQUID flux qubits with constant always-on coupling. The coupled flux qubits comprise two SQUIDs coupled inductively through their mutual inductance $`M`$. Each SQUID consists of a superconducting loop of inductance $`L`$ interrupted by a Josephson tunnel junction characterized by its critical current $`I_c`$ and shunt capacitance $`C`$ Danilov et al. (1983). A flux-biased SQUID with total magnetic flux $`\mathrm{\Phi }`$ enclosed in the loop is analogous to a “flux” particle of mass $`m=C\mathrm{\Phi }_0^2`$ moving in a one-dimensional potential, where $`\mathrm{\Phi }_0`$ $`=h/2e`$ is the flux quantum. For simplicity, we assume that the two SQUIDs are identical. In this case, $`C_i=C`$, $`L_i=L`$, and $`I_{ci}=I_c`$ for $`i=1,2`$. The Hamiltonian of the coupled qubits is Mooij et al. (1999); Zhou et al. (2005, in press) $`H(x_1,x_2)=H_0\left(x_1\right)+H_0\left(x_2\right)+H_{12}(x_1,x_2)`$, where $`H_0\left(x_i\right)`$ is Hamiltonian of the $`i`$th single qubit given by $`H_0\left(x_i\right)=p_i^2/2m+m\omega _{LC}^2\left(x_ix_{ei}\right)^2/2E_J\mathrm{cos}\left(2\pi x_i\right)`$ and $`H_{12}`$ is the interaction between the SQUIDs given by $`H_{12}=m\omega _{LC}^2\kappa \left(x_1x_{e1}\right)\left(x_2x_{e2}\right)/2`$. Here, $`x_i=\mathrm{\Phi }_i/\mathrm{\Phi }_0`$ is the canonical coordinate of the $`i`$th “flux” particle and $`p_i=i\mathrm{}/x_i`$ is the canonical momentum conjugate to $`x_i`$, $`x_{ei}=\mathrm{\Phi }_{ei}/\mathrm{\Phi }_0`$ is the normalized external flux bias of the $`i`$th qubit, $`E_J=m\omega _{LC}^2\beta _L/4\pi ^2`$ is the Josephson coupling energy, $`\beta _L=2\pi LI_c/\mathrm{\Phi }_0`$ is the potential shape parameter, $`\omega _{LC}=1/\sqrt{LC}`$ is the characteristic frequency of the SQUID, and $`\kappa =2M/L`$ is the coupling strength. The coupled SQUID qubits are a multi-level system. The eigenenergies $`E_n`$ and eigenstates $`|n)`$ are obtained by numerically solving the eigenvalue equation of $`H(x_1,x_2)`$ using the two-dimensional Fourier-grid Hamiltonian method Chu (1990). When $`x_{e1}`$ and $`x_{e2}0.5`$ the coupled SQUID qubits have four wells in the potential energy surface Zhou et al. (2005, in press). The four computational states are chosen to be the lowest eigenstate of each well.
To realize single-qubit and two-qubit gates in the coupled SQUID qubits, we apply microwave pulses $`x_C`$ and $`x_T`$ to $`𝒞q`$ and $`𝒯q`$, respectively. The interaction of the coupled qubits and pulses is $`V(x_1,x_2,t)=d_1\left(x_1x_{e1}\right)+d_2\left(x_2x_{e2}\right)+d_{12}`$ with $`d_1=m\omega _{LC}^2\left(x_C+\kappa x_T/2\right)`$, $`d_2=m\omega _{LC}^2\left(x_T+\kappa x_C/2\right)`$, and $`d_{12}=m\omega _{LC}^2\left(x_C^2+x_T^2+\kappa x_Cx_T\right)/2`$. The time-dependent wave functions of the coupled SQUID qubits, $`\psi (x_1,x_2,t)`$, are governed by the time-dependent Schrödinger equation $`i\mathrm{}\psi /t=\left[H(x_1,x_2)+V(x_1,x_2,t)\right]\psi `$. To solve this equation we expand $`\psi `$ in terms of the first 20 eigenstates of the coupled qubits: $`\psi =_nc_n\left(t\right)|n)`$. The expansion coefficients $`c_n(\tau )`$ are calculated by numerically solving the matrix equation $`ic_n\left(\tau \right)/\tau =_n^{}H_{nn^{}}^R\left(\tau \right)c_n^{}\left(\tau \right)`$ using the split-operator method Hermann and Fleck, Jr. (1988), where $`\tau =\omega _{LC}t`$ is the dimensionless time and $`H_{nn^{}}^R=\left[E_n\delta _{nn^{}}+\left(n\left|V\right|n^{}\right)\right]/\mathrm{}\omega _{LC}`$ is the reduced Hamiltonian matrix element. The probability of being in the state $`|n)`$ is thus $`\left|c_n(\tau )\right|^2`$. For single-qubit gates two resonant pulses, $`x_{C1}=x_{C10}\mathrm{cos}\left(\omega _{C1}t\right)`$ with $`\omega _{C1}=\mathrm{\Delta }E_{13}/\mathrm{}`$ and $`x_{C2}=x_{C20}\mathrm{cos}\left(\omega _{C2}t\right)`$ with $`\omega _{C2}=\mathrm{\Delta }E_{24}/\mathrm{}`$, are applied simultaneously to the $`𝒞q`$ so that $`x_C=x_{C1}+x_{C2}`$. For controlled two-qubit gates such as the CNOT gate one resonant pulse $`x_T=x_{T0}\mathrm{cos}\left(\omega _Tt\right)`$ with $`\omega _T=\mathrm{\Delta }E_{34}/\mathrm{}`$ is applied to the $`𝒯q`$. To minimize the possible intrinsic gate errors caused by leakage to unintended states, we select the parameters of the coupled SQUID qubits using the independent transition approach Zhou et al. (submitted). For SQUIDs with $`L=100`$ pH, $`C=40`$ fF, and $`\beta _L=1.2`$ one set of the better working parameters is $`x_{e1}=0.499`$, $`x_{e2}=0.4998`$, and $`\kappa =5\times 10^4`$, which will be used in the calculation.
Most of the single-qubit gates, such as the NOT gate and Hadamard gate, can be described by rotations of the state vector on the Bloch sphere Nielsen and Chuang (2000). Thus we demonstrate how to realize an arbitrary single-qubit rotation of an angle $`\theta `$ about an axis perpendicular to the $`z`$ axis, $`R^{\left(2\right)}(\theta )`$, in the coupled SQUID qubits. Based on Eq.(1), the rotation $`R^{\left(2\right)}(\theta )`$ on $`𝒞q`$ is accomplished via two conditional two-qubit rotations $`R_0^{\left(2\right)}(\theta )`$ and $`R_1^{\left(2\right)}\left(\theta \right)`$ as shown in FIG. 2(a). They are realized by applying two microwave pulses $`x_{C1}`$ and $`x_{C2}`$ to $`𝒞q`$ simultaneously. The amplitudes and frequencies of $`x_{C1}`$ and $`x_{C2}`$ are $`x_{C10}=5\times 10^5`$, $`\omega _{C1}=0.239\omega _{LC}=2\pi \times 19.0`$ GHz, $`x_{C20}=5.14\times 10^5`$, and $`\omega _{C2}=0.259\omega _{LC}=2\pi \times 20.6`$ GHz. In Fig. 2 (b) and (c), we plot the rotation angle $`\theta `$ and the populations of the states $`|00`$ and $`|10`$ (populations of the other states remain essentially at zero) as a function of pulse width when the initial state of the coupled qubits is $`|00`$. It is shown in Fig. 2 (b) that the rotation angle $`\theta `$ is essentially a linear function of pulse width. This indicates that the state of $`𝒞q`$ undergoes Rabi oscillations for which the phase angle $`\mathrm{\Omega }\tau `$ is a linear function of pulse width, where $`\mathrm{\Omega }`$ is the Rabi frequency. It is also shown in Fig. 2 (c) that from the initial state $`|00`$ the coupled qubits evolve into the state $`\left(|00+|10\right)/\sqrt{2}`$ after the $`\pi /2`$ pulses and into the state $`|10`$ after the $`\pi `$ pulses. We have also computed the rotation angles and populations for the coupled qubits with the initial states $`|01`$, $`|10`$, and $`|11`$ using the same pulse sequence. In each case, the state is transformed from $`|ij`$ to $`R^{\left(2\right)}(\theta )|ij=\left[R^{\left(2\right)}(\theta )|i\right]|j`$ for $`i,j=0,1`$ with $`\theta =\mathrm{\Omega }\tau `$ as expected.
Entanglement is one of the most profound characteristics of quantum systems which plays a crucial role in quantum information processing and communication Nielsen and Chuang (2000). The maximally entangled two-qubit states are referred to as the Bell states. To create the Bell states from a state $`|ij`$, a Hadamard gate on $`𝒞q`$, $`H^{(2)}`$, which is decomposed into two conditional two-qubit Hadamard gates $`H_0^{(2)}`$ and $`H_1^{(2)}`$, and a CNOT gate are commonly used \[see FIG. 3(a)\]. The two conditional two-qubit Hadamard gates are implemented by applying two $`\pi /2`$ pulses $`x_{C1}`$ and $`x_{C2}`$ to $`𝒞q`$ and the following CNOT gate is implemented by applying a $`\pi `$ pulse $`x_T`$ to $`𝒯q`$, as shown in FIG. 3(b). The amplitudes and frequencies of $`x_{C1}`$ and $`x_{C2}`$ are the same as those used in FIG. 2 and those of $`x_T`$ are $`x_{T0}=5\times 10^5`$ and $`\omega _T=0.0592\omega _{LC}=2\pi \times 4.7`$ GHz. In FIG. 3(c), we plot the population evolution of the computational states when the initial state is $`|00`$. It is shown clearly that the coupled qubits evolve first into a product state $`\left(|00+|10\right)/\sqrt{2}`$ from the initial state $`|00`$ after the $`\pi /2`$ pulses and then into a Bell state $`\left(|00+|11\right)/\sqrt{2}`$ after the subsequent $`\pi `$ pulse. We have also calculated the population evolution for the coupled qubits being initially in $`|01`$, $`|10`$, and $`|11`$, respectively. The final state in each case is also one of the expected Bell states.
In summary, we showed that in a coupled two-qubit system any single-qubit gate can be realized via two conditional two-qubit gates and that any conditional two-qubit gate can be implemented with a manipulation analogous to that used for a controlled two-qubit gate. Based on this universal property of single-qubit gates we present a general approach to implement the universal single-qubit and two-qubit gates in the same coupled two-qubit system with fixed always-on coupling. This approach is demonstrated by using a unit of two SQUID flux qubits with realistic device parameters and constant always-on coupling. Compared to other methods our approach has the following characteristics and advantages: (1) The approach is universal as long as each qubit can be locally addressed; (2) No additional decoherence from the hardware added to control inter-qubit coupling; (3) Gate error induced by the population propagation from one qubit to another is completely eliminated; (4) The architecture for both hardware (circuits) and software (pulse sequence) is much simplified. This approach can be readily extended to multi-qubit systems and other types of solid-state qubit systems in which each qubit can be individually addressed. Therefore, it is very promising for realizing universal gates with minimum resource and complexity and maximum efficiency.
This work is supported in part by the NSF (DMR-0325551) and by AFOSR, NSA, and ARDA through DURINT grant (F49620-01-1-0439). |
warning/0506/cond-mat0506219.html | ar5iv | text | # Correlation effects in the ground state of trapped atomic Bose gases
## I Introduction
The many-body physics in trapped Bose gases has drawn intense interest since the experimental realization of Bose-Einstein condensation (BEC) in ultracold, dilute alkali atoms Anderson et al. (1995). The systems are “clean” and highly controllable experimentally. The dominant interactions are simple and well-understood, and the strength of the interatomic interactions can be readily tuned by means of Feshbach resonances Cornish et al. (2000). With the recent realization of degenerate Fermi gases DeMarco and Jin (1999); Schreck et al. (2001); Truscott et al. (2001), these ultracold systems provide an ideal “laboratory” for studying many-body physics.
In the weakly-interacting regime, mean-field theories work quite well, for instance the Gross-Pitaevskii (GP) equation Gross (1961, 1963); Pitaevskii (1961) for boson ground states. Much work has been done to study the ground state of the Bose atomic gases beyond mean field. For example, a modified GP equation was proposed Braaten and Nieto (1997) by inclusion of one-loop quantum corrections and the use of local-density approximation. Esry Esry (1997) developed a Hartree-Fock theory as a means of including the correlation effects in the BEC many-body calculations. Mazzanti and co-workers Mazzanti et al. (2003) applied a correlated basis theory Fantoni and Fabrocini (1998) to study the detailed structure of dilute hard- and soft-sphere Bose gases. A comparative study for the modified GP and correlated basis approaches is presented in LABEL:Fabrocini1999. Recently, McKinney and co-workers McKinney et al. (2004) used a many-body dimensional perturbation theory to compute the ground-state energy and breathing-mode frequency of spherically trapped gases at different interaction strengths.
Semianalytic methods are versatile and generally very easy to extend to realistic systems with large number of particles. However, they are approximate and each has its own limitations, especially in the strongly-interacting regime. Computational methods such as quantum Monte Carlo (QMC) provide a useful, complementary alternative. A variety of such calculations have been carried out for atomic boson systems, including variational Monte Carlo DuBois and Glyde (2001) and the exact diffusion Monte Carlo (DMC) Ceperley and Alder (1980); Umrigar et al. (1993) studies on both the homogenous Giorgini et al. (1999) and trapped gases Blume and Greene (2001); Sakhel et al. (2002); DuBois and Glyde (2003).
We have recently developed an auxiliary-field quantum Monte Carlo (AF QMC) method Purwanto and Zhang (2004) for the ground state of many-boson systems. While the standard DMC works in real space with particle configurations, our method works in the second-quantized formalism, which automatically accounts for particle permutation statistics. The calculation can be carried out in any single-particle basis. Conceptually, the method provides a way to systematically improve upon mean field while retaining its basic machinery, capturing correlation effects with a stochastic, coherent ensemble of independent-particle solutions. Various observables and correlation functions can be calculated relatively straightforwardly.
The initial motivation of this study was to use the AF QMC method to quantify, by direct comparison with GP, the effects of interactions in trapped Bose gases, and to provide additional precise numerical data where they were not available. (Although the method is not exact for bosons with repulsive interactions, the systematic errors are very small in the parameter region of interest, as we show below.) In particular, we were interested in the behavior of the system as a function of the interaction strength, which, unlike in typical condensed matter systems, can be probed directly in experiments. We found that GP yielded significant errors in the energetics in the Feshbach resonance regime, which resulted in a qualitatively incorrect behavior of the kinetic energy in GP as a function of the scattering length. To study the origin of these errors, we carried out additional calculations using first-order Bogoliubov results under a local-density approximation (LDA). The purpose of this paper is thus to present our QMC data, and discuss the behavior of the GP, modified GP, and Bogoliubov-LDA methods as benchmarked by QMC.
The rest of the paper is organized as follows. In Sec. II, we describe the many-body Hamiltonian. Our QMC method is summarized in Sec. III, as are the procedures of our GP and Bogoliubov-LDA calculations. Results from QMC, GP, and first-order Bogoliubov-LDA methods are presented in Sec. IV, where we study the energetics of the gas in three dimensions as a function of the number of particle $`N`$ and the $`s`$-wave scattering length $`a_s`$, and examine the density profile and momentum distribution. Our study extends to the strongly-interacting regime achieveable by Feshbach resonances. In Sec. V, we discuss the implications of our comparisons between GP, modified GP, Bogoliubov-LDA and QMC. In addition, we also discuss the influence of the details of the two-body potential. Concluding remarks are given in Sec. VI. Finally, in the appendix, we describe additional details on our Bogoliubov and QMC calculations, including benchmark results on the systematic errors in our QMC.
## II Modified Bose-Hubbard model
We consider $`N`$ Bose particles in a three-dimensional cube of side length $`2r_b`$, under the periodic boundary condition. Similar to our earlier work Purwanto and Zhang (2004), we use the Bose-Hubbard model as the discrete representation of the many-body Hamiltonian on a real-space lattice:
$$\begin{array}{cc}\hfill \widehat{H}& =\frac{\mathrm{}^2}{2m}\underset{𝐤}{}k^2\widehat{\phi }^{}(𝐤)\widehat{\phi }(𝐤)\hfill \\ & +\frac{1}{2}m\omega _0^2d^3𝐫r^2\widehat{\psi }^{}(𝐫)\widehat{\psi }(𝐫)\hfill \\ & +\frac{1}{2}\left(\frac{4\pi a_s\mathrm{}^2}{m}\right)\hfill \\ & \times d^3𝐫_1d^3𝐫_2\widehat{\psi }^{}(𝐫_1)\widehat{\psi }^{}(𝐫_2)\delta (𝐫_1𝐫_2)\widehat{\psi }(𝐫_2)\widehat{\psi }(𝐫_1),\hfill \end{array}$$
(1)
where the kinetic energy operator is modified from the Bose-Hubbard form we used earlier, and is expressed in momentum space instead, with
$$\begin{array}{cc}\hfill \widehat{\phi }(𝐤)& =\frac{1}{(2r_b)^{3/2}}𝑑𝐫\widehat{\psi }(𝐫)e^{\mathrm{i}𝐤𝐫}.\hfill \end{array}$$
(2)
The sum over $`𝐤`$ is taken over all the (discretized) momentum coordinates. Equation (1) describes both the homogenous and trapped Bose gases. For a homogenous gas, $`\omega _0=0`$. In both cases, we use a large enough $`r_b`$ to minimize the boundary effects. We will set $`\mathrm{}=m=1`$ throughout this paper.
We discretize the cubic simulation box into an $`L\times L\times L`$ lattice. The lattice spacing is $`\varsigma =2r_b/L`$. We enumerate the real-space sites using an integral index $`i`$ ranging from 1 through $`L^3`$. The coordinate of the $`i`$-th site is given by $`𝐫_i`$. The periodic boundary condition restricts the values for the momentum coordinates $`𝐤=(k_1,\mathrm{},k_3)`$ to $`k_i=\pi n_i/r_b`$, where $`n_i`$ is an integer in the range $`L/2n_i<L/2`$. We will use the index $`q=1,2,\mathrm{},L^3`$ to enumerate the points in the momentum space; correspondingly, $`𝐤_q`$ is the momentum vector of the $`q`$-th point.
The field operators on the lattice are defined to be
$`c_i`$ $`\varsigma ^{3/2}\widehat{\psi }(𝐫_i),`$ (3)
$`b_q`$ $`\widehat{\phi }(𝐤_q).`$ (4)
The discretized Hamiltonian is therefore given by
$$\begin{array}{cc}\hfill \widehat{H}& =\frac{1}{2}\underset{q}{}𝐤_q^2b_q^{}b_q+\frac{1}{2}\left(\frac{\kappa }{\varsigma ^2}\right)\underset{i}{}|𝐫_i𝐫_0|^2c_i^{}c_i\hfill \\ & +\frac{1}{2}U\underset{i}{}\left(c_i^{}c_ic_i^{}c_ic_i^{}c_i\right),\hfill \end{array}$$
(5)
where
$`U`$ $`={\displaystyle \frac{4\pi a_s}{\varsigma ^3}},`$ (6)
$`\kappa `$ $`={\displaystyle \frac{\varsigma ^2}{a_{\mathrm{ho}}^4}},`$ (7)
and $`a_{\mathrm{ho}}\sqrt{\mathrm{}/m\omega _0}`$ is the harmonic oscillator length scale. The representation of the kinetic energy in Eq. (5) reproduces the continuum spectrum more faithfully than the real-space finite-difference form in the original Hubbard form, and allows quicker convergence with the size of the grid, $`L`$.
The contact two-body potential in the continuum is ill-defined Proukakis et al. (1998); Esry and Greene (1999) because of the ultraviolet divergence. The momentum-space interaction strength,
$`\stackrel{~}{V}_{2B}(𝐪){\displaystyle 𝑑𝐫V_{2B}(𝐫)e^{i𝐪𝐫}},`$
is uniform for any $`|𝐪|`$. The discretized Hamiltonian alleviates the problem to a large degree by introducing a mometum space cut-off $`k_c`$ and replacing the $`\delta `$-potential by an on-site interaction parameterized by the scattering length, $`a_s`$. However, the discretized two-body potential in Eq. (5) must be adjusted in order to yield the correct two-body scattering length, and $`a_s`$ in Eq. (6) must be replaced by an appropriate $`a_s^{}`$ for the lattice. Following the standard treatment, we obtain the regularized $`a_s^{}`$, which for a 3D lattice is Castin (2004)
$$\begin{array}{c}\hfill a_s^{}\frac{a_s}{12.442749a_s/\varsigma }.\end{array}$$
(8)
For the system to be in the dilute limit and the form of our two-body potential to be valid, we need the density at the lattice sites to satisfy $`\widehat{n}_i1`$.
## III Computational methods
### III.1 Quantum Monte Carlo method
#### III.1.1 General formalism for many-boson boson ground states
We briefly describe our method of computing the ground state of many bosons. A detailed account can be found in LABEL:Purwanto2004. We project the ground-state wave function $`|\mathrm{\Phi }_0`$ from a trial wave function $`|\mathrm{\Psi }_\mathrm{T}`$,
$$\begin{array}{cc}\hfill (𝒫_{\mathrm{gs}})^n|\mathrm{\Psi }_\mathrm{T}& \stackrel{n\mathrm{}}{}|\mathrm{\Phi }_0,\hfill \end{array}$$
(9)
where $`|\mathrm{\Psi }_\mathrm{T}`$ in this study is the GP solution (see Sec. III.2 for details). The projector
$`𝒫_{\mathrm{gs}}`$ $`e^{\mathrm{\Delta }\tau E_\mathrm{T}}e^{\mathrm{\Delta }\tau \widehat{H}}`$ (10)
$`=e^{\mathrm{\Delta }\tau E_\mathrm{T}}e^{\frac{1}{2}\mathrm{\Delta }\tau \widehat{K}}e^{\mathrm{\Delta }\tau \widehat{V}}e^{\frac{1}{2}\mathrm{\Delta }\tau \widehat{K}}+𝒪(\mathrm{\Delta }\tau ^2)`$ (11)
is evaluated stochastically by rewriting the two-body part into a multidimensional integral.
The two-body part of the potential in Eq. (5) can be written as a sum of the squares of one-body operators $`\widehat{V}=\frac{1}{2}_i\widehat{v}_i^2`$, where $`\widehat{v}_i\sqrt{U}c_i^{}c_i`$ is essentially the density operator. We use the following Gaussian integral identity to rewrite $`e^{\mathrm{\Delta }\tau \widehat{V}}`$ in terms of the one-body operators:
$$\begin{array}{c}\hfill e^{\frac{1}{2}\mathrm{\Delta }\tau \widehat{v}^2}=\frac{1}{\sqrt{2\pi }}_{\mathrm{}}^{\mathrm{}}𝑑\sigma e^{\frac{1}{2}\sigma ^2}e^{\sigma \overline{\sigma }\frac{1}{2}\overline{\sigma }^2}e^{\sqrt{\mathrm{\Delta }\tau }(\sigma \overline{\sigma })\widehat{v}},\end{array}$$
(12)
where the constant $`\overline{\sigma }`$ is determined below. We use an importance sampling scheme to sample the ground-state wave function, so that
$$\begin{array}{c}\hfill |\mathrm{\Phi }_0\underset{\{\varphi \}}{}w_\varphi \frac{|\varphi }{\mathrm{\Psi }_\mathrm{T}|\varphi },\end{array}$$
(13)
where each $`|\varphi `$ is a mean-field solution, i.e., a permanent consisting of identical single-particle orbitals. In practice, this means that each $`|\varphi `$ is represented by a single-particle orbital.
The projection in Eq. (9) is then realized by random walks in a manifold of mean-field solutions Zhang et al. (1997); Purwanto and Zhang (2004), which are governed by the following equation Zhang and Krakauer (2003); Purwanto and Zhang (2004):
$$\begin{array}{cc}\hfill |\varphi ^{}& =𝑑𝝈p(𝝈)\widehat{B}(𝝈\overline{𝝈})W(𝝈,\varphi )|\varphi ,\hfill \end{array}$$
(14)
where
$`p(𝝈)`$ $`={\displaystyle \underset{i}{}}{\displaystyle \frac{e^{\frac{1}{2}\sigma _i^2}}{\sqrt{2\pi }}},`$ (15)
$`\begin{array}{cc}\hfill \widehat{B}(𝝈\overline{𝝈})& =e^{\mathrm{\Delta }\tau E_\mathrm{T}}e^{\frac{1}{2}\mathrm{\Delta }\tau \widehat{K}}\left\{{\displaystyle \underset{i}{}}e^{\sqrt{\mathrm{\Delta }\tau }(\sigma _i\overline{\sigma }_i)\widehat{v}_i}\right\}\hfill \\ & \times e^{\frac{1}{2}\mathrm{\Delta }\tau \widehat{K}},\hfill \end{array}`$ (16)
$`W(𝝈,\varphi )`$ $`={\displaystyle \frac{\mathrm{\Psi }_\mathrm{T}|\widehat{B}(𝝈\overline{𝝈})|\varphi }{\mathrm{\Psi }_\mathrm{T}|\varphi }}e^{𝝈\overline{𝝈}\frac{1}{2}\overline{𝝈}\overline{𝝈}}.`$ (17)
The optimal choice of the constant vector $`\overline{𝝈}`$ is Zhang and Krakauer (2003); Purwanto and Zhang (2004):
$$\begin{array}{c}\hfill \overline{\sigma }_i=\sqrt{\mathrm{\Delta }\tau }\frac{\mathrm{\Psi }_\mathrm{T}|\widehat{v}_i|\varphi }{\mathrm{\Psi }_\mathrm{T}|\varphi }\sqrt{\mathrm{\Delta }\tau }\overline{v}_i.\end{array}$$
(18)
With this choice, the weight factor in Eq. (14) can be written in the so-called local energy form Zhang and Krakauer (2003); Purwanto and Zhang (2004):
$$\begin{array}{c}\hfill W(𝝈,\varphi )e^{\mathrm{\Delta }\tau \mathrm{\Psi }_\mathrm{T}|\widehat{H}|\varphi /\mathrm{\Psi }_\mathrm{T}|\varphi }e^{\mathrm{\Delta }\tau E_\mathrm{L}(\varphi )}.\end{array}$$
(19)
In practice, whether the local-energy or the hybrid form in Eq. (17) is more efficient will depend on the system. For the calculations in this paper, we have mostly used the local-energy form.
We initialize a population $`\{|\varphi \}`$ to mean-field solutions, e.g., $`|\mathrm{\Psi }_\mathrm{T}`$. A single random-walk step for each walker consists of updating the orbital and its associated weight $`w_\varphi `$,
$`|\varphi ^{}`$ $`\widehat{B}(𝝈\overline{𝝈})|\varphi `$ (20a)
$`w_\varphi ^{}`$ $`W(𝝈,\varphi )w_\varphi .`$ (20b)
where the auxiliary fields $`\{\sigma _i\}`$ are drawn from the Gaussian probability density function $`p(𝝈)`$.
The computation of observables is done using the back-propagation estimator Zhang et al. (1997); Purwanto and Zhang (2004),
$$\begin{array}{c}\hfill \widehat{A}_{\mathrm{bp}}=\frac{\mathrm{\Psi }_\mathrm{T}|e^{\tau _{\mathrm{bp}}\widehat{H}}\widehat{A}|\mathrm{\Phi }_0}{\mathrm{\Psi }_\mathrm{T}|e^{\tau _{\mathrm{bp}}\widehat{H}}|\mathrm{\Phi }_0},\end{array}$$
(21)
which for large enough $`\tau _{\mathrm{bp}}`$ yields the ground-state expectation value for any observable.
#### III.1.2 Phaseless approximation
The formalism above is exact. For repulsive interactions, unfortunately, $`\widehat{v}_i`$ in Eq. (12) becomes imaginary. This is similar to the phase problem in fermionic systems Zhang and Krakauer (2003), and we apply the recently developed *phaseless approximation*, which has been shown to work well in electronic-structure calculations Zhang and Krakauer (2003). This method eliminates the phase problem at the cost of a systematic bias which is dependent on the trial wave function. As we will demonstrate in benchmark calculations in Appendix A, the bias is relatively small for the bosonic systems we study here. Indeed, for all but the largest values of $`a_s`$, it is possible to perform unconstrained calculations with fixed imaginary-time, $`\beta =n\mathrm{\Delta }\tau `$, in which $`\beta `$ can be made sufficiently long that essentially exact ground-state values are obtained. Comparison with these results shows that the systematic error in the phaseless approximation is small (see Appendix A).
In the phaseless approximation, the weights $`\{w_\varphi \}`$ are restricted to real, positive values. We define the phase rotation angle $`\mathrm{\Delta }\theta `$ by
$$\begin{array}{c}\hfill \mathrm{\Delta }\theta \mathrm{}\mathrm{ln}\left(\frac{\mathrm{\Psi }_\mathrm{T}|\varphi ^{}}{\mathrm{\Psi }_\mathrm{T}|\varphi }\right).\end{array}$$
(22)
This is the complex-phase rotation of the walker’s overlap with the trial wave function as a result of the application of $`\widehat{B}(𝝈\overline{𝝈})`$ to $`|\varphi `$. In the phaseless approximation, the evolution of $`w_\varphi `$ is altered to
$$\begin{array}{c}\hfill w_\varphi ^{}\{\begin{array}{cc}\mathrm{cos}(\mathrm{\Delta }\theta )|W(𝝈,\varphi )|w_\varphi ,\hfill & |\mathrm{\Delta }\theta |<\pi /2\hfill \\ 0,\hfill & \text{otherwise}\hfill \end{array},\end{array}$$
(23)
which prevents the walkers from reaching the origin of the $`\mathrm{\Psi }_\mathrm{T}|\varphi `$-complex-plane. Equations (20a) and (23) define the algorithm of the phaseless QMC method.
In invoking the phaseless approximation, it is helpful to rearrange the two-body interaction term in $`\widehat{H}`$ such that a mean-field background is subtracted:
$$\begin{array}{c}\hfill \widehat{V}=\frac{1}{2}\underset{i}{}(\widehat{v}_i\widehat{v}_i)^2\underset{i}{}\widehat{v}_i\widehat{v}_i+\frac{1}{2}\underset{i}{}\widehat{v}_i^2,\end{array}$$
(24)
where the constant $`\widehat{v}_i`$ is the mean-field expectation value, e.g., with respect to $`|\mathrm{\Psi }_\mathrm{T}`$. The residual term involving $`(\widehat{v}_i\widehat{v}_i)`$ is then used in Eq. (12). This would have *no effect* in the exact formalism above, where, as we discussed in LABEL:Purwanto2004, the importance sampling transformation effectively introduces the background subtraction even if the bare form of $`\widehat{V}`$ is used. With the phaseless approximation, however, the rotation angle is controlled by the mixed-estimate of $`\widehat{v}_i`$. Reducing its average by subtracting the mean-field background will thus help reduce the rotation, and improve the behavior of the approximation in Eq. (23).
We note that the presence of phaseless approximation breaks the time-reversal symmetry of the ground-state projector. The forward, phaseless propagator $`(e^{\tau _{\mathrm{bp}}\widehat{H}})_{\mathrm{ph}}`$ is no longer formally equivalent to the back-propagated, phaseless propagator $`(e^{\tau _{\mathrm{bp}}\widehat{H}})_{\mathrm{ph}}^{}`$ \[see Eq. (21)\]. This results in an additional systematic error in the back-propagation estimator. The expectation value of an operator $`\widehat{A}`$ computed from back-propagation is $`\mathrm{\Phi }_0^{\prime \prime }|\widehat{A}|\mathrm{\Phi }_0^{}`$, where $`|\mathrm{\Phi }_0^{}`$ and $`|\mathrm{\Phi }_0^{\prime \prime }`$ are the approximate ground-state wave functions (normalized) in the forward- and backward-direction, respectively, and they are in general not the same. This is similarly the case in the constrained-path Monte Carlo for fermion lattice models Zhang et al. (1997); Carlson et al. (1999). It was shown Carlson et al. (1999) that the error vanishes linearly as $`|\mathrm{\Psi }_\mathrm{T}|\mathrm{\Phi }_0`$. We will further discuss the effect of the phaseless constraint in Sec. V and Appendix A.
### III.2 GP self-consistent projection and QMC trial wave functions
We solve the GP equation on the same lattice defined for QMC, using a self-consistent projection with the GP propagator $`\mathrm{exp}(\mathrm{\Delta }\tau \widehat{H}_{\mathrm{GP}})`$ Purwanto and Zhang (2004). Aside from a factor $`(N1)/N`$ in front of the interaction terms, the one-body Hamiltonian $`\widehat{H}_{\mathrm{GP}}`$ is simply Eq. (5) with the replacement
$$\begin{array}{c}\hfill c_i^{}c_i^{}c_ic_i2c_i^{}c_ic_i^{}c_ic_i^{}c_i^2,\end{array}$$
(25)
where the expectation is with respect to the GP solution. As discussed in LABEL:Purwanto2004, our QMC can be thought of as stochastically carrying out the functional integral, while GP is the saddle-point approximation.
The $`U`$ parameter in the GP calculations is given by the bare $`a_s`$ rather than the regularized $`a_s^{}`$ using Eq. (8), because the shape-independent $`\delta `$ potential has become a mean-field potential in the GP approximation. It is these GP results that we compare with.
For our QMC calculations, the trial wave function $`\mathrm{\Psi }_\mathrm{T}`$ is taken to be the solution of the GP-like projection, but with the regularized $`a_s^{}`$. This wave function is different from the correct GP solution above, which is obtained using the bare $`a_s`$. Each value of the discretization parameter $`\varsigma `$ corresponds to a different renormalized $`a_s^{}`$ \[see Eq. (8)\], and gives rise to distinctly different results, while the correct GP solution converges rapidly with $`\varsigma `$ (see Fig. 8). As the trial wave function, however, we argue that the optimal choice is the best variational solution, which is given by the corresponding mean-field calculation with the same $`a_s^{}`$.
### III.3 Bogoliubov approximation
In the Bogoliubov approximation Bogoliubov (1947); Lee et al. (1957); Wu (1959), correlation effects are treated by means of perturbation, where the zeroth-order term is the GP mean-field solution. The approach was first formulated for a homogenous Bose gas. It assumes a macroscopic occupancy of the lowest energy state ($`𝐤=\mathrm{𝟎}`$), namely $`(NN_0)N`$. For each density $`\rho =N/\mathrm{\Omega }`$ and interaction strength, the total energy per particle $`E_{\mathrm{Bog}}/N`$, momentum distribution $`\pi _{\mathrm{Bog}}(𝐤)`$, and condensate fraction $`N_0/N`$ can be written down analytically in the thermodynamic limit. The corrections to the mean-field GP results are expressed in terms of the gas parameter $`\rho a_s^3`$, which gives a measure of the deviation from the mean-field picture. Note that the bare $`a_s`$ should be used, since the regularization of the scattering length is implicitly done in the Bogoliubov approximation as is in GP.
It is important to truncate the summation over $`𝐤`$ in computing the momentum distribution and kinetic energy. This stems from the incorrect behavior of the Bogoliubov $`\pi (𝐤)`$ at large momenta: $`\pi _{\mathrm{Bog}}(𝐤)1/|𝐤|^4`$ as $`|𝐤|\mathrm{}`$. Physically, the form of the two-body potential requires that $`|𝐤|a_s1`$, therefore the contribution from $`|𝐤|`$ larger than a cutoff momentum $`k_c`$ should be excluded. We use an explicit numerical summation with the same $`𝐤`$-space grid as in QMC. This automatically limits the sum to the reciprocal lattice (excluding $`𝐤=\mathrm{𝟎}`$). In addition, it helps to correlate the finite-size effects in the two calculations, and allows for a more direct comparison of the results between Bogoliubov and QMC.
We extend the Bogoliubov approach to the inhomogeneous case using a local-density approximation (LDA), by treating each lattice site as a locally homogenous Bose gas. This is similar to the LDA approximation for electronic systems under density functional theory Kohn (1999), and we refer to it as *Bogoliubov-LDA*. The approximation is expected to be reasonable if the density is smooth and slowly varying, which is fulfilled in our dilute Bose gas systems.
The kinetic energy, for example, is a sum of two contributions under this approach: one from the curvature (inhomogeneity) of the density profile, and the other from Bogoliubov correction. Given the real-space density $`\rho (𝐫)`$, it is
$$\begin{array}{cc}\hfill \widehat{T}_{\mathrm{Bog}\text{-}\mathrm{LDA}}& =\frac{1}{2}d^3𝐫\sqrt{\rho (𝐫)}^2\sqrt{\rho (𝐫)}\hfill \\ & +d^3𝐫\stackrel{~}{T}_{\mathrm{Bog}}[\rho (𝐫)]\rho (𝐫),\hfill \end{array}$$
(26)
where the functional $`\stackrel{~}{T}_{\mathrm{Bog}}[\rho (𝐫)]`$ is the Bogoliubov kinetic energy per particle for a gas with uniform density $`\rho =\rho (𝐫)`$. More details on our Bogoliubov-LDA procedure are provided in Appendix B.)
## IV Results
In this section, we present results on the energetics, condensate fraction, density profile, and momentum distribution. Individual energy terms are computed: $`\widehat{T}`$ is the kinetic energy, $`\widehat{V}_{2\mathrm{B}}`$ the two-body interaction energy, and $`\widehat{V}_{\mathrm{trap}}`$ the external trapping potential.
In the calculations, we typically use a $`24\times 24\times 24`$ lattice, with a simulation box of linear dimension $`2r_b=14a_{\mathrm{ho}}`$. This gives us a lattice constant of $`\varsigma =0.583a_{\mathrm{ho}}`$. Our trap length is $`a_{\mathrm{ho}}=8546\text{Å}`$, which gives typical peak densities of about $`10`$ to $`40\mu \mathrm{m}^3`$ for $`100`$ to $`1000`$ particles in the trap. The lattice constant $`\varsigma `$ is large compared to our scattering lengths (up to $`a_s1000\text{Å}`$), which is consistent with the assumption in neglecting the details of the two-body potential.
### IV.1 Density Profile
Figure 1 shows the density profiles of 100 trapped bosons for three different scattering lengths. To make a connection with experiments, we show the column density
$$\begin{array}{c}\hfill \rho _y(x,z)𝑑y\rho (x,y,z),\end{array}$$
(27)
that is, the density integrated along a particular direction (e.g., the $`y`$-axis), which can be observed through optical measurements Andrews et al. (1996, 1997); Hau et al. (1998). As we increase $`a_s`$, the condensate expands due to the increasing repulsive interactions. Similarly, as we add more particles into the gas, the density profiles expands, as shown in Fig. 2.
Compared to GP, the QMC peak density is always lowered, and the QMC overall density profile is more extended. For $`a_s=80\text{Å}`$, the peak column density is lowered by $`0.5\%`$ from GP. For $`a_s=500\text{Å}`$, this difference is about $`7\%`$. Earlier many-body calculations using the correlated basis approach Fabrocini and Polls (1999, 2001) and DMC Blume and Greene (2001); DuBois and Glyde (2003) also showed the same qualitative behavior. Below we will further discuss these in connection with the energetics and momentum distribution.
### IV.2 Energetics
Figure 3 shows the ground-state energy and its individual components as a function of the interaction strength. We see that, as the scattering length $`a_s`$ is increased, the total energy increases as expected. Both GP and Bogoliubov-LDA energies are in reasonable agreement with QMC, deviating more at larger $`a_s`$. The GP energy is slightly lower than the exact results (no variational principle due to regularization), while Bogoliubov-LDA is higher. The external potential energy, $`\widehat{V}_{\mathrm{trap}}`$, also increases with $`a_s`$, which is a consequence of the expansion of the density profile with interaction, as shown in Fig. 1. The GP trap energy is lower than QMC, consistent with the result in Fig. 1 that QMC density profiles are more extended.
The kinetic and interaction energies are shown in the bottom panels of Fig. 3. The discrepancy between GP and QMC is more pronounced. In particular, the GP kinetic energy decreases monotonically with $`a_s`$, because the density profile expands and the system becomes less confined. The QMC kinetic energy, on the other hand, shows a *nonmonotonic behavior*. For small $`a_s`$, the kinetic energy decreases as $`a_s`$ is increased, tracking the GP result. At $`a_s400\text{Å}`$, however, the kinetic energy curves up and increases with $`a_s`$. The QMC interaction energy is significantly lower than the mean-field interaction energy at large $`a_s`$, and the GP result increases much more rapidly with $`a_s`$ than QMC. Indeed the QMC curve appears to turn downward at the last point, but our data is not sufficient to establish this, as it is possible that a larger systematic error from the phaseless approximation may have contributed to make the QMC result smaller (see the benchmark results in Appendix A).
From a single-particle picture, we would expect the QMC kinetic energy to be lower than that of GP, since the QMC density profiles are more extended. In reality, correlation effects become more important as $`a_s`$ increases, which raises the kinetic energy with interaction. This is illustrated clearly by considering the uniform Bose gas, for which we show corresponding results in Fig. 4. The GP ground state is a zero-momentum condensate. In the many-body ground state, interactions excite particles into higher-momentum single-particle states, raising the kinetic energy as a result.
The QMC results in the trapped gas are thus the outcome of the competition between mean-field and correlation effects.
The Bogoliubov-LDA calculations, whose results are also shown in Fig. 3, help to quantify this picture further. We use QMC density profiles in the calculation (hence the exact agreement between the Bogoliubov-LDA and QMC estimates of the trap energy in Fig. 3), although we have verified that the physics is qualitatively unchanged if the GP densities are used instead. The result shows good agreement with the full many-body calculation. In particular, the Bogoliubov kinetic energy shows an increase similar to the QMC prediction. The corresponding interaction energy is also reduced, although not as much as in QMC. Overall, the Bogoliubov results capture the basic picture and confirm that correlations are an important ingredient in the energetics of the gas.
### IV.3 Condensate fraction and momentum distribution
Figure 5 shows the condensate fraction as a function of interaction strength. GP by definition gives 100%. We see that the actual depletion is about $`4\%`$ at $`800\text{Å}`$. Again, the Bogoliubov result agrees well with QMC.
Figure 6 shows the momentum distribution for two scattering lengths: $`a_s=200\text{Å}`$ and $`500\text{Å}`$. The QMC’s momentum distribution is more peaked than GP. This translates in the real space to a more extended density profile for QMC, as is observed in Fig. 1.
The graph also shows the contribution to the kinetic energy from various $`k|𝐤|`$ regions, since the kinetic energy is related to the momentum distribution through
$$\begin{array}{cc}\hfill \widehat{T}& k^2𝑑k\pi (k)k^2.\hfill \end{array}$$
(28)
Relative to GP, the QMC distribution is depleted in the medium-$`k`$ regime, around $`ka_{\mathrm{ho}}^1`$. Part of this depletion goes to the low-momentum region near $`k=0`$, and the other to the high-$`k`$ region. At a higher $`a_s`$, the depletion shifts toward the smaller $`k`$ region. It is clear that the enhancement in the high-$`k`$ region results in the increase of the kinetic energy. The kinetic energy is strongly enhanced in the larger $`a_s`$ cases, which results in the upturn of the kinetic energy curve in Fig. 3.
A precision measurement of the momentum distribution would be useful to reveal the detailed structure of the many-body correlations in the Bose gas. Our results from a lattice do not have enough resolution to reveal whether there are finer structures in the momentum- or real-space density. (A fine structure in the density profile was predicted by the DMC calculations DuBois and Glyde (2003).) In the auxiliary-field QMC framework, a better resolution in the density profile may be obtained by choosing a more suitable basis set, such as Hartree-Fock states Esry (1997), whereby the GP solution becomes the lowest-energy state in this basis set, and also the leading solution in the ground-state wave function.
## V Discussions
### V.1 GP, modified GP, and Bogoliubov-LDA approaches
We have shown that the many-body correlations qualitatively change the behavior of the kinetic energy in the trapped Bose gas. The Bogoliubov approximation Bogoliubov (1947); Lee et al. (1957); Wu (1959) under the local density approximation (LDA), which we refer to as Bogoliubov-LDA, captures this trend quite well. The LDA provides a good way to include the correlation effects based on the homogenous Bose gas results. This is perhaps not surprising, given the diluteness of the gas.
In contrast, the mean-field GP method by construction approximates the kinetic energy only by the part that arises from the inhomogeneity of the gas, missing the portion from many-body effects. The separation of these two portions is especially clear in the homogeneous gas, as we illustrated in Sec. IV.2. This appears to be a rather generic feature of independent-particle approaches. The same would apply to the modified GP (MGP) method Braaten and Nieto (1997); Nunes (1999); Fabrocini and Polls (1999); Banerjee and Singh (2001); Fu et al. (2003), which can be viewed as the bosonic counterpart of the standard electronic structure method of LDA under density-functional theory (DFT). In that framework, the MGP equation is an outcome of using the Bogoliubov results for the uniform Bose gas as the “exchange-correlation” (xc) functional, i.e., LDA+Bogoliubov (*as opposed to* the Bogoliubov-LDA above). This method has a great advantage in that it allows self-consistent calculations. For example, the real-space density can be calculated directly and would not need to be imported as was done with the Bogoliubov-LDA. Further, it is of course possible to use exact QMC results on the uniform gas to further improve the MGP equation, and make it more like DFT-LDA. For the kinetic energy, however, the MGP would give the same qualitative results as GP, even when the exact xc-functional is used and the exact density is obtained, because the “kinetic energy” that is explicitly defined in the MGP framework is incomplete. In fact, the same would seem to apply to DFT-LDA for electronic systems. This is an important conceptual difference between MGP and Bogoliubov-LDA approach, although they are closely related and lead to the same total energy results.
### V.2 Finite-size effects and limitations of the on-site potential
There are two kinds of finite-size errors in our calculation: the error due to finite simulation box size, and the discretization error due to finite lattice constant. The first kind is easily reduced, by increasing the simulation box size, $`r_b`$. In the trapped boson calculations with $`N=100`$ particles, we have checked that $`r_b5a_{\mathrm{ho}}`$ is sufficient for $`a_s1000`$. For calculations with large values of $`N`$, we use $`r_b=7a_{\mathrm{ho}}`$ to allow simulations of large enough condensate while keeping the finite-size errors much less than our statistical error.
The discretization error from the finite lattice constant, $`\varsigma `$, is more subtle. On the one hand, sufficiently small $`\varsigma `$ should be used so the results converge to the continuum values. Figure 7 shows the convergence of the total energy. It also illustrates the effect of regularizing the scattering length, as discussed in Sec. II. In Fig. 8, we show the convergence of the density profile.
On the other hand, the lattice constant is also coupled to the on-site potential that we use, which in turn affects the detailed energetics of the system. The on-site potential effectively has finite range and strength which depend on $`\varsigma `$. This is equivalent to setting the cutoff momentum $`k_c1/\varsigma `$ in the interaction matrix elements. Figure 9 shows the total and kinetic energies as $`\varsigma `$ is varied. The total energy is less sensitive to the details of the interaction potential, as are the real-space density (see Fig. 8) and the trap energy. The dependence on $`\varsigma `$ in the kinetic and interaction energies, however, is not negligible. (This dependence is consistent with the observation of Mazzanti and co-workers Mazzanti et al. (2003) when they varied the range of their soft-sphere repulsive potential.) It is important to note that the nonmonotonic behavior of the kinetic energy is observed at all $`\varsigma `$ values. As $`\varsigma `$ is reduced, the upturn is more enhanced, indicating a stronger effect from the interactions as the potential is made narrower and harder.
Ideally, we would like to decouple the basis-size error (due to finite lattice spacing) from the effect of the details of the potential. For this purpose, the on-site pseudopotential is inadequate. The $`\delta `$-function potential is meant to be used with the short-distance contributions already “integrated out” Leggett (2001). The effects above represent corrections from the details of the interaction potential as defined by the on-site form, which change as we vary $`\varsigma `$. It is easy to see that in the limit of $`\varsigma 0`$, the gas is trivially noninteracting in the exact many-body picture Braaten and Nieto (1997), since the range of the interaction potential is zero. However, if the conditions for the validity of the potential are maintained, the corrections should be small and not affect essential properties, as we have illustrated. A better pseudopotential should have an intrinsic decay in momentum space with well-defined convergence properties.
### V.3 Bias due to phaseless approximation
The phaseless approximation, as demonstrated by the benchmarks in Appendix A, gives an excellent approximation to the true many-body ground state for weak to moderate interaction strengths. Nevertheless, systematic errors on the computed observables are expected. For example, the variational principle, that the total energy computed by QMC is an upper bound to the exact energy, is not guaranteed in the presence of phaseless approximation Zhang and Krakauer (2003); Carlson et al. (1999). We even observe this bias in the $`a_s=500\text{Å}`$ results shown in Table 2.
The systematic bias is noticeable, but remains quite small up to the largest scattering lengths we study, as can be seen from the benchmark data. It is interesting to compare the phaseless and unconstrained QMC energies in Table 2. At a large $`a_s=500\text{Å}`$, the phaseless approximation lowers the kinetic energy (as well as the interaction energy) compared to the unconstrained result. This trend is observed for all $`a_s`$ values. Since the phaseless bias increases with the interaction strength, it should lead to an *underestimation* of the upturn of the kinetic energy. Thus the nonmonotonic behavior of the kinetic energy should actually be slightly stronger than shown by QMC.
We have shown in LABEL:Purwanto2004 that the QMC results is independent of the input trial wave function $`\mathrm{\Psi }_\mathrm{T}`$. This is no longer the case in the presence of the phaseless approximation. The approximation imposes a constraint based on the overlap $`\mathrm{\Psi }_\mathrm{T}|\varphi `$, and each $`\mathrm{\Psi }_\mathrm{T}`$ in principle has different constraining properties. This dependence is very weak, however, as we observed in our calculations among trial wave functions of the same general form (GP-like).
The phaseless approximation can also affect the Trotter error, which arises from the use of a finite time step $`\mathrm{\Delta }\tau `$ in Eq. (11). This error is controllable, and can be extrapolated away by running at different values of $`\mathrm{\Delta }\tau `$. Because the rotation angle in the random walk is proportional to $`\sqrt{\mathrm{\Delta }\tau U}`$, the severity of the phaseless projection is affected by $`\mathrm{\Delta }\tau `$, as is the extent of the population fluctuation. The latter is important in back-propagation, where it is highly desirable to keep branching to a minimum. If phaseless projection causes significant loss of the population, the Trotter error will be increased. Procedures that reduce the extent of the phase projection, for example, subtracting the mean-field background shown in Eq. (24), will thus improve computational efficiency (in addition to possibly reducing the systematic error).
## VI Conclusions
We have studied the ground state of realistic systems of trapped interacting Bose atomic gases using a many-body auxiliary-field QMC method, as well as GP and the Bogoliubov method under a local density approximation. We observed the effect of correlations in the energetics, condensate fraction, real-space density profiles, and momentum distribution. The density profile is more expanded compared to the GP prediction. The momentum distribution shows enhancement in the occupation of the low- and high-momentum states. The kinetic energy, contrary to the GP estimate, is *not* monotonic with the scattering length $`a_s`$. The Bogoliubov method is able to reproduce this trend qualitatively. Additional calculations on the uniform Bose gas were performed to help understand and quantify our results.
Through this study we also further tested and developed our QMC method. We found that the phaseless approximation developed for electronic systems Zhang and Krakauer (2003) worked quite well in the context of boson calculations with repulsive calculations. Because of the simplicity of these bosonic systems compared to electronic systems, they have provided an ideal testbed and allowed us to carry out more benchmark calculations and gain additional insights on controlling the phase problem, which is crucial for making QMC more useful for a wide variety of problems. It is hoped that the formalism we developed will allow the study of many interacting Bose, Fermi, and mixed-species systems. The method can also account for different external experiment environments (1-D or 2-D, rotations, anisotropic traps, optical lattices, etc.) quite straightforwardly.
###### Acknowledgements.
It is a pleasure to thank Markus Holzmann and Henry Krakauer for stimulating discussions. We gratefully acknowledge financial support from NSF and ONR. We also thank the College of William and Mary’s Computational Science cluster (SciClone) project and the Center of Piezoelectric by Design for computing support.
## Appendix A Benchmark results on the phaseless approximation in QMC
In this appendix, we show benchmark results on the phaseless approximation in dealing the complex-phase problem, as discussed in Sec. III.1.2. We first show results on a small system for which exact diagonalization can be done. We choose a one-dimensional Bose-Hubbard system. The corresponding Gross-Pitaevskii calculation is also done at the same Hubbard parameters $`t`$, $`U`$, and $`\kappa `$. (Here $`U`$ is a fixed parameter which is the same in QMC and GP.) Table 1 compares the energetics and condensate fraction obtained using various methods: exact diagonalization, our QMC with the phaseless approximation (ph-QMC), and GP self-consistent projection.
The ph-QMC improves over GP, and in general agrees well with exact diagonalization. The bias due to the phaseless approximation is visible in the trap energy $`\widehat{V}_{\mathrm{trap}}`$. In our phaseless QMC calculation, the mean-field background was subtracted in the Hamiltonian, as shown in Eq. (24). Applying the phaseless approximation directly leads to more population fluctuations in the random walk and larger systematic errors in $`\widehat{V}_{\mathrm{trap}}`$ and $`\widehat{V}_{2\mathrm{B}}`$.
We now show calculations on a large system with realistic $`a_s`$ values. We use the unconstrained QMC (u-QMC) as the reference. For weak to moderate interaction strength, the unconstrained QMC can be carried out for a short period of time $`\tau `$ before the signal is completely lost in large Monte Carlo fluctuations. To obtain the desired accuracy, we perform many short QMC runs and average the results. For each scattering lengths, we verified that the short runs have reached convergence with respect to the projection time. The severity of the phase problem grows rapidly with $`a_s`$, and such runs are not possible for large values of $`a_s`$.
Table 2 shows the phaseless QMC with the local-energy approximation \[Eq. (19)\] for 3D trapped gas of $`100`$ atoms with $`a_s=80\text{Å}`$ and $`500\text{Å}`$. The first case represents a typical situation in the trapped atomic gas experiments far from Feshbach resonances, while the second is a medium-strength interaction deep into the range of $`a_s`$ we study. The $`\mathrm{\Delta }\tau `$ parameter was adjusted so that the Trotter error is similar to or smaller than the statistical error.
We see that the agreement between the phaseless and unconstrained calculations is good.
As a further check, we compare our QMC result on the uniform Bose gas with an earlier diffusion Monte Carlo (DMC) calculation by Giorgini and co-workers Giorgini et al. (1999), which is exact. We use their results for the soft sphere potential with large radius of $`R=5a_s`$, which best matches our situation, namely $`\varsigma 2R10a_s`$. As we show in the left panel of Fig. 10, our results agree well with their DMC energies.
## Appendix B Bogoliubov ground state
The Bogoliubov approximation for the *homogenous Bose gas* assumes a macroscopic occupancy of the lowest energy state ($`𝐤=\mathrm{𝟎}`$), namely $`(NN_0)N`$. We will work in the thermodynamic limit, $`N\mathrm{}`$ and $`\mathrm{\Omega }\mathrm{}`$, keeping the density $`\rho =N/\mathrm{\Omega }`$ finite. The creation and annihilation operators for the zero-momentum state are approximated as scalars,
$$\begin{array}{cc}\hfill \widehat{\phi }^{}(\mathrm{𝟎})& \widehat{\phi }(\mathrm{𝟎})\sqrt{N_0}.\hfill \end{array}$$
(29)
We then ignore all terms higher than quadratic in the remaining creation and annihilation operators. The form of the two-body potential also requires that $`ka_s1`$. Within this approximation, the energy per particle is given by Huang (1987)
$`\stackrel{~}{E}_{\mathrm{Bog}}`$ $`E_{\mathrm{Bog}}/N`$ (30)
$`={\displaystyle \frac{4\pi \rho a_s}{2N}}\left[N{\displaystyle \underset{𝐤0}{}}\left(\alpha _𝐤^2{\displaystyle \frac{1}{2x_𝐤^2}}\right)\right]`$ (31)
$`=2\pi \rho a_s\left(1+{\displaystyle \frac{128}{15\sqrt{\pi }}}\sqrt{\rho a_s^3}\right),`$ (32)
and the occupation of the $`𝐤`$ momentum state by <sup>1</sup><sup>1</sup>1 The continuum momentum distribution, which is often referred to in the main text, *per particle*, is given by $`\pi (𝐤)/N\stackrel{~}{\pi }(𝐤)=n(𝐤)/[(2\pi )^3\rho ]`$.
$$\begin{array}{c}\hfill n_{\mathrm{Bog}}(𝐤)=\frac{\alpha _𝐤^2}{1\alpha _𝐤^2}\text{(}𝐤\mathrm{𝟎}\text{)},\end{array}$$
(33)
where
$`x_𝐤`$ $`{\displaystyle \frac{|𝐤|}{(8\pi \rho a_s)^{1/2}}}\xi |𝐤|,`$ (34)
$`\alpha _𝐤`$ $`1+x_𝐤^2x_𝐤\sqrt{x_𝐤^2+2}.`$ (35)
The quantity $`\xi (8\pi \rho a_s)^{1/2}`$ is the healing length Leggett (2001) of the condensate. The condensate fraction is given by
$`{\displaystyle \frac{N_0}{N}}`$ $`=1{\displaystyle \frac{1}{N}}{\displaystyle \underset{𝐤\mathrm{𝟎}}{}}n(𝐤)`$ (36)
$`=1{\displaystyle \frac{8}{3}}\sqrt{{\displaystyle \frac{\rho a_s^3}{\pi }}}.`$ (37)
The kinetic energy per particle can be computed through
$$\begin{array}{c}\hfill \stackrel{~}{T}_{\mathrm{Bog}}=\frac{1}{2N}\underset{𝐤0}{}|𝐤|^2n(𝐤).\end{array}$$
(38)
The summation, however, must be performed with care, as mentioned in Sec. III.3. The analytic results for the energy and condensate fraction, Eqs. (32) and (37), are obtained by extending the summation variable to infinity, because the contribution from outside the $`ka_s1`$ region is assumed to be small. This assumption does not hold for the kinetic energy, since the sum diverges due to the unphysical nature of $`n(𝐤)`$ at large $`|𝐤|`$.
To benchmark our Bogoliubov approach, we perform QMC and Bogoliubov calculations for a homogenous Bose gas at different scattering lengths, as shown in Fig. 10. We compute the energetics and condensate fraction using three different methods: GP, Bogoliubov, and QMC. As we see here, the first-order Bogoliubov approximation estimates the energetics and condensate fraction very well for a small enough gas parameter (here $`\rho a_s^310^4`$).
Note that the condensate fraction estimated by Bogoliubov with the truncation in the sum in $`k`$-space agrees much better with QMC than the analytic Bogoliubov. The analytic Bogoliubov estimate is off, as discussed above, because it is extrapolated to an infinite box size, and it includes contributions from very high momentum states.
We note that the kinetic energy, which is very small in the small $`a_s`$ regime, is no longer negligible for larger $`a_s`$ values. For $`a_s=600\text{Å}`$, or equivalently $`\rho a_s^3=1.2\times 10^4`$, the kinetic energy (see Fig. 4) is about $`37\%`$ of the total energy. This is consistent with our discussion in Sec. IV on the balance between the mean-field and correlation effects.
We can extend the Bogoliubov analysis above to deal with the case of a *inhomogeneous, trapped gas*. We use the so-called local-density approximation (LDA) by treating each lattice site as a locally homogenous gas. The density profile $`\rho (𝐫)`$ can be estimated using GP or any other methods which provides a good approximation to the density profile. Using the same $`k`$-space lattice as QMC, we compute the “local” energetics (per particle) and condensate fraction. The density is then used to weight-average the local contributions. The Bogoliubov-LDA estimate of the kinetic energy is
$$\begin{array}{cc}\hfill \widehat{T}_{\mathrm{Bog}\text{-}\mathrm{LDA}}& =\frac{1}{2}d^3𝐫\sqrt{\rho (𝐫)}^2\sqrt{\rho (𝐫)}\hfill \\ & +d^3𝐫\stackrel{~}{T}_{\mathrm{Bog}}[\rho (𝐫)]\rho (𝐫).\hfill \end{array}$$
(39)
The interaction energy is given by
$$\begin{array}{cc}\hfill \widehat{V}_{2\mathrm{B}}_{\mathrm{Bog}\text{-}\mathrm{LDA}}& =d^3𝐫\left(\stackrel{~}{E}_{\mathrm{Bog}}[\rho (𝐫)]\stackrel{~}{T}[\rho (𝐫)]_{\mathrm{Bog}}\right)\rho (𝐫).\hfill \end{array}$$
(40)
The external trap energy is straightforward to compute, namely
$$\begin{array}{c}\hfill \widehat{V}_{\mathrm{trap}}_{\mathrm{Bog}\text{-}\mathrm{LDA}}=d^3𝐫V_{\mathrm{trap}}(𝐫)\rho (𝐫).\end{array}$$
(41) |
warning/0506/physics0506072.html | ar5iv | text | # 1 Introduction
## 1 Introduction
One of the fundamental problems of quantitative finance is to develop a description of collective dynamical properties of market prices of an ensemble of financial instruments.
Let us stress that a problem of working out an economic description of the properties of market prices is not completely solved even at the level of individual securities. A simple and very popular dynamical picture allowing transparent analytical treatment, that of a random walk, assumes a) normal distribution for the price increments and b) independence of price increments corresponding to different time intervals. Starting from the studies of Mandelbrot in the 60’th through more recent analysis there accumulated a large body of evidence that real price dynamics for individual securities reveals substantial deviations from both assumptions.
Obviously we discover a much higher complexity when moving from a single financial security to a basket of securities. At this level we expect to deal with such novel effects as a) specific non-gaussian properties of the multivariate distribution of price increments ; b) temporal autocorrelations in price changes of single securities mixed with simultaneous cross-correlations between price increments of different basket ingredients . In fact, collective price dynamics is characterized by pronounced non-gaussian properties and a complicated web of interdependencies.
In this study we apply a conditional distribution approach to scrutinize the dependence structure within an ensemble of financial instruments and the related nongaussian effects. Analysis of the value of ”response” conditioned on the ”input” having a certain magnitude enables to explicitly quantify the dependencies in the market data. Generically, dealing with a set of securities and following its temporal evolution, we can identify two types of conditional distributions which are of interest to us: a) distribution of a future price increment given that past price increments of all securities lie in a certain range; b) distribution of a price increment given that all other price increments in the same time interval lie in a certain range. A simple example of the phenomenon described by the former distribution is the lagged autocorrelation, of the latter - the simultaneous cross-sectional correlation. Let us stress that separating time-lagged dependencies (”horizontal” for further reference) from simultaneously existing ones (”vertical” for further reference) is a simplification of the generic picture which allows, however, to discuss various types of dependencies in a simple setting. A generic dependence pattern is a ”product” of both: past evolution of a subset of securities may influence future evolution of another subset. An importance of these generic ”non-diagonal” contributions was studied, in the context of profitability of a simple contrarian strategy, in . Let us also mention the recent studies of lagged conditional distributions of daily returns , in the latter reference - in relation to a particular stochastic volatility model.
Analyzing, in terms of conditional distributions, the market data on intraday price increments of a large set of liquid stocks traded in NYSE and NASDAQ we have found pronounced specific effects characterizing the conditional dynamics of price increments for both lagged and simultaneous types of dependence. Most spectacular is a relationship between the volatility of the ”response” increment and the magnitude of the ”input” one which can in simple terms be described as a dependence-induced volatility smile (”D”-smile). Another striking feature seen in the data is a dramatic reduction of the kurtosis of the conditional distribution of the ”response” increments.
To give a quantitative interpretation of these results we have developed a model description of the corresponding conditional distributions based on a multivariate non-gaussian t-Student distribution depending on both past and future price increments. Let us note that a multivariate t-Student distribution is a popular choice for analyzing the simultaneous and lagged correlations in financial dynamics. The non-gaussian nature of the model turned out to be a key element enabling to explain the dependence structures observed in the market data. In particular, conditional volatility smile and decrease of kurtosis take place even in complete absence of linear correlations. The above-described effects completely disappear, however, if one uses a multivariate gaussian distribution depending on the corresponding matrix of covariances (correlations) instead of the fat-tailed multivariate t-Student distribution.
## 2 Observed features
The object of our study is a dynamical evolution of a group of $`N=100`$ most liquid stocks from S&P 500 <sup>2</sup><sup>2</sup>2A list of stocks is given in the Appendix within a two-year time period from January 1, 2003 through December 31, 2004, characterized by the price increments $`\delta p(\tau )`$ in the time interval of length $`\tau `$. In our analysis we use two intervals of length $`\tau =6\mathrm{min}`$ and $`\tau =60\mathrm{min}`$. For an interval $`[t,t+\tau ]`$ we thus have a configuration of $`N`$ price increments $`\{\delta p^j(t)p^j(t+\tau )p^j(t)\}`$, $`j=1\mathrm{}N`$, evolving in time. Most interesting are, of course, the features of this evolution distinguishing it from that of a group of independent objects. Such cohesion can be of both simultaneous (interrelations between the values of price increments of different stocks in the same time interval) and lagged (interrelations between the price increments of the same or different stocks in different time intervals) nature.
Below we shall concentrate on the two simplest types of dependencies:
1. Interrelation between the price increments in consecutive time intervals for the same stock (”horizontal” case)
2. Interrelations between the price increments of different stocks in the same time interval (”vertical” case)
Let us start with ”horizontal” case and consider all pairs { $`\delta p^j(t),\delta p^j(t+\tau )\}`$ of stock price increments in two consecutive time intervals for some given j-th stock. Our goal is to describe probabilistic properties of the set of increments at time $`t+\tau `$ conditioned on the sign and magnitude of the increments at preceding time $`t`$. These properties are characterized by the corresponding conditional distribution constructed as follows:
* First, we normalize the price increments $`\delta p^j(t)`$ in the first interval of the pair by their unconditional standard deviation $`\sigma _{\mathrm{tot}}^j`$, $`\delta p^j(t)x^j(t)=\delta p^j(t)/\sigma _{\mathrm{tot}}^j`$
* Second, we divide the set of thus normalized increments into subintervals $`\mathrm{\Delta }_i`$ having the fixed length $`0.5`$. The total interval we consider is $`\mathrm{\Delta }=[3.25,3.25]`$. The subinterval $`\mathrm{\Delta }_1`$ thus corresponds (for j-th stock) to $`x^j[3.25,2.75]`$, etc.
* For a pair with $`x^j`$ belonging to some fixed subinterval $`\mathrm{\Delta }_i`$ we study the conditional distribution $`𝒫_{\mathrm{\Delta }_i}(y^j)`$ of the normalized price increments $`y^j=\delta p^j(t+\tau )/\sigma _{\mathrm{tot}}^j`$ in the second interval of the pair
$$𝒫_{\mathrm{\Delta }_i}(y^j)𝒫(y^j|x^j\mathrm{\Delta }_i)$$
(1)
The distribution (1) is then a ”horizontal” coarse-grained conditional distribution<sup>3</sup><sup>3</sup>3Coarse graining refers to conditioned variable $`x`$ belonging to some fixed interval $`\mathrm{\Delta }_i`$: $`x\mathrm{\Delta }_i`$.
The basic properties of the conditional distribution $`𝒫_{\mathrm{\Delta }_i}(y)`$ are conveniently summarized by the values of its lowest moments - mean $`\mu _{\mathrm{cond}}`$, standard deviation $`\sigma _{\mathrm{cond}}`$, anomalous kurtosis $`\kappa _{\mathrm{cond}}`$, etc. . In this paper we shall study the correspondingly normalized conditional mean, conditional standard deviation and conditional anomalous kurtosis. The above-described normalization allows to consider all stocks simultaneously. The normalized mean $`\mu _{\mathrm{cond}}/\sigma _{\mathrm{tot}}`$, standard deviation $`\sigma _{\mathrm{cond}}/\sigma _{\mathrm{tot}}`$ and anomalous kurtosis $`\kappa _{\mathrm{cond}}/\kappa _{\mathrm{tot}}`$, where $`\kappa _{\mathrm{tot}}`$ is an unconditional anomalous kurtosis of the increments’ distribution, of the ”horizontal” coarse-grained conditional distribution (1) (i.e. that characterizing the set of all adjacent 6-min. intervals for each stock) are plotted as a function of the rescaled initial push $`\delta p/\sigma _{\mathrm{tot}}`$ in Fig. 1 .
Let us now turn to the analysis of the ”vertical” interrelations between simultaneous price increments of different stocks The corresponding coarse-grained conditional distribution is constructed in complete analogy with the above-described ”horizontal” case:
$$𝒫_{\mathrm{\Delta }_i}(y^j)𝒫(y^j|x^k\mathrm{\Delta }_i),$$
(2)
where the conditioned variable $`x^k\delta p^k(t)/\sigma _{\mathrm{tot}}^k`$ refers to the $`k`$-th stock, and the response variable $`y^j\delta p^j(t)/\sigma _{\mathrm{tot}}^j`$ \- to the $`j`$-th one.
In Fig. 2 we show the normalized conditional mean, standard deviation and kurtosis for 6-min. intervals for the ”vertical”case.
In Fig. 3 we plot the medians of the scatterplots for the normalized conditional mean, standard deviation and kurtosis for 6-min. and 60-min. intervals, combining the ”horizontal” and ”vertical” quantities.
The analysis of Figs. 1, 2 and 3 leads to the following conclusions:
* The resulting plots for conditional mean $`\mu _{\mathrm{cond}}`$ in the ”horizontal” case are too noisy to allow unambiguous interpretation. In the ”vertical” case one observes, for both cases of $`\tau =6\mathrm{min}`$ and $`\tau =60\mathrm{min}`$, a picture consistent with that of conditional mean generated through the presence of positive correlation, see below Eqs. (3) and (3).
* The plots of the relative conditional standard deviation $`\sigma _{\mathrm{cond}}`$ in ”horizontal” and ”vertical” case are, for the both cases of $`\tau =6\mathrm{min}`$ and $`\tau =60\mathrm{min}`$, strikingly similar. For $`\tau =6\mathrm{min}`$ we observe a pronounced conditional volatility smile, or dependence-induced volatility smile (D-smile) (see a more detailed discussion of this phenomenon in the next section), such that at small $`x`$ the standard deviation of the response is smaller than the unconditional standard deviation, while in the tails it is, on contrary, larger. For $`\tau =60\mathrm{min}`$ the smile is noticeably flatter than for $`\tau =6\mathrm{min}`$. This effect can be explained by the decay of anomalous kurtosis of price increments with growing $`\tau `$, see below discussion after Eq. (3).
* The median conditional kurtosis is noticeably smaller than the unconditional one.
We see, that in both ”vertical” and ”horizontal” cases the data shows, for both scales of $`\tau =6\mathrm{min}`$ and $`\tau =60\mathrm{min}`$, the same rather nontrivial patterns: conditional volatility smile and decrease of conditional kurtosis. The origin of the first effect is discussed in the next section. We shall argue, that it is in the probabilistic dependence of the adjacent price increments, whereas the role of linear correlation effects is in fact minor.
## 3 Model
Let us now present a model that explains the phenomenona of dependence-induced volatility smile and kurtosis reduction in the coarse-grained conditional distributions described in the previous section.
At the fundamental level of description the model describing the behavior of $`N`$ securities in two adjacent time intervals is fully specified by a $`2N`$ \- dimensional probability distribution. The focus of our study is on the properties of the conditional distributions constructed from this basic enveloping $`2N`$-dimensional distribution. Generically conditional distributions are obtained by restricting the values of a subset of variables. Let us collectively denote these variables by $`𝐱`$, where $`𝐱`$ is a $`N_𝐱`$ \- dimensional vector. Generically the vector $`𝐱`$ can include increments belonging to different time intervals. We are thus dealing with a conditional distribution depending on $`N_𝐲2NN_𝐱`$ variables. If we stay within the class of elliptical distributions, the multivariate probability distribution is a function of a quadratic form $`𝒦`$ constructed from the vector $`𝐳^{}=(𝐲,𝐱)`$ and the generalized covariance matrix $`\mathrm{\Sigma }`$, $`𝒦=𝐳^{}\mathrm{\Sigma }^1𝐳`$. The covariance matrix $`\mathrm{\Sigma }`$ includes the $`N_𝐱\times N_𝐱`$ covariance matrix $`C_𝐱`$ describing the correlations within the subset of conditioned variables $`𝐱`$, the $`N_𝐲\times N_𝐲`$ covariance matrix $`C_𝐲`$ describing the correlations within the subset of the variables $`𝐲`$ and the $`N_𝐱\times N_𝐲`$ covariance matrix $`C_{\mathrm{𝐱𝐲}}`$ describing the cross-covariances between the two groups:
$$\mathrm{\Sigma }=\left(\begin{array}{cc}C_𝐲& C_{\mathrm{𝐱𝐲}}\\ C_{\mathrm{𝐱𝐲}}^{}& C_𝐱\end{array}\right)$$
(3)
At this stage we have to give an explicit description of the multivariate distribution containing the covariance matrix $`\mathrm{\Sigma }`$. As will be elucidated below, a simplest choice of a gaussian multivariate distribution does not allow to explain the phenomena of D-smile and kurtosis reduction. There is, therefore, a clear need of taking into account the non-gaussian effects. The simplest possibility of keeping a fat-tailed nature of the probability distributions of individual increments is to construct a multivariate distribution from the fat-tailed marginals. Recombination of these marginals into a multivariate distribution requires constructing an appropriate copula. This construction is not unique, so the choice is guided by simplicity and ability to reproduce basic features of market data . In what follows we will show that a multivariate t-Student distribution makes a good job in this respect, while the Gaussian multivariate distribution fails to reproduce the properties of conditional distributions observed in market data.
Let us consider a $`2N`$-dimensional t-Student distribution
$$P_S^{(2N)}=\frac{1}{\sqrt{(\pi \mu )^{2N}\xi _\mu ^{2N}\mathrm{det}\mathrm{\Sigma }}}\frac{\mathrm{\Gamma }\left(\frac{\mu +2N}{2}\right)}{\mathrm{\Gamma }\left(\frac{\mu }{2}\right)}\left[1+\frac{1}{\mu }\frac{1}{\xi _\mu }𝐳^{}\mathrm{\Sigma }^1𝐳\right]^{\frac{\mu +2N}{2}}.$$
(4)
where $`\xi _\mu =(\mu 2)/\mu `$ is a normalization factor ensuring, in particular, that the covariances computed with the distribution (4) are equal to the corresponding matrix elements of the matrix $`\mathrm{\Sigma }`$.
Fixing some particular configuration of the ”initial” increments $`𝐱=𝐱_\mathrm{𝟎}`$ leads to the conditional distribution (see, e.g., ):
$`P_S^{(N)}(𝐲|𝐱_\mathrm{𝟎})`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{(\pi (\mu +N_𝐱))^{N_𝐲}\xi _{\mu +N_𝐱}^{N_𝐲}\mathrm{det}\mathrm{\Sigma }_{𝐲|𝐱_\mathrm{𝟎}}}}}{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{\mu +N_𝐱+N_𝐲}{2}\right)}{\mathrm{\Gamma }\left(\frac{\mu }{2}\right)}}`$ (5)
$`\times `$ $`\left[1+{\displaystyle \frac{1}{\mu +N_𝐱}}{\displaystyle \frac{1}{\xi _{\mu +N_𝐱}}}(𝐲𝐲_{𝐱_\mathrm{𝟎}})^{}\mathrm{\Sigma }_{𝐲|𝐱_\mathrm{𝟎}}^1(𝐲𝐲_{𝐱_\mathrm{𝟎}})\right]^{\frac{\mu +N_𝐱+N_𝐲}{2}}.`$
The conditional distribution (5) is a multivariate $`N_𝐲`$ \- dimensional t-Student distribution with the index $`\mu +N_𝐱`$ and the following expected mean and covariance matrix:
$`𝐲_{𝐱_\mathrm{𝟎}}`$ $`=`$ $`C_{\mathrm{𝐱𝐲}}C_𝐱^1𝐱_\mathrm{𝟎}`$
$`\mathrm{\Sigma }_{𝐲|𝐱_\mathrm{𝟎}}`$ $`=`$ $`\left(C_𝐲C_{\mathrm{𝐱𝐲}}C_𝐱^1C_{\mathrm{𝐱𝐲}}^{}\right)\left[{\displaystyle \frac{\mu 2}{\mu +N_𝐱2}}\right]\left[1+{\displaystyle \frac{1}{\mu }}{\displaystyle \frac{1}{\xi _\mu }}𝒦_{𝐱_\mathrm{𝟎}}\right],`$ (6)
where $`𝒦_{𝐱_\mathrm{𝟎}}=𝐱_{\mathrm{𝟎}}^{}{}_{}{}^{}C_𝐱^1𝐱_\mathrm{𝟎}`$. Let us note, that if we had used the Gaussian multivariate distribution for constructing the conditional distribution analogous to (5), we would obtain a Gaussian conditional distribution with the following expected mean and covariance matrix:
$`𝐲_{𝐱_\mathrm{𝟎}}^G`$ $`=`$ $`C_{\mathrm{𝐱𝐲}}C_𝐱^1𝐱_\mathrm{𝟎}`$
$`\mathrm{\Sigma }_{𝐲|𝐱_\mathrm{𝟎}}^G`$ $`=`$ $`C_𝐲C_{\mathrm{𝐱𝐲}}C_𝐱^1C_{\mathrm{𝐱𝐲}}^{}`$ (7)
Comparing Eqs. (3) and (3) we see, that the expected mean is in both cases the same, whereas the expected variance in the t-Student case is a product of the gaussian expression and a $`\mu `$\- and $`𝒦_𝐱`$ \- dependent factor. An additional important phenomenon in the case of a t-Student distribution is an increase of the tail exponent determining the fat-tailedness of the distribution: $`\mu \mu +N_𝐱`$ that thereby reduces the anomalous kurtosis<sup>4</sup><sup>4</sup>4Note that the extent of this ”gaussization” depends on the number of conditioned variables which in the considered example is equal to $`N_𝐱`$.
$$\kappa =\frac{6}{\mu 4}\kappa =\frac{6}{\mu +N_𝐱4}$$
(8)
To describe the conditional volatility smile phenomenon discussed in the previous section, one clearly needs initial conditions’ depending covariances. From the formula (3) we see that in the Gaussian case this effect is absent, whereas for t-Student distribution the required dependence is manifest (see the second expression in (3) containing the factor of $`\left[1+\frac{1}{\mu }\frac{1}{\xi _\mu }𝒦_{𝐱_\mathrm{𝟎}}\right]`$). Of course, one should still prove that this dependence allows to describe the market data, see below. Nevertheless, already at this stage of our analysis, one can conclude that the phenomenon of conditional volatility smile can be explained only by non-gaussian effects - simply because the gaussian formalism does not have room for its description.
The conditional distribution Eq. (5) summarizes the impact the ”initial” configuration $`𝐱_\mathrm{𝟎}`$ has on the ”final” one $`𝐲`$.
The explanation of the conditional volatility smile and kurtosis reduction effects described in the previous section requires a simpler $`2`$-dimensional version of (4) with one-dimensional $`y`$ and $`x`$. Let us thus consider two price increments in the two consecutive time intervals for the same stock for the ”horizontal” case (or the simultaneous increments of two stocks for the ”vertical” case) and introduce the corresponding bivariate distribution
$$P_S^{(2)}(y,x)=\frac{1}{\sqrt{(\pi \mu )^2\xi _\mu ^2\mathrm{det}\mathrm{\Sigma }}}\frac{\mathrm{\Gamma }\left(\frac{\mu +2}{2}\right)}{\mathrm{\Gamma }\left(\frac{\mu }{2}\right)}\frac{1}{\left(1+\frac{1}{\mu }\frac{1}{\xi _\mu }K_\mathrm{\Sigma }(x,y)\right)^{\frac{\mu +2}{2}}}$$
(9)
Here $`\mathrm{\Sigma }`$ is a covariance matrix
$$\mathrm{\Sigma }=\left(\begin{array}{cc}\sigma _y^2& \sigma _x\sigma _yr\\ \sigma _x\sigma _yr& \sigma _x^2\end{array}\right)$$
(10)
and $`K_\mathrm{\Sigma }=(y,x)\mathrm{\Sigma }^1(y,x)^{}`$ The conditional distribution $`𝒫(y|x=x_0)`$ corresponding to the above distribution is again a t-Student distribution with the tail exponent $`\mu +1`$, conditional mean $`y_{x_0}`$ and conditional $`x_0`$ \- dependent variance $`\sigma _{y|x_0}^2`$
$`y_{x_0}`$ $`=`$ $`rx_0`$
$`\sigma _{y|x_0}^2`$ $`=`$ $`\sigma _y^2(1r^2){\displaystyle \frac{\mu 2}{\mu 1}}\left(1+{\displaystyle \frac{1}{\mu }}{\displaystyle \frac{1}{\xi _\mu }}{\displaystyle \frac{x_0^2}{\sigma _x^2}}\right)`$ (11)
Therefore the conditional distribution is more gaussian (the ratio of its anomalous kurtosis to the unconditional one is equal to $`(\mu 4)/(\mu 3)<1`$), but its standard deviation can be smaller or larger than the unconditional value $`\sigma _y`$ depending on the value of the conditioned variable $`x_0`$. The parabolic dependence of the conditional volatility on the initial push $`x_0`$ is just the feature we need to explain the D-smiles in Figs. 1, 2 and 3. The fine structure we have observed – namely, the flattening of the D-smile with growing $`\tau `$, can also be explained with the help of Eq. (3). Indeed, the coefficient at $`x_0^2`$ is equal to $`1/(\mu \xi _\mu )1/(\mu 2)`$. Now the data shows (see below Fig. 4) that for larger time intervals the unconditional anomalous kurtosis $`\kappa _\tau `$ is smaller, and the tail index $`\mu _\tau =4+6/\kappa _\tau `$ is, correspondingly, larger, leading to the desired flattening of the smile. The unconditional anomalous kurtosis $`\kappa _\tau `$ and the corresponding tail index $`\mu _\tau `$ are plotted for the ensemble of $`N=100`$ stocks considered in the paper for several intraday time intervals, in Fig. 4.
Let us also note that in the gaussian case one has $`y_x=rx_0`$ and $`\sigma _{y|x_0}^2=\sigma _y^2(1r^2)`$ so, as has been already mentioned, the gaussian probabilistic link between the price increments does not leave room for $`x_0`$ \- dependent effects in the conditional covariance matrix.
To make the correspondence with the market data quantitative we should, however, introduce a coarse-grained version of the conditional distribution $`𝒫(y|x\mathrm{\Delta })`$, where the variable $`x`$ belongs to a certain subinterval $`\mathrm{\Delta }`$:
$$𝒫(y|x\mathrm{\Delta })=\frac{_{x\mathrm{\Delta }}𝑑xP_S^{(2)}(y,x)}{_{x\mathrm{\Delta }}𝑑xP_S^{(1)}(x)}$$
(12)
We have computed the normalized mean, relative standard deviation and anomalous kurtosis of a set of conditional distributions corresponding to the same coarse-graining of the increments $`x\delta p(t)/\sigma _{\mathrm{tot}}`$ as used in the analysis of the market data in the previous section, tail index $`\mu =5`$ and a set of correlation coefficients $`r=0.25,0.5,0.75`$. The conditional mean is, of course, simply proportional to $`x_0`$. The conditional kurtosis drops to the expected $`\kappa =3`$, with small deviations. Most interesting is, of course, the behavior of the conditional standard deviation shown in Fig. 5.
We see that the model reproduces the conditional volatility smile with characteristics very similar to those observed in the market data.
A crucial point in the correct interpretation of the above result is that linear correlation (present through the correlation coefficient $`r`$) shows itself only via setting the absolute scale for the variance, see the second of Eq. (3). It is clear,that the conditional volatility smile would be present even in the complete absence of correlations ($`r=0`$). Therefore it is really appropriate to call the volatility dependence in question a dependence-induced volatility smile (D-smile). Considering for instance the ”horizontal” case, the probabilistic dependence between the increments $`\delta p(t)`$ and $`\delta p(t+1)`$ can be manifestly demonstrated by computing, e.g., the correlator of their absolute values $`G(1)=|\delta p(t)||\delta p(t+1)|_t(|\delta p|_t)^2)`$. Calculating this correlator for the bivariate t-Student distribution (9) and its Gaussian counterpart gives
$`G_G(1)`$ $`=`$ $`{\displaystyle \frac{2}{\pi }}\sigma _{\mathrm{tot}}^2r^2\left(\sqrt{1r^2}+{\displaystyle \frac{\mathrm{Arcsin}\mathrm{r}}{r}}\right)`$ (13)
$`G_S(1)`$ $`=`$ $`{\displaystyle \frac{\mu }{\pi }}\sigma _{\mathrm{tot}}^2\left[{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{\mu 2}{2}\right)\mathrm{\Gamma }\left(\frac{\mu }{2}\right)\mathrm{\Gamma }^2\left(\frac{\mu 1}{2}\right)}{\mathrm{\Gamma }^2\left(\frac{\mu }{2}\right)}}\right]+`$ (14)
$`{\displaystyle \frac{\mu }{\pi }}\sigma _{\mathrm{tot}}^2{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{\mu 2}{2}\right)}{\mathrm{\Gamma }\left(\frac{\mu }{2}\right)}}r^2\left[\sqrt{1r^2}+{\displaystyle \frac{\mathrm{Arcsin}\mathrm{r}}{r}}\right]`$
In the gaussian case the (linearly) uncorrelated variables are also independent and, indeed, the correlator (13) vanishes as $`r^2`$ at $`r0`$. In the case of t-Student distribution the correlator (14) is, on contrary, nonzero at $`r=0`$, so increments are in this case probabilistically dependent. Let us stress that this dependence is in fact imposed by the form of the unconditional distribution we have chosen. One crucial feature is that the t-Student distribution ensures, in agreement with observations, that the corresponding marginal distributions are fat-tailed. The t-Student copula we have used provides a framework in which the dependence effects are present even in the complete absence of linear correlations.
## 4 Discussion
There still remains a number of important issues related to the questions discussed in the paper that we leave for the future analysis .
First, one would like to generalize the binary-level description of simultaneous ”vertical” interdependence of stock price increments to the fully multivariate case of the influence of the $`n`$-point ”trigger” configuration $`\{\delta p_j(t)\}`$ on the move of the $`k`$-th stock in the next time interval $`\delta p_j(t+\tau )`$.
Second, perhaps more difficult issue is studying the properties of the conditional distributions for arbitrary separation of corresponding time intervals. Preliminary analysis of the market data shows the dependence of D-smile on this separation (”maturity”). This is to be expected from the fact that volatility autocorrelations decay, albeit slowly, with time. This forces to generalize the formalism we have used<sup>5</sup><sup>5</sup>5For an example of a construction of this sort see .. In any case, a big goal is to establish connection with the explicit models of volatility dynamics, see e.g. , including the leverage effects .
Finally, we would like to analyze in more details application of the nonlinear patterns we have described to portfolio optimization problems. Expected mean, volatility and degree of fat-tailedness are crucial ingredients of portfolio optimization schemes , so specific effects related to them are of clear interest in this context.
## 5 Conclusion
Let us summarize the main results of the present paper.
The focus of our analysis is on the properties of conditional distributions characterizing the probabilistic behavior of an ensemble of financial instruments. The analysis of market data in the simplest case of a binary probabilistic dependence has revealed two major effects:
* The smile-shaped dependence of conditional volatility on the magnitude of the input due to non-gaussian nature of the enveloping t-Student distribution
* A noticeable reduction of the conditional anomalous kurtosis as compared to the unconditional one
Let us also mention the flattening of the D-smile with growing time interval on which the price increments are computed.
We have constructed an explicit model characterizing the collective probabilistic pattern of an ensemble of price increments that gives a natural explanation of the above-listed phenomena. The model is based on a multinomial t-Student distribution. This theoretical framework allows to unambiguously relate the effects of a dependence-induced volatility smile and kurtosis reduction to the non-gaussian nature of the eneveloping distribution.
Acknowledgements
The authors are very grateful to Eugene Pinsky for discussions and comments.
The work of A.L. was supported by the RFBR Grant 04-02-16880, and the Scientific school support grant 1936.2003.02
## 6 Appendix
Below we give a list of stocks studied in the paper:
A, AA, ABS, ABT, ADI, ADM, AIG, ALTR, AMGN, AMD, AOC, APA, APOL, AV, AVP, AXP, BA, BBBY, BBY, BHI, BIIB, BJS, BK, BLS, BR, BSX, CA, CAH, CAT, CC, CCL, CCU, CIT, CL, COP, CTXS, CVS, CZN, DG, DE, EDS, EK, EOP, EXC, FCX, FD, FDX, FE, FISV, FITB, FRE, GENZ, GIS, HDI, HIG, HMA, HOT, HUM, JBL, JWN, INTU, KG, KMB, KMG, LH, LPX, LXK, MAT, MAS, MEL, MHS, MMM, MO, MVT, MX, MYG, NI, NKE, NTRS, PBG, PCAR, PFG, PGN, PNC, PX, RHI, ROK, SOV, SPG, STI, SUN, T, TE, TMO, TRB, TSG, UNP, UST, WHR, WY |
warning/0506/quant-ph0506197.html | ar5iv | text | # Estimating the spectrum of a density matrix with LOCC
## 1 Introduction
Estimating a mixed state density matrix optimally, when one has $`N`$ copies of it available, is a difficult problem. The problem has been solved for qubits by \[Vidal et al., 1999\], \[Bagan et al., 2004\] and by \[Hayashi and Matsumoto, 2004\] and it is known that optimal collective measurements perform strictly better than any measurement which can be implemented with local operations and classical communication (LOCC). For mixed qudits, i.e., mixed states on a Hilbert space of dimension $`d`$, not much work on finding optimal collective measurements has been done. In the present work a simpler case is studied, the estimation of the spectrum of a qudit density matrix. This problem has already been studied from the large deviation point of view by \[Keyl and Werner, 2001\] and for the qubit case by \[Bagan et al., 2005\].
In addition to being interesting in itself, spectrum estimation is useful because other problems can be reduced to it:
* Estimation of bipartite pure state entanglement. This problem has been studied for $`d=2`$ by \[Sancho and Huelga, 2000\] and by \[Acín et al., 2000\].
* Estimation of generalized Pauli channel. This problem has been studied by \[Fujiwara and Imai, 2003\] and the depolarizing channel (special case of Pauli channel) by \[Sasaki et al., 2002\].
In the present paper, an LOCC asymptotically optimal<sup>1</sup><sup>1</sup>1i.e. it performs asymptotically as well as any other measurement strategy strategy will be described. The optimality of this LOCC strategy will be established by showing that it asymptotically satisfies the quantum Cramér-Rao bound (QCRB), stated by \[Helstrom, 1976\]. The QCRB is a bound on the mean square error of “reasonable” estimators.
This paper is organized as follows. In section 2 the necessary concepts are introduced and it is specified what is meant by optimality. In section 3 the estimation strategy is described and the main result is stated more precisely (equation (4)). In section 4 the conditional mean square error matrix (MSE) is calculated, this is needed for the next two sections. A heuristic argument supporting the main result is given in section 5 and a proof will be given in section 6 (theorem 3).
## 2 Preliminaries
The density matrix $`\rho `$ ($`\rho 0,\mathrm{tr}\rho =1`$) will be parametrized in the following way:
$$\rho (p)=\underset{k=1}{\overset{d1}{}}p_k|kk|+(1\underset{l=1}{\overset{d1}{}}p_l)|dd|,$$
where $`p\mathrm{\Theta }^{d1}`$ is the parameter of interest,
$`\mathrm{\Theta }=\{(p_1,\mathrm{},p_{d1}):0p_k1,{\displaystyle \underset{k=1}{\overset{d1}{}}}p_k1\}`$
is the set of possible values of the parameter, and $`\{|1,\mathrm{},|d\}`$ is a basis of eigenvectors.
The quantum estimation problem that will be studied in this paper is that, given $`N`$ copies of a completely unknown $`\rho `$, one is only interested in estimating its eigenvalues. Some of the needed concepts and results will be introduced next for the $`N=1`$ case.
Let $`M`$ be a measurement with outcomes in a finite set $`\mathrm{\Omega }`$, i.e., a collection of matrices $`\{M_\xi :\xi \mathrm{\Omega }\}`$ satisfying $`M_\xi 0`$ and $`_{\xi \mathrm{\Omega }}M_\xi =\mathrm{𝟙}`$, and let $`\widehat{p}=(\widehat{p}_1,\mathrm{},\widehat{p}_{d1})`$ be an estimator of $`p`$, i.e., a map from $`\mathrm{\Omega }`$ to $`\mathrm{\Theta }`$. The performance of such a measurement-estimator pair will be quantified by the MSE
$$MSE(\widehat{p},p,M)_{kl}=𝔼[(\widehat{p}_kp_k)(\widehat{p}_lp_l)]=\underset{\xi \mathrm{\Omega }}{}\mathrm{tr}[\rho (p)M_\xi ](\widehat{p}_{\xi k}p_k)(\widehat{p}_{\xi l}p_l),$$
where $`𝔼f`$ means expectation of $`f`$.
The QCRB states that any unbiased<sup>2</sup><sup>2</sup>2Unbiased means that
$$𝔼\widehat{p}_k=\underset{\xi \mathrm{\Omega }}{}\mathrm{tr}[\rho (p)M_\xi ]\widehat{p}_{\xi k}=p_k.$$
measurement-estimator pair $`(\widehat{p},M)`$ of $`p`$ satisfies
$$MSE(\widehat{p},p,M)H(p)^1,$$
where $`H`$ is the quantum Fisher information (QFI) (see for example \[Helstrom, 1976\] or \[Holevo, 1982\]). The QFI can be defined as the matrix with elements
$$H(p)_{kl}=Re\mathrm{tr}[\rho (p)\lambda _k(p)\lambda _l(p)],$$
where $`\{\lambda _1(p),\mathrm{},\lambda _{d1}(p)\}`$ are the symmetric logarithmic derivatives (SLD). The SLD are defined as selfadjoint solutions to the equation
$`_k\rho (p)={\displaystyle \frac{\rho (p)\lambda _k(p)+\lambda _k(p)\rho (p)}{2}},`$ (1)
where $`_k`$ means partial derivative with respect to $`p_k`$.
The SLD for the model studied in this paper are easy to calculate, indeed, writing (1) on the basis of eigenvectors we get
$$i|[|kk||dd|]|j=\frac{p_i+p_j}{2}i|\lambda _k(p)|j,$$
or
$$\lambda _k(p)=\frac{|kk|}{p_k}\frac{|dd|}{p_d}.$$
From the SLD one can then calculate the QFI to get:
$`H(p)_{kl}={\displaystyle \frac{\delta _{kl}}{p_k}}+{\displaystyle \frac{1}{p_d}},k,l\{1,\mathrm{},d1\}`$
where $`p_d=1_{l=1}^{d1}p_l`$, the inverse of $`H`$ is
$`H(p)_{kl}^1=p_k\delta _{kl}p_kp_l,k,l\{1,\mathrm{},d1\}.`$ (2)
When one has $`N`$ copies of $`\rho `$, i.e., the model is of the form $`\rho (p)^N`$ the QCRB becomes
$$MSE(\widehat{p},p,M)^{(N)}\frac{H(p)^1}{N},$$
and this bound is valid for *any* measurement $`M`$ (i.e. LOCC or not), as long as the measurement-estimator pair $`(\widehat{p},M)`$ is unbiased.
The class of unbiased estimators, however, is too restrictive since in most practical situations one deals with biased ones. \[Gill and Levit, 1995\] used a multivariate extension of an inequality due to \[van Trees, 1968\] to prove a more general bound. From their result and an inequality due to \[Braunstein and Caves, 1994\], it can be shown that, under some regularity conditions, if $`\sqrt{N}(\widehat{p}p)\stackrel{D}{}Z(p)`$ then
$`VarZ(p)H(p)^1,`$ (3)
where “$`\stackrel{D}{}`$” means convergence in distribution. This means that the variance of the limiting distribution of any regular estimator satisfies the QCRB.
## 3 Estimation strategy
Suppose now that one knows the basis of eigenvectors, and let us consider the measurement with elements $`M_k=|kk|`$. For this measurement the probability of outcome $`k`$ is
$$\mathrm{tr}[\rho (p)M_k]=p_k.$$
Now suppose this measurement is performed on $`N`$ copies of $`\rho `$, let $`N_k`$ be the number of times that outcome $`k`$ was observed, then $`\{N_1,\mathrm{},N_{d1}\}`$ have a multinomial distribution, i.e.,
$$\mathrm{Pr}(N_1=n_1,\mathrm{},N_{d1}=n_{d1})=\frac{N!}{_{k=1}^dn_k!}\underset{k=1}{\overset{d}{}}p_k^{n_k},$$
where $`n_d=N_{k=1}^{d1}n_k`$. The estimator
$$\widehat{p}_k=\frac{N_k}{N}$$
is unbiased and a simple calculation shows that its MSE equals the inverse of the QFI divided by $`N`$ which means that it saturates the QCRB and therefore it is optimal.
This would be the whole story, except for the fact that we have assumed that the eigenbasis of $`\rho `$ is known. If the eigenbasis is not known one can try to use a two-step adaptive strategy such as the one considered by \[Gill and Massar, 2000\]. The idea is to make an initial rough estimate of $`\rho `$ on an asymptotically vanishing fraction of the copies, e.g., $`N^\mu `$ with $`0<\mu <1`$. Let $`\sigma `$ be that initial estimate of $`\rho `$ and $`|\psi _k`$ be its (not necessarily unique) eigenbasis. On the rest of the copies ($`NN^\mu `$) of $`\rho `$, the measurement with elements $`M_k=|\psi _k\psi _k|`$ is performed.
In the rest of this paper, it will be shown that this method is asymptotically optimal, i.e., it asymptotically achieves the QCRB:
$`\underset{N\mathrm{}}{lim}NMSE(\widehat{p},p,M)^{(N)}=H(p)^1,`$ (4)
provided $`\mu `$ is chosen strictly larger than $`1/2`$.
## 4 The MSE in the adaptive scheme
Let $`N_i=N^\mu `$ and $`N_f=NN^\mu `$ be the sample sizes for the first and second stages respectively. In the second stage, the probability of outcome $`k`$, given the initial estimate $`\sigma `$, is
$$q_k=\mathrm{tr}M_k\rho (p)=\psi _k|\rho (p)|\psi _k.$$
These probabilities are also a random variable.
Next the MSE of the second stage (i.e. assuming fixed $`q`$’s) will be calculated. A condition for obtaining (4) will be derived from it.
Just as before, let $`N_k`$ be the number of times that outcome $`k`$ is observed and let us estimate $`p_k`$ as
$$\widehat{p}_k=\frac{N_k}{N_f}.$$
The expectation of this estimator conditioned on $`\sigma `$is
$$𝔼[\widehat{p}_k|\sigma ]=q_k,$$
so that in general it is a biased estimator. A simple calculation shows that the MSE conditioned on the first rough estimate of $`\rho `$ is
$`𝔼[(\widehat{p}_kp_k)(\widehat{p}_lp_l)|\sigma ]={\displaystyle \frac{q_k\delta _{kl}q_kq_l}{N_f}}+(p_kq_k)(p_lq_l),`$ (5)
the second term is the square of the bias, the MSE itself is
$`MSE(\widehat{p},p,M)^{(N)}=𝔼[𝔼[(\widehat{p}_kp_k)(\widehat{p}_lp_l)|\sigma ]].`$
Comparing (2) and (5) and using the fact that $`N/N_f1`$ as $`N\mathrm{}`$, it is easy to see that in order to get (4) it is sufficient that
$`\underset{N\mathrm{}}{lim}𝔼[N(q_kp_k)(q_lp_l)]=0.`$ (6)
Indeed, if this is true, then it also holds that $`𝔼[q_k]p_k`$ and $`𝔼[q_kq_l]p_kp_l`$.
## 5 Heuristic argument
Suppose for simplicity, that all eigenvalues of $`\rho `$ are different, then one expects that after the first estimate, the eigenbasis of $`\rho `$ and the eigenbasis of $`\sigma `$ are related by a unitary matrix which is very close to the identity, i.e.,
$$|\psi _k=U|k,$$
with
$$U=\mathrm{exp}\left(i\underset{\alpha =1}{\overset{d^21}{}}\eta _\alpha T_\alpha \right)=\mathrm{}^{i\eta T},$$
where $`\{T_1,\mathrm{},T_{d^21}\}`$ is a basis of $`𝔰𝔲(d)`$ satisfying $`\mathrm{tr}T_\alpha T_\beta =\delta _{\alpha \beta }`$, $`\eta ^{d^21}`$ and $`\eta `$ is small. One can then expand $`U`$ in Taylor series about $`\eta =0`$,
$$U=\mathrm{𝟙}+i\eta T\frac{1}{2}\left(\eta T\right)^2+\mathrm{o}(\eta ^2).$$
For any decent initial estimation strategy, $`\eta `$ is expected to go to $`0`$ as $`N\mathrm{}`$ at a rate of $`N_i^{1/2}=N^{\mu /2}`$.
The expression for $`q_k`$ is
$$q_k=\underset{l}{}p_l|l|U|k|^2,$$
and
$$|l|U|k|^2=\delta _{kl}+l|\eta T|kk|\eta T|l\delta _{kl}k|(\eta T)^2|k+\mathrm{o}(\eta ^2),$$
therefore
$$q_kp_k=k|(\eta T)\rho (\eta T)|kp_kk|(\eta T)^2|k+\mathrm{o}(\eta ^2).$$
From the previous expression and the fact that $`\eta `$ goes to zero at the rate $`N^{\mu /2}`$ one can expect that
$$𝔼(q_kp_k)^2=\frac{c}{N^{2\mu }}+\mathrm{o}(N^{2\mu }),$$
where $`c`$ is a constant possibly depending on $`p`$. From the previous equation it follows that
$`\underset{n\mathrm{}}{lim}N𝔼(q_kp_k)^2=0,`$ (7)
if and only if $`\mu >1/2`$. Now, using (7) together with the Cauchy-Schwartz inequality
$$\left(𝔼[N(q_kp_k)(q_lp_l)]\right)^2𝔼[N(q_kp_k)^2]𝔼[N(q_lp_l)^2],$$
(6) follows. As pointed out before, the desired result (4) is a consequence of (6).
## 6 Rigorous argument
### 6.1 Some intermediate results
If $`\rho =\mathrm{𝟙}/d`$, then any basis chosen for the second stage will give $`(q_kp_k)=0`$, so in what follows it is assumed that $`\rho \mathrm{𝟙}/d`$, i.e., $`\rho `$ has at least two different eigenvalues.
The following intermediate result will be needed. Basically it states that if $`\rho `$ and $`\sigma `$ are close to each other, then so will be their eigenvalues and eigenspaces.
###### Lemma 1.
Let
$`\rho `$ $`={\displaystyle \underset{a=1}{\overset{n}{}}}p_a\mathrm{\Pi }_a,`$
$`\sigma `$ $`={\displaystyle \underset{k=1}{\overset{d}{}}}s_k|\psi _k\psi _k|,`$
where $`p_ap_b`$ for $`ab`$, $`2nd`$ is the number of different eigenvalues and $`\mathrm{\Pi }_a`$ is a projector onto the eigenspace corresponding to eigenvalue $`p_a`$, and let $`d_a=\mathrm{tr}\mathrm{\Pi }_a`$ be the degeneracy of $`p_a`$, also let
$$\mathrm{\Delta }=\underset{a}{\mathrm{min}}\underset{ba}{\mathrm{min}}|p_ap_b|>0.$$
If
$$d_{HS}(\rho ,\sigma )=\sqrt{\mathrm{tr}(\rho \sigma )^2}\delta <\frac{\mathrm{\Delta }}{1+\sqrt{d}},$$
then
1. $`a,k`$
$$|p_as_k|\sqrt{\psi _k|\mathrm{\Pi }_a|\psi _k}\delta ,$$
i.e., either $`p_a`$ is close to $`s_k`$ or $`|\psi _k`$ is almost orthogonal to the eigenspace corresponding to $`p_a`$.
2. $`a`$ $`k`$ such that $`|p_as_k|\delta `$ and $`k`$ $`a`$ such that $`|p_as_k|\delta `$, i.e., every eigenvalue of $`\sigma `$ is close to an eigenvalue of $`\rho `$ and vice versa. Let $`M_a=\{k:|p_as_k|\delta \}`$ and $`m_a=|M_a|>0`$. Note that $`M_aM_b=\mathrm{}`$ for $`ab`$.
3. Let $`ab`$, then if $`kM_b`$, then $`|p_as_k|\mathrm{\Delta }\delta `$ and
$$\sqrt{\psi _k|\mathrm{\Pi }_a|\psi _k}\frac{\delta }{\mathrm{\Delta }\delta },$$
i.e., if $`s_k`$ is within a distance $`\delta `$ of $`p_bp_a`$, then $`|\psi _k`$ is almost orthogonal to the eigenspace corresponding to $`p_a`$.
4. $`m_a=d_a`$, i.e., for $`\delta `$ small enough, the number of eigenvalues of $`\sigma `$ within a distance $`\delta `$ from $`p_a`$ is equal to the degeneracy of $`p_a`$.
5. $`kM_a`$,
$$|p_a\psi _k|\rho |\psi _k|c(\rho )\delta ^2,$$
where
$$c(\rho )=\frac{4(d1)}{\mathrm{\Delta }}.$$
The proof of this lemma is given in A.
Now the way in which the first rough estimation is done will be specified. For this part it is convenient to represent $`\rho `$ and $`\sigma `$ in the following way
$`\rho `$ $`={\displaystyle \frac{\mathrm{𝟙}}{d}}+\theta T,`$
$`\sigma `$ $`={\displaystyle \frac{\mathrm{𝟙}}{d}}+\widehat{\theta }T.`$
The initial measurement strategy (which will be called *plain tomography*) is to divide the initial number of copies $`N_i`$ in $`d^21`$ groups of size $`N_0=N_i/(d^21)`$, and in group $`\alpha \{1,\mathrm{},d^21\}`$ perform the measurement
$$M_\pm ^{(\alpha )}=\frac{\mathrm{𝟙}\pm T_\alpha }{2}.$$
The probabilities are
$$p_\pm ^{(\alpha )}=\frac{1\pm \theta _\alpha }{2}.$$
Let $`w_{\alpha +}`$ be the number of times that outcome $`+`$ was obtained, it is binomially distributed $`w_{\alpha +}\text{Bin}(N_0,(1+\theta _\alpha )/2)`$. The estimator for $`\theta _\alpha `$ is taken to be
$$\widehat{\theta }_\alpha =2\frac{w_{\alpha +}}{N_0}1.$$
The following result holds:
###### Lemma 2.
If $`\mu >1/2`$ then $`ϵ>0`$ and $`h0`$
$`\underset{N\mathrm{}}{lim}\left(N^h\mathrm{Pr}\left[\sqrt{N}|q_kp_k|ϵ\right]\right)=0.`$ (8)
The proof of this lemma is given in B.
### 6.2 Proof of the main result
###### Theorem 3.
If $`\mu >1/2`$ then (4) holds.
###### Proof.
Let $`X_k^{(N)}=\sqrt{N}(q_kp_k)`$, clearly $`(X_k^{(N)})^2N`$. All that needs to be proven is that
$$\underset{N\mathrm{}}{lim}𝔼[X_k^{(N)}X_l^{(N)}]=0.$$
We have that
$`|𝔼[X_k^{(N)}X_l^{(N)}]|𝔼[|X_k^{(N)}X_l^{(N)}|]\sqrt{𝔼[(X_k^{(N)})^2]𝔼[(X_l^{(N)})^2]},`$ (9)
where in the second inequality the Cauchy-Schwartz inequality has been used. Now choose any $`ϵ>0`$,
$`𝔼[(X_k^{(N)})^2]`$ $`={\displaystyle \underset{x0}{}}x\mathrm{Pr}[(X_k^{(N)})^2=x]`$
$`={\displaystyle \underset{0x<ϵ^2}{}}x\mathrm{Pr}[(X_k^{(N)})^2=x]+{\displaystyle \underset{x>ϵ^2}{}}x\mathrm{Pr}[(X_k^{(N)})^2=x]`$
$`ϵ^2\mathrm{Pr}[(X_k^{(N)})^2<ϵ^2]+N\mathrm{Pr}[(X_k^{(N)})^2ϵ^2]`$
$`ϵ^2+N\mathrm{Pr}[|X_k^{(N)}|ϵ],`$
using now (8) one gets that $`ϵ>0`$,
$$\underset{N\mathrm{}}{lim}𝔼[(X_k^{(N)})^2]ϵ^2,$$
which implies that it must be zero; this fact and (9) imply (6) and therefore the desired result (4). ∎
We have proven something about the limit of the MSE, but (3) is a bound to the variance of the limiting distribution. However, since the limit of the MSE cannot be smaller than the variance of the limiting distribution (which in this case can easily be proven to be Gaussian) it follows that our estimator achieves the bound (3).
## 7 Estimation of bipartite pure state entanglement
A bipartite entangled pure state $`|\psi _{AB}_A_B`$ can be written as (Schmidt’s decomposition)
$$|\psi _{AB}=\underset{k=1}{\overset{d}{}}\sqrt{p_k}|k|e_k,$$
where $`\{|k\}`$ and $`\{|e_k\}`$ are orthonormal basis of $`_A`$ and $`_B`$ which are both of dimension $`d`$.
The entanglement of $`|\psi _{AB}`$ can be calculated as the entropy of one of the reduced states,
$$E(|\psi _{AB})=\mathrm{tr}(\rho _A\mathrm{log}_2\rho _A)=\underset{k=1}{\overset{d}{}}p_k\mathrm{log}_2p_k,$$
where $`\rho _A=\mathrm{tr}_B|\psi _{AB}\psi _{AB}|`$, i.e., entanglement is a function of the eigenvalues of the reduced density matrix. This means that entanglement can be estimated by performing measurements on $`\rho _A`$ only, in order to estimate its spectrum. The question is whether this procedure is optimal. A quick calculation of the QFI for the parameters $`p_k`$ in the model given by $`|\psi _{AB}`$ shows that indeed the entanglement of $`|\psi _{AB}`$ can be optimally estimated by estimating the spectrum of $`\rho _A`$ using the procedure described above in this paper.
The same result<sup>3</sup><sup>3</sup>3That entanglement can be optimally estimated by estimating the spectrum of the reduced density matrix. was obtained by \[Acín et al., 2000\] for $`d=2`$ using other tools.
## 8 Conclusions
The estimation of the spectrum of a finite dimensional density matrix has been analyzed. The following LOCC procedure has been studied:
1. Perform the so called plain tomography on $`N^\mu `$ copies where $`\mu >1/2`$ and $`N`$ is the total number of copies. From this one gets an initial estimate of the whole density matrix, call it $`\sigma `$. Let $`|\psi _1,\mathrm{},|\psi _d`$ be a set of eigenvectors of $`\sigma `$.
2. Perform the measurement with elements $`M_k=|\psi _k\psi _k|`$ on the remaining $`NN^\mu `$ copies and estimate $`p_k`$ as the number of times the outcome $`k`$ was obtained divided by $`N`$.
It has been shown that the above procedure performs asymptotically as well as any measurement (including collective ones). This means that in the asymptotic regime there is no need to perform the more complicated collective measurements for the estimation of the spectrum of a density matrix (or pure bipartite entanglement).
I would like to thank Richard Gill, Madalin Guţă and Igor Grubis̆ić for their very useful comments. This research was funded by the Netherlands Organization for Scientific Research (NWO), support from the RESQ (IST-2001-37559) project of the IST-FET programme of the European Union is also acknowledged.
## Appendix A Proof of lemma 1
1. The square of the distance between $`\rho `$ and $`\sigma `$ can be written as
$`d_{HS}(\rho ,\sigma )^2`$ $`={\displaystyle \underset{k=1}{\overset{d}{}}}{\displaystyle \underset{a=1}{\overset{n}{}}}\psi _k|(\rho \sigma )\mathrm{\Pi }_a|(\rho \sigma )|\psi _k`$
$`={\displaystyle \underset{k=1}{\overset{d}{}}}{\displaystyle \underset{a=1}{\overset{n}{}}}(p_as_k)^2\psi _k|\mathrm{\Pi }_a|\psi _k\delta ^2.`$
Since all terms are nonnegative, this implies that all of them are less than or equal to $`\delta `$ and this implies point (1).
2. For point (2), only the first statement will be proven, the proof of the second is almost identical. Suppose that the opposite is true, i.e., that $`a`$ such that $`k`$ $`|p_as_k|>\delta `$ then
$`d_{HS}(\rho ,\sigma )^2`$ $`={\displaystyle \underset{k=1}{\overset{d}{}}}{\displaystyle \underset{b=1}{\overset{n}{}}}(p_bs_k)^2\psi _k|\mathrm{\Pi }_b|\psi _k`$
$`{\displaystyle \underset{k=1}{\overset{d}{}}}(p_as_k)^2\psi _k|\mathrm{\Pi }_a|\psi _k`$
$`>\delta ^2\mathrm{tr}\mathrm{\Pi }_a\delta ^2,`$
i.e., $`d_{HS}(\rho ,\sigma )>\delta `$ which is a contradiction.
3. $`|p_as_k|=|(p_ap_b)+(p_bs_k)||p_ap_b||p_bs_k|\mathrm{\Delta }\delta `$, the second statement follows from the previous inequality and point (1).
4. $`m_a`$ $`={\displaystyle \underset{kM_a}{}}\psi _k|\psi _k{\displaystyle \underset{kM_a}{}}\psi _k|\mathrm{\Pi }_a|\psi _k`$
$`=\mathrm{tr}\mathrm{\Pi }_a{\displaystyle \underset{kM_a}{}}\psi _k|\mathrm{\Pi }_a|\psi _k`$
$`\mathrm{tr}\mathrm{\Pi }_a{\displaystyle \underset{kM_a}{}}\left({\displaystyle \frac{\delta }{\mathrm{\Delta }\delta }}\right)^2`$
$`\mathrm{tr}\mathrm{\Pi }_ad\left({\displaystyle \frac{\delta }{\mathrm{\Delta }\delta }}\right)^2,`$
where point (3) has been used. Now, since $`d_a=\mathrm{tr}\mathrm{\Pi }_a`$, we get
$$m_ad_ad\left(\frac{\delta }{\mathrm{\Delta }\delta }\right)^2.$$
Since $`\delta <\mathrm{\Delta }/(1+\sqrt{d})`$,
$$d\left(\frac{\delta }{\mathrm{\Delta }\delta }\right)^2<1,$$
and since $`m_a`$ is an integer, we have that $`m_ad_a`$. Using the fact that $`_am_a=_ad_a=d`$, we get that $`m_a=d_a`$.
5. Let $`ab`$, and $`kM_a`$
$`|p_ap_b|\sqrt{\psi _k|\mathrm{\Pi }_b|\psi _k}`$ $`=|(p_as_k)+(s_kp_b)|\sqrt{\psi _k|\mathrm{\Pi }_b|\psi _k}`$
$`[|p_as_k|+|s_kp_b|]\sqrt{\psi _k|\mathrm{\Pi }_b|\psi _k}`$
$`\left[\delta \sqrt{\psi _k|\mathrm{\Pi }_b|\psi _k}+|s_kp_b|\sqrt{\psi _k|\mathrm{\Pi }_b|\psi _k}\right]`$
$`\left[\delta \sqrt{\psi _k|\mathrm{\Pi }_b|\psi _k}+\delta \right]2\delta ,`$
where points (1) and (2) have been used. Thus, we have that
$`\psi _k|\mathrm{\Pi }_b|\psi _k{\displaystyle \frac{4\delta ^2}{(p_ap_b)^2}}.`$
Now I turn to the quantity of interest,
$`|p_a\psi _k|\rho |\psi _k|`$ $`=\left|p_a{\displaystyle \underset{b}{}}p_b\psi _k|\mathrm{\Pi }_b|\psi _k\right|`$
$`=\left|{\displaystyle \underset{b}{}}(p_ap_b)\psi _k|\mathrm{\Pi }_b|\psi _k\right|`$
$`{\displaystyle \underset{b}{}}|p_ap_b|\psi _k|\mathrm{\Pi }_b|\psi _k`$
$`={\displaystyle \underset{ba}{}}|p_ap_b|\psi _k|\mathrm{\Pi }_b|\psi _k`$
$`4{\displaystyle \underset{ba}{}}{\displaystyle \frac{1}{|p_ap_b|}}\delta ^2`$
$`{\displaystyle \frac{4(d1)}{\mathrm{\Delta }}}\delta ^2=c(\rho )\delta ^2.\mathrm{}`$
## Appendix B Proof of lemma 2
Now we enumerate the eigenvalues of $`\rho `$ from $`1`$ to $`d`$ again, with some of them possibly equal. Points (2) and (4) of lemma 1, take care that for every eigenvalue of $`\rho `$, the right number of eigenvalues of $`\sigma `$ will satisfy point (5). From point (5) of lemma 1 we get that $`|q_kp_k|c(\rho )\delta ^2`$ implies $`d(\rho ,\sigma )^2\delta ^2`$, we have
$`\mathrm{Pr}\left[|q_kp_k|c(\rho )\delta ^2\right]`$ $`\mathrm{Pr}\left[d(\rho ,\sigma )^2\delta ^2\right]`$
$`=\mathrm{Pr}\left[{\displaystyle \underset{\alpha =1}{\overset{d^21}{}}}(\theta _\alpha \widehat{\theta }_\alpha )^2\delta ^2\right].`$
Since
$$\underset{\alpha =1}{\overset{d^21}{}}(\theta _\alpha \widehat{\theta }_\alpha )^2\delta ^2$$
implies that for at least one $`\alpha `$
$$(\theta _\alpha \widehat{\theta }_\alpha )^2\frac{\delta ^2}{d^21},$$
it follows that
$`\mathrm{Pr}\left[{\displaystyle \underset{\alpha =1}{\overset{d^21}{}}}(\theta _\alpha \widehat{\theta }_\alpha )^2\delta ^2\right]`$ $`1\mathrm{Pr}[\alpha ,(\theta _\alpha \widehat{\theta }_\alpha )^2<{\displaystyle \frac{\delta ^2}{d^21}}]`$
$`=1{\displaystyle \underset{\alpha =1}{\overset{d^21}{}}}\mathrm{Pr}\left[|\theta _\alpha \widehat{\theta }_\alpha |<{\displaystyle \frac{\delta }{\sqrt{d^21}}}\right]`$
$`=1{\displaystyle \underset{\alpha =1}{\overset{d^21}{}}}\mathrm{Pr}\left[\left|w_{\alpha +}{\displaystyle \frac{1+\theta _\alpha }{2}}N_0\right|<{\displaystyle \frac{N_0}{2}}{\displaystyle \frac{\delta }{\sqrt{d^21}}}\right]`$
$`=1{\displaystyle \underset{\alpha =1}{\overset{d^21}{}}}\left(1\mathrm{Pr}\left[\left|w_{\alpha +}{\displaystyle \frac{1+\theta _\alpha }{2}}N_0\right|{\displaystyle \frac{N_0}{2}}{\displaystyle \frac{\delta }{\sqrt{d^21}}}\right]\right)`$
$`1\left(12\mathrm{exp}\left[{\displaystyle \frac{\delta ^2}{2(d^21)}}N_0\right]\right)^{d^21}.`$
In the last inequality we have used a form of the Chernoff bound<sup>4</sup><sup>4</sup>4If $`X\text{Bin}(n,p)`$ then $`\mathrm{Pr}[|Xnp|\lambda ]2\mathrm{exp}(2\lambda ^2/n).`$. Thus, we finally have that
$$\mathrm{Pr}\left[|q_kp_k|c(\rho )\delta ^2\right]1\left(12\mathrm{exp}\left[\frac{\delta ^2}{2(d^21)}N_0\right]\right)^{d^21},$$
now let $`c(\rho )\delta ^2=ϵN^{1/2}`$ and substitute $`N_0`$ by its value, $`N^\mu /(d^21)`$, the result is
$`\mathrm{Pr}\left[\sqrt{N}|q_kp_k|ϵ\right]`$ $`1\left(12\mathrm{exp}\left[{\displaystyle \frac{ϵN^{\mu 1/2}}{2c(\rho )(d^21)^2}}\right]\right)^{d^21}`$
$`2(d^21)\mathrm{exp}\left[{\displaystyle \frac{ϵN^{\mu 1/2}}{2c(\rho )(d^21)^2}}\right],`$
multiplying by $`N^h`$, taking $`\mu >1/2`$ and $`N\mathrm{}`$, we get the desired result (8).
## References |
warning/0506/hep-ph0506129.html | ar5iv | text | # Gravitino dark matter in gauge mediated supersymmetry breaking
## I Introduction
Supersymmetric extensions of the standard electroweak model of particle physics come with a variety of appealing byproducts, such as the stabilization of the Higgs mass against radiative corrections, radiatively induced electro-weak symmetry breaking at the electro-weak scale, and the possibility of $`SU(3)_c\times SU(2)_L\times U(1)_Y`$ gauge coupling unification at a sufficiently high energy scale. Most of these features occur rather naturally after supersymmetry (SUSY) has been broken softly, and it is believed that ultimately such a breaking has to occur spontaneously within some theory describing all four fundamental interactions.
Despite the lack of a particularly compelling model, there is a number of proposals with various theoretical and phenomenological merits, where the spontaneous SUSY breaking usually takes place dynamically in a (hidden) sector of the theory which does not contain the standard model particles. Models can be classified according to the origin of SUSY breaking and of the soft terms, i.e. how the breaking of supersymmetry is transmitted to the low energy (visible) sector. In the so-called gauge mediated supersymmetry breaking (GMSB) models, this transmission is induced by renormalizable gauge interactions GMSB0 ; GMSB (see GR99 for a review). Particularly attractive features of these scenarios are the natural suppression of neutral current flavor changing interactions as well as a highly predictive mass spectrum that will be put to test in forthcoming collider experiments.
In theories with gauge mediation, the lightest supersymmetric particle (LSP) is the gravitino and its mass can lie anywhere in the range $`m_{3/2}1\mathrm{eV}1\mathrm{GeV}`$. Such a light gravitino (i.e. when $`m_{3/2}1\mathrm{keV}`$) is traditionally associated with a cosmological catastrophe. Indeed many studies have examined the production of light gravitinos in the early Universe in order to place stringent bounds on the post-inflationary reheating temperature (see TMY93 ; M95 ; CHKL99 and references therein). Few studies have contemplated the possibility that this gravitino LSP could make up the dark matter of the Universe BM03 ; FY02 ; FY02b ; FIY04 ; IY04 ; AD05 , mainly because in the most naive model one needs to adjust the reheating temperature as a function of the gravitino mass in order to obtain the required dark matter relic density. However Refs. BM03 ; FY02 have recognized the important cosmological rôle of the messenger particles that are part of the spectrum of all GMSB theories. In particular, it has been shown that the late decay of the lightest messenger to visible sector particles can induce a substantial amount of entropy production which would result in the dilution of the predicted gravitino abundance. As a result, the light gravitino problem could be turned into a light gravitino blessing, i.e. one would obtain suitable gravitino dark matter for arbitrarily high reheating temperatures.
These studies BM03 ; FY02 have focused on two specific couplings between the messenger and visible sectors. This motivates us to examine in more generality the possibility of producing the right amount of gravitino dark matter in GMSB scenarios. We do so in the present paper by considering all messenger-matter interactions allowed by the gauge symmetries of the theory and by phenomenology and by considering their impact on the gravitino abundance. The present study thus aims at being more exhaustive than prior investigations; on the way we will also improve on some results previously obtained. In particular we show that the coupling introduced in Ref. BM03 does not appear in minimal GMSB models and that multi-goldstino production channels modify substantially the results of Ref. FY02 . Finally we also take into account stringent constraints from big-bang nucleosynthesis and large scale structure formation.
Our study is similar in spirit to those conducted for neutralino dark matter in minimal supergravity neutralino . However since the present paper is of an exploratory nature, we approximate the mass spectrum of GMSB models by three parameters: the messenger mass scale $`M_X`$, the supersymmetric particles mass scale $`M_{\mathrm{SUSY}}1`$TeV and the gravitino mass $`m_{3/2}`$. In particular, we treat $`m_{3/2}`$ and $`M_X`$ as the fundamental parameters in our search for gravitino dark matter. We also discuss the influence of the nature and mass of the next-to-lightest supersymmetric particle (NLSP). Furthermore we calculate the velocity of the dark matter gravitinos in order to examine whether this dark matter is hot, warm or cold. For simplicity we assume that $`R`$parity holds. However, we will show that, when $`m_{3/2}10\mathrm{MeV}`$, our results remain valid even if $`R`$parity is violated. In this range the gravitino lifetime becomes much longer than the age of the Universe, so that it can be considered as stable on our cosmological time and constraints from diffuse background distortions are eluded TY00 ; MC02 . Finally we note that Ref. EOSS03 ; FST04 has discussed the possibility of gravitino dark matter in minimal supergravity models (with gravity mediation of supersymmetry breaking). The cosmology of these models is different from that of GMSB scenarios as the gravitino is heavier ($`m_{3/2}10`$GeV) and there are no messenger particles. It is found in these studies that only a limited region of parameter space can satisfy the big-bang nucleosynthesis constraints and that the reheating temperature must be tuned in order to obtain the required dark matter relic density.
The plan of this paper is as follows. In Section II we review the basics of gauge mediation models and discuss the nature of the lightest messenger particle which plays a crucial rôle in our analysis. In Section III we discuss the numerous channels of light gravitino production in the early Universe as well as gravitino dilution due to late decay of a massive particle and the relevant cosmological constraints. In Section IV we survey the various renormalizable and non-renormalizable messenger number violating interactions allowed by the gauge symmetry of the theory and discuss their consequences with respect to the light gravitino problem and gravitino dark matter. Finally in Section V, we discuss various perspectives, notably with respect to the case of $`SO(10)`$ grand unification. We restrict ourselves to GMSB scenarios in the framework of $`N=1`$ $`D=4`$ supergravity and use natural units $`\mathrm{}=c=k_\mathrm{B}=1`$; $`m_{\mathrm{Pl}}2.4\times 10^{18}`$GeV denotes the reduced Planck mass.
## II Messenger sector
### II.1 Gauge mediation of SUSY breaking
Gauge-mediated SUSY breaking (GMSB) GMSB ; GR99 , is usually implemented by adding a term
$$W=S\mathrm{\Phi }_M\overline{\mathrm{\Phi }}_M+\mathrm{\Delta }W(S,Z_i)$$
(1)
to the superpotential where $`\mathrm{\Phi }_M`$ and $`\overline{\mathrm{\Phi }}_M`$ are messenger left chiral superfields with $`SU(3)\times SU(2)\times U(1)`$ quantum numbers, whereas the spurion left chiral superfield $`S`$ and the secluded sector $`Z_i`$ fields are electroweak- and strong- interactions singlets. Upon the development of a non-vanishing vev $`S`$ of the scalar component of the spurion superfield and a SUSY-breaking expectation value of the spurion auxiliary field $`F_S`$, due to unspecified dynamics in the secluded sector $`\mathrm{\Delta }W(S,Z)`$, fermionic messengers combine into Dirac fermions of mass $`M_{X,1/2}=M_XS`$, whereas their bosonic partners mix in a mass matrix of eigenstates $`\varphi `$ and $`\overline{\varphi }`$ with masses $`M_\varphi =M_X(1F_S/M_X^2)^{1/2}`$ and $`M_{\overline{\varphi }}=M_X(1+F_S/M_X^2)^{1/2}`$. In terms of the messengers bosonic components, $`\varphi =(\mathrm{\Phi }_M^{}+\overline{\mathrm{\Phi }}_M)/\sqrt{2}`$ and $`\overline{\varphi }=(\mathrm{\Phi }_M+\overline{\mathrm{\Phi }}_M^{})/\sqrt{2}`$. Note that $`\varphi `$ denotes a set of scalar fields transforming under some representation of the grand unified gauge group, and that the mass degeneracy of this multiplet is lifted by $`D`$terms and radiative corrections.
Since the messengers share the standard model gauge interactions, the gaugino and scalar spartners acquire mass at the one- and two-loop levels respectively:
$$\stackrel{~}{m}_{1/2}\left(\frac{\alpha }{4\pi }\right)\frac{F_S}{M_X},\stackrel{~}{m}_0^2\left(\frac{\alpha }{4\pi }\right)^2\left(\frac{F_S}{M_X}\right)^2,$$
(2)
hence the quantity $`\mathrm{\Lambda }=F_S/M_X`$ is the supersymmetry breaking scale in the visible sector. Provided $`F_S/M_X100`$TeV, this generates the required order of magnitude for the soft parameters<sup>1</sup><sup>1</sup>1In the present exploratory study, we do not address the detailed features of the GMSB mass spectrum and the related electroweak symmetry breaking and fine-tuning issues.. Note that $`F_S<M_X^2`$ is mandatory otherwise one of the messengers bosons acquires a negative mass squared. This also implies $`M_X100`$TeV, and for $`M_X100`$TeV, $`F_SM_X^2`$, hence $`M_X`$ sets the mass scale for the messenger sector. In particular, $`M_\varphi M_X`$, and in the following no distinction will be made between these two mass scales, except where otherwise noted.
The gravitino mass is related to the fundamental SUSY breaking scale, $`m_{3/2}F/\sqrt{3}m_{\mathrm{Pl}}`$, with $`FF_S+_iF_{Z_i}`$ the sum of $`F`$terms in the secluded sector. We define the parameter $`kF_S/F1`$, so that $`m_{3/2}=F_S/(k\sqrt{3}m_{\mathrm{Pl}})`$. In direct gauge mediated scenarios, one expects $`k1`$, whereas in scenarios in which the transmission of supersymmetry breaking to the messenger sector is loop suppressed one may find $`k1`$. Note that one can also relate the parameters $`k`$, $`M_X`$ and $`m_{3/2}`$ via the following formula: $`m_{3/2}=\mathrm{\Lambda }M_X/\left(k\sqrt{3}m_{\mathrm{Pl}}\right)`$. Since $`\mathrm{\Lambda }`$ is tied to the electroweak scale, the latter equation allows to eliminate one parameter, which we choose to be $`k`$, in terms of $`m_{3/2}`$ and $`M_X`$, which we will treat as the fundamental parameters.
### II.2 Lightest messenger
Taken at face value, GMSB scenarios generically lead to a cosmological catastrophe, as they predict that the lightest messenger should overclose the Universe<sup>2</sup><sup>2</sup>2Obviously, similar problems can in principle arise also for the $`Z_i`$ fields present in $`\mathrm{\Delta }W`$ if the lightest $`Z_i`$ mass is larger than $`100`$TeV DGP96 . However, the issue becomes much more model-dependent here, and we will thus assume for simplicity that secluded sector fields can decay rapidly to the spurion field $`S`$. The latter field is free from such cosmological problems: even though it can be either heavier or lighter than the messengers \[see the discussion following Eq.(3)\], in the first case it decays at tree-level to messengers, while in the second it decays to gauge bosons and gauginos fairly quickly through one-loop effects which are not suppressed by supersymmetry.. In effect, messenger gauge interactions as well as those derived from Eq. (1) conserve messenger number so that the lightest messenger (a boson) is stable in this minimal version of the theory. As messengers can be easily produced in the primordial plasma thanks to their gauge interactions, their present day abundance is given by the result of a thermal freeze-out of messenger annihilation (akin to the well-known neutralino LSP freeze-out in gravity mediated SUSY breaking).
Through explicit computation, one can show that the lightest messenger generically overcloses the Universe unless its mass $`M_X10^4`$GeV DGP96 . By lightest messenger, it is understood the lightest component of $`\varphi `$ after taking into account $`D`$terms and radiative corrections. Henceforth we denote this component by $`X`$. It has been shown that if the messengers sit in $`\mathrm{𝟓}+\overline{\mathrm{𝟓}}`$ representations<sup>3</sup><sup>3</sup>3Messengers sitting in complete GUT representations preserve automatically gauge coupling unification. of $`SU(5)`$, the lightest messenger carries the gauge charges of a sneutrino $`\stackrel{~}{\nu }_L`$, and its relic abundance would be of the right order of magnitude provided its mass $`1030`$TeV DGP96 ; HMR04 . Note that the mass scale in the messenger sector is a priori unconstrained, since phenomenology constrains the ratio $`F_S/M_X`$ as discussed previously. If the messengers sit in $`\mathrm{𝟏𝟎}+\overline{\mathrm{𝟏𝟎}}`$ representations of $`SU(5)`$, the lightest messenger carries the gauge charges of a selectron $`\stackrel{~}{e}_R`$. Charged dark matter is forbidden by cosmology CDM and moreover this messenger would overclose the Universe for typical values of $`M_X`$. Finally, if the messengers sit in $`\mathrm{𝟏𝟔}+\overline{\mathrm{𝟏𝟔}}`$ representations of $`SO(10)`$, the lightest messenger is a $`\stackrel{~}{\nu }_R`$like $`SU(3)\times SU(2)\times U(1)`$ singlet. The next-to-lightest messenger can only decay to the lightest messenger by GUT scale suppressed interactions, and its lifetime of order $`10^{10}\mathrm{yrs}(M_X/100\mathrm{T}\mathrm{e}\mathrm{V})^5`$ is so long that its decay produces unacceptable distortions of the diffuse backgrounds DGP96 .
Cosmology thus forbids the lightest messenger to be stable unless messengers can be diluted to a very low abundance or the lightest messenger happens to have mass $`M_X1030`$TeV. If the post-inflationary reheating temperature is larger than $`M_X`$ and no late-time entropy production occurs, then messenger number must be violated, i.e. the Lagrangian of the theory must contain additional messenger-matter interactions. However such terms can spoil the phenomenological successes of the minimal model, in particular the absence of flavor changing neutral currents or an adequate pattern of electroweak symmetry breaking. One is thus tempted to believe that such further messenger interactions with visible sector particles are rather weak, possibly resulting from non-renormalizable operators. This will be discussed in more detail below.
Delayed messenger decay can have dramatic consequences for the gravitino problem and/or the possibility of gravitino dark matter. If a non-relativistic species comes to dominate the energy density of the early Universe and subsequently decays into visible sector particles, a secondary epoch of reheating results and is concommitant with the dilution of any pre-existing relics, such as gravitinos. The abundance of these relics may then, even for “arbritrarily” high primary reheat temperatures of the Universe after an inflationary epoch, be in accord with current observational constraints BM03 ; FY02 ; FY02b .
A crucial element in this analysis is the messenger abundance before decay. This is given, as mentioned earlier, by the thermal freeze-out of messengers, hence by their annihilation cross-section. The mass splitting in the messenger multiplet is generally small, of order $`F_S/M_XM_X`$, hence one should in principle consider the various co-annihilation channels. This task is however left to a future more refined study; we note that over most of parameter space the inclusion of co-annihilation channels should modify the relic abundance of the lightest messenger by at most a factor of order unity since the the various particles that would co-annihilate have comparable annihilation cross-sections DGP96 .
#### II.2.1 Annihilation cross-section
##### Annihilation through gauge interactions.
In the case of $`SU(5)`$ unification, Dimopoulos et al. DGP96 have calculated the annihilation cross-section of the lightest messenger through gauge interactions for $`\mathrm{𝟓}+\overline{\mathrm{𝟓}}`$ representations and parametrized it as:
$$\sigma _{XX^{}}v\frac{1}{M_X^2}\left(A+\frac{B}{x}\right),$$
(3)
with $`xM_X/T`$ and $`AB310^3`$. This calculation takes into account all annihilation channels (mediated by gauge interactions) to gauge bosons, Higgses, neutralinos, charginos, fermions and sfermions, see Ref. DGP96 .
##### Annihilation into two goldstinos.
The above calculation neglects annihilation into two goldstinos $`XX^{}\stackrel{~}{G}\stackrel{~}{G}`$. This latter occurs through a variety of diagrams: in the $`t`$ and $`u`$ channels, the annihilation takes place through the exchange of the fermionic mass eigenstate partner of $`X`$. In the $`s`$ channel, the annihilation occurs through the exchange of a graviton, a spurion and other secluded sector scalar particles. Finally, annihilation also occurs through four-point contact interactions $`XX^{}\stackrel{~}{G}\stackrel{~}{G}`$. These various contributions are triggered by various operators in the supergravity Lagrangian, taking into account the goldstino component of the gravitino $`\mathrm{\Psi }_\mu `$ and the fermionic spurion $`\psi _S`$ fields after supersymmetry breaking as follows:
$`\mathrm{\Psi }_\mu =i\sqrt{{\displaystyle \frac{2}{3}}}{\displaystyle \frac{_\mu \stackrel{~}{G}}{m_{3/2}}}+\mathrm{}`$ (4)
$`\psi _S={\displaystyle \frac{F_S}{F}}\stackrel{~}{G}+\mathrm{}.`$ (5)
For instance, the four-point contact interactions $`XX^{}\stackrel{~}{G}\stackrel{~}{G}`$ derive from the gravitino mass term in the supergravity Lagrangian after expanding the exponential of the Kähler function to second order in the scalar around their vev. The Yukawa coupling between the goldstino and the hidden sector SUSY breaking scalars $`Z_i`$, which enters the $`s`$channel exchange diagrams, also derives from this expansion. The coupling $`XX^{}Z_i`$ derives from the scalar potential trilinear couplings. Finally, the coupling between $`X`$, its fermions mass eigenstate partner and $`\stackrel{~}{G}`$ is obtained directly from the supergravity Lagrangian coupling between the gravitino and a pair of fermion-boson partners, taking the appropriate linear combination to express it in terms of the mass eigenstates after SUSY breaking.
The annihilation cross-section into two goldstinos must be calculated with care since some leading high energy contributions are expected to cancel out BR88 ; G96 ; G98 . It will be important to distinguish between the cases where the spurion mass $`M_S`$ is larger or smaller than that of the lightest messenger. Both configurations are dynamically possible: for instance, one finds in the simplest models GMSB that $`M_X^2=(\kappa \sqrt{3}\lambda )(\kappa /\lambda ^2)M_S^2`$, where $`\lambda `$ and $`\kappa `$ denote respectively the spurion self-coupling and its coupling to the messenger fields in the superpotential, and where we have neglected here the effect of the spurion coupling to the secluded sector fields following the study of DDR97 for the stability and local minima conditions. In this case the spurion is heavier than the lightest messenger when 3κ/λ
<
[-0.07cm]
2.23𝜅𝜆
<
[-0.07cm]
2.2\sqrt{3}\leq\kappa/\lambda\raisebox{-3.69899pt}{~{}\shortstack{$<$ \\
[-0.07cm] $\sim$}}~{}2.2 and lighter when κ/λ
>
[-0.07cm]
2.2𝜅𝜆
>
[-0.07cm]
2.2\kappa/\lambda\raisebox{-3.69899pt}{~{}\shortstack{$>$ \\
[-0.07cm] $\sim$}}~{}2.2.
In the light scalars and non-relativistic limit, $`s4M_X^2(1+3/x)m_{Z_i}^2`$ with $`m_{Z_i}`$ mass of the secluded sector scalar $`Z_i`$, and $`x=M_X/T`$ as above, the cross-section reads:
$$\sigma _{XX^{}\stackrel{~}{G}\stackrel{~}{G}}v=\left(\frac{F_S}{F}\right)^4\frac{1}{32\pi M_X^2}\left(1\frac{15}{4x}\right).$$
(6)
In effect, performing an expansion in $`F_S/s`$ in the matrix element, one finds that both terms of order 0 and 1 cancel among the various contributions, yielding a cross-section $`F_S^4`$. Its high energy limit $`sM_X^2,m_{Z_i}^2`$ is actually that of spurion-spurion annihilation into two goldstinos G98 . In practice, Eq. (6) applies if the spurion is much lighter than the lightest messenger since its coupling to $`XX^{}`$ dominates that of secluded sector scalars as a result from the tree-level coupling in the superpotential between $`S`$, $`\mathrm{\Phi }_M`$ and $`\overline{\mathrm{\Phi }}_M`$. However, in this limit one must also account for annihilation of $`X,X^{}`$ into spurions; this will be discussed further below.
If secluded sector scalars are heavier than the lightest messenger, the limit $`sm_{Z_i}^2`$ applies and in this case, one can neglect the $`s`$channel exchange of these scalars. The cancellation between the various diagrams occurs only to order 0 in $`(F_S/s)`$, leaving a cross-section $`F_S^2`$:
$$\sigma _{XX^{}\stackrel{~}{G}\stackrel{~}{G}}v\frac{1}{8\pi }\frac{F_S^2M_X^2}{F^4}\left(1\frac{3}{2x}\right).$$
(7)
The $`s`$channel exchange graphs of the secluded sector scalars are suppressed by $`s/m_{Z_i}^2`$. Note that the term $`F_S`$ in these expressions should be understood as the mass squared difference between the fermion and boson components of the lightest messenger multiplet, rather than as the vev of the auxiliary component of $`S`$. These two quantities differ if the superpotential includes a coupling constant $`\kappa `$, $`W\kappa S\mathrm{\Phi }\overline{\mathrm{\Phi }}`$; our choice here is $`\kappa =1`$.
The annihilation cross-section in the heavy spurion limit \[Eq. (7)\] violates the unitarity bound $`\sigma _{\mathrm{ann}}v8\pi /M_X^2`$ GK90 (in the non-relativistic regime) for $`M_X1.6310^7\mathrm{GeV}(m_{3/2}/1\mathrm{keV})^{2/3}`$. Beyond this limit the effective Lagrangian is no longer valid, and one expects sizeable contributions from multi-goldstinos production. Hence the results obtained hereafter in the region where the unitarity bound is violated are highly uncertain and model-dependent. In what follows, we assume that the cross-section saturates at the unitarity limit in this region $`M_X1.610^{11}\mathrm{GeV}(m_{3/2}/1\mathrm{GeV})^{2/3}`$ if the spurion is heavier than the lightest messenger. We also consider the other possible limit in which the cross-section follows Eq. (6), so that the comparison of these two cases will allow us to assess the impact ofthe above effects on the relic gravitino abundance.
Since phenomenology requires that $`M_XF_S/10^5\mathrm{GeV}`$ (see Section II.A), it is easy to see that annihilation into goldstinos dominates the cross-section in the heavy scalars limit for $`M_X310^6\mathrm{GeV}(m_{3/2}/1\mathrm{keV})^{2/3}`$. Hence the inclusion of this channel in the present calculation modifies rather drastically the relic abundance of the lightest messenger in this part of parameter space.
##### Annihilation into two spurions.
If $`S`$ is lighter than $`X`$, there is no problem associated with unitarity, and one can safely use Eq. (6) all throughout parameter space. One must nonetheless account for $`XX^{}SS`$ annihilation, whose cross-section reads
$$\sigma _{XX^{}SS}v=\frac{1}{64\pi M_X^2}\left(1\frac{1}{x}+(\frac{3}{2}+\frac{23}{4x})r\right)$$
(8)
to first order in $`x^1`$ and in $`r(M_S^2/M_X^2)`$, and where we neglected for simplicity the contribution of a $`\lambda S^3`$ term in the superpotential, assuming that $`\lambda \kappa 1`$.
This cross-section is comparable to the annihilation cross-section through gauge interactions given in Eq. (3) and results in the decrease of the relic abundance of the lightest messenger by a factor of order 2. For direct GMSB models in which $`F_SF`$, the annihilation channel into goldstinos becomes dominant and must be taken into account.
#### II.2.2 Relic abundance
Freeze-out of the lightest messenger annihilations occurs at a value $`x_f`$:
$$x_f\mathrm{log}\left[Q_f\left(1+\frac{B/A}{\mathrm{log}(Q_f)}\right)\frac{1}{\sqrt{\mathrm{log}(Q_f)}}\right],$$
(9)
where $`Q_f6.110^{10}(M_X/10^6\mathrm{GeV})^1A`$, and the values of $`A`$ and $`B`$ accounts for the various possible channels depicted above. In terms of this freeze-out value $`x_f`$, the relic abundance is then given by:
$$Y_X2.110^{14}\left(\frac{M_X}{10^6\mathrm{GeV}}\right)\frac{x_f}{A+B/2x_f}.$$
(10)
In the case of messenger sitting in $`\mathrm{𝟏𝟎}+\overline{\mathrm{𝟏𝟎}}`$ representations of $`SU(5)`$, the relic abundance is expected to be similar to the above to within a factor of a few, since the lightest messenger carries hypercharge. For simplicity, we thus assume that its relic abundance is also given by Eq. (10) above.
Finally, in the case of $`SO(10)`$ grand unification, the lightest messenger is a singlet under the standard model gauge interactions. As argued in Ref. LMJ05 , it can annihilate through one-loop diagrams (which dominate the exchange of tree level GUT mass bosons considered in Ref. DGP96 ) and into two goldstinos at tree level as above. This case has been discussed in some detail in the low $`M_X`$ region in Ref. LMJ05 . In Section V we sketch briefly the parameter space of $`SO(10)`$ GMSB scenarios using order of magnitude estimates of these diagrams. In the main discussion of this paper, we thus focus on $`SU(5)`$ grand unification with messengers either in $`\mathrm{𝟓}+\overline{\mathrm{𝟓}}`$ or $`\mathrm{𝟏𝟎}+\overline{\mathrm{𝟏𝟎}}`$ representations.
The cosmological scenario we have in mind is the following. As the Universe reheats to high temperature after inflation, radiation along with gravitinos and messengers are produced. As the temperature decreases, the lightest messenger annihilations cease and its abundance freezes-out. This non-relativistic lightest messenger may come to dominate the energy density if its decay to the visible sector is sufficiently delayed and its relic abundance sufficiently large. In particular, messengers come to dominate the energy density when the background temperature
$$T_{\mathrm{dom}}\frac{4}{3}M_XY_X$$
(11)
(provided $`T_{\mathrm{dom}}T_{\mathrm{dec}}`$, with $`T_{\mathrm{dec}}`$ the temperature at which messengers decay, see below). During decay of the lightest messenger, gravitinos and possibly sparticles may be produced, but the pre-existing gravitino and NLSP abundances are diluted by entropy production. Finally, at late cosmological times, the NLSP decays to the gravitino. In some cases the NLSP may decay before the lightest messenger. The final gravitino abundance is then the sum of gravitinos produced by sparticle and messenger interactions at early times and diluted later by the appropriate factor due to messenger decay, plus gravitinos produced during the secondary reheating induced by messenger decay as well as gravitinos produced in the decay of the NLSP. Then, for a given messenger decay width, which characterizes the dilution factor, one may find in the $`m_{3/2}M_X`$ plane the region where satisfactory abundances for gravitinos are found.
If $`T_{\mathrm{dom}}<T_{\mathrm{dec}}`$, the lightest messenger never comes to dominate the energy density and no entropy production ensues. This is what has been generally assumed in previous studies that derived upper limits on the post-inflationary reheating temperature from the upper limit $`\mathrm{\Omega }_{3/2}<1`$, e.g. TMY93 ; CHKL99 . This case is discussed in Section IV.A.1.
## III Gravitino production
### III.1 Production channels
The fractional contribution to the present critical density of non-relativistic gravitinos $`\stackrel{~}{G}`$ with number-to-entropy ratio $`Y_{3/2}=n/s`$ is given by
$$\mathrm{\Omega }_{3/2}h^2=2.8110^8\left(\frac{m_{3/2}}{1\mathrm{GeV}}\right)Y_{3/2}$$
(12)
where $`h`$ is the Hubble constant in units of 100 km s$`^1`$Mpc<sup>-1</sup>. The number-to-entropy ratio $`Y_{3/2}`$ is found by following the Boltzmann equation describing gravitino production:
$$\frac{\mathrm{d}Y}{\mathrm{d}T}=\frac{1}{sHT}\left(\underset{i}{}\mathrm{\Gamma }_{i\stackrel{~}{G}+\mathrm{}}n_i+\underset{i,j}{}\sigma _{i+j\stackrel{~}{G}+\mathrm{}}v_{ij}n_in_j\right),$$
(13)
which includes production by sparticle decays $`\mathrm{\Gamma }_{i\stackrel{~}{G}+\mathrm{}}`$ and scatterings $`\sigma _{i+j\stackrel{~}{G}+\mathrm{}}v_{ij}n_in_j`$, neglecting three-body and higher order interactions. In principle, one should take into account gravitino annihilation as well. However it is sufficient to approximate possible gravitino losses by imposing that the gravitino number-to-entropy ratio never exceeds its thermal equilibrium value $`Y_{\mathrm{eq}}3.7\times 10^3(g_{}/230)^1`$, with $`g_{}`$ the number of relativistic degrees of freedom in the thermal bath<sup>4</sup><sup>4</sup>4Note that this equilibrium abundance only includes spin 1/2 goldstinos, with the population of 3/2 components of the gravitino assumed to be negligible..
The production of gravitinos in the early Universe is dominated by the production of the helicity $`\pm 1/2`$ component if $`m_{3/2}M`$, where $`M`$ denotes the mass scale of particles leading to gravitino production. This helicity $`\pm 1/2`$ component is related to the goldstino through Eq. (4), associated with the breaking of local SUSY, where $`\mathrm{\Psi }_\mu `$ and $`\stackrel{~}{G}`$ denote respectively the helicity $`\pm 1/2`$ gravitino components and the goldstino F79 . In scenarios of gauge mediated SUSY breaking in which $`m_{3/2}1`$GeV$`M_{\mathrm{SUSY}}1`$TeV, the correspondence $`\mathrm{\Psi }_\mu i\sqrt{2/3}_\mu \stackrel{~}{G}/m_{3/2}`$, is generically satisfied. The production of goldstinos may then be decomposed into the contributions from particle scatterings, decay of particles before the freeze-out of the NLSP, and the possible contribution of NLSP decay.
The decay and scattering contributions of visible sector fields have been calculated in various studies using the above gravitino-goldstino equivalence (see e.g., Ref. M95 for a detailed discussion and references). Later on it has been shown that the contribution of messengers is quite significant CHKL99 . Indeed messengers couple directly to the goldstino $`\stackrel{~}{G}`$ through the superpotential term $`S\mathrm{\Phi }_M\overline{\mathrm{\Phi }}_M`$ and the fraction of goldstino $`\stackrel{~}{G}`$ comprised in the fermionic component of $`S`$, Eq. (5). We use the decay widths and cross-sections for messengers interactions leading to gravitino production given in Ref. CHKL99 .
These authors have also argued that the goldstino decouples from the visible sector fields at energies $`M_XE\sqrt{F}`$ claiming that the effective gravitino-particle-sparticle vertex is induced by loop diagrams involving messengers Lee00 . However this analysis has been performed in the limit of global supersymmetry and it thus ignores the tree level fermion-sfermion-goldstino $`(1/\sqrt{2}m_{\mathrm{Pl}})g_{ij^{}}_\mu \mathrm{\Phi }^j\overline{\chi }_R\gamma ^\nu \gamma ^\mu \mathrm{\Psi }_\nu +\mathrm{h}.\mathrm{c}.`$ and gauge boson-gaugino-goldstino $`(i/4m_{\mathrm{Pl}})\overline{\mathrm{\Psi }}_\mu \sigma ^{\rho \sigma }\gamma ^\mu \lambda _L^{(a)}[F_{\rho \sigma }^{(a)}+\stackrel{~}{F}_{\rho \sigma }^{(a)}]`$ WBook interaction terms which appear in local supersymmetry. Therefore the goldstino does not decouple from visible sector fields at energies $`M_XE\sqrt{F}`$, at variance with Ref. CHKL99 . Nevertheless, it is true, following Refs. Lee00 ; CHKL99 that production channels should include the loop-induced messenger contribution to the effective particle-sparticle-goldstino vertex at energies $`EM_X`$. This is done in the present study; we do not include thermal corrections to the cross-sections since they are found to be negligible ENOR96 .
The various contributions to $`Y_{3/2}`$ have different temperature dependences. The largest fraction of the decay contributions results from cosmic epochs when the temperature falls below the mass of the corresponding particles. In contrast the dimension-5 operators associated to visible sector particle – sparticle scatterings are most effective at high temperatures; the contribution from messengers peaks at temperatures $`TM_X`$. A simple fit to the results of Ref. CHKL99 , with the modifications according to the remarks above, gives an estimate of the amount of gravitinos produced by scatterings and decays of sparticles and messengers for a reheating temperature $`T_{RH}`$:
$`\mathrm{\Omega }_{3/2}h_{70}^2`$ $``$ $`210^4\left({\displaystyle \frac{m_{3/2}}{1\mathrm{GeV}}}\right)^1[0.6\left({\displaystyle \frac{M_3}{10^3\mathrm{GeV}}}\right)^3`$ (14)
$`+2.4\left({\displaystyle \frac{M_3}{10^3\mathrm{GeV}}}\right)^2\left({\displaystyle \frac{T_{\mathrm{RH}}}{10^5\mathrm{GeV}}}\right)]`$
$`+\mathrm{\hspace{0.17em}2}\left({\displaystyle \frac{m_{3/2}}{1\mathrm{GeV}}}\right)^1[3\left({\displaystyle \frac{M_3}{10^3\mathrm{GeV}}}\right)^4\left({\displaystyle \frac{M_X}{10^5\mathrm{GeV}}}\right)^1`$
$`+0.2\left({\displaystyle \frac{M_3}{10^3\mathrm{GeV}}}\right)^2\left({\displaystyle \frac{M_X}{10^5\mathrm{GeV}}}\right)].`$
The first two lines correspond to sparticle decays and particle-sparticle scatterings, respectively, while the last two correspond to interactions involving messenger fields. Hence the latter should only be included if $`T_{\mathrm{RH}}M_X`$. In these equations, $`M_3`$ denotes the gluino mass scale, which controls the amount of gravitinos produced. Indeed, the coupling between gauginos, gauge bosons and golstinos dominates the particle-sparticle-goldstino contributions at high temperatures and scales as the gauge coupling constant; its contribution is represented by the second term on the r.h.s. of Eq. (14). Decay contributions to $`\mathrm{\Omega }_{3/2}h^2`$ scale as $`M^3`$, with $`M`$ the mass of the decaying sparticle, see the first term on the r.h.s. of Eq. (14). The gluino appears in this term due to the larger degree of freedom for colored particles than for color singlets as well as due to colored sparticles being heavier than color singlets in the ratio $`\alpha _3/\alpha _2`$ or $`\alpha _3/\alpha _1`$ \[see Eq. (2)\]. Strictly speaking this relation applies at the messenger scale and must be corrected by renormalization group running at lower energy scales. Nevertheless the mass spectrum of GMSB models is dominated by squarks and gluinos whose masses are comparable at scale $`M_X`$. As regards the contribution of messengers, given in Eq. (14) for $`T_{\mathrm{RH}}M_X`$, the gluino mass appears only in the combination $`(4\pi /\alpha _3)M_3F_S/M_X`$, i.e. as the normalization of $`F_S`$ as a function of $`M_X`$ and the electroweak scale.
We provide Eq. (14) in order to assist the reader in interpreting the figures that follow. The results shown below are obtained from the integration of the Boltzmann equation. Also recall that Eq. (14) does not take into account all contributions to the gravitino abundance. One must notably add the contribution of gravitinos produced by annihilations of the lightest messenger, as well as lightest messenger decay and NLSP decay.
The NLSP decays into final states including one gravitino with width:
$$\mathrm{\Gamma }_{\mathrm{NLSP}\stackrel{~}{\mathrm{G}}}\frac{1}{48\pi }\frac{M_{\mathrm{NLSP}}^5}{m_{3/2}^2m_{\mathrm{Pl}}^2}.$$
(15)
The background temperature at NLSP decay is then $`T_{\mathrm{NLSP}\stackrel{~}{\mathrm{G}}}5\mathrm{MeV}(M_{\mathrm{NLSP}}/100\mathrm{G}\mathrm{e}\mathrm{V})^{5/2}(m_{3/2}/1\mathrm{M}\mathrm{e}\mathrm{V})^1`$. This decay occurs late and consequently big-bang nucleosynthesis constraints on hadronic or electromagnetic energy injection at time $`10^210^8`$sec lead to the exclusion of a significant part of parameter space, as will be seen in the following. Since one NLSP produces one gravitino, the gravitino yield of NLSP decays in terms of entropy density is simply $`Y_{\mathrm{NLSP}}`$.
In principle, the NLSPs result from a freeze-out of thermal equilibrium, hence $`Y_{\mathrm{NLSP}}`$ can be calculated in the same way as the relic abundance of the lightest messenger. However, one must not forget the possible production of NLSPs during the decay of the lightest messenger. If the decay temperature of the latter is larger than the NLSP freeze-out temperature, this contribution is washed out by NLSP annihilations. However if decay occurs later, NLSPs are regenerated to a level which depends on the time of messenger decay and on the annihilation cross-section of the lightest messenger. In more detail, if the lightest messenger decays to NLSPs before NLSPs decay in turn to gravitinos, the NLSPs have time to annihilate. In Ref. FY02 , a procedure was outlined to calculate the number density of NLSPs remaining after these further annihilations. If the lightest messenger decays to NLSPs after pre-existing NLSPs have decayed to gravitinos, the calculation of the final number of gravitinos produced is more involved and necessitates the integration of coupled Boltzmann equations. For simplicity, we assume that all NLSPs produced in messenger decay then decay instantaneously to gravitinos without annihilating. This maximizes the number of gravitinos produced hence reinforces the constraints from relic density arguments, and, in this sense, this assumption is conservative. Keeping track of NLSP annihilation before decay to gravitinos would not affect our results significantly as it is marginal in most of parameter space FY02 .
Finally there exist other potential channels of gravitino production. One is that of helicity $`\pm 3/2`$ production, which for large gravitino mass and high reheating temperature may become important. The helicity $`\pm 3/2`$ modes interact with gravitational strength only and are produced by interactions in the thermal bath in abundance $`\mathrm{\Omega }_{3/2}^{\pm 3/2}h^2(m_{3/2}/1\mathrm{G}\mathrm{e}\mathrm{V})(T_{\mathrm{RH}}/10^{14}\mathrm{GeV})`$ M95 . Therefore this contribution does not dominate in most of the $`m_{3/2}M_XT_{\mathrm{RH}}`$ parameter space. We do not take into account possible non-thermal production channels of helicity $`\pm 3/2`$ gravitinos during inflation non-th . The amount of gravitinos produced in this way depends strongly on the underlying model of inflation, eventhough it may exceed the amount of helicity $`\pm 3/2`$ modes produced by scatterings in the thermal bath in particular models.
### III.2 Gravitino Dilution
The delayed decay of non-relativistic messengers may have dramatic consequences on the abundance of any pre-existing species, such as gravitinos. In case delayed decay results in the temporary matter-domination by messenger rest mass, i.e. $`\rho _X\rho _\mathrm{r}`$ where $`\rho `$ denote energy densities and subscripts $`X`$ and “$`\mathrm{r}`$” refer to messenger and radiation, respectively, entropy production is significant and results in the severe dilution of any pre-existing number-to-entropy ratio $`Y`$. In this case the post-messenger-decay cosmic radiation temperature $`T_{\mathrm{dec}}^>`$ is substantially larger than the pre-decay temperature $`T_{\mathrm{dec}}^<`$, akin of a second reheat. Approximating decay to be instantaneous when the Hubble scale equals the decay width of the lightest messenger, $`H\mathrm{\Gamma }_X`$, one finds
$$T_{\mathrm{dec}}^>(g_>\pi ^2/90)^{1/4}\sqrt{\mathrm{\Gamma }_Xm_{\mathrm{Pl}}},$$
(16)
where $`g_>`$ denotes the number of relativistic d.o.f. at temperature $`T_{\mathrm{dec}}^>`$. If the particle decays into the visible and into an invisible sector, the decay width in Eq. (16) above should be multiplied by $`B_{\mathrm{visible}}`$, with $`B_{\mathrm{visible}}`$ the branching ratio into visible sector particles. This may be of relevance notably when the lightest messenger decays into visible sector particles and into gravitinos which do not share their energy density with the visible sector afterwards.
By equating the pre- and post-decay energy densities, the pre-decay radiation temperature $`T_{\mathrm{dec}}^<`$ is obtained in terms of $`T_{\mathrm{dec}}^>`$ and $`T_{\mathrm{dom}}`$ \[see Eq.(11)\] at which $`X`$ comes to dominate the energy density, as:
$$T_{\mathrm{dec}}^<T_{\mathrm{dec}}^>\mathrm{min}[1,\left(\frac{g_>}{g_<}\right)^{1/3}\left(\frac{T_{\mathrm{dec}}^>}{T_{\mathrm{dom}}}\right)^{1/3}].$$
(17)
Obviously, if $`X`$ does not dominate the energy density before decaying, $`T_{\mathrm{dec}}^<T_{\mathrm{dec}}^>`$, while if $`T_{\mathrm{dec}}^>T_{\mathrm{dom}}`$, one finds $`T_{\mathrm{dec}}^<T_{\mathrm{dec}}^>`$ and entropy production is very substantial. In effect the ratio of pre-decay and post-decay entropy densities, gives the entropy release $`\mathrm{\Delta }_Xs_>/s_<=g_>T_>^3/g_<T_<^3`$:
$$\mathrm{\Delta }_X\mathrm{max}[1,\frac{T_{\mathrm{dom}}}{T_{\mathrm{dec}}^>}],$$
(18)
The values of $`T_{\mathrm{dec}}^>`$ and $`T_{\mathrm{dom}}`$ are given in terms of $`Y_X`$, $`M_X`$ and $`\mathrm{\Gamma }_X`$ through Eqs. (16) and (11).
Such entropy release dilutes pre-existing densities according to: $`Y_>=Y_</\mathrm{\Delta }_X`$. Nevertheless, it should be borne in mind that, in case of high second reheat temperatures $`T_{\mathrm{dec}}^>`$, the regeneration of diluted species may occur. This effect is taken into account in our calculations by treating messenger decay as a second reheat. Note that substantial entropy release after BBN is unacceptable, and $`T_{\mathrm{dec}}^>1`$MeV is required. Eq. (16) may thus be employed to infer a fairly strict lower limit of $`\mathrm{\Gamma }_X4.3\times 10^{25}`$GeV on abundant and slowly decaying particle species in the early Universe.
Particularly interesting to cosmology is the case of significant entropy dilution. In this limit one may use Eqs. (18) and (17) to derive the entropy dilution factor for an abritrary species with mass $`M_X`$, decay width $`\mathrm{\Gamma }_X`$ and abundance $`Y_X`$, in the limit $`\mathrm{\Delta }_X1`$:
$`\mathrm{\Delta }_X`$ $``$ $`0.77g_>^{1/4}Y_X\mathrm{\Gamma }_X^{1/2}m_{\mathrm{pl}}^{1/2}M_X`$ (19)
$``$ $`28\left({\displaystyle \frac{M_X}{10^8\mathrm{GeV}}}\right)\left({\displaystyle \frac{Y_X}{10^{10}}}\right)\left({\displaystyle \frac{\mathrm{\Gamma }_X}{10^{25}\mathrm{GeV}}}\right)^{1/2}\left({\displaystyle \frac{g_>}{10}}\right)^{1/4},`$
where it is understood that if $`\mathrm{\Delta }_X1`$ is found, it ought to be substituted by $`\mathrm{\Delta }_X=1`$.
If there exists a whole tower of $`N`$ unstable, but long-lived particles, with abundances $`Y_i`$, masses $`M_i`$ and decay widths $`\mathrm{\Gamma }_i`$ for particle $`i`$, the final dilution factor is determined solely by the properties of the slowest decaying messenger. In particular, all the equations above may be employed as if any prior decays had not occurred.
### III.3 Cosmological constraints
#### III.3.1 Hot, warm or cold dark matter
Collisionless damping during the radiation era leads to the erasure of power in density fluctuations below a length scale (the free-streaming scale) that is mostly determined by the time at which dark matter particles become non-relativistic, or equivalently by their velocity extrapolated to zero redshift. A particle of mass $`m`$ that thermally decouples from the plasma when relativistic has a present-day velocity:
$$v_00.018\mathrm{km}\mathrm{s}^1(g_{,\mathrm{dec}}/230)^{1/3}(m/1\mathrm{k}\mathrm{e}\mathrm{V})^1,$$
(20)
with $`g_{,\mathrm{dec}}`$ the number of d.o.f. at decoupling. Cosmological data on the power spectrum of density fluctuations allow to place constraints on the mass of the particle. One finds $`m1`$keV from the requirement that the Universe has reionized by $`z6`$ BHO01 or from the measurement of the power spectrum in the Lyman $`\alpha `$ forest DDT03 . If reionization has occurred as early as $`z17`$, as suggested by the recent WMAP data WMAP , then a mass larger than $`10`$keV seems required YSHS03 .
These constraints are important to our analysis and we keep track of the average velocity extrapolated to $`z=0`$. Obviously the limits between hot, warm and cold matter are fuzzy, and we choose to qualify as warm dark matter particles with velocity $`0.0018\mathrm{km}\mathrm{s}^1v_00.054\mathrm{km}\mathrm{s}^1`$, corresponding to particle masses $`0.3\mathrm{keV}m10\mathrm{keV}`$ (for freeze-out from thermal equilibrium as above). For velocities above the upper limit or below the lower limit, we mean hot or cold dark matter respectively. It is important to note that entropy production after decoupling of the gravitinos cools down these dark matter particles according to: $`v_0v_0/\mathrm{\Delta }_X^{1/3}`$.
For gravitinos produced by out-of-equilibrium processes, notably by the late decay of a massive particle, the above relation between mass and velocity is modifed. Assuming the outgoing gravitino carries a momentum of half the mass $`M`$ of the decaying particle, the present velocity reads:
$$v_0(M/2m_{3/2})(3.91/g_{,\mathrm{dec}})^{1/3}(T_0/T_{\mathrm{dec}}),$$
(21)
with $`g_{,\mathrm{dec}}`$ the number of d.o.f. at decay, $`T_0`$ the present cosmic background temperature and $`T_{\mathrm{dec}}`$ the temperature at decay. Note that $`T_{\mathrm{dec}}`$ generally depends on $`M`$ so that the dependence between the nature of gravitino dark matter (cold/warm/hot) and the mass of the decaying particle is not necessarily trivial. For instance, one can show that a decaying NLSP produces hot/warm dark matter if its mass $`500`$GeV.
If the decay occurs at temperatures sufficiently high that the gravitino can interact and thermalize, one should rather use Eq. (20). However at temperatures $`T100`$GeV the gravitino has decoupled from the thermal plasma G96 , mainly because the sparticles have decoupled themselves. Therefore the decay of the NLSP or of the lightest messenger (if sufficiently late) generally produces highly relativistic gravitinos.
#### III.3.2 Big-bang nucleosynthesis constraints
Due to the different channels of gravitino production, one generally finds that gravitinos are made of two generic sub-populations: one that has been produced by equilibrium processes and another made of hot gravitinos produced by out-of-equilibrium decays. It may be that the latter are so highly relativistic that they form a hot dark matter component. However their impact on the formation of large scale structure may be negligible if their contribution to the gravitino energy density is negligible. In this particular case, constraints from big-bang nucleosynthesis on extra degrees of freedom may apply and constrain this population. We take the BBN constraints on extra degrees of freedom to be $`\delta g1.8`$ (corresponding to 1 extra neutrino family allowed) BBN-dof . Gravitinos that are relativistic at the time of BBN and carry energy density $`\rho _{3/2}^\mathrm{r}`$ contribute to the level $`\delta g_{\mathrm{BBN}}/g_{\mathrm{BBN}}^s\rho _{3/2}^\mathrm{r}/\rho _\mathrm{r}^s`$, with $`\rho _\mathrm{r}^s`$ the standard radiation energy density at the onset of BBN with $`g_{\mathrm{BBN}}^s=10.75`$. Gravitinos that were once in thermal equilibrium or that were produced by equilibrium processes at temperature $`T`$ (d.o.f. $`g`$) carry characteristic momentum $`p_{3/2}3T_{\mathrm{BBN}}(g_{\mathrm{BBN}}^s/g^s)^{1/3}`$ at BBN, with $`T_{\mathrm{BBN}}1`$MeV. Assuming $`g^s230`$, these gravitinos are relativistic if $`m_{3/2}1`$MeV and their contribution to the energy density is $`\delta g_{\mathrm{BBN}}/g_{\mathrm{BBN}}^sY_{3/2}^\mathrm{r}`$ with $`Y_{3/2}^\mathrm{r}=n_{3/2}^\mathrm{r}/s`$, hence it is negligible due to the upper bound on $`Y_{3/2}`$ resulting from thermal equilibrium.
However most relativistic gravitinos at the time of BBN result from out-of-equilibrium decays, e.g. from NLSP or from the lightest messenger decay. Given that, immediately after decay the outgoing gravitinos carry a fraction $`B_{3/2}`$ of the rest mass energy of the decaying particle of mass $`M`$ as kinetic energy, their contribution to the energy density at the onset of BBN reads: $`\delta g_{\mathrm{BBN}}/g_{\mathrm{BBN}}^sB_{3/2}MY_M(4/3)(g_{\mathrm{BBN}}^s/g^s)^{1/3}/T_\mathrm{d}`$, where $`T_\mathrm{d}`$ is the decay temperature and $`Y_M`$ the number-to-entropy ratio of the parent at decay. Here it has been assumed that $`T_\mathrm{d}4MY_M/3`$. In the opposite limit, i.e. in the case of significant entropy production at decay, the above relation becomes $`\delta g_{\mathrm{BBN}}/g_{\mathrm{BBN}}^sB_{3/2}(g_{\mathrm{BBN}}^s/g^s)^{1/3}`$ (assuming $`B_{3/2}1`$), and it will be this limit which results in the strongest BBN constraints.
The time of decay of NLSPs to gravitinos is also strongly constrained by big-bang nucleosynthesis limits on hadronic and electromagnetic energy injection at times $`10^210^8`$sec. These constraints have been examined in Refs. BBN1 ; BBN2 ; BBN4 ; FST04 , while Ref. GGR99 has translated these bounds on the messenger scale $`M_X`$ of GMSB models assuming $`k=1`$ (which is equivalent to setting an upper bound on $`m_{3/2}`$). In the present analysis, we use the latest constraints from hadronic and electromagnetic energy injection from Ref. BBN2 .
In GMSB scenarios, the NLSP is generically a neutralino (mainly bino) or a stau. The former decays predominantly into a photon and a goldstino; the fraction of energy spent with branching ratio $`B_{\mathrm{em}}1`$ and by three body decays into a pair of quarks and goldstino with hadronic branching ratio $`B_{\mathrm{had}}10^3`$; if its decay to $`Z`$ bosons is not suppressed by phase space, i.e. $`(M_{\mathrm{NLSP}}M_Z)/M_Z1`$, the hadronic branching ratio $`B_{\mathrm{had}}0.15`$. For simplicity, we use this latter value, which is conservative with respect to the constraints inferred. Concerning the annihilation cross-section of the bino NLSP, we use the value $`\sigma _{\mathrm{NLSP}}v=10^9\mathrm{GeV}^2(M_{\mathrm{NLSP}}/100\mathrm{GeV})^2`$, which corresponds to the bulk region of minimal supergravity FST04 . We will comment on the dependence of our results on the choices made when discussing the results shown in Fig. 1 below.
A stau NLSP may produce in its decay electromagnetic and hadronic showers. About 100% of the energy is converted to electromagnetically interacting particles. In 70% of all decays, a stau NLSP produces hadrons, but these are mesons whose lifetimes are so short that they do not have time to interact before decaying if they were emitted at times $`10^2`$sec. Hence, we use $`B_{\mathrm{had}}=0.7`$ for stau decay timescales shorter than $`10^2`$sec and $`B_{\mathrm{had}}=10^3`$ for longer decay times. The stau annihilation cross-section is not as model dependent as that of the bino, $`\sigma _{\mathrm{NLSP}}v10^7\mathrm{GeV}^2(M_{\mathrm{NLSP}}/100\mathrm{GeV})^2`$. Given its large annihilation cross-section, the stau has a small relic abundance, and consequently the BBN constraints are comparatively weaker.
Overall big-bang nucleosynthesis constraints apply to a combination of $`Y_{\mathrm{NLSP}}`$ (relic abundance) and decay timescale $`M_{\mathrm{NLSP}}^5m_{3/2}^2`$. Note that for a decay timescale $`\tau 10^3`$sec, interesting modifications to BBN may result BBN4 . We also note that in a very limited part of parameter space of GMSB theories, the NLSP can be a sneutrino GR99 , for which the BBN constraints would be largely reduced KM95 ; FST04 .
Finally, since the mass of the NLSP enters the BBN constraints while the gravitino yield is controled by the gluino mass, it is necessary to schematize the mass spectrum of GMSB scenarios. We do so by assuming a mass ratio $`M_3/M_{\mathrm{NLSP}}6`$ GR99 ; DTW97 and fiducial values $`M_{\mathrm{NLSP}}=150`$GeV and $`M_3=1`$TeV. Where relevant we mention the possible influence of these values on our results.
## IV Messenger couplings to matter and gravitino dark matter
In this section, we investigate the possible solutions for gravitino dark matter for various messenger number violating interactions added to the Lagrangian. As argued in Section II, such interactions are mandatory if no substantial entropy production occurs at temperatures $`M_X`$ (other than due to lightest messenger domination and decay) in order to avoid the cosmological problems that would result from the stability of the lightest messenger. For definiteness, we adopt the notations of Ref. WBook including four component spinors, in the general supergravity Lagrangian.
### IV.1 Renormalizable couplings
#### IV.1.1 Superpotential couplings
Renormalizable couplings, beyond those of the required messenger gauge interactions, may exist, though they are constrained by considerations of flavor changing neutral currents as well as the potential development of charge- and color- breaking minima, among other issues. Dine et al. DNS97 have analysed viable extensions of the minimal GMSB scenario in this direction, introducing couplings of the form
$$Wy_l^iH_D\mathrm{\Phi }_M^l\overline{e}_i+y_q^iH_DQ_i\overline{\mathrm{\Phi }}_M^q,$$
(22)
where $`\mathrm{\Phi }_M^l`$ and $`\overline{\mathrm{\Phi }}_M^q`$ denote lepton- and quark- like messengers. Here $`\mathrm{\Phi }_l`$ denotes an SU(2) messenger doublet, the $`y_i`$’s are Yukawa couplings with family index $`i`$, and $`H_D`$, $`\overline{e}`$, and $`Q_i`$ are standard model down-type Higgs, right-handed lepton, and quark doublet, respectively.
Additional SUSY-breaking mass splittings are generated by these types of interactions via one-loop contributions yielding, for example, negative relative mass contributions of order $`\delta m_{\stackrel{~}{e}}/m_{\stackrel{~}{e}}10^3y_l^2F_S^2/M_X^4`$ to slepton masses. Flavor changing neutral currents may place potentially restrictive limits on such couplings, as due to the experimentally verified weakness of such processes, mass splittings between 1st- and 2nd- generation sleptons (and squarks) are constrained to be smaller than $`(m_{\stackrel{~}{e}_1}^2m_{\stackrel{~}{e}_2}^2)/m_{\stackrel{~}{e}}^210^3`$. Assuming conservatively, $`y_l^10`$ and $`y_l^2=y_l^3=0`$ one may thus infer a limit $`y_l^1(M_X/10^8\text{GeV})`$ on this extra Yukawa coupling.
Interactions induced by the superpotential Eq. (22) also induce the decay of messengers, in particular $`X\stackrel{~}{H}^{}e^+`$ assuming $`X`$ carries the same gauge charges as a $`\stackrel{~}{\nu }_L`$ (see Section II). Hence the decay width $`\mathrm{\Gamma }_X=y^2M_X/8\pi `$. Though the limit on Yukawa couplings as inferred above may be quite severe, entropy production due to delayed messenger decay is absent when terms of Eq. (22) are included into $`W`$, unless the extra Yukawa coupling is extremely small, $`y10^{15}(M_X/10^7\mathrm{GeV})^{3/2}`$, assuming $`SU(5)`$ grand unification.
It is nevertheless instructive to study the influence of such “fast” decay of the lightest messenger on the possiblity of having gravitino dark matter. As mentioned earlier, this case (without entropy production) has been implicitly assumed in previous studies that have drawn upper bounds on the post-inflationary reheating temperature from the upper bound on the gravitino density $`\mathrm{\Omega }_{3/2}<1`$ TMY93 ; CHKL99 (and references therein). In order to provide a point of comparison with this previous litterature, we plot in Fig. 1 the results of the calculation of $`\mathrm{\Omega }_{3/2}`$ in the plane $`T_{\mathrm{RH}}m_{3/2}`$, using the techniques developed in the previous section. We assume “fast” decay with width $`\mathrm{\Gamma }_X10^9M_X`$ corresponding to $`y10^4`$, which ensures that phenomenological constraints are satisfied for all values of $`M_X`$. The results shown are insensitive to the exact value of $`y`$, provided it is not so tiny that substantial entropy production would occur.
The shade (color) coding in this figure and all subsequent figures is as follows: lightest (yellow) corresponds to $`\mathrm{\Omega }_{3/2}<0.01`$ (no gravitino problem but no dark matter), and the increasingly darker (respectively green, red and blue) areas indicate respectively the regions of cold, warm and hot dark matter in which $`0.01<\mathrm{\Omega }_{3/2}<1`$. The area shaded by lines oriented NE-SW at the right of each figure corresponds to the region excluded by BBN constraints on NLSP to gravitino decay. White color indicates $`\mathrm{\Omega }_{3/2}>1`$, i.e., overclosure of the Universe by gravitinos. Finally the area marked with horizontal lines is unphysical as it corresponds to $`F_S>F`$ ($`k>1`$).
We choose $`0.01`$ and $`1`$ as lower and upper bounds respectively to delimit where the gravitino can account for dark matter, eventhough cosmological data restrict this to a much smaller range. However the calculations presented here contain intrinsic uncertainties of factors of a few that were mentioned in the previous sections, and therefore the green, red and blue areas should be understood as indicative of the region in which one can find solutions for gravitino dark matter.
As indicated in the caption of Fig. 1, the left panels correspond to $`M_X=10^5`$GeV while the right panels correspond to $`M_X=10^{10}`$GeV (for which the condition $`k1`$ translates in $`m_{3/2}250`$keV). The upper and lower panels correspond respectively to a stau and a bino NLSP. As anticipated in the previous Section, the BBN constraints on hadronic and electromagnetic energy injection do not apply to the stau at these small values of $`m_{3/2}`$, as a result of the low stau relic abundance. In fact, for a stau of mass $`M_{\mathrm{NLSP}}150`$GeV, a gravitino as heavy as $`m_{3/2}10`$GeV is allowed by BBN FST04 . However, the constraints are quite stringent for the case of the bino NLSP, and result in an upper bound $`m_{3/2}10`$MeV for $`M_{\mathrm{NLSP}}=150`$GeV. At a fixed value of the NLSP relic abundance, the BBN constraints give an upper bound on the decay timescale; hence, the above limit on $`m_{3/2}`$ scales as $`M_{\mathrm{NLSP}}^{5/2}`$, see Eq. (15). The BBN limit on $`m_{3/2}`$ evolves as follows with respect to the bino annihilation cross-section; for reference, we recall that the cross-section used in the calculations reported in Fig. 1 is $`\sigma _{\mathrm{NLSP}}v\sigma _0(M_{\mathrm{NLSP}}/100\mathrm{GeV})^2`$ with $`\sigma _0=10^9`$GeV<sup>-2</sup> and $`M_{\mathrm{NLSP}}=150`$GeV. If $`\sigma _0`$ is decreased by a factor 10 to 100, the upper bound on $`m_{3/2}`$ shifts to $`100`$MeV; if, conversely, the cross-section is increased by a factor 10 to 100, the upper bound on $`m_{3/2}`$ shifts to $`3`$MeV. As mentioned in the previous section, we have implicitly assumed a branching ratio to hadronic decay $`B_{\mathrm{had}}=0.15`$; if the bino is nearly degenerate in mass with $`Z`$, the hadronic decay mode is suppressed, with a value possibly as small as $`B_{\mathrm{had}}10^3`$. In this case, for $`\sigma _0`$ chosen as above, the bound on $`m_{3/2}`$ would be $`100`$MeV, increasing to $`10`$GeV if $`\sigma _0`$ is increased by a factor 100, and remaining constant if $`\sigma _0`$ is decreased by a factor as large as 100. Overall the BBN constraints result in a bound $`m_{3/2}10\mathrm{MeV}1\mathrm{GeV}`$ depending on the bino mass, annihilation cross-section and hadronic branching ratio.
Figure 1 illustrates the so-called light gravitino problem: if $`m_{3/2}1`$keV, gravitinos overclose the Universe and/or disrupt BBN unless the reheating temperature is low, $`T_{\mathrm{RH}}\mathrm{min}[10^8m_{3/2}(M_3/1\mathrm{T}\mathrm{e}\mathrm{V})^2,M_X/10]`$, and gravitinos are not too heavy, $`m_{3/2}10\mathrm{MeV}1`$GeV. Although it is not impossible to achieve such small reheating temperatures, either by low-scale inflation or a late phase of thermal inflation, it is not particularly attractive as it puts further non-trivial requirements on the model and may pose problems for a successful genesis of baryon or lepton asymmetry. Moreover it is necessary to tune the reheating temperature to the gravitino and/or messenger mass, eventhough these quantities derive from sectors of the theory that are a priori unrelated.
The region $`m_{3/2}1`$keV is devoid of constraints on $`T_{\mathrm{RH}}`$ since the gravitino is so light that even at thermal equilibrium, it cannot overclose the Universe. However such light gravitinos make up dark matter that is too warm to reproduce existing data on the large scale structures. Heavier gravitinos $`m_{3/2}10`$MeV are excluded by big-bang nucleosynthesis constraints if the NLSP is a bino. Since NLSPs decay at time $`\tau _{\mathrm{NLSP}}6\times 10^4\mathrm{sec}(M_{\mathrm{NLSP}}/100\mathrm{G}\mathrm{e}\mathrm{V})^5(m_{3/2}/1\mathrm{G}\mathrm{e}\mathrm{V})^2`$, the heavier the gravitino the later and the more constrained the decay. Figure 1 differs from those shown in Refs. TMY93 ; CHKL99 because we have included constraints from GMSB phenomenology (namely $`k1`$) as well as updated constraints from BBN and structure formation, and a more accurate calculation of the gravitino relic abundance. Overall, one finds that the range of allowed $`m_{3/2}`$ is severely restricted when compared to previous studies.
In particular, the present conclusion is at variance with Ref. CHKL99 which argued that for a sufficiently small messenger mass scale $`M_X10^5`$GeV and sufficiently large gravitino mass $`m_{3/2}2`$GeV, it is possible to find solutions to the gravitino problem for arbitrarily high reheating temperatures. The discrepancy with CHKL99 is tied to the neglect in that study of the SUGRA induced MSSM particles contribution to the gravitino abundance at large reheating temperatures, as discussed in Section III.A., as well as of the big-bang nucleosynthesis constraints on NLSP decay. As can be seen in the left panels of Fig. 1, the upper bound on $`T_{\mathrm{RH}}`$ does indeed shift upwards, albeit for larger $`m_{3/2}`$ than expected in CHKL99 , i.e. $`m_{3/2}10`$GeV for $`M_X10^5`$GeV. This $`m_{3/2}`$ region is however forbidden by BBN constraints on energy injection, for both stau and bino NLSP. If the gluino mass scale is smaller than the fiducial value of $`1`$TeV, say $`M_3=300`$GeV, this region of parameter space where the bound on $`T_{\mathrm{RH}}`$ is relaxed, moves to smaller $`m_{3/2}`$, i.e. $`m_{3/2}0.41`$GeV. However, the BBN constraints also move to smaller $`m_{3/2}`$ since the NLSP mass is reduced as the gluino mass. Finally, since the MSSM particle contribution to $`\mathrm{\Omega }_{3/2}`$ contains a dependence on $`m_{3/2}`$ and $`T_{\mathrm{RH}}`$, we find that $`T_{\mathrm{RH}}`$ is always bounded from above due to the BBN bound on $`m_{3/2}`$; for instance, for $`M_3=300`$GeV and for a stau NLSP, the maximal reheating temperature where $`\mathrm{\Omega }_{3/2}<1`$, is $`10^9`$GeV. In the end it turns out that some fine-tuning betwen $`M_3`$, $`M_{\mathrm{NLSP}}`$ and $`m_{3/2}`$ is required to find a region in which $`T_{\mathrm{RH}}`$ can become as large as $`10^9`$GeV.
There exist other potential renormalizable interaction terms that violate messenger number by one unit. (Such operators lead to $`1/m_{\mathrm{Pl}}^2`$ suppressed proton decay DGP96 , due to gauge coupling of the messenger fields.) In the following we assume for definiteness that messengers come in complete representations of $`SU(5)`$, in particular, as $`\mathrm{𝟓}_M+\overline{\mathrm{𝟓}}_M`$ (or $`\mathrm{𝟏𝟎}_M+\overline{\mathrm{𝟏𝟎}}_M`$), while the visible sector superfields are denoted by $`\overline{\mathrm{𝟓}}_F+\mathrm{𝟏𝟎}_F`$, and the Higgses sit in one pair of $`\mathrm{𝟓}_H+\overline{\mathrm{𝟓}}_H`$ and one $`\mathrm{𝟐𝟒}_H`$ supermultiplets. Gauge symmetry limits those interaction terms between messengers and visible sectors superfields in the superpotential to the following:
$`W_{\mathrm{ren}}`$ $`\{`$ $`\overline{\mathrm{𝟓}}_M\overline{\mathrm{𝟓}}_{F,H}\mathrm{𝟏𝟎}_F,\mathbf{\hspace{0.17em}5}_M\mathrm{𝟏𝟎}_F\mathrm{𝟏𝟎}_F,\mathbf{\hspace{0.17em}5}_M\overline{\mathrm{𝟓}}_{F,H}\mathrm{𝟐𝟒}_H,`$ (23)
$`\overline{\mathrm{𝟓}}_M\mathrm{𝟓}_H\mathrm{𝟐𝟒}_H,\overline{\mathrm{𝟏𝟎}}_M\mathrm{𝟓}_H\mathrm{𝟓}_H,\mathbf{\hspace{0.17em}10}_M\overline{\mathrm{𝟓}}_{H,F}\overline{\mathrm{𝟓}}_{H,F},`$
$`\mathrm{𝟏𝟎}_M\mathrm{𝟏𝟎}_F\mathrm{𝟓}_H,\mathbf{\hspace{0.17em}10}_F\overline{\mathrm{𝟏𝟎}}_M\mathrm{𝟐𝟒}_H\}.`$
The interaction terms considered in Ref. DNS97 are contained in $`\overline{\mathrm{𝟓}}_M\overline{\mathrm{𝟓}}_H\mathrm{𝟏𝟎}_F`$ and $`\mathrm{𝟏𝟎}_M\mathrm{𝟏𝟎}_F\mathrm{𝟓}_H`$, and lead to fast messenger decay. Slow decay may, in principle, result from operators in Eq. (23) which involve particles with GUT scale masses. However, renormalizable couplings to a $`\mathrm{𝟐𝟒}_H`$ must be excluded as they lead to $`M_{\mathrm{GUT}}`$ mass mixings between messengers and visible sectors particles, hence they would spoil the phenomenology at the electroweak scale. The colored triplet GUT Higgs in $`\mathrm{𝟓}_H`$ and $`\overline{\mathrm{𝟓}}_H`$ carry masses of order $`M_{\mathrm{GUT}}`$ (but no vev), but they do not couple to the lightest messenger, which is either the $`5^{\mathrm{th}}`$ component of a $`\mathrm{𝟓}_M`$ or the $`(4,5)`$ component of a $`\mathrm{𝟏𝟎}_M`$. Hence the colored Higgses do not lead to suppression of the decay width. In Ref. BM03 it was proposed that delayed messenger decay could occur if the decay of the lightest messenger was suppressed by the mediation of a particle of mass $`10^{12}`$GeV in a renormalizable interaction. Our present discussion shows that this model is not natural in the sense that it requires a new particle with both the required mass $`10^{12}`$GeV and gauge charges such that the required renormalizable interaction could occur. In the above list of possible interactions terms, such coupling does not appear for the minimal content of $`SU(5)`$.
One may also argue that the required fine-tuning to avoid flavor changing neutral currents (in particular for light messengers) may actually indicate the absence of those renormalizable interactions, and that messenger number violation occurs via further suppressed interactions. Such couplings will be discussed further below.
#### IV.1.2 Renormalizable couplings in the Kähler function
Fujii & Yanagida FY02 have proposed messenger-matter mixing due to a correction in the superpotential $`\delta W(W/m_{\mathrm{Pl}}^2)\mathrm{𝟓}_M\overline{\mathrm{𝟓}}_F`$, where $`Wm_{3/2}m_{\mathrm{Pl}}^2`$. In the framework of supergravity, a possible origin for such a superpotential term can be highlighted by adding to the minimal Kähler potential $`K_0=_i\mathrm{\Phi }_i^{}\mathrm{\Phi }_i`$, ($`i`$ running over all superfields $`\mathrm{\Phi }_i`$) a non-minimal part $`\delta K`$ given by,
$$\delta K=\mathrm{𝟓}_M\overline{\mathrm{𝟓}}_F+h.c.$$
(24)
$`\delta K`$ is allowed by gauge symmetries (and possibly by an R-symmetry as well, for conveniently chosen R charges). Then, making use of the usual invariance of the supergravity Lagrangian under Kähler transformations $`KK+F(\mathrm{\Phi })+F^{}(\mathrm{\Phi }^{})`$, followed by superpotential $`We^FW`$ (and super-Weyl) scalings, the above $`\delta W`$ is obtained for $`F(\mathrm{\Phi })=\delta K`$ to the lowest order in $`1/m_{\mathrm{Pl}}^2`$, provided that a constant is added to the superpotential to fine-tune the cosmological constant after supersymmetry breaking (whence $`Wm_{3/2}m_{\mathrm{Pl}}^2`$)<sup>5</sup><sup>5</sup>5Altogether this is very reminiscent of the Giudice-Masiero mechanism which provides a solution to the so-called $`\mu `$-problem MG . . In our notations, the lightest messenger $`\varphi `$ is a linear combination of the lightest scalar components of $`\mathrm{𝟓}_M`$ and $`\overline{\mathrm{𝟓}}_M^{}`$, see Section 2.1. The mixing between $`\mathrm{𝟓}_M`$ and $`\overline{\mathrm{𝟓}}_F`$ generated from $`\delta Wm_{3/2}\mathrm{𝟓}_M\overline{\mathrm{𝟓}}_F`$ thus leads to the decay of $`X`$ into a lepton and a gaugino with width $`\mathrm{\Gamma }_{Xl\lambda }(g^2/16\pi )m_{3/2}^2/M_X`$ FY02 .
We should stress here that, starting as we do from $`\delta K`$ rather than $`\delta W`$ of FY02 , one expects further contributions to the messenger decay or annihilation, with possibly important effects on the final gravitino abundance. Indeed, other contributions to the decay into visible sector particles originate from the supergravity scalar potential WBook
$`V_B=e^{K/m_{\mathrm{Pl}}^2}`$ $`[`$ $`K^{ij^{}}\left(W{\displaystyle \frac{K_i}{m_{\mathrm{Pl}}^2}}+W_i\right)\left(W^{}{\displaystyle \frac{K_j^{}}{m_{\mathrm{Pl}}^2}}+W_j^{}^{}\right)`$ (25)
$`{\displaystyle \frac{3WW^{}}{m_{\mathrm{Pl}}^2}}]`$
where, $`i,j^{}`$ label the full set of scalar fields $`\varphi _i,\varphi _j^{}^{}`$; $`K^{ij^{}}`$ denotes the inverse of the matrix $`K/\varphi _i\varphi _j^{}^{}`$, and $`W_i=W/\varphi _i`$, $`W_j^{}^{}=W^{}/\varphi _j^{}^{}`$. From $`KK_0+\delta K`$ and taking for illustration the case of $`\mathrm{𝟓}+\overline{\mathrm{𝟓}}`$ messengers with $`WS\mathrm{𝟓}_M\overline{\mathrm{𝟓}}_M+y\overline{\mathrm{𝟓}}_F\overline{\mathrm{𝟓}}_H\mathrm{𝟏𝟎}_F+W`$ one finds the leading contributions to the potential which induce the decay of the lightest messenger,
$`V_B`$ $``$ $`m_{3/2}S^{}\overline{\mathrm{𝟓}}_F\overline{\mathrm{𝟓}}_M^{}+ym_{3/2}\mathrm{𝟓}_M\overline{\mathrm{𝟓}}_H^{}10_F^{}2m_{3/2}^2\mathrm{𝟓}_M\overline{\mathrm{𝟓}}_F`$ (26)
$`+{\displaystyle \frac{1}{m_{\mathrm{Pl}}^2}}\{\mathrm{𝟓}_M\overline{\mathrm{𝟓}}_F[|S|^2(\mathrm{𝟓}_M\mathrm{𝟓}_M^{}+\overline{\mathrm{𝟓}}_M\overline{\mathrm{𝟓}}_M^{})`$
$`+\mathrm{𝟓}_M\overline{\mathrm{𝟓}}_M|^2+y^2|W_{y|i}|^2]`$
$`+(S^{}\overline{\mathrm{𝟓}}_M^{}\overline{\mathrm{𝟓}}_F+y\mathrm{𝟓}_M\overline{\mathrm{𝟓}}_H^{}\mathrm{𝟏𝟎}_F^{})(S\mathrm{𝟓}_M\overline{\mathrm{𝟓}}_M+y\overline{\mathrm{𝟓}}_F\overline{\mathrm{𝟓}}_H\mathrm{𝟏𝟎}_F)\}`$
$`+h.c.`$
where we have neglected terms suppressed by higher powers of $`1/m_{\mathrm{Pl}}^2`$ or of order $`m_{3/2}/m_{\mathrm{Pl}}^2`$ and smaller. Other operators which do not induce messenger decay into standard model particles are not shown, being irrelevant for the present discussion. After the spurion scalar field $`S`$ has developed a supersymmetric vev, the first term of Eq. (26) leads to the bilinear operator $`m_{3/2}M_X\overline{\mathrm{𝟓}}_F\overline{\mathrm{𝟓}}_M^{}`$ contributing to the decay considered in FY02 (actually a similar mixing between the fermionic partners is also generated, see below). Note that the second operator in Eq. (26) can mediate an equally efficient decay of the lightest messenger to a (colorless) Higgs and a slepton $`\mathrm{\Gamma }_{XH\stackrel{~}{f}}(y^2/32\pi )m_{3/2}^2/M_X`$, provided that the Yukawa coupling is large (i.e. if $`\delta K`$ involves the third generation). A similar decay to a Higgsino and a standard model fermion occurs as well, triggered by the MSSM Yukawa vertex and $`m_{3/2}M_X\overline{\mathrm{𝟓}}_F\overline{\mathrm{𝟓}}_M^{}`$. In any case, these new contributions do not lead to a significant change in the analysis of the gravitino relic density, being of the same order as $`\mathrm{\Gamma }_{Xl\lambda }`$.
Other contributions in Eq. (26) can potentially lead to important modifications when loop effects are considered. This is due on one hand to the supersymmetry preserving non-renormalizable operators of the form $`\mathrm{𝟓}_M\mathrm{}\varphi \varphi ^{}\mathrm{}\overline{\mathrm{𝟓}}_F/m_{\mathrm{Pl}}^2`$, with $`\varphi `$ denoting an arbritrary field, and on the other hand to renormalizable operators which induce hard supersymmetry-breaking after cancellation of the cosmological constant. (In connection with the latter operators, the presence of the spurion, a visible sector singlet, could destabilize the hierarchy of the messenger and/or electroweak scales JBPR .) A thorough discussion of these issues which are relevant for the theoretical consistency of the effective supergravity model is out of the scope of the present paper. Hereafter we give only a partial and brief discussion of the two types of operators.
After the vev shift $`SS+M_X`$, the term of order $`m_{\mathrm{Pl}}^2`$ in Eq. (26) gives the operators
$`\left(\frac{M_X^2}{m_{\mathrm{Pl}}^2}\right)\times (\mathrm{𝟓}_M\overline{\mathrm{𝟓}}_M\overline{\mathrm{𝟓}}_M^{}\overline{\mathrm{𝟓}}_F,\mathrm{𝟓}_M\overline{\mathrm{𝟓}}_F\overline{\mathrm{𝟓}}_M\overline{\mathrm{𝟓}}_M^{},\mathrm{𝟓}_M\overline{\mathrm{𝟓}}_F\mathrm{𝟓}_M\mathrm{𝟓}_M^{})`$. These operators generate potentially very large corrections to the $`\mathrm{𝟓}_M\overline{\mathrm{𝟓}}_F`$ mixing, through one-loop contributions of the $`\mathrm{𝟓}_M`$ and $`\overline{\mathrm{𝟓}}_M`$, which are of the order of $`\mathrm{\Lambda }^2m_{\mathrm{Pl}}^2`$, where $`\mathrm{\Lambda }`$ is the cut-off scale of the physics underlying the effective supergravity Lagrangian. Even more, the operators containing $`SS^{}`$ instead of $`M_X^2`$ lead to more dangerous corrections of order $`\mathrm{\Lambda }^4`$ due to two-loop diagrams. The same diagrams can also induce direct decays of the lightest messenger into MSSM particles, such as $`X\gamma (Z)+\stackrel{~}{l}`$. It is thus important to assess the supersymmetric cancellations which would keep those contributions under control, eventhough they originate from the gravitational non-renormalizable sector. The companion operators involving the messenger fermions are contained in the $`O(1/m_{\mathrm{Pl}}^2)`$ part of
$`_F`$ $``$ $`e^{K/2m_{\mathrm{Pl}}^2}[{\displaystyle \frac{1}{2}}𝒟_iD_jW\overline{\chi }_R^i\chi _L^j+h.c.]`$ (27)
where $`𝒟_iD_jW=W_{ij}+(K_{ij}/m_{\mathrm{Pl}}^2)W+(K_i/m_{\mathrm{Pl}}^2)D_jW+(K_j/m_{\mathrm{Pl}}^2)D_iW(K_iK_j/m_{\mathrm{Pl}}^4)W\mathrm{\Gamma }_{ij}^kD_kW`$
(see WBook ). They read
$`_F{\displaystyle \frac{1}{2m_{\mathrm{Pl}}^2}}`$ $`[`$ $`\overline{\psi }_R\psi _LS(\mathrm{𝟓}_M\overline{\mathrm{𝟓}}_F+h.c.)`$ (28)
$`+\overline{\psi }_R\mathrm{𝟓}_M\psi _L\overline{\mathrm{𝟓}}_FS]+h.c.`$
where $`\psi `$ denotes the Dirac field which combines the two fermionic components of the $`\mathrm{𝟓}+\overline{\mathrm{𝟓}}`$ messenger superfields. \[The order of occurrence of the fields in Eq. (28) indicates how they are combined into $`SU(5)`$ invariants .\] In the supersymmetric limit, the one-loop contributions to the $`5_M\overline{5}_F`$ mixing induced by $`_F`$ with $`SS`$ cancel exactly the ones from the scalar loops discussed above. After SUSY breaking through the vev of the auxiliary field $`F_S`$, no dependence on the ultra-violet cut-off $`\mathrm{\Lambda }`$ is reintroduced. In particular, even $`\mathrm{log}\mathrm{\Lambda }`$ terms cancel out (though they would not have altered the size of the mixing) yielding a correction of order $`F_S^2/m_{\mathrm{Pl}}^2`$ in the limit $`F_SM_X^2`$, and of order $`M_X^4/m_{\mathrm{Pl}}^2`$ in the limit $`F_SM_X^2`$, which remains negligible when compared to the tree-level mixing magnitude $`m_{3/2}M_X`$, Eq. (26), in the parameter space region relevant for gravitino dark matter.<sup>6</sup><sup>6</sup>6We checked also for cancellations to two-loop order considering subclasses of Feynman diagrams which involve the spurion scalar and fermion virtual contributions. A detailed study is outside the scope of the present paper. Finally, as mentioned before, some hard SUSY breaking operators are generated in Eq. (27) from $`m_{3/2}e^{K/2m_{\mathrm{Pl}}^2}(\frac{1}{2}\overline{\chi }_R^i\chi _L^j\delta K_{ij}+h.c.)`$, leading to bilinear matter fermion mixing between the messengers and the MSSM particles, $`m_{3/2}(\overline{f}_R\psi _L+h.c.)`$. These can potentially lead to quadratic divergences which would destabilise the scale of the mixing $`\mathrm{𝟓}_M\overline{\mathrm{𝟓}}_F`$ between the lightest messenger and the MSSM scalar particles. However, leading one-loop (tadpole) effects with one mass insertion occur as corrections to $`\overline{\mathrm{𝟓}}_M^{}\overline{\mathrm{𝟓}}_F`$ and turn out to be at worse $`O((\mathrm{\Lambda }^2/m_{\mathrm{Pl}}^2)m_{3/2}M_X)`$, thus harmless for a cut-off of order $`m_{\mathrm{Pl}}`$. (Note also that the matter fermion kinetic terms induce fermionic tadpole contributions which cancel among themselves due to the derivative couplings.) To summarize, in the sector of the supergravity Lagrangian not involving the gravitino the $`\delta K`$ piece of the Kähler potential leads essentially to the same features as discussed in FY02 , at least to leading order in $`1/m_{\mathrm{Pl}}^2`$ and up to one-loop. One exception is the messenger decay induced by first term in Eq. (26) which we will consider later on.
The relic abundance for the gravitino can be calculated using the techniques developed in Section III and the results are shown in Fig. 2. This figure uses the color shading as in Fig. 1 but is plotted in the plane $`m_{3/2}M_X`$ instead of $`T_{\mathrm{RH}}`$, which was taken to be $`T_{\mathrm{RH}}=10^{12}`$GeV. The plot shown in the left panel assumes that the spurion $`S`$ is much heavier than the lightest messenger $`X`$, in which case annihilation $`XX^{}\stackrel{~}{G}\stackrel{~}{G}`$ takes place with the cross-section given in Eq. (7). As discussed in Section II.B, this cross-section violates unitarity for $`M_X10^7\mathrm{GeV}(m_{3/2}/1\mathrm{keV})^{2/3}`$, i.e. in the region above the dashed thick line in the left panel of Fig. 2. There is no solution in this case for gravitino dark matter, at variance with the conclusions of Ref. FY02 . The main reason is that annihilation into goldstinos has not been accounted for in Ref. FY02 , yet the solution for gravitino dark matter proposed by these authors lies in the region in which unitarity is violated. The annihilation cross-section at its unitarity bound is much larger than that used in Ref. FY02 for annihilation through gauge interactions, hence the messenger relic abundance and the amount of gravitino dilution are correspondingly smaller. At the very least, since one cannot predict the cross-section in this region where multi-goldstino production violates unitarity, one can conclude that the results for the scenario of Ref. FY02 in this region are model dependent in that they require contributions from the hidden sector to bring down multi-goldstino production to a negligible level.
In the right panel of Fig. 2, it is assumed that $`S`$ is much lighter than $`X`$. Annihilation into goldstinos scales with the effective Yukawa coupling $`F_S/F`$ of the messenger to the goldstino component of the fermionic partner of $`S`$, as in Eq. (6). This annihilation channel thus contributes only in the region of direct GMSB scenarios where $`kF_S/F1`$, which happens to be that where gravitino dark matter can be found for the mixing term $`\mathrm{𝟓}_M\overline{\mathrm{𝟓}}_F`$ proposed in Ref. FY02 . The lightest messenger can also annihilate into a pair of spurions, as discussed in Section II.B. Furthermore, one must also consider the possible decay of the lightest messenger into one visible sector sneutrino and two goldstinos, which is induced by the above mixing term and $`XX^{}\stackrel{~}{G}\stackrel{~}{G}`$ four-point vertices similar to the ones discussed after Eq. (3). For instance, the latter vertex is induced by terms of the form $`e^{K/2m_{\mathrm{Pl}}^2}W\overline{\mathrm{\Psi }}_L^\mu \sigma _{\mu \nu }\mathrm{\Psi }_R^\nu /m_{\mathrm{Pl}}^2`$, WBook , using Eqs. (24, 4) and $`WW`$.
One expects this decay width to scale as $`\mathrm{\Gamma }_{X\stackrel{~}{\nu }\stackrel{~}{G}\stackrel{~}{G}}0.1(F_S/F)^4\mathrm{\Gamma }_{Xl\lambda }`$, with the prefactor of $`0.1`$ accounting for the enlargement of phase space. This decay channel produces highly relativistic gravitinos which mix with the “cold” gravitinos produced by other channels, see Section III.A, resulting in mixed dark matter. The small area around $`m_{3/2}10100`$keV is shown in red (medium shading) in Fig. 2, indicating that the hot gravitinos contribute to more than 10% of the gravitino energy density and that the averaged velocity corresponds to warm dark matter. Finally one must also include decay into a spurion and a sneutrino induced by the three scalar coupling, first term in Eq. (26), with width $`\mathrm{\Gamma }_{XS\stackrel{~}{\nu }}(1/16\pi )m_{3/2}^2/M_X`$.
The inclusion of these various annihilation and decay channels modify the region of parameter space where one can find gravitino dark matter with respect to the conclusions of Ref. FY02 . The various effects add up in pushing this region to higher values of $`M_X`$. In effect, the increased lightest messenger annihilation cross-section decreases their relic abundance hence the amount of entropy production; similarly, the increased lightest messenger decay width increases their decay temperature hence also decreases the entropy production. These effects can be compensated, at a given value of $`m_{3/2}`$, by increasing $`M_X`$ since $`\sigma _{XX^{}\mathrm{}}v1/M_X^2`$ and $`\mathrm{\Gamma }_X\mathrm{}1/M_X`$.
Note that Fig. 2 assumes that the NLSP is a bino; as seen previously, the BBN bounds would be relaxed is the NLSP turned out to be a stau.
Finally, the influence of the choice of the post-inflationary reheating temperature (taken as $`T_{\mathrm{RH}}=10^{12}`$GeV here) is as follows. For this chosen value of $`T_{\mathrm{RH}}`$, gravitinos (i.e. their 1/2 component) abundances are brought to equilibrium for all $`m_{3/2}100`$MeV. For smaller values of $`T_{\mathrm{RH}}`$, the plot would thus look similar in the region where $`M_XT_{\mathrm{RH}}`$, unless $`T_{\mathrm{RH}}`$ is that low $`T_{\mathrm{RH}}10^{12}\mathrm{GeV}(m_{3/2}/0.1\mathrm{GeV})^2`$ that gravitinos are initially not in chemical equilibrium. For $`M_XT_{\mathrm{RH}}`$, messengers are not produced at the post-inflationary reheating, there is no entropy production (gravitino dilution) and the gravitino abundance in the plane $`m_{3/2}T_{\mathrm{RH}}`$ then resembles that shown in the right panels of Fig. 1.
If the messengers sit in $`\mathrm{𝟏𝟎}+\overline{\mathrm{𝟏𝟎}}`$ representations of $`SU(5)`$, a similar mixing term $`\overline{\mathrm{𝟏𝟎}}_M\mathrm{𝟏𝟎}_F`$ can be induced, and leads to similar effects. Note that the mixing $`\mathrm{𝟓}_M\overline{\mathrm{𝟓}}_F`$ and $`\overline{\mathrm{𝟏𝟎}}_M\mathrm{𝟏𝟎}_F`$ are constrained by $`R`$symmetry charge assignments. If a non-holomorphic mixing $`\overline{\mathrm{𝟓}}_M^{}\overline{\mathrm{𝟓}}_F`$ (or $`\overline{\mathrm{𝟏𝟎}}_M^{}\overline{\mathrm{𝟏𝟎}}_F`$) is allowed, the situation is quite different. In effect, the lightest messenger now mixes in the kinetic terms with the sneutrino and its decay is no longer suppressed by $`m_{3/2}/M_X`$. As a consequence, decay is much faster, $`\mathrm{\Gamma }𝒪(M_X)`$, and no entropy production occurs; the situation is then as shown in Fig. 1.
### IV.2 Non-renormalizable couplings
#### IV.2.1 Superpotential interactions
It is possible that interactions of the type Eq. (22) are forbidden and that messengers may only decay via non-renormalizable operators in $`W`$, as discussed previously. Operators which may arise due to unknown Planck-scale physics, which respect the $`SU(5)`$ gauge symmetry and that violate the messenger number by one unit are given to leading order in $`1/m_{\mathrm{Pl}}`$, by:
$`W_{\mathrm{non}\mathrm{ren}}`$ $`{\displaystyle \frac{1}{m_{\mathrm{Pl}}}}\{`$ $`\overline{\mathrm{𝟓}}_M\mathrm{𝟏𝟎}_F\mathrm{𝟏𝟎}_F\mathrm{𝟏𝟎}_F,\mathbf{\hspace{0.17em}5}_M\mathrm{𝟓}_H\overline{\mathrm{𝟓}}_{H,F}\overline{\mathrm{𝟓}}_{H,F},`$ (29)
$`\overline{\mathrm{𝟓}}_M\mathrm{𝟓}_H\mathrm{𝟓}_H\overline{\mathrm{𝟓}}_{H,F},\mathbf{\hspace{0.17em}5}_M\mathrm{𝟓}_H\mathrm{𝟓}_H\mathrm{𝟏𝟎}_F,`$
$`\overline{\mathrm{𝟓}}_M\overline{\mathrm{𝟓}}_H\mathrm{𝟏𝟎}_F\mathrm{𝟐𝟒}_H,\mathbf{\hspace{0.17em}5}_M\overline{\mathrm{𝟓}}_{H,F}\mathrm{𝟐𝟒}_H\mathrm{𝟐𝟒}_H,`$
$`\overline{\mathrm{𝟓}}_M\mathrm{𝟓}_H\mathrm{𝟐𝟒}_H\mathrm{𝟐𝟒}_H,\mathbf{\hspace{0.17em}10}_F\overline{\mathrm{𝟏𝟎}}_M\mathrm{𝟓}_H\overline{\mathrm{𝟓}}_{H,F},`$
$`\overline{\mathrm{𝟏𝟎}}_M\overline{\mathrm{𝟓}}_{H,F}\overline{\mathrm{𝟓}}_{H,F}\overline{\mathrm{𝟓}}_{H,F},\mathbf{\hspace{0.17em}10}_M\mathrm{𝟓}_H\mathrm{𝟓}_H\mathrm{𝟓}_H,`$
$`\mathrm{𝟏𝟎}_M\overline{\mathrm{𝟓}}_{H,F}\mathrm{𝟏𝟎}_F\mathrm{𝟏𝟎}_F,\mathbf{\hspace{0.17em}10}_M\mathrm{𝟏𝟎}_F\mathrm{𝟓}_H\mathrm{𝟐𝟒}_H,`$
$`\overline{\mathrm{𝟏𝟎}}_M\mathrm{𝟓}_H\mathrm{𝟓}_H\mathrm{𝟐𝟒}_H,\mathbf{\hspace{0.17em}10}_M\overline{\mathrm{𝟓}}_{H,F}\overline{\mathrm{𝟓}}_{H,F}\mathrm{𝟐𝟒}_H,`$
$`\overline{\mathrm{𝟏𝟎}}_M\mathrm{𝟏𝟎}_F\mathrm{𝟐𝟒}_H\mathrm{𝟐𝟒}_H,\mathbf{\hspace{0.17em}5}_M\overline{\mathrm{𝟓}}_M\mathrm{𝟓}_M\overline{\mathrm{𝟓}}_F,`$
$`\mathrm{𝟏𝟎}_M\overline{\mathrm{𝟏𝟎}}_M\mathrm{𝟏𝟎}_M\mathrm{𝟏𝟎}_F\}`$
All terms in Eq. (29) which involve couplings of one lightest messenger to $`\overline{\mathrm{𝟓}}_F`$ or $`\mathrm{𝟏𝟎}_F`$ but not Higgses lead to decay into three-body final states with decay width $`10^4M_X^3/m_{\mathrm{Pl}}^2`$ . It is easy to see, using Eq. (19) that entropy production is not sufficient to dilute the gravitinos to the required abundance for a high post-inflationary reheating temperature $`T_{\mathrm{RH}}10^8`$GeV (and $`T_{\mathrm{RH}}M_X`$). Admittedly this is a drawback of the present scenario since those terms are the most generic.
Let us now consider the terms involving couplings to Higgses. For terms involving one $`\mathrm{𝟐𝟒}_H`$ acquiring vevs of order $`M_{\mathrm{GUT}}`$, say $`X\mathrm{\Phi }_1\mathrm{\Phi }_2H_{\mathrm{𝟐}4}/m_{\mathrm{Pl}}`$, the Lagrangian contains the effective Yukawa interaction $`(M_{\mathrm{GUT}}/m_{\mathrm{Pl}})X\overline{\chi }_1\chi _2`$ from the fermionic part of the Lagrangian contained in Eq. (27), with $`\chi _1`$ and $`\chi _2`$ the Weyl spinors of $`\mathrm{\Phi }_1`$ and $`\mathrm{\Phi }_2`$. This Yukawa interaction between the lightest messenger and two fermions with effective coupling constant $`10^3`$ leads to fast decay if $`\chi _1`$ and $`\chi _2`$ have electroweak scale masses. Other terms in the Lagrangian lead to similar or only somewhat smaller partial widths. Consequently the conclusions of the previous section with regards to gravitino dark matter with ”fast” decaying messengers apply. From Eq. (29) one can check that all possible above combinations involving one $`\mathrm{𝟐𝟒}_H`$ contain a coupling of $`X`$ to particles with electroweak masses except $`\mathrm{𝟏𝟎}_M\mathrm{𝟏𝟎}_F\mathrm{𝟓}_H\mathrm{𝟐𝟒}_H`$. However, even for the latter term, the scalar potential contains the interaction $`\mathrm{𝟏𝟎}_M\mathrm{𝟏𝟎}_F\mathrm{𝟐𝟒}_Hy\mathrm{𝟏𝟎}_F^{}\mathrm{𝟏𝟎}_F^{}/m_{\mathrm{Pl}}`$ generated by $`|W/\mathrm{𝟓}_H|^2`$, with $`y`$ the third family Yukawa coupling. This interaction gives a decay width $`\mathrm{\Gamma }10^5y^2(M_{\mathrm{GUT}}/m_{\mathrm{Pl}})^2M_X`$, which is too large to allow solutions for gravitino dark matter.
Terms involving two $`\mathrm{𝟐𝟒}`$ should be excluded as they lead to unacceptable mass mixings between messengers and visible sectors superfields.
Consider now terms of the form $`X\mathrm{\Phi }_1\mathrm{\Phi }_2\mathrm{\Phi }_3/m_{\mathrm{Pl}}`$ containing at least one $`\mathrm{𝟓}_H`$ but no $`\mathrm{𝟐𝟒}_H`$. It can be checked that all terms of this form in Eq. (29) contain couplings of $`X`$ to particles with electroweak masses, hence lead to fast decay as above, except for $`\mathrm{𝟓}_M\mathrm{𝟓}_H\mathrm{𝟓}_H\mathrm{𝟏𝟎}_F/m_{\mathrm{Pl}}`$, $`\mathrm{𝟏𝟎}_M\mathrm{𝟓}_H\mathrm{𝟓}_H\mathrm{𝟓}_H/m_{\mathrm{Pl}}`$ and $`\overline{\mathrm{𝟏𝟎}}_M\overline{\mathrm{𝟓}}_{H,F}\overline{\mathrm{𝟓}}_{H,F}\overline{\mathrm{𝟓}}_{H,F}/m_{\mathrm{Pl}}`$. The first of these terms, when written for the lightest messenger, contains at least one colored Higgs, say $`H_i`$. However, the scalar potential term $`|W/\mathrm{𝟓}_{H_i}|^2`$ generates here as well an interaction between the lightest messenger with 4 particles of electroweak masses, leading to decay width $`\mathrm{\Gamma }10^5y^2M_X^3/m_{\mathrm{Pl}}^2`$, still too large for dark matter solutions. The term $`\mathrm{𝟏𝟎}_M\mathrm{𝟓}_H\mathrm{𝟓}_H\mathrm{𝟓}_H/m_{\mathrm{Pl}}`$ couples $`X`$ to the three colored Higgses, hence its decay is too highly suppressed both by the GUT scale and phase space, $`\mathrm{\Gamma }10^{12}(M_X/M_{\mathrm{GUT}})^8M_X^3/m_{\mathrm{Pl}}^2`$, and cannot lead to decay before BBN for $`M_X10^{14}`$GeV. Note that the big-bang nucleosynthesis constraints on NLSP decay can be translated into an extreme upper bound $`M_X10^{12}`$GeV (see Fig. 2 and GGR99 ). Finally, for the term $`\overline{\mathrm{𝟏𝟎}}_M\overline{\mathrm{𝟓}}_{H,F}\overline{\mathrm{𝟓}}_{H,F}\overline{\mathrm{𝟓}}_{H,F}/m_{\mathrm{Pl}}`$, similar conclusions apply if three Higgses are involved. For couplings involving two colored Higgses, decay is also too highly suppressed. In effect, the lightest messenger can then decay into 5-body final state with mediation by a GUT mass Higgs, leading to $`\mathrm{\Gamma }10^6y^2(M_X/M_{\mathrm{GUT}})^4M_X^3/m_{\mathrm{Pl}}^2`$. It cannot decay before BBN if $`M_X10^{12}`$GeV, as can be checked using Eq. (16). Finally, if only one Higgs is involved in the coupling, the messenger can decay into four particle final states with decay width $`\mathrm{\Gamma }10^5y^2M_X^3/m_{\mathrm{Pl}}^2`$ and yet provide no solution for gravitino dark matter.
The last two operators of Eq. (29) together with the first term of Eq (1) induce at one-loop order a mixing between the lightest messenger and the MSSM matter fields. This mixing is non-vanishing only after supersymmetry breaking and is found to be of magnitude $`(\sqrt{2}M_X/m_{\mathrm{Pl}})F_S\mathrm{log}(\mathrm{\Lambda }/M_X)`$. As discussed in Section IV.A.2, $`\mathrm{\Lambda }`$ is identified with a physical cut-off of order $`m_{\mathrm{Pl}}`$, however, in contrast with the results of that section, there is here no (accidental) cancellation of the cut-off dependence and no suppression by the gravitino mass. The mixing is so large that decay is not accompanied by entropy production, and consequently there is no solution for gravitino dark matter.
In summary, non-renormalizable interaction terms in the superpotential for the minimal content of GMSB scenarios in $`SU(5)`$ grand unification do not allow for natural solutions leading to gravitino dark matter (unless the post-inflationary reheating temperature is tuned as before).
#### IV.2.2 Non-renormalizable interactions in the Kähler function
Interaction terms in the Kähler potential $`K`$ can either be holomorphic or not in the superfields, leading to different phenomenologies, which we explore in turn. To leading order in $`1/m_{\mathrm{Pl}}`$, non-renormalizable holomorphic operators have the same form as those shown in Eq. (23),
$$K_{\mathrm{hol}}=\frac{W_{\mathrm{ren}}}{m_{\mathrm{Pl}}}+h.c.$$
(30)
We write generically these operators as $`X\varphi _1\varphi _2/m_{\mathrm{Pl}}`$. Let us assume for the moment that both $`\varphi _1`$ and $`\varphi _2`$ have masses of order of the electroweak scale.
The Kähler $`U(1)`$ connection $`K_j_\mu \varphi _jh.c.`$, with $`K_jK/\varphi _j`$ induces the following coupling to the fermionic components $`\psi ^i`$ of all the chiral superfields of the model WBook ,
$$_F\frac{1}{2m_{\mathrm{Pl}}^3}\overline{\psi }_L^i\gamma ^\mu \psi _L^i\mathrm{Im}_\mu (X\varphi _1\varphi _2),$$
(31)
which leads to a partial decay width $`M_X^7/m_{\mathrm{Pl}}^6`$. From Eq. (27) an effective Yukawa coupling is generated,
$$\frac{^2K}{\varphi _1\varphi _2}\frac{W}{m_{\mathrm{Pl}}^2}\overline{\psi }_R^1\psi _L^2\frac{m_{3/2}}{m_{\mathrm{Pl}}}X\overline{\psi }_R^1\psi _L^2,$$
(32)
which leads to a highly suppressed partial decay width $`(m_{3/2}/m_{\mathrm{Pl}})^2M_X^3/m_{\mathrm{Pl}}^2`$. Finally, couplings to goldstinos are generated notably by the gravitino mass term and gravitino kinetic terms:
$$\frac{1}{2m_{\mathrm{Pl}}}X\varphi _1\varphi _2\frac{W}{m_{\mathrm{Pl}}^2}\mathrm{\Psi }_\mu \sigma ^{\mu \nu }\mathrm{\Psi }_\nu +ϵ^{\mu \nu \rho \sigma }\frac{1}{m_{\mathrm{Pl}}^3}_\rho (X\varphi _1\varphi _2)\overline{\mathrm{\Psi }}_\mu \overline{\sigma }_\nu \mathrm{\Psi }_\sigma .$$
(33)
with the replacement $`\mathrm{\Psi }_\mu i\sqrt{2/3}_\mu G/m_{3/2}`$.
All terms lead to extremely slow decay if $`\varphi _1`$ and $`\varphi _2`$ have electroweak scale masses, and must be forbidden in order for the lightest messenger not to decay after BBN. However if one $`\mathrm{𝟐𝟒}_H`$ is present, the replacement of this field by its vev in the Kähler function shows that one recovers a mixing term as proposed in Ref. FY02 albeit with effective coupling $`M_{\mathrm{GUT}}/m_{\mathrm{Pl}}`$. This mixing term then leads to the same decay widths into one sfermion and one gaugino, or one sfermion and two goldstinos, as discussed before, albeit decreased by $`(M_{\mathrm{GUT}}/m_{\mathrm{Pl}})^210^5`$. The consequences for gravitino dark matter are shown in Fig. 3. As discussed before, we should consider the cases where the spurion is heavier or lighter than the lightest messenger. The left panel shows the case where the spurion is heavier than the lightest messenger and the annihilation cross-section into goldstinos increases with increasing $`M_X`$ to saturate at the unitarity bound above the dashed line. Only a small portion of parameter space allows for gravitino dark matter in this case; it is actually an amputated part of the solution shown in the right panel, see below. Above the dashed line, the results are highly uncertain since multi-goldstino is not well controled. Moreover, the lightest messenger can decay into a sneutrino and a pair of goldstinos due to the above mixing, with a decay width which is expected to scale as $`\mathrm{\Gamma }_{X\stackrel{~}{\nu }\stackrel{~}{G}\stackrel{~}{G}}(F_S^2M_X^4/F^4)\mathrm{\Gamma }_{Xl\lambda }`$. The prefactor $`(F_S^2M_X^4/F^4)`$ denotes the effective coupling of $`XX^{}`$ to a pair of goldstinos in the heavy spurion limit. In the region in which unitarity is violated, it has been assumed that this effective coupling saturates at the value reached \[$`𝒪(1)`$\] when the annihilation cross-section reaches the unitarity bound. It is then comparable to the decay width into a lepton and a gaugino and produces highly relativistic gravitinos. The energy density contained in these gravitinos exceeds the BBN bounds on additional relativistic degrees of freedom so that most of this region is excluded, as indicated by the NW-SE oriented dashed lines in Fig. 3. The SW-NE oriented dashed lines exclude the part of parameter space at small $`m_{3/2}`$ in which the lightest messenger decay occurs so late that it is forbidden by BBN constraints on energy injection. If the NLSP were a stau, this region would still be forbidden but the constraints at large $`m_{3/2}`$ would be relaxed.
In the right panel of Fig. 3, the spurion is assumed to be lighter than the lightest messenger and annihilation into goldstinos is less effective. One finds a solution for gravitino dark matter in a large part of parameter space, $`m_{3/2}10\mathrm{keV}1\mathrm{MeV}`$, for scenarios of indirect gauge mediation, i.e. $`k1`$. This solution is the same as that discussed in Section IV.A.2, albeit shifted to smaller values of $`M_X`$; this can be understood from the fact that for a same relic abundance of $`X`$, entropy production is larger in the present case since the decay width is further suppressed. Hence one can tolerate a smaller relic abundance, or equivalently a higher annihilation cross-section, i.e. a smaller $`M_X`$.
Finally consider now non-holomorphic non-renormalizable couplings between $`X`$ and visible sector particles in $`K`$. Such couplings can take the form:
$`K`$ $`{\displaystyle \frac{1}{m_{\mathrm{Pl}}}}\{`$ $`\mathrm{𝟓}_M^{}\overline{\mathrm{𝟓}}_{H,F}\mathrm{𝟏𝟎}_F,\overline{\mathrm{𝟓}}_M\mathrm{𝟓}_H^{}\mathrm{𝟏𝟎}_F,\overline{\mathrm{𝟓}}_M^{}\mathrm{𝟏𝟎}_F\mathrm{𝟏𝟎}_F,`$ (34)
$`\mathrm{𝟓}_M^{}\mathrm{𝟓}_H\mathrm{𝟐𝟒}_H,\mathbf{\hspace{0.17em}5}_M\mathrm{𝟓}_H^{}\mathrm{𝟐𝟒}_H,\overline{\mathrm{𝟓}}_M\overline{\mathrm{𝟓}}_{H,F}^{}\mathrm{𝟐𝟒}_H,`$
$`\overline{\mathrm{𝟓}}_M^{}\overline{\mathrm{𝟓}}_{H,F}\mathrm{𝟐𝟒}_H,\mathbf{\hspace{0.17em}10}_M^{}\mathrm{𝟓}_H\mathrm{𝟓}_H,\overline{\mathrm{𝟏𝟎}}_M\overline{\mathrm{𝟓}}_{H,F}^{}\mathrm{𝟓}_H,`$
$`\overline{\mathrm{𝟏𝟎}}_M^{}\overline{\mathrm{𝟓}}_{H,F}\overline{\mathrm{𝟓}}_{H,F},\overline{\mathrm{𝟏𝟎}}_M^{}\mathrm{𝟏𝟎}_F\mathrm{𝟓}_H,\mathbf{\hspace{0.17em}10}_M\mathrm{𝟏𝟎}_F\overline{\mathrm{𝟓}}_{H,F}^{},`$
$`\mathrm{𝟏𝟎}_M^{}\mathrm{𝟏𝟎}_F\mathrm{𝟐𝟒}_H,\mathbf{\hspace{0.17em}10}_M\mathrm{𝟏𝟎}_F^{}\mathrm{𝟐𝟒}_H,+\mathrm{h}.\mathrm{c}.\}`$
As before, we write this coupling as $`X^{}\varphi _1\varphi _2/m_{\mathrm{Pl}}`$ and assume for the moment that $`\varphi _1`$ and $`\varphi _2`$ carry electroweak scale masses. Then decay into $`\varphi _1`$ and $`\varphi _2`$ with width $`\mathrm{\Gamma }10^2M_X^3/m_{\mathrm{Pl}}^2`$ occurs via the mixing of kinetic terms between $`\varphi _1`$ and $`X`$ or between $`\varphi _2`$ and $`X`$. As seen before, such a decay width does not lead to solutions for gravitino dark matter as entropy production is not significant. Inspection of Eq. (34) reveals that all terms fall in the above category except those involving one $`\mathrm{𝟐𝟒}_H`$ as well as $`\overline{\mathrm{𝟏𝟎}}_M^{}\mathrm{𝟏𝟎}_F\mathrm{𝟓}_H`$ and $`\mathrm{𝟏𝟎}_M\mathrm{𝟏𝟎}_F\overline{\mathrm{𝟓}}_H^{}`$.
The latter two terms necessarily contain one colored Higgs with GUT mass but no vev, which we assume to be $`\varphi _2`$. Then the scalar potential term involving the inverse of the Kähler metric $`g`$ contains a coupling of $`X`$ to visible sector particles: $`g^{\varphi _2M^{}}D_{\varphi _2}WD_M^{}W^{}(\varphi _1/m_{\mathrm{Pl}})yM_X\mathrm{𝟏𝟎}_F\mathrm{𝟏𝟎}_FX`$. This leads to decay into three-body final state with width $`\mathrm{\Gamma }10^4y^2M_X^3/m_{\mathrm{Pl}}^2`$, again too large to yield solutions for gravitino dark matter.
Finally, if coupling to one $`\mathrm{𝟐𝟒}_H`$ occurs, say $`\varphi _2`$, mass mixing of order $`M_{\mathrm{GUT}}M_X`$ occurs between $`\varphi _1`$ and $`X`$ and leads to one negative mass squared eigenstate; this coupling must therefore be excluded.
## V Discussion
In this section we explore qualitatively other possible avenues which may help reconcile gravitino dark matter with more generic GMSB scenarios. Indeed, the previous discussion has shown that interesting solutions for gravitino dark matter and/or the gravitino problem in GMSB with $`SU(5)`$ grand unification with a high post-inflationary reheating temperature, can only be found for some very specific couplings between the messenger and visible sector and in some restricted regions of parameter space. It is furthermore necessary to assume that the spurion is much lighter than the lightest messenger so that multi-goldstino production remains at a safe level. A gravitational decay width $`\mathrm{\Gamma }10^3M_X^3/m_{\mathrm{Pl}}^2`$, which is generic in the sense that it is generated by most allowed non-renormalizable messenger-matter interaction terms, does not lead to satisfying solutions for gravitino dark matter.
It is instructive to consider the case of decay widths with the same scaling but whose prefactor is much smaller, $`\mathrm{\Gamma }ϵM_X^3/m_{\mathrm{Pl}}^2`$ with $`ϵ10^3`$. Figure 4 shows the solution for $`ϵ=10^{10}`$. Such decay width can be achieved by terms of the form $`WX\mathrm{\Phi }_1\mathrm{\Phi }_2\mathrm{\Phi }_3H_{24}/m_{\mathrm{Pl}}^2`$, with $`H_{24}`$ designing a Higgs with non-zero vev in $`\mathrm{𝟐𝟒}_H`$, or by most non-renormalizable couplings to order $`1/m_{\mathrm{Pl}}`$ discussed in the previous section, provided they are further suppressed by a factor $`10^7`$. In this figure, one finds that in the left panel, where $`S`$ is assumed heavier than $`X`$, i.e. where multi-goldstino production plays a significant role, there is room for gravitino dark matter only in a very limited region of parameter space. In this area, furthermore, the post-decay reheating temperature is quite close to $`1`$MeV. On the contrary, in the right panel one recovers a solution for gravitino dark matter with mass $`m_{3/2}10\mathrm{keV}10\mathrm{MeV}`$ in direct gauge mediated scenarios $`k1`$. Higher values of the coupling $`ϵ`$ lead to solutions shifted to higher $`M_X`$, with the dashed region due to late messenger decay shifting downwards. Lower values of $`ϵ`$ lead to solutions shifted to smaller values of $`M_X`$, but with the excluded dashed region moving upwards in $`M_X`$.
Dangerous operators involving GUT-scale vev’s may also be present. However, global continuous R-symmetries are expected to play an important rôle in scenarios of supersymmetry breaking BKN94 , and in a generic setting they could control the absence of unwanted operators for properly chosen R-charges SN94 . Discrete Z-symmetries, motivated by the need to improve the fine-tuning issues AG97 , can also play a selective rôle. In particular, the spurion gauge singlet superfield $`S`$ can be present or not in non-renormalizable operators in the Kähler or superpotential depending on its $`R`$ and $`Z`$ charges attributions. Such operators, provided that they do not destabilize the mass hierarchies Abel96 , JBPR , would lead to tree-level suppressions of the form $`(S/m_{\mathrm{Pl}})^{2n}(M_X/m_{\mathrm{Pl}})^{2n}`$. However the decay width is now too suppressed to yield reheating before BBN except in a very narrow region centered on $`m_{3/2}1`$MeV and $`M_X10^{10}`$GeV.
In contrast, a larger number of seemingly generic solutions to the gravitino problem/gravitino dark matter may be found if one considers $`SO(10)`$ grand unification and $`M_S>M_X`$. This case is studied analytically in a companion paper LMJ05 for $`M_X10^6`$GeV. In Section II, it has been shown that the amount of entropy produced depends directly on the relic abundance of the lightest messenger, which in turn depends directly on its nature. In $`SO(10)`$, the lightest messenger is a $`SU(3)\times SU(2)\times U(1)`$ singlet ($`\stackrel{~}{\nu }_R`$like), hence its annihilation cross-section is suppressed: it may either annihilate through one-loop diagrams or at tree level into goldstinos, and (at tree level) through suppressed diagrams of GUT mass bosons exchange. One may estimate the relic abundance of the lightest messenger in this case by using the dimensional estimate $`\sigma _{XX\mathrm{}}v(\alpha /4\pi )^4/M_X^2`$ for one-loop diagrams and the annihilation cross-section into goldstinos computed in Section II.B (see also Ref. LMJ05 ). One finds that the relic abundance is larger than for the $`SU(5)`$ case, hence the amount of entropy production is expected to be correspondingly larger. One then finds that a generic gravitational decay width $`\mathrm{\Gamma }10^3M_X^3/m_{\mathrm{Pl}}^2`$ can lead to natural solutions for gravitino dark matter in a significant part of parameter space, as shown in Fig. 5.
This solution is attractive for several reasons. As seen in the left panel of Figs. 5, gravitino cold dark matter with the right relic abundance occurs in a rather large part of parameter space, hence it appears “natural” in this sense. Moreover the decay width to sparticles assumed $`\mathrm{\Gamma }10^3M_X^3/m_{\mathrm{Pl}}^2`$ is generic as it is predicted by most non-renormalizable operators that violate messenger number by one unit. In this region of parameter space, the factor $`kF_S/F10^510^2`$, hence gravitino dark matter would be obtained for indirect gauge mediated scenarios; in Ref. LMJ05 it is argued that this case can be naturally incorporated in the simplest indirect GMSB scenarios GMSB . Furthermore, the solution for gravitino dark matter occurs in the predictive region in which multi-goldstino production satisfies the unitarity bound, unlike the solutions seen hitherto for $`SU(5)`$. Finally, Fig. 5 assumes, as the previous figures, that the NLSP is a bino; for a stau NLSP, the BBN constraints at large $`m_{3/2}`$ would be relaxed, and this would enlarge in turn the space of solutions for $`\mathrm{\Omega }_{3/2}`$.
One can show LMJ05 that the amount of entropy production does not hinder successful leptogenesis at high reheating temperatures; this is all the more interesting as leptogenesis scenarios typically operate in models of $`SO(10)`$ grand unification rather than $`SU(5)`$. Strictly speaking, the above scenario requires the spurion field to be heavier than the lightest messenger. This issue is however model-dependent as was briefly discussed in section II.B. For completeness, we illustrate in the left panel of Fig. 5 the opposite configuration where the lightest messenger annihilation into a pair of spurion fields is controlled by Eq.(8). This annihilation leads to a too low messenger relic density for the entropy dilution mechansim to work.
Finally we note that the results obtained in the present study remain valid when $`R`$parity is violated. In effect, in this case the gravitino lifetime is $`\tau _{3/2}10^{20}\mathrm{sec}(m_{3/2}/1\mathrm{GeV})^3`$ for trilinear $`R`$parity violating terms MC02 or $`\tau _{3/2}10^{27}\mathrm{sec}(m_{3/2}/1\mathrm{GeV})^3`$ for bilinear $`R`$parity violating terms TY00 . Hence the gravitino is sufficiently long-lived that it can be considered as stable dark matter with respect to the formation of large-scale structure. If the gravitino lifetime $`10^{27}\mathrm{sec}`$ one also finds that distortions of the diffuse backgrounds due to gravitino decay are evading observational constraints. For trilinear $`R`$parity violating terms, this requires $`m_{3/2}10`$MeV, while for bilinear terms, $`m_{3/2}1`$GeV is sufficient. With regards to the NLSP, its decay can proceed into visible sector particles on a short timescale and BBN constraints can be evaded, albeit they are replaced with constraints on diffuse background distortions. Hence the plots in parameter space would look similar to what has been found above.
## VI Conclusions
We have presented an exploratory though detailed investigation of relic LSP gravitino abundances in scenarios of gauge mediated supersymmetry breaking (GMSB). This study focuses on the possibility of gravitino dark matter and on solving the light gravitino overproduction problem for reheating temperatures after inflation that are ”arbritrarily” high. GMSB scenarios contain intermediate mass scale $`10^4\mathrm{GeV}M_X10^{12}`$GeV messenger fields which by virtue of their gauge interactions are easily produced in the primordial plasma. Cosmology requires these particles to subsequently decay as they would otherwise overclose the Universe (except for a lightest messenger with $`M_X1030`$TeV). Flavor-changing neutral currents impose somewhat restrictive limits on messenger number violating Yukawa interactions, possibly arguing for such messenger number violation to be rather weak. If so, the delayed decay of messengers may subsequently dilute any pre-existing gravitino abundances in accord with cosmological constraints.
We have thus investigated a fairly complete set of renormalizable and non-renormalizable messenger number violating operators within supersymmetric unification in $`SU(5)`$ (as well as some within $`SO(10)`$) and their impact on relic gravitino abundances. Results are shown for a variety of operators and imposing relevant constraints on NLSP decay and messenger decay from BBN, as well as constraints on the ”warmness” of gravitino dark matter from the required sucessful formation of large-scale structure. With respect to prior, less detailed, studies DGP96 ; BM03 ; FY02 ; FY02b ; FIY04 ; TY00 , we have uncovered a number of significant changes, notably the importance of messenger-messenger annihilation into two goldstinos in part of the $`M_X`$ \- $`m_{3/2}`$ parameter space, which modifies the messenger pre-decay freeze-out abundances in $`SU(5)`$ and $`SO(10)`$ grand unification.
In general, we have found that gravitino dark matter in $`SU(5)`$ grand unification in scenarios with high post-inflationary reheating temperatures $`T_{\mathrm{RH}}`$ is only possible for a few specific messenger-matter couplings. Furthermore we have shown that these models predict gravitino dark matter in regions of parameter space in which messengers annihilation to goldstinos violates unitarity unless one makes specific assumptions on the mass spectrum of GMSB models, and in particular, that the spurion $`S`$ be much lighter than the lightest messenger.
In contrast, in $`SO(10)`$ grand unification gravitino dark matter may be obtained for a variety of generic operators and in the predictive region of parameter space where multi-goldstino production is under control, as long as renormalizable messenger number violating interactions in the superpotential are absent LMJ05 . We thus believe that gravitino dark matter in GMSB scenarios is a viable alternative to neutralino (and gravitino) dark matter in supergravity scenarios, and as such deserves further detailed study. |
warning/0506/hep-th0506019.html | ar5iv | text | # On non-uniform smeared black branes
## 1 Introduction
Resolving the structure of vacuum solutions of Kaluza-Klein theory and finding the final fate of the Gregory-Laflamme (GL) instability of black objects has been an intriguing subject in the last decade. Recent progress in the subject has revealed many aspects of the problems, but at the same time several important issues remain to be resolved .
Perturbative analysis of the vacuum system on a circle has been made extensively, in particular for black holes recently. For uniform black strings, perturbative non-uniform deformation is allowed only for the critical uniform string that has the static GL mode . However, the most interesting state of black objects on a circle is in the regime beyond perturbations, and so far only a numerical approach has succeeded. The fully nonlinear static solutions with nontrivial horizon geometry have been constructed in 6 dimensions (and, partially, 5 dimensions) . The numerical solutions are constructed by employing the conformal gauge that is a general metric ansatz under axisymmetry and advantageous to numerics. The constructed two branches of a vacuum state, black hole and black string, help us to make a phase diagram in this theory .
To push forward our understanding of vacuum states in many directions, it might be useful to employ a metric ansatz proposed by Harmark and Obers (HO) , which involves just two undetermined functions. Although the metric ansatz was originally hypothetical because its proposers could not provide a proof of self-consistency, it has been shown that their simple ansatz is equivalent to the general conformal ansatz by an appropriate coordinate transformation . A key point is that one can also show the consistency of the associated boundary conditions.
While vacuum solutions represent a rich vacuum structure, charged black objects on a circle are also an interesting subject. A system of charged black objects has in general a thermodynamically stable parameter regime, and related to this fact, there is a conjecture that correlates dynamical stability and thermodynamic stability. The conjecture is known as the Gubser-Mitra (GM) conjecture (or the correlated stability conjecture ), which states that for systems with a translational symmetry and an infinite extent the dynamical GL instabilities arise precisely when the system is thermodynamically unstable. There is a lot of evidence for this conjecture, and no fully demonstrated counterexamples have been discovered so far (see, e.g., ). Another aspect of introducing charges into the system on a circle is connections with non-gravitational theories . Besides, existence of charged non-uniform black branes having higher entropy than their uniform counterparts is conjectured in Ref. .
In this paper, we investigate charged dilatonic black branes smeared on a transverse circle, which we call smeared black branes. Following the perturbation theory developed by Gubser , we study the stability and possible non-uniform deformation of the smeared black branes. To perform the perturbations, we employ the HO metric ansatz, which greatly simplifies the analysis in many respects. Based on the solutions of non-uniform smeared black branes, which include vacuum case, we study the thermodynamic properties of the system in detail. We think that the perturbative analysis developed in this paper will give a good theoretical basis for constructing fully nonlinear solutions of the HO ansatz, as the original perturbative solution played an important role in constructing nonlinear solutions.
The organization of this paper is as follows. In the next section, we review the smeared black branes in the HO metric ansatz and derive thermodynamical quantities. In Sec. 3, we perform a perturbative analysis up to the third order. The thermodynamics of the smeared black branes is discussed in Sec. 4, and Sec. 5 is devoted to a conclusion.
## 2 Smeared black branes
### 2.1 Harmark-Obers ansatz
The part of the classical supergravity action relevant for the considerations in this paper is
$`I_D={\displaystyle \frac{1}{16\pi G_D}}{\displaystyle d^Dx\sqrt{g}\left[R\frac{1}{2}(\varphi )^2\frac{1}{2(p+2)!}e^{\overline{a}\varphi }F_{p+2}^2\right]},`$ (1)
where $`F_{p+2}`$ is a $`(p+2)`$-form field and $`\varphi `$ is a dilaton field. The coupling constant $`\overline{a}`$ is $`\overline{a}^2=42(p+1)(d2)/(D2)`$ with $`D=p+d+1`$. From this action, we obtain the following equations:
$`R_{\mu \nu }={\displaystyle \frac{1}{2}}_\mu \varphi _\nu \varphi +{\displaystyle \frac{e^{\overline{a}\varphi }}{2(p+2)!}}\left[(p+2)F_\mu ^{\mu _2\mathrm{}\mu _{p+2}}F_{\nu \mu _2\mathrm{}\mu _{p+2}}{\displaystyle \frac{p+1}{D2}}g_{\mu \nu }F_{p+2}^2\right],`$ (2)
$`\mathrm{}\varphi ={\displaystyle \frac{\overline{a}}{2(p+2)!}}e^{\overline{a}\varphi }F_{p+2}^2,`$
$`_\mu \left(\sqrt{g}e^{\overline{a}\varphi }F^{\mu \mu _2\mathrm{}\mu _{p+2}}\right)=0.`$
The form field must satisfy the Bianchi identity, $`_{[\mu }F_{\mu _1\mathrm{}\mu _{p+2}]}=0.`$
The ansatz proposed in for a charged dilatonic black $`p`$-brane with transverse space $`^{d1}\times S^1`$ is
$`ds_D^2=H^{\frac{d2}{D2}}\left[fdt^2+{\displaystyle \underset{i=1}{\overset{p}{}}}\left(dx^i\right)^2+H\left({\displaystyle \frac{L}{2\pi }}\right)^2\left({\displaystyle \frac{A}{f}}dR^2+{\displaystyle \frac{A}{B^{d2}}}dv^2+BR^2d\mathrm{\Omega }_{d2}^2\right)\right],`$
$`f=1\left({\displaystyle \frac{R_0}{R}}\right)^{d3},H=1+\left({\displaystyle \frac{R_0}{R}}\right)^{d3}\mathrm{sinh}^2\alpha ,`$
$`e^{2\varphi }=H^{\overline{a}},𝒜_{tx^1\mathrm{}x^p}=(1H^1)\mathrm{coth}\alpha ,`$ (3)
where $`A=A(R,v)`$ and $`B=B(R,v)`$ are undetermined functions, and $`𝒜`$ is the potential of the form field $`F=d𝒜`$. The constant $`\alpha `$ is a parameter of electric charge. The coordinate $`v`$ has the periodicity $`v=v+2\pi `$, and the asymptotic size of the circle is $`L`$. The spacetime has an event horizon at $`R=R_0`$. The uniform smeared black $`p`$-brane is given by setting $`A=B=1`$. The consistency between this ansatz and the conformal form for the vacuum case and the charged non-dilatonic case are given in and , respectively. The consistency for the most general (charged dilatonic) case is discussed in by applying a Lorentz boost and U-duality.
The remarkable character of this metric is that the EOMs for the two unknown functions $`A`$ and $`B`$ are independent of the value of $`p`$ and the charge parameter $`\alpha `$. Thus the EOMs are the same as those of the neutral non-dilatonic black branes on $`^{d1}\times S^1`$, which we represent in Appendix A. Consequently, the problem of finding solutions of charged non-uniform smeared branes is mapped to the problem of neutral non-uniform strings with the effective spacetime dimension $`𝒟`$<sup>3</sup><sup>3</sup>3 Generating (near-extremal) charged solutions from uncharged solutions via the M-theory lift-boost-reduction is possible. This kind of solution-generating technique was used in .
$`𝒟Dp=d+1.`$ (4)
For example, the non-dilatonic charged cases are $`(D,p,d;\overline{a})=`$ (11, 5, 5; 0) and (10, 3, 6; 0), and their effective dimensions are $`𝒟=6`$ and $`7`$, respectively. The two cases correspond to M5- and D3-branes. In this paper, we take $`𝒟`$ as a parameter, and we study perturbations in a wide range of $`𝒟`$.
### 2.2 Thermodynamic quantities
Thermodynamic quantities of the smeared black branes are calculated from the metric (3). The mass, temperature, entropy and charge are given by
$`M`$ $`=`$ $`𝒢\left({\displaystyle \frac{LR_0}{2\pi }}\right)^{d3}(d3)\left[{\displaystyle \frac{d2}{d3}}\chi +\mathrm{sinh}^2\alpha \right],`$
$`T`$ $`=`$ $`{\displaystyle \frac{d3}{2R_0L\sqrt{A_h}\mathrm{cosh}\alpha }},`$
$`S`$ $`=`$ $`4\pi 𝒢\sqrt{A_h}\left({\displaystyle \frac{LR_0}{2\pi }}\right)^{d2}\mathrm{cosh}\alpha ,`$
$`Q`$ $`=`$ $`𝒢\left({\displaystyle \frac{LR_0}{2\pi }}\right)^{d3}(d3)\mathrm{cosh}\alpha \mathrm{sinh}\alpha ,`$ (5)
where we have used $`𝒢LV_p\mathrm{\Omega }_{d2}/16\pi G_D`$. $`A_h(R_0)A(R_0,v)`$ and $`\chi =\chi (R_0)`$ are constants for a given horizon radius. The latter is defined by the asymptotic behavior of the metric functions:
$`A,B1\chi \left({\displaystyle \frac{R_0}{R}}\right)^{d3},\text{for}RR_0.`$ (6)
$`V_p`$ is the volume of $`x^i`$ ($`i=1,2,\mathrm{},p`$) directions. It is not necessary to take the coordinate $`x^i`$ to be periodic, but we do it formally so that we may present a finite expression of mass and charge. The volume of the unit $`n`$-dimensional sphere is $`\mathrm{\Omega }_n=2\pi ^{(n+1)/2}/\mathrm{\Gamma }\left[(n+2)/2\right]`$. For a regular event horizon, $`A_h`$ is independent of $`v`$ so that the zeroth law of thermodynamics holds. The tensions, which are asymptotic charges accompanied by the $`S^1`$ compactification, are given by
$`L𝕋`$ $`=`$ $`𝒢\left({\displaystyle \frac{LR_0}{2\pi }}\right)^{d3}\left[1(d3)(d2)\chi \right],`$
$`L_i𝕋_i`$ $`=`$ $`𝒢\left({\displaystyle \frac{LR_0}{2\pi }}\right)^{d3}\left[1+(d3)(\mathrm{sinh}^2\alpha \chi )\right],`$ (7)
where $`𝕋`$ and $`𝕋_i`$ are the tensions along the $`v`$-direction and along the $`i`$-th world-volume with circle length $`L_i`$, respectively. Defining the dimensionless tension $`n`$ by
$`n{\displaystyle \frac{L𝕋}{MM_{el}}},`$ (8)
where $`M_{el}\left[(d1)L_i𝕋_iL𝕋M\right]/(d2)`$ is “electric” mass , the Smarr formula then takes the form
$`TS={\displaystyle \frac{d2n}{d1}}(MM_{el}).`$ (9)
Note that for the neutral case ($`\alpha =0`$), we have $`M_{el}=0`$ and that then the dimensionless tension becomes the one commonly employed for vacuum black branes.
The specific heat $`C_Q`$ and the isothermal permittivity $`ϵ_T`$ are computed as
$`C_Q`$ $``$ $`\left({\displaystyle \frac{M}{T}}\right)_Q={\displaystyle \frac{V_p\mathrm{\Omega }_{d2}L^{d1}}{4G}}\left({\displaystyle \frac{R_0}{2\pi }}\right)^{d2}{\displaystyle \frac{\left[d3+(d1)\mathrm{cosh}2\alpha \right]\mathrm{cosh}\alpha }{(d3)+(d5)\mathrm{cosh}2\alpha }},`$ (10)
$`ϵ_T`$ $``$ $`\left({\displaystyle \frac{Q}{\mathrm{\Phi }}}\right)_T={\displaystyle \frac{V_p\mathrm{\Omega }_{d2}L^{d2}}{32\pi G}}\left({\displaystyle \frac{R_0}{2\pi }}\right)^{d3}(d3)\left[(d3)(d5)\mathrm{cosh}2\alpha \right]\mathrm{cosh}^2\alpha ,`$
where
$`\mathrm{\Phi }=\mathrm{tanh}\alpha `$ (11)
is an electric potential energy at the horizon. The specific heat changes the sign from minus to plus at a critical value of $`\alpha `$, if $`d6`$. For a fixed charge $`Q`$, this implies that a thermodynamically stable region exists for the smeared black branes in the parameter region in which black branes suffer from dynamical GL instability . However, we should consider the thermodynamic ensemble that allows the charge to change. In other words, we should also take into account the isothermal permittivity, which probes the thermodynamic stability under changes of the charge . One finds that the specific heat and isothermal permittivity always have opposite signs from the above expression. Consequently, the two conditions $`C_Q>0`$ and $`ϵ_T>0`$, which guarantee thermodynamic stability in the grand canonical ensemble, cannot hold simultaneously. That is to say, the smeared black brane background is never thermodynamically stable. This is perfectly compatible with the existence of a static mode that is a signal of the onset of unstable modes. We discuss the static mode in detail in the next section.<sup>4</sup><sup>4</sup>4 A partial preliminary analysis of dynamical s-wave perturbations is found in .
## 3 Static perturbation
According to the perturbation theory in Ref. , we perform static perturbations of the smeared black $`p`$-branes.<sup>5</sup><sup>5</sup>5 Readers who are interested in the details in this section can examine a Mathematica notebook In general, the static perturbations are easier to perform than dynamical perturbations, and they allows us to construct a non-uniform solution as well as to determine the GL critical waves. For convenience, we transform the metric functions and the coordinates as
$`A=e^a,B=e^b,`$
$`x={\displaystyle \frac{L}{2\pi }}v,y={\displaystyle \frac{L}{2\pi }}R.`$ (12)
It is possible to rescale the entire metric so that the event horizon can be located at $`y=1`$ in the new coordinates. In the following analysis we adopt this normalization.
We expand the metric function $`X(x,y)`$ ($`X=a,b`$) around the uniform solution as follows:
$`X(x,y)={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}ϵ^nX_n(y)\mathrm{cos}nKx.`$ (13)
Here, $`ϵ`$ is a small parameter of expansion. The function $`X_n`$ and wave number $`K=2\pi /L`$ are also expanded as
$`X_n(y)={\displaystyle \underset{p=0}{\overset{\mathrm{}}{}}}ϵ^{2p}X_{n,p}(y),K={\displaystyle \underset{q=0}{\overset{\mathrm{}}{}}}ϵ^{2q}k_q,`$ (14)
where $`X_{0,0}(y)=0`$ is imposed. Some leading order terms are
$`X(x,y)`$ $`=`$ $`ϵX_{1,0}(y)\mathrm{cos}k_0x+ϵ^2\left[X_{0,1}(y)+X_{2,0}(y)\mathrm{cos}2k_0x\right]`$ (15)
$`+ϵ^3\left[X_{1,1}(y)\mathrm{cos}k_0xX_{1,0}(y)k_1x\mathrm{sin}k_0x+X_{3,0}(y)\mathrm{cos}3k_0x\right]+O(ϵ^4).`$
Having expanded the metric functions as above, all we have to do is to solve the ordinary differential equations (ODEs) for $`X_{n,p}(y)`$ at each order of $`ϵ`$. Before entering into detail, however, it would be suitable to outline our calculations. First of all we should notice that in this perturbation theory the asymptotic circle length (or the wave number) is not fixed at all, and we have to determine it by solving the perturbations. At the first order $`O(ϵ)`$, we look for a solution $`X_{1,0}`$, which corresponds to the GL static mode. Thus, we can show the instability of the uniform solution by finding the GL critical wave number $`k_0`$. At the second order $`O(ϵ^2)`$, a back-reaction leads to nontrivial $`X_{0,1}`$ and $`X_{2,0}`$. One of them is the Kaluza-Klein (KK) massive mode, which falls off exponentially for a large $`y`$. The other is a massless mode, which determines the mass of smeared branes as an asymptotic charge. Therefore, one might consider the second-order calculation to be enough to see a thermodynamical symptom of the non-uniform deformation. It is known, however, that to calculate entropy difference between non-uniform and uniform solutions for the same mass one needs third-order perturbations, at least in the uncharged case. This is also the case for our general situation.
It is important to note the advantage of working in the HO metric. The metric ansatz reduces the number of equations of motion (EOMs) greatly. As discussed in the original paper , the EOMs do not contain the charge parameter $`\alpha `$ and the spatial world-volume $`p`$. The EOMs for most general charged diatonic cases are equivalent to those for neutral black strings on a cylinder $`^{𝒟2}\times S^1`$ with $`𝒟=Dp=d+1`$. That is, we do not need to perturb the form and dilation fields at any order.
It is worthwhile to consider the meaning of the parameter $`ϵ`$. In a leading order, this parameter $`ϵ`$ can be related to the parameter used in the literature ,
$`\lambda {\displaystyle \frac{1}{2}}\left({\displaystyle \frac{R_{max}}{R_{min}}}1\right).`$ (16)
Here, $`R_{max}`$ and $`R_{min}`$ are the maximum and minimum of the proper horizon radius of the black object so that $`\lambda `$ means geometrically the non-uniformity of the black brane. Substituting the expansion (13) into the metric (3), the parameter $`ϵ`$ is identical to $`\lambda `$ at the linear order if we set $`b_{1,0}(1)=2`$, which is possible due to the linearity of the differential equation at this order. Consequently, the expansion parameter $`ϵ`$ describes in the above sense how the black brane is non-uniform. However, the relation between $`ϵ`$ and $`\lambda `$ becomes nonlinear when we take into account higher-order perturbations. Because of the nonlinear corrections, the naive identification between $`ϵ`$ and $`\lambda `$ is apparently invalid in the regime. After all, $`ϵ`$ is just an expansion parameter and is no more and no less than the parameter controlling the amplitude of the perturbations.
### 3.1 First-order perturbation $`X_{1,0}`$
Let us begin with the first-order perturbation. Substituting the expansion (15) into the Einstein equation (2), we obtain a master equation for $`b_{1,0}`$, which determines the GL critical wave number $`k_0`$,
$`_1b_{1,0}=0,`$ (17)
where
$`_1={\displaystyle \frac{d^2}{dy^2}}+U{\displaystyle \frac{d}{dy}}+V_{(21)},`$
$`U(y)={\displaystyle \frac{(d1)f^2(d3)(3d8)f+(d3)^2}{yf[d3+(d1)f]}},`$
$`V_{(ij)}(y)={\displaystyle \frac{i(d3)^2(f1)+jk_0^2y^2[d3+(d1)f]}{y^2f[d3+(d1)f]}}.`$ (18)
After some manipulation, one recognizes that the problem is the same as solving the Schrödinger-type equation with energy $`k_0^2`$, although the potential is rather complicated in the present case.<sup>6</sup><sup>6</sup>6 Another type of single linear second-order ODE was discussed in , based on perturbations of Euclidean Schwarzschild black holes. However, since the equation has a singular point between the horizon and the infinity, it is clearly less tractable compared to our master equation.
One can solve algebraically for $`a_{1,0}`$, with the result
$`a_{1,0}(y)={\displaystyle \frac{2(d2)f(b_{1,0}+yb_{1,0}^{})}{d3+(d1)f}},`$ (19)
where the prime denotes the derivative with respect to $`y`$.
The horizon is a regular singular point of the differential equation, and demanding the regularity of the perturbations on the horizon, an appropriate boundary condition is required at the horizon:
$`b_{1,0}^{}(1)+\left(2{\displaystyle \frac{k_0^2}{d3}}\right)b_{1,0}(1)=0.`$ (20)
According to the discussion in the beginning of this section, we take
$`b_{1,0}(1)=2.`$ (21)
From the zeroth law of thermodynamics, all Kaluza-Klein modes of $`a(x,y)`$ must vanish at the horizon. So we have $`a_{1,0}(1)=0`$, as is also evident from (19) and the finiteness of $`b_{1,0}(1)`$.
The perturbation equation has two independent solutions, and at infinity the respective solutions are exponentially growing or decaying. The asymptotic flatness implies the growing mode is absent at infinity. The asymptotic behavior of the decaying mode is
$`b_{1,0}e^{k_0y}y^{(d4)/2}.`$ (22)
Thus, finding the solution of Eq. (17) is a one parameter shooting problem, which is easy to carry out. The GL critical wave number $`k_0`$ is determined for several dimensions and is summarized in Table 1.<sup>7</sup><sup>7</sup>7 Note that we have performed the calculations up to $`𝒟=20`$, although we show only a part of them in the Tables. The values of $`k_0`$ do of course accord with previous results . The numerical plots of the solutions are given in Fig. 1.
Let us take a closer look at the solution near the horizon. We assume the following series expansion for $`b_{1,0}`$,
$`b_{1,0}=e^{k_0(y1)}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\beta _n(y1)^n.`$ (23)
Substituting this into Eq. (17), we obtain a five-term recurrence equation. The recurrence equation can be solved by employing the continued fraction method, which determines the critical wavelength. However, the shooting method is more tractable than the fraction method, and thus we refrain from presenting the latter. The first few terms are still useful to observing the near horizon behavior:
$`\beta _0=2,`$
$`\beta _1=2\left(k_02+{\displaystyle \frac{k_0^2}{d4}}\right),`$
$`\beta _2=82d+{\displaystyle \frac{k_0}{2}}\left[{\displaystyle \frac{k_0^2(k_0+4d12)}{(d3)^2}}8+6k_0\right],`$
$`\beta _3=\mathrm{}.`$ (24)
If we truncate the series expansion at some level and match the truncated solution to the asymptotic solution at an appropriate radius, we might be able to obtain an approximate solution of the critical wavelength. For the truncation at $`O(\beta _3)`$ and matching at $`y2`$, we find a fourth-order algebraic equation,
$$\begin{array}{c}k_0^4+(d3)\hfill \\ \hfill \times [4k_0^3+2\frac{\left\{138+(1+d)\left[49+3(1+d)\right]\right\}k_0^2}{d12}\frac{4(d16)(d3)k_0}{d12}4(d3)^2]0.\end{array}$$
(25)
A real solution of this algebraic equation provides a good approximate solution of the GL critical wave for $`𝒟9`$ with small errors (a few percent). A more accurate and readily available algebraic equation could be obtained by pushing the truncation level up to $`O(\beta _5)`$. But the algebraic equation becomes eighth order and it has no analytic solution in general.
We have seen the existence of the GL critical wave number $`k_0`$ for an arbitrary value of $`\alpha `$, which does not appear in the master equation (17). In other words, the smeared black branes suffer from the GL instability, irrespective of their charge. As mentioned in Sec. 2.2, however, the non-extremal smeared black brane cannot be locally thermodynamically stable for any value of $`\alpha `$. Therefore, the existence of the GL critical wave number does not imply the breakdown of the GM conjecture.
### 3.2 Second-order perturbation
The first-order perturbations determine the GL critical wave number $`k_0`$, and they are localized near the horizon due to the rapid decay at the far region. The localized modes produce sources for higher-order perturbations that give back-reaction to the background geometry. At the second order $`O(ϵ^2)`$, two types of mode appear: the KK modes, which decay exponentially, and the zero modes, which have power-law decay at the asymptotics. Since the power-law decay of the zero mode gives finite contributions to the asymptotic charges, the zero modes have rather different characters compared with the KK modes.
#### 3.2.1 Zero modes $`X_{0,1}`$
The equations for the zero modes are rather complicated. The ODEs for $`a_{0,1}`$ and $`b_{0,1}`$ are
$`b_{0,1}^{\prime \prime }+{\displaystyle \frac{3y^21}{y(y^21)}}b_{0,1}^{}{\displaystyle \frac{4}{y^21}}(a_{0,1}b_{0,1})=S_{b_{0,1}}(b_{1,0};k_0),`$
$`a_{0,1}^{\prime \prime }+{\displaystyle \frac{3y^21}{y(y^21)}}a_{0,1}^{}{\displaystyle \frac{12}{y^21}}(a_{0,1}b_{0,1})=S_{a_{0,1}}(b_{1,0};k_0),`$ (26)
where we abbreviate the source term to $`S_{X_{0,1}}`$ which consists of a quadratic expression in $`b_{1,0}`$ and its first and second derivatives. All second derivatives are eliminated from the source term using the equation of the first-order perturbations. The critical wave number appears in the source term.
To have decoupled equations, it is convenient to introduce new variables,
$`c_{0,1}b_{0,1}a_{0,1},d_{0,1}b_{0,1}+a_{0,1}.`$ (27)
Then the equations are
$`c_{0,1}^{\prime \prime }+{\displaystyle \frac{f+d3}{yf}}c_{0,1}^{}{\displaystyle \frac{2(d3)^2}{y^2f}}c_{0,1}`$ $`=`$ $`S_{c_{0,1}}(b_{1,0};k_0),`$
$`d_{0,1}^{\prime \prime }+{\displaystyle \frac{f+d3}{yf}}d_{0,1}^{}+{\displaystyle \frac{2(d3)(d1)}{y^2f}}c_{0,1}`$ $`=`$ $`S_{d_{0,1}}(b_{1,0};k_0),`$ (28)
where
$`S_{c_{0,1}}(b_{1,0};k_0)={\displaystyle \frac{1}{4y^2f}}\left\{2(d3)^2(a_{1,0}b_{1,0})^2+(d2)k_0^2y^2b_{1,0}\left[2a_{1,0}+(d3)b_{1,0}\right]\right\},`$
$`S_{d_{0,1}}(b_{1,0};k_0)={\displaystyle \frac{1}{4y^2f}}\left\{2(d3)(d1)(a_{1,0}b_{1,0})^2+(d2)k_0^2y^2b_{1,0}\left[2a_{1,0}+(d3)b_{1,0}\right]\right\}.`$
Note that from the Einstein equations we also have an additional equation which takes schematically the form $`d_{0,1}^{}=F(b_{0,1},c_{0,1},b_{0,1}^{},c_{0,1}^{};k_0)`$. It is not an independent equation; the second-order equation (28) for $`d_{0,1}`$ follows from the first-order equations.
If we assume the regularity of the perturbed quantities, we obtain the horizon boundary conditions by setting $`y=1`$:
$`c_{0,1}^{}(1)=2(d3)c_{0,1}{\displaystyle \frac{1}{4}}\left[2(d3)+(d2)k_0^2\right]b_{1,0}^2,`$
$`d_{0,1}^{}(1)=2(d1)c_{0,1}+{\displaystyle \frac{1}{4}}[2(d1)+(d2)k_0^2)]b_{1,0}^2.`$ (30)
Since at the asymptotic region the source terms decay as $`O(e^{2k_0y})`$, we can neglect the right-hand side of (28). Then the leading asymptotic behavior of $`d_{0,1}`$ and $`c_{0,1}`$ are $`d_{0,1}1/y^{(d3)}`$ and $`c_{0,1}1/y^{2(d3)}`$. Interestingly, the power-law decay of $`c_{0,1}`$ is so rapid that it does not contribute to the asymptotic charges, due to Eq. (6).
Now let us discuss how we can solve these equations. The shooting parameter of this system is $`c_{0,1}(1)`$. The value of $`d_{0,1}(1)`$ at the horizon is not a shooting parameter. It is an arbitrary constant. This is because $`d_{0,1}`$ has a shift symmetry in the differential equation and the boundary condition,
$`d_{0,1}d_{0,1}+\mathrm{const}.`$ (31)
Clearly, we cannot choose the constant in Eq. (31) arbitrarily: If we shift the value of $`d_{0,1}`$, the value of $`a_{0,1}`$ on the horizon changes, which means that the temperature and entropy of the black brane changes. The constant should be fixed by imposing the asymptotic flatness on $`d_{0,1}`$, after one obtains a solution.
The numerical plot of the solutions are presented in Fig. 1, in which one can see the power-law decay of the solution. If we expand the constant $`\chi `$ in Eq. (6) as
$`\chi =\chi _1ϵ^2+O(ϵ^4),`$ (32)
the asymptotic behaviors of $`a_{0,1}`$ and $`b_{0,1}`$ are obtained from $`d_{0,1}`$ as
$`a_{0,1},b_{0,1}{\displaystyle \frac{\chi _1}{y^{d3}}}.`$ (33)
It is worthwhile to note that from the first law of the thermodynamics the following relation holes,
$`\chi _1={\displaystyle \frac{a_{0,1}(1)}{2}}.`$ (34)
The results of numerics for various effective dimensions $`𝒟`$ are shown in Table 2.
#### 3.2.2 KK modes $`X_{2,0}`$
The KK modes at this order are $`a_{2,0}`$ and $`b_{2,0}`$. The function $`a_{2,0}`$ can be solved algebraically,
$$\begin{array}{c}a_{2,0}=\frac{2f(d2)(yb_{2,0})^{}}{d3+(d1)f}\frac{a_{1,0}\left[a_{1,0}(d1)b_{1,0}\right]}{8}\frac{(d2)fb_{1,0}}{4\left[d3+(d1)f\right]^2}\hfill \\ \hfill \times \left(2(d3)(d2)a_{1,0}+b_{1,0}\left\{2(d2)k_0^2y^2(d3)^2+f\left[5+(d4)d\right]\right\}\right).\end{array}$$
(35)
Here the source term consists of $`b_{1,0},b_{2,0}`$ (and also $`a_{1,0}`$). A master equation for $`b_{2,0}`$ is similar to the first-order perturbation:
$`b_{2,0}^{\prime \prime }+Ub_{2,0}^{}+V_{(12)}b_{2,0}=S_{b_{2,0}}(b_{1,0};k_0),`$ (36)
where the source term $`S_{b_{2,0}}`$ is quadratic in $`b_{1,0}`$ and its first derivative. The explicit form is
$$\begin{array}{c}S_{b_{2,0}}=\frac{d3}{4y^2f}[a_{1,0}^2+a_{1,0}b_{1,0}\{(d5)\frac{4f(d3)(d2)^2}{\left[d3+(d1)f\right]^2}\}\frac{2b_{1,0}^2}{\left[d3+(d1)f\right]^2}\hfill \\ \hfill \times \{(d3)[3d2(d2)k_0^2y^2]f\{d[136d+d^2+2(d2)k_0^2y^2]12\}\\ \hfill +\frac{f^2}{d3}[3356d+36d^210d^3+d^42(d2)(d1)^2k_0^2y^2]\}],\end{array}$$
(37)
where $`a_{1,0}`$ in the source term can be eliminated by Eq. (19).
The asymptotic behavior of these KK modes are $`b_{2,0},a_{2,0}\mathrm{exp}(2k_0y)`$, and the horizon boundary conditions are
$`b_{2,0}^{}(1)+\left(2{\displaystyle \frac{k_0^2}{d3}}\right)b_{2,0}=\left[{\displaystyle \frac{1}{2}}+{\displaystyle \frac{(d2)k_0^2}{d3}}\right]b_{1,0}^2`$ (38)
and $`a_{2,0}(1)=0`$. Since the critical wave number $`k_0`$ is given at the first order, this ODE is a shooting problem with one shooting parameter, i.e., amplitude at the horizon $`b_{2,0}(1)`$. Applying the same numerics as for the first order, we can easily obtain the solution. The parameter that gives a regular solution is summarized in Table 2.
### 3.3 Third-order perturbation $`X_{1,1}`$
As we see later, the second-order perturbations are not enough to determine the difference of the entropy between a non-uniform black brane and a uniform one of the same mass. Knowledge is required of the correction to the critical wave number, $`k_1`$, which appears at this third order. At the third order $`O(ϵ^3)`$, new KK modes $`X_{1,1}`$ and $`X_{3,0}`$ appear independently. The equations of the latter modes do not contain $`k_1`$, and so we are not interested in these modes. The KK modes $`X_{1,1}`$ are the higher-order counterparts of $`X_{1,0}`$; the form of the homogeneous equation is identical to Eq. (17). Then the master equation for $`b_{1,1}`$ takes the form of
$`_1b_{1,1}`$ $`=`$ $`{\displaystyle \frac{2k_0k_1}{f}}b_{1,0}+S_{b_{1,1}}(b_{0,1},b_{1,0},b_{2,0};k_0).`$ (39)
The function $`a_{1,1}(y)`$ can be solved algebraically as before. The source term $`S_{b_{1,1}}`$ is independent of $`k_1`$, and the algebraic equation for $`a_{1,1}`$ does not explicitly contain $`k_1`$. The dependence of $`a_{1,1}`$ on the critical wave number $`k_1`$ is only through $`b_{1,1}`$. The horizon boundary condition is given by
$$\begin{array}{c}b_{1,1}^{}+\frac{2d6k_0^2}{d3}b_{1,1}=\frac{2k_1k_0}{d3}b_{1,0}+\frac{b_{1,0}}{8(d3)}\hfill \\ \hfill \times \left\{2(d3)(8a_{0,1}8b_{0,1}+b_{1,0}^24b_{2,0})+(d2)k_0^2\left[8b_{0,1}+(d2)b_{1,0}^2+4b_{2,0}\right]\right\},\end{array}$$
(40)
where all functions are evaluated at $`y=1`$. We can solve this second-order equation in the same manner used for the first-order perturbations. The main difference is that the present case is a two-parameter shooting problem; $`b_{1,1}(1)`$ is not arbitrary. A regular solution exists only for an appropriate set of $`b_{1,1}(1)`$ and $`k_1`$. We have performed the two-parameter shooting, and the results for various dimensions are summarized in Table 1 and 2 (see also Fig. 1). We should notice that the sign of $`k_1`$ changes above $`𝒟=13`$:
$`k_1>0(𝒟13),`$
$`k_1<0(𝒟>13).`$ (41)
This would be a precursor to the existence of a critical dimension.
## 4 Thermodynamics
We are now ready to calculate the thermodynamical quantities of the perturbative non-uniform solutions. From Eq. (5), we obtain the mass $`M`$, temperature $`T`$, entropy $`S`$ and charge $`Q`$ for a non-uniform smeared black brane. The differences for these quantities between the non-uniform and the uniform critical solutions are
$`{\displaystyle \frac{\delta M}{M}}={\displaystyle \frac{\chi _1ϵ^2}{(d2)/(d3)+\mathrm{sinh}^2\alpha }}+O(ϵ^4),`$
$`{\displaystyle \frac{\delta T}{T}}={\displaystyle \frac{1}{\sqrt{A_h}}}1={\displaystyle \frac{1}{2}}a_{0,1}(1)ϵ^2+O(ϵ^4),`$
$`{\displaystyle \frac{\delta S}{S}}=\sqrt{A_h}1={\displaystyle \frac{1}{2}}a_{0,1}(1)ϵ^2+O(ϵ^4),`$
$`{\displaystyle \frac{\delta Q}{Q}}=0.`$ (42)
These physical quantities are expressed in terms of the perturbative quantities, which are calculated without fixing the asymptotic periodicity of the circle. It is easy to compare the obtained physical quantities with others if they are given in a frame where the asymptotic length of the circle is fixed. Without transforming physical quantities into such a frame, we introduce variables that are invariant under rigid rescalings of the entire solution. Once physical quantities are expressed in terms of the invariant variables, we do not need to care about the variation of the asymptotic periodicity.
The invariant quantities can be obtained by multiplying $`K`$ by suitable powers:
$`{\displaystyle \frac{\delta \mu }{\mu }}{\displaystyle \frac{\delta M}{M}}+(d3){\displaystyle \frac{\delta K}{K}}=\mu _1ϵ^2+O(ϵ^4),`$
$`{\displaystyle \frac{\delta \tau }{\tau }}{\displaystyle \frac{\delta T}{T}}{\displaystyle \frac{\delta K}{K}}=\tau _1ϵ^2+O(ϵ^4),`$
$`{\displaystyle \frac{\delta s}{s}}{\displaystyle \frac{\delta S}{S}}+(d2){\displaystyle \frac{\delta K}{K}}=s_1ϵ^2+O(ϵ^4),`$
$`{\displaystyle \frac{\delta q}{q}}{\displaystyle \frac{\delta Q}{Q}}+(d3){\displaystyle \frac{\delta K}{K}}=q_1ϵ^2+O(ϵ^4),`$ (43)
where the leading-order coefficients are given by
$`\mu _1={\displaystyle \frac{\chi _1}{(d2)/(d3)+\mathrm{sinh}^2\alpha }}+(d3){\displaystyle \frac{k_1}{k_0}},`$
$`\tau _1={\displaystyle \frac{1}{2}}a_{0,1}(1){\displaystyle \frac{k_1}{k_0}},`$
$`s_1={\displaystyle \frac{1}{2}}a_{0,1}(1)+(d2){\displaystyle \frac{k_1}{k_0}},`$
$`q_1=(d3){\displaystyle \frac{k_1}{k_0}}.`$ (44)
Here we see the reason why we have integrated the perturbations up to order $`O(ϵ^3)`$. The correction to the wave number $`k_1`$ is necessary to obtaining invariant quantities.
Note that non-extremal solutions obey the first law,
$`dM(R_0,\alpha )=T(R_0,\alpha )dS(R_0,\alpha )+\mathrm{\Phi }(\alpha )dQ(R_0,\alpha ),`$ (45)
provided the asymptotic circle length is fixed. We can confirm this directly from Eq. (5) by utilizing the first law for the uncharged system, which gives
$`{\displaystyle \frac{1}{2A_h}}{\displaystyle \frac{dA_h}{dR_0}}={\displaystyle \frac{d3}{R_0}}\chi {\displaystyle \frac{d\chi }{dR_0}}.`$ (46)
### 4.1 Vacuum black branes
To begin, let us first consider the thermodynamics of neutral black branes for vacuum spacetime.<sup>8</sup><sup>8</sup>8 Since there is no charge in this system, the metric (3) is identical to that of the (neutral vacuum) black $`(p+1)`$-brane. What we are interested in is the difference between the entropy of a uniform brane and that of a non-uniform one of the same mass. The difference is evaluated as follows:
$`{\displaystyle \frac{S_{\mathrm{NU}}S_\mathrm{U}}{S_\mathrm{U}}}=\sigma _1ϵ^2+\sigma _2ϵ^4+O(ϵ^6),`$
$`\sigma _1=\mu _1{\displaystyle \frac{d3}{d2}}s_1,\sigma _2={\displaystyle \frac{d2}{2(d3)}}\left(\tau _1+{\displaystyle \frac{1}{d3}}\mu _1\right)\mu _1,`$ (47)
where we have used the Smarr formula $`TS/M=(d2n)/(d1)`$ and the tension of uniform strings, $`n=1/(d2)`$. This equation can be easily derived by starting from the frame with $`\delta K=0`$ and re-expressing the results in terms of the invariant quantities (43). See Appendix B for more details of the calculation.
As discussed above, the relation (34) holds for $`\chi _1`$ and $`a_{1,0}(1)`$ from the first law, and then $`\sigma _1`$ vanishes. As a result, the entropy difference arises at $`O(ϵ^4)`$. The correction $`\sigma _2`$ becomes $`\sigma _2=a_{0,1}\mu _1/4`$. Since the numerical value of $`a_{0,1}`$ is negative (Table 2), the sign of $`\sigma _2`$ is determined by $`\mu _1`$. Our numerical results for $`𝒟=415`$ are summarized in Table 2, and they are consistent with previously obtained results ; $`\sigma _2`$ is negative for $`𝒟13`$, hence the non-uniform phase is not entropically favorable and it implies a first-order phase transition. For $`𝒟>13`$, $`\sigma _2`$ becomes positive, implying a higher-order phase transition as observed in Ref. . These are summarized as
$`\sigma _2<0(𝒟13),`$
$`\sigma _2>0(𝒟>13).`$ (48)
The vanishing of $`\sigma _1`$ can be used as a check of numerics. In fact, our numerical results give vanishingly small $`\sigma _1`$ for relatively low dimensions ($`𝒟9`$). For much higher dimensions, however, it is hard to read off $`\chi _1`$ numerically because of the rapid decay of the zero modes. Thus, we have calculated $`\chi _1`$ by using the relation (34), which allows very precise determination of physical quantities because $`a_{0,1}`$ is a local quantity. Moreover, the results obtained here are very simple and concise compared to those obtained by the general perturbation theory ; the invariant quantities (43) under rigid rescalings are given only by 3 variables, i.e., $`a_{0,1},k_0`$, and $`k_1`$, whereas for the general perturbation theory many quantities enter into the estimation of the scheme independent quantities.
Next, we consider a canonical ensemble, i.e., fixed temperature. Helmholtz free energy $`F=MTS`$ for the neutral solutions is given by
$`F=𝒢\left({\displaystyle \frac{LR_0}{2\pi }}\right)^{d3}\left[1(d3)\chi \right].`$ (49)
The free energy of the uniform strings is $`F_U=𝒢\left[(d3)/4\pi T\right]^{d3}`$. The difference of the free energy between the uniform and the non-uniform solutions for the same temperature is evaluated as
$`{\displaystyle \frac{F_{\mathrm{NU}}F_\mathrm{U}}{F_\mathrm{U}}}=\rho _2ϵ^4+O(ϵ^6),\rho _2={\displaystyle \frac{1}{2}}(d3)(d2)\left(\tau _1+{\displaystyle \frac{1}{d3}}\mu _1\right)\tau _1.`$ (50)
See Appendix B for the derivation of Eq. (50). The structure of the coefficients $`\rho _2`$ is quite similar to that in the microcanonical ensemble (47), except the overall coefficient, which is $`\tau _1`$ in the present case. It is interesting that the invariant variable $`\tau _1`$ changes its sign at the different dimension $`𝒟`$ from that for $`\mu _1`$. $`\tau _1`$ is negative for $`𝒟12`$ and positive for $`𝒟>12`$ (Table 2). Consequently, the free energy of the non-uniform phase becomes favorable at $`𝒟>12`$;
$`\rho _2>0(𝒟12),`$
$`\rho _2<0(𝒟>12).`$ (51)
It is remarkable that $`𝒟=13`$, which is the critical dimension in the microcanonical ensemble, is not special in the canonical ensemble.
### 4.2 Weakly charged smeared black branes
General thermodynamic properties of the charged (dilatonic) black branes are rather complicated compared with the vacuum cases. Besides, the general formula that is given in terms of the invariant quantities is untractable, so we focus attention on two limiting cases: a weakly charged case and a near-extremal case. The latter will be discussed in the next subsection.
Firstly, we consider weakly charged smeared black branes in a microcanonical ensemble (fixed mass and charge). Expanding the entropy difference for small charge $`QM`$, we find
$`{\displaystyle \frac{S_{\mathrm{NU}}S_\mathrm{U}}{S_\mathrm{U}}}\left[\sigma _2+\left({\displaystyle \frac{Q}{M_{\chi =0}}}\right)^2\delta \sigma _2+O(Q^4)\right]ϵ^4,`$
$`\delta \sigma _2={\displaystyle \frac{(d2)^2}{2(d3)^2}}\left\{\left[\tau _1+q_1{\displaystyle \frac{2(d4)}{d3}}\mu _1\right]q_1+\left({\displaystyle \frac{d5}{d3}}\mu _1\tau _1\right)\mu _1\right\},`$ (52)
where $`\mu _1`$ and $`\tau _1`$ are evaluated at $`\alpha =0`$, and we note that $`Q/M`$ is invariant under rigid rescalings. In Appendix B, we will show how Eq. (52) is derived. As in the vacuum case, the leading entropy difference appears at $`O(ϵ^4)`$, and no correction appears at a lower order. This is not due to an accidental cancellation under the approximation of weak charge, but we can show after a tedious calculation that leading corrections are always $`O(ϵ^4)`$ for general charge $`Q`$. The correction term $`\delta \sigma _2`$ is positive with numerical values of $`O(1)`$ for $`𝒟15`$, whereas $`\delta \sigma _2`$ becomes negative for a very large value of $`𝒟`$. (We have confirmed this up to $`𝒟=20`$.) In any case, the new correction is sufficiently small compared to $`\sigma _2`$, as long as the assumption of weak change is valid, and hence the charge and dilaton do not change the thermodynamic phase structure. One might think that the phase structure is drastically affected by the bulk fields with an arbitrary charge, because the correction $`\delta \sigma _2`$ counteracts the leading term $`\sigma _2`$ in the above case. As we will see below, however, the same result of (48) holds even for the near-extremal limit in a microcanonical ensemble.
Next, we turn to a canonical ensemble for weakly charged branes. The difference of the free energy for the same temperature and charge is
$`{\displaystyle \frac{F_{NU}F_U}{F_U}}\left[\rho _2+\left({\displaystyle \frac{Q}{M_{\chi =0}}}\right)^2\delta \rho _2\right]ϵ^4,`$
$`\delta \rho _2={\displaystyle \frac{(d2)^2}{2(d3)}}\left\{(d3)(s_1+\tau _1)\tau _1\left[q_1+2(d3)\tau _1\right]q_1\right\}.`$ (53)
See Appendix B for the derivation of Eq. (53). The dimensional dependence of $`\delta \rho _2`$ is depicted in Fig. 2. From the figure we see that $`\delta \rho _2`$ is negative for $`𝒟13`$, hence, the correction counteracts $`\rho _2`$. The correction in this weakly charged case is too small to change the phase structure; interestingly, the effect becomes evident near the extremality, as we will see in the next section.
### 4.3 Near-extremal smeared black branes
Let us finally consider the thermodynamic behavior near the extremal state. A near-extremal limit can be realized by taking $`\alpha `$ to infinity. We define energy above extremality $``$, “reduced temperature” $`𝒯`$ and “reduced entropy” $`𝒮`$ as follows:
$`𝒯\underset{\alpha \mathrm{}}{lim}\sqrt{Q}T=\left[{\displaystyle \frac{(d3)^3𝒢}{4A_hL^2R_0^2}}\left({\displaystyle \frac{LR_0}{2\pi }}\right)^{d3}\right]^{1/2},`$
$`𝒮\underset{\alpha \mathrm{}}{lim}{\displaystyle \frac{S}{\sqrt{Q}}}=\left[{\displaystyle \frac{16\pi ^2A_h𝒢}{(d3)}}\left({\displaystyle \frac{LR_0}{2\pi }}\right)^{d1}\right]^{1/2},`$
$`\underset{\alpha \mathrm{}}{lim}(MQ)=𝒢\left({\displaystyle \frac{LR_0}{2\pi }}\right)^{d3}\left[{\displaystyle \frac{d1}{2}}(d3)\chi \right],`$
$`_{\mathrm{el}}\underset{\alpha \mathrm{}}{lim}(M_{el}Q)=𝒢\left({\displaystyle \frac{LR_0}{2\pi }}\right)^{d3}{\displaystyle \frac{d3}{2}}.`$ (54)
Note that $`_{\mathrm{el}}=M|_{\alpha =0}`$ and $`𝒯𝒮=TS|_{\alpha =0}`$ hold. Using Eq. (46), we can rewrite the first law in terms of the reduced quantities near the extremal state:
$`d(R_0)=𝒯(R_0)d𝒮(R_0).`$ (55)
In addition, we notice that the Smarr formula also takes the same form as that for the uncharged case.
The invariant fractional changes of the reduced quantities are given by
$`{\displaystyle \frac{\delta \overline{\mu }}{\overline{\mu }}}{\displaystyle \frac{\delta }{}}+(d3){\displaystyle \frac{\delta K}{K}}=\overline{\mu }_1ϵ^2+O(ϵ^4),`$
$`{\displaystyle \frac{\delta \overline{\tau }}{\overline{\tau }}}{\displaystyle \frac{\delta 𝒯}{𝒯}}+{\displaystyle \frac{1}{2}}(d5){\displaystyle \frac{\delta K}{K}}=\overline{\tau }_1ϵ^2+O(ϵ^4),`$
$`{\displaystyle \frac{\delta \overline{s}}{\overline{s}}}{\displaystyle \frac{\delta 𝒮}{𝒮}}+{\displaystyle \frac{1}{2}}(d1){\displaystyle \frac{\delta K}{K}}=\overline{s}_1ϵ^2+O(ϵ^4),`$ (56)
where the leading-order coefficients are given by
$`\overline{\mu }_1`$ $`=`$ $`{\displaystyle \frac{2(d3)}{d1}}\chi _1+(d3){\displaystyle \frac{k_1}{k_0}},`$
$`\overline{\tau }_1`$ $`=`$ $`{\displaystyle \frac{1}{2}}a_{0,1}(1)+{\displaystyle \frac{1}{2}}(d5){\displaystyle \frac{k_1}{k_0}},`$
$`\overline{s}_1`$ $`=`$ $`{\displaystyle \frac{1}{2}}a_{0,1}(1)+{\displaystyle \frac{1}{2}}(d1){\displaystyle \frac{k_1}{k_0}}.`$ (57)
We consider first the thermodynamics in a microcanonical ensemble (fixed energy $``$). The entropy difference between the uniform and non-uniform solutions for the same energy $``$ is computed with the first law (55) as
$`{\displaystyle \frac{𝒮_{\mathrm{NU}}𝒮_\mathrm{U}}{𝒮_\mathrm{U}}}=\overline{\sigma }_2ϵ^4+O(ϵ^6),\overline{\sigma }_2={\displaystyle \frac{d1}{4(d3)}}\left[{\displaystyle \frac{d5}{2(d3)}}\overline{\mu }_1\overline{\tau }_1\right]\overline{\mu }_1.`$ (58)
The numerical coefficients of $`\overline{\mu }_1`$ and $`\overline{\tau }_1`$ differ from those of (47), due to the “corrections” $`_{\mathrm{el}}`$ to the mass in the Smarr formula (9), i.e., $`𝒯𝒮/\left(_{\mathrm{el}}\right)=2(d3)/(d1)`$ for the uniform phase ($`\chi =0`$). The entropy difference again arises at $`O(ϵ^4)`$ as the neutral case. As we can see from Table 3, the entropy of the non-uniform phase is less than the entropy of the uniform phase for a given energy $``$ for $`𝒟13`$. In other words, the non-uniform phase is disfavored thermodynamically over the uniform phase. For $`𝒟>13`$, on the other hand, the non-uniform phase is favorable over the uniform phase:
$`\overline{\sigma }_2<0(𝒟13),`$
$`\overline{\sigma }_2>0(𝒟>13).`$ (59)
The critical dimension at which the sign of $`\overline{\sigma }_2`$ changes is same as the one for the uncharged case.
We now turn to a canonical ensemble (fixed temperature $`𝒯`$). We define the “free energy” of the solution as
$`𝒯𝒮.`$ (60)
The free energy of the uniform solution is
$`_U(R_0)={\displaystyle \frac{d5}{d1}}(R_0),`$ (61)
which vanishes for $`d=5`$ and is negative (positive) for $`d>5`$ ($`d<5`$). The difference of free energy between the uniform and non-uniform solutions for the same temperature is given by
$`{\displaystyle \frac{_{\mathrm{NU}}_\mathrm{U}}{_\mathrm{U}}}=\overline{\rho }_2ϵ^4+O(ϵ^6),\overline{\rho }_2={\displaystyle \frac{d1}{d5}}\left({\displaystyle \frac{1}{2}}\overline{\mu }_1{\displaystyle \frac{d3}{d5}}\overline{\tau }_1\right)\overline{\tau }_1,`$ (62)
where $`𝒟6`$ ($`d5`$). The exceptional case of $`𝒟=6`$ has to be treated separately because the free energy of the uniform solution vanishes. This case will be discussed elsewhere (see Ref. ). Table 3 summarizes our results. We see that $`\overline{\rho }_2`$ is positive for $`𝒟=5`$ and $`𝒟>14`$. This change of sign at the lower dimensions is in contrast to the previously studied cases in which the change of sign occurs only at relatively higher dimensions. However, we should recall that the free energy $`_\mathrm{U}`$ of the uniform solution changes its sign depending on the effective dimension. Consequently, the non-uniform phase becomes thermodynamically favorable over the uniform phase only at $`𝒟>14`$:
$`_\mathrm{U}<_{\mathrm{NU}}(𝒟14),`$
$`_\mathrm{U}>_{\mathrm{NU}}(𝒟>14).`$ (63)
This result differs from the vacuum case (51). We have observed in the weakly charged case that the correction due to the charge has a tendency to change the phase structure. The large correction near the extremality shifts the critical dimension upward, although the correction cannot change the phase structure significantly.
## 5 Conclusion
We have investigated the uniform and non-uniform black branes smeared on a transverse circle by constructing perturbative solutions explicitly. We have made use of the HO metric ansatz in our perturbation analysis. In the original perturbation scheme based on the general metric ansatz, analysis requires a good deal of subtle maneuvering. Making use of the HO metric has made the analysis enormously simple and clear in many respects. First of all, the EOMs of the charged dilatonic black branes are the same as those of the vacuum black branes on a circle. Consequently, we do not need to perturb the dilaton and form fields at any order, and the task is reduced to solving only one master equation at each order. Owing to such simplicity, the asymptotic behavior of the perturbation and the (analytically approximate) GL critical wave number are easily obtained. Besides, all physical quantities are characterized by a single parameter $`a_{0,1}`$.
The existence of a static mode means an onset of unstable modes, hence, the smeared black branes suffer from the GL instability irrespective of their charge. However, this does not mean that they provide a counterexample of the GM conjecture since there is no parameter region in which local thermodynamic stability is realized, except at the extremality.
Having performed the perturbation up to the third order, we obtain the phase structure around the critical uniform solution. We have begun with studying vacuum black branes. In the microcanonical ensemble, the thermodynamic system has the critical dimension $`𝒟_{}=Dp=13`$, at which the nature of phase transition between the uniform and non-uniform branes changes . On the other hand, we found that the critical dimension changes if we change the ensemble. In the canonical ensemble, the critical dimension is $`𝒟_{}=12`$. Furthermore, if the system is near the extremal state, the critical dimension in the canonical ensemble becomes $`𝒟_{}=14`$, whereas the critical dimension in the microcanonical ensemble does not change. We think that a specific value for a critical dimension is not universal. It depends on matters in the bulk and type of ensemble.
While the third-order perturbations have been performed to obtain meaningful leading-order corrections, the next-order corrections are also interesting to study. To obtain such corrections, we need to perform the perturbations up to the fifth order, and it would appear possible to obtain them. We think that the perturbative results in this paper give a good theoretical basis for understanding and constructing fully nonlinear solutions of the HO ansatz. Fully nonlinear solutions will allow us to explore the thermodynamic phase structure for an arbitrary charge, which might have significant consequences. We will discuss the fully nonlinear (numerical) solutions in a forthcoming paper. It is clearly interesting if a critical dimension appears at much lower dimensions, because there is no sensible quantum theory of gravity for $`D>10,11`$. Such a situation might be possible if we consider other types of bulk fields or configurations.
###### Acknowledgments.
H. K. is supported by the JSPS. U. M. is partially supported by a Grant for The 21st Century COE Program (Holistic Research and Education Center for Physics Self-Organization Systems) at Waseda University.
## Appendix A Vacuum Einstein equations
In this appendix we give explicit representation of the Einstein equations. From the ansatz (3), the metric of the ($`d+1`$)-dimensional vacuum black branes on a circle is
$`ds_{d+1}^2=fdt^2+\left({\displaystyle \frac{L}{2\pi }}\right)^2\left[{\displaystyle \frac{e^a}{f}}dR^2+e^{a(d2)b}dv^2+e^bR^2d\mathrm{\Omega }_{d2}^2\right],`$ (64)
where we have taken $`A(R,v)=e^a`$ and $`B(R,v)=e^b`$. The nontrivial components of the vacuum Einstein equations $`_j^i=0`$ are given by
$`_v^v`$ $`=`$ $`{\displaystyle \frac{e^a}{2}}\left({\displaystyle \frac{2\pi }{L}}\right)^2(f[(d2)b^{\prime \prime }a^{\prime \prime }]e^{(d2)b}[(d2)\ddot{b}+\ddot{a}]`$ (66)
$`+{\displaystyle \frac{1}{2R}}\{2(d3+f)[(d2)b^{}a^{}]e^{(d2)b}R(d2)(d1)\dot{b}^2\}),`$
$`_R^R`$ $`=`$ $`{\displaystyle \frac{e^a}{2}}\left({\displaystyle \frac{2\pi }{L}}\right)^2(fa^{\prime \prime }+\ddot{a}e^{(d2)b}`$ (68)
$`+(d2)\dot{a}\dot{b}e^{(d2)b}+(d2)f\{[{\displaystyle \frac{2}{R}}a^{}+{\displaystyle \frac{(d1)b^{}}{2}}]b^{}{\displaystyle \frac{a^{}}{R}}\}),`$
$`_\theta ^\theta `$ $`=`$ $`{\displaystyle \frac{e^a}{2}}\left({\displaystyle \frac{2\pi }{L}}\right)^2[{\displaystyle \frac{2(d3)}{R^2}}(1e^{ab})`$ (69)
$`+{\displaystyle \frac{1}{R}}(d3+f)b^{}+e^{(d2)b}(d2)\dot{b}^2+fb^{\prime \prime }+e^{(d2)b}\ddot{b}],`$
$`_v^R`$ $`=`$ $`{\displaystyle \frac{e^a}{4}}\left({\displaystyle \frac{2\pi }{L}}\right)^2({\displaystyle \frac{1}{R}}\{d3+f[d1+(d2)Rb^{}]\}\dot{a}`$ (71)
$`+(d2)f[a^{}{\displaystyle \frac{2}{R}}(d1)b^{}]\dot{b}2f(d2)\dot{b}^{}),`$
where the prime and dot stand for the derivatives with respect to $`R`$ and $`v`$, respectively. Since the angular component $`_\theta ^\theta `$ does not contain any derivatives of $`a(R,v)`$, $`a(R,v)`$ can be given in terms of $`b(R,v)`$ and its derivatives.
## Appendix B Entropy and free energy difference
In this appendix, we derive the formulae of the entropy difference and free energy difference, i.e., Eqs. (47), (50), (52) and (53).
First, we expand the thermodynamical quantities by $`ϵ`$ as follows:
$`M={\displaystyle \underset{p=0}{\overset{\mathrm{}}{}}}M_pϵ^{2p},T={\displaystyle \underset{p=0}{\overset{\mathrm{}}{}}}T_pϵ^{2p},S={\displaystyle \underset{p=0}{\overset{\mathrm{}}{}}}S_pϵ^{2p},Q={\displaystyle \underset{p=0}{\overset{\mathrm{}}{}}}Q_pϵ^{2p}.`$ (72)
The entropy difference between the uniform and non-uniform branes of the same mass and charge is expressed as
$`{\displaystyle \frac{S_{\mathrm{NU}}S_\mathrm{U}}{S_\mathrm{U}}}`$ $`=`$ $`{\displaystyle \frac{_{p=0}^{\mathrm{}}S_pϵ^{2p}}{S_0+\mathrm{\Delta }S_0}}1`$ (73)
$`=`$ $`{\displaystyle \frac{S_1}{S_0}}ϵ^2+{\displaystyle \frac{S_2}{S_0}}ϵ^4{\displaystyle \frac{S_1}{S_0}}{\displaystyle \frac{\mathrm{\Delta }S_0}{S_0}}ϵ^2{\displaystyle \frac{\mathrm{\Delta }S_0}{S_0}}+\left({\displaystyle \frac{\mathrm{\Delta }S_0}{S_0}}\right)^2+O\left(ϵ^6\right),`$
where $`\mathrm{\Delta }S_0`$ is the entropy change of a uniform brane due to the change of the mass and charge, which we denote by $`\mathrm{\Delta }M_0`$ and $`\mathrm{\Delta }Q_0`$, respectively.
Let us focus on the quantities $`S_1/S_0`$ and $`S_2/S_0`$ in Eq. (73). From the first law with a fixed scale of the circle ($`\delta L=0`$), we have
$`M_1=T_0S_1+Q_1\mathrm{tanh}\alpha ,M_2=T_0S_2+{\displaystyle \frac{1}{2}}T_1S_1+Q_2\mathrm{tanh}\alpha .`$ (74)
From Eq. (5), the mass of a uniform brane is written in two ways:
$`M_0=\left({\displaystyle \frac{d2}{d3}}+\mathrm{sinh}^2\alpha \right)T_0S_0\text{and}M_0=\left[1+{\displaystyle \frac{d2}{(d3)\mathrm{sinh}^2\alpha }}\right]Q_0\mathrm{tanh}\alpha .`$ (75)
From these equations (74) and (75), we can write the entropy difference, $`S_1/S_0`$ and $`S_2/S_0`$, in terms of $`M_1/M_0`$, $`M_2/M_0`$, $`Q_1/Q_0`$, $`Q_2/Q_0`$ and $`T_1/T_0`$.
Next, we focus on the quantity $`\mathrm{\Delta }S_0/S_0`$ in Eq. (73). To express the entropy of uniform brane in terms of mass and charge, we solve Eq. (5) for $`\alpha `$ and $`R_0`$ with $`\chi =0`$,
$`\mathrm{sinh}^2\alpha `$ $`=`$ $`{\displaystyle \frac{1}{2(1Q_0^2/M_0^2)}}\left[1+2{\displaystyle \frac{d2}{d3}}\left({\displaystyle \frac{Q_0}{M_0}}\right)^2+\sqrt{1+{\displaystyle \frac{4(d2)}{(d3)^2}}\left({\displaystyle \frac{Q_0}{M_0}}\right)^2}\right],`$
$`{\displaystyle \frac{LR_0}{2\pi }}`$ $`=`$ $`\left[{\displaystyle \frac{𝒢(d3)}{M_0}}\left({\displaystyle \frac{d2}{d3}}+\mathrm{sinh}^2\alpha \right)\right]^{1/(d3)}.`$ (76)
Substituting above two relations into the entropy of the uniform brane ($`A_h=1`$) in Eq. (5), one obtains in the leading order of $`Q_0/M_0`$ ($`1`$)
$`S_0(M_0,Q_0)4\pi 𝒢\left[{\displaystyle \frac{M_0}{(d2)𝒢}}\right]^{(d2)/(d3)}\left[1{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{d2}{d3}}\right)^2\left({\displaystyle \frac{Q_0}{M_0}}\right)^2\right].`$ (77)
From this expression, we can compute $`\mathrm{\Delta }S_0=S(M_0+\mathrm{\Delta }M_0,S_0+\mathrm{\Delta }S_0)S(M_0,S_0)`$. Taking the mass and charge differences as $`\mathrm{\Delta }M_0M_1ϵ^2+M_2ϵ^4`$ and $`\mathrm{\Delta }Q_0Q_1ϵ^2+Q_2ϵ^4`$, we obtain $`\mathrm{\Delta }S_0/S_0`$ in terms of $`M_1/M_0`$, $`Q_1/Q_0`$ and so on.
Finally, replacing such gauge dependent quantities with gauge independent quantities such as $`\mu _1`$ with $`\delta K=0`$, we obtain the expressions (47) and (52).
In a similar way, we can derive the formula of the free energy difference, Eqs. (50) and (53). The free energy difference between uniform and non-uniform branes of the same temperature and charge is given by
$`{\displaystyle \frac{F_{\mathrm{NU}}F_\mathrm{U}}{F_\mathrm{U}}}`$ $`=`$ $`{\displaystyle \frac{_{p=0}^{\mathrm{}}F_pϵ^{2p}}{F_0+\mathrm{\Delta }F_0}}1`$ (78)
$`=`$ $`{\displaystyle \frac{F_1}{F_0}}ϵ^2+{\displaystyle \frac{F_2}{F_0}}ϵ^4{\displaystyle \frac{F_1}{F_0}}{\displaystyle \frac{\mathrm{\Delta }F_0}{F_0}}ϵ^2{\displaystyle \frac{\mathrm{\Delta }F_0}{F_0}}+\left({\displaystyle \frac{\mathrm{\Delta }F_0}{F_0}}\right)^2+O\left(ϵ^6\right),`$
where $`\mathrm{\Delta }F_0`$ means the free energy change of a uniform brane due to the change of temperature and charge. From the first law with a fixed scale of the circle, we have
$`F_1=S_0T_1+Q_1\mathrm{tanh}\alpha ,F_2=S_0T_2+{\displaystyle \frac{1}{2}}S_1T_1+Q_2\mathrm{tanh}\alpha .`$ (79)
From Eq. (75), we can express the free energy of a uniform brane in two ways:
$`F_0\left({\displaystyle \frac{1}{d3}}+\alpha ^2\right)T_0S_0\text{and}F_0\alpha ^1\left({\displaystyle \frac{1}{d3}}+\alpha ^2\right)Q_0(\alpha 1).`$ (80)
From Eqs. (79) and (80), we can write the free energy difference, $`F_1/F_0`$ and $`F_2/F_0`$, in terms of $`T_1/T_0`$, $`T_2/T_0`$, $`Q_1/Q_0`$, $`Q_2/Q_0`$ and $`S_0/S_1`$.
To calculate $`\mathrm{\Delta }F_0/F_0`$, we have to express the free energy of uniform branes as a function of temperature and charge. From Eq. (5), we have
$`\alpha `$ $`=`$ $`{\displaystyle \frac{Q_0}{(d3)𝒢}}\left({\displaystyle \frac{4\pi T_0}{d3}}\right)^{d3}+O\left(Q_0^3\right),`$
$`{\displaystyle \frac{LR_0}{2\pi }}`$ $`=`$ $`{\displaystyle \frac{d3}{4\pi }}\left[1{\displaystyle \frac{1}{2}}{\displaystyle \frac{Q_0^2}{(d3)^2𝒢^2}}\left({\displaystyle \frac{4\pi T_0}{d3}}\right)^{2(d3)}\right]+O\left(Q_0^4\right).`$ (81)
With these relations, the free energy of a uniform brane is given by
$`F_0(T_0,Q_0)=𝒢\left({\displaystyle \frac{d3}{4\pi T_0}}\right)^{d3}\left[1+{\displaystyle \frac{Q_0^2}{2(d3)𝒢^2}}\left({\displaystyle \frac{4\pi T_0}{d3}}\right)^{2(d3)}\right].`$ (82)
From this, one can compute the free energy difference of a uniform brane as $`\mathrm{\Delta }F_0=F_0(T_0+\mathrm{\Delta }T_0,Q_0+\mathrm{\Delta }Q_0)F_0(T_0,Q_0)`$. Taking the change of temperature and charge as $`\mathrm{\Delta }T_0T_1ϵ^2+T_2ϵ^4`$, $`\mathrm{\Delta }Q_0Q_1ϵ^2+Q_2ϵ^4`$, we can write $`\mathrm{\Delta }F_0/F_0`$ in terms of $`S_1/S_0`$, $`T_1/T_0`$ and so on. Therefore, replacing these with gauge independent quantities, we obtain Eqs. (50) and (53). |
warning/0506/astro-ph0506705.html | ar5iv | text | # A Medium Resolution Near-Infrared Spectral Atlas of O and Early B Stars
## 1 INTRODUCTION
Kleinmann & Hall (1986) were the first to present reasonably-high resolution, high signal-to-noise (S/N) near-infrared (NIR) spectra for cool stars. The first NIR spectral atlas of hot stars was given by Lançon & Rocca-Volmerage (1992). Designed for use in stellar population synthesis models, the Lançon & Rocca-Volmerage atlas lacked adequate resolution for applications in many stellar and galactic programs. A few years later, Dallier et al. (1996) and Hanson et al. (1996) presented $`H`$-band and $`K`$-band spectral atlases, respectively, which included OB stars with significantly higher resolution and S/N. Numerous NIR atlases of OB stars have been published since that time (Wallace & Hinkle 1997, Lenorzer et al. 2002, for a recent review of all NIR spectral atlases, see Ivanov et al. 2004). The utility of a NIR spectral classification scheme, for hot stars in particular, has proved exceedingly useful for a variety of applications, including studies of very young star forming regions (Bik et al. 2003) and the galactic center region (Najarro et al. 1997, Ghez et al. 2003; Najarro et al. 2004) as well as distant massive clusters through out the Galaxy (Hanson, Conti & Howarth 1997, Blum, Damineli & Conti 1999; Figer et al. 2005). Furthermore, researchers studying heavily reddened high-mass X-ray binary systems (Clark et al. 2003; Morel & Grosdidier 2005) and microquasars (Mirabel et al. 1997, Martì et al. 2000) have found NIR spectral classification to be uniquely valuable.
In light of these successes, our group wishes to push NIR spectral studies of OB stars to a new, more sophisticated level. Our goal is to obtain new, higher resolution and S/N NIR OB spectra to test and guide existing quantitative atmospheric models for OB stars in the NIR regime (Najarro et al. 1998; Kudritzki & Puls 2000 and references therein). In turn, once the atmospheric models have been calibrated to properly predict stellar characteristics based on the NIR spectra of known, UV- and optically-studied stars, it is our hope that they may be used to provide accurate constraints to the characteristics of stars observable only in the NIR. This NIR atlas of well known, optically visible OB stars makes up the sample of high-quality spectra which are being used by our group for a successful NIR qualitative analysis (Repolust et al. submitted).
## 2 OBSERVATIONS
A list of the stars used for the survey, their position, and salient details of the observations is given in Table 1. When our observations were carried out, reasonably-high-resolution ($`R>8000`$) spectrometers which allowed for sufficient spectral coverage were only available on 8- and 10-meter class telescopes. Our targets are exceedingly bright for such a large aperture system. However, these observations are absolutely necessary for the sake of developing and testing quantitative model atmospheres. Furthermore, we needed to start with optically visible, well known O and early-B stars. The OB sample was selected to give reasonable coverage of the temperature and luminosity range of O and early-B stars. The temperature and luminosity range sampled is illustrated in the spectral type versus luminosity class presentation given in Table 2.
There presently are no classification standards in the NIR. Until sufficient numbers of OB stars have been observed in the NIR, it will be too soon to claim any star as a classification standard. The stars selected in this survey are not being promoted as standard stars for spectral classification. The targets for this program were selected based on entirely different criteria. The OB supergiant stars observed with the VLT were obtained as part of a study of OB supergiants near the galactic center (Fickenscher, Hanson & Puls, 2004). For ease of the observing run, this required them to be in the vicinity of the galactic center fields being observed. Most of the stars observed with Subaru were hand selected by one of us (J.P.). The Subaru sample was selected to cover a reasonable sampling of effective temperature and gravity, and most importantly, with the expressed desire to model their spectral profiles using modifications to the atmospheric code FASTWIND (Puls et al. 2005). Most of the stars selected have already undergone significant previous spectroscopic modeling in the optical and UV, and did not show any serious irregularities in those analyzes. The very high resolution of these observations are well beyond those typically used for classification in the optical (R $``$ few thousand). Spectra for the purpose of classification, no matter the wavelength, are best obtained at more moderate resolutions (see further discussion of this point in §5.2).
### 2.1 The VLT-ISAAC Spectra
The first of our data came from the Infrared Spectrometer and Array Camera (ISAAC). The instrument is mounted at Nasmyth focus on the 8.2m Unit 1 telescope of the European Southern Observatory’s (ESO) Very Large Telescope (VLT), located on Cerro Paranal in Atacama, Chile (Moorwood 1997). ISAAC employs a Rockwell Hawaii 1024<sup>2</sup> array and a single grating. The resolution is set by the slit width. In May and April of 2001, long slit (120”) $`H`$ and $`K`$band spectra were obtained for a number of optically visible late-O and early-B supergiants. All data for this program were obtained in the queue observing mode, spread out over about two months and using approximately 12 hours of VLT queue time. Typically 8 slit positions were obtained, with total on-source integration times of one to a few minutes. A slit-width of 0.3” was used, giving a spectral resolving power of $`R10,000`$ in the $`H`$band and $`R8,000`$ in the $`K`$band. Three grating settings, 1.710, 2.085, and 2.166 $`\mu `$m, were used with ISAAC.
### 2.2 The Subaru-IRCS Spectra
Later that year, we obtained additional spectra at the 8.2-m Subaru Telescope, operated by the National Astronomical Observatory of Japan (Tokunaga et al. 1998) and located at the top of Mauna Kea in Hawaii. We employed the Infrared Camera and Spectrograph (IRCS) installed at Cassegrain focus. IRCS uses two 1024<sup>2</sup> ALADDIN arrays and offers a cross-dispersed echelle mode for high resolution work. With the present arrays, the cross-dispersed mode does not allow for full spectral coverage, leaving unobserved spectral regions between the echelle orders. For full coverage, two settings can be completed. However, nearly all the important lines for our survey (except the $`\lambda 2.0581\mu `$m HeI line) were observed with a single grating setting.
To achieve the highest resolution, we used the most narrow slit setting (0.15”). This provided a resolution of about 0.5 Å/pixel in $`K`$, and about 0.4 Å/pixel in $`H`$. In practice, we achieved an approximate FWHM of 3.5 pixels, measured through the OH sky emission features, resulting in a final spectral resolution of approximately $`R12,000`$. These spectra were obtained over two separate runs, first in November 2001, then later in July 2002, based on the targets’ right ascension. Two nights were granted for each run. The first night was used to obtain just the $`K`$band spectra, the second night was reserved for the $`H`$band spectra for the same stars. The weather in November 2001 was rather poor; clouds on the mountain had closed just about every other optical observatory. Because our sources are so bright (with $`K`$ magnitudes as large as 2), we managed to get sufficient counts despite the weather. Telluric corrections did prove to be more problematic, however, because of the poor sky conditions. Conditions in July 2002 were also less than ideal, but again, given the brightness of our targets and the 8-m aperture we obtained sufficient counts. Intermittent heavy clouds did lead to some poor telluric corrections in the final spectra for some sources. We failed to get a follow up $`H`$ band spectrum of HD 15558 during the second night. Also, the $`H`$band spectrum of $`\tau `$ Sco, HD 149757, was somehow corrupted. Regrettably, a trustworthy spectrum proved irretrievable from our raw data.
## 3 REDUCTION OF THE SPECTRA
All VLT-ISAAC spectra were reduced using IRAF<sup>1</sup><sup>1</sup>1IRAF is distributed by the National Optical Astronomy Observatories, which are operated by the Association of Universities for Research in Astronomy, Inc., under cooperative agreement with the National Science Foundation. routines. Subsequent analysis of the Subaru spectra was done using routines written in PerlDL<sup>2</sup><sup>2</sup>2homepage at http://pdl.perl.org/.
The OB stars and a few telluric standards (more on that presently) were observed with the traditional long-slit AB “nodding” pattern on both telescopes (with a small, few arcsecond jitter to better sample the detector). This allows quick and effective subtraction of background emission when the A slit position is subtracted from the B slit position several arcseconds away (and vice versa). For all stars observed with the VLT, four AB pairs were obtained, giving us 8 individual spectra. For the November 2001 Subaru run, three ABBA sets were used. This was reduced to just 2 ABBA sets for the July 2002 run. Each two-dimensional spectral image was bias (and other instrumental effects) subtracted before flat fielding. For all spectra, regardless of telescope, dark frames and flat field frames were averaged together to form a master dark and flat frame. Then AB pairs are subtracted to remove the last of the sky emission. Finally, each individual spectrum from the set were extracted, weighted and scaled before being averaged, using a 3-sigma rejection. Telluric OH emission lines (Rousselot et al. 2000), ubiquitous throughout the spectral range of our data, were used for wavelength calibrations.
The Earth’s atmosphere introduces a myriad of absorption lines into ground-based, NIR spectra, and these need to be removed (Fig. 1). Because the features are complex and continuously changing, in situ measures, via a telluric standard star, are the best means of constraining their character. Owing to their near featureless continuum (with the very big exception of the Brackett series of Hydrogen), late-B and early-A dwarf stars are frequently used in this regard. For each target star observed with the VLT-ISAAC, a telluric standard star was observed either just before or after the target object. This telluric star was selected to exhibit the same airmass (to within a few hundredths usually) as the target object, and to lie in the same general sky direction of the target object. AB nodding along the slit was used to observe the telluric stars. Identical procedures were used for the reduction of the telluric stars as was used for the OB stars as outlined above. For the VLT run, we observed a late-B dwarf, matched to each OB target star and put into individual observing blocks (as required for VLT queue observing). We also obtained spectra of a few early-G dwarfs in some of the observing blocks. These early-G dwarfs are used with the NIR spectrum of the Sun to help constrain the Brackett features in the B-dwarf star. Once the Brackett features have been modeled in the late-B dwarf, they can be divided from the late-B dwarf spectrum to create the telluric spectrum. The telluric spectrum is then finally divided from the raw OB target star spectrum (see the Appendix in Hanson et al. 1996).
For the Subaru run, observations were taken in standard visitor mode. A number of early-A dwarfs were observed through out the night on the assumption that they would be used in the removal of telluric features. In the end, however, they had very limited use. For OB stellar lines far from the influence of any Brackett series transitions (HeII, most of the HeI lines, etc.) the A stars were useful to derive the telluric spectrum, but the Subaru spectra had unprecedented S/N and even higher resolution than the VLT-ISAAC spectra. Consequently, any estimated fit for the Hydrogen Brackett lines in the A dwarf stars could not be adequately constrained to derive the OB star hydrogen lines. Ironically, the most ideal telluric standards would be very hot stars, because they have the fewest, weakest lines of all normal stars. In the end, we were faced with using the OB stars themselves, taken throughout the night, as telluric standards for each other in deriving the Hydrogen Brackett lines. In the infrared, every spectrum is like an equation with two unknowns: the telluric spectrum and the stellar spectrum. Taking additional spectra of other stars does not solve this, even if the telluric spectrum is the same, since every new spectrum is a new equation, but adds an additional unknown (the underlying spectrum of the new star). This is the case ONLY in the Brackett region because every star shows a unique Brackett feature. At some point, one is forced to make an assumption about the Brackett line profile in one of the stars (it is our experience that this is preferred over trying to make an assumption about the strength and shape of the numerous telluric features that span the Brackett spectral range). This reduces the number of unknowns to equal the number of spectra (‘equations’) and allows for a solution. Typically, a Voigt profile was used as a first, reasonable guess for the most well behaved stars observed during the night: late-O or early-B dwarfs. This ‘solution’ was then propagated through all the stars observed that night, dividing one against another until one finally resolves the underlying spectrum for all the stars from the night. If problems or inconsistencies are revealed, then a new assumption can be made and propagated through the stars for that night.
### 3.1 The need for synthesized spectra
Is there any way to know if our assumptions made to model our most ’normal’ stars are reasonable? What is our likely error? This would be very difficult to determine if we had no idea of the expected profile shapes for our OB stars. Luckily, in conjunction (and as a parallel goal) with the development of this high S/N atlas, our group is also expanding existing sophisticated models of OB stellar atmospheres to predict and produce line profiles in the NIR (Repolust et al. 2005). Once a best guess was made and the underlying spectrum of all the stars was completed, these were compared with synthesized spectra our group created. If our assumptions about the Brackett profile are wrong, then very obvious patterns of inconsistencies would appear in the comparison of data and models. Indeed, it was found that the Br10 line of Hydrogen appeared to be consistently too weak (as compared to the model predictions) for all stars observed in November 2001. We returned to the data, used a stronger line for the primary assumption, and then this solution was propagated through out the night. We found that a more satisfactory fit was obtained for all the star in this particular line after having done this.
But still, what is the expected error in our profiles? Given the brightness of our target (and telluric) stars, flux is not so much a limit to our confidence (indeed we have very high S/N) it is the accuracy of the telluric corrections and the assumptions made in fitting the Brackett features. For the November 2001 ‘inconsistency’, the strength of the Br10 line was increased by approximately 1.5% compared to the continuum. This made a significant difference in the Br10 line as seen in the O4 supergiant, HD 14947 (it was the weakness of the Br10 line in HD 14947 which exposed the problem). But such a small difference was essentially undetectable in the Br10 spectrum of the B2 dwarf, HD 36166. Given the goodness of the fits to the synthesized spectra and the reproducibility of the features using the telluric A-star standards, we expect the HeI and HeII lines to be good to at least 0.5% of the continuum. The Brackett lines pose special problems, as we’ve outlined. We estimate those lines to be good to between 0.5% and 1% through out the wings, and perhaps only good to 1% to even as bad as 2% of the continuum in the line core.
Telluric corrections in the short-K band provide additional challenges. While there is no strong hydrogen line residing here, making early-A dwarfs near perfect telluric tracers between 2 and 2.15 $`\mu `$m, the Earth’s atmosphere becomes increasing problematic below 2.08 $`\mu `$m. Here, very strong and quickly changing telluric absorption is found, predominately due to the ro-vibrational transitions of CO<sub>2</sub>. Just the start of this very strong absorption is illustrated in Figure 1. The absorption continues to get deeper at shorter wavelengths not given in Figure 1. For a thorough discussion of the special problems posed in this spectral region, see Kenworthy & Hanson (2004). The S/N measures listed in Table 1 represent the average through out most of the spectral range covered for the spectrum, and does not consider that over the CO<sub>2</sub> region, the S/N can drop by 30% or more. Similar telluric noise can also be seen in the $`H`$band region, where telluric features become increasingly strong and problematic longward of 1.72 $`\mu `$m (Fig. 1).
## 4 THE SPECTRA
This atlas contains $`H`$\- and $`K`$-band spectra for 37 OB stars, ranging from O3 to B3, and sampling most luminosity classes over that spectral range. The spectra have been displayed in two ways. All spectra in the atlas are presented in Figures 2 through 8. Here the spectra have been arranged based on temperature so the reader can see the temperature-dependent variations. In Figures 9 through 12, we have presented the spectra grouped by similar temperature, but arranged along varying luminosity. Here the reader can see the luminosity-dependent variations for four differing temperature regimes.
### 4.1 Line Identifications
The NIR is home to the Brackett Hydrogen series, beginning with Brackett $`\alpha `$ at 4.052 $`\mu `$m (all wavelengths listed are given for air). Over the spectral range covered in this atlas, we see $`2.1661\mu `$m (4-7) Brackett $`\gamma `$ (Br$`\gamma `$), 1.736 $`\mu `$m (4-10) Br10, $`1.681\mu `$m (4-11) Br11 and $`1.641\mu `$m (4-12) Br12.
Lines of HeI and HeII both exist in the $`H`$\- and the $`K`$-band regions. Within the $`H`$-band, one finds absorption due to HeII $`\lambda `$1.693 (7-12) and HeI $`\lambda `$1.700 (3$`p^3`$P$`{}_{}{}^{o}4d^3`$D, triplet). Within the $`K`$-band, there exists the lines, HeII $`\lambda `$2.188 (7-10), HeI $`\lambda `$2.1127 ($`3p^3`$P$`{}_{}{}^{o}4s^3`$S, triplet) and $`\lambda `$2.1138 ($`3p^1`$P$`{}_{}{}^{o}4s^1`$S, singlet). Though our spectral coverage from Subaru does not include this line, still another important He line is the singlet HeI transition, $`\lambda `$2.0581 ($`2s^1`$S$`2p^1`$P<sup>o</sup>). This line is seen in several of our late-O, early-B supergiants observed with ISAAC. The higher resolution and S/N obtained in this atlas allows additional HeI lines to be identified. In late-O and early-B supergiants (Fig. 5, 7, 8, 11 & 12) we see HeI $`\lambda `$2.161 ($`4d^3`$D$`7f^3`$F, triplet) and at $`\lambda `$2.162 ($`4d^1`$D$`7f^1`$F, singlet) begins to appear in the blue wing of Br$`\gamma `$.
After Helium and Hydrogen, most OB stars show relatively few spectral features in their NIR spectra. Among the hottest stars, the triplet due to CIV ($`3d^2D3p^2P^o`$) at 2.069, 2.078 and 2.083 is seen. In our O3 supergiant, Cyg OB2 #7, the leading line in this set, at 2.078$`\mu `$m, is seen in absorption! A second important metal line seen in the NIR spectra of hot stars is the broad emission feature found at 2.1155$`\mu `$m, tentatively identified by Hanson et al. (1996) as NIII (7-8), though may instead (or also) be due to the very similar transition of CIII (7-8). The difficulty in firmly identifying this features comes from the fact that the atoms share a similar structure and the transition identified here originates between very high lying levels. It is expected, based on arguments of relative abundances between N and C, that the feature seen is dominated by NIII, particularly in more evolved stars. Stars where CIII may dominate the profile, would include hotter, less evolved stars which also exhibit strong CIV emission. One can find additional discussion of the 2.1155$`\mu `$m feature in Repolust et al. (2005).
Additional lines previously cataloged in hot stars and seen in our spectra include: $`\lambda `$2.100$`\mu `$m NV seen in HD 64568, an O3 V((f)), $`\lambda 2.137/2.143\mu `$m MgII in early-B supergiants, and numerous HeI lines at $`\lambda 2.150\mu `$m, $`\lambda 2.161/2\mu `$m, and $`\lambda 2.184\mu `$m, all seen in late-O and early-B supergiants (see Fig. 8). Because of the high resolution and S/N of our spectra, a few new, unidentified lines have been found. These include absorption lines seen in mid and late-O stars at $`\lambda 1.649,1.651\mu `$m, an emission feature seen in the O5 If+ star HD 14947 at $`\lambda 2.1035\mu `$m, and several emission features seen clustered around the Br10 line of the O3 If\* star, Cyg OB2 # 7. The strongest of these, centered around $`\lambda 1.735\mu `$m, might also be seen in the O4 I(n)f star HD 66811.
### 4.2 Classification Criteria for Temperature and Luminosity
As with the optical (Walborn & Fitzpatrick 1990), the principal temperature classification criteria for O stars in the NIR is found within the behaviors of the HeI and HeII lines. The strength of the He lines as a function of temperature has been well established based on earlier, lower-resolution work (Hanson et al. 1996; Blum et al. 1997; Hanson et al. 1998), which have quantified its behavior. Ionized Helium, the strongest of the NIR transitions being the HeII (7-10) line at $`\lambda `$2.1885, is present in absorption in essentially all O stars, regardless of luminosity class (though it is so weak in O9 V stars it would go undetected in lower quality spectra). By mid-O, neutral Helium emerges (see Fig. 2, 4) and is retained until early-B for dwarf stars (see Fig. 3), but as late as B7 or B8 in supergiants (see Hanson et al. 1996). When both HeII and HeI are present, temperature estimates are at their best.
The unique HeI line at $`\lambda `$2.0581$`\mu `$m is highly sensitive to temperature and wind properties. It first appears in absorption in mid- to late-O stars (see Fig. 5, 7), and is frequently seen in strong emission in early-B supergiants (see Fig. 8). Unfortunately, this particular study is unable to shed significant light on its behavior due to the small number of stars observed here. It should be noted that this line lies within a region of the Earth’s atmosphere where interfering telluric absorption is large, and accurate, high S/N profiles are challenging to obtain from the ground.
Unlike the Helium transitions which are unique to hot stars, hydrogen lines are instead ubiquitous to all stars of cosmic abundance. Despite being common, the hydrogen lines serve an important role in constraining characteristics of OB stars. Over the temperature range of this study, the cores of the Brackett transitions, particularly Br10 and Br11, are excellent indicators of gravity. Among stars of similar rotational velocity, stars with lower gravity (higher luminosity), are seen to have much deeper cores. This leads to a stronger, larger equivalent width in the Brackett lines in high luminosity O and early-B stars (see Figs. 9, 10, 11 and 12). While some of the HeI lines appear to follow a similar route, their profiles becoming deeper and sharper as luminosity increases in OB stars, much of this profile change is due to a change in rotational velocity (see for example, Fig. 11). This behavior of deepening of the cores in the upper level Hydrogen lines, is entirely explained theoretically. Repolust et al. (2005) explain that the core is simply responding to the Stark-profiles which are a strong function of electron density. Such an effect is only seen in the cores of hydrogen lines with large upper principal quantum numbers. For Br$`\gamma `$, with upper quantum number n=7, the effect is less sensitive and it is entirely insensitive in H$`\alpha `$ where the cores are instead dominated by Doppler-broadening. Thus the profile shape of the cores in Br10 and Br 11 provide a sensitive indicator of gravity (luminosity), provided broadening from rotation is already constrained <sup>3</sup><sup>3</sup>3Unfortunately, a higher luminosity implies a larger mass-loss rate, so that the profiles might become refilled by wind-emission and thus weaker again. One such example is given in Fig. 9, if one compares Br12/11/10 from Cyg OB 8c (O5 If) with those from HD 14947 (O5 If+).. In high gravity OB stars, there is greater absorption in the wings of the higher order Brackett profiles (see in particular Fig. 12). It is important to recognize that without adequate resolution these subtle changes in Brackett $`\gamma `$ and Helium line profiles which trace luminosity would be lost.
The metal lines of CIV and NIII can be used as vague temperature indicators and are particularly useful with low S/N spectra or in spectra plagued with very strong nebular emission lines rendering the helium lines unusable. CIV appears in emission starting around O4-O5, and is seen down to about O7 in dwarf stars (though it is weak and may be missed with low quality spectra). This atlas lacks adequate sampling of the O8 spectral sequence in luminosity. However, the existence of CIV in O8 supergiants, and its increasing strength with increasing luminosity, was shown in Hanson et al. (1996). NIII is seen in all early-O stars, down to about O7 in the dwarfs, and possibly as late as O8 in supergiants, though it is heavily blended with the HeI $`\lambda `$2.112/3 line. In very luminous late-O stars, the triplet $`\lambda `$2.1127 ($`3p^3`$P$`{}_{}{}^{o}4s^3`$S) remains in absorption, while the singlet $`\lambda `$2.1138 ($`3p^1`$P$`{}_{}{}^{o}4s^1`$S) line goes into emission (Najarro et al. 1994). As previously discussed, a confident identification of the line at 2.1155 $`\mu `$m (whether it be CIII or NIII) is still lacking. For purposes of spectral classification, this is not an issue. For purposes of spectral analysis, it will be.
## 5 DISCUSSION
In the mid-90s, several papers presented low and moderate resolution NIR spectra of stars. Previous to spectroscopic studies, the only tool for determining the characteristics of heavily reddened stars was through broadband NIR colors. While such measures are effective for constraining cool stars, NIR photometric colors becoming increasingly degenerate for hot stars. It was pointed out by Massey et al. (1995) that UBV colors were degenerate for O stars, necessitating the need for MK classification for proper determinations of mass functions within OB clusters. For JHK colors, the degeneracy begins at A stars. The development of even a low resolution classification system for the NIR has proved enormously useful for those studying hot stars behind significant interstellar extinction.
However, with near-infrared spectrometers becoming more sophisticated, it was time to push the observations and theory to a complementary level. Our groups theoretical work shows that with increased resolution and S/N rather accurate physical characteristics of hot stars can be derived with NIR spectra alone (Repolust et al. 2005). But this does require an analysis be used in conjunction with theoretically derived line profiles to achieve such accurate characteristics.
### 5.1 What level of S/N and resolution will be required for a quantitative analysis?
It’s unlikely that typical researchers will have the luxury that this study possessed: using an 8-m class telescope to observe exceedingly bright, single sources, void of nebular contamination. In many astrophysical situations, nebular emission can contaminate important diagnostic lines (most usually Br$`\gamma `$ and $`\lambda 2.058\mu `$m HeI), rendering them useless. In the most extreme situations, where the extinction to the massive star is exasperated by thermal emission in the NIR from nebular or circumstellar material, all absorption lines may be undetectable. It is clear that both NIR spectral analysis and classification will have very real and frustrating limits in their application to massive young stellar objects (Hanson et al. 1997; Blum et al. 2001; Bik et al. 2005a). In incidences where there is significant contamination from thermal emission, this has an effect much like reducing the S/N. Imagine if one obtained a stellar spectrum with a S/N of 200, but half of the photons originate in a featureless continuum generated by a disk. The depth and strength of the stellar lines would be reduced by about one half. This would be roughly equivalent to obtaining the spectrum of an undiluted star, but with half the S/N. If it is believed that continuum contamination is roughly equal to the stellar flux, then improving S/N might allow the detection of the stellar lines. Spectroscopic systems designed to be used with adaptive optics would reduce the contamination from extended emission and increase the likelihood of detecting stellar features. In cases where there may be significant thermal contamination, classification should not be based on equivalent width measures alone, but by the relative behavior of critical line pairs (see §5.2).
Ignoring for a moment contamination from nebular or circumstellar emission, what advice might we give to those interested in doing a more accurate spectral analysis in the NIR? Unfortunately, the resolution and S/N limits required are a function of the stars spectral type, luminosity class and rotational velocity (something the researcher does not know at the start!) and whether one wishes to use a simple classification scheme or a full blown profile analysis. This is realized when one looks closely at the line strengths among the O dwarfs. In Figure 13, we plot the central region of the $`H`$band, which happens to contain a Hydrogen line as well as neutral and ionized Helium. To create this figure, we created a set of low S/N line-free spectral regions from our data. These have been multiplied against are true spectra to illustrate the effect of reduced S/N on the detection of these weak stellar features. All the spectral features shown are $`<`$ 2 Å in strength, the strongest line being the HeI line at 1.70 $`\mu `$m in the O9 V star HD 149757 (e.w. = 1.8Å). Once the S/N was reduced to 150, the weak (0.45 Å) HeII feature at 1.693 $`\mu `$m in HD 217086, is no longer confidently detected. For nearly all early O stars, we suggest a S/N $`>`$ 150 for a quantitative analysis, to detect the very weak features and properly match the wings in the line. For late-O and B stars, such as in HD 149757, keeping S/N $`>`$ 100 should be sufficient for their slightly stronger lines. Crude classification can still be obtained with a S/N under 100, but just barely. Only the strongest lines, EW $`>`$ 2.0 Å can be confidently detected once the S/N drops to 50 (Fig. 13).
For a proper profile analysis (and to resolve blends important in diagnosing luminosity classes) a resolution of at least $`\lambda /\delta \lambda >5000`$ should be adhered to. Such a resolution is well matched for many if not most OB stars, which typically possess fairly high rotational rates ($`V\mathrm{sin}i>100`$ km/s). However, for OB stars with low rotation ($`V\mathrm{sin}i<100`$ km/s), a resolution of perhaps 8,000 or more will be needed to resolve their underlying profile (see, for instance, the slow rotator, HD 149438, Tau Sco, in Fig. 3). While seemingly high (by NIR standards!), these resolutions are really very low. Similar studies on cooler stars require much higher resolutions of several to many tens of thousands for a proper analysis (e.g., Luck & Heiter 2005).
### 5.2 The use of equivalent width for classification
In nearly all previous NIR spectral atlases, enormous tables and figures are given showing the equivalent width of strategic lines as a function of spectral type or luminosity class. Such measures have been critical to the development of a classification tool for the NIR. However, we do not present such tables here. The spectra in this atlas differ greatly from previous atlases, and thus dictates a different mind set for their use. We suggest that when working with high S/N and moderate to high resolution spectra, even without the use of quantitative profile analysis, one must adopt the philosophy outlined decades ago for the classification of stars in the optical. Here classification is based on the comparison of spectra; that between stars of known spectral and luminosity class and stars which are unknown. Important for a proper comparison, the resolution and approximate S/N of the stars being compared must be the same. What has been sorrily missing, however, in most NIR spectral classification programs has been spectroscopic standards. There isn’t a single, conscientious optical spectroscopist doing classification work that would dream of skipping the step of taking classification standards. This is for a field for which optical classification spectra must number in the millions. In the infrared, researchers are relying on the scant, 100 or so O stars observed to date, to make crude comparisons to their own, independently obtained NIR spectra. This level of crude judgment to derive temperature and luminosity should not be tolerated if NIR classification is to become a robust technique. Moreover, as mentioned in §5.1, thermal contamination, something most optical astronomers do not typically deal with, makes classification via equivalent width strengths a very risky method for some applications.
The importance of obtaining ones own set of classification standards was recently highlighted in a study by Bik et al. (2005b). They obtained high S/N, moderate resolution NIR spectra of ionizing stars of IRAS selected young star forming regions. Within the realm of this study, the researchers obtained spectra for a number of known O and early-B dwarf stars, using the identical spectroscopic set up and typical S/N obtained for their heavily reddened ionizing sources. They found a small but significant difference in the line strengths of their standard stars compared to previous lower resolution, lower S/N NIR atlases. Without the ability to calibrate with their own set of classification standards, they would have miss-classified most of their sources.
Looking again to the optical for guidance, most classifications criteria are based on the comparison of two lines in a spectrum. Typically, sub-classes are defined by when a set of lines are of equal strength, one line is stronger than the other, or when a line first appears or disappears, etc. (Jaschek & Jaschek 1987). While at first this method may appear crude, it in fact makes for a highly repeatable (from person to person) evaluation and is the crux to the success of the optical classification system. In the NIR we do not have such a wealth of diagnostic lines to allow for comparisons between oppositely behaving spectral features over all temperature and luminosity ranges of interest. The He lines do offer such a system to constrain temperature once O stars get cool enough to show HeI, and before they get too cool to show HeII (over the very narrow range of O7 to about O9.5). However, when including comparisons with other lines, CIV, Br$`\gamma `$, reasonable estimations appear possible. We have tried to outline just those comparisons within this paper (§4.2), though we admit, the field of NIR spectral classification is still very young. Still more observations are need to establish “typical” spectroscopic behavior in the NIR.
## 6 CONCLUSIONS
We present intermediate resolution (R $``$ 8,000 - 12,000) high S/N $`H`$\- and $`K`$-band spectroscopy of a sample of optically visible O and early-B stars. The purpose of this study is to better characterize OB spectral profiles in the NIR. We have also established some observational limits for researchers preparing to use a quantitative analysis to derive stellar temperature and luminosity with NIR spectra alone. In order to directly determine effective temperatures and log$`g`$ for individual stars, one needs to work with atmospheric model codes which rely on profile fits. Such programs have recently been developed and show great promise (Lenorzer et al. 2004; Repolust et al. 2005). For most stars, a S/N $`>100`$ and resolution of $`r5000`$ should be just sufficient. However, if the targets of interest turn out to be early-O, or have a very low rotational velocity, a higher S/N ($`>150`$) or resolution (R $`>8000`$), respectively, will need to be obtained.
When NIR spectral classification is sought, we strongly encourage researchers to obtain spectra of known OB stars for direct comparisons to their target stars. Until a time when equivalent “MK” standards are developed in the NIR, optically studied and thus known to be ’well-behaved’ stars which bracket the expected spectral and luminosity range of the targets should be sufficient for this purpose.
We gratefully acknowledge the Subaru and VLT Observatories for their support of our program. This research has made use of the NASA’s Astrophysics Data System Bibliographic Services and the SIMBAD database operated at CDS (Strasbourg, France). MMH and MAK gratefully acknowledge support for this program from the National Science Foundation under Grant AST-0094050 to the University of Cincinnati. Facilities: VLT:Antu (ISAAC), Subaru (IRCS)
## Appendix A Appendix
While not part of the original science goals of this program, the extraordinary spectral resolution and signal-to-noise achieved on all stars in this study, including the standard stars, seemed too good to go unpublished. Both for illustrative purposes, as well as for use by those looking to model and remove Brackett series transitions from their A-dwarf telluric standard stars, we provide spectra for all the A dwarf stars observed as part of the Subaru-IRCS program. Table A.1 gives the stars and spectral types, Figure 14 shows the spectra. Because we intend for these spectra to be used as templates, we have removed most of the high order noise features artificially. This was not done in the OB spectra because we did not want to inadvertently alter their very weak profiles.
The stars in Figure 14 have been arranged first by spectral class (all are dwarfs), and next by rotational velocity. Over even this small sample, its clear that in modeling (to remove) Brackett features in A-dwarf stars, knowing the rotational velocity is at least as important as knowing the spectral type. The profile differences are small compared to the differences seen as a function of rotational velocity. |
warning/0506/math0506238.html | ar5iv | text | # A characterization of Prym varieties
## 1 Introduction
The problem of the characterization of the Prymians among principally polarized abelian varieties is almost as old as the famous Riemann-Schottky problem on the characterization of the Jacobian locus. Until now despite all the efforts it has remained unsolved. Analogs of quite a few geometrical characterizations of the Jacobians for the case of Prym varieties are either unproved or known to be invalid (see reviews and references therein).
The first effective solution of the Riemann-Schottky problem was obtained by T.Shiota (), who proved Novikov’s conjecture: the Jacobains of curves are exactly the indecomposable principally polarized abelian varieties whose theta-functions provide explicit solutions of the KP equation. Attempts to prove the analog of Novikov’s conjecture for the case of Prym varieties were made in . In it was shown that Novikov-Veselov (NV) equation provides local solution of the characterization problem. In the characterizations of the Prym varieties in terms of BKP and NV equations were proved only under certain additional assumptions. Note, that in a counter example showing that BKP equation has theta-functional solutions which do not correspond to the Prym varieties was constructed.
The goal of this work is to solve the characterization problem of the Prym varieties using the new approach proposed in the author’s previous work , where it was shown that KP equation contains excessive information and the Jacobian locus can be characterized in terms only of one of its auxiliary linear equations.
###### Theorem 1.1
() An indecomposable symmetric matrix $`B`$ with positive definite imaginary part is the matrix of the $`b`$-periods of normalized holomorphic differentials on a curve of genus g if and only if there exist $`g`$-dimensional vectors $`U0,V,A`$ such that the equation
$$\left(_y_x^2+u\right)\psi =0$$
(1.1)
is satisfied with
$$u=2_x^2\mathrm{ln}\theta (Ux+Vy+Z)\mathrm{and}\psi =\frac{\theta (A+Ux+Vy+Z)}{\theta (Ux+Vy+Z)}e^{px+Ey},$$
(1.2)
where $`p,E`$ are constants.
Here $`\theta (z)=\theta (z|B),z=(z_1,z_2,\mathrm{},z_g)`$ is the Riemann theta-function, defined by the formula
$$\theta (z)=\underset{m^g}{}e^{2\pi i(z,m)+\pi i(Bm,m)},(z,m)=m_1z_1+\mathrm{}+m_gz_g$$
(1.3)
The addition formula for the Riemann theta-function directly implies that equation (1.1) with $`u`$ and $`\psi `$ as in (1.2) is in fact equivalent to the system of equations
$$_V\mathrm{\Theta }[\epsilon ,0](A/2)_U^2\mathrm{\Theta }[\epsilon ,0](A/2)2p_U\mathrm{\Theta }[\epsilon ,0](A/2)+(Ep^2)\mathrm{\Theta }[\epsilon ,0](A/2)=0,$$
(1.4)
where $`\mathrm{\Theta }[\epsilon ,0](z)=\theta [\epsilon ,0](2z|2B)`$ are level two theta-functions with half-integer characteristics $`\epsilon \frac{1}{2}Z_2^g`$.
The characterization of the Jacobian locus given by Theorem 1.1 is stronger than that given in terms of the KP equation (see details in , where Theorem 1.1 was proved under the assumption that the closure $`A`$ of the subgroup of $`X`$ generated by $`A`$ is irreducible). In terms of the Kummer map,
$$\kappa :ZX\{\mathrm{\Theta }[\epsilon _1,0](Z):\mathrm{}:\mathrm{\Theta }[\epsilon _{2^g},0](Z)\}^{2^g1},$$
(1.5)
the statement of Theorem 1.1 is equivalent to the characterization of the Jacobians via flexes of the Kummer varieties, which is a particular case of the trisecant conjecture, first formulated in .
The Prym variety of a smooth algebraic curve $`\mathrm{\Gamma }`$ with involution $`\sigma :\mathrm{\Gamma }\mathrm{\Gamma }`$ is defined as the odd subspace $`𝒫(\mathrm{\Gamma })J(\mathrm{\Gamma })`$ of the Jacobian with respect to the involution $`\sigma ^{}:J(\mathrm{\Gamma })J(\mathrm{\Gamma })`$ induced by $`\sigma `$. It is principally polarized only if $`\sigma `$ has no fixed points or has two fixed points $`P_\pm `$. In this work we consider only the second case.
Let $`\mathrm{\Gamma }`$ be a smooth algebraic curve with involution $`\sigma `$ having two fixed points $`P_\pm `$. From the Riemann-Hurwitz formula it follows that if the genus of the factor-curve $`\mathrm{\Gamma }_0=\mathrm{\Gamma }/\sigma `$ equals $`g`$, then the genus $`\mathrm{\Gamma }`$ is equal to $`2g`$. It is known that on $`\mathrm{\Gamma }`$ there exists a basis of cycles $`a_i,b_i`$ with the canonical matrix of intersections $`a_ia_j=b_ib_j=0,a_ib_j=\delta _{ij},1i,j2g,`$ such that $`\sigma (a_k)=a_{g+k},\sigma (b_k)=b_{g+k},1kg`$. If $`d\omega _i`$ are normalized holomorphic differentials on $`\mathrm{\Gamma }`$, then the differential $`du_k=d\omega _k+d\omega _{g+k}`$ are odd $`\sigma ^{}(du_k)=du_k`$. By definition they are called the normalized holomorphic Prym differentials. The matrix of their $`b`$-periods
$$\mathrm{\Pi }_{kj}=_{b_k}𝑑u_j,1k,jg,$$
(1.6)
is symmetric, has positive definite imaginary part, and defines the Prym theta-function $`\theta _{Pr}=\theta (z|\mathrm{\Pi })`$.
Before presenting our main result it is necessary to mention that the Prym variety remains non-degenerate (compact) under certain degenerations of the curve. No characterization of Prym varieties given in terms of equations for the matrix $`\mathrm{\Pi }`$ of periods of the Prym differentials can single out the possibility of such degenerations. An algebraic curve $`\mathrm{\Gamma }`$ that is smooth outside fixed points $`P_+,P_{},Q_1,Q_2,\mathrm{},Q_k`$ of its involution $`\sigma `$, where $`P_\pm `$ are smooth and $`Q_k`$ are simple double points at which $`\sigma `$ does not permute branches of $`\mathrm{\Gamma }`$, will be denoted below by $`\{\mathrm{\Gamma },\sigma ,P_\pm ,Q_k\}`$.
###### Theorem 1.2
An indecomposable principally polarized abelian variety $`(X,\theta )`$ is the Prym variety of a curve of type $`\{\mathrm{\Gamma },\sigma ,P_\pm ,Q_k\}`$, if and only if there exist $`g`$-dimensional vectors $`U0,V0,A`$ such that one of the following equivalent conditions holds:
$`(A)`$ The equation
$$\left(_x_t+u\right)\psi =0$$
(1.7)
is satisfied with
$$u=2_{xt}^2\mathrm{ln}\theta (Ux+Vt+Z)+C,\mathrm{and}\psi =\frac{\theta (A+Ux+Vt+Z)}{\theta (Ux+Vt+Z)}e^{px+Et},$$
(1.8)
where $`C,p,E`$ are constants.
$`(B)`$ The equations
$$_{UV}^2\mathrm{\Theta }[\epsilon ,0](A/2)+p_V\mathrm{\Theta }[\epsilon ,0](A/2)+E_U\mathrm{\Theta }[\epsilon ,0](A/2)+C\mathrm{\Theta }[\epsilon ,0](A/2)=0$$
(1.9)
are satisfied for all $`\epsilon \frac{1}{2}Z_2^g`$.
$`(C)`$ The equation
$$_U\theta _V\theta \left(_U^2_V^2\theta \right)+_U^2\theta _V^2\theta \left(_U_V\theta \right)_U^2\theta _V\theta \left(_U_V^2\theta \right)_U\theta _V^2\theta \left(_U^2_V\theta \right)|_\mathrm{\Theta }=0$$
(1.10)
is valid on the theta-divisor $`\{Z\mathrm{\Theta }:\theta (Z)=0\}`$.
The equivalence of $`(A)`$ and $`(B)`$ is a direct corollary of the addition formula for the theta-function. The ”if” part of $`(A)`$ follows from the construction of integrable $`2D`$ Schrödinger operators given in . This construction is presented in the next section.
The statement $`(C)`$ is actually what we use for the proof of the theorem. It is stronger than $`(A)`$. The implication $`(A)(C)`$ does not require the explicit theta-functional form of $`\psi `$. It is enough to require only that equation (1.7) with $`u`$ as in (1.8) has local meromorphic in $`x`$ (or $`t`$) solutions which are holomorphic outside the divisor $`\theta (Ux+Vt+Z)=0`$.
To put it more precisely, let us consider a function $`\tau (x,t)`$ which is a holomorphic function of $`x`$ in some domain where the equation $`\tau (x,t)=0`$ has a simple root $`\eta (t)`$. It turns out that equation (1.7) with the potential $`u=2_{xt}^2\mathrm{ln}\tau +C`$, where $`C`$ is a constant, has a meromorphic solution in $`D`$, if this root satisfies the equation
$$\ddot{\eta }v\dot{\eta }\dot{v}+2\dot{\eta }^2w=0,$$
(1.11)
where $`v=v(t),w=w(t)`$ are the first coefficients of the Laurent expansion of $`u`$ at $`\eta `$
$$u(x,t)=\frac{2\dot{\eta }}{(x\eta )^2}+v+w(x\eta )+\mathrm{}$$
(1.12)
and ”dots” stand for $`t`$-derivatives. Straightforward but tedious computations with expansion of $`\theta `$ at the generic points of its divisor $`\mathrm{\Theta }`$ show that equation (1.11) in the case when $`\tau =\theta (Ux+Vt+Z)`$ is equivalent to equation (1.10).
Note that equations (1.11) are analogues of the equations derived in and called in the formal Calogero-Moser system. In a similar way, if we represent an entire function $`\tau `$ as a product
$$\tau (x,t)=c(t)\underset{i}{}(xx_i(t)),$$
(1.13)
then equation (1.11) takes the form
$$\underset{ji}{}\left[\frac{\ddot{x}_i\dot{x}_j\dot{x}_i\ddot{x}_j}{(x_ix_j)^2}\frac{2\dot{x}_i\dot{x}_j(\dot{x}_i+\dot{x}_j)}{(x_ix_j)^3}\right]=0.$$
(1.14)
At the moment the only reason for presenting equations (1.14) is to show that in the case when $`\tau `$ is a rational, trigonometric or elliptic polynomial the system (1.11) gives well-defined equations of motion for a multi-particle system. <sup>2</sup><sup>2</sup>2A.Zotov noticed that equation (1.14) is equivalent to the equations $`\ddot{x}_i=2_{ji}\dot{x}_i\dot{x}_j/(x_ix_j)`$ and, therefore, can be regarded as a limiting case of the Ruijesenaars-Schneider system.
At the beginning of section 3 we derive equation (1.11) and show that equation (1.10) is sufficient for the local existence of wave solutions of (1.7) having the form
$$\psi _\pm (x,t,k)=e^{kt_\pm }\left(1+\underset{s=1}{\overset{\mathrm{}}{}}\xi _s^\pm (x,t)k^s\right),t_+=x,t_{}=t,$$
(1.15)
and such that
$$\xi _s^\pm =\frac{\tau _s^\pm (Ux+Vt+Z,t_{})}{\theta (Ux+Vt+Z)},Z\mathrm{\Sigma }_\pm ,$$
(1.16)
where $`\tau _s^\pm (Z,t_{})`$, as a function of $`Z`$, is holomorphic in some open domain in $`^g`$. Here and below $`\mathrm{\Sigma }_\pm \mathrm{\Theta }`$ are subsets of the theta-divisor invariant under the shifts along constant vector fields $`_U`$ or $`_V`$, respectively.
The coefficients $`\xi _s^\pm `$ of the wave solutions are defined recurrently by the equations $`_{}\xi _{s+1}^\pm =_{xt}\xi _su\xi _s`$. The local existence of meromorphic solutions requires vanishing of the residues of the nonhomogeneous terms. That is controlled by equation (1.11). At the local level the main problem is to find the translational invariant normalization of $`\xi _s^\pm `$ which defines wave solutions uniquely up to a $`(x,t)`$-independent factor.
Following the ideas of and we fix such a normalization using extensions of $`\xi _s^\pm `$ along the affine subspaces $`Z+_\pm ^d`$, where $`_\pm ^d`$ are universal covers of the abelian subvarieties $`Y_\pm X`$ which are closures of the subgroups $`Ux`$ and $`Vt`$ in $`X`$, respectively. The corresponding wave solutions are called $`\lambda `$-periodic.
In the last section we show that for each $`Z\mathrm{\Sigma }_\pm `$ a local $`\lambda `$-periodic wave solution is the common eigenfunction of a commutative ring $`𝒜_\pm ^Z`$ of ordinary differential operators. The coefficients of these operators are independent of ambiguities in the construction of $`\psi `$. The theory of commuting differential operators implies then that the correspondence $`Z𝒜_\pm ^Z`$ defines a map $`j`$ of $`X\mathrm{\Sigma }_\pm `$ into the space $`\overline{\mathrm{Pic}(\mathrm{\Gamma })}`$ of torsion-free rang 1 sheafs $``$ on $`Z`$-independent spectral curve of $`𝒜^Z`$. That allows us to make the next crucial step and prove the global existence of the wave function. The global existence of the wave function implies that for the generic $`Z\mathrm{\Sigma }_\pm `$ the orbit of $`𝒜^Z`$ under the NV flows defines an imbedding $`i_Z`$ of the Prym variety $`𝒫(\mathrm{\Gamma })`$ of the spectral curve into $`X`$. Therefore, the Prym variety is compact. That implies the explicit description of possible types of singular points of $`\mathrm{\Gamma }`$. The final step in the proof of the main theorem is to show that there are not singular points of the multiplicity bigger then 2.
## 2 Integrable 2D Schrödinger operators
In this section we present necessary facts from the theory of integrable $`2D`$ Schrödinger equations and related hierarchies.
Let $`\mathrm{\Gamma }`$ be a smooth algebraic curve of genus $`g`$ with fixed local coordinates $`k_\pm ^1`$ at punctures $`P_\pm `$ and let $`t^{(\pm )}=\{t_i^{(\pm )}\}`$ be finite sets of complex variables. Then according to the general construction of the multi-point Baker-Akhiezer functions () for each non-special effective divisor $`D=\{\gamma _1,\mathrm{},\gamma _g\}`$ of degree $`g`$ there exists a unique function $`\psi _0(t^{(+)},t^{()},P)`$, which, as a function of the variable $`P\mathrm{\Gamma }`$, is meromorphic on $`\mathrm{\Gamma }P_\pm `$, where it has poles at $`\gamma _s`$ of degree not greater then the multiplicity of $`\gamma _s`$ in $`D`$. In the neighborhood of $`P_\pm `$ the function $`\psi _0`$ has the form
$$\psi _0=e^{_ik^it_i^{(\pm )}}\left(\underset{s=0}{\overset{\mathrm{}}{}}\xi _s^\pm (t)k^s\right),\xi _0^+=1,$$
(2.1)
where $`k=k_\pm ^1(P)`$ and $`t=\{t^{(+)},t^{()}\}`$.
The uniqueness of $`\psi _0`$ implies that for each positive integer $`n`$ there exists unique differential operators $`B_n^\pm `$ in the variables $`t_1^\pm `$
$$B_n^\pm =_\pm ^n+\underset{i=0}{\overset{n1}{}}v_{n,i}^\pm (t)_\pm ^{ni},_\pm =/t_1^\pm ,$$
(2.2)
such that
$$\left(_{t_n^\pm }B_n^\pm \right)\psi _0=0.$$
(2.3)
Equations (2.3) directly imply
$$[_{t_n^\pm }B_n^\pm ,_{t_m^\pm }B_m^\pm ]=0.$$
(2.4)
In other words, the operators $`B_n^\pm `$ satisfy zero-curvature equations which define two copies of the KP hierarchy with respect to the times $`t_n^\pm `$.
The two-point Baker-Akhiezer function with separated variables was introduced in where it was proved that in addition to (2.3) it satisfies the equation
$$H\psi _0=\left(_+_{}+w_++u\right)\psi _0=0,$$
(2.5)
where
$$w=_{}\mathrm{ln}\xi _0^{},u=_+_{}\mathrm{ln}\xi _1^+.$$
(2.6)
The operator $`H`$ defined in the left hand side of (2.5) ”couples” two copies of the KP hierarchy corresponding to the punctures $`P_\pm `$ via the equation
$$[_{t_n^+}B_n^+,_{t_m^{}}B_m^{}]=D_{nm}H.$$
(2.7)
The sense of (2.7) is as follows. Each differential operator $`𝒟`$ in the two variables $`t_1^\pm `$ can be uniquely represented in the form
$$𝒟=DH+D^++D^{},$$
(2.8)
where $`D^\pm `$ are ordinary differential operators in the variables $`t_1^\pm `$, respectively. The equation (2.7) is just the statement that the second and the third terms in the corresponding representation of the left hand side of (2.7) are equal to zero. This implies $`n+m1`$ equations on $`n+m1`$ unknown functions (the coefficients of $`B_n^+`$ and $`B_m^{}`$ ). Therefore, the operator equation (2.8), is equivalent to the well-defined system of non-linear partial differential equations.
Explicit theta-functional formulae for the solutions of these equations follow from the theta-functional formula for the Baker-Akhiezer function
$$\psi _0=\frac{\theta (A(P)+\underset{i}{}\left(U_i^+t_i^++U_i^{}t_i^{}\right)+Z)\theta (A(P_+)+Z)}{\theta (A(P_+)+_i\left(U_i^+t_i^++U_i^{}t_i^{}\right)+Z)\theta (A(P)+Z)}e^{_i\left(t_i^+\mathrm{\Omega }_i^+(P)+t_i^{}\mathrm{\Omega }_i^{}(P)\right)}$$
(2.9)
Here:
a) $`\theta (z)=\theta (z|B)`$ is the Riemann theta-function defined by the matrix $`B`$ of $`b`$-periods of normalized holomorphic differentials $`d\omega _k`$ on $`\mathrm{\Gamma }`$.
b) $`\mathrm{\Omega }_i^\pm (P)=^P𝑑\mathrm{\Omega }_i^\pm `$ is the Abelian integral corresponding to the normalized, $`_{a_k}𝑑\mathrm{\Omega }_i^\pm =0,`$ meromorphic differential on $`\mathrm{\Gamma }`$ with the only pole of the form
$$d\mathrm{\Omega }_i^\pm =dk_\pm ^i(1+O(k_\pm ^{i1}))$$
(2.10)
at the puncture $`P_\pm `$;
c) $`2\pi iU_j^\pm `$ is a vector of $`b`$-periods of the differential $`d\mathrm{\Omega }_j^\pm `$ with the coordinates
$$U_{j,k}^\pm =\frac{1}{2\pi i}_{b_k}𝑑\mathrm{\Omega }_j^\pm ;$$
(2.11)
d) $`A(P)`$ is the Abel transform, i.e. it is a vector with the coordinates $`A_k(P)=^P𝑑\omega _k`$;
e) $`Z`$ is an arbitrary vector (it corresponds to the divisor of poles of Baker-Akhiezer function).
Taking the evaluation at $`P_{}`$ and the expansion at $`P_+`$ of the regular factor in (2.9), one gets theta-functional formulae for the coefficients (2.6) of the corresponding 2D Schrödinger operator:
$`w`$ $`=`$ $`_{}\mathrm{ln}\left({\displaystyle \frac{\theta (A(P_{})+\underset{i}{}\left(U_i^+t_i^++U_i^{}t_i^{}\right)+Z)}{\theta (A(P_+)+_i\left(U_i^+t_i^++U_i^{}t_i^{}\right)+Z)}}\right),`$ (2.12)
$`u`$ $`=`$ $`_+_{}\mathrm{ln}\left(\theta (A(P_+)+{\displaystyle \underset{i}{}}\left(U_i^+t_i^++U_i^{}t_i^{}\right)+Z)\right)+C,`$ (2.13)
where the constant $`C`$ is equal to $`C=\mathrm{res}_{P_+}\mathrm{\Omega }_{}d\mathrm{\Omega }_+`$. Note, that the second factors in the numerator and denominator of the formula (2.9) are $`t`$-independent. Therefore, the function $`\psi `$ given by the following formula
$$\psi =\frac{\theta (A(P)+\underset{i}{}\left(U_i^+t_i^++U_i^{}t_i^{}\right)+Z)}{\theta (A(P_+)+_i\left(U_i^+t_i^++U_i^{}t_i^{}\right)+Z)}e^{_i\left(t_i^+\mathrm{\Omega }_i^+(P)+t_i^{}\mathrm{\Omega }_i^{}(P)\right)}$$
(2.14)
is a solution of the same linear equations as $`\psi _0`$. Below $`\psi `$ given by (2.14) will be called non-normalized Baker-Akhiezer function.
Potential operators. From now on we will consider only potential Schrödinger operators $`H=_+_{}+u`$. The reduction of the above described algebraic-geometrical construction to the potential case was found in . The corresponding algebraic-geometrical data are singled out by the following constraints:
$`(i)`$ The curve $`\mathrm{\Gamma }`$ should be a curve with involution $`\sigma :\mathrm{\Gamma }\mathrm{\Gamma }`$ which has two fixed points $`P_\pm `$.
$`(ii)`$ The equivalence class $`[D]J(\mathrm{\Gamma })`$ of the divisor $`D`$ should satisfy the equation
$$[D]+[\sigma (D)]=K+P_++P_{}J(\mathrm{\Gamma }),$$
(2.15)
where $`K`$ is the canonical class, i.e. the equivalence class of the zero-divisor of a holomorphic differential on $`\mathrm{\Gamma }`$.
Equation (2.15) is equivalent to the condition that the divisor $`D+\sigma (D)`$ is the zero divisor of a meromorphic differential $`d\mathrm{\Omega }`$ on $`\mathrm{\Gamma }`$ with simple poles at the punctures $`P_\pm `$. The differential $`d\mathrm{\Omega }`$ is even with respect to the involution and descends to a meromorphic differential on the factor-curve $`\mathrm{\Gamma }_0=\mathrm{\Gamma }/\sigma `$. The projection $`\pi :\mathrm{\Gamma }\mathrm{\Gamma }_0=\mathrm{\Gamma }/\sigma `$ represents $`\mathrm{\Gamma }`$ as a two-sheet covering of $`\mathrm{\Gamma }_0`$ with $`2`$ branch points $`P_\pm `$. In this realization the involution $`\sigma `$ is a permutation of the sheets. For $`P\mathrm{\Gamma }`$ we denote the point $`\sigma (P)`$ by $`P^\sigma `$. From the Riemann-Hurwitz formula it follows that the genus $`g`$ of $`\mathrm{\Gamma }`$ equals $`g=2g_0`$, where $`g_0`$ is the genus of $`\mathrm{\Gamma }_0`$. Note that the divisors that satisfy (2.15) are parameterized by the points $`Z_0`$ of the Prym variety $`𝒫(\mathrm{\Gamma })J(\mathrm{\Gamma })`$.
###### Theorem 2.1
(). Let a smooth algebraic curve $`\mathrm{\Gamma }`$ and an effective divisor $`D`$ satisfy the constraints $`(i),(ii)`$. Let $`k_\pm ^1(P)`$ be odd local coordinates in the neighborhoods of the fixed points $`P_\pm `$, i.e. $`k_\pm (P)=k_\pm (\sigma (P))`$, and let all the even times vanish, i.e. $`t_{2i}^\pm =0`$. Then the corresponding $`2D`$ Schrödinger operator is potential, i.e. $`w=0`$.
In ) it was also found that for the potential operators the formulae (2.9) and (2.13) can be expressed in terms of the Prym theta-function. For further use it is enough to present these formulae for the case of only two nontrivial variables $`x=t_1^+,t=t_1^{}`$:
$$\psi =\frac{\theta _{Pr}(A^{Pr}(P)+Ux+Vt+Z)}{\theta _{Pr}(A^{Pr}(P_+)+Ux+Vt+Z)}e^{x\mathrm{\Omega }_1^++t\mathrm{\Omega }_1^{}}$$
(2.16)
$$u=2_+_{}\mathrm{ln}\theta _{Pr}(A^{Pr}(P_+)+xU+tV+Z)+C.$$
(2.17)
Here $`A^{Pr}:\mathrm{\Gamma }𝒫(\mathrm{\Gamma })`$ is the Abel-Prym map defined by the Prym differentials, i.e. $`A^{Pr}(P)`$ is a vector with the coordinates $`A_k^{Pr}(P)=^P𝑑u_k`$.
In it was proved that for the case of smooth periodic potentials $`u(x,t)`$ (considered as a function of real variables $`x,t`$) the conditions found by Novikov and Veselov are sufficient and necessary.
## 3 $`\lambda `$-periodic wave solutions
To begin with, let us show that equations (1.11) are the necessary condition of the existence of a meromorphic solution to equation (1.7).
Let $`\tau (x,t)`$ be a smooth $`t`$-parametric family of holomorphic functions of the variable $`x`$ in some open domain $`D`$. Suppose that in $`D`$ the function $`\tau `$ has a simple zero,
$$\tau (\eta (t),t)=0,\tau _x(\eta (t),t)0.$$
(3.1)
###### Lemma 3.1
If equation (1.7) with the potential $`u=2_{xt}^2\mathrm{ln}\tau (x,t)+C`$, where $`C`$ is a constant, has a meromorphic solution $`\psi _0(x,t)`$, then equation (1.11) holds.
Proof. Consider the Laurent expansions of $`\psi _0`$ and $`u`$ in the neighborhood of $`\eta `$:
$$u=\frac{2\dot{\eta }}{(x\eta )^2}+v+w(x\eta )+\mathrm{}$$
(3.2)
$$\psi _0=\frac{\alpha }{x\eta }+\beta +\gamma (x\eta )+\mathrm{}$$
(3.3)
(All the coefficients in these expansions are smooth functions of the variable $`t`$). Substitution of (3.2,3.3) into (1.7) gives a system of equations. The first three of them are
$$\dot{\alpha }2\dot{\eta }\beta =0,$$
(3.4)
$$2\dot{\eta }\gamma +\alpha v=0,$$
(3.5)
$$\dot{\gamma }+v\beta +\alpha w=0.$$
(3.6)
Taking the $`t`$-derivative of the second equation and using two others we get (1.11).
Let us show that equations (1.11) are sufficient for the existence of meromorphic wave solutions.
###### Lemma 3.2
Suppose that equation (1.11) for the zero of $`\tau (x,t)`$ holds. Then equation (1.7) has wave solutions of the form
$$\psi =e^{kt}\left(1+\underset{s=1}{\overset{\mathrm{}}{}}\xi _s(x,t)k^s\right)$$
(3.7)
such that the coefficients $`\xi _s`$ have simple poles at $`\eta `$ and are holomorphic everywhere else in $`D`$.
Proof. Substitution of (3.7) into (1.7) gives a recurrent system of equations
$$\xi _{s+1}^{}=_{xt}^2\xi _su\xi _s.$$
(3.8)
We are going to prove by induction that this system has meromorphic solutions with simple poles at $`\eta `$.
Let us expand $`\xi _s`$ at $`\eta `$:
$$\xi _s=\frac{r_s}{x\eta }+r_{s0}+r_{s1}(x\eta )+\mathrm{}$$
(3.9)
Suppose that $`\xi _s`$ is defined, and equation (3.8) has a meromorphic solution. Then the right hand side of (3.8) has the zero residue at $`x=\eta `$, i.e.,
$$\mathrm{res}_\eta \left(_{xt}^2\xi _s+u\xi _s\right)=vr_s+2\dot{\eta }r_{s1}=0.$$
(3.10)
We need to show that the residue of the next equation also vanishes. From (3.8) it follows that the coefficients of the Laurent expansion for $`\xi _{s+1}`$ are equal to
$$r_{s+1}=\dot{r}_s+2\dot{\eta }r_{s0},$$
(3.11)
$$r_{s+1,1}=vr_{s0}wr_s\dot{r}_{s1}.$$
(3.12)
These equations and equation (3.10) imply
$$(vr_{s+1}+2\dot{\eta }r_{s+1,1})=2\dot{\eta }wr_s+v\dot{r}_s+2\dot{\eta }\dot{r}_{s,1}=_t\left(vr_s+2\dot{\eta }r_{s,1}\right)\left(\dot{v}\frac{\ddot{\eta }}{\dot{\eta }}v2\dot{\eta }w\right)r_s=0,$$
(3.13)
and the lemma is proved.
Our next goal is to fix a translation-invariant normalization of $`\xi _s`$ which defines wave functions uniquely up to a $`(x,t)`$-independent factor. It is instructive to consider first the case of the periodic potentials $`u(x+1,t)=u(x,t)`$ (compare with ).
Equations (3.8) are solved recursively by the formulae
$$\xi _{s+1}(x,t)=c_{s+1}(t)+\xi _{s+1}^0(x,t),$$
(3.14)
$$\xi _{s+1}^0(x,t)=_t\xi _s_{x_0}^xu\xi _s𝑑x,$$
(3.15)
where $`c_s(t)`$ are arbitrary functions of the variable $`t`$. Let us show that the periodicity condition $`\xi _s(x+1,t)=\xi _s(x,t)`$ defines these functions uniquely up to constants. Assume that $`\xi _{s1}`$ is known and satisfies the condition that the corresponding function $`\xi _s^0`$ is periodic. The choice of the function $`c_s(t)`$ does not affect the periodicity property of $`\xi _s`$, but it does affect the periodicity in $`x`$ of the function $`\xi _{s+1}^0(x,t)`$. In order to make $`\xi _{s+1}^0(x,t)`$ periodic, the function $`c_s(t)`$ should satisfy the linear differential equation
$$_tc_s(t)+_{x_0}^{x_0+1}u(x,t)(c_s(t)+\xi _s^0(x,t))𝑑x=0.$$
(3.16)
This defines $`c_s`$ uniquely up to a constant.
In the general case, when $`u`$ given by (1.8) is quasi-periodic, the normalization of the wave functions is defined along the same lines.
Let $`\mathrm{\Theta }_1`$ be defined by the equations $`\mathrm{\Theta }_1=\{Z:\theta (Z)=_U\theta (Z)=0\}`$, where $`_U`$ is a constant vector-field on $`^g`$, corresponding to the vector $`U`$ in (1.8). The $`_U`$-invariant subset $`\mathrm{\Sigma }`$ of $`\mathrm{\Theta }_1`$ will be called the singular locus.
Consider the closure $`Y_U=Ux`$ of the group $`Ux`$ in $`X`$. Shifting $`Y_U`$ if needed, we may assume, without loss of generality, that $`Y_U`$ is not in the singular locus, $`Y_U\mathrm{\Sigma }`$. Then, for a sufficiently small $`t`$, we have $`Y_U+Vt\mathrm{\Sigma }`$ as well. Consider the restriction of the theta-function onto the affine subspace $`^d+Vt`$, where $`^d=\pi ^1(Y_U)`$, and $`\pi :^gX=^g/\mathrm{\Lambda }`$ is the universal cover of $`X`$:
$$\tau (z,t)=\theta (z+Vt),z^d.$$
(3.17)
The function $`u(z,t)=2_U_t\mathrm{ln}\tau +C`$ is periodic with respect to the lattice $`\mathrm{\Lambda }_U=\mathrm{\Lambda }^d`$ and, for fixed $`t`$, has a double pole along the divisor $`\mathrm{\Theta }^U(t)=\left(\mathrm{\Theta }Vt\right)^d`$.
###### Lemma 3.3
Let equation (1.10) hold and let $`\lambda `$ be a vector of the sublattice $`\mathrm{\Lambda }_U=\mathrm{\Lambda }^d^g`$. Then:
(i) equation (1.7) with the potential $`u(Ux+z,t)`$ has a wave solution of the form $`\psi =e^{kt+bxk^1}\varphi (Ux+z,t,k)`$ such that the coefficients $`\xi _s(z,t)`$ of the formal series
$$\varphi (z,t,k)=1+\underset{s=1}{\overset{\mathrm{}}{}}\stackrel{~}{\xi }_s(z,t)k^s$$
(3.18)
are $`\lambda `$-periodic meromorphic functions of the variable $`z^d`$ with a simple pole at the divisor $`\mathrm{\Theta }^U(t)`$, i.e.
$$\stackrel{~}{\xi }_s(z+\lambda ,t)=\stackrel{~}{\xi }_s(z,t)=\frac{\tau _s(z,t)}{\tau (z,t)};$$
(3.19)
(ii) $`\varphi (z,t,k)`$ is unique up to a factor $`\rho (z,k)`$ that is $`t`$-independent, $`_U`$-invariant and holomorphic in $`z`$,
$$\varphi _1(z,t,k)=\varphi (z,t,k)\rho (z,k),_U\rho =0.$$
(3.20)
Proof. The functions $`\stackrel{~}{\xi }_s(z)`$ are defined recursively by the equations
$$_U\stackrel{~}{\xi }_{s+1}=_U_t\stackrel{~}{\xi }_s(u+b)\stackrel{~}{\xi }_sb_t\stackrel{~}{\xi }_{s1}.$$
(3.21)
A particular solution of the first equation $`_U\stackrel{~}{\xi }_1=ub`$ is given by the formula
$$\stackrel{~}{\xi }_1^0=2_t\mathrm{ln}\tau (l,z)(b+C),$$
(3.22)
where $`(l,z)`$ is a linear form on $`^d`$ given by the scalar product of $`z`$ with a vector $`l^d`$ such that $`(l,U)=1`$. By definition, the vector $`\lambda `$ is in $`Y_U`$. Therefore, $`(l,\lambda )0`$. The periodicity condition for $`\stackrel{~}{\xi }_1^0`$ defines the constant $`b`$, which depends only on a choice of the lattice vector $`\lambda `$. From the monodromy properties of $`\theta `$ it follows that without loss of generality we may assume that $`\lambda `$ is chosen such that the corresponding constant $`b`$ is not equal to zero, i.e.
$$b=C+(l,\lambda )^1(2_t\mathrm{ln}\tau (z,t)2_t\mathrm{ln}\tau (z+\lambda ,t))0,$$
(3.23)
Note, that the second factor in (3.7) and the series $`\varphi `$ in (3.18) differ by the factor $`e^{bxk^1}`$, which does not affect the results of the previous lemma. Therefore, equations (1.11) are sufficient for the local solvability of (3.21) in any domain, where $`\tau (z+Ux,t)`$ has simple zeros, i.e. outside the set $`\mathrm{\Theta }_1^U(t)=\left(\mathrm{\Theta }_1Vt\right)^d`$. Recall that $`\mathrm{\Theta }_1=\mathrm{\Theta }_U\mathrm{\Theta }`$. This set does not contain a $`_U`$-invariant line because such line is dense in $`Y_U`$. Therefore, the sheaf $`𝒱_0`$ of $`_U`$-invariant meromorphic functions on $`^d\mathrm{\Theta }_1^U(t)`$ with poles along the divisor $`\mathrm{\Theta }^U(t)`$ coincides with the sheaf of holomorphic $`_U`$-invariant functions. That implies the vanishing of $`H^1(C^d\mathrm{\Theta }_1^U(t),𝒱_0)`$ and the existence of global meromorphic solutions $`\xi _s^0`$ of (3.21) which have a simple pole at the divisor $`\mathrm{\Theta }^U(t)`$ (see details in ).
Let us assume, as in the example above, that a $`\lambda `$-periodic solution $`\stackrel{~}{\xi }_{s1}`$ is known and that it satisfies the condition that there exists a $`\lambda `$-periodic solution $`\stackrel{~}{\xi }_s^0`$ of the next equation such that the equation
$$_U\chi _s=(u+b)\stackrel{~}{\xi }_s^0b_t\stackrel{~}{\xi }_{s1}$$
(3.24)
has a $`\lambda `$-periodic solution. If $`\stackrel{~}{\xi }_s^0`$ and a particular solution $`\chi _s^{}`$ of (3.24) are fixed, then $`\stackrel{~}{\xi }_{s+1}^{}=_t\stackrel{~}{\xi }_s^0+\chi _s^{}`$ is a $`\lambda `$-periodic solution of (3.21) for $`\stackrel{~}{\xi }_s^0`$.
A choice of a $`\lambda `$-periodic $`_U`$-invariant function $`c_s(z,t)`$ does not affect the periodicity property of $`\stackrel{~}{\xi }_s=c_s+\stackrel{~}{\xi }_s^0`$. It changes the right hand side of (3.24). A particular solution of the new equation is given by the formula $`\chi _s^0=\chi _s^{}+c_s\stackrel{~}{\xi }_1^0`$. Therefore, $`\stackrel{~}{\xi }_{s+1}^0=_t\stackrel{~}{\xi }_s+\chi _s^0`$ is a $`\lambda `$-periodic solution of (3.21) for $`\stackrel{~}{\xi }_s`$. The choice of $`c_s`$ does affect the existence of periodic solutions of the equation
$$_U\chi _{s+1}=(u+b)\stackrel{~}{\xi }_{s+1}^0b_t\stackrel{~}{\xi }_s$$
(3.25)
Let $`\stackrel{~}{\chi }_{s+1}`$ be a solution of the equation
$$_U\stackrel{~}{\chi }_{s+1}=(u+b)\stackrel{~}{\xi }_{s+1}^{}b_t\stackrel{~}{\xi }_s^0.$$
(3.26)
Then the function
$$\chi _{s+1}(z,t)=\stackrel{~}{\chi }_{s+1}(z,t)+\frac{1}{2}c_s(z,t)(\xi _1^0(z,t))^2+(\xi _1^0(z,t)(l,z)b)_tc_s(z,t),$$
(3.27)
is a solution of (3.25). In order to make $`\chi _{s+1}`$ periodic, the function $`c_s(z,t)`$ should satisfy the linear differential equation
$$_tc_s(z,t)=((l,\lambda )b)^1(\stackrel{~}{\chi }_{s+1}(z+\lambda ,t)\stackrel{~}{\chi }_{s+1}(z,t)).$$
(3.28)
This equation, together with the initial condition $`c_s(z)=c_s(z,0)`$ uniquely defines $`c_s(z,t)`$. The induction step is then completed. We have shown that the ratio of two periodic formal series $`\varphi _1`$ and $`\varphi `$ is $`t`$-independent. Therefore, equation (3.20), where $`\rho (z,k)`$ is defined by the evaluation of both the sides at $`t=0`$, holds. The lemma is thus proven.
###### Corollary 3.1
Let $`\lambda _1,\mathrm{},\lambda _d`$ be a set of linear independent vectors of the lattice $`\mathrm{\Lambda }_U`$ and let $`z_0`$ be a point of $`^d`$. Then, under the assumptions of the previous lemma, there is a unique wave solution of equation (1.7) such that the corresponding formal series $`\varphi (z,t,k;z_0)`$ is quasi-periodic with respect to $`\mathrm{\Lambda }_U`$, i.e. for $`\lambda \mathrm{\Lambda }_U`$
$$\varphi (z+\lambda ,t,k;z_0)=\varphi (z,t,k;z_0)\mu _\lambda (k)$$
(3.29)
and satisfies the normalization conditions
$$\mu _{\lambda _i}(k)=1,\varphi (z_0,0,k;z_0)=1.$$
(3.30)
The proof is identical to that in the part (b) of Lemma 12 in (compare with the proof of the corollary in ).
## 4 The spectral curve
In this section we show that $`\lambda `$-periodic wave solutions of equation (1.7), with $`u`$ as in (1.8), are common eigenfunctions of rings of commuting operators and identify $`X`$ with the Prym variety of the spectral curve of these rings.
Note that a simple shift $`zz+Z`$, where $`Z\mathrm{\Sigma },`$ gives $`\lambda `$-periodic wave solutions with meromorphic coefficients along the affine subspaces $`Z+^d`$. These $`\lambda `$-periodic wave solutions are related to each other by $`t`$-independent, $`_U`$-invariant factor. Therefore choosing in the neighborhood of any $`Z\mathrm{\Sigma },`$ a hyperplane orthogonal to the vector $`U`$ and fixing initial data on this hyperplane at $`t=0,`$ we define the corresponding series $`\varphi (z+Z,t,k)`$ as a local meromorphic function of $`Z`$ and the global meromorphic function of $`z`$.
###### Lemma 4.1
Let equation (1.10) hold. Then there is a unique pseudo-differential operator
$$(Z,_t)=_t+\underset{s=1}{\overset{\mathrm{}}{}}w_s(Z)_t^s$$
(4.1)
such that for $`Z+Vt\mathrm{\Sigma }`$
$$(Ux+Vt+Z,_t)\psi =k\psi ,$$
(4.2)
where $`\psi =e^{(kt+bxk^1)}\varphi (Ux+Z,t,k)`$ is a $`\lambda `$-periodic solution of (1.7). The coefficients $`w_s(Z)`$ of $``$ are meromorphic functions on the abelian variety $`X`$ with poles along the divisor $`\mathrm{\Theta }`$.
Proof. The construction of $``$ is standard for the KP theory. First, we define $``$ as a pseudo-differential operator with coefficients $`w_s(Z,t)`$, which are functions of $`Z`$ and $`t`$.
Let $`\psi `$ be a $`\lambda `$-periodic wave solution. The substitution of (3.18) into (4.2) gives a system of equations that recursively define $`w_s(Z,t)`$ as differential polynomials in $`\stackrel{~}{\xi }_s(Z,t)`$. The coefficients of $`\psi `$ are local meromorphic functions of $`Z`$, but the coefficients of $``$ are well-defined global meromorphic functions of on $`^g\mathrm{\Sigma }`$, because different $`\lambda `$-periodic wave solutions are related to each other by $`t`$-independent factor, which does not affect $``$. The singular locus is of codimension $`2`$. Then Hartogs’ holomorphic extension theorem implies that $`w_s(Z,t)`$ can be extended to a global meromorpic function on $`^g`$.
The translational invariance of $`u`$ implies the translational invariance of the $`\lambda `$-periodic wave solutions. Indeed, for any constant $`s`$, the series $`\varphi (Vs+Z,ts,k)`$ and $`\varphi (Z,t,k)`$ correspond to $`\lambda `$-periodic solutions of the same equation. Therefore, they coincide up to a $`t`$-independent, $`_U`$-invariant factor. This factor does not affect $``$. Hence, $`w_s(Z,t)=w_s(Vt+Z)`$.
The $`\lambda `$-periodic wave functions corresponding to $`Z`$ and $`Z+\lambda ^{}`$ for any $`\lambda ^{}\mathrm{\Lambda }`$ are also related to each other by a $`t`$-independent, $`_U`$-invariant factor. Hence, $`w_s`$ are periodic with respect to $`\mathrm{\Lambda }`$ and therefore are meromorphic functions on the abelian variety $`X`$. The lemma is proved.
###### Lemma 4.2
Let $``$ be a pseudo-differential operator corresponding to $`\lambda `$-periodic solution and $`^{}`$ be its formal adjoint operator. Then the following equation
$$^{}=_t_t^1$$
(4.3)
holds.
Recall, that the operator which is formally adjoint to $`(w^i)`$ is the operator $`()^iw`$, where $`w`$ stands for the operator of multiplication by the function $`w`$. Below we will use the notion of the left action of an operator which is identical to the formal adjoint action, i.e. by definition we assume that for a function $`f`$ the identity
$$(f𝒟)=𝒟^{}f$$
(4.4)
holds.
Proof. If $`\psi `$ is as in Lemma 3.3, then there exists a unique pseudo-differential operator $`\mathrm{\Phi }`$ such that
$$\psi =\mathrm{\Phi }e^{kt},\mathrm{\Phi }=1+\underset{s=1}{\overset{\mathrm{}}{}}\phi _s(Ux+Z,t)_t^s.$$
(4.5)
The coefficients of $`\mathrm{\Phi }`$ are universal differential polynomials on $`\stackrel{~}{\xi }_s`$. Therefore, $`\phi _s(z+Z,t)`$ is a global meromorphic function of $`zC^d`$ and a local meromorphic function of $`Z\mathrm{\Sigma }`$. Note that $`=\mathrm{\Phi }(_t)\mathrm{\Phi }^1`$, and the equation $`H\psi =0`$ is equivalent to the operator equation
$$_t\mathrm{\Phi }_x+u\mathrm{\Phi }=0,$$
(4.6)
where $`\mathrm{\Phi }_x`$ is the pseudo-differential operator with the coefficients $`_x\phi `$. Note that (4.6) implies
$$_x=[_t^1u,].$$
(4.7)
Let us define the dual wave function $`\psi ^{}`$ by the formula
$$\psi ^{}=\left(e^{kt}_t\mathrm{\Phi }^1_t^1\right)=\left(_t^1\left(\mathrm{\Phi }^1\right)^{}_t\right)e^{kt}.$$
(4.8)
Equation (4.6) implies $`H\psi ^{}=0`$. The dual wave function $`\psi ^{}`$ is $`\lambda `$-periodic. Therefore, the same arguments as used above show that if equation (1.11) is satisfied, then the dual wave function is of the form $`\psi ^{}=e^{(kt+bxk^1)}\varphi ^{}(Ux+Z,t,k)`$, where the coefficients $`\stackrel{~}{\xi }_s^{}(z+Z,t)`$ of the formal series
$$\varphi ^{}(z+Z,t,k)=1+\underset{s=1}{\overset{\mathrm{}}{}}\stackrel{~}{\xi }_s^{}(z+Z,t)k^s$$
(4.9)
have simple poles at the divisor $`\mathrm{\Theta }^U(t)`$. They are $`\lambda `$-periodic. Therefore,
$$\varphi ^{}(z+Z,t,k)=\varphi (z+Z,t,k)\rho (z+Z,k),$$
(4.10)
where $`\rho `$ is a $`t`$-independent, $`_U`$-invariant factor. Equation (4.10) implies (4.3) and the lemma is proved.
Commuting differential operators. Let as denote strictly positive differential part of the pseudo-differential operator $`^m`$ by $`_+^m`$, i.e. if $`^m=_{i=m}^{\mathrm{}}F_m^{(i)}_t^i`$, then <sup>3</sup><sup>3</sup>3 Note that this definition differs from the one used in the KP theory, where plus subscript denotes nonnegative part of a pseudo-differential operator
$$_+^m=\underset{i=1}{\overset{m}{}}F_m^{(i)}_t^i,_{}^m=^m_+^m=F_m^{(0)}+F_m^{(1)}_t^1+O(^2).$$
(4.11)
By definition of the residue of a pseudo-differential operator, the first leading coefficients of $`_{}^m`$ are
$$F_m^{(0)}=\mathrm{res}_{}\left(^m_t^1\right),F_m^{(1)}=\mathrm{res}_{}^m.$$
(4.12)
###### Lemma 4.3
The operators $`_+^m`$ satisfy the equations
$`H_+^{2m}`$ $`=`$ $`F_{2m,x}^{(0)}_t{\displaystyle \frac{1}{2}}F_{2m,xt}^{(0)}+B_{2m}H,`$ (4.13)
$`H_+^{2m+1}`$ $`=`$ $`F_{2m+1,x}^{(1)}+B_{2m+1}H,`$ (4.14)
where $`B_m`$ is a pseudo-differential operator in the variable $`t`$.
Proof. First, we prove the equation
$$H_+^m=F_{m,x}^{(0)}_t\left(F_{m,xt}^{(0)}+F_{m,x}^{(1)}\right)+B_mH.$$
(4.15)
Each operator $`𝒟`$ of the form $`𝒟=_{i=N}^{\mathrm{}}(a+b_x)_t^i`$ can be uniquely represented in the form $`𝒟=D_1+D_2H`$, where $`D_{1,2}`$ are pseudo-differential operators in the variable $`t`$. Consider such a representation for the operator $`H^m=D_1+D_2H`$. From the definition of $``$ it follows that $`H^m\psi =0`$. That implies $`D_1=0`$ or the equation
$$H^m=D_2H.$$
(4.16)
We have the identity
$$[_x_t+u,_+^m]=_{+,xt}^m+_{+,x}^m_t+[u,_+^m]_{+,t}^m_t^1u+_{+,t}^m_t^1H.$$
(4.17)
The first three terms are differential operators in the $`t`$ variable. By definition of $`_+^m`$ the fourth term is also a differential operator. Therefore, the pseudo-differential operator $`D_1`$ in the decomposition $`H_+^m=D_{m,1}+B_mH`$ is a differential operator.
In the same way we get the equation
$$H_{}^m=\stackrel{~}{D}_{m,1}+\stackrel{~}{B}_{2m}H,$$
(4.18)
where
$$\stackrel{~}{D}_{m,1}=_{,xt}^m+_{,x}^m_t+[u,_{}^m]_{,t}^m_t^1u$$
(4.19)
By definition of $`_{}^m`$ the operator $`\stackrel{~}{D}_{m,1}`$ is a pseudo-differential operator of order not greater than $`1`$. Equation (4.16) implies $`H_+^m=H_{}^m+D_2H`$. Hence, $`D_{m,1}=\stackrel{~}{D}_{m,1}`$ is a differential operator of the order $`1`$, i.e. has the form $`a_t+b`$. The coefficients of this operator can be easily found from the leading coefficients of the right hand side of (4.19). Direct computations give equation (4.15).
Now in order to complete the proof of (4.13) and (4.14) it is enough to use (4.3). From equation (4.3) and the relation $`\mathrm{res}_{}D=\mathrm{res}_{}D^{}`$ it follows that
$$F_{2m}^{(1)}=\mathrm{res}_{}(^{})^{2m}=\mathrm{res}_{}\left(_t^m_t^1\right)=F_{2m,t}^{(0)}F_{2m}^{(1)}.$$
(4.20)
In the same way we get
$$F_{2m+1}^{(0)}=\mathrm{res}_{}(^{2m+1}_t^1)=\mathrm{res}_{}\left(^{2m+1}_t^1\right)^{}=F_{2m+1}^{(0)}=0.$$
(4.21)
Equations (4.15, 4.20, 4.21) imply (4.13) and (4.14). The lemma is proved.
The following statement is a direct corollary of equations (4.20, 4.21).
###### Corollary 4.1
The operators $`_+^m`$ satisfy the relation
$$\left(_+^m\right)^{}=(1)^m_t_+^m_t^1$$
(4.22)
The next step is crucial for the construction of commuting operators.
###### Lemma 4.4
The functions $`F_{2m}^{(0)},F_{2m+1}^{(1)}`$ have at most second order pole on the divisor $`\mathrm{\Theta }`$.
Proof. The ambiguity in the definition of $`\psi `$ does not affect the product
$$\psi ^{}\psi =\left(e^{kt}_t\mathrm{\Phi }^1_t^1\right)\left(\mathrm{\Phi }e^{kt}\right).$$
(4.23)
Therefore, although each factor is only a local meromorphic function on $`^g\mathrm{\Sigma }`$, the coefficients $`J_s^{(0)}`$ of the product
$$\psi ^{}\psi =\varphi ^{}(Z,t,k)\varphi (Z,t,k)=1+\underset{s=2}{\overset{\mathrm{}}{}}J_s^{(0)}(Z,t)k^s$$
(4.24)
are global meromorphic functions of $`Z`$. Moreover, the translational invariance of $`u`$ implies that they have the form $`J_s(Z,t)=J_s(Z+Vt)`$. Each of the factors in the left hand side of (4.24) has a simple pole on $`\mathrm{\Theta }Vt`$. Hence, $`J_s(Z)`$ is a meromorphic function on $`X`$ with a second order pole at $`\mathrm{\Theta }`$.
From the definition of $``$ it follows that
$$\mathrm{res}_k\left(\psi ^{}(^n\psi )\right)k^1dk=\mathrm{res}_k\left(\psi ^{}k^n\psi \right)k^1dk=J_n^{(0)}.$$
(4.25)
On the other hand, using the identity
$$\mathrm{res}_k\left(e^{kx}𝒟_1\right)\left(𝒟_2e^{kx}\right)dk=\mathrm{res}_{}\left(𝒟_2𝒟_1\right),$$
(4.26)
we get
$$\mathrm{res}_k\left(\psi ^{}\left(^n\psi \right)\right)k^1dk=\mathrm{res}_k\left(e^{kt}_t\mathrm{\Phi }^1_t^1\right)\left(^n\mathrm{\Phi }_t^1e^{kt}\right)dk=\mathrm{res}_{}\left(^n_t^1\right)=F_n^{(0)}.$$
(4.27)
Therefore, $`F_n^{(0)}=J_n^{(0)}`$ has the second order pole at $`\mathrm{\Theta }`$.
Consider now the coefficients $`J_s^{(1)}`$ of the series
$$\psi ^{}\psi _t\psi _t^{}\psi =2k+\underset{s=1}{\overset{\mathrm{}}{}}2J_s^{(1)}(Z,t)k^s.$$
(4.28)
They are meromorphic functions on $`X`$ with the second order pole at $`\mathrm{\Theta }`$. We have
$$2J_n^{(1)}=\mathrm{res}_k\left((\psi ^{}^n)\psi _t\psi _t^{}(^n\psi )\right)k^1dk=\mathrm{res}_{}(^n+_t^n_t^1)=2F_n^{(1)}+F_{n,t}^{(0)}$$
(4.29)
Then from equation (4.21) it follows that $`F_{2m+1}^{(1)}=J_{2m+1}^{(1)}`$ and the lemma is proved.
Let $`𝐅`$ be a direct sum of the linear spaces $`\widehat{𝐅}^\alpha ,\alpha =0,1,`$ spanned by $`\{F_{2m+\alpha }^{(\alpha )},m=0,1,\mathrm{}\}`$. They are subspaces of the $`2^g`$-dimensional space of the abelian functions with at most second order pole at $`\mathrm{\Theta }`$. Therefore, for all but $`\widehat{g}^\alpha =\mathrm{dim}\widehat{𝐅}^\alpha `$ positive integers $`2n+\alpha `$, there exist constants $`c_{i,n}^{(\alpha )}`$ such that
$$F_{2n+\alpha }^{(\alpha )}(Z)+\underset{i=1}{\overset{n}{}}c_{i,n}^{(\alpha )}F_{2n2i+\alpha }^{(\alpha )}(Z)=0.$$
(4.30)
Let $`I^{(\alpha )}`$ denote the subset of integers $`2n+\alpha `$ for which none of such constants exist. We call the union $`I=I^{(0)}I^{(1)}`$ the gap sequence.
###### Lemma 4.5
Let $``$ be the pseudo-differential operator corresponding to a $`\lambda `$-periodic wave function $`\psi `$ constructed above. Then, for the differential operators
$$L_{2n+\alpha }=_+^{2n+\alpha }+\underset{i=1}{\overset{n}{}}c_{i,n}^{(\alpha )}_+^{2n+\alpha 2i},2n+\alpha I^\alpha ,$$
(4.31)
the equations
$$L_{2n+\alpha }\psi =a_{2n+\alpha }(k)\psi ,a_{2n+\alpha }(k)=k^{2n+\alpha }+\underset{s=1}{\overset{\mathrm{}}{}}a_{s,n}k^{2n+\alpha s},$$
(4.32)
where $`a_{s,n}`$ are constants, hold.
Proof. First note, that from (4.13, 4.14) it follows that $`HL_{2n+\alpha }\psi =0`$. Hence, if $`\psi `$ is a $`\lambda `$-periodic wave solution of (1.7) corresponding to $`Z\mathrm{\Sigma }`$, then $`L_{2n+\alpha }\psi `$ is also a formal solution of the same equation. That implies the equation $`L_{2n+\alpha }\psi =a_{2n+\alpha }(Z,k)\psi `$, where $`a_{2n+\alpha }(Z,k)`$ is $`t`$-independent and $`_U`$-invariant function of the variable $`Z`$. The ambiguity in the definition of $`\psi `$ does not affect $`a_{2n+\alpha }`$. Therefore, the coefficients of $`a_{2n+\alpha }`$ are well-defined global meromorphic functions on $`^g\mathrm{\Sigma }`$. The $`_U`$\- invariance of $`a_{2n+\alpha }`$ implies that $`a_{2n+\alpha }`$, as a function of $`Z`$, is holomorphic outside the locus. Hence it has an extension to a holomorphic function on $`^g`$. The $`\lambda `$-periodic wave functions corresponding to $`Z`$ and $`Z+\lambda ^{}`$ for any $`\lambda ^{}\mathrm{\Lambda }`$ are related to each other by a $`t`$-independent, $`_U`$-invariant factor. Hence, $`a_{2n+\alpha }`$ is periodic with respect to $`\mathrm{\Lambda }`$ and therefore is $`Z`$-independent. Note that $`a_{2s+1,n}=0`$ and $`a_{2s,n}=c_{s,n}`$ if $`sn`$. The lemma is proved.
The operator $`L_m`$ can be regarded as a $`(Z,x)`$-parametric family of ordinary differential operators $`L_m^Z`$ whose coefficients have the form
$$L_m^{Z,x}=_t^m+\underset{i=1}{\overset{m}{}}u_{i,m}(Ux+Vt+Z)_t^{mi},mI.$$
(4.33)
where $`u_{i,m}(Z)`$ are abelian function regular outside of $`\mathrm{\Theta }`$. For $`Z+Ux\mathrm{\Sigma }_{}`$ the coefficients of $`L_m^{Z,x}`$ are meromorphic functions of the variable $`t`$, which are not identically equal infinity. Recall, that $`\mathrm{\Sigma }_{}`$ is a $`_V`$-invariant set of $`\mathrm{\Theta }`$.
###### Corollary 4.2
The operators $`L_m^{Z,x}`$ commute with each other,
$$[L_n^{Z,x},L_m^{Z,x}]=0.$$
(4.34)
From (4.32) it follows that $`[L_n^{Z,x},L_m^{Z,x}]\psi =0`$. The commutator is an ordinary differential operator. Hence, the last equation implies (4.34).
###### Lemma 4.6
Let $`𝒜^{Z,x},Z+Ux\mathrm{\Sigma }_{},`$ be a commutative ring of ordinary differential operators spanned by the operators $`L_n^{Z,x}`$. Then there is an irreducible algebraic curve $`\mathrm{\Gamma }`$ such that $`𝒜^{Z,x}`$ is isomorphic to the ring $`A_{}(\mathrm{\Gamma },P_+,P_{})`$ of the meromorphic functions on $`\mathrm{\Gamma }`$ with the only pole at a smooth point $`P_{}`$ vanishing at another smooth point $`P_+`$. The correspondence $`Z𝒜^{Z,\mathrm{\hspace{0.17em}0}}`$ defines a holomorphic map of $`X\mathrm{\Sigma }_{}`$ into the space of torsion-free rank 1 sheaves $``$ on $`\mathrm{\Gamma }`$
$$j:X\backslash \mathrm{\Sigma }_{}\overline{\mathrm{Pic}}(\mathrm{\Gamma }).$$
(4.35)
On an open set the map $`j`$ is an imbedding.
The proof of the lemma is almost identical to the proof of lemma 3.4 in . It is the fundamental fact of the theory of commuting linear ordinary differential operators () that there is a natural correspondence
$$𝒜\{\mathrm{\Gamma },P_{},[k^1]_1,\}$$
(4.36)
between regular at $`t=0`$ commutative rings $`𝒜`$ of ordinary linear differential operators in the variable $`t`$, containing a pair of monic operators of co-prime orders, and sets of algebraic-geometrical data $`\{\mathrm{\Gamma },P_{},[k^1]_1,\}`$, where $`\mathrm{\Gamma }`$ is an algebraic curve with a fixed first jet $`[k^1]_1`$ of a local coordinate $`k^1`$ in the neighborhood of a smooth point $`P_{}\mathrm{\Gamma }`$ and $``$ is a torsion-free rank 1 sheaf on $`\mathrm{\Gamma }`$ such that
$$H^0(\mathrm{\Gamma },)=H^1(\mathrm{\Gamma },)=0.$$
(4.37)
The correspondence becomes one-to-one if the rings $`𝒜`$ are considered modulo conjugation $`𝒜^{}=g(t)𝒜g^1(t)`$.
Note, that in the main attention was paid to the generic case of the commutative rings corresponding to smooth algebraic curves. The invariant formulation of the correspondence given above is due to Mumford .
The algebraic curve $`\mathrm{\Gamma }`$ is called the spectral curve of $`𝒜`$. The ring $`𝒜`$ is isomorphic to the ring $`A(\mathrm{\Gamma },P_{})`$ of meromorphic functions on $`\mathrm{\Gamma }`$ with the only pole at the puncture $`P_{}`$. The isomorphism is defined by the equation
$$L_a\psi _0=a\psi _0,L_a𝒜,aA(\mathrm{\Gamma },P_{}).$$
(4.38)
Here $`\psi _0`$ is a common eigenfunction of the commuting operators. At $`t=0`$ it is a section of the sheaf $`𝒪(P_{})`$.
Important remark. The construction of the correspondence (4.36) depends on a choice of the initial point $`t_0=0`$. The spectral curve and the sheaf $``$ are defined by the evaluations of the coefficients of generators of $`𝒜`$ and a finite number of their derivatives at the initial point. In fact, the spectral curve is independent on the choice of $`t_0`$, but the sheaf does depend on it, i.e. $`=_{t_0}`$.
Using the shift of the initial point it is easy to show that the correspondence (4.36) extends to the commutative rings of operators whose coefficients are meromorphic functions of $`t`$ at $`t=0`$. The rings of operators having poles at $`t=0`$ correspond to sheaves for which the condition (4.37) is violated.
As it was mentioned above, the operators $`L_n`$, $`L_m`$ can be seen as a $`(Z,x)`$-parametric family of commuting ordinary differential operators in the variable $`t`$. Let $`\mathrm{\Gamma }^{Z,x}`$ be the corresponding spectral curve. The eigenvalues $`a_n(k)`$ of the operators $`L_n^{Z,x}`$ defined in (4.32) coincide with the Laurent expansions at $`P_{}`$ of the meromorphic functions $`a_nA(\mathrm{\Gamma }^{Z,x},P_{})`$. They are $`(Z,x)`$-independent. Hence, the spectral curve is $`(Z,x)`$-independent, as well, $`\mathrm{\Gamma }=\mathrm{\Gamma }^{Z,x}`$.
Equations (1.10), which are equivalent to (1.11) and are sufficient for the construction of the $`\lambda `$-periodic wave solutions, are symmetric with respect to $`x`$ and $`t`$. Therefore, the simple interchange of the variables $`x`$ and $`t`$ shows that if equations (1.10) hold then there exist commuting ordinary differential operators $`L_m^+`$ of the form
$$L_m^+=_x^m+\underset{i=1}{\overset{m}{}}u_{i,m}^+(Ux+Vt+Z)_x^{mi},mI^+,$$
(4.39)
where $`I^+`$ is the gap sequence associated with the variable $`x`$. These operators satisfy the equations
$$L_m^+H=B_m^+H,$$
(4.40)
where $`B_m^+`$ are differential operators in the variable $`x`$.
Let $`\psi `$ be a $`\lambda `$-periodic solution of (1.7). Then the same arguments as in the proof of Lemma 4.5 show that equations (4.40) imply
$$L_n^+\psi =a_n^+(k)\psi ,a_n^+=\underset{s=1}{\overset{\mathrm{}}{}}a_{s,n}^+k^s,$$
(4.41)
where $`a_{s,n}^+`$ are constants. From (4.41) it follows that the operators $`L_n,L_m^+`$ satisfy the equation
$$[L_n,L_m^+]=BH,$$
(4.42)
where $`B`$ is a differential operator in the variables $`x,t`$. Equation (4.41) also implies that there exists a polynomial $`\stackrel{~}{R}`$ such that $`\stackrel{~}{R}(L_n,L_m^+)\psi =0`$, i.e. eigenvalues $`a_n,a_m^+`$ of $`L_n`$ and $`L_m^+`$ satisfy the equation $`\stackrel{~}{R}(a_n,a_m^+)=0`$. Therefore, the spectral curves of commutative rings $`𝒜^Z`$ and $`𝒜_+^Z`$ coincide. Note, that ((4.41) implies that $`a_m^+`$ vanishes at $`P_{}`$. The symmetry between $`x`$ and $`t`$ variables implies that the ring $`𝒜_+^Z`$ is isomorphic to the ring $`A_+(\mathrm{\Gamma },P_+,P_{})`$ of meromorphic functions on $`\mathrm{\Gamma }`$ with the only pole at $`P_+`$ that vanish at $`P_{}`$.
Let us fix $`x=0`$ and consider the commuting operators $`L_n^Z=L_n^{Z,0}`$. The construction of the correspondence (4.36) implies that if the coefficients of the operators in $`𝒜`$ holomorphically depend on parameters, then the algebraic-geometrical spectral data are also holomorphic functions of the parameters.
Therefore, $`j`$ is holomorphic out of $`\mathrm{\Theta }`$. Then, using the shift of the initial point and the fact, that $`_{t_0}`$ holomorphically depends on $`t_0`$, we get that $`j`$ holomorphically extends on $`\mathrm{\Theta }\mathrm{\Sigma }_{}`$, as well.
The theta-divisor is not invariant under the shifts by constant vectors. Hence, for the generic $`Z`$ and $`Z^{}`$ the operators in $`𝒜^Z`$ and $`𝒜^Z^{}`$ have different poles. Hence, those rings do not coincide. Thus, the map $`j`$ is an imbedding on an open set. The lemma is proved.
It implies the global existence of the wave function.
###### Lemma 4.7
Let equations (1.10) hold. Then there exists a common eigenfunction of the operators $`L_n^Z`$ of the form $`\psi =e^{kt}\varphi (Vt+Z,k)`$ such that the coefficients of the formal series
$$\varphi (Z,k)=1+\underset{s=1}{\overset{\mathrm{}}{}}\xi _s(Z)k^s$$
(4.43)
are global meromorphic functions with a simple pole at $`\mathrm{\Theta }`$.
The proof of the lemma is identical to the proof of lemma 3.5 in . The function $`\psi `$ is first defined for $`Z\mathrm{\Sigma }_{}`$ as the inverse image $`\psi =j^{}\widehat{\psi }_{BA}`$ of the Baker-Akhiezer function, which is known to be globally defined on $`\overline{\mathrm{Pic}}(\mathrm{\Gamma })`$. Then, Hartogs’ extension theorem implies that $`\psi `$ has a meromorphic extension on $`C^g`$. The Baker-Akhiezer function is regular out of divisor corresponding to the commutative rings of operators whose coefficients have poles at $`t=0`$. Hence, $`\psi `$ is holomorphic out of $`\mathrm{\Theta }`$.
Let us show now that the correspondence $`\psi \psi ^{}`$ defines an involution of $`\mathrm{\Gamma }`$ under which $`P_\pm `$ are fixed.
###### Lemma 4.8
The eigenvalues $`a_{2n+\alpha }`$ of the commuting operators $`L_{2n+\alpha }`$ satisfy the relation
$$a_{2n+\alpha }(k)=(1)^\alpha a_{2n+\alpha }(k).$$
(4.44)
Proof. From equations (4.10,4.22,4.31) it follows that
$$\psi _t^{}(L_{2n+\alpha }\psi )=a_{2n+\alpha }(k)(\psi _t^{}\psi ),$$
(4.45)
$$(\psi _t^{}L_{2n+\alpha })\psi =\left((L_{2n+\alpha }^{}\psi _t^{})\psi \right)=(1)^\alpha a_{2n+\alpha }(k)(\psi _t^{}\psi ).$$
(4.46)
The left and right action of pseudo-differential operators are formally adjoint, i.e., for any two operators the equality $`\left(e^{kt}𝒟_1\right)\left(𝒟_2e^{kt}\right)=e^{kt}\left(𝒟_1𝒟_2e^{kt}\right)+_t\left(e^{kt}\left(𝒟_3e^{kt}\right)\right)`$ holds. Here $`𝒟_3`$ is a pseudo-differential operator whose coefficients are differential polynomials in the coefficients of $`𝒟_1`$ and $`𝒟_2`$. Therefore, equations (4.45,4.46) imply
$$\left(a_{2n+\alpha }(k)(1)^\alpha a_{2n+\alpha }(k)\right)(\psi _t^{}\psi )=_tQ_{2n+\alpha }.$$
(4.47)
The coefficients of the series $`Q_{2n+\alpha }`$ are differential polynomials on the coefficients of the wave operator $`\mathrm{\Phi }`$ defined by equation (4.5). For the globally defined wave function $`\psi `$, which exists according to the previous lemma, the coefficients of the wave operator are global meromorphic functions. Hence,
$$Q_{2n+\alpha }=\underset{s=1}{\overset{\mathrm{}}{}}Q_{2n+\alpha ,s}(Vt+Z),$$
(4.48)
where $`Q_{2n+\alpha ,s}(Z)`$ are meromorphic functions regular out of $`\mathrm{\Theta }`$.
In a similar way we have
$$\psi _t^{}\psi =\left(e^{kt}_t\mathrm{\Phi }^1\right)\left(\mathrm{\Phi }e^{kt}\right)=k+_tQ^{(1)}.$$
(4.49)
The series $`Q^{(1)}`$ has the form
$$Q^{(1)}=\underset{s=1}{\overset{\mathrm{}}{}}Q_s^{(1)}(Vt+Z),$$
(4.50)
where $`Q_s^{(1)}(Z)`$ are meromorphic functions regular out of $`\mathrm{\Theta }`$.
Let us fix a neighborhood of the theta-divisor in $`X`$. It defines the neighborhood $`S`$ of $`\mathrm{\Theta }`$ in $`^{\widehat{g}}`$. Outside of $`S`$ the functions $`Q_s^{(1)}`$ are bounded. Consider a sequence of real numbers $`l_i\mathrm{}`$ such that $`Z\pm Vl_i`$ is not in $`S`$. Then, from (4.49) it follows that
$$\psi _t^{}\psi =\underset{l_i\mathrm{}}{lim}\frac{1}{2l_i}_{l_i}^{l_i}(\psi _t^{}\psi )𝑑t=k.$$
(4.51)
The integration (4.51) is taken along a curve connecting points $`Z+\pm Vl_i`$ and which does not intersect $`\mathrm{\Theta }`$.
The same arguments imply that under ”averaging” in $`t`$ the right hand side of (4.47) vanishes. Hence, (4.47) implies (4.44). The lemma is proved.
The series $`a_n(k)`$ are the expansions at $`P_{}`$ of meromorphic functions on $`\mathrm{\Gamma }`$. Therefore, from (4.44) it follows there exists a holomorphic involution $`\sigma :\mathrm{\Gamma }\mathrm{\Gamma }`$ of the spectral curve such that
$$a_n^\sigma =a_n(\sigma (P))=(1)^na_n(P)$$
(4.52)
The point $`P_{}`$ is fixed under $`\sigma `$ and the local parameter is odd with respect to $`\sigma `$, i.e. $`\sigma ^{}k=k`$. In the same way using $`x`$ variable instead of $`t`$ we get that the second puncture $`P_+\mathrm{\Gamma }`$ is also fixed under $`\sigma `$.
The involution $`\sigma `$ induces an involution on the generalized Jacobian $`J(\mathrm{\Gamma })`$ which is by definition is the group of the equivalence classes of zero-degree divisors on $`\mathrm{\Gamma }`$, i.e. $`J(\mathrm{\Gamma })=\mathrm{Pic}^0(\mathrm{\Gamma })`$. The odd subgroup of $`J(\mathrm{\Gamma })`$ with respect to the induced involution $`\sigma ^{}`$ is the Prym variety of the spectral curve, $`𝒫(\mathrm{\Gamma })=\mathrm{ker}(1+\sigma ^{})`$. Our next goal is to show that $`𝒫(\mathrm{\Gamma })`$ of the spectral curve is compact.
###### Lemma 4.9
There exist $`g`$-dimensional vectors $`V_{2m+1}=\{V_{2m+1,k}\}`$ and constants $`v_{2m+1}`$ such that
$$F_{2m+1}^{(1)}(Z)=\underset{k=1}{\overset{g}{}}V_{2m+1,k}_Vh_k(Z)+v_{2m+1},$$
(4.53)
where $`F_{2m+1}^{(1)}=\mathrm{res}_{}^{2m+1}`$ and $`h_k=_{z_k}\mathrm{ln}\theta (Z)`$.
Proof. From equations (4.28, 4.29) and (4.49) it follows that $`F_{2m+1}^{(1)}=2_VQ_{2m+1}^{(1)}`$, where $`Q_{2m+1}^{(1)}`$ is a meromorphic function with a pole along $`\mathrm{\Theta }`$. The function $`F_{2m+1}^{(1)}`$ is an abelian function. Hence, for any vector $`\lambda `$ in the period lattice $`Q_{2m+1}^{(1)}(Z+\lambda )=Q_{2m+1}^{(1)}(Z)+c_{m,\lambda }`$. There is no abelian function with a simple pole on $`\mathrm{\Theta }`$. Hence, there exists a constant $`q_n`$ and $`g`$-dimensional vectors $`l_m`$ and $`V_{2m+1}`$, such that $`Q_{2m+1}^{(1)}=q_m+(l_m,Z)+(V_{2m+1},h(Z))`$, where $`h(Z)`$ is a vector with the coordinates $`h_k(Z)`$. Therefore, $`F_{2m+1}^{(1)}=(l_m,V)+(V_{2m+1},_Vh(Z))`$. The lemma is proved.
In order to complete the proof of our main result we need few more facts of the KP theory: flows of the KP hierarchy define deformations of the commutative rings $`𝒜`$ of ordinary linear differential operators. The spectral curve is invariant under these flows. If a commutative ring $`𝒜`$ of linear ordinary differential operators is maximal, i.e., it is not contained in any bigger commutative ring, then the KP orbit of $`𝒜`$ is isomorphic to the generalized Jacobian $`J(\mathrm{\Gamma })=\mathrm{Pic}^0(\mathrm{\Gamma })`$ of the spectral curve of $`𝒜`$ (see details in ).
The KP hierarchy in the Sato form is a system of commuting differential equation for a pseudo-differential operator $``$
$$_{t_n}=[_+^n,].$$
(4.54)
If the operator $``$ is as above. i.e., if it is defined by $`\lambda `$-periodic wave solutions of equation (1.7), then equation (4.14) implies that for odd $`n`$ equations (4.54) are equivalent to the equations
$$_{t_{2n+1}}u=_xF_{2m+1}^{(1)}(Ux+Vt+Z).$$
(4.55)
The first time of the hierarchy is identified with the variable $`t_1=t`$.
Equations (4.55) identify the space generated by the functions $`_UF_{2m+1}^{(1)}`$ with the tangent space at $`𝒜^Z`$ of the orbit of the part of the NV hierarchy associated with the puncture $`P_{}`$. In terms of $`u`$ the deformation with respect to $`z_i`$ is given by the equation
$$_{z_i}u=_x_Vh_i,h_i=_{z_i}\mathrm{ln}\theta (Z).$$
(4.56)
Equations (4.55, 4.56) and (4.53) imply
$$_{t_{2n+1}}=_{V_{2n+1}}=\underset{k=1}{\overset{g}{}}V_{2n+1,k}_{z_k}.$$
(4.57)
Hence, the orbit of $`𝒜^Z`$ is isomorphic to the factor of $`Z+Y/T(Z)`$ of the affine subvariety $`Z+YX`$, where $`Y`$ is the closure in $`X`$ of the subgroup $`_nV_{2n+1}t_{2n+1}`$, and $`T(Z)`$ is a lattice in the universal cover of $`Y`$.
###### Lemma 4.10
For the generic $`Z\mathrm{\Sigma }_{}`$, the orbit of $`𝒜^Z`$ under the NV flows defines an isomorphism:
$$i_Z:𝒫(\mathrm{\Gamma })Z+YX.$$
(4.58)
Proof. Recall, that according to the NV orbit of a maximal commutative ring is isomorphic to the Prym variety of the corresponding spectral curve. The arguments showing that $`𝒜^Z`$ is maximal for the generic $`Z`$ are identical to those used in . Indeed, suppose that $`𝒜^Z`$ is not maximal for all $`Z`$. Then there exits $`2n+\alpha I`$, where $`I`$ is the gap sequence defined above, such that for each $`Z\mathrm{\Sigma }_{}`$ there exists an operator $`L_{2n+\alpha }^Z`$ of order $`2n+\alpha `$ which commutes with all the operators $`L_m^Z𝒜^Z`$. Therefore, it commutes with $``$. That implies the equality
$$F_{2n+\alpha }^{(\alpha )}(Z)+\underset{i=1}{\overset{n}{}}c_{i,n}^{(\alpha )}(Z)F_{2n2i+\alpha }^{(\alpha )}(Z)=0.$$
(4.59)
Note the difference between (4.30) and (4.59). In the first equation the coefficients $`c_{i,n}^{(\alpha )}`$ are constants.
The $`\lambda `$-periodic wave solution of equation (1.7) is a common eigenfunction of all commuting operators, i.e. $`L_{2n+\alpha }\psi =a_{2n+\alpha }(Z,k)\psi `$, where is $`_V`$-invariant. The compactness of $`X`$ implies that $`a_{2n+\alpha }`$ is $`Z`$-independent. The first $`n`$ coefficients of $`a_{2n+\alpha }`$ coincide with the coefficients in (4.59). Hence, these coefficients are $`Z`$-independent. That contradicts the assumption that $`2n+\alpha I`$.
The map $`j`$ defined in Lemma 4.6 restricted to $`Z+YX`$ is inverse to $`i_Z`$. For the generic $`Z`$ it is an imbedding. Hence for the generic $`Z`$ the lattice $`T(Z)`$ is trivial. The lemma is thus proven.
###### Corollary 4.3
The Prym variety $`𝒫(\mathrm{\Gamma })`$ of the spectral curve $`\mathrm{\Gamma }`$ is compact.
The compactness of the Prym variety is not as restrictive, as the compactness of the Jacobian (see ). Nevertheless, it implies an explicit description of the the singular points of the spectral curve. The proof of the following statement is due to Robert Friedman and is presented in the Appendix.
###### Corollary 4.4
(R. Friedman) The spectral curve $`\mathrm{\Gamma }`$ is smooth outside of fixed points $`P_\pm ,Q_k`$ of the involution $`\sigma `$. The branches of $`\mathrm{\Gamma }`$ at $`Q_k`$ are linear and are not permuted by $`\sigma `$.
An equivalent formulation of the corollary is as follows: there is a a smooth algebraic curve $`\stackrel{~}{\mathrm{\Gamma }}`$ with involution $`\stackrel{~}{\sigma }`$ and a regular equivariant map $`p:\stackrel{~}{\mathrm{\Gamma }}\mathrm{\Gamma }`$ which is one-to-one out of preimages $`Q_k^i,i=1\mathrm{},\nu _k,`$ on $`\stackrel{~}{\mathrm{\Gamma }}`$ of the singular points $`Q_k`$.
The common eigenfunction of commuting differential operators is well-defined up to a constant factor for all smooth points of the spectral curve. It can be analytically extended along the branches of the spectral curve passing through the singular points, i.e. the preimage $`\stackrel{~}{\psi }`$ of the Baker-Akhiezer on $`\stackrel{~}{\mathrm{\Gamma }}`$ can be regarded as a section of a line bundle on $`\stackrel{~}{\mathrm{\Gamma }}`$. From the construction of the correspondence (4.36) it follows that the evaluations of $`\stackrel{~}{\psi }`$ at the preimages of the singular points $`Q_k`$ satisfy linear relations
$$\underset{i=1}{\overset{\nu _k}{}}c_{k,j}^i\stackrel{~}{\psi }(t,Q_k^i)=0,j=1,\mathrm{},n_k.$$
(4.60)
The coefficients of these relations and the zero divisor $`D`$ of $`\stackrel{~}{\psi }(0,\stackrel{~}{P})`$ can be regarded the data defining the corresponding sheaf $``$. The divisor $`D`$ is the pole divisor of the normalized eigenfunction $`\stackrel{~}{\psi }_0(t,\stackrel{~}{P})=\stackrel{~}{\psi }(t,\stackrel{~}{P})/\stackrel{~}{\psi }(0,\stackrel{~}{P})`$.
The following theta-functional formula (4.64) for $`\stackrel{~}{\psi }`$ is crucial for the final steps of the proof. First note that using the transformation $`\psi e^{(l(k),Z)}\psi `$, where $`l(k)`$ is a series such that $`(l(k),V)=0`$, we may assume without loss of generality that the series $`\varphi `$ in (4.43) satisfies the following monodromy properties:
$$\varphi (Z+e_j,k)=\varphi (Z),\varphi (Z+B_j)=\varphi (Z)\rho _j(k)$$
(4.61)
where $`e_j`$ are the basis vectors in $`^g`$ and $`B_j`$ are vectors defined by the columns of the matrix $`B`$, corresponding to the principle polarization of $`X`$.
Equations (4.61) and the fact the the coefficients of $`\varphi `$ are meromorphic functions with simple poles along $`\mathrm{\Theta }`$ imply that there is a series $`A(k)`$ such that
$$\varphi =\frac{\theta (A(k)+Z)}{\theta (Z)}$$
(4.62)
The series $`A(k)`$ defines an imbedding of the neighborhood of $`P_{}`$ into $`X`$.
The same arguments show that there is a holomorphic map
$$\stackrel{~}{A}:\stackrel{~}{\mathrm{\Gamma }}X$$
(4.63)
such that the function $`\stackrel{~}{\psi }(t,\stackrel{~}{P}),\stackrel{~}{P}\stackrel{~}{\mathrm{\Gamma }},`$ given by the formula
$$\stackrel{~}{\psi }=\frac{\theta (\stackrel{~}{A}(\stackrel{~}{P})+Vt+Z)}{\theta (Vt+Z)}e^{t\mathrm{\Omega }(\stackrel{~}{P})},$$
(4.64)
is the common eigenfunction of the operators in $`𝒜^Z`$. Here $`\mathrm{\Omega }(\stackrel{~}{P})`$ is an abelian integral on $`\stackrel{~}{\mathrm{\Gamma }}`$ having the form $`\mathrm{\Omega }=k+O(k^1)`$ at $`P_{}`$. Then the normalized eigenfunction of the commuting operators is given by the formula
$$\stackrel{~}{\psi }_0=\frac{\theta (\stackrel{~}{A}(\stackrel{~}{P})+Vt+Z)\theta (Z)}{\theta (\stackrel{~}{A}(\stackrel{~}{P})+Z)\theta (Vt+Z)}e^{t\mathrm{\Omega }(\stackrel{~}{P})}.$$
(4.65)
Our next goal is to show that the pole divisor $`D`$ of $`\stackrel{~}{\psi }_0`$ satisfies the condition analogous to (2.15) found by Novikov and Veselov in the case of the smooth spectral curves.
###### Lemma 4.11
The equivalence class of $`[D]J(\stackrel{~}{\mathrm{\Gamma }})`$ of the divisor $`D`$ satisfies the equation
$$[D]+[\stackrel{~}{\sigma }(D)]=K+P_++P_{}+\underset{k,i}{}Q_k^iJ(\stackrel{~}{\mathrm{\Gamma }}),$$
(4.66)
where $`K`$ is the canonical class, i.e. the equivalence class of the zero-divisor of a holomorphic differential on $`\stackrel{~}{\mathrm{\Gamma }}`$.
Proof. Equation (4.66) is equivalent to the condition that the divisor $`D+\sigma (D)`$ is the zero divisor of a meromorphic differential $`d\mathrm{\Omega }`$ on $`\stackrel{~}{\mathrm{\Gamma }}`$ with simple poles at the punctures $`P_\pm `$ and the points $`Q_k^i`$. The differential $`d\mathrm{\Omega }`$ is even with respect to the involution and descends to a meromorphic differential on the factor-curve $`\mathrm{\Gamma }_0`$.
The existence of such differential can be proved almost identically to the proof of the statement that the conditions (2.15) are necessary conditions for the potential reduction of the $`2D`$ Schrödinger operators given in (Theorem 3.1).
Let $`\stackrel{~}{\psi }_0(x,t,P)`$ be the normalized solution of the Schrödinger operator. It is obtained by the deformation along the $`x`$-flow from the normalized eigenfunction of the operators in $`𝒜^Z`$ considered above. Therefore, it has the form (2.16) with $`\theta _{pr}`$ and $`A^{pr}`$ replaced by $`\theta `$ and $`\stackrel{~}{A}`$, respectively. Following we present another real form of $`\stackrel{~}{\psi }_0`$. Let us introduce real coordinates of a complex vector $`Z^g`$ by the formula $`Z=\zeta ^{}+B\zeta ^{\prime \prime }`$, where $`\zeta ^{},\zeta ^{\prime \prime }`$ are $`g`$-dimensional real vectors, and $`B`$ is the matrix of $`b`$-periods of the normalized holomorphic differentials on $`\mathrm{\Gamma }`$. Then the absolute value $`|\varphi |`$ of the function $`\varphi (\zeta ,P),\zeta =(\zeta ^{},\zeta ^{\prime \prime })`$ given by the formula
$$\varphi (\zeta ,P)=\frac{\theta (\stackrel{~}{A}(P)+Z)}{\theta (Z)}e^{2\pi i(\stackrel{~}{A}(P),\zeta ^{\prime \prime })}$$
(4.67)
is a periodic function of the coordinates $`\zeta _k^{},\zeta _k^{\prime \prime }`$. For real $`x,t`$ the function $`\psi `$ can be represented in the form
$$\stackrel{~}{\psi }_0=\frac{\varphi (\widehat{U}x+\widehat{V}t+\zeta ,P)}{\varphi (\zeta ,P)}e^{tp_{}+xp+}.$$
(4.68)
Here $`\widehat{U}=(U_1^+^{},U_1^{+^{\prime \prime }}),\widehat{V}=(U_1^{^{}},U_1^{^{\prime \prime }})`$ are $`2g`$-dimensional real vectors corresponding to the complex vectors $`U,V`$; $`dp_\pm `$ are meromorphic differentials on $`\stackrel{~}{\mathrm{\Gamma }}`$ with poles of the second order at $`P_\pm `$ and whose periods are pure imaginary.
The differential $`d\stackrel{~}{\psi }_0`$ is also a solution of the same Schrödinger equation. That implies the equality
$$_x(_t\stackrel{~}{\psi }_0^\sigma d\stackrel{~}{\psi }_0\stackrel{~}{\psi }_0^\sigma _td\stackrel{~}{\psi }_0)=_t(_x\stackrel{~}{\psi }_0^\sigma d\stackrel{~}{\psi }_0\stackrel{~}{\psi }_0^\sigma _xd\stackrel{~}{\psi }_0)$$
(4.69)
The ”averaging” of this equation in the variables $`x,t`$ gives the equation
$$_t\stackrel{~}{\psi }_0^\sigma \stackrel{~}{\psi }_0\stackrel{~}{\psi }_0^\sigma _t\stackrel{~}{\psi }_0_xdp_{}=_x\stackrel{~}{\psi }_0^\sigma \stackrel{~}{\psi }_0\stackrel{~}{\psi }_0^\sigma _x\stackrel{~}{\psi }_0_tdp_+.$$
(4.70)
Here $`_t`$ stands for the mean value in $`t`$ defined as in (4.51), and $`_x`$ stands for the mean value in $`x`$ defined in a similar way. The same arguments as in , show that the differential
$$d\mathrm{\Omega }=\frac{dp_+}{_t\stackrel{~}{\psi }_0^\sigma \stackrel{~}{\psi }_0\stackrel{~}{\psi }_0^\sigma _t\stackrel{~}{\psi }_0_x}=\frac{dp_{}}{_x\stackrel{~}{\psi }_0^\sigma \stackrel{~}{\psi }_0\stackrel{~}{\psi }_0^\sigma _x\stackrel{~}{\psi }_0_t}$$
(4.71)
is holomorphic on $`\stackrel{~}{\mathrm{\Gamma }}`$ except at the branch points where it has simple poles. It has zeros at the poles of $`\stackrel{~}{\psi }_0`$ and $`\stackrel{~}{\psi }_0^\sigma `$. The lemma is proved.
The differential $`\stackrel{~}{\psi }_0\stackrel{~}{\psi }_0^\sigma d\mathrm{\Omega }`$ is a meromorphic differential on $`\stackrel{~}{\mathrm{\Gamma }}`$. Its residues at the points $`P_\pm `$ are equal to $`\pm 1`$, respectively. Therefore, sum of its residues at the points $`Q_k^i`$ equals zero, i.e.,
$$\underset{i,k}{}\stackrel{~}{r}_k^i\stackrel{~}{\psi }_0(t,Q_k^i)\stackrel{~}{\psi }_0^\sigma (t,Q_k^i)=0,\stackrel{~}{r}_k^i=\mathrm{res}_{Q_k^i}d\mathrm{\Omega }.$$
(4.72)
Note, that equation (4.72) is sufficient for the potential reduction of the Schrödinger operator, which is equivalent to the equation $`\stackrel{~}{\psi }_0(t,P_+)=1`$.
From (4.60) it follows that the evaluations of $`\stackrel{~}{\psi }_0`$ at the points $`Q_k^i`$ satisfy linear equations
$$\underset{i=1}{\overset{\nu _k}{}}\stackrel{~}{c}_{k,j}^i\stackrel{~}{\psi }_0(t,Q_k^i)=0,j=1,\mathrm{},n_k.$$
(4.73)
Note, that the normalization of $`\psi _0`$ at $`t=0`$ implies
$$\underset{i=1}{\overset{\nu _k}{}}\stackrel{~}{c}_{k,j}^i=0.$$
(4.74)
The coefficients $`\stackrel{~}{c}_{k,j}^i`$ in (4.73) are unique up to the transformations $`\stackrel{~}{c}_{k,j}^i_lg_{k,j}^l\stackrel{~}{c}_{k,l}^i`$, where $`g_k=\{g_{k,j}^l\}`$ are $`t`$-independent non-degenerate matrices. In what follows we normalize $`\stackrel{~}{c}_k`$ by the condition
$$\stackrel{~}{c}_{k,j}^i=\delta _j^i,i=1,\mathrm{},n_k.$$
(4.75)
In that gauge the matrix elements $`\stackrel{~}{c}_{k,j}^i(Z),i>n_k`$ become well-defined abelian functions on $`X`$.
Equations (4.72) should follow from equations (4.73). The evaluations of $`\stackrel{~}{\psi }_0`$ at the preimages of any two distinct singular points $`Q_k,Q_k^{},kk^{}`$ are independent. That implies the following orthogonality relations
$$\underset{j=1}{\overset{n_k}{}}\stackrel{~}{r}_k^j\stackrel{~}{c}_{k,j}^i\stackrel{~}{c}_{k,j}^i=\stackrel{~}{r}_k^i,\underset{j=1}{\overset{n_k}{}}\stackrel{~}{r}_k^j\stackrel{~}{c}_{k,j}^i\stackrel{~}{c}_{k,j}^i^{}=0,n_k<ii^{}\nu _k.$$
(4.76)
(Compare (4.76) with the orthogonality conditions established in ()).
###### Corollary 4.5
The multiplicity $`\nu _k`$ of the singular point $`Q_k`$ of the spectral curve is equal to $`\nu _k=2n_k`$, where $`n_k`$ is the number of the linear relations in (4.73).
Proof. As it was shown above, for the generic $`Z`$ the ring $`𝒜^Z`$ is maximal. The ring $`𝒜^Z`$ is maximal if and only if for each $`k`$ the linear subspace $`W_k^{2\nu _k}`$ defined by the equations $`_i\stackrel{~}{c}_{k,j}^i(Z)w_k^i=0`$ is invariant under the multiplication by a diagonal matrix $`H_k`$ only if $`H_k`$ is a scalar matrix. The last condition implies that each $`(n_k\times n_k)`$ minor of the matrix $`c_k=\{\stackrel{~}{c}_{k,j}^i\}`$ with $`n_k<i\nu _k`$ is non-degenerate. The columns of the matrix $`\stackrel{~}{c}_{k,j}^i`$ with $`i>n_k`$ are ”orthogonal” to each other. Then, non-degeneracy of all the corresponding minors implies that the number $`\nu _kn_k`$ of such columns is not bigger than the dimension $`n_k`$ of the column vectors, i.e. $`\nu _k2n_k`$.
From (4.66) it follows that the degree of the pole divisor $`D`$ equals $`\mathrm{deg}D=\stackrel{~}{g}+1/2_k\nu _k`$, where $`\stackrel{~}{g}`$ is the genus of $`\stackrel{~}{\mathrm{\Gamma }}`$. The uniqueness of the function $`\psi _0`$ defined by $`D`$ and the relations (4.73), imply that $`\mathrm{deg}D=\stackrel{~}{g}+_kn_k`$. The latter equations imply $`_k(2n_k\nu _k)=0`$. As shown above, each term of the sum is non-negative. Hence, $`\nu _k=2n_k`$ and the corollary is thus proven.
###### Lemma 4.12
There exist constants $`r_k^i`$ such that the equation
$$\underset{i=1}{\overset{2n_k}{}}r_k^i\theta (A_k^i+Z)\theta (A_k^iZ)=0,A_k^i=\stackrel{~}{A}(Q_k^i),$$
(4.77)
holds.
Proof. Taking the square of (4.74) and using (4.76) we get the equation
$$\underset{i=1}{\overset{2n_k}{}}\stackrel{~}{r}_k^i=0.$$
(4.78)
The residues of the differential $`d\mathrm{\Omega }`$ are well-defined abelian functions $`\stackrel{~}{r}_k^i(Z)`$ on $`X`$. The pole divisor of all the residues coincides with the zero divisor of $`\theta `$. The residue of $`d\mathrm{\Omega }`$ at $`Q_k^i`$ equals zero, when the pole divisor of $`\stackrel{~}{\psi }_0`$ contains the puncture $`P_{}`$. Therefore, from (4.65) it follows that $`\stackrel{~}{r}_k^i`$ has the form
$$\stackrel{~}{r}_k^i(Z)=r_k^i\frac{\theta (A_k^i+Z)\theta (A_k^iZ)}{\theta ^2(Z)},$$
(4.79)
where $`r_k^i`$ are constants. Equations (4.78) and (4.79) imply (4.77) and the lemma is proved.
Our next and the final goal is to show that $`n_k=1`$, i.e. all of the singular points of $`\mathrm{\Gamma }`$ are simple double points, as it is stated in the main theorem.
If $`n_k>1`$, then from indecomposability of the matrix $`\stackrel{~}{c}_{k,j}^i(Z)`$ at the generic $`Z`$ it follows that all the points $`A_k^i`$ are distinct, $`A_k^iA_k^j`$. That and the formula (4.65) for $`\stackrel{~}{\psi }_0`$ imply that in the gauge (4.75) the coefficient $`\stackrel{~}{c}_{k,j}^i`$ for $`i>n_k`$ has pole at the divisor $`\theta (A_k^j+Z)=0`$ and zero at the divisor $`\theta (A_k^i+Z)=0`$.
Let us fix a pair of indices $`m,l>n_k`$ and define a set $`𝒟_k^{m,l}X`$ by the equations:
$$\stackrel{~}{c}_{k,1}^i(Z)=0,n_k<im,l.$$
(4.80)
On $`𝒟_k^{m,l}`$ equation (4.74) takes the form
$$1+\stackrel{~}{c}_{k,1}^m(Z)+\stackrel{~}{c}_{k,1}^l(Z)=0,Z𝒟_k^{m,l}.$$
(4.81)
From (4.80) and the orthogonality conditions (4.76) it follows that on $`𝒟_k^{m,l}`$ the equation $`\stackrel{~}{c}_{k,1}^l(Z)=0`$ implies
$$\stackrel{~}{c}_{k,j}^m(Z)=0,j=n_k+2,\mathrm{},2n_k.$$
(4.82)
Then, from (4.82) it follows that
$$Z𝒟_k^{m,l},\stackrel{~}{c}_{k,1}^l(Z)=0\stackrel{~}{r}_k^1+\stackrel{~}{r}_k^m=0.$$
(4.83)
Hence, $`\stackrel{~}{c}_{k,1}^l`$, restricted to $`𝒟_k^{m,l}`$, is of the form
$$\stackrel{~}{c}_{k,1}^l=\frac{\stackrel{~}{r}_k^1+\stackrel{~}{r}_k^m}{h(Z)}\frac{\theta (A_k^l+Z)}{\theta (A_k^1+Z)}\theta ^2(Z),Z𝒟_k^{m,l},$$
(4.84)
where $`h`$ is a holomorphic section of the line bundle of $`|2\mathrm{\Theta }A_k^l+A_k^1|`$ restricted to $`𝒟_k^{m,l}`$. (Recall, that $`\stackrel{~}{c}_{k,j}^i`$ and $`\stackrel{~}{r}_k^i`$ are abelian functions.)
The same arguments imply that on $`𝒟_k^{m,l}`$ zeros of $`\stackrel{~}{c}_{k,1}^m`$ are in the zero divisor of $`\stackrel{~}{r}_k^1+\stackrel{~}{r}_k^l`$. Then using (4.81) we get
$$\stackrel{~}{c}_{k,1}^m=g\frac{\stackrel{~}{r}_k^1+\stackrel{~}{r}_k^l}{h(Z)}\frac{\theta (A_k^m+Z)}{\theta (A_k^1+Z)}\theta ^2(Z),Z𝒟_k^{m,l},$$
(4.85)
where $`g`$ is a constant. Therefore, $`h`$ is a section of the restriction to $`𝒟_k^{m,l}`$ of the line bundle $`|2\mathrm{\Theta }A_k^m+A_k^1|`$. Hence, $`A_k^m=A_k^l`$. The choice of points $`A_k^m,A_k^l`$ was arbitrary, Therefore, we have proved that all the points $`A_k^i=A_k`$ do coincide. In that case, equations (4.73) are equivalent to $`(2n_k1)`$ equations $`\stackrel{~}{\psi }_0(t,Q_k^i)=\stackrel{~}{\psi }_0(t,Q_k^j)`$. That implies $`2n_k1=n_k=1`$, i.e., $`\mathrm{\Gamma }`$ has only simple double singular points. For such a curve all the sheafs $``$ are line bundles. Therefore, the map $`j`$ in (4.35) is inverse to $`i_Z`$ in (4.58) and the main theorem is thus proven.
## 5 Appendix.
###### Theorem 5.1
(R.Friedman) Let $`\mathrm{\Gamma }`$ be an irreducible projective curve, with an involution $`\sigma `$, and suppose that the generalized Prym variety $`P(\mathrm{\Gamma },\sigma )`$ is compact. Then every singular point $`x`$ of $`\mathrm{\Gamma }`$ is a fixed point of $`\sigma `$, the singularity at $`x`$ is locally analytically isomorphic to a union of coordinate axes in a neighborhood of the origin in $`^N`$ for some $`N`$, and in a neighborhood of such a singular point $`\sigma `$ fixes each of the local analytic branches.
Proof Let $`p:\stackrel{~}{\mathrm{\Gamma }}\mathrm{\Gamma }`$ be the normalization of $`\mathrm{\Gamma }`$. The involution $`\sigma `$ lifts to an involution on $`\stackrel{~}{\mathrm{\Gamma }}`$, also denoted $`\sigma `$. The sheaf $`p_{}(𝒪_{\stackrel{~}{\mathrm{\Gamma }}})/𝒪_\mathrm{\Gamma }`$ is supported at the finitely many points of $`\mathrm{\Gamma }_{\text{sing}}`$. The cohomology long exact sequence for
$$0𝒪_\mathrm{\Gamma }p_{}(𝒪_{\stackrel{~}{\mathrm{\Gamma }}})p_{}(𝒪_{\stackrel{~}{\mathrm{\Gamma }}})/𝒪_\mathrm{\Gamma }0$$
yields a long exact sequence
$`0H^0(\mathrm{\Gamma };𝒪_\mathrm{\Gamma })H^0(\mathrm{\Gamma };p_{}(𝒪_{\stackrel{~}{\mathrm{\Gamma }}}))=H^0(\stackrel{~}{\mathrm{\Gamma }};𝒪_{\stackrel{~}{\mathrm{\Gamma }}})H^0(\mathrm{\Gamma };p_{}(𝒪_{\stackrel{~}{\mathrm{\Gamma }}})/𝒪_\mathrm{\Gamma })`$
$`H^1(\mathrm{\Gamma };𝒪_\mathrm{\Gamma })H^1(p_{}(𝒪_{\stackrel{~}{\mathrm{\Gamma }}}))=H^1(\stackrel{~}{\mathrm{\Gamma }};𝒪_{\stackrel{~}{\mathrm{\Gamma }}})0.`$
Since $`\stackrel{~}{\mathrm{\Gamma }}`$ is irreducible, $`H^0(\stackrel{~}{\mathrm{\Gamma }};𝒪_{\stackrel{~}{\mathrm{\Gamma }}})=`$, so that there is an exact sequence
$$0H^0(\mathrm{\Gamma };p_{}(𝒪_{\stackrel{~}{\mathrm{\Gamma }}})/𝒪_\mathrm{\Gamma })H^1(\mathrm{\Gamma };𝒪_\mathrm{\Gamma })H^1(\stackrel{~}{\mathrm{\Gamma }};𝒪_{\stackrel{~}{\mathrm{\Gamma }}})0.$$
Here $`H^1(\mathrm{\Gamma };𝒪_\mathrm{\Gamma })`$ is the tangent space to the generalized Jacobian of $`\mathrm{\Gamma }`$ and the subspace $`H^0(\mathrm{\Gamma };p_{}(𝒪_{\stackrel{~}{\mathrm{\Gamma }}})/𝒪_\mathrm{\Gamma })`$ is the tangent space to its noncompact part. If $`V`$ is a vector space on which $`\sigma `$ acts, let $`V^{}`$ denotes the anti-invariant part of $`V`$, i.e. the $`(1)`$-eigenspace. Then the tangent space $`T_{P(\mathrm{\Gamma },\sigma )}`$ to $`P(\mathrm{\Gamma },\sigma )`$ fits into an exact sequence
$$0H^0(\mathrm{\Gamma };p_{}(𝒪_{\stackrel{~}{\mathrm{\Gamma }}})/𝒪_\mathrm{\Gamma })^{}T_{P(\mathrm{\Gamma },\sigma )}T_{P(\stackrel{~}{\mathrm{\Gamma }},\sigma )}0.$$
It follows that $`P(\mathrm{\Gamma },\sigma )`$ is compact if and only if $`H^0(\mathrm{\Gamma };p_{}(𝒪_{\stackrel{~}{\mathrm{\Gamma }}})/𝒪_\mathrm{\Gamma })^{}=0`$.
First let us show that, for all $`x\mathrm{\Gamma }_{\text{sing}}`$, $`\sigma (x)=x`$. There is an isomorphism
$$H^0(\mathrm{\Gamma };p_{}(𝒪_{\stackrel{~}{\mathrm{\Gamma }}})/𝒪_\mathrm{\Gamma })=\underset{x\mathrm{\Gamma }_{\mathrm{sing}}}{}\stackrel{~}{R}_x/R_x,$$
where $`R_x`$ is the local ring $`𝒪_{\mathrm{\Gamma },x}`$ and $`\stackrel{~}{R}_x`$ is its normalization. Clearly, $`\sigma `$ induces an isomorphism $`\stackrel{~}{R}_x/R_x\stackrel{~}{R}_{\sigma (x)}/R_{\sigma (x)}`$. If $`\sigma (x)x`$, and $`\alpha `$ is a nonzero element of $`\stackrel{~}{R}_x/R_x`$, then $`\alpha \sigma (\alpha )H^0(\mathrm{\Gamma };p_{}(𝒪_{\stackrel{~}{\mathrm{\Gamma }}})/𝒪_\mathrm{\Gamma })`$ is nonzero, a contradiction. Hence $`\sigma (x)=x`$. Moreover, for all $`\alpha \stackrel{~}{R}_x/R_x`$, $`\sigma (\alpha )=\alpha `$.
We now fix attention on a given $`x\mathrm{\Gamma }_{\text{sing}}`$, and write $`R=R_x`$ and $`\stackrel{~}{R}=\stackrel{~}{R}_x`$. Note that, if $`y_1,\mathrm{},y_n`$ are the preimages of $`x`$ in $`\stackrel{~}{\mathrm{\Gamma }}`$, and $`t_i`$ is a local analytic coordinate for $`\stackrel{~}{\mathrm{\Gamma }}`$ at $`y_i`$, then $`\stackrel{~}{R}_i\{t_i\}`$. Moreover, $`R`$ is a subalgebra of $`\stackrel{~}{R}`$, and $`\mathrm{dim}_{}(\stackrel{~}{R}/R)<\mathrm{}`$; on particular, $`\stackrel{~}{R}`$ is a finite $`R`$-module. Let $`𝔪_i=t_i\{t_i\}`$. Clearly, $`R`$ is contained in the subalgebra $`_i𝔪_i`$, and the claim about the analytic nature of the singularities is just the statement that $`R=_i𝔪_i`$.
Next we claim that $`\sigma `$ does not permute the analytic branches through $`x`$. If it did, then the action of $`\sigma `$ on $`\stackrel{~}{R}`$ would exchange two factors $`\{t_i\}`$ and $`\{t_j\}`$ for $`ji`$. In this case, let $`e_i`$ be the image of $`1\{t_i\}`$ in $`\stackrel{~}{R}`$. Then $`\sigma (e_i)e_i(\stackrel{~}{R})^{}`$, and $`\sigma (e_i)e_iR`$. Thus $`\sigma (e_i)e_i`$ is a nonzero class in $`(\stackrel{~}{R}/R)^{}=(\stackrel{~}{R})^{}/R^{}`$, a contradiction. Thus, for every $`i`$, $`\sigma `$ fixes $`\{t_i\}`$ and induces a holomorphic involution on the corresponding branch of $`\stackrel{~}{\mathrm{\Gamma }}`$.
We can thus choose the coordinate $`t_i`$ so that $`\sigma (t_i)=t_i`$. Since $`(\stackrel{~}{R}/R)^{}=(\stackrel{~}{R})^{}/R^{}=0`$, $`t_iR`$ for every $`i`$. Clearly, $`t_i𝔪`$, where $`𝔪`$ is the maximal ideal of the local ring $`R`$. Now let $`\widehat{R}=_i𝔪_i\stackrel{~}{R}`$. As noted above, $`R\widehat{R}`$ and we must show that $`R=\widehat{R}`$. In any case, $`\widehat{R}`$ is a finite $`R`$-module. Given $`r\widehat{R}`$, there exists a $`cR\widehat{R}`$ such that $`rc_i𝔪_i=(t_1,\mathrm{},t_n)\widehat{R}𝔪\widehat{R}`$. Thus $`\widehat{R}=R+𝔪\widehat{R}`$, and so by Nakayama’s lemma $`R=\widehat{R}`$.
Acknowledgments. The author wishes to thank Sam Grushevski for very useful conversation on the subject of this paper, Robert Friedman for the communication of the proof of Corollary 4.4. The author is grateful to Takahiro Shiota whose remarks helped the author to clarify some missing arguments in the first version of the paper. |
warning/0506/hep-th0506254.html | ar5iv | text | # Noncommutative 6D Gauge Higgs Unification Models
## I Introduction
A renewed interest in theories in 6D has recently emerged cgp . An anomaly free gauged supergravity in $`D=6`$, the Salam-Sezgin model ss , has been considered. This model is compactified on a 2-sphere and in four dimensions gives a $`SU(2)\times U(1)`$ gauge theory rss . In particular, it has been argued that these theories with 3-Branes could point out towards solving the cosmological constant problem quevedo . Also, in brandenberger it is shown that chaotic inflation consistent with constraints coming from the amplitude of the cosmic microwave anisotropies can be naturally realized.
In the search for a unified theory of elementary particles, the incorporation of the Higgs field in the standard model (SM) of electroweak interactions has motivated various proposals in 6D Fairlie . These are 6D pure gauge theories, in which after dimensional reduction the Higgs field naturally arises. Recently new proposals have been made, considering orbifold compactifications, in antoniadis2 , a $`U(3)\times U(3)`$ model has been considered. In this work the Higgs mass term is generated radiatively, with a finite value at one loop as the quadratic divergences are suppressed by the six dimensional gauge symmetry. Further, a $`SU(3)`$ model of this type has been developed in scruca ; Wulzer , with one Higgs doublet and a predicted W-boson mass of half the Higgs mass. In this case the weak angle has a non realistic value, although it can be improved by an extended gauge group as in antoniadis2 , or by the introduction of an $`U(1)`$ factor as done in scruca .
Noncommutativity in field theories has been the subject of an important number of works in the last few years. In particular, the Seiberg-Witten construction SW and its generalization for any gauge group W2 have been studied. This construction allows to express the noncommutative gauge fields in terms of the usual ones and their derivatives, maintaining the same degrees of freedom. It has been extended for noncommutative matter fields, which can also be generated in terms of the commutative matter fields and the gauge fields of interest W2 . By this procedure, noncommutative versions of the standard model and consequently the electroweak interaction sector have been given W3 (see also Aschieri ). As a consequence new interactions among the fields of the theory are predicted.
In this work, we will investigate the noncommutative generalization of the bosonic sector of Gauge Higgs unification models in 6D based on the $`SU(3)`$ gauge group compactified on $`T^2/Z_N`$ Wulzer . The noncommutative extension is obtained by means of the Seiberg-Witten map. We calculate, for the bosonic sector, the resulting first order corrections and compare them with the results obtained in other works. Assuming noncommutativity only between the extra dimensions, the phenomenological consequences are considered in the framework of the effective theories technique. First we compute the new physics effects corrections to the $`S`$ and $`T`$ oblique parameters. Further, from the experimental constraints we get a bound for the noncommutativity parameter $`\theta ^{45}`$. This bound allows us to calculate the correction to the decay width for the rare decay of the Higgs boson into two photons.
In section 2 we review the model of reference Wulzer , in section 3 the Seiberg-Witten map and its generalization to nonabelian groups is presented in some detail. In section 4 we present the noncommutative formulation of the model Wulzer and show that our results differ from those calculated directly in 4D. The phenomenological consequences of the noncommutativity between extra dimensions are discussed in section 5. Section 6 is devoted to conclusions.
## II The 6-Dimensional Model
### II.1 Gauge Fields in a 6-Dimensional Space-Time
Let us consider a Yang-Mills theory in 6-dimensional space-time with a $`SU(3)`$ gauge group, the Lagrangian of the theory is
$$=\frac{1}{2}\mathrm{Tr}F_{mn}F^{mn},$$
the field strength tensor is defined by
$$F_{mn}=_mA_n_nA_mig_6[A_m,A_n],$$
and $`g_6`$ is the coupling constant in 6D. This action is interpreted by a dimensional reduction on an orbifold $`T^2/Z_N`$ for $`N=3,4,6`$ Wulzer , by a separation of the connection in its 4-dimensional space-time part $`A_\mu `$, and the other two components, $`A_z`$ and $`A_{\overline{z}}`$ which in 4D will play the role of scalars, with $`z=\frac{1}{\sqrt{2}}\left(x^4+ix^5\right)`$ and $`\overline{z}=\frac{1}{\sqrt{2}}\left(x^4ix^5\right)`$. These fields are the zero modes of Kaluza-Klein and depend only on the four space-time coordinates $`x^\mu `$. The result of this reduction is given by
$$=\frac{1}{2}\mathrm{Tr}F_{\mu \nu }F^{\mu \nu }+2\mathrm{T}\mathrm{r}D_\mu A_{\overline{z}}D^\mu A_zg^2\mathrm{Tr}[A_z,A_{\overline{z}}]^2,$$
(1)
where $`g=g_6\sqrt{V}`$ is the gauge coupling of the 4-dimensional effective theory, $`V`$ is the volume of the two extra dimensions and
$`F_{\mu \nu }`$ $`=_\mu A_\nu _\nu A_\mu ig[A_\mu ,A_\nu ],`$
$`D_\mu A_z`$ $`=_\mu A_zig[A_\mu ,A_z]=F_{\mu z},`$ (2)
$`D_\mu A_{\overline{z}}`$ $`=_\mu A_{\overline{z}}ig[A_\mu ,A_{\overline{z}}]=F_{\mu \overline{z}}.`$
The orbifold reduction Wulzer for the gauge fields $`A_m`$ leads to: the 4-dimensional $`A_\mu ,`$ that containts four electroweak bosons, $`W_\mu SU\left(2\right)`$ and $`B_\mu U\left(1\right)`$,
$$A_\mu =\left(\begin{array}{cc}W_\mu & 0\\ 0& 0\end{array}\right)+\frac{1}{2\sqrt{3}}\left(\begin{array}{cc}B_\mu I& 0\\ 0& 2B_\mu \end{array}\right),$$
and the two complex components of the scalar boson doublet (Higgs), which are contained in the $`A_z`$ and $`A_{\overline{z}}`$ gauge fields,
$$A_z=\frac{1}{\sqrt{2}}\left(\begin{array}{cc}0\hfill & \varphi \hfill \\ 0\hfill & 0\hfill \end{array}\right),A_{\overline{z}}=\frac{1}{\sqrt{2}}\left(\begin{array}{cc}0\hfill & 0\hfill \\ \varphi ^{}\hfill & 0\hfill \end{array}\right).$$
Substituting these expressions in the Lagrangian (1) we find
$$=\frac{1}{2}\mathrm{Tr}F_{\mu \nu }(W)F^{\mu \nu }(W)\frac{1}{4}F_{\mu \nu }(B)F^{\mu \nu }(B)+\left(D_\mu \varphi \right)^{}(D^\mu \varphi )V\left(\varphi \right),$$
(3)
where $`D_\mu \varphi =\left(_\mu \frac{1}{2}igW_\mu ^a\tau _a\frac{1}{2}ig\mathrm{tan}\theta _WB_\mu \right)\varphi ,`$ $`\mathrm{tan}\theta _W=\sqrt{3}`$ and $`V\left(\varphi \right)=\frac{g^2}{2}\left|\varphi \right|^4.`$
Thus, this Lagrangian has a $`SU(2)\times U(1)`$ invariance, with a scalar massless doublet with a quartic potential. However, as shown in antoniadis1 , quantum fluctuations induce corrections to the potential $`V\left(\varphi \right)`$ which can trigger radiative symmetry breaking. The leading terms in the one-loop effective potential for the Higgs are,
$$V_{eff}\left(\varphi \right)=\mu ^2\left|\varphi \right|^2+\lambda \left|\varphi \right|^4.$$
Assuming $`\mu ^2>0,`$ so that electroweak symmetry breaking can occur, we have that $`\left|\varphi \right|=\nu /\sqrt{2}`$ with $`\nu =\mu /\sqrt{\lambda }`$. Note that the value of the electroweak angle in (3) is too large. However, as mentioned in the introduction, the model can be extended in such a way that it correctly reproduces the SM value.
## III Noncommutative Gauge Theories
### III.1 Noncommutative Space-Time
Noncommutative space-time incorporates coordinates $`\widehat{x}^\mu `$, given by operators that satisfy the following relations,
$$[\widehat{x}^\mu ,\widehat{x}^\nu ]=i\theta ^{\mu \nu },$$
(4)
where $`\theta ^{\mu \nu }=\theta ^{\nu \mu }`$ are real numbers. To construct a field theory in this space, it is more convenient to consider usual fields, which are functions. This is allowed by the Weyl-Wigner-Moyal correspondence, which establishes an equivalence between the Heisenberg algebra of the operators $`\widehat{x}^\mu `$ and the function algebra in $`^m`$. It has an associative and noncommutative star product, the Moyal $``$-product, given by,
$`f(x)g(x)`$ $`\left[\mathrm{exp}\left({\displaystyle \frac{i}{2}}{\displaystyle \frac{}{\epsilon ^\alpha }}\theta ^{\alpha \beta }{\displaystyle \frac{}{\eta ^\beta }}\right)f(x+\epsilon )g(y+\eta )\right]_{\epsilon =\eta =0}`$ (5)
$`=fg+{\displaystyle \frac{i}{2}}\theta ^{\alpha \beta }_\alpha f_\beta g+𝒪(\theta ^2).`$
Under complex conjugation it satisfies $`\left(\overline{fg}\right)=\overline{g}\overline{f}`$. Since we will work with a nonabelian gauge group, our functions are matrix valued, and the corresponding matrix Moyal product is denoted by an $``$. In this case the hermitian conjugation is given by $`\left(fg\right)^{}=g^{}f^{}`$. Under the integral, for closed manifolds, this product has the cyclic property $`\mathrm{Tr}f_1f_2\mathrm{}f_n=\mathrm{Tr}f_nf_1f_2\mathrm{}f_{n1}`$. In particular $`\mathrm{Tr}f_1f_2=\mathrm{Tr}f_1f_2.`$ Therefore, a theory on the noncommutative space of the $`\widehat{x}`$, is equivalent to a theory of usual fields, where the function product is substituted by the Moyal $``$-product.
This suggests that any theory can be converted into a noncommutative one by replacing the ordinary function product with the $``$-product.
### III.2 The Seiberg-Witten Map
In order to build noncommutative Yang-Mills theories it is necessary, first of all, that a commutative limit exists and that a perturbative study of the noncommutative theory is possible. In this case the solutions of such a theory must depend on the noncommutativity parameter $`\theta `$ in the form of a power series expansion.
For an ordinary Yang-Mills theory, the gauge field and the strength field tensor transformations can be written as:
$`\delta _\lambda A_\mu `$ $`=_\mu \lambda +i\lambda A_\mu iA_\mu \lambda ,`$ (6)
$`F_{\mu \nu }`$ $`=_\mu A_\nu _\nu A_\mu iA_\mu A_\nu +iA_\nu A_\mu ,`$
$`\delta _\lambda F_{\mu \nu }`$ $`=i\lambda F_{\mu \nu }iF_{\mu \nu }\lambda .`$
For the noncommutative gauge theory, we use the same equations in the gauge field and strength field tensor transformations, except that the matrix multiplications are replaced by the $``$ product. Then the gauge field and the strength field tensor transformations are SW :
$`\widehat{\delta }_{\widehat{\lambda }}\widehat{A}_\mu `$ $`=_\mu \widehat{\lambda }+i\widehat{\lambda }\widehat{A}_\mu i\widehat{A}_\mu \widehat{\lambda },`$
$`\widehat{F}_{\mu \nu }`$ $`=_\mu \widehat{A}_\nu _\nu \widehat{A}_\mu i\widehat{A}_\mu \widehat{A}_\nu +i\widehat{A}_\nu \widehat{A}_\mu ,`$ (7)
$`\widehat{\delta }_{\widehat{\lambda }}\widehat{F}_{\mu \nu }`$ $`=i\widehat{\lambda }\widehat{F}_{\mu \nu }i\widehat{F}_{\mu \nu }\widehat{\lambda },`$
from which the original Yang-Mills theory (6) results in the limit $`\theta 0`$. Notice that equations (7) are valid even for abelian gauge fields. Due to the form of the Moyal product (5), the noncommutative theory has the structure of a nonlocal theory. However, if we consider it as an effective theory, its energy scale gives us a cutoff and nonlocality is not a problem. Further, as shown by Kontsevich Kontsevich , at the level of the physical degrees of freedom there is a one to one relation between the commutative and the noncommutative theories. However both theories are quite different, as noncommutativity generates new couplings.
Let us consider the noncommutative gauge transformations of an abelian theory,
$$\delta _{\widehat{\lambda }}\widehat{A}_\mu =_\mu \widehat{\lambda }+i\widehat{\lambda }\widehat{A}_\mu i\widehat{A}_\mu \widehat{\lambda },$$
(8)
we see that they look like nonabelian ones, although they continue to depend on only one generator. For nonabelian groups, things are more complicated W2 ,
$`\delta _{\widehat{\lambda }}\widehat{A}_\mu `$ $`=_\mu \widehat{\lambda }+i\widehat{\lambda }\widehat{A}_\mu i\widehat{A}_\mu \widehat{\lambda }`$
$`=_\mu \widehat{\lambda }^a\mathrm{\Lambda }_a+i\widehat{\lambda }^a\mathrm{\Lambda }_a\widehat{A}_\mu ^b\mathrm{\Lambda }_bi\widehat{A}_\mu ^b\mathrm{\Lambda }_b\widehat{\lambda }^a\mathrm{\Lambda }_a`$
$`=_\mu \widehat{\lambda }^a\mathrm{\Lambda }_a+{\displaystyle \frac{i}{2}}\left\{\widehat{\lambda }^a\stackrel{}{,}\widehat{A}_\mu ^b\right\}[\mathrm{\Lambda }_a,\mathrm{\Lambda }_b]+{\displaystyle \frac{i}{2}}\left[\widehat{\lambda }^a\stackrel{}{,}\widehat{A}_\mu ^b\right]\{\mathrm{\Lambda }_a,\mathrm{\Lambda }_b\}.`$ (9)
Now the transformation algebra is generated by commutators and anticommutators, which amounts to the universal enveloping algebra of the original algebra $`U(g,R)`$, where $`R`$ is the corresponding representation. The generators of this algebra satisfy,
$$[\mathrm{\Lambda }_A,\mathrm{\Lambda }_B]=if_{ABC}\mathrm{\Lambda }_C,\{\mathrm{\Lambda }_A,\mathrm{\Lambda }_B\}=d_{ABC}\mathrm{\Lambda }_C,$$
(10)
where $`f_{ABC}=f_{BAC}`$ and $`d_{ABC}=d_{BAC}`$ are the structure constants.
These transformations are satisfied order by order on $`\theta `$, and all coefficients of the higher terms can be used to fix the gauge degrees of freedom of $`\widehat{A}_\mu `$. In such a gauge fixing, the only remaining freedom of the transformation parameters $`\widehat{\lambda }`$ are the ones of the commutative theory, so they should depend only on $`\lambda ^a`$ and their derivatives. In this case consistency implies that an infinitesimal commutative gauge transformation $`\delta _\lambda A_\mu =_\mu \lambda +i\lambda A_\mu iA_\mu \lambda `$, will induce the noncommutative one,
$$\widehat{A}_\mu (A+\delta _\lambda A)=\widehat{A}_\mu (A)+\widehat{\delta }_{\widehat{\lambda }}\widehat{A}_\mu (A).$$
(11)
This is the so called Seiberg-Witten map.
The solution to (11) can be obtained by setting $`\widehat{A}_\mu =A_\mu +A_\mu ^{}(A)`$ and $`\widehat{\lambda }=\lambda +\lambda ^{}(\lambda ,A)`$, where $`A_\mu ^{}`$ and $`\lambda ^{}`$ are local functions of $`\lambda `$ and $`A_\mu `$ of first order in $`\theta `$. Then substituting in (11) and expanding to first order,
$$A_\mu ^{}(A+\delta _\lambda A)A_\mu ^{}(A)_\mu \lambda ^{}i[\lambda ^{},A_\mu ]i[\lambda ,A_\mu ^{}]=\frac{1}{2}\theta ^{\alpha \beta }(_\alpha \lambda _\beta A_\mu +_\beta A_\mu _\alpha \lambda ).$$
(12)
One solution of this equation is given by SW ,
$`\widehat{A}_\mu (A)`$ $`=A_\mu +A_\mu ^{}(A)=A_\mu {\displaystyle \frac{1}{4}}\theta ^{\alpha \beta }\{A_\alpha ,_\beta A_\mu +F_{\beta \mu }\}+𝒪(\theta ^2),`$ (13)
$`\widehat{\lambda }(\lambda ,A)`$ $`=\lambda +\lambda ^{}(\lambda ,A)=\lambda +{\displaystyle \frac{1}{4}}\theta ^{\alpha \beta }\{_\alpha \lambda ,A_\beta \}+𝒪(\theta ^2),`$ (14)
from which it turns out that,
$$\widehat{F}_{\mu \nu }=F_{\mu \nu }+\frac{1}{4}\theta ^{\alpha \beta }\left(2\{F_{\mu \alpha },F_{\nu \beta }\}\{A_\alpha ,\left(D_\beta +_\beta \right)F_{\mu \nu }\}\right)+𝒪(\theta ^2).$$
(15)
These equations (13, 14, 15) are the explicit form of the Seiberg-Witten map, which in this way can be constructed for any Lie algebra of transformations W2 .
As shown the noncommutative generators $`\widehat{\lambda }`$ take values in the enveloping algebra. In the case of the fundamental representation of unitary groups $`U(N)`$, they coincide with their enveloping algebras. For the algebra of $`SU(N)`$ in the fundamental representation, the enveloping algebra incorporates, through the anticommutators of the generators, the identity matrix $`\mathrm{\Lambda }_0=\frac{1}{\sqrt{2N}}𝕀_{N\times N}`$, and is then given by $`U(N)`$.
## IV The Noncommutative Model
As previously mentioned, our purpose is the construction of a noncommutative version of the 6-dimensional $`SU(3)`$ gauge theory presented in Section 2. The fact that we are considering noncommutativity in 6D, means that we only need the Seiberg-Witten map for gauge fields. Thus the effects of noncommutativity on the Higgs field and its interactions will arise after dimensional reduction, in particular from the Seiberg-Witten map of the gauge fields $`A_z`$ and $`A_{\overline{z}}`$.
The noncommutative action is given by:
$$\widehat{𝒮}_{NC}=\frac{1}{2}\mathrm{Tr}d^6x\widehat{F}_{mn}\widehat{F}^{mn},$$
(16)
where
$$\widehat{F}_{mn}=F_{mn}+\frac{1}{4}\theta ^{kl}\left(2\{F_{mk},F_{nl}\}\{A_k,\left(D_l+_l\right)F_{mn}\}\right)+𝒪(\theta ^2).$$
(17)
Here the indexes $`m,n,k`$ and $`l`$ take the values $`0,\mathrm{},3,z`$ and $`\overline{z}`$. Thus the noncommutative parameter $`\theta ^{kl}`$ can be: $`\theta ^{\mu \nu }`$ (noncommutativity among the 4-dimensional space-time coordinates), $`\theta ^{\mu z},`$ $`\theta ^{\mu \overline{z}}`$ (noncommutativity among the 4-dimensional space-time coordinates and the extra dimensions coordinates) and $`\theta ^{z\overline{z}}`$ (noncommutativity between the extra dimensions). Therefore, after inserting the noncommutative field strength (17) into (16), the noncommutative action gets the following first order corrections,
$`{\displaystyle \frac{\theta ^{\alpha \beta }}{4}}\mathrm{Tr}\{[2\{F_{\mu \alpha },F_{\nu \beta }\}\{A_\alpha ,(D_\beta +_\beta )F_{\mu \nu }\}]F^{\mu \nu }`$
$`+2\left[2\{F_{\mu \alpha },F_{z\beta }\}\{A_\alpha ,\left(D_\beta +_\beta \right)F_{\mu z}\}\right]F^{\mu z}`$
$`+2\left[2\{F_{\mu \alpha },F_{\overline{z}\beta }\}\{A_\alpha ,\left(D_\beta +_\beta \right)F_{\mu \overline{z}}\}\right]F^{\mu \overline{z}}`$
$`+2[2\{F_{z\alpha },F_{\overline{z}\beta }\}\{A_\alpha ,(D_\beta +_\beta )F_{z\overline{z}}\}]F^{z\overline{z}}\}`$
$`{\displaystyle \frac{\theta ^{\alpha i}}{4}}\mathrm{Tr}\{[4\{F_{\mu \alpha },F_{\nu i}\}\{A_\alpha ,(D_i+_i)F_{\mu \nu }\}+\{A_i,(D_\alpha +_\alpha )F_{\mu \nu }\}]F^{\mu \nu }`$
$`+2\left[2\{F_{\mu \alpha },F_{ji}\}2\{F_{j\alpha },F_{\mu i}\}\{A_\alpha ,\left(D_i+_i\right)F_{\mu j}\}+\{A_i,\left(D_\alpha +_\alpha \right)F_{\mu j}\}\right]F^{\mu j}`$
$`+2[2\{F_{i\alpha },F_{\overline{z}z}\}\{A_\alpha ,(D_i+_i)F_{z\overline{z}}\}+\{A_i,(D_\alpha +_\alpha )F_{z\overline{z}}\}]F^{z\overline{z}}\}`$
$`{\displaystyle \frac{\theta ^{z\overline{z}}}{4}}\mathrm{Tr}\{[4\{F_{\mu z},F_{\nu \overline{z}}\}\{A_z,D_{\overline{z}}F_{\mu \nu }\}+\{A_{\overline{z}},D_zF_{\mu \nu }\}]F^{\mu \nu }`$
$`+2\left[2\{F_{\mu z},F_{z\overline{z}}\}\{A_z,D_{\overline{z}}F_{\mu z}\}+\{A_{\overline{z}},D_zF_{\mu z}\}\right]F^{\mu z}`$
$`+2\left[2\{F_{z\overline{z}},F_{\mu \overline{z}}\}\{A_z,D_{\overline{z}}F_{\mu \overline{z}}\}+\{A_{\overline{z}},D_zF_{\mu \overline{z}}\}\right]F^{\mu \overline{z}}`$
$`+2[2\{F_{z\overline{z}},F_{z\overline{z}}\}\{A_z,D_{\overline{z}}F_{z\overline{z}}\}+\{A_{\overline{z}},D_zF_{z\overline{z}}\}]F^{z\overline{z}}\},`$ (18)
where $`\mu ,\nu ,\alpha ,\beta =0,\mathrm{},3`$, $`i=z,\overline{z}`$.
After somewhat cumbersome computations, we obtain the following expression for these corrections in terms of the $`SU(2)`$ and $`U(1)`$ field strengths $`W^{\mu \nu }`$ and $`B^{\mu \nu }`$ respectively, the corresponding gauge fields $`W^\mu `$ and $`B^\mu `$ and the Higgs field $`\varphi `$,
$`\widehat{}_{NC}`$ $`={\displaystyle \frac{1}{2}}\mathrm{Tr}W_{\mu \nu }W^{\mu \nu }{\displaystyle \frac{1}{4}}B_{\mu \nu }B^{\mu \nu }+\left(D_\mu \varphi \right)^{}(D^\mu \varphi ){\displaystyle \frac{g^2}{2}}\left|\varphi \right|^4`$
$`{\displaystyle \frac{1}{4}}\theta ^{\alpha \beta }\{{\displaystyle \frac{1}{2\sqrt{3}}}\mathrm{Tr}[4\{W_{\mu \alpha },B_{\nu \beta }\}W^{\mu \nu }+2\{W_{\mu \alpha },W_{\nu \beta }\}B^{\mu \nu }I\{W_\alpha ,D_\beta W_{\mu \nu }\}B^{\mu \nu }I`$
$`\{B_\alpha ,D_\beta W_{\mu \nu }\}W^{\mu \nu }]+{\displaystyle \frac{1}{2\sqrt{3}}}B_\alpha _\beta B_{\mu \nu }B^{\mu \nu }{\displaystyle \frac{1}{2\sqrt{3}}}B_{\mu \alpha }B_{\nu \beta }B^{\mu \nu }`$
$`+2(D^\mu \varphi )^{}\left(W_{\mu \alpha }{\displaystyle \frac{1}{2\sqrt{3}}}B_{\mu \alpha }I\right)(D_\beta \varphi )+h.c.`$
$`+(D^\mu \varphi )^{}\left(W_\alpha {\displaystyle \frac{1}{2\sqrt{3}}}B_\alpha I\right)\left(\stackrel{}{}_\beta +\stackrel{}{D}_\beta \right)(D_\mu \varphi )`$
$`+(D^\mu \varphi )^{}\left(\stackrel{}{}_\beta +\stackrel{}{D}_\beta \right)\left(W_\alpha {\displaystyle \frac{1}{2\sqrt{3}}}B_\alpha I\right)(D_\mu \varphi )`$
$`+ig\left[\varphi ^{}(D_\alpha \varphi )(D_\beta \varphi )^{}\varphi (D_\beta \varphi )^{}(D_\alpha \varphi )\varphi ^{}\varphi \right]ig^3\varphi ^{}\varphi \varphi ^{}W_\beta W_\alpha \varphi `$
$`g^2[\varphi ^{}(W_\alpha +{\displaystyle \frac{1}{2\sqrt{3}}}B_\alpha I)_\beta (\varphi \varphi ^{})\varphi {\displaystyle \frac{2}{\sqrt{3}}}B_\alpha _\beta (\varphi \varphi ^{})\varphi ^{}\varphi `$
$`+\varphi ^{}_\beta (\varphi \varphi ^{})(W_\alpha +{\displaystyle \frac{1}{2\sqrt{3}}}B_\alpha I)\varphi ]\}`$
$`+{\displaystyle \frac{i}{2}}\theta ^{z\overline{z}}\{2i(D_\mu \varphi )^{}(W^{\mu \nu }+{\displaystyle \frac{1}{\sqrt{3}}}B^{\mu \nu }I)(D_\nu \varphi )`$
$`+{\displaystyle \frac{g}{2}}\left[\varphi ^{}\varphi (D_\mu \varphi )^{}(D^\mu \varphi )(D_\mu \varphi )^{}\varphi \varphi ^{}(D^\mu \varphi )\right]`$
$`g\varphi ^{}(W_{\mu \nu }W^{\mu \nu }+{\displaystyle \frac{1}{\sqrt{3}}}W_{\mu \nu }B^{\mu \nu }{\displaystyle \frac{1}{4}}B_{\mu \nu }B^{\mu \nu })\varphi \},`$ (19)
In this equation there are new interactions with respect to the ones found in the $`4D`$ noncommutative formulations of the SM W3 ; Chaichian1 . For instance the interactions between the weak gauge fields and the electromagnetic field which appear in the first terms that multiply the four dimensional noncommutativity parameter $`\theta ^{\alpha \beta }`$. Note that there are not corrections linear in the noncommutativity parameter $`\theta ^{i\alpha }`$, they turn out to be identically zero, as a consequence of the orbifold symmetries. Of particular interest are the corrections corresponding to noncommutativity between the extra dimensions, i.e. the terms multiplied by $`\theta ^{z\overline{z}}`$, given by interactions among the Higgs and the gauge bosons, and also higher order Higgs self-interactions. Considering only these sort of corrections, we have,
$`\widehat{}_{NC}=`$ $`{\displaystyle \frac{1}{2}}\mathrm{Tr}W_{\mu \nu }W^{\mu \nu }{\displaystyle \frac{1}{4}}B_{\mu \nu }B^{\mu \nu }+\left(D_\mu \varphi \right)^{}(D^\mu \varphi ){\displaystyle \frac{g^2}{2}}\left|\varphi \right|^4`$
$`+{\displaystyle \frac{i}{2}}\theta ^{z\overline{z}}\{2i\left(D_\mu \varphi \right)^{}(W^{\mu \nu }+{\displaystyle \frac{1}{\sqrt{3}}}B^{\mu \nu }I)(D_\nu \varphi )`$
$`+{\displaystyle \frac{g}{2}}\left[\varphi ^{}\varphi \left(D_\mu \varphi \right)^{}(D^\mu \varphi )\left(D_\mu \varphi \right)^{}\varphi \varphi ^{}(D^\mu \varphi )\right]`$
$`g\varphi ^{}(W_{\mu \nu }W^{\mu \nu }+{\displaystyle \frac{1}{\sqrt{3}}}W_{\mu \nu }B^{\mu \nu }{\displaystyle \frac{1}{4}}B_{\mu \nu }B^{\mu \nu })\varphi \}.`$ (20)
The noncommutative corrections in this Lagrangian are dimension-six operators, well known from the electroweak effective Lagrangian technique EL , a scheme in which the effects of these terms can be studied in a model-independent manner. It is interesting to consider the extension of the $`SU(3)`$ six dimensional gauge group by an $`U(1)`$ factor. Its gauge field is invariant under the orbifold symmetries scruca and as it does not mix with the $`SU(3)`$ gauge field, the noncommutative corrections (17) and in (20) will not be affected.
## V Phenomenological implications
In this section we will analyze the phenomenological implications of noncommutativity between extra dimensions, in the case of Lagrangian (20). Thus, we fix our attention on the terms multiplied by $`\theta ^{45}`$ $`\left(\theta ^{45}=i\theta ^{z\overline{z}}\right)`$. These terms contain new interactions relative to the SM and its noncommutative versions W3 ; Chaichian1 . As mentioned, the model we are considering predicts a too large weak angle. However we can expect that the kind of noncommutative terms considered here will also arise from a more realistic theory, that would reproduce the SM in the limit of vanishing $`\theta `$ parameters. For instance, as mentioned end of the last section, in the case of an $`U(1)`$ extension, the noncommutative corrections which affect the electroweak gauge fields are not modified. Accordingly, in the following we will consider these terms as representing deviations of the genuine standard electroweak Lagrangian. In order to analyze their effects, we observe that these new interactions are given by well–known dimension–six operators, which have been already studied in, by example, the electroweak effective Lagrangian approach EL , a scheme appropriate to investigate in a model–independent manner physics lying beyond the Fermi scale. For this purpose, it is convenient to divide the above operators into three sets, as follows
$`𝒪_{\varphi W}={\displaystyle \frac{g}{2}}\theta ^{45}(\varphi ^{}W_{\mu \nu }W^{\mu \nu }\varphi ),`$ (21)
$`𝒪_{\varphi B}={\displaystyle \frac{g}{8}}\theta ^{45}(\varphi ^{}B_{\mu \nu }B^{\mu \nu }\varphi ),`$ (22)
$`𝒪_{WB}={\displaystyle \frac{g}{2\sqrt{3}}}\theta ^{45}(\varphi ^{}W_{\mu \nu }B^{\mu \nu }\varphi ),`$ (23)
$`𝒪_\varphi ^{(1)}=g\theta ^{45}(\varphi ^{}\varphi )(D_\mu \varphi )^{}(D^\mu \varphi ),`$ (24)
$`𝒪_\varphi ^{(3)}=g\theta ^{45}[(D_\mu \varphi )^{}\varphi ][\varphi ^{}(D^\mu \varphi )],`$ (25)
$`𝒪_{DW}=i\theta ^{45}(D_\mu \varphi )^{}W^{\mu \nu }(D_\nu \varphi ),`$ (26)
$`𝒪_{DB}={\displaystyle \frac{i}{\sqrt{3}}}\theta ^{45}(D_\mu \varphi )^{}B^{\mu \nu }(D_\nu \varphi ).`$ (27)
First we observe that there are potential modifications induced by these operators on the quadratic SM Lagrangian, as they can alter some tree level relations which are experimentally constrained, such as the kinetic energy part of the $`W`$ and $`Z`$ bosons. In particular, they can give tree–level contributions to the $`S`$ and $`T`$ oblique parametersPT . Let us focus our attention on those interactions which affect the quadratic part of the SM gauge sector.
After spontaneous symmetry breaking, all the above operators induce new nonrenormalizable interactions, as well as renormalizable ones, which modify those predicted by the dimension–four theory. In particular, the first two sets of operators induce bilinear terms that can eventually modify the SM parameters EL ; Wudka2 . On the other hand, although the last set of operators are potentially interesting from the phenomenological point of view TH , they are not important for our purposes, as they do not introduce modifications in the quadratic Lagrangian. Concerning the first set, it is easy to see that $`𝒪_{\varphi W}`$ and $`𝒪_{\varphi B}`$ modify the canonical form of the kinetic terms $`W_{\mu \nu }W^{\mu \nu }`$ and $`B_{\mu \nu }B^{\mu \nu }`$, respectively. However, these effects are unobservable indeed, since they can be absorbed in a finite renormalization of the gauge fields and the coupling constant $`g`$. As to the $`𝒪_{WB}`$, $`𝒪_\varphi ^{(1)}`$, and $`𝒪_\varphi ^{(3)}`$ operators, they introduce nontrivial modifications in the quadratic Lagrangian. In particular, as we will see below, the first and the last of these operators are sensitive to the low–energy data, as they contribute to the $`S`$ and $`T`$ parameters at the tree level. Up to some surface terms, the quadratic part of the effective Lagrangian, i.e., the SM and new contributions, can conveniently be written as
$`_{Kinetic}`$ $`=`$ $`{\displaystyle \frac{1}{2}}W^{a\mu }\left\{\left[\mathrm{}+{\displaystyle \frac{g^2v^2}{4}}\left(1{\displaystyle \frac{\alpha _\varphi ^{(1)}}{2}}\right)\right]g_{\mu \nu }_\mu _\nu \right\}W^{a\nu }`$ (28)
$`+{\displaystyle \frac{1}{2}}B^\mu \left\{\left[\mathrm{}+{\displaystyle \frac{g^2v^2}{4}}\left(1{\displaystyle \frac{\alpha _\varphi ^{(1)}}{2}}\right)\right]g_{\mu \nu }_\mu _\nu \right\}B^\nu `$
$`+W^{3\mu }\left({\displaystyle \frac{g^2v^2}{16}}\alpha _\varphi ^{(3)}\right)g_{\mu \nu }W^{3\nu }+W^{3\mu }\left[\alpha _{WB}\left(\mathrm{}g_{\mu \nu }_\mu _\nu \right)\right]B^\nu ,`$
where the unobservable effects arising from the $`𝒪_{\varphi W}`$ and $`𝒪_{\varphi B}`$ operators were ignored. In addition, in order to identify the origin of each contribution, we have introduced the definitions: $`\alpha _{WB}=gv^2\theta ^{45}/2\sqrt{3}`$ and $`\alpha _\varphi ^{(1)}=\alpha _\varphi ^{(3)}=gv^2\theta ^{45}`$, with $`v`$ the Fermi scale. The new ingredients in this expression with respect to the standard result, is the mixing between the field strengths $`W_{\mu \nu }^3`$ and $`B_{\mu \nu }`$ induced by the $`𝒪_{WB}`$ operator, as well as the presence of a quadratic term in $`W_\mu ^3`$ generated by the $`𝒪_\varphi ^{(3)}`$ operator. As we will see below, the $`W_{\mu \nu }^3B^{\mu \nu }`$ mixing given by the $`𝒪_{WB}`$ operator contribute to the $`S`$ parameter at the tree level, as it involves derivatives. Also, it is important to notice that while $`𝒪_\varphi ^{(1)}`$ affects with the same intensity both the $`W`$ and $`Z`$ masses (see first and second terms in (28)), the $`𝒪_\varphi ^{(3)}`$ operator modifies only the $`Z`$–mass, as it is evident from the term proportional to $`W_\mu ^3W_\nu ^3`$. This asymmetric contribution to the gauge field masses is the responsible for deviations from the SM value of the $`\rho `$–parameter $`\rho =\alpha T`$EL ; DHT . This contributions is associated with a violation of the custodial $`SU(2)`$ symmetry Custodial , which as it is well known, guarantees the tree level value $`\rho =1`$ in the SM. The diagonalization of the resultant kinetic energy sector and its impact on the SM parameters have been studied in the literature in the more general context of electroweak effective Lagrangians EL ; Wudka2 . The tree level contribution of the $`𝒪_{WB}`$ and $`𝒪_\varphi ^{(3)}`$ operators to the $`S`$ and $`T`$ parameters has already been studied in this more general context by Hagiwara et al. H . For our purposes, it is convenient to follow the approach introduced by these authors. The oblique parameters characterize the influence of physics beyond the Fermi scale. They are given as linear combinations of the transverse components of the gauge–boson vacuum polarizations
$$\mathrm{\Pi }_{ij}^{\mu \nu }(p)=\mathrm{\Pi }_{ij}(p^2)g^{\mu \nu }+(p^\mu p^\nu \mathrm{terms}),$$
(29)
where $`ij`$ stands for $`aa`$, $`YY`$, and $`3Y`$, where $`(a,Y)`$ are $`SU(2)\times U_Y(1)`$ indices. In particular, the $`S`$ and $`T`$ parameters are defined byDKS
$`\alpha S`$ $`=`$ $`{\displaystyle \frac{2e^2}{m_Z^2}}[\mathrm{\Pi }_{3Y}(0)\mathrm{\Pi }_{3Y}(m_Z^2)],`$ (30)
$`\alpha T`$ $`=`$ $`{\displaystyle \frac{2e^2}{m_W^2}}\mathrm{Re}[\mathrm{\Pi }_{11}(0)\mathrm{\Pi }_{33}(0)].`$ (31)
It is not difficult to see from (28), that $`𝒪_{WB}`$ contributes to $`S`$ but not to $`T`$, whereas $`𝒪_\varphi ^{(3)}`$ contributes to $`T`$ but not to $`S`$. It is also evident from (28) that $`𝒪_\varphi ^{(1)}`$ does not contribute to these parameters. Considering only the new physics contribution, one finds,
$`S_{NP}`$ $`=`$ $`2c_W\sqrt{{\displaystyle \frac{\pi }{3\alpha }}}(v^2\theta ^{45}),`$ (32)
$`T_{NP}`$ $`=`$ $`{\displaystyle \frac{1}{s_W}}\sqrt{{\displaystyle \frac{\pi }{\alpha }}}(v^2\theta ^{45}),`$ (33)
where $`s_W`$ and $`c_W`$ stand for $`sin\theta _W`$ and $`cos\theta _W`$, respectively. On the other hand, the current experimental data PDG give the next values of $`S`$ and $`T`$ which can be induced by new physics effects
$`S_{NP}^{Exp}=0.13\pm 0.10(0.08),`$ (34)
$`T_{NP}^{Exp}=0.17\pm 0.12(+0.09),`$ (35)
where the central values assume $`m_H=117`$ GeV. The change for $`m_H=300`$ GeV is shown in parenthesis. Assuming the first value for the Higgs mass and using the values for the SM parameters reported in PDG , one finds at $`95\%`$ C.L.
$`5.63\times 10^8GeV^2`$ $`<\theta ^{45}<`$ $`1.91\times 10^7GeV^2,`$ (36)
$`3.18\times 10^7GeV^2`$ $`<\theta ^{45}<`$ $`1.06\times 10^7GeV^2,`$ (37)
which leads to the following bound for the new physics scale
$$5.63\times 10^8GeV^2<\theta ^{45}<1.06\times 10^7GeV^2.$$
(38)
To conclude this part, it is worth comparing this bound with those obtained in the literature for the noncommutativity scale of four dimensional theories. To this respect, bounds of order of one TeV have been obtained from collider physics BCP . On the other hand, more stringent bounds of order $`(10TeV)^2`$ Carroll or higher Carlson have been derived from low–energy tests of Lorentz violation. However, as it has been recently argued Calmet , that these bounds are extremely model dependent and should be taken with some care.
We are now in position to discuss some phenomenological implications of these new interactions. Motivated by the fact that new physics effects would be more evident in those processes which are forbidden or strongly suppressed in the SM, we will consider the rare Higgs boson decay into two photons, which is an one–loop prediction of the model and thus is naturally suppressed HSM . Due to its phenomenological importance, this decay has been the subject of permanent interest in the literature. Apart from providing a good signature for the Higgs boson search at hadron colliders with mass in the intermediate range $`120\mathrm{Gev}<m_H<2m_Z`$ DHC , the decay width of this process is also of great interest because it determines the cross section for Higgs production in $`\gamma \gamma `$ collisions PPP . Due to the fact that the $`H\gamma \gamma `$ coupling is generated by loop effects of charged particles, its sensitivity to new heavy charged particles has been studied in many well motivated extensions of the SM, as the two Higgs doublet model (THDM) HHG , the minimal supersymmetric standard model (MSSM) SUSY , the left–right symmetric models (LRM) LR , and the Littlest Higgs model (LHM) Han . Many of its properties have also been studied in a model–independent manner using the effective Lagrangian framework TH ; HEL . In our model, the $`H\gamma \gamma `$ vertex (as well as $`H\gamma Z`$ one) is induced at the tree level by the set of operators given in eqs.(21-23). The corresponding Lagrangian can be written as follows:
$$_{H\gamma \gamma }=\frac{\alpha _W}{4}m_W\theta ^{45}HF_{\mu \nu }F^{\mu \nu },$$
(39)
where $`\alpha _W=4s_W^2c_W^22s_{2W}/\sqrt{3}`$. The total decay width $`\mathrm{\Gamma }_{NC}`$ can be conveniently written in terms of the SM width $`\mathrm{\Gamma }_{SM}`$, as follows:
$$\mathrm{\Gamma }_{NC}(H\gamma \gamma )=\mathrm{\Gamma }_{SM}(H\gamma \gamma )\left|1+\frac{𝒜_{NC}}{𝒜_{SM}}\right|^2,$$
(40)
where $`𝒜_{NC}=\alpha _Wm_W^2\theta ^{45}`$, whereas $`𝒜_{SM}`$ represents the charged fermion and $`W`$ boson loop contributions, which is given by
$$𝒜_{SM}=\frac{\alpha ^{3/2}}{\sqrt{4\pi }s_W}\left[\underset{f}{}N_{Cf}Q_f^2F_f+F_W\right].$$
(41)
In this expression, $`f`$ stands for quarks or leptons, $`N_{Cf}`$ is the color index, and $`Q_f`$ is the electric charge of the fermion in units of the charge of the positron. In addition, $`F_f`$ and $`F_W`$ are the fermion and $`W`$ boson loop amplitudes, respectively, which can be found in ref. HHG .
Though the bound for $`\theta ^{45}`$ was estimated for a value $`m_H=117`$ GeV, for illustration purposes we will present results that contemplate larger values of the Higgs mass. In Fig.1, the variation of the normalized decay width $`R=\mathrm{\Gamma }_{NC}/\mathrm{\Gamma }_{SM}`$ is displayed as function of the Higgs mass. From this figure, it can be appreciated that $`R`$ is sensitive to the sign of the $`\theta ^{45}`$ parameter. A constructives or destructives effect corresponds to $`\theta ^{45}<0(\theta ^{45}>0)`$, which can increase or decrease the standard model prediction $`\mathrm{\Gamma }_{SM}`$ up to by $`27\%`$ and $`35\%`$, respectively, for $`m_H`$ in the range $`120200`$ GeV.
It is interesting to compare this result with those obtained from other models. In general, theories beyond the SM require more complicated Higgs sectors, i.e. they incorporate new neutral and charged Higgs bosons. However, in most cases, it is always possible to identify in an appropriate limit a SM–like Higgs boson, that is a CP–even neutral scalar whose couplings to pairs of $`W`$ and $`Z`$ bosons coincide with those given in the minimal SM. Furthermore, new contributions coming from new charged scalar, fermion, and vector particles are expected. We briefly review the results for a SM–like Higgs boson decaying into two photons, for the models mentioned previously. In the THDM, the only new charged particle is the $`H^\pm `$ Higgs boson, but its loop contribution to this decay is very small compared with the dominant $`W`$ contribution. Thus, in this model, the $`\gamma \gamma `$ decay width is essentially the same as for the SM HHG . In the case of the MSSM, this decay gets new contributions from superpartner loops. The $`\gamma \gamma `$ width tends to be lower than the one of the SM due to cancelations between the $`W`$ loop and the supersymmetric chargino loops. In this case, if the charginos are heavy, the decay width can be quite near to the SM width SUSY . Further, it was found in ref. LR that the contribution of a new $`W`$ boson, like the one predicted by LRM, is quite suppressed, because the corresponding loop amplitude is related to the SM $`W`$ boson amplitude as $`F_{W_R}=(m_{W_L}/m_{W_R})^2F_{W_L}`$. Taking into account the existing $`W_R`$–mass bounds PDG , the $`\gamma \gamma `$ width can be enhanced up to $`5\%`$ at best. As to the prediction of the Littlest Higgs model is concerned, in Han it was found that the $`\gamma \gamma `$ width is reduced by $`57\%`$ compared to the SM value. From these results, we can see that the impact of noncommutative extra dimensions on the $`H\gamma \gamma `$ decay may be significantly more important than that predicted by some of the most popular SM extensions.
## VI Conclusions
In this work we explore the consequences of noncommutativity in a 6-dimensional model, by means of the Seiberg-Witten map. We consider the $`SU(3)`$ gauge Higgs unification model of the electroweak interactions of antoniadis1 and Wulzer , compactified to 4D on an orbifold $`T^2/Z_N`$ for $`N=3,4,6`$. We analyze noncommutativity among all the 6-dimensional coordinates. As a consequence of the orbifold symmetries, it turns out that there are no corrections to the model due to noncommutativity among the 4D coordinates and the two-extra dimensions. We find that the corrections we obtain corresponding to noncommutativity among the 4D coordinates, differ from the ones of noncommutative models calculated directly in 4D, also by means of the Seiberg-Witten map W3 . On the other side, the corrections corresponding to noncommutativity of the extra dimensions have interesting phenomenological consequences, as we emphasize below.
As well as in the commutative model, the spontaneous symmetry breaking should arise dynamically, from first order quantum corrections. This step can be done in the noncommutative theory, as far as the expected Higgs mass is much less than the noncommutativity scale. Thus it would be interesting to include matter and to study the corresponding noncommutative corrections, which could be done following W2 , progress in this direction will be reported elsewhere.
As mentioned in the introduction, the model we are considering here has a too high value for the weak angle. However, as noted in Wulzer , there are various ways to solve this problem, in particular by an extension by a $`U(1)`$ factor. Furthermore, the noncommutative Seiberg-Witten map of the corresponding gauge field will not mix with the already present noncommutative corrections. Thus we can expect that in a noncommutative version of this extended model, the kind of corrections presented here will still be present, in particular those corresponding to noncommutativity between extra dimensions. Under this working assumption, we have studied the corrections due to noncommutativity of the extra dimensions by means of the effective lagrangian techniques. First, by the observation that these terms are quite sensitive to the $`S`$ and $`T`$ oblique parameters, we could obtain the bound to $`\theta ^{45}`$ given by (38). In four dimensional noncommutative models, there are bounds obtained e.g. from low–energy tests of Lorentz violation, which are extremely model dependent Calmet . We think that, in the framework of our working assumption, our bound has a less speculative nature, as it was obtained directly from the experimental constraints on the oblique parameters, without additional assumptions. With this bound established, we have looked at the impact of our corrections on the rare Higgs decay into two photons. It turns out that the effect depends on the signature of the noncommutativity parameter, increasing or decreasing respectively the value of the SM decay width $`\mathrm{\Gamma }_{SM}`$, with a net effect which could be significantly more important than that predicted by some of the most popular SM extensions PDG ; SUSY ; LR ; Han .
Finally, from the results of the particular noncommutative model we started with, which could be interesting on its own, we can conclude that noncommutativity in higher dimensional models can have interesting consequences and phenomenological effects beyond those of four dimensional noncommutative theories. The study of more realistic models, including matter fields, is in progress.
###### Acknowledgements.
We thank H. García-Compeán for discussions. This work has been supported by CONACYT grant 47641 and Projects by PROMEP, UG and VIEP-BUAP 13/I/EXC/05. |
warning/0506/math0506531.html | ar5iv | text | # Ultrametric Logarithm Laws I.
## 1. Introduction
The notion of shrinking targets for dynamical systems is a much studied and extremely useful concept (, and the references therein) especially for geometric and number theoretic applications. Shrinking target properties for group actions on homogeneous spaces were studied in an important paper of D. Kleinbock and G. A. Margulis. We first describe their results. Let $`G`$ denote a connected semisimple Lie group without compact factors, and $`\mathrm{\Gamma }`$ denote a non-uniform lattice in $`G`$. Let $`K`$ be a maximal compact subgroup of $`G`$, and let $`Y=K\backslash G/\mathrm{\Gamma }`$ denote the associated non-compact irreducible locally symmetric space of finite volume. In , D. Kleinbock and G. A. Margulis studied the phenomenon of shrinking target properties for the geodesic flow, generalizing an earlier work of D. Sullivan and established the following theorem, commonly called a logarithm law.
For $`xY`$, let $`T_x^1(Y)`$ denote the unit tangent space at $`x`$. Let $`\nu `$ denote the Haar measure on $`T_x^1(Y)`$. For $`\theta T_x^1(Y)`$, and $`t`$ let $`g_t(x,\theta )`$ denote the image of $`(x,\theta )`$ under geodesic flow for time $`t`$. Let $`d_Y`$ denote a metric on $`Y`$, obtained from a bi-$`K`$ invariant metric $`d`$ on $`G`$.
###### Theorem 1.1.
(Kleinbock-Margulis ) There exists a $`k=k(Y)>0`$ such that the following holds: for all $`x,yY`$, and almost every $`\theta T_x^1(Y)`$,
$$\underset{t\mathrm{}}{lim\; sup}\frac{d_Y(g_t(x,\theta ),y)}{\mathrm{log}t}=1/k.$$
The Theorem thus studies the statistical properties of geodesic excursions into smaller and smaller cuspidal neighborhoods (i.e. the complements of these neighborhoods are larger and larger compacta) of $`Y`$. In fact, this phenomenon seems to be prevalent in many dynamical systems, and a very general result to the effect was obtained in , which can then be used to obtain Theorem 1.1. To describe this result, we need some more definitons, all taken from loc. cit. Let $`(X,\mu )`$ be a probability space, $``$ be a family of measurable subsets of $`X`$ and let $`F=\{f_t\}`$ denote a sequence of $`\mu `$-preserving transformations of $`X`$. Then,
###### Definition 1.2.
$``$ is called Borel-Cantelli for $`F`$ if for every sequence $`\{A_t|t\}`$ of sets from $``$, the following holds:
$$\mu (\{xX|f_t(x)A_t\text{for infinitely many}t\})$$
$$=\{\begin{array}{ccc}0& \text{if }_{t=1}^{\mathrm{}}\mu (A_t)\text{ converges},& \\ & & \\ 1& \text{if }_{t=1}^{\mathrm{}}\mu (A_t)\text{ diverges}.& \end{array}$$
For a function $`\delta `$ on $`X`$, denote by $`\mathrm{\Phi }_\delta `$, the tail distribution function, defined by:
(1.3)
$$\mathrm{\Phi }_\delta (z)=\mu (\{x|\delta (x)z\}).$$
And for $`\kappa >0`$, we say that $`\delta `$ is $`\kappa DL`$ (an abbreviation for $`\kappa `$distance like) if it is uniformly continuous and
(1.4)
$$C_1,C_2>0,\text{such that}C_1e^{\kappa z}\mathrm{\Phi }_\delta (z)C_2e^{\kappa z}z,$$
and $`DL`$ if there exists $`\kappa >0`$ such that (1.4) holds. The following is then Theorem $`1.8`$ in .
###### Theorem 1.5.
Let $`G`$, $`\mathrm{\Gamma }`$ and $`X`$ be as above, $`F=\{f_t|t\}`$ be a sequence of elements of $`G`$ satisfying
(1.6)
$$\underset{t}{sup}\underset{s=1}{\overset{\mathrm{}}{}}e^{\beta d(f_sf_t^1,e)}<\mathrm{}\beta >0,$$
and let $`\delta `$ be a $`DL`$ function on $`X`$. Then
(1.7)
$$(\delta )\stackrel{def}{=}\left\{\{xX|\delta (x)r\}\right|r\}$$
is Borel-Cantelli for $`F`$.
Sequences which satisfy (1.6) above are referred to as “exponentially divergent” (abbreviated $`ED`$) in . To derive Theorem 1.1 from Theorem 1.5, the authors use a philosophy of F.Mautner to realize the geodesic flow on the unit tangent bundle of $`Y`$ as a one-parameter flow on $`G/\mathrm{\Gamma }`$. One then needs to check the $`ED`$ condition for this flow, as well as the $`DL`$ condition for the metric on $`G/\mathrm{\Gamma }`$ described above. In this paper and its sequel , we will obtain $`S`$-arithmetic analogues of Theorem 1.5.
To motivate our results, we start with an analogue of Theorem 1.1. Let $`𝔽_s`$ denote the finite field of $`s=p^\vartheta `$ elements, $`𝐤`$ denote the global function field of rational functions with coefficients in $`𝔽_s`$ and $`k`$ denote the completion of $`𝐤`$ at the infinite place, identified naturally with the field of Laurent series $`𝔽_s((X^1))`$ with coefficients in $`𝔽_s`$. On $`𝔽_s((X^1))`$, we will assume the usual norm and ultrametric, (cf.) for details. We denote this norm $``$ and when used in the context of $`k^r`$, it will refer to the supremum norm. Let $`𝔾`$ denote a simple, isotropic, linear algebraic group defined and split over $`𝔽_s`$, $`\mathrm{\Gamma }`$ a non-uniform lattice in $`𝔾(k)`$<sup>3</sup><sup>3</sup>3It is well-known $`G`$ contains many such lattices., and $`K`$ denote a parahoric subgroup of $`𝔾(k)`$. We fix a bi $`K`$-invariant metric on $`𝔾(k)`$, and hence also on $`𝔾(k)/\mathrm{\Gamma }`$ and call this metric $`d(,)`$. When the field in question is evident, we will sometimes refer to $`𝔾(k)`$ simply as $`G`$.
A natural geometric object on which $`G`$ acts is the Bruhat-Tits building $`Y`$ of $`G`$. This is a Euclidean building which comes with a natural metric, equipped with which it becomes a CAT-$`0`$ space. The stabilizers of vertices in $`Y`$ are the maximal parahoric subgroups $`K`$ of $`G`$. Therefore the quotient $`Y/\mathrm{\Gamma }`$ can be naturally identified with $`K\backslash G/\mathrm{\Gamma }`$. Let $`Y`$ denote the geodesic boundary of $`Y`$ and $`\pi :YY/\mathrm{\Gamma }`$ be the natural projection. Let $`\{g_t(x,\theta )\}_{t0}`$ denote the geodesic starting at the vertex $`xY`$ in direction $`\theta Y`$. There is a natural measure class on $`Y`$. A function field analogue of Theorem 1.1 would therefore be:
###### Theorem 1.8.
There is a $`l=l(Y/\mathrm{\Gamma })`$ such that for any $`xY`$, $`yY/\mathrm{\Gamma }`$ and almost all $`\theta Y`$,
$$\underset{t\mathrm{}}{lim\; sup}\frac{\mathrm{log}d_{Y/\mathrm{\Gamma }}(\pi (g_t(x,\theta )),y)}{\mathrm{log}t}=1/l,$$
where $`d_Y`$ is the metric on the quotient $`Y/\mathrm{\Gamma }`$.
Remark: This result was obtained by S. Hersonsky and F. Paulin in in the case of quotients of a regular tree $`Y`$ by any non-uniform lattice in $`Aut(Y)`$. Their approach is more geometric, since $`Aut(Y)`$ is a locally compact group but does not have any obvious algebraic structure, and in some sense is in the spirit of Sullivan . A natural question would be to ask if an analogue of Theorem 1.5 holds in a more general algebraic setting. We answer this in the affirmative.
Let $`𝐤`$ be a global field, $`S`$ a finite set of places of $`𝐤`$ (containing the infinite ones, in case $`𝐤`$ is a number field), and for each $`sS`$, let $`k_s`$ denote the completion of $`𝐤`$ at the place $`s`$, $`𝔾`$ be a connected, simple, algebraic $`𝐤`$-group without anisotropic factors, set $`G_s=𝔾(k_s)`$ and $`G_S=_{sS}G_s`$. Let $`Y_s`$ denote the Bruhat-Tits building of $`G_s`$ (or the symmetric space as the case may be) and $`Y_S=_{sS}Y_s`$. Let $`\mathrm{\Gamma }`$ be a non-uniform lattice in $`G_S`$ and set $`X_S=G_S/\mathrm{\Gamma }`$. In , we prove:
###### Theorem 1.9.
With notation as above, let $`F`$ be an exponentially divergent sequence for $`G_S`$ and $`\delta `$ be a $`DL`$ function on $`X_S`$. Then $`(\delta )`$ is Borel-Cantelli for $`F`$.
Remark 1. The reader will notice that Theorem 1.5 is stated in somewhat greater generality, i.e. in the context of semisimple groups. However, the principal ideas involved are already present in the case of simple groups. In where the proof of the above theorem will be presented in detail, we will point out the extra details required to extend Theorem 1.9 to semisimple groups.
Remark 2. As a corollary, we will obtain generalizations of Theorems 1.1 and 1.8 for $`S`$-isotropic symmetric spaces.
We now turn to some examples : Let $`𝒪_S`$ be the ring of $`S`$-integers of $`𝐤`$, i.e.
(1.10)
$$𝒪_S=\{x𝐤||x|_s1\text{for every finite place}sS\}.$$
Then, by theorems of Borel-Harish Chandra and Behr-Harder in the function field case (cf.), $`\mathrm{\Gamma }\stackrel{def}{=}𝔾(𝒪_S)`$ is a lattice<sup>4</sup><sup>4</sup>4This, and commensurable lattices are called arithmetic. in $`G_S`$. Let $`K_s`$ denote a maximal parahoric subgroup (maximal compact if $`s`$ is an infinite place) of $`G_s`$, and $`d_s`$ denote the $`K_s`$-bi invariant metric on $`G_s`$. This gives rise to a product metric $`d`$ on $`X_S`$. In , using reduction theory we show that $`d`$ is a $`DL`$ function on $`X_S`$ for *any* non-uniform lattice $`\mathrm{\Gamma }`$, thereby establishing Theorem 1.8 by an application of Theorem 1.9.
For an application to number theory, take $`G=\mathrm{SL}(n,),\mathrm{\Gamma }=\mathrm{SL}(n,)`$, so the space
(1.11)
$$X=\mathrm{SL}(n,)/\mathrm{SL}(n,)$$
can be identified with the space of unimodular lattices in $`^n`$. In , Theorem 1.5 and the function
(1.12)
$$\delta (\mathrm{\Lambda })=\underset{𝐱\mathrm{\Lambda }\backslash \{0\}}{inf}𝐱$$
which turns out to be $`DL`$, was used to prove Khintchine’s theorem, a cornerstone of metric Diophantine approximation, and stronger variants of the it. Now let $`𝐤=`$ and following , define
(1.13)
$$\mathrm{GL}^1(n,_S)\stackrel{def}{=}\left\{g=(g_s)_{sS}\mathrm{GL}(n,_S)\right|\underset{sS}{}det(g_s)_{sS}=1\}$$
Then $`\mathrm{GL}^1(n,_S)/\mathrm{GL}(n,_S)`$ is the space of unimodular lattices in $`_S`$. In , a dynamical interpretation of Diophantine approximation, originally introduced by D. Kleinbock and G. Margulis in was developed in the $`S`$-arithmetic setting by D. Kleinbock and G. Tomanov to prove $`S`$-arithmetic analogues of Mahler’s conjectures, see also for function field analogues. In , we use Theorem 1.9 in this context with specific, analogous choices of $`DL`$ functions to establish Khintchine’s theorems and multiplicative versions in the number and function field cases. We illustrate with an example. Let $`l`$ denote the cardinality of $`S`$, and assume, for simplicity that $`S`$ contains the infinite valuation<sup>5</sup><sup>5</sup>5This is just a convenience because the definitions are slightly more involved cf.. This restriction will be removed in . Let $`\psi :_+`$ be a non-increasing function, and let $`\mathrm{Mat}(_S)`$ denote the set of $`m\times n`$-matrices with entries in $`_S`$. Let $`𝒲(\psi )`$ denote the set of $`A\mathrm{Mat}(_S)`$ for which there exist infinitely many $`𝐪^n`$ such that
(1.14)
$$(𝐩+A𝐪^l)^m\psi (𝐪_{\mathrm{}}^n)\text{for some}𝐩^m.$$
Here the norms $``$ and $`_{\mathrm{}}`$ are defined as follows:
For $`𝐱=(x_1^{(v)},x_2^{(v)},\mathrm{},x_m^{(v)})_v^m`$, we define
(1.15)
$$𝐱_v=\underset{i}{\mathrm{max}}|x_i^{(v)}|_v$$
and finally, for $`𝐱_S^m`$
(1.16)
$$𝐱=\underset{v}{\mathrm{max}}𝐱_v.$$
For a very nice motivation of the theory of $`S`$-arithmetic Diophantine approximation, and especially the correct normalization, we refer the reader to Section $`10`$ in . The $`S`$-arithmetic analogue of the Khintchine-Groshev theorem would then be:
###### Theorem 1.17.
Let $`𝒲(\psi )`$ denote the set of $`\psi `$-approximable matrices as above. Then,
1. Almost every $`A`$ belongs to $`𝒲(\psi )`$ if $`\psi (x)𝑑x`$ diverges.
2. Almost no $`A`$ belongs to $`𝒲(\psi )`$ if $`\psi (x)𝑑x`$ converges.
This theorem will be proved in by adapting the method of Kleinbock-Margulis to this set-up. In addition, we will prove variations and strengthenings of the above theorem, including:
1. The more general multiplicative case, where the norm $``$ is replaced by a more general function. This will prove $`S`$-arithmetic analogues of a conjecture of Skriganov .
2. Function field analogues of Theorem 1.17. This is elaborated upon in the next section.
This paper is intended as an announcement of our results above. In particular, we will not present complete proofs, but rather try to convey the main ideas involved in the proof. To do this, we will focus on the case of a single local field of positive characteristic. Namely, we will outline the proofs of Theorems 1.8 and 1.17 in the case where $`k`$ is a local field of positive characteristic. We would like to stress that our strategy is substantially similar to that of Kleinbock-Margulis. The primary content of our work is in generalizing their methods to much wider settings. However, in , we will also provide a direct, “non-spectral” proof of Theorem 1.8 in the case where $`𝔾`$ is a rank $`1`$ group, which has been explained to us by S. Mozes. In this case, we use a symbolic description of the geodesic flow on trees, developed in . We are indebted to him for sharing his insight.
### Acknowledgements.
A. G. announced these results at an AMS regional conference in Chicago, Ill. and would like to thank Marian Gidea for his hospitality. J. A. and A. G. thank D. Kleinbock and the participants of the Workshop on Shrinking Target Properties, Feb. 2008 for helpful conversations. We also thank the Clay Mathematical Institute for sponsoring the workshop. We thank the referee for several helpful comments which have improved the exposition of the paper.
## 2. Borel-Cantelli lemmata and applications.
Theorem 1.9, like Theorem 1.5 are probabilistic in nature-in fact they are strongly reminiscent of $`01`$ laws in probability theory, especially the elementary Borel-Cantelli lemma, which we now recall:
###### Lemma 2.1.
(Borel-Cantelli) Let $`\{X_n\}_{n=0}^{\mathrm{}}`$ be a sequence of $`01`$ random variables, with $`P(X_n=1)=:p_n`$. Then
1. If $`_{n=0}^{\mathrm{}}p_n<\mathrm{}`$, then $`P(_{n=0}^{\mathrm{}}X_n=\mathrm{})=0`$
2. If the $`X_n`$’s are pairwise independent ,i.e.
$$p_{nm}:=P(X_nX_m=1)=p_np_mm,n,$$
and $`_{n=0}^{\mathrm{}}p_n=\mathrm{}`$, then $`P(_{n=0}^{\mathrm{}}X_n=\mathrm{})=1`$.
The first part of the above Lemma easily allows one to derive the convergence halves of the various $`01`$ laws we have in mind. Unfortunately, it is typically hopeless to expect that the dynamical random variables (or events) we are interested in are independent, i.e. satisfy $`(2)`$ in the above theorem. In the context of logarithm laws, for instance, one is only able to show that geodesic excursions to shrinking cusp neighborhoods are relatively (sometimes referred to as “quasi”) independent events. That this suffices for applications is due to the following strengthening of the Borel-Cantelli lemma, abstracted from the works of W. Schmidt, by V. Sprindzhuk . We first set up some of the notation, taken from . For a function $`f`$ on a probability space $`(X,\mu )`$, we will denote
$$\mu (f)\stackrel{def}{=}_Xf𝑑\mu .$$
and for $`N\{\mathrm{}\}`$, a family of functions $`=\{h_t|t\}`$ on $`X`$,
(2.2)
$$S_{,N}\stackrel{def}{=}\underset{t=1}{\overset{N}{}}h_t(x)\text{and}E_{,N}\stackrel{def}{=}\underset{t=1}{\overset{N}{}}\mu (h_t).$$
###### Lemma 2.3.
Let $`(X,\mu )`$ denote a probability space, and let $`=\{h_t|t\}`$ denote a sequence of functions on $`X`$ which satisfy:
(2.4)
$$\mu (h_t)1\text{for every}t.$$
Assume also that there exists $`C>0`$ such that
(2.5)
$$\underset{s,t=M}{\overset{N}{}}\left(\mu (h_sh_t)\mu (h_s)\mu (h_t)\right)C.\underset{t=M}{\overset{N}{}}\mu (h_t)\text{for every}N>M1.$$
Then for every $`ϵ>0`$,
(2.6)
$$S_{,N}=E_{,N}+O\left(E_{,N}^{1/2}\mathrm{log}^{3/2+ϵ}E_{,N}\right)$$
for $`\mu `$ almost every $`xX`$. In particular, for $`\mu `$ almost every $`x`$,
(2.7)
$$\underset{N\mathrm{}}{lim}\frac{S_{,N}(x)}{E_{,N}}=1,$$
whenever $`_{t=1}^{\mathrm{}}\mu (h_t)`$ diverges.
So, in order to use Sprindzhuk’s lemma to prove Theorem 1.9, we need to show that the “quasi-independence” condition (2.5) holds in the context of our dynamical systems. We start with the Khintchine-Groshev theorem and its generalizations, focussing on the function field case. Our plan is to first describe a dynamical system to which the above Lemma can be applied to derive the desired number theoretic results, and then to use quantitative mixing bounds for the dynamical system to ensure that condition (2.5) is satisfied.
Accordingly, let us take $`G=\mathrm{SL}(n,𝔽_s((X^1)))`$ and $`\mathrm{\Gamma }=\mathrm{SL}(n,𝔽_s[X])`$. Then $`\mathrm{\Gamma }`$ is a non-compact lattice in $`G`$. Let $`X=G/\mathrm{\Gamma }`$ and $`\mu `$ denote the normalized, projected Haar measure on $`X`$. Diophantine approximation in the function field setting is naturally analogous<sup>6</sup><sup>6</sup>6This is true to some extent. There are important and much studied distinctions, especially in the approximation of algebraic elements. We will not address this. to the real case, namely the role of $``$ is played by $`𝔽_s((X^1))`$, while that of $``$ is played by the polynomial ring $`𝔽_s[X]`$, and there is an analogous continued fraction decomposition for every element of $`𝔽_s((X^1))`$ (see, for example, ). As would be expected, it is possible to read off Diophantine properties of Laurent series from their continued fraction expansions, and in fact one can provide a proof of Khintchine’s theorem in one dimension using a careful study of these continued fractions.
More generally, the set $`𝒲(\psi )`$ of $`\psi `$-approximable $`(m\times n)`$ matrices with entries in $`𝔽_s((X^1))`$ is naturally defined to be those $`A\mathrm{Mat}_{m,n}(𝔽_s((X^1)))`$ for which there exist infinitely many $`𝐪𝔽_s[X]^n`$ such that
(2.8)
$$A𝐪+𝐩^m<\psi (𝐪^n)\text{for some}𝐩𝔽_s[X]^m.$$
And the function field analogue of the Khintchine-Groshev Theorem would precisely be Theorem 1.17 with the above definition of $`𝒲(\psi )`$.
An ingenious scheme, due to Dani in a special case, and developed in full generality by Kleinbock-Margulis translates the above problem into one of shrinking target properties for certain homogeneous flows. Let us briefly describe this correspondence. The group $`\mathrm{SL}_{m+n}(𝔽_s((X^1)))`$ acts transitively on the space of unimodular (i.e. co-volume $`1`$) lattices $`\mathrm{\Omega }_{m+n}`$ of $`𝔽_s((X^1))^{m+n}`$ and the stabilizer of $`𝔽_s[X]^{m+n}`$ is $`\mathrm{SL}_{m+n}(𝔽_s[X])`$.
The space $`\mathrm{SL}_{m+n}(𝔽_s((X^1)))/\mathrm{SL}_{m+n}(𝔽_s[X])`$ can thus be naturally identified with $`\mathrm{\Omega }_{m+n}`$. Given $`A\mathrm{Mat}_{m,n}(𝔽_s((X^1)))`$, we associate to it the following lattice:
$$A\mathrm{\Lambda }_A\stackrel{def}{=}\left(\begin{array}{cc}I_m& A\\ 0& I_n\end{array}\right)𝔽_s[X]^{m+n}$$
where $`I_i`$ denotes the square identity matrix of dimension $`i`$. Good rational approximations to $`A`$ can be shown to correspond to small vectors in the lattice $`\mathrm{\Lambda }_A`$ corresponding to $`A`$. Let $`\pi `$ denote the uniformizer of $`k`$ and set
(2.9)
$$g_t\stackrel{def}{=}diag(\underset{m\text{times}}{\underset{}{\pi ^{nt},\mathrm{},\pi ^{nt}}}\underset{n\text{times}}{\underset{}{\pi ^{mt},\mathrm{},\pi ^{mt}}}).$$
and let $`\delta `$ denote the following function which measures small vectors in lattices:
(2.10)
$$\delta =\underset{𝐯\mathrm{\Lambda }\backslash \{0\}}{\mathrm{min}}𝐯.$$
By the above identification of
$$\mathrm{\Omega }_{m+n}\text{with}\mathrm{SL}(m+n,𝔽_s((X^1)))/\mathrm{SL}(m+n,𝔽_s[X])$$
and Mahler’s compactness criterion, those $`t`$ for which $`\delta (g_t\mathrm{\Lambda }_A)`$ is small correspond to $`t`$ for which the $`g_t`$ trajectory of $`\mathrm{\Lambda }_A`$ has cusp excursions. To define these excursions, we set
(2.11)
$$\mathrm{Cusp}(t,\delta )\stackrel{def}{=}\{\mathrm{\Lambda }\mathrm{\Omega }_{m+n}|\delta (\mathrm{\Lambda })|\pi |^t\}.$$
The following lemma, proven<sup>7</sup><sup>7</sup>7In the real case, but the function field proof goes through with obvious modifications. An $`S`$-arithmetic version will be provided in . in then establishes the precise connection between Diophantine properties and cusp excursions.
###### Lemma 2.12.
There exists an explicit function $`r(t)`$ which depends only on $`\psi ,m`$ and $`n`$ and satisfies
$$\psi (x)𝑑x<\mathrm{}|\pi |^{(m+n)r(t)}𝑑t<\mathrm{}.$$
Moreover, $`A𝒲(\psi )`$ if and only if
(2.13)
$$\text{infinitely many}t\text{such that}g_t\mathrm{\Lambda }_A\mathrm{Cusp}(r(t),\delta ).$$
So in order to prove Theorem 1.17 for function fields, it clearly suffices to prove that:
###### Theorem 2.14.
With notation as above,
1. Almost every $`\mathrm{\Lambda }_A`$ satisfies (2.13) if $`|\pi |^{(m+n)r(t)}𝑑t`$ diverges.
2. Almost no $`\mathrm{\Lambda }_A`$ satisfies (2.13) if $`|\pi |^{(m+n)r(t)}𝑑t`$ converges.
Similarly, in order to prove logarithm laws for $`K\backslash 𝔾(k)/\mathrm{\Gamma }`$, where $`𝔾`$ is a simple group as before, we will use the philosophy of F. Mautner to realize the geodesic flow on the unit tangent bundle of $`K\backslash 𝔾(k)/\mathrm{\Gamma }`$ as a one-parameter flow $`g_t`$ on $`𝔾(k)/\mathrm{\Gamma }`$. Fix $`x_0𝔾(k)/\mathrm{\Gamma }`$, and set
(2.15)
$$\mathrm{Cusp}(t,d)=\{y𝔾(k)/\mathrm{\Gamma }|d(x_0,y)>t\}.$$
The object of study then becomes the excursions of $`g_t`$ orbits into $`\mathrm{Cusp}(t,d)`$.
We now tie these themes up with Sprindzhuk’s lemma, focussing for convenience on Theorem 2.14. Accordingly, let $`h_t`$ denote the characteristic function of $`\mathrm{Cusp}(r(t),\delta )`$, and let
(2.16)
$$^𝒢=\{g_t^1h_t|t\}.$$
Then to derive Theorem 2.14 from Sprindzhuk’s Lemma, we need to ensure that
(2.17)
$$\underset{s,t=M}{\overset{N}{}}\left((g_s^1h_s,g_t^1h_t)\mu (h_s)\mu (h_t)\right)$$
is small, in other words we want to show that the excursions of $`g_t`$ orbits into $`\mathrm{Cusp}(r(t),\delta )`$ are quasi-independent. The key to this is the spectral gap of the $`G`$-action on $`G/\mathrm{\Gamma }`$.
## 3. Spectral Gap.
We retain the notation of the previous section, i.e. $`G`$ is the group of $`k`$-points of a simple $`𝔾`$ as before, and $`\mathrm{\Gamma }`$ is a non-uniform lattice in $`G`$. A natural tool to study the ergodic properties of the action of $`G`$ on $`G/\mathrm{\Gamma }`$ is the spectral properties of the action of $`G`$ on $`L^2(G/\mathrm{\Gamma })`$. Let $`L_0^2(G/\mathrm{\Gamma })`$ denote the subspace of $`G/\mathrm{\Gamma }`$ orthogonal to constants. Let $`K`$ denote a maximal compact open subgroup of $`G`$. We call $`\varphi L_0^2(G/\mathrm{\Gamma })`$, $`K`$-finite if its $`K`$-span is finite dimensional. On $`G`$ there is an important function which controls the rate of decay of matrix coefficients, the Harish-Chandra function $`\mathrm{\Xi }`$. Let $`G=KAN`$ denote the Iwasawa decomposition of $`G`$. The Harish-Chandra function of $`G`$ is then defined by:
(3.1)
$$\mathrm{\Xi }(g)=_K\delta ^{1/2}(gk)𝑑k$$
where $`\delta `$ is the left modular function of $`AN`$ defined by:
(3.2)
$$d\mu _G=d\mu _K\delta (an)d\mu _{AN}$$
where $`\mu _{}`$ denotes Haar measure on the locally compact group $``$. The following estimate for decay of matrix coefficients is then known to hold by work of several authors. We refer to Section 5.1 in as a convenient reference.
###### Theorem 3.3.
Let $`G,K`$ and $`\mathrm{\Gamma }`$ be as above. Then there exist positive constants $`C,\chi `$ such that for every $`K`$-finite $`\varphi ,\psi L_0^2(G/\mathrm{\Gamma })`$ and any $`gG`$,
(3.4)
$$|<g\varphi ,\psi >|C\varphi \psi \mathrm{\Xi }(g)^{1/\chi }.$$
It turns out that the above estimate, or more precisely a version of this estimate which caters to smooth functions (the $`L^2`$-norm is then replaced by an appropriate Sobolev norm) is enough to ensure the quasi-independence condition in Sprindzhuk’s lemma. The passage from $`L^2`$ to smooth functions is quite standard and will be elaborated upon in the $`S`$-arithmetic setting in . It remains to show that the characteristic functions of $`\mathrm{Cusp}(r(t),\delta )`$ (resp. $`\mathrm{Cusp}(t,d)`$) can be approximated by appropriate smooth functions, which is precisely showing that the functions in question are $`DL`$. This “smoothing of cusp neighborhoods” is carried out using reduction theory, which we elaborate upon in the next section.
We return briefly to the proof of Theorem 3.3. It turns out that these bounds are closely related to the notion of spectral gap for the $`G`$ action. Recall that the $`G`$ action on a probability space $`X`$ has spectral gap if the regular representation $`\rho _0`$ of $`G`$ on the space $`L_0^2(X)`$, the subspace of $`L^2(X)`$ orthogonal to constants, is isolated in the Fell topology from the identity (or trivial) representation. If $`G`$ is an almost direct product of groups $`G_i`$, then the $`G`$ action on $`X`$ has strong spectral gap if the restriction of $`\rho _0`$ to any factor $`G_i`$ is isolated from the trivial representation.
Establishing strong spectral gap turns out to be quite difficult in general. In , the authors proved that if $`G`$ is a connected semisimple, center-free Lie group without compact factors, $`\mathrm{\Gamma }`$ is an irreducible non-uniform lattice in $`G`$, then the $`G`$ action on $`G/\mathrm{\Gamma }`$ has strong spectral gap. The analogous question for uniform lattices does not seem to be known, however see the recent work . Of course, if all the factors $`G_i`$ have Kazhdan’s property T, strong spectral gap is immediate.
However, since rank $`1`$-groups defined over non-Archimedean local fields act on their Bruhat-Tits trees without fixed points, they cannot have property $`T`$. A weaker and very useful property is property $`\tau `$ as defined by Lubotzky and Zimmer , which means that $`\rho _0`$ is isolated in the Automorphic Spectrum of $`G`$. In , we will gather various tools from representation theory namely the restriction technique of Burger-Sarnak , as extended to finite places and function fields by Clozel-Ulmo (cf. and ) and property $`\tau `$ for congruence subgroups of $`\mathrm{SL}_2`$ as established by Selberg , Gelbart-Jacquet , Clozel and Drinfeld in various contexts, to record a proof of (the $`S`$-arithmetic generalization) of Theorem 3.3.
We note that in many cases uniform, optimal estimates for decay of matrix coefficients are known due to H.Oh and these have several important applications. While the use of spectral gap is a powerful tool, we remark that very little is known in the context of the automorphism group of a tree or more generally, a building and this poses a potential hurdle to obtain logarithm laws in these settings using spectral methods.
The reader will notice that while the existence of the constant $`\chi `$ in Theorem 3.3 is crucial as well as sufficient for the purposes of this paper, we do not say anything regarding bounds, i.e. the so-called extent of temperedness of the representation. These bounds are related to bounds towards the Generalized Ramanujan Conjecture and play a crucial role in many Diophantine problems. In , the second named author and A. Gorodnik study the connection of Diophantine approximation on symmetric spaces and bounds on temperedness in greater detail. This problem involves studying shrinking target properties for targets which are balls around a fixed point in the symmetric space, and needs new techniques.
## 4. Adèlic Reduction Theory.
We now present the tools required to obtain smoothing of the cusp neighborhoods, i.e. to show that the characteristic functions of super-level sets of the various cusp neighborhoods are $`DL`$. The main tool is adèlic reduction theory, as developed in and from where we also borrow notation. Let $`𝐤,𝔾`$ etc. be as before and let $`𝔸`$ denote the adèle ring of $`𝐤`$. Let $`B`$ be a fixed Borel subgroup of $`𝔾(k)`$ containing a maximal split torus $`T`$ of $`G`$. Let $`X_{}(T)`$ denote the lattice of cocharacters of $`T`$. Let $`X_{}(T)^{++}`$ denote the subset of $`X_{}(T)`$ consisting of dominant cocharacters. Let $`W=N_GT(𝔽_s)/T(𝔽_s)`$ denote the Weyl group of $`G`$. Define the map $`\varphi _\pi :W\times X_{}(T)G(\overline{k})`$ by
(4.1)
$$\varphi _\pi (w,\mu )=w\mu (\pi ).$$
Let $`\rho `$ denote the sum of all roots that are positive with respect to $`B`$. For $`\tau >0`$, we denote by $`X_\tau `$ the subset of $`𝔾(𝔸)/𝔾(k)`$ consisting of the union of $`𝔾(k)\varphi _{t^1}(e^\lambda )K`$ where $`\lambda `$ ranges over the dominant cocharacters of $`T`$ for which $`\rho ,\lambda \tau `$. Then, the following estimate, which follows from the reduction theory developed in and allows us to show the requisite $`DL`$ property for $`d(x_0,)`$.
###### Theorem 4.2.
Let $`r`$ denote the $`k`$-rank of $`𝔾(k)`$. Then, there exist positive real numbers $`C_1`$ and $`C_2`$ such that for every $`T>0`$,
(4.3)
$$C_1<\frac{\mu (X_T)}{_{lT}s^ll^{r1}}<C_2.$$
We now turn to Khintchine’s theorem , where in light of the previous discussion, we have to develop the reduction theory in a slightly different setting-namely on the space $`\mathrm{SL}(d,𝔽_s((X^1)))/\mathrm{SL}(d,𝔽_s[X])`$ and for the function $`\delta `$. The key to this, as shown in , lies in the so-called Siegel integral formula , specifically to a multi-dimensional version of it. We will again proceed in the adèlic setting. Set:
(4.4)
$$\mathrm{GL}^1(d,𝔸)\stackrel{def}{=}\{g\mathrm{GL}(d,𝔸)|\underset{vV(𝐤)}{}|detg|_v=1\}$$
where $`V(𝐤)`$ denotes the set of inequivalent places of the global field $`𝐤`$. Let $`fL^1(𝔸^d)`$, $`g\mathrm{GL}^1(d,𝔸)`$ and define
(4.5)
$$\stackrel{~}{f}\stackrel{def}{=}\underset{𝐱g^d}{}f(𝐱)$$
Then the adèlic version of Siegel’s integral formula, as proved by Weil in states that:
###### Theorem 4.6.
$$_{𝔸^d}f𝑑\mu _𝔸=C(d)_{\mathrm{GL}^1(d,𝔸)/\mathrm{GL}(d,)}\stackrel{~}{f}𝑑\stackrel{~}{\mu }.$$
where $`C(d)`$ is a constant depending on the field in question (in this case $``$). We remark that the above Theorem was proved by Weil for arbitrary number fields as well as function fields. A multidimensional generalization of this theorem is developed in and is shown to imply a local version of this formula which can then be used to show that the function $`\delta `$ is $`DL`$. |
warning/0506/physics0506073.html | ar5iv | text | # Variational principle for linear stability of moving magnetized plasma.
## Abstract
The variational principle for linear stability of three-dimensional, inhomogenious, compressible, moving magnetized plasma is suggested. The principle is “softer” (easier to be satisfied) than all previously known variational stability conditions. The key point of the analysis is a conservation in variations of new integrals inherent in the linearized equation of the motion that was not earlier discussed in the literature.
PACS46.15.Cc, 52.30.Cv
It is well known that stability of the static equilibrium of magnetized plasma can be described by so-called “energy principle” . The principle claims that if the second variation of the potential energy, $`W`$, of the system “plasma-magnetic field” is positive definite near the equilibrium point, then this point is stable. Sufficiency of the claim follows from the Lyapunov stability theorem, and necessity can also be proved . Note that the above mentioned second variation of the potential energy corresponds exactly to the potential energy of the linearized equation system.
The main drawback of the principle is that there always may be neutral perturbations which do not perturb any physical quantity – and, therefore, $`W`$ as well. Thus, the second variation of the potential energy can be guaranteed to be only positive semi-definite. In other words, using the energy principle , one can talk about spectral stability only, namely, about presence or absence of imaginary frequencies in the spectrum of the linearized force operator (nonlinear stability needs an analysis of neutral perturbations – see, e.g., ). Continuing this logic, we restrict ourselves with linearized equations.
The attempt of using the similar approach to investigate stability of moving plasma performed by Frieman and Rotenberg was not so lucky, although the energy principle was formally obtained. Their result can be briefly described as follows. Consider the linearized equation of motion for plasma displacement $`𝝃`$ in the frame of ideal one-fluid magnetohydrodynamics,
$$\rho \ddot{𝝃}+2\rho (𝐕\mathbf{})\dot{𝝃}𝐅(𝝃)=0,$$
(1)
where the linearized force operator,
$`𝐅(𝝃)`$ $`=`$ $`\delta \rho (𝐕\mathbf{})𝐕\rho (\delta 𝐕\mathbf{})𝐕\rho (𝐕\mathbf{})\delta 𝐕`$
$``$ $`\mathbf{}\delta p+(\mathbf{}\times \delta 𝐁)\times 𝐁+(\mathbf{}\times 𝐁)\times \delta 𝐁,`$
is combined of usual perturbed quantities,
$$\delta \rho =\mathbf{}(\rho 𝝃),\delta 𝐕=(𝐕\mathbf{})𝝃(𝝃\mathbf{})𝐕,$$
$$\delta p=𝝃\mathbf{}p\gamma p\mathbf{}𝝃,\delta 𝐁=\mathbf{}\times (𝝃\times 𝐁)$$
(note that $`\delta 𝐕`$ denotes here only the part of full Eulerian velocity perturbation – the part, which survives even for time-independent displacements, $`𝝃`$). Stationary plasma density, $`\rho `$, velocity, $`𝐕`$, pressure, $`p`$, and magnetic field, $`𝐁`$, satisfy the following equilibrium conditions:
$`\rho (𝐕\mathbf{})𝐕+\mathbf{}p=(\mathbf{}\times 𝐁)\times 𝐁,`$
$`\mathbf{}(\rho 𝐕)=0,`$
$`𝐕\mathbf{}p+\gamma p\mathbf{}𝐕=0,`$
$`\mathbf{}\times (𝐕\times 𝐁)=0.`$
Dot means a partial time-derivative, $`\gamma `$ means the adiabatic exponent. Force operator is proved to be self-adjoint in the following sense,
$$𝜼𝐅(𝝃)d^3r=𝝃𝐅(𝜼)d^3r,$$
while the second term in (1) is obviously antisymmetric:
$$𝜼\rho (𝐕\mathbf{})𝝃d^3r=𝝃\rho (𝐕\mathbf{})𝜼d^3r.$$
Multiplying Eq. (1) by $`\dot{𝝃}`$ and integrating over the whole space, we found the energy conservation in the form $`\dot{E}=0`$, where
$$E(t)=\left(\rho \frac{\dot{𝝃}^2}{2}\frac{𝝃𝐅(𝝃)}{2}\right)d^3r.$$
(2)
Minimizing $`E`$ over $`\dot{𝝃}`$, we approach to the energy principle by Frieman-Rotenberg,
$$𝝃𝐅(𝝃)d^3r0.$$
(3)
Contrary to the static case $`(𝐕=0)`$, in which condition (3) appears to be both sufficient and necessary for linear stability, in the case of $`𝐕0`$, condition (3) is normally too strong, and never can be satisfied except of field-aligned flows $`(𝐕𝐁)`$ or of those which may be reduced to the field-aligned flows (see, e.g., ).
Energy principle (3) may be improved by use of the Arnold conjecture -, following which we have to add to the energy (2) the set of other known integrals of the motion. Speaking in other words, variables $`\dot{𝝃}`$ and $`𝝃`$ in (2) are not absolutely independent but subject to the constraints resulting from conservation of other integrals of the motion.
Such an improved principle was derived by Ilgisonis and Pastukhov , then it was verified by Hameiri . It was also re-obtained with help of Pfirsch & Morrison’s method of dynamically accessed perturbations. That stability condition is currently the best among the known ones, although it is still not appropriate for arbitrary stationary plasma flow.
For the linearized equation (1), the Ilgisonis & Pastukhov extra invariants can be written in terms of neutral perturbation $`𝝃_N`$:
$$𝐅(𝝃_N)=0,_t𝝃_N=0.$$
Multiplying Eq. (1) by $`𝝃_N`$ and integrating again over the space, we have $`\dot{I}=0`$, where
$$I=(\rho \dot{𝝃}𝝃_N+2\rho 𝝃_N(𝐕\mathbf{})𝝃)d^3r.$$
(4)
For the system with nested set of magnetic surfaces, $`\psi (𝐫)=\mathrm{const}`$, $`𝝃_N`$ may be generally represented as
$$𝝃_N=\lambda _u(\psi )𝐮+\lambda _v(\psi )𝐯,$$
(5)
where $`𝐮=𝐁/\rho ,𝐯=𝐃/\rho `$, and $`𝐃`$ is a divergence-free frozen-in-plasma vector, tangential to the same magnetic surfaces, $`\psi (𝐫)=\mathrm{const}`$, but different from $`𝐁`$,
$$𝐁\times 𝐃=\rho \mathbf{}\psi $$
(6)
– see for explanations of how $`𝐃`$ can be built-up. For $`\lambda _v=0`$ in (5), conservation of $`I`$ corresponds to the cross-helicity invariance, which is well known contrary to more general quantity $`I`$. Note that taking into account (5), (6), the second term under the integral in Eq. (4) can also be written as
$$2𝝃\rho (𝐕\mathbf{})𝝃_N,$$
or as
$$\rho 𝝃_N(𝐕\mathbf{})𝝃+\rho 𝐕(𝝃_N\mathbf{})𝝃,$$
for absolutely arbitrary functions $`\lambda _{u,v}(\psi ).`$
It is very important that Eq. (1) allows for some extra set of invariants different from (2) and (4). Indeed, differentiating Eq. (1) with respect to time, then multiplying it by $`\ddot{𝝃}`$ and integrating, we found – like in the case of energy but using once again original equation of motion (1) – that the following quantity,
$$E_2=\frac{1}{2}\left\{\frac{1}{\rho }\left(F(𝝃)2\rho 𝐕\mathbf{}\dot{𝝃}\right)^2\dot{𝝃}𝐅(\dot{𝝃})\right\}d^3r,$$
(7)
is conserved. This invariant is exact for linearized dynamics (1), and cannot be reduced to the conservation of energy (2). In principle, we may continue the procedure and get in the same manner an infinite set of similar invariants. However, to investigate a stability, it might be sufficient to involve into our analysis only finite number of the invariants. Here we show that taking into account even one of them, $`E_2`$, it appears to be possible to improve the stability condition significantly.
Note that being varied separately (i.e., when $`\dot{𝝃}`$ is independent), invariant (7) results in the same Frieman-Rotenberg condition of semi-positive definiteness of quadratic form (3) based on $`𝐅`$-operator. To improve stability condition, we will consider the functional
$$U=E+\mu E_2I$$
(8)
to be varied over $`𝝃`$ and $`\dot{𝝃}`$, subject to the independent conservation of the integrals, $`E_2`$ and $`I`$. Explicitly,
$`U={\displaystyle }`$ $`\{{\displaystyle \frac{\rho \dot{𝝃}^2}{2}}{\displaystyle \frac{𝝃𝐅(𝝃)}{2}}+{\displaystyle \frac{\mu }{2\rho }}(𝐅(𝝃)2\rho (𝐕\mathbf{})\dot{𝝃})^2{\displaystyle \frac{\mu }{2}}\dot{𝝃}𝐅(\dot{𝝃})`$ (9)
$`+(2\rho 𝝃(𝐕\mathbf{})\rho \dot{𝝃})(\lambda _u𝐮+\lambda _v𝐯)\}d^3r,`$
where the constant, $`\mu `$, and 1-D functions, $`\lambda _{u,v}(\psi )`$, play roles of the Lagrangian multipliers; we have to choose them to provide the integrals $`E_2,I`$ be equal to their equilibrium values, i.e., to zero.
Functional $`U`$ is minimized by $`\dot{𝝃}`$:
$$\dot{𝝃}=\underset{𝝃_N}{\underset{}{\lambda _u𝐮+\lambda _v𝐯}}+\mu \frac{𝐅(\dot{𝝃})}{\rho }+2\mu 𝐕\mathbf{}\left(2(𝐕\mathbf{})\dot{𝝃}\frac{𝐅(𝝃)}{\rho }\right).$$
(10)
Putting $`\mu 0`$ in Eqs. (9), (10), we approach to the Ilgisonis-Pastukhov-Hameiri condition . Indeed, we have at $`\mu 0`$:
$$\dot{𝝃}𝝃_N,UU_{IPH}=d^3r(\frac{\rho 𝝃_N^2}{2}\frac{𝝃𝐅(𝝃)}{2}),$$
where $`𝝃_N`$ satisfies (5), and $`\lambda _{u,v}`$:
$$(\rho 𝝃_N^22\rho 𝝃(𝐕\mathbf{})𝝃_N)d^3r=0.$$
In this limit, we have not really used the condition of $`E_2`$-conservation. Note that $`U_{IPH}𝝃𝐅(𝝃)d^3r/2`$, and, therefore, the Ilgisonis-Pastukhov-Hameiri condition, $`U_{IPH}0,`$ is ”softer” than the condition (3) by Frieman and Rotenberg. However, it is still not appropriate for arbitrary flow, in which $`𝐕`$ is not parallel to $`𝐁`$. The sign-indefinite term, $`𝝃𝐅(𝝃),`$ contains the high-order $`𝝃`$-derivatives and, therefore, can always prevail on the positive term $`\rho 𝝃_N^2`$ (see for more details).
Now let us account for small but finite $`\mu `$. Solving Eq. (10) by iterations in $`\mu `$, we found
$$\dot{𝝃}𝝃_N2\mu (𝐕\mathbf{})\ddot{𝝃}_0,$$
(11)
where we used the notation
$$\ddot{𝝃}_0=𝐅(𝝃)/\rho 2(𝐕\mathbf{})𝝃_N.$$
(12)
Stability condition is expressed again by functional $`U`$ depending on $`𝝃`$:
$$U=\left\{\frac{\rho }{2}(𝝃_N2\mu (𝐕\mathbf{})\ddot{𝝃}_0)^2\frac{1}{2}𝝃𝐅(𝝃)\right\}d^3r0.$$
(13)
Here $`𝝃_N`$ and $`\ddot{𝝃}_0`$ are defined by Eqs. (5), (12), and also depend on $`𝝃`$.
Lagrangian multipliers have to be found by substituting Eq. (11) into the conditions
$$E_2,I(\dot{𝝃},𝝃)|_{\dot{𝝃}=𝝃_N2\mu (𝐕\mathbf{})\ddot{𝝃}_0}0.$$
They are:
$$\mu =\frac{1}{8}\frac{\rho \ddot{𝝃}_0^2d^3r}{\rho ((𝐕\mathbf{})\ddot{𝝃}_0)^2d^3r},$$
(14)
$$\lambda _u=\frac{A_vD_uA_0D_v}{A_uA_vA_0^2},\lambda _v=\frac{A_uD_vA_0D_u}{A_uA_vA_0^2}.$$
(15)
Here
$$A_{w=u,v}=<4\mu \rho ((𝐕\mathbf{})𝐰)^2\rho 𝐰^2>,$$
$$A_0=<4\mu \rho ((𝐕\mathbf{})𝐮)((𝐕\mathbf{})𝐯)\rho 𝐮𝐯>,$$
$$D_{w=u,v}=<2(\mu 𝐅(𝝃)\rho 𝝃)(𝐕\mathbf{})𝐰>,$$
and angular brackets mean the averaging over magnetic surface.
Note that the left-hand-side of the condition (13) contains the high-order derivatives of $`𝝃`$ in the first (non-negative) term, hence, the second (sign-indefinite) term is no more critical. It is the main advantage of the condition (13) with respect to previous one, $`U_{IPH}0`$, that it has a sense for arbitrary (not only field-aligned) flow and, therefore, may have a practical merit.
This work was partially supported by the Human Capital Foundation Grant No. 41. |
warning/0506/quant-ph0506230.html | ar5iv | text | # Bell inequalities for three particles
## I Introduction
That no local and realistic theory agrees with all predictions of quantum mechanics was shown by Bell in 1964 bell through the violations of certain constraints. Local realism imposes constraints in the form of Bell inequalities on statistical correlations of measurements on multiparticles. Quantum mechanics predicts violations of such Bell inequalities. The original Bell inequality and the subsequent famous CHSH inequality chsh , the latter being cast into a form more amenable for experimental verification, were formulated for the simplest composite quantum system, namely, a system of two qubits.
Bell inequalities which eliminate local realistic description are of importance not only on fundamental research but on identifying ultimate resources for quantum information processing. It was shown that there is a direct link between the security of quantum cryptography and the violation of Bell inequalities Artur ; Acin2 . Collins et al Gisin found a tight Bell inequality for two arbitrary $`d`$-dimensional systems (or two qudits) in terms of joint probabilities, hereafter we call it as the CGLMP inequality. For $`d=2`$, the CGLMP inequality reduces to the CHSH inequality. Alternatively, for $`N`$ particles of two dimensions (called $`N`$-qubit), it was shown that there exists a general Bell inequality which is a sufficient and necessary condition for $`N`$-body correlations to be describable in a local and realistic theory based on two local settings for each observer ZB . However, Bell inequalities for $`N`$ ($`N>2`$) entangled $`d`$-dimensional ($`d>2`$) quantum systems are not so well formulated as those for two-qudit or $`N`$-qubit. Only recently, the problems have been solved partly in the case of three three-dimensional particles in Acin . The authors developed a coincidence Bell inequality in terms of probabilities. For general $`N`$ ($`N>2`$) entangled $`d`$-dimensional quantum systems, no Bell inequality for either probabilities or correlation functions has been presented until now, although GHZ paradox has been generalized to $`N`$ qudits systems paradox .
Even for $`N`$-qubit systems, there is one problem on Bell inequality remains, that is “do all pure entangled states violate Bell inequalities for correlation functions”? In other words, it is that whether the theorem of Gisin gisin ; popescu can be generalized to $`N`$ qubits or not. It is found that there is a family of pure entangled states of $`N`$ qubits which do not violate all Bell inequalities Zukowski . For three qubits, we have proposed a Bell inequality to solve the problem 3qubit .
In this work, we present the tight Bell inequalities expressed by probabilities for three four- and five-dimensional systems. The tight structure of Bell inequalities for three $`d`$-dimensional systems (qudits) is proposed. Some interesting Bell inequalities for three qubits reduced from those of three qudits are also studied.
## II Bell inequalities for three $`d`$-level systems
Local realism cannot exhibit arbitrary correlations. The constraints that local realistic correlations must obey can be written in the form of Bell inequalities. For a three $`d`$-dimensional system with an arbitrary value of $`d`$, some efforts have been given to develop Bell inequalities recently. The first step came in 1990 with a paper of Mermin Mermin in which he derived a Bell inequality for arbitrary $`N`$-qubit states; quantum mechanics violates this inequality by an amount that grows with $`N`$. This result clearly gives us a first three-qubit Bell inequality in a correlation form,
$`Q_{112}+Q_{121}+Q_{211}Q_{222}2.`$ (1)
It can be expressed in terms of probabilities,
$`P(a_1+b_1+c_2=0)P(a_1+b_1+c_2=1)+P(a_1+b_2+c_1=0)P(a_1+b_2+c_1=1)`$
$`+P(a_2+b_1+c_1=0)P(a_2+b_1+c_1=1)P(a_2+b_2+c_2=0)+P(a_2+b_2+c_2=1)2.`$ (2)
This Bell inequality is maximally violated by the three-qubit GHZ state $`|\psi _2=\frac{1}{\sqrt{2}}(|000+|111)`$. But for the generalized GHZ states $`|\psi _2_{GHZ}=\mathrm{cos}\xi |000+\mathrm{sin}\xi |111`$, there exists one region $`\xi (0,\pi /12]`$ in which the Bell inequality is not violated.
The second step is due to Ref. Acin . The authors developed a three-qutrit Bell inequality involving probabilities which can be given in an alternative form,
$`P(a_1+b_1+c_1=0)P(a_1+b_1+c_1=1)+2P(a_1+b_1+c_1=2)+P(a_1+b_1+c_2=0)`$
$`2P(a_1+b_1+c_2=1)+P(a_1+b_1+c_2=2)+P(a_1+b_2+c_1=0)2P(a_1+b_2+c_1=1)`$
$`+P(a_1+b_2+c_1=2)+P(a_2+b_1+c_1=0)2P(a_2+b_1+c_1=1)+P(a_2+b_1+c_1=2)`$
$`+2P(a_1+b_2+c_2=0)P(a_1+b_2+c_2=1)P(a_1+b_2+c_2=2)+2P(a_2+b_1+c_2=0)`$
$`P(a_2+b_1+c_2=1)P(a_2+b_1+c_2=2)+2P(a_2+b_2+c_1=0)P(a_2+b_2+c_1=1)`$
$`P(a_2+b_2+c_1=2)2P(a_2+b_2+c_2=0)2P(a_2+b_2+c_2=1)+4P(a_2+b_2+c_2=2)6.`$ (3)
The above inequality is maximally violated by the three-qutrit GHZ state $`|\psi _3=\frac{1}{\sqrt{3}}(|000+|111+|222)`$. It is worthy of mentioning that both inequalities (II) and (3) are tight.
Here the third step is coming. Our approach of constructing a new Bell inequality for tripartite four-dimensional systems is based on the Gedanken experiment. There are three separated observers, denoted by A, B, and C hereafter, each can carry out two possible local measurements, $`A_1`$ or $`A_2`$ for A, $`B_1`$ or $`B_2`$ for B and $`C_1`$ or $`C_2`$ for C respectively. Each measurement may have four possible outcomes, labeled by 0, 1, 2 and 3. We denote the observable $`X_i`$ measured by party $`X`$ and the outcome $`x_i`$ with $`X=A,B,C(x=a,b,c)`$. A local realistic theory can be described by $`8\times 56`$ probabilities. Here we denote the joint probability $`P(a_i+b_j+c_k=r)`$ that the measurements $`A_i`$, $`B_j`$ and $`C_k`$ have outcomes that differ, modulo four, by $`r`$:
$`P(a_i+b_j+c_k=r)={\displaystyle \underset{a,b=0,1,2,3}{}}P(a_i=a,b_j=b,c_k=rab).`$ (4)
Some of the local realistic constraints are trivial, such as normalization and the no-signaling conditions which are not violated by quantum predictions. Only the non-trivial inequality, which is not true for quantum mechanics, is of use for checking whether we can describe quantum correlations by a classical model. The new Bell inequality for three four-dimensional systems reads
$`5P(a_1+b_1+c_1=0)+P(a_1+b_1+c_1=1)+3P(a_1+b_1+c_1=2)+P(a_1+b_1+c_1=3)`$
$`+3P(a_1+b_1+c_2=0)7P(a_1+b_1+c_2=1)+3P(a_1+b_1+c_2=2)+P(a_1+b_1+c_2=3)`$
$`+3P(a_1+b_2+c_1=0)7P(a_1+b_2+c_1=1)+3P(a_1+b_2+c_1=2)+P(a_1+b_2+c_1=3)`$
$`+3P(a_2+b_1+c_1=0)7P(a_2+b_1+c_1=1)+3P(a_2+b_1+c_1=2)+P(a_2+b_1+c_1=3)`$
$`+3P(a_1+b_2+c_2=0)+P(a_1+b_2+c_2=1)5P(a_1+b_2+c_2=2)+P(a_1+b_2+c_2=3)`$
$`+3P(a_2+b_1+c_2=0)+P(a_2+b_1+c_2=1)5P(a_2+b_1+c_2=2)+P(a_2+b_1+c_2=3)`$
$`+3P(a_2+b_2+c_1=0)+P(a_2+b_2+c_1=1)5P(a_2+b_2+c_1=2)+P(a_2+b_2+c_1=3)`$
$`P(a_2+b_2+c_2=0)3P(a_2+b_2+c_2=1)P(a_2+b_2+c_2=2)+5P(a_2+b_2+c_2=3))12.`$ (5)
That the maximum value of the left hand side of inequality (5) for local theories is 12 can be given in the following sense. By using $`_{r=0}^3P(a_i+b_j+c_k=r)=1`$, the inequality (5) is reformed as
$`3P(a_1+b_1+c_1=0)+P(a_1+b_1+c_1=2)5P(a_1+b_1+c_2=1)`$
$`P(a_1+b_1+c_2=3)5P(a_1+b_2+c_1=1)P(a_1+b_2+c_1=3)`$
$`5P(a_2+b_1+c_1=1)P(a_2+b_1+c_1=3)+P(a_1+b_2+c_2=0)`$
$`3P(a_1+b_2+c_2=2)+P(a_2+b_1+c_2=0)3P(a_2+b_1+c_2=2)`$
$`+P(a_2+b_2+c_1=0)3P(a_2+b_2+c_1=2)P(a_2+b_2+c_2=1)`$
$`+3P(a_2+b_2+c_2=3))0.`$ (6)
To beat the bound $`0`$, terms $`P(a_1+b_1+c_1=2)`$ and $`P(a_2+b_2+c_2=3)`$ are taken equal to one first. This means that $`a_1+b_1+c_1+a_2+b_2+c_2=5`$. Among the remained terms, we take $`P(a_1+b_1+c_2=3)`$, $`P(a_1+b_2+c_1=3)`$ and $`P(a_2+b_1+c_1=3)`$ equal to one to maximize the value of left hand side of the inequality (6). As a result, $`a_2+b_2+c_1=2`$, $`a_2+b_1+c_2=2`$ and $`a_1+b_2+c_2=2`$ according to the constraint $`a_1+b_1+c_1+a_2+b_2+c_2=5`$. So $`P(a_1+b_1+c_2=1)=P(a_1+b_2+c_1=1)=P(a_2+b_1+c_1=1)=0`$, $`P(a_2+b_2+c_1=2)=P(a_2+b_1+c_2=2)=P(a_1+b_2+c_2=2)=1`$ and $`P(a_2+b_2+c_1=0)=P(a_2+b_1+c_2=0)=P(a_1+b_2+c_2=0)=0`$. Therefore we have $`0+1010101+03+03+030+3=80`$. If initially terms $`P(a_1+b_1+c_1=0)`$ and $`P(a_2+b_2+c_2=3)`$ are taken equal to one first. This means that $`a_1+b_1+c_1+a_2+b_2+c_2=3`$. Among the remained terms, we take $`P(a_1+b_1+c_2=3)`$, $`P(a_1+b_2+c_1=3)`$ and $`P(a_2+b_1+c_1=3)`$ equal to one to maximize the value of left hand side of the inequality (6). As a result, $`a_2+b_2+c_1=0`$, $`a_2+b_1+c_2=0`$ and $`a_1+b_2+c_2=0`$ according to the constraint $`a_1+b_1+c_1+a_2+b_2+c_2=3`$. So $`P(a_1+b_1+c_2=1)=P(a_1+b_2+c_1=1)=P(a_2+b_1+c_1=1)=0`$, $`P(a_2+b_2+c_1=2)=P(a_2+b_1+c_2=2)=P(a_1+b_2+c_2=2)=0`$ and $`P(a_2+b_2+c_1=0)=P(a_2+b_1+c_2=0)=P(a_1+b_2+c_2=0)=1`$. Therefore we have $`3+0010101+10+10+100+3=00`$. Therefore, after some simple and patient calculations, it can be shown that the inequality (6) is always bounded by $`0`$ in a local realistic model. Furthermore, the Bell inequality (5) is a tight inequality for three four-dimensional systems Acin3 .
Let us now consider the maximum value that can be attained for the inequality (5) for quantum measurements on an entangled quantum state. First, we specify the quantum state and measurement. The initial state is a natural generalization of bipartite maximally entangled state to three four-level systems,
$`|\psi _4={\displaystyle \frac{1}{2}}(|000+|111+|222+|333).`$ (7)
Consider a Gedanken experiment in which A, B and C measure observables defined by unbiased symmetric multi-port beam splitters Zukowski2 on $`|\psi `$. The unbiased symmetric multi-port beam splitter is an optical device with $`d`$ input and $`d`$ output ports. In front of every input port there is a phase shifter that changes the phase of the photon entering the given port. If a phase shifter in some input port is set to zero and a photon enters the device through this port then it has an equal chance of leaving the device through any output port. The phase shifters can be changed by the observers; they represent the local macroscopic parameters available to the observers. The matrix elements of an unbiased symmetric multi-port beam splitter are given by $`U_{kl}(\stackrel{}{\varphi })=\frac{1}{\sqrt{d}}\alpha ^{kl}\mathrm{exp}(i\varphi ^l)`$, where $`\alpha =\mathrm{exp}(\frac{2i\pi }{d})`$ and $`\varphi ^l(l=0,1,2,\mathrm{},d1)`$ are the settings of the appropriate phase shifters, for convenience we denote them as a $`d`$ dimensional vector $`\stackrel{}{\varphi }=(\varphi ^0,\varphi ^1,\varphi ^2,\mathrm{},\varphi ^{d1})`$. For four-dimensional systems, $`d=4`$.
The quantum prediction for the probabilities of obtaining the outcome $`(a,b,c)`$ is then given as
$`P(a_i=a,b_j=b,c_k=c)=|abc|U(\stackrel{}{\varphi _A})U(\stackrel{}{\varphi _B})U(\stackrel{}{\varphi _C})|\psi _4|^2.`$ (8)
Thus the quantum analogue of the joint probability can be easily calculated
$`P(a_i+b_j+c_k=r)`$
$`={\displaystyle \frac{1}{16}}(4+2\mathrm{cos}(\phi ^1\phi ^0+{\displaystyle \frac{\pi }{2}}r)+2\mathrm{cos}(\phi ^2\phi ^0+\pi r+2\mathrm{cos}(\phi ^2\phi ^1+{\displaystyle \frac{\pi }{2}}r)`$
$`+2\mathrm{cos}(\phi ^3\phi ^0+{\displaystyle \frac{3\pi }{2}}r)+2\mathrm{cos}(\phi ^3\phi ^1+\pi r)+2\mathrm{cos}(\phi ^3\phi ^2+{\displaystyle \frac{\pi }{2}}r)),`$ (9)
where $`\phi ^i=\varphi _A^i+\varphi _B^i+\varphi _C^i,(i=0,1,2,3)`$. In order to look for the maximal violation of the inequality, we choose the optimal settings as the following: $`\stackrel{}{\varphi }_{A1}=\stackrel{}{\varphi }_{B1}=\stackrel{}{\varphi }_{C1}=(0,\frac{1}{3}\mathrm{arccos}(\frac{1}{3}),\frac{1}{3}\mathrm{arccos}(\frac{1}{3})\frac{\pi }{3},\frac{\pi }{3})`$, $`\stackrel{}{\varphi }_{A2}=\stackrel{}{\varphi }_{B2}=\stackrel{}{\varphi }_{C2}=(0,\frac{1}{3}\mathrm{arcsin}\frac{7}{9},\frac{1}{3}\mathrm{arcsin}\frac{7}{9}+\frac{\pi }{6},\frac{\pi }{6})`$. Numerical results show that for this choice, all the probability terms have definite values as listed in Table 1.
Putting them into the left hand side of the inequality (5), we arrive at $`2\frac{1}{6}+3\frac{2}{3}+3(3\frac{1}{2}+3\frac{1}{2})+3(3\frac{2}{3}+2\frac{1}{6})2\frac{1}{18}+5\frac{8}{9}=\frac{68}{3}>12`$.
In Ref. v , a proposal was made to measure the strength of violation of local realism by the minimal amount of noise that must be added to the system in order to hide the non-classical character of the observed correlations. This is equivalent to a replacement of the pure state $`|\psi \psi |`$ by the mixed state $`\rho (F)`$ of the form $`\rho (F)=(1F)|\psi \psi |+\frac{F}{56}III`$, where $`I`$ is an identity matrix and $`F(0F1)`$ is the amount of noise present in the system. For $`F=0`$, local realistic description does not exist, whereas it does for $`F=1`$. Therefore, there exists some threshold value of $`F`$, denoted by $`F_{thr}`$, such that for every $`FF_{thr}`$, local and realistic description does not exist. The threshold fidelity for the three four-level systems is determined by $`(1F_{thr})\frac{68}{3}=12`$, namely $`F_{thr}=\frac{8}{17}=0.4706`$.
Similarly, we propose a Bell inequality for three five-dimensional systems based on the Gedanken experiment:
$`2P(a_1+b_1+c_1=0)+P(a_1+b_1+c_1=1)+P(a_1+b_1+c_1=4)`$
$`+P(a_1+b_1+c_2=0)2P(a_1+b_1+c_2=2)+P(a_1+b_1+c_2=4)`$
$`+P(a_1+b_2+c_1=0)2P(a_1+b_2+c_1=2)+P(a_1+b_2+c_1=4)`$
$`+P(a_2+b_1+c_1=0)2P(a_2+b_1+c_1=2)+P(a_2+b_1+c_1=4)`$
$`+P(a_1+b_2+c_2=0)+P(a_1+b_2+c_2=3)2P(a_1+b_2+c_2=4)`$
$`+P(a_2+b_1+c_2=0)+P(a_2+b_1+c_2=3)2P(a_2+b_1+c_2=4)`$
$`+P(a_2+b_2+c_1=0)+P(a_2+b_2+c_1=3)2P(a_2+b_2+c_1=4)`$
$`2P(a_2+b_2+c_2=1)+P(a_2+b_2+c_2=3)+P(a_2+b_2+c_2=4)4,`$ (10)
which is satisfied by the local and realistic theories. By the way, the Bell inequality (10) is also tight Acin3 .
Using specified quantum state and measurement, we calculate the maximum value that can be attained for the inequality (10). The considered state is a natural generalization of bipartite maximally entangled state to three five-level systems,
$`|\psi _5={\displaystyle \frac{1}{\sqrt{5}}}(|000+|111+|222+|333+|444).`$ (11)
The measurement is also based on unbiased symmetric multi-port beam splitter with $`d=5`$ and the quantum prediction for the probabilities of obtaining the outcome $`(a,b,c)`$ is then given as
$`P(a_i=a,b_j=b,c_k=c)=|abc|U(\stackrel{}{\varphi _A})U(\stackrel{}{\varphi _B})U(\stackrel{}{\varphi _C})|\psi _5|^2.`$ (12)
Thus the quantum analogue of the joint probability is given as
$`P(a_i+b_j+c_k=r)`$
$`={\displaystyle \frac{1}{25}}(5+2\mathrm{cos}(\phi ^1\phi ^0+{\displaystyle \frac{2\pi }{5}}r)+2\mathrm{cos}(\phi ^2\phi ^0+{\displaystyle \frac{4\pi }{5}}r)+2\mathrm{cos}(\phi ^2\phi ^1+{\displaystyle \frac{2\pi }{5}}r)`$
$`+2\mathrm{cos}(\phi ^3\phi ^0+{\displaystyle \frac{6\pi }{5}}r)+2\mathrm{cos}(\phi ^3\phi ^1+{\displaystyle \frac{4\pi }{5}}r)+2\mathrm{cos}(\phi ^3\phi ^2+{\displaystyle \frac{2\pi }{5}}r)`$
$`+2\mathrm{cos}(\phi ^4\phi ^0+{\displaystyle \frac{8\pi }{5}}r)+2\mathrm{cos}(\phi ^4\phi ^1+{\displaystyle \frac{6\pi }{5}}r)+2\mathrm{cos}(\phi ^4\phi ^2+{\displaystyle \frac{4\pi }{5}}r)+2\mathrm{cos}(\phi ^4\phi ^3+{\displaystyle \frac{2\pi }{5}}r)),`$ (13)
where $`\phi ^i=\varphi _A^i+\varphi _B^i+\varphi _C^i(i=0,1,2,3,4)`$. Numerical results show that for the following angle settings $`\stackrel{}{\varphi }_{A1}=\stackrel{}{\varphi }_{B1}=\stackrel{}{\varphi }_{C1}=(0,\beta _1,\beta _2,\beta _2,\beta _1)`$, $`\stackrel{}{\varphi }_{A2}=\stackrel{}{\varphi }_{B2}=\stackrel{}{\varphi }_{C2}=(0,\beta _1+\frac{\pi }{5},\beta _2+\frac{2\pi }{5},\beta _2\frac{2\pi }{5},\beta _1\frac{\pi }{5})`$, where $`\mathrm{cos}(3\beta _1)\mathrm{cos}(3\beta _2)=\frac{1}{2}`$, the maximum value of the left hand side of inequality (10) attained is $`6.72216`$. The threshold fidelity for the three five-level systems is determined by $`(1F_{thr})\times 6.72216=4`$, namely $`F_{thr}=0.40495`$.
Given the above known Bell inequalities for three particles with $`d=2,3,4,5`$, it is worthy of noting that such tight inequalities exhibit perfect symmetries. One symmetry is that these Bell inequalities are symmetric under the permutations of subsystems $`A,B,`$ and $`C`$. The second symmetry is that they can be expressed in a general form. Based on which, the structure of tight Bell inequalities is suggested as
$`{\displaystyle \frac{1}{d(d1)}}{\displaystyle \underset{ijk}{}}{\displaystyle \underset{r=0}{\overset{d1}{}}}f_{d;ijk}^rP(a_i+b_j+c_k=r)1,`$ (14)
with $`1\frac{f_{d;ijk}^r}{d(d1)}1`$. It is worthy of noting that the tight Bell inequalities for two qudits (i.e., the CGLMP inequality, see Gisin ) can be written in a similar form,
$`{\displaystyle \frac{1}{d1}}{\displaystyle \underset{ij}{}}{\displaystyle \underset{r=0}{\overset{d1}{}}}f_{d;ij}^rP(a_i+b_j=r)1,`$ (15)
with $`1\frac{f_{d;ij}^r}{d1}1`$. The coefficients $`f_{d;ijk}^r`$ and $`f_{d;ij}^r`$ are integers or half integers. Note that these inequalities are fulfilled with displacement of probabilities. In other words, for a known number $`m`$, where $`md1`$, the above inequalities are still true for local realistic description by replacing $`P(a_i+b_j+c_k=r)`$ with $`P(a_i+b_j+c_k=r+m)`$, which is the third symmetry of the set of Bell inequalities for three particles.
## III Interesting Bell inequalities of three-qubit reduced from those of three-qudit
In 1991 Gisin presented a theorem, which states that any pure entangled state of two particles violates a Bell inequality for two-particle correlation functions gisin ; popescu . Recent investigations show a surprising result that there exists a family of pure entangled $`N`$-qubit states that does not violate any Bell inequality for $`N`$-particle correlations for the case of a standard Bell experiment on $`N`$ qubits scarani . This family is the generalized GHZ states given by
$`|\psi _{GHZ}=\mathrm{cos}\xi |0\mathrm{}0+\mathrm{sin}\xi |1\mathrm{}1,`$ (16)
with $`0\xi \pi /4`$. The usual GHZ states ghz are for $`\xi =\pi /4`$. For a three-qubit system, whose corresponding generalized GHZ state reads $`|\psi _{GHZ}=\mathrm{cos}\xi |000+\mathrm{sin}\xi |111`$, it has been shown that for the region $`\xi (0,\pi /12]`$, the inequalities given in ZB are not violated based on the standard Bell experiment. Recently, we developed a three-qubit Bell inequality which is a solution to such a problem. That is all pure entangled states of three qubits violates the Bell inequality given in 3qubit . Indeed Bell inequalities are sensitive to the presence of noise and above a certain amount of noise the Bell inequalities will cease to be violated by a quantum system. However, it seems that the inequality in Ref. 3qubit is not good enough to the resistance of noise. For the three-qubit GHZ state, the threshold visibility is $`V_{GHZ}=4\sqrt{3}/9=0.7698`$ and for the W state, the threshold visibility is $`V_W=0.7312`$. Tne inequality in Ref. 3qubit can be derived from three-qutrit Bell inequality (3). Actually any Bell inequality for tripartite $`d(d>2)`$-level systems may reduce to a Bell inequality for three qubits when one considers only two outcomes of measurement.
Here we present a new Bell inequality for three-qubit systems which is reduced from inequality (5)
$`3P(a_1+b_1+c_1=0)+P(a_1+b_1+c_1=1)5P(a_1+b_1+c_1=2)+P(a_1+b_1+c_1=3)`$
$`+3P(a_1+b_1+c_2=0)+P(a_1+b_1+c_2=1)+3P(a_1+b_1+c_2=2)7P(a_1+b_1+c_2=3)`$
$`+3P(a_1+b_2+c_1=0)+P(a_1+b_2+c_1=1)+3P(a_1+b_2+c_1=2)7P(a_1+b_2+c_1=3)`$
$`+3P(a_2+b_1+c_1=0)+P(a_2+b_1+c_1=1)+3P(a_2+b_1+c_1=2)7P(a_2+b_1+c_1=3)`$
$`5P(a_1+b_2+c_2=0)+P(a_1+b_2+c_2=1)+3P(a_1+b_2+c_2=2)+P(a_1+b_2+c_2=3)`$
$`5P(a_2+b_1+c_2=0)+P(a_2+b_1+c_2=1)+3P(a_2+b_1+c_2=2)+P(a_2+b_1+c_2=3)`$
$`5P(a_2+b_2+c_1=0)+P(a_2+b_2+c_1=1)+3P(a_2+b_2+c_1=2)+P(a_2+b_2+c_1=3)`$
$`P(a_2+b_2+c_2=0)+5P(a_2+b_2+c_2=1)P(a_2+b_2+c_2=2)3P(a_2+b_2+c_2=3))12.`$ (17)
The inequality can be expressed in terms of correlation functions
$`E(A_1B_1C_1)+E(A_1B_1C_2)+E(A_1B_2C_1)+E(A_2B_1C_1)E(A_2B_2C_2)`$
$`E(A_1B_2)E(A_2B_1)E(A_2B_2)E(A_1C_2)E(A_2C_1)E(A_2C_2)`$
$`E(B_1C_2)E(B_2C_1)E(B_2C_2)+E(A_1)+E(B_1)+E(C_1)3.`$ (18)
The above inequality (18) includes the terms of single correlation functions, it is symmetric under the permutations of $`A_j,B_j`$ and $`C_j`$. Quantum mechanically, the above inequality is violated by all pure entangled states of three qubits. Pure states of three qubits constitute a five-parameter family, with equivalence up to local unitary transformations. This family has the following representation Acin00
$`|\psi `$ $`=`$ $`\sqrt{\mu _0}|000+\sqrt{\mu _1}e^{i\varphi }|100+\sqrt{\mu _2}|101`$ (19)
$`+\sqrt{\mu _3}|110+\sqrt{\mu _4}|111,`$
with $`\mu _i0`$, $`_i\mu _i=1`$, and $`0\varphi \pi `$. Numerical results show that this Bell inequality for probabilities is violated by all pure entangled states of three-qubit systems. However, no analytical proof of the conclusion has been given. In the following, some special cases will be given to show the inequality (18) is violated by all pure entangled states. The first quantum state considered is generalized GHZ state $`|\psi _2_{\mathrm{GHZ}}=\mathrm{cos}\xi |000+\mathrm{sin}\xi |111`$. The inequality (18) is violated by the generalized GHZ states for the whole region except $`\xi =0,\pi /2`$. For the GHZ state with $`\xi =\pi /4`$, the quantum violation reaches its maximum value $`4.40367`$. The variation of the violation with $`\xi `$ is shown in Fig. 1. Another state considered is generalized W state $`|\psi _W=\mathrm{sin}\beta \mathrm{cos}\xi |100+\mathrm{sin}\beta \mathrm{sin}\xi |010+\mathrm{cos}\beta |001`$. By fixing the value of $`\beta `$, quantum violation of the inequality (18) varies with $`\xi `$ (see Fig. 2). The inequality (18) is violated by the generalized W states except the cases with $`\beta =\frac{\pi }{2}`$, $`\xi =0`$ and $`\xi =\frac{\pi }{2}`$. The states in these cases are direct-product states which do not violated any Bell inequality. For the standard W state, quantum violation of the inequality (18) approaches $`4.54086`$.
Hence the inequality (18) is also one candidate to generalize the theorem of Gisin to three-qubit systems. One of the features of the new inequality for three qubits is that it is highly resistant to noise. The inequality (18) is violated by the generalized GHZ state $`|\psi _2_{GHZ}=\mathrm{cos}\xi |000+\mathrm{sin}\xi |111`$ for the whole region, the threshold visibility is $`V_{thr}^{GHZ}=0.68125`$. The inequality (18) is also violated by the W state, the threshold visibility is $`V_{thr}^W=0.660668`$. We plot the variation of quantum violation for the generalized GHZ states with angle $`\xi `$ for inequality (18) and inequality given in Ref. 3qubit , see Fig. 3. In plotting the figure, we reform the expressions of these two inequalities as
$`{\displaystyle \frac{1}{4}}`$ $`[`$ $`Q(A_1B_1C_1)Q(A_1B_2C_2)Q(A_2B_1C_2)Q(A_2B_2C_1)+2Q(A_2B_2C_2)`$ (20)
$`Q(A_1B_1)Q(A_1B_2)Q(A_2B_1)Q(A_2B_2)+Q(A_1C_1)+Q(A_1C_2)`$
$`+Q(A_2C_1)+Q(A_2C_2)+Q(B_1C_1)+Q(B_1C_2)+Q(B_2C_1)+Q(B_2C_2)]1,`$
$`{\displaystyle \frac{1}{3}}`$ $`[`$ $`Q(A_1B_1C_1)+Q(A_1B_1C_2)+Q(A_1B_2C_1)+Q(A_2B_1C_1)Q(A_2B_2C_2)`$ (21)
$`Q(A_1B_2)Q(A_2B_1)Q(A_2B_2)Q(A_1C_2)Q(A_2C_1)Q(A_2C_2)`$
$`Q(B_1C_2)Q(B_2C_1)Q(B_2C_2)+Q(A_1)+Q(B_1)+Q(C_1)]1,`$
respectively. By the reformation, the violation degrees of the two inequalities can be compared directly. Comparing the results of the inequality given in Ref. 3qubit , the new inequality (18) is indeed more resistant to noise. It seems that we could derive some new three-qubit Bell inequalities, which would be more highly resistance to noise, if we get other Bell inequalities of tripartite $`d`$-dimensional ($`d>4`$) quantum systems.
When setting $`C_1=1,C_2=1`$, the inequality (18) reduces directly to the CHSH inequality for two-qubit
$$E(A_1B_1)E(A_1B_2)E(A_2B_1)E(A_2B_2)2,$$
which is equivalent to
$$E(A_1B_1)+E(A_1B_2)+E(A_2B_1)E(A_2B_2)2.$$
Starting from the Bell inequality for three five-level systems, another three-qubit Bell inequality can be obtained as
$`P(a_1+b_1+c_1=1)2P(a_1+b_1+c_1=2)+P(a_1+b_1+c_1=3)+P(a_1+b_1+c_2=1)+P(a_1+b_1+c_2=2)`$
$`+P(a_1+b_2+c_1=1)+P(a_1+b_2+c_1=2)+P(a_2+b_1+c_1=1)+P(a_2+b_1+c_1=2)`$
$`+P(a_1+b_2+c_2=0)2P(a_1+b_2+c_2=1)+P(a_1+b_2+c_2=2)+P(a_2+b_1+c_2=0)`$
$`2P(a_2+b_1+c_2=1)+P(a_2+b_1+c_2=2)+P(a_2+b_2+c_1=0)2P(a_2+b_2+c_1=1)`$
$`+P(a_2+b_2+c_1=2)+P(a_2+b_2+c_2=0)+P(a_2+b_2+c_2=1)2P(a_2+b_2+c_2=3)4.`$ (22)
The above inequality is just the Mermin inequality when it is expressed in terms of correlation functions:
$`Q_{111}+Q_{122}+Q_{212}+Q_{221}2.`$ (23)
## IV Summary
To summarize, we have presented the tight Bell inequalities expressed by probabilities for three four- and five-dimensional systems. The tight structure of Bell inequalities for three $`d`$-dimensional systems (qudits) is also proposed. Moreover, a new Bell inequality for three qubits is derived (or say, reduced) from the inequality for three four-level systems. The inequality (18) is violated by any pure three-qubit entangled state, and it is more resistant to noise compared with the one given in Ref. 3qubit . Furthermore, the tight Mermin inequality for thre-qubit can be obtained by reducing the Bell inequality for three five-level systems \[see Eq. (10) and Eq. (22)\].
ACKNOWLEDGMENTS The authors thank A. Acin for many useful discussions. This work was supported by NUS research Grant No. WBS: R-144-000-123-112. J.L.C acknowledges financial supports from NSF of China (Grant No. 10605013), Program for New Century Excellent Talents in University, and the Project-sponsored by SRF for ROCS, SEM. C.F.W. acknowledges financial support from Singapore Millennium Foundation. |
warning/0506/cs0506041.html | ar5iv | text | # Competitive on-line learning with a convex loss function
## 1 Introduction
In the simple problem of sequential decision making that we consider in this paper, the loss $`\lambda (y_n,\gamma _n)`$ (maybe negative) caused by a decision $`\gamma _n`$ depends only on the following binary observation $`y_n`$. All relevant information available to the decision maker by the time he makes his decision is collected in what we call the datum, $`x_n`$. For example, in time series applications the datum may contain all or the most recent observations; in pattern recognition, where the observations are the true classes of patterns, the datum may be the vector of a pattern’s attributes.
The traditional approach to this problem assumes a statistical model for the sequence of pairs $`(x_n,y_n)`$; e.g., statistical learning theory () assumes that the $`(x_n,y_n)`$ are generated independently from the same probability distribution. A more recent approach, known in learning theory as “prediction with expert advice” (e.g., ) and in information theory as “universal prediction” (e.g., ), avoids making assumptions about the way the observations and data are generated. Instead, the goal of the decision maker is to compete with a more or less general benchmark class of decision rules, mapping the $`x`$s to the $`y`$s (the framework of prediction with expert advice is usually even more general). We will use the phrase “competitive on-line” to refer to this area (as in , emphasizing similarities to competitive on-line algorithms in computation theory).
First papers on competitive on-line learning with general loss functions (e.g., ) dealt with countable (often finite) benchmark classes. The next step was to consider finite-dimensional benchmark classes (e.g., ). This paper continues with infinite-dimensional classes. (Such classes were considered earlier by Kimber and Long , who, however, assumed that the benchmark class contains a perfect decision rule.) To get an idea of our central results, the reader is advised to start from Corollaries 13.
Our implicit assumption, common with other work in competitive on-line learning, is that the decision maker is “small”: his decisions do not affect the future observations. This is not a mathematical assumption: as already mentioned, we do not make any assumptions at all about the way observations are generated; however, interpretation of our results becomes problematic if the decision maker is not small. “Big” decision makers can still use our algorithms for prediction (cf. Remarks 1 and 3 below).
In conclusion of this section we will briefly describe the content of the paper. Our main result is stated in §2, and several examples are given in §3. It is proved in §6; in §4 we describe the main ideas behind the proof and in §5 we prove some preparatory results for §6. Our decision algorithm is explicitly described in §7. We conclude with a short list of directions of further research (§8). A preliminary version of this paper is to appear as . In this new version we made the title of the paper more specific (the old title was even somewhat misleading: in prediction with expert advice, the experts are usually completely free in making their decisions).
## 2 Main result
Our decision protocol is:
FOR $`n=1,2,\mathrm{}`$:
Reality announces $`x_n𝐗`$.
Decision Maker announces $`\gamma _n\mathrm{\Gamma }`$.
Reality announces $`y_n\{0,1\}`$.
END FOR.
At each step (or *round*) $`n`$ Decision Maker makes a decision $`\gamma _n`$ whose consequences depend on the *observation* $`y_n\{0,1\}`$ chosen by Reality. All relevant information available to Decision Maker by the time he makes his decision is collected in $`x_n`$, called the *datum*. We assume that the data $`x_n`$ are elements of a *data space* $`𝐗`$ and that the decisions are elements of a *decision space* $`\mathrm{\Gamma }`$ (both sets assumed non-empty).
###### Remark 1
In this paper we are interested, first of all, in prediction of future observations. However, our framework allows a fairly wide class of loss functions, not all of which can be interpreted in terms of predictions (such as, e.g., Cover’s and long-short games, in the terminology of , §2). This is the main reason why we prefer to talk about decision making in general; another reason is that in §5 we will deal with a very different kind of prediction (for which we reserve the term “forecasting”).
A *decision strategy* is a strategy for Decision Maker in this protocol (explicitly defined specific strategies will also be called “decision algorithms”). Its performance is measured with a *loss function* $`\lambda :\{0,1\}\times \mathrm{\Gamma }`$, and so its cumulative loss over the first $`N`$ rounds is
$$\underset{n=1}{\overset{N}{}}\lambda (y_n,\gamma _n).$$
The pair $`(\mathrm{\Gamma },\lambda )`$ is the *game* being played. Decision Maker will compete against a class $``$, called the *benchmark class*, of functions $`D:𝐗\mathrm{\Gamma }`$ considered as decision rules; the cumulative loss suffered by such a decision rule is
$$\underset{n=1}{\overset{N}{}}\lambda (y_n,D(x_n)).$$
Before stating our main result we define some useful notions connected with the two main components of our decision framework, the game $`(\mathrm{\Gamma },\lambda )`$ and the benchmark class $``$. The reader might want in parallel to read the next section, which describes some important examples of games and benchmark classes.
### Games
The *exposure* $`Exp_\lambda (\gamma )`$ of a decision $`\gamma \mathrm{\Gamma }`$ is
$$Exp_\lambda (\gamma ):=\lambda (1,\gamma )\lambda (0,\gamma )$$
and the *exposure* $`Exp_{\lambda ,D}:𝐗`$ of a decision rule $`D`$ at a point $`x𝐗`$ is
$$Exp_{\lambda ,D}(x):=\lambda (1,D(x))\lambda (0,D(x)).$$
Let $`\lambda (p,\gamma )`$ be the expected loss caused by taking a decision $`\gamma `$ when the probability of $`1`$ is $`p`$:
$$\lambda (p,\gamma ):=p\lambda (1,\gamma )+(1p)\lambda (0,\gamma ).$$
(1)
We only consider games $`(\mathrm{\Gamma },\lambda )`$ such that
$$C_0:=\underset{\gamma \mathrm{\Gamma }}{inf}\lambda (0,\gamma ),C_1:=\underset{\gamma \mathrm{\Gamma }}{inf}\lambda (1,\gamma )$$
(2)
are finite. It is convenient (see, e.g., ) to summarize a game by its *superdecision set*
$$\mathrm{\Sigma }:=\{(x,y)^2|\gamma \mathrm{\Gamma }:x\lambda (0,\gamma )\text{ and }y\lambda (1,\gamma )\};$$
(3)
elements of this set will be called *superdecisions*. Superdecisions of the form $`(\lambda (0,\gamma ),\lambda (1,\gamma ))`$ will sometimes be called *decisions*. We will assume, additionally, that the set $`\mathrm{\Sigma }^2`$ is convex and closed. The *Eastern tail* of the game is the function
$$\begin{array}{cc}\hfill f:[C_0,\mathrm{})& \{\mathrm{}\}\hfill \\ \hfill x& inf\{y|(x,y)\mathrm{\Sigma }\}C_1\hfill \end{array}$$
(4)
and its *Northern tail* is
$$\begin{array}{cc}\hfill g:[C_1,\mathrm{})& \{\mathrm{}\}\hfill \\ \hfill y& inf\{x|(x,y)\mathrm{\Sigma }\}C_0,\hfill \end{array}$$
(5)
where, as usual, $`inf\mathrm{}:=\mathrm{}`$; it is clear that $`f`$ and $`g`$ are nonnegative everywhere and finite on $`(C_0,\mathrm{})`$ and $`(C_1,\mathrm{})`$, respectively.
### The theorem
A *reproducing kernel Hilbert space* (RKHS) on $`𝐗`$ is a Hilbert space $``$ of real-valued functions on $`𝐗`$ such that the evaluation functional $`ff(x)`$ is continuous for each $`x𝐗`$. By the Riesz–Fischer theorem, for each $`x𝐗`$ there exists a function $`𝐊_x`$ such that
$$f(x)=𝐊_x,f_{},f.$$
Let
$$𝐜_{}:=\underset{x𝐗}{sup}𝐊_x_{};$$
(6)
we will be interested in the case $`𝐜_{}<\mathrm{}`$. With each game $`(\mathrm{\Gamma },\lambda )`$ and each RKHS $``$ we associate the non-negative (but maybe infinite) constant $`𝐜_{\lambda ,}`$ defined by
$`𝐜_{\lambda ,}^2`$ $`:=\underset{p(0,1)}{sup}\underset{\gamma \mathrm{\Gamma }_p}{sup}\underset{x𝐗}{sup}p(1p)\left(Exp_\lambda ^2(\gamma )+𝐊_x_{}^2\right)`$ (7)
$`=\underset{p(0,1)}{sup}\underset{\gamma \mathrm{\Gamma }_p}{sup}p(1p)\left(Exp_\lambda ^2(\gamma )+𝐜_{}^2\right),`$
where $`\mathrm{\Gamma }_p:=\mathrm{arg}\mathrm{min}_{\gamma \mathrm{\Gamma }}\lambda (p,\gamma )`$ (and $`\lambda (p,\gamma )`$ is defined by (1)).
The following is our main result.
###### Theorem 1
Let the game $`(\mathrm{\Gamma },\lambda )`$ be such that (2) are finite, the superdecision set $`\mathrm{\Sigma }`$ is convex and closed, and the tails $`f`$ and $`g`$ satisfy
$$f_+^{}(t)=O(t^2),g_+^{}(t)=O(t^2)$$
(8)
as $`t\mathrm{}`$, where $`f_+^{}`$ and $`g_+^{}`$ stand for the right derivatives (see, e.g., , §23) of $`f`$ and $`g`$. Let $``$ be an RKHS on $`𝐗`$ and $`𝐜_{}`$, $`𝐜_{\lambda ,}`$ be defined by (6) and (7). Suppose $`𝐜_{}<\mathrm{}`$. Then $`𝐜_{\lambda ,}<\mathrm{}`$ and there is a decision strategy which guarantees that
$$\underset{n=1}{\overset{N}{}}\lambda (y_n,\gamma _n)\underset{n=1}{\overset{N}{}}\lambda (y_n,D(x_n))+𝐜_{\lambda ,}\left(Exp_{\lambda ,D}_{}+1\right)\sqrt{N}$$
(9)
for all $`N=1,2,\mathrm{}`$ and all $`D:𝐗\mathrm{\Gamma }`$ with $`Exp_{\lambda ,D}`$.
###### Remark 2
If the loss function $`\lambda `$ is bounded, (8) holds trivially. The right derivatives in (8) can be replaced by the corresponding left derivatives, since $`\left|f_+^{}\right|\left|f_{}^{}\right|`$ and $`\left|g_+^{}\right|\left|g_{}^{}\right|`$ (see, e.g., , Theorem 24.1). Condition (8) can be interpreted as saying that the tails should shrink fast enough. The case $`f(t)=g(t)=t^1`$ can be considered borderline; Theorem 1 is still applicable in this case, but it ceases to be applicable for tails that shrink less fast.
## 3 Examples
In this section we first define a specific RKHS and then describe three important games.
### Kernels as source of RKHS
We start by describing an equivalent language for talking about RKHS. The *kernel* of an RKHS $``$ on $`𝐗`$ is
$$𝐊(x,x^{}):=𝐊_x,𝐊_x^{}_{}$$
(equivalently, we could define $`𝐊(x,x^{})`$ as $`𝐊_x(x^{})`$ or as $`𝐊_x^{}(x)`$). There is a simple internal characterization of the kernels $`𝐊`$ of RKHS.
It is easy to check that the function $`𝐊(x,x^{})`$, as we defined it, is symmetric ($`𝐊(x,x^{})=𝐊(x^{},x)`$ for all $`x,x^{}𝐗`$) and positive definite ($`_{i=1}^m_{j=1}^m\alpha _i\alpha _j𝐊(x_i,x_j)0`$ for all $`m=1,2,\mathrm{}`$, all $`(\alpha _1,\mathrm{},\alpha _m)^m`$, and all $`(x_1,\mathrm{},x_m)𝐗^m`$). On the other hand, for every symmetric and positive definite $`𝐊:𝐗^2`$ there exists a unique RKHS $``$ such that $`𝐊`$ is the kernel of $``$ (, Théorème 2).
We can see that the notions of a kernel of RKHS and of a symmetric positive definite function on $`𝐗^2`$ have the same content, and we will sometimes say “kernel on $`𝐗`$” to mean a symmetric positive definite function on $`𝐗^2`$. Kernels in this sense are the main source of RKHS in learning theory; see, e.g., , , and for numerous examples. Every kernel on $`𝐗`$ is a valid parameter for our decision algorithm; to apply Theorem 1 we can use the equivalent definition of $`𝐜_{}`$,
$$𝐜_{}:=\underset{x𝐗}{sup}\sqrt{𝐊(x,x)}.$$
A long list of RKHS together with their kernels is given in , §7.4. For concreteness, in this section we will use the Sobolev space $`𝒮`$ of absolutely continuous functions $`f`$ on $``$ with finite norm
$$f_𝒮:=\sqrt{_{\mathrm{}}^{\mathrm{}}f^2(x)dx+_{\mathrm{}}^{\mathrm{}}(f^{}(x))^2dx};$$
(10)
its kernel is
$$𝐊(x,x^{})=\frac{1}{2}\mathrm{exp}\left(\left|xx^{}\right|\right)$$
(see or , §7.4, Example 24). From the last equation we can see that $`𝐜_𝒮=1/\sqrt{2}`$.
### The square loss game
For the square loss game, $`\mathrm{\Gamma }=[0,1]`$ and $`\lambda (y,\gamma )=(y\gamma )^2`$, and so we have
$$Exp_\lambda (\gamma )=\lambda (1,\gamma )\lambda (0,\gamma )=(1\gamma )^2\gamma ^2=12\gamma ,$$
(11)
$$\lambda (p,\gamma )=p(1\gamma )^2+(1p)\gamma ^2=p(1p)+(\gamma p)^2,$$
and
$$\mathrm{\Gamma }_p=\{p\}.$$
(12)
Therefore,
$$𝐜_{\lambda ,}=\{\begin{array}{cc}𝐜_{}/2\hfill & \text{if }𝐜_{}1\hfill \\ (1+𝐜_{}^2)/4\hfill & \text{if }𝐜_{}<1;\hfill \end{array}$$
in particular, $`𝐜_{\lambda ,𝒮}=3/8`$ for the Sobolev space (10), and Theorem 1 implies
###### Corollary 1
Suppose the decision space is $`𝐗=`$. There is a decision strategy that guarantees that, for all $`N`$ and all decision rules $`D𝒮`$,
$$\underset{n=1}{\overset{N}{}}(y_n\gamma _n)^2\underset{n=1}{\overset{N}{}}(y_nD(x_n))^2+\frac{3}{8}\left(2D1_𝒮+1\right)\sqrt{N}$$
($`2D1`$ is the decision rule “normalized” to take values in $`[1,1]`$).
###### Remark 3
The games of this section illustrate Remark 1: here the decisions $`\gamma _n`$ are best interpreted as predictions of $`y_n`$. Loss functions $`\lambda `$ satisfying (12) are called *proper scoring rules*. Such loss functions “encourage honesty”: it is optimal to predict with the true probability (provided it is known). We will later see another loss function of this type (the log loss function).
To illustrate Corollary 1, suppose there are constants $`c>1`$ and $`d>1`$ and a good absolutely continuous decision rule $`D:[0,1]`$ such that $`|x_n|c`$, $`n=1,2,\mathrm{}`$, and $`|D^{}(x)|d`$ for all $`x𝐗`$. At rounds $`Ncd^2`$ the average loss of our decision algorithm will be almost as good as (or better than) the loss of $`D`$. We refrain from giving similar illustrations for the other corollaries in this section.
### The absolute loss game
In this game, $`\lambda (y,\gamma )=|y\gamma |`$ with $`\mathrm{\Gamma }=[0,1]`$. We find:
$$Exp_\lambda (\gamma )=\lambda (1,\gamma )\lambda (0,\gamma )=(1\gamma )\gamma =12\gamma $$
(the same as in the square loss case, (11)),
$$\lambda (p,\gamma )=p(1\gamma )+(1p)\gamma =p+(12p)\gamma ,$$
and
$$\mathrm{\Gamma }_p=\{\begin{array}{cc}\{0\}\hfill & \text{if }p<1/2\hfill \\ \{1\}\hfill & \text{if }p>1/2\hfill \\ [0,1]\hfill & \text{if }p=1/2.\hfill \end{array}$$
Therefore,
$$𝐜_{\lambda ,}=\frac{1}{2}\sqrt{1+𝐜_{}^2}$$
(in particular, $`𝐜_{\lambda ,𝒮}=\sqrt{6}/4`$), and we have the following corollary of Theorem 1.
###### Corollary 2
Let $`𝐗=`$. There is a decision strategy that produces decisions $`\gamma _n`$ such that, for all $`N`$ and all $`D𝒮`$,
$$\underset{n=1}{\overset{N}{}}\left|y_n\gamma _n\right|\underset{n=1}{\overset{N}{}}\left|y_nD(x_n)\right|+\frac{\sqrt{6}}{4}\left(2D1_𝒮+1\right)\sqrt{N}.$$
(13)
### The log loss game
For the log loss game, $`\mathrm{\Gamma }=(0,1)`$ and
$$\lambda (y,\gamma )=y\mathrm{ln}\gamma (1y)\mathrm{ln}(1\gamma ).$$
For this game, $`𝐜_{\lambda ,}<\mathrm{}`$ (assuming $`𝐜_{}<\mathrm{}`$) since its tails satisfy
$$f^{}(t)=g^{}(t)=\frac{1}{e^t1}e^t=O(t^2);$$
this will be also clear from the following direct calculation. Since
$$Exp_\lambda (\gamma )=\lambda (1,\gamma )\lambda (0,\gamma )=\mathrm{ln}\gamma +\mathrm{ln}(1\gamma )=\mathrm{ln}\frac{1\gamma }{\gamma },$$
$$\lambda (p,\gamma )=p\mathrm{ln}\gamma (1p)\mathrm{ln}(1\gamma )=\lambda (p,p)+D(p,\gamma )$$
(where $`D(p,\gamma ):=p\mathrm{ln}\frac{p}{\gamma }+(1p)\mathrm{ln}\frac{1p}{1\gamma }`$ is the Kullback distance between $`p`$ and $`\gamma `$, known to take its minimal value in $`\gamma `$ at $`\gamma =p`$), and $`\mathrm{\Gamma }_p=\{p\}`$, we can bound $`𝐜_{\lambda ,}`$ from above as follows:
$$\begin{array}{c}𝐜_{\lambda ,}^2=\underset{p(0,1)}{sup}p(1p)\left(\left(\mathrm{ln}\frac{1p}{p}\right)^2+𝐜_{}^2\right)\hfill \\ \hfill 𝐜_{}^2/4+\underset{p(0,1)}{sup}p(1p)\left(\mathrm{ln}\frac{1p}{p}\right)^2\\ \hfill 𝐜_{}^2/4+0.439𝐜_{}^2/4+0.44.\end{array}$$
Of course, for specific values of $`𝐜_{}`$ it is better to find the $`sup_{p(0,1)}`$ directly, without using this bound. Such a direct calculation shows that $`𝐜_{\lambda ,𝒮}0.6930.7`$, and Theorem 1 now implies the following.
###### Corollary 3
Some decision strategy in the log loss game with $`𝐗=`$ produces decisions $`\gamma _n`$ such that, for all $`N`$ and all $`D:𝐗(0,1)`$ with the log-likelihood ratio $`\mathrm{ln}\frac{D}{1D}`$ in $`𝒮`$,
$$\underset{n=1}{\overset{N}{}}\lambda (y_n,\gamma _n)\underset{n=1}{\overset{N}{}}\lambda (y_n,D(x_n))+0.7\left(\mathrm{ln}\frac{D}{1D}_𝒮+1\right)\sqrt{N}.$$
## 4 Idea of the proof of Theorem 1
This section describes the intuition behind the proof. The following sections, which carry out the proof, are formally independent of this section. We will also describe a general research program that may lead, it can be hoped, to many other results.
### Game-theoretic probability
Our proof technique is based on a game-theoretic alternative to the standard measure-theoretic axioms of probability (). Many of the standard laws of probability, including the weak and strong laws of large numbers, the central limit theorem, and the law of the iterated logarithm, can be restated in terms of perfect information games involving three key players: Reality, Forecaster, and Skeptic. A typical game-theoretic law of probability states that Skeptic has a strategy which, without risking bankruptcy, greatly enriches him if the law is violated. All such strategies for Skeptic were explicitly constructed continuous functions; game-theoretic laws of probability with a continuous strategy for Skeptic will be called “continuous laws of probability”.
Game-theoretic probability as developed in was to a large degree parallel to measure-theoretic probability. Following and the literature that this paper spawned, paper pointed out a surprising feature of game-theoretic probability: for any continuous law of probability, Forecaster has a strategy that prevents Skeptic’s capital from growing (cf. Lemma 1 below). In other words, for any continuous law of probability there is a forecasting strategy that is perfect as far as this law is concerned (we will say “perfect relative to” this law). This result was obtained in for binary forecasting, and in it was extended to more general protocols. Forecasting strategies obtained in this way from various laws of probability were called “defensive forecasting” strategies.
### General procedure
Now we are ready to describe a general procedure whose implementation leads, in the most straightforward case, to Theorem 1.
Choose a goal which could be achieved if you knew the true probabilities generating the observations. It is important that this goal should be “practical”, in the sense of being stated in terms of observable quantities, such as data, decisions, and observations. The goal is not allowed to contain theoretical quantities, such as the true probabilities themselves, and it should be achievable no matter what the true probabilities are. Construct a decision strategy which, using the true probabilities, leads to the goal.
Realistically, however, we do not know the true probabilities. To get rid of them, isolate the law of probability on which the proof that your decision strategy achieves the goal depends; typically, this law can be stated as a continuous game-theoretic law of probability. (If the proof depends on several laws, they should first be merged into a single law.) There is a forecasting strategy whose forecasts are at least as good as (and often better than) the true probabilities, as far as the law you have just isolated is concerned. It remains to feed your decision strategy with those forecasts.
Implementing this procedure for various interesting goals appears to be a promising research program.
### Introduction to the proof
In this paper our goal is to achieve (9), which we roughly rewrite as
$$\underset{n=1}{\overset{N}{}}\lambda (y_n,\gamma _n)\underset{n=1}{\overset{N}{}}\lambda (y_n,D(x_n)),$$
where the informal notation $``$ is used to mean that the left-hand side does not exceed the right-hand side plus a quantity small as compared to $`N`$. The goal is stated in terms of the observables.
Let us see how our goal could be achieved if we knew the true probabilities $`p_n`$ that $`y_n=1`$ (slightly more formally, $`p_n`$ is the conditional probability that $`y_n=1`$ given the available information). By the law of large numbers (see, e.g., , Theorem VII.5.4, for a suitable measure-theoretic statement and , Theorem 4.1, for its game-theoretic counterpart), we expect
$$\left|\underset{n=1}{\overset{N}{}}f(p_n,x_n)(y_np_n)\right|N$$
(14)
if $`f`$ is a bounded function (assumed measurable in the measure-theoretic case). If $`f`$ is allowed to range over a function class $``$ that is not excessively wide, (14) will still continue to hold uniformly in $`f`$.
Suppose, for simplicity, that $`\mathrm{\Gamma }_p`$ is a singleton for all $`p[0,1]`$; the only element of $`\mathrm{\Gamma }_p`$ will be denoted $`G(p)`$. Our decision strategy will make the decision $`G(p_n)`$ at round $`n`$, i.e., the decision that leads to the smallest expected loss. We will sometimes say that $`G`$ is our “choice function”.
Notice that
$$\lambda (y,\gamma )\lambda (p,\gamma )=(yp)\left(\lambda (1,\gamma )\lambda (0,\gamma )\right)$$
always holds (this can be checked by subtracting (1) from $`\lambda (y,\gamma ):=y\lambda (1,\gamma )+(1y)\lambda (0,\gamma )`$). In conjunction with the law of large numbers (14) this implies
$$\begin{array}{c}\underset{n=1}{\overset{N}{}}\lambda (y_n,\gamma _n)=\underset{n=1}{\overset{N}{}}\lambda (y_n,G(p_n))\hfill \\ \hfill =\underset{n=1}{\overset{N}{}}\lambda (p_n,G(p_n))+\underset{n=1}{\overset{N}{}}\left(\lambda (y_n,G(p_n))\lambda (p_n,G(p_n))\right)\\ \hfill =\underset{n=1}{\overset{N}{}}\lambda (p_n,G(p_n))+\underset{n=1}{\overset{N}{}}(y_np_n)\left(\lambda (1,G(p_n))\lambda (0,G(p_n))\right)\underset{n=1}{\overset{N}{}}\lambda (p_n,G(p_n))\\ \hfill \underset{n=1}{\overset{N}{}}\lambda (p_n,D(x_n))=\underset{n=1}{\overset{N}{}}\lambda (y_n,D(x_n))\underset{n=1}{\overset{N}{}}\left(\lambda (y_n,D(x_n))\lambda (p_n,D(x_n))\right)\\ \hfill =\underset{n=1}{\overset{N}{}}\lambda (y_n,D(x_n))\underset{n=1}{\overset{N}{}}(y_np_n)\left(\lambda (1,D(x_n))\lambda (0,D(x_n))\right)\\ \hfill \underset{n=1}{\overset{N}{}}\lambda (y_n,D(x_n)).\end{array}$$
(15)
This shows that we can achieve our goal if we know the true probabilities, and it remains to replace the true probabilities with the forecasts that are perfect relative to the law of large numbers.
For clarity, let us summarize the idea of the proof expressed by (15). To show that the actual loss of our decision strategy does not exceed the actual loss of a decision rule $`D`$ by much, we notice that:
* the actual loss $`_{n=1}^N\lambda (y_n,G(p_n))`$ of our decision strategy is approximately equal, by the law of large numbers, to the (one-step-ahead conditional) expected loss $`_{n=1}^N\lambda (p_n,G(p_n))`$ of our strategy;
* since we used the expected loss minimization principle, the expected loss of our strategy does not exceed the expected loss of $`D`$;
* the expected loss $`_{n=1}^N\lambda (p_n,D(x_n))`$ of $`D`$ is approximately equal to its actual loss $`_{n=1}^N\lambda (y_n,D(x_n))`$ (by the law of large numbers).
To get the strongest possible result, we will have to use more specific laws of probability than the general law of large numbers. It will be convenient to use the following informal terminology introduced in . Let $`p_n`$ be the forecasts output by some forecasting strategy (rather than the true probabilities). We say that the forecasting strategy has good *calibration-cum-resolution* if the left-hand side of (14) is much less than $`N`$ for a relatively wide class of functions $`f:[0,1]\times 𝐗`$ and large $`N`$. We say that the strategy has good *calibration* if
$$\left|\underset{n=1}{\overset{N}{}}(y_np_n)f(p_n)\right|N$$
for a wide class of functions $`f:[0,1]`$ and large $`N`$. Finally, we say that the strategy has good *resolution* if
$$\left|\underset{n=1}{\overset{N}{}}(y_np_n)f(x_n)\right|N$$
for a wide class of $`f:𝐗`$ and for large $`N`$. For a detailed discussion and examples, see .
Notice that in applying the law of large numbers to establishing the two approximate inequalities in (15) we need not general $`f=f(p,x)`$ but only $`f=f(p)`$ (known in advance) and $`f=f(x)`$. In particular, we only need calibration and resolution separately, not calibration-cum-resolution. These are the two specific probability laws we will be concerned with.
The requirement that $`\mathrm{\Gamma }_p`$ should always be a singleton (in fact, we will even need the function $`G(p)`$ to be continuous) is restrictive: for example, it is not satisfied for the absolute loss function. To deal with this problem, we will have to consider forecasting strategies that output extended forecasts $`(p_n,q_n)[0,1]^2`$, where $`p_n`$ is the forecast of $`y_n`$ and the extra component $`q_n`$ will play a more technical role.
The next section is devoted to constructing a perfect forecasting strategy relative to the law of large numbers. In the following section we will be able to prove Theorem 1.
## 5 The algorithm of large numbers
This section is the core of our proof of Theorem 1. First we describe a forecasting protocol in which Forecaster tries to predict the observations chosen by Reality. Following , we introduce another player, Skeptic, who is allowed to bet at the odds implied by Forecaster’s moves.
Binary Forecasting Game I
Players: Reality, Forecaster, Skeptic
Protocol:
FOR $`n=1,2,\mathrm{}`$:
Reality announces $`x_n𝐗`$.
Forecaster announces $`(p_n,q_n)[0,1]^2`$.
Skeptic announces $`s_n`$.
Reality announces $`y_n\{0,1\}`$.
$`𝒦_n:=𝒦_{n1}+s_n(y_np_n)`$.
END FOR.
The real forecast is $`p_n`$ (the “probability” that $`y_n=1`$), which is interpreted as the price Forecaster charges for a ticket paying $`y_n`$; $`s_n`$ is the number of tickets Skeptic decides to buy. The protocol describes not only the players’ moves but also the changes in Skeptic’s capital $`𝒦_n`$; its initial value $`𝒦_0`$ can be an arbitrary real number. Skeptic demonstrates that the forecasts are poor if he manages to multiply his initial capital (assumed positive) manyfold without risking bankruptcy (i.e., $`𝒦_n`$ becoming negative). Forecaster also provides an additional number $`q_n[0,1]`$ which does not affect Skeptic’s capital; intuitively, the role of $`q_n`$ is to help those of Forecaster’s customers who find themselves in a position of Buridan’s ass (find two or more actions equally attractive in view of the forecast $`p_n`$) to break the tie.
The main difference between our decision protocol (stated at the beginning of §2) and the protocols of this section is that in the latter Forecaster implicitly claims (by pricing the tickets) that he has the fullest possible knowledge of the way Reality chooses the observations, and Skeptic tries to prove him wrong by gambling against him. In the decision protocol, Decision Maker does no make any such claims and simply tries to minimize his losses.
It will be convenient to make the set $`[0,1]^2`$ from which the forecasts $`(p_n,q_n)`$ are chosen into a topological space. The *lexicographic square* $`\mathrm{\pounds }`$ is defined to be the set $`[0,1]^2`$ equipped with the following linear order: if $`(x_1,y_1)`$ and $`(x_2,y_2)`$ are two points in $`\mathrm{\pounds }`$, $`(x_1,y_1)<(x_2,y_2)`$ means that either $`x_1<x_2`$ or $`x_1=x_2,y_1<y_2`$. (Cf. , Problem 3.12.3(d).) The topology on the lexicographic square is, as usual, generated by the open intervals
$$(a,b):=\left\{u\mathrm{\pounds }|a<u<b\right\},$$
$`a`$ and $`b`$ ranging over $`\mathrm{\pounds }`$. As a topological space, the lexicographic square is normal (, Problem 1.7.4(d)), compact (, Problem 3.12.3(a), , Problem 5.C), and connected (, Problem 6.3.2(a), , Problem 1.I(d)).
As in , we will see that for any continuous strategy for Skeptic there exists a strategy for Forecaster that does not allow Skeptic’s capital to grow, regardless of what Reality is doing. To state this observation in its strongest form, we make Skeptic announce his strategy for each round before Forecaster’s move on that round rather than announce his full strategy at the beginning of the game. Therefore, we consider the following perfect-information game:
Binary Forecasting Game II
Players: Reality, Forecaster, Skeptic
Protocol:
FOR $`n=1,2,\mathrm{}`$:
Reality announces $`x_n𝐗`$.
Skeptic announces continuous $`S_n:\mathrm{\pounds }`$.
Forecaster announces $`(p_n,q_n)\mathrm{\pounds }`$.
Reality announces $`y_n\{0,1\}`$.
$`𝒦_n:=𝒦_{n1}+S_n(p_n,q_n)(y_np_n)`$.
END FOR.
###### Lemma 1
Forecaster has a strategy in Binary Forecasting Game II that ensures $`𝒦_0𝒦_1𝒦_2\mathrm{}`$.
Before proving this lemma, we will need another lemma, which will play the role of the Intermediate Value Theorem, used in .
###### Lemma 2
If a continuous function $`f:\mathrm{\pounds }`$ takes both positive and negative values, there exists $`x\mathrm{\pounds }`$ such that $`f(x)=0`$.
* A continuous image of a connected compact set is connected (, Theorem 6.1.4) and compact (, Theorem 3.1.10). Therefore, $`f(\mathrm{\pounds })`$ is a closed interval.
* Forecaster can now use the following strategy to ensure $`𝒦_0𝒦_1\mathrm{}`$:
+ if the function $`S_n(p,q)`$ takes value 0, choose $`(p_n,q_n)`$ such that $`S_n(p_n,q_n)=0`$;
+ if $`S_n`$ is always positive, take $`p_n:=1`$ and choose $`q_n[0,1]`$ arbitrarily;
+ if $`S_n`$ is always negative, take $`p_n:=0`$ and choose $`q_n[0,1]`$ arbitrarily.
A kernel $`𝐊`$ on $`\mathrm{\pounds }\times 𝐗`$ is *forecast-continuous* if the function $`𝐊((p,q,x),(p^{},q^{},x^{}))`$ is continuous in $`(p,q,p^{},q^{})\mathrm{\pounds }^2`$, for each fixed $`(x,x^{})𝐗^2`$. (Kernels on $`\mathrm{\pounds }\times 𝐗`$ are defined analogously to kernels on $`𝐗`$.) For such a kernel the function
$$S_n(p,q):=\underset{i=1}{\overset{n1}{}}𝐊((p,q,x_n),(p_i,q_i,x_i))(y_ip_i)+\frac{1}{2}𝐊((p,q,x_n),(p,q,x_n))(12p)$$
(16)
is continuous in $`(p,q)\mathrm{\pounds }`$.
The lexicographic algorithm of large numbers ($`\mathrm{\pounds }`$ALN)
Parameter: forecast-continuous kernel $`𝐊`$ on $`\mathrm{\pounds }\times 𝐗`$
FOR $`n=1,2,\mathrm{}`$:
Read $`x_n𝐗`$.
Define $`S_n(p,q)`$ by (16), $`(p,q)\mathrm{\pounds }`$.
Output any root $`(p,q)`$ of $`S_n(p,q)=0`$ as $`(p_n,q_n)`$;
if there are no roots,
set $`p_n:=(1+signS_n)/2`$ and set $`q_n`$ to any number in $`[0,1]`$.
Read $`y_n\{0,1\}`$.
END FOR.
(Notice that $`signS_n`$ is well defined by Lemma 2.) It is well known that there exists a function $`\mathrm{\Phi }:\mathrm{\pounds }\times 𝐗`$ (a *feature mapping* taking values in a Hilbert space $``$) such that
$$𝐊(a,b)=\mathrm{\Phi }(a)\mathrm{\Phi }(b),a,b\mathrm{\pounds }\times 𝐗.$$
(17)
(For example, we can take the RKHS on $`\mathrm{\pounds }\times 𝐗`$ with kernel $`𝐊`$ as $``$ and take $`a𝐊_a`$ as the feature mapping $`\mathrm{\Phi }`$; there are, however, easier and more transparent constructions.) It can be shown that $`\mathrm{\Phi }(p,q,x)`$ is *forecast-continuous*, i.e., continuous in $`(p,q)\mathrm{\pounds }`$ for each fixed $`x𝐗`$, if and only if the kernel $`𝐊`$ defined by (17) is forecast-continuous (see, e.g., , Appendix B).
###### Theorem 2
Let $`𝐊`$ be the kernel defined by (17) for a forecast-continuous feature mapping $`\mathrm{\Phi }:\mathrm{\pounds }\times 𝐗`$. The lexicographic algorithm of large numbers with parameter $`𝐊`$ outputs $`(p_n,q_n)`$ such that
$$\underset{n=1}{\overset{N}{}}(y_np_n)\mathrm{\Phi }(p_n,q_n,x_n)^2\underset{n=1}{\overset{N}{}}p_n(1p_n)\mathrm{\Phi }(p_n,q_n,x_n)^2$$
(18)
always holds for all $`N=1,2,`$.
* Following $`\mathrm{\pounds }`$ALN Forecaster ensures that Skeptic will never increase his capital with the strategy
$$\begin{array}{c}s_n:=\underset{i=1}{\overset{n1}{}}𝐊((p_n,q_n,x_n),(p_i,q_i,x_i))(y_ip_i)\hfill \\ \hfill +\frac{1}{2}𝐊((p_n,q_n,x_n),(p_n,q_n,x_n))(12p_n).\end{array}$$
(19)
Using the formula
$$(y_np_n)^2=p_n(1p_n)+(12p_n)(y_np_n)$$
(which can be checked by setting $`y_n:=0`$ and $`y_n:=1`$), we can see that the increase in Skeptic’s capital when he follows (19) is
$`𝒦_N𝒦_0`$ $`={\displaystyle \underset{n=1}{\overset{N}{}}}s_n(y_np_n)`$
$`={\displaystyle \underset{n=1}{\overset{N}{}}}{\displaystyle \underset{i=1}{\overset{n1}{}}}𝐊((p_n,q_n,x_n),(p_i,q_i,x_i))(y_np_n)(y_ip_i)`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{n=1}{\overset{N}{}}}𝐊((p_n,q_n,x_n),(p_n,q_n,x_n))(12p_n)(y_np_n)`$
$`={\displaystyle \frac{1}{2}}{\displaystyle \underset{n=1}{\overset{N}{}}}{\displaystyle \underset{i=1}{\overset{N}{}}}𝐊((p_n,q_n,x_n),(p_i,q_i,x_i))(y_np_n)(y_ip_i)`$
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{n=1}{\overset{N}{}}}𝐊((p_n,q_n,x_n),(p_n,q_n,x_n))(y_np_n)^2`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{n=1}{\overset{N}{}}}𝐊((p_n,q_n,x_n),(p_n,q_n,x_n))(12p_n)(y_np_n)`$
$`={\displaystyle \frac{1}{2}}{\displaystyle \underset{n=1}{\overset{N}{}}}{\displaystyle \underset{i=1}{\overset{N}{}}}𝐊((p_n,q_n,x_n),(p_i,q_i,x_i))(y_np_n)(y_ip_i)`$
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{n=1}{\overset{N}{}}}𝐊((p_n,q_n,x_n),(p_n,q_n,x_n))p_n(1p_n)`$
$`={\displaystyle \frac{1}{2}}{\displaystyle \underset{n=1}{\overset{N}{}}}(y_np_n)\mathrm{\Phi }(p_n,q_n,x_n)^2{\displaystyle \frac{1}{2}}{\displaystyle \underset{n=1}{\overset{N}{}}}p_n(1p_n)\mathrm{\Phi }(p_n,q_n,x_n)^2,`$
which immediately implies (18).
### Resolution
This subsection makes the next step in our proof of Theorem 1. Its forecasting protocol is:
FOR $`n=1,2,\mathrm{}`$:
Reality announces $`x_n𝐗`$.
Forecaster announces $`(p_n,q_n)[0,1]`$.
Reality announces $`y_n\{0,1\}`$.
END FOR.
Our goal is to prove the following result (although in §6 we will need a slight modification of this result rather than the result itself).
###### Theorem 3
Let $``$ be an RKHS on $`𝐗`$. The forecasts $`(p_n,q_n)`$ output by $`\mathrm{\pounds }`$ALN always satisfy
$$\left|\underset{n=1}{\overset{N}{}}(y_np_n)f(x_n)\right|\frac{𝐜_{}}{2}f_{}\sqrt{N}$$
for all $`N`$ and all functions $`f`$.
* Applying ALN to the feature mapping $`x𝐗𝐊_x`$ and using (18), we obtain
$$\begin{array}{c}\left|\underset{n=1}{\overset{N}{}}(y_np_n)f(x_n)\right|=\left|\underset{n=1}{\overset{N}{}}(y_np_n)𝐊_{x_n},f_{}\right|\hfill \\ \hfill =\left|\underset{n=1}{\overset{N}{}}(y_np_n)𝐊_{x_n},f_{}\right|\underset{n=1}{\overset{N}{}}(y_np_n)𝐊_{x_n}_{}f_{}\\ \hfill f_{}\sqrt{\underset{n=1}{\overset{N}{}}p_n(1p_n)𝐊(x_n,x_n)}𝐜_{}f_{}\sqrt{\underset{n=1}{\overset{N}{}}p_n(1p_n)}\end{array}$$
(20)
for any $`f`$.
###### Remark 4
In the terminology introduced in the previous section, Theorem 3 is about resolution. This is sufficient for the purpose of this paper, but it is easy to see that similar statements hold for calibration-cum-resolution and calibration. For example, let $``$ be an RKHS on $`\mathrm{\pounds }\times 𝐗`$. The forecasts $`(p_n,q_n)`$ output by $`\mathrm{\pounds }`$ALN always satisfy
$$\left|\underset{n=1}{\overset{N}{}}(y_np_n)f(p_n,q_n,x_n)\right|\frac{𝐜_{}}{2}f_{}\sqrt{N}$$
for all $`N`$ and all functions $`f`$.
## 6 Proof of Theorem 1
Before starting the proof proper, we need to discuss two topics: choosing a suitable choice function and “mixing” different feature mappings.
### The canonical choice function
Let us say that a straight line $`(1p)x+py=c`$ in the $`(x,y)`$-plane, where $`p[0,1]`$ and $`c`$, is *southwest of* the superdecision set $`\mathrm{\Sigma }`$ (defined by (3)) if
$$(x,y)\mathrm{\Sigma }:(1p)x+pyc.$$
For each $`p[0,1]`$ let $`c(p)`$ be the largest $`c`$ (which obviously exists) such that the line $`(1p)x+py=c`$ is southwest of $`\mathrm{\Sigma }`$. It is clear that, for $`p(0,1)`$, the line $`(1p)x+py=c(p)`$ intersects $`\mathrm{\Sigma }`$ and the intersection, being compact and convex, has the form $`[A(p),B(p)]`$, where $`A(p)`$ and $`B(p)`$ are points (perhaps $`A(p)=B(p)`$) on the line. For concreteness, let $`A(p)`$ be northwest of $`B(p)`$ (i.e., if $`A(p)=(A_0,A_1)`$ and $`B(p)=(B_0,B_1)`$, we assume that $`A_0B_0`$ and $`A_1B_1`$). Now we can define the *canonical choice function* $`G`$ associated with $`(\mathrm{\Gamma },\lambda )`$ as follows:
* if $`0<p<1`$ and $`q[0,1]`$, $`G(p,q)`$ is defined to be any $`\gamma \mathrm{\Gamma }`$ satisfying
$$(\lambda (0,\gamma ),\lambda (1,\gamma ))=(1q)A(p)+qB(p);$$
the existence of such a $`\gamma `$ is obvious;
* if $`p=0`$ and $`q[0,1]`$, $`G(p,q)`$ is defined to be any fixed $`\gamma _0\mathrm{\Gamma }`$ satisfying
$$(\lambda (0,\gamma _0),\lambda (1,\gamma _0))=(C_0,f(C_0))$$
($`C_0`$ and $`f`$ are defined in (2) and (4)); if $`f(C_0)=\mathrm{}`$, such a $`\gamma _0`$ does not exist and $`G(p,q)`$ is undefined;
* if $`p=1`$ and $`q[0,1]`$, $`G(p,q)`$ is defined to be any fixed $`\gamma _1\mathrm{\Gamma }`$ such that
$$(\lambda (0,\gamma _1),\lambda (1,\gamma _1))=(g(C_1),C_1)$$
($`C_1`$ and $`g`$ are defined in (2) and (5)); if $`g(C_1)=\mathrm{}`$, such a $`\gamma _1`$ does not exist and $`G(p,q)`$ is undefined.
It is easy to see that the function $`(\lambda (0,G(p,q)),\lambda (1,G(p,q)))`$ is continuous in $`(p,q)domG`$ and, therefore, $`Exp_{\lambda ,G}(p,q):=Exp_\lambda (G(p,q))`$ is continuous in $`(p,q)domG`$. We defined $`G`$ in such a way that it is a “perfect” choice function: $`\lambda (p,G(p,q))=inf_{\gamma \mathrm{\Gamma }}\lambda (p,\gamma )`$ for virtually all $`(p,q)`$ (in any case, for all $`(p,q)domG`$).
### Mixing
In the proof of Theorem 1 we will mix the feature mapping $`\mathrm{\Phi }_0(p,q,x):=Exp_{\lambda ,G}(p,q)`$ (into $`_0:=`$) and the feature mapping $`\mathrm{\Phi }_1(p,q,x):=𝐊_x`$ used in the proof of Theorem 3 (as discussed in §4, we will have to achieve two goals simultaneously, only one of them connected with resolution). This can be done using the following corollary of Theorem 2.
###### Corollary 4
Let $`\mathrm{\Phi }_j:\mathrm{\pounds }\times 𝐗_j`$, $`j=0,1`$, be forecast-continuous mappings from $`\mathrm{\pounds }\times 𝐗`$ to Hilbert spaces $`_j`$. The forecasts output by $`\mathrm{\pounds }`$ALN with a suitable kernel parameter always satisfy
$$\begin{array}{c}\underset{n=1}{\overset{N}{}}(y_np_n)\mathrm{\Phi }_j(p_n,q_n,x_n)__j^2\hfill \\ \hfill \underset{n=1}{\overset{N}{}}p_n(1p_n)\left(\mathrm{\Phi }_0(p_n,q_n,x_n)__0^2+\mathrm{\Phi }_1(p_n,q_n,x_n)__1^2\right)\end{array}$$
for all $`N`$ and for both $`j=0`$ and $`j=1`$.
* Define the direct sum $``$ of $`_0`$ and $`_1`$ as the Cartesian product $`_0\times _1`$ equipped with the inner product
$$g,g^{}_{}=(g_0,g_1),(g_0^{},g_1^{})_{}:=\underset{j=0}{\overset{1}{}}g_j,g_j^{}__j.$$
Now we can define $`\mathrm{\Phi }:\mathrm{\pounds }\times 𝐗`$ by
$$\mathrm{\Phi }(p,q,x):=(\mathrm{\Phi }_0(p,q,x),\mathrm{\Phi }_1(p,q,x));$$
the corresponding kernel is
$$\begin{array}{c}𝐊((p,q,x),(p^{},q^{},x^{})):=\mathrm{\Phi }(p,q,x),\mathrm{\Phi }(p^{},q^{},x^{})_{}\hfill \\ \hfill =\underset{j=0}{\overset{1}{}}\mathrm{\Phi }_j(p,q,x),\mathrm{\Phi }_j(p^{},q^{},x^{})__j=\underset{j=0}{\overset{1}{}}𝐊_j((p,q,x),(p^{},q^{},x^{})),\end{array}$$
where $`𝐊_0`$ and $`𝐊_1`$ are the kernels corresponding to $`\mathrm{\Phi }_0`$ and $`\mathrm{\Phi }_1`$, respectively. It is clear that this kernel is forecast-continuous. Applying $`\mathrm{\pounds }`$ALN to it and using (18), we obtain
$$\begin{array}{c}\underset{n=1}{\overset{N}{}}(y_np_n)\mathrm{\Phi }_j(p_n,q_n,x_n)__j^2\hfill \\ \hfill (\underset{n=1}{\overset{N}{}}(y_np_n)\mathrm{\Phi }_0(p_n,q_n,x_n),\underset{n=1}{\overset{N}{}}(y_np_n)\mathrm{\Phi }_1(p_n,q_n,x_n))_{}^2\\ \hfill =\underset{n=1}{\overset{N}{}}(y_np_n)\mathrm{\Phi }(p_n,q_n,x_n)_{}^2\underset{n=1}{\overset{N}{}}p_n(1p_n)\mathrm{\Phi }(p_n,q_n,x_n)_{}^2\\ \hfill =\underset{n=1}{\overset{N}{}}p_n(1p_n)\underset{j=0}{\overset{1}{}}\mathrm{\Phi }_j(p_n,q_n,x_n)__j^2.\end{array}$$
Merging $`\mathrm{\Phi }_0`$ and $`\mathrm{\Phi }_1`$ by Corollary 4, we obtain
$$\begin{array}{c}\left|\underset{n=1}{\overset{N}{}}(y_np_n)Exp_{\lambda ,G}(p_n,q_n)\right|=\underset{n=1}{\overset{N}{}}(y_np_n)\mathrm{\Phi }_0(p_n,q_n,x_n)_{}\hfill \\ \hfill \sqrt{\underset{n=1}{\overset{N}{}}p_n(1p_n)\left(Exp_{\lambda ,G}^2(p_n,q_n)+𝐊(x_n,x_n)\right)}\end{array}$$
(21)
and, using (20),
$$\begin{array}{c}\left|\underset{n=1}{\overset{N}{}}(y_np_n)f(x_n)\right|\underset{n=1}{\overset{N}{}}(y_np_n)𝐊_{x_n}_{}f_{}\hfill \\ \hfill =\underset{n=1}{\overset{N}{}}(y_np_n)\mathrm{\Phi }_1(p_n,q_n,x_n)_{}f_{}\\ \hfill f_{}\sqrt{\underset{n=1}{\overset{N}{}}p_n(1p_n)\left(Exp_{\lambda ,G}^2(p_n,q_n)+𝐊(x_n,x_n)\right)},\end{array}$$
(22)
for each function $`f`$.
### Proof: Part I
In this subsection we will assume that $`domG=\mathrm{\pounds }`$. Subtracting (1) from $`\lambda (y,\gamma )=y\lambda (1,\gamma )+(1y)\lambda (0,\gamma )`$, we obtain
$$\lambda (y,\gamma )\lambda (p,\gamma )=(yp)\left(\lambda (1,\gamma )\lambda (0,\gamma )\right)=(yp)Exp_\lambda (\gamma )$$
(23)
(we already did this in §4, but we promised that the rest of the paper would be formally independent of §4). Using the last equality and (21)–(22), we obtain for the decision strategy $`\gamma _n:=G(p_n,q_n)`$ based on the $`(p_n,q_n)`$ output by $`\mathrm{\pounds }`$ALN with the merged kernel as parameter:
$`{\displaystyle \underset{n=1}{\overset{N}{}}}\lambda (y_n,\gamma _n)={\displaystyle \underset{n=1}{\overset{N}{}}}\lambda (y_n,G(p_n,q_n))`$
$`={\displaystyle \underset{n=1}{\overset{N}{}}}\lambda (p_n,G(p_n,q_n))+{\displaystyle \underset{n=1}{\overset{N}{}}}\left(\lambda (y_n,G(p_n,q_n))\lambda (p_n,G(p_n,q_n))\right)`$
$`={\displaystyle \underset{n=1}{\overset{N}{}}}\lambda (p_n,G(p_n,q_n))+{\displaystyle \underset{n=1}{\overset{N}{}}}(y_np_n)Exp_{\lambda ,G}(p_n,q_n)`$
$`{\displaystyle \underset{n=1}{\overset{N}{}}}\lambda (p_n,G(p_n,q_n))+\sqrt{{\displaystyle \underset{n=1}{\overset{N}{}}}p_n(1p_n)\left(Exp_{\lambda ,G}^2(p_n,q_n)+𝐊(x_n,x_n)\right)}`$
$`{\displaystyle \underset{n=1}{\overset{N}{}}}\lambda (p_n,D(x_n))+\sqrt{{\displaystyle \underset{n=1}{\overset{N}{}}}p_n(1p_n)\left(Exp_{\lambda ,G}^2(p_n,q_n)+𝐊(x_n,x_n)\right)}`$
$`={\displaystyle \underset{n=1}{\overset{N}{}}}\lambda (y_n,D(x_n)){\displaystyle \underset{n=1}{\overset{N}{}}}\left(\lambda (y_n,D(x_n))\lambda (p_n,D(x_n))\right)`$
$`+\sqrt{{\displaystyle \underset{n=1}{\overset{N}{}}}p_n(1p_n)\left(Exp_{\lambda ,G}^2(p_n,q_n)+𝐊(x_n,x_n)\right)}`$
$`={\displaystyle \underset{n=1}{\overset{N}{}}}\lambda (y_n,D(x_n)){\displaystyle \underset{n=1}{\overset{N}{}}}(y_np_n)Exp_{\lambda ,D}(x_n))`$
$`+\sqrt{{\displaystyle \underset{n=1}{\overset{N}{}}}p_n(1p_n)\left(Exp_{\lambda ,G}^2(p_n,q_n)+𝐊(x_n,x_n)\right)}`$
$`{\displaystyle \underset{n=1}{\overset{N}{}}}\lambda (y_n,D(x_n))+Exp_{\lambda ,D}_{}\sqrt{{\displaystyle \underset{n=1}{\overset{N}{}}}p_n(1p_n)\left(Exp_{\lambda ,G}^2(p_n,q_n)+𝐊(x_n,x_n)\right)}`$
$`+\sqrt{{\displaystyle \underset{n=1}{\overset{N}{}}}p_n(1p_n)\left(Exp_{\lambda ,G}^2(p_n,q_n)+𝐊(x_n,x_n)\right)}`$
$`={\displaystyle \underset{n=1}{\overset{N}{}}}\lambda (y_n,D(x_n))+\left(Exp_{\lambda ,D}_{}+1\right)\sqrt{{\displaystyle \underset{n=1}{\overset{N}{}}}p_n(1p_n)\left(Exp_{\lambda ,G}^2(p_n,q_n)+𝐊(x_n,x_n)\right)}`$
$`{\displaystyle \underset{n=1}{\overset{N}{}}}\lambda (y_n,D(x_n))+\left(Exp_{\lambda ,D}_{}+1\right)𝐜_{\lambda ,}\sqrt{N}.`$
It remains to show that $`𝐜_{\lambda ,}<\mathrm{}`$ (assuming $`𝐜_{}<\mathrm{}`$, here and in the rest of this section). In this case, $`domG=\mathrm{\pounds }`$, this is easy: essentially, this is the case of a bounded loss function (the reservation “essentially” is needed since $`\mathrm{\Gamma }`$ can contain “litter”—decisions dominated by other decisions in $`\mathrm{\Gamma }`$). Since $`Exp_{\lambda ,G}`$ is continuous and $`\mathrm{\pounds }`$ is compact,
$$\underset{(p,q)\mathrm{\pounds }}{sup}p(1p)\left(Exp_{\lambda ,G}^2(p,q)+𝐜_{}^2\right)<\mathrm{}.$$
### Proof: Part II
The *stripped lexicographic square* is the subset
$${}_{}{}^{}\mathrm{\pounds }_{}^{}:=(0,1)\times [0,1]$$
of $`\mathrm{\pounds }`$. In this subsection we consider the case $`domG={}_{}{}^{}\mathrm{\pounds }_{}^{}`$.
The order and topology on $`{}_{}{}^{}\mathrm{\pounds }_{}^{}`$ are inherited from $`\mathrm{\pounds }`$. The following analogue of Lemma 2 still holds.
###### Lemma 3
If a continuous function $`f:{}_{}{}^{}\mathrm{\pounds }_{}^{}`$ takes both positive and negative values, it also takes the value $`0`$.
* See the proof of Lemma 2; $`f({}_{}{}^{}\mathrm{\pounds }_{}^{})`$ is still a connected set in $``$.
A kernel $`𝐊`$ on $`{}_{}{}^{}\mathrm{\pounds }_{}^{}\times 𝐗`$ is *forecast-continuous* if the function $`𝐊((p,q,x),(p^{},q^{},x^{}))`$ is continuous in $`(p,q,p^{},q^{})({}_{}{}^{}\mathrm{\pounds }_{}^{})^2`$. The function (16) is then continuous in $`(p,q){}_{}{}^{}\mathrm{\pounds }_{}^{}`$, and for our current kernel
$$𝐊((p,q,x),(p^{},q^{},x^{}))=Exp_{\lambda ,G}(p,q)Exp_{\lambda ,G}(p^{},q^{})+𝐊_x,𝐊_x^{}_{}$$
(24)
it equals
$$\begin{array}{c}S_n(p,q)=\underset{i=1}{\overset{n1}{}}\left(Exp_{\lambda ,G}(p,q)Exp_{\lambda ,G}(p_i,q_i)+𝐊_{x_n},𝐊_{x_i}_{}\right)(y_ip_i)\hfill \\ \hfill +\frac{1}{2}\left(Exp_{\lambda ,G}^2(p,q)+𝐊_{x_n}_{}^2\right)(12p)\\ \hfill =AExp_{\lambda ,G}(p,q)+B+\frac{1}{2}Exp_{\lambda ,G}^2(p,q)(12p)+Cp,\end{array}$$
(25)
where $`A`$, $`B`$, and $`C`$ do not depend on $`(p,q)`$. Since $`domG={}_{}{}^{}\mathrm{\pounds }_{}^{}`$, $`\left|Exp_{\lambda ,G}(p,q)\right|\mathrm{}`$ as $`p0`$ or $`p1`$, and so
$$\underset{\begin{array}{c}(p,q)(1,0)\\ (p,q){}_{}{}^{}\mathrm{\pounds }_{}^{}\end{array}}{lim}S_n(p,q)=\mathrm{}$$
(26)
and
$$\underset{\begin{array}{c}(p,q)(0,1)\\ (p,q){}_{}{}^{}\mathrm{\pounds }_{}^{}\end{array}}{lim}S_n(p,q)=\mathrm{}.$$
(27)
The *stripped lexicographic ALN* (or, briefly, $`{}_{}{}^{}\mathrm{\pounds }_{}^{}`$ALN) is defined as the lexicographic ALN except that:
* its parameter is a forecast-continuous kernel $`𝐊`$ on $`{}_{}{}^{}\mathrm{\pounds }_{}^{}\times 𝐗`$;
* it outputs a root $`(p,q)`$ (an element of $`{}_{}{}^{}\mathrm{\pounds }_{}^{}=domS_n`$) of the equation $`S_n(p,q)=0`$ as $`(p_n,q_n)`$ and crashes if this equation does not have roots (this will never happen for the kernel (24)).
Because of (26) and (27), $`{}_{}{}^{}\mathrm{\pounds }_{}^{}`$ALN applied to the kernel (24) on $`{}_{}{}^{}\mathrm{\pounds }_{}^{}\times 𝐗`$ still ensures that (18) holds for our feature mapping $`(\mathrm{\Phi }_0,\mathrm{\Phi }_1)`$; this algorithm never crashes and, of course, never outputs $`(p_n,q_n)`$ with $`p_n\{0,1\}`$. We can see that the proof of (9) given in the previous subsection still works.
Let us now prove that $`𝐜_{\lambda ,}<\mathrm{}`$ when $`domG={}_{}{}^{}\mathrm{\pounds }_{}^{}`$. It suffices to check that
$$\underset{\begin{array}{c}(p,q)(0,1)\\ (p,q){}_{}{}^{}\mathrm{\pounds }_{}^{}\end{array}}{lim\; sup}pExp_{\lambda ,G}^2(p,q)<\mathrm{}$$
(28)
and
$$\underset{\begin{array}{c}(p,q)(1,0)\\ (p,q){}_{}{}^{}\mathrm{\pounds }_{}^{}\end{array}}{lim\; sup}(1p)Exp_{\lambda ,G}^2(p,q)<\mathrm{}.$$
(29)
For example, let us demonstrate (29). Without loss of generality, we replace (8) with
$$f_{}^{}(t)=O(t^2),g_{}^{}(t)=O(t^2)$$
(30)
(this can be done since $`f_{}^{}(t)f_+^{}(t)f_{}^{}(t+1)`$ and $`g_{}^{}(t)g_+^{}(t)g_{}^{}(t+1)`$). Consider the decision
$$(X,Y):=(\lambda (0,G(p,q)),\lambda (1,G(p,q))).$$
Since $`\frac{1p}{p}`$ is a subgradient (see, e.g., , Section 23) of $`f(x)`$ at $`X`$, (30) implies that $`1p=O(X^2)`$, i.e., $`(1p)X^2=O(1)`$. Since $`|Exp_{\lambda ,G}(p,q)|=XYXC_1`$ for $`(p,q)<(1,0)`$ sufficiently close to $`(1,0)`$, (29) indeed holds.
### Proof: Part III
In this subsection we consider the remaining possibilities for $`domG`$. Let us define the *left-stripped lexicographic ALN* ($`{}_{}{}^{}\mathrm{\pounds }`$ALN for brief) as the lexicographic ALN except that:
* its parameter is a forecast-continuous kernel $`𝐊`$ on $`{}_{}{}^{}\mathrm{\pounds }\times 𝐗`$, where the *left-stripped lexicographic square*
$${}_{}{}^{}\mathrm{\pounds }:=(0,1]\times [0,1]$$
is equipped with the order and topology inherited from $`\mathrm{\pounds }`$;
* it outputs a root $`(p,q){}_{}{}^{}\mathrm{\pounds }`$ of the equation $`S_n(p,q)=0`$ as $`(p_n,q_n)`$; if this equation does not have roots in $`{}_{}{}^{}\mathrm{\pounds }`$, we set $`p_n:=1`$ and set $`q_n[0,1]`$ arbitrarily (we will make sure that this happens only when $`S_n`$ is everywhere positive).
In a similar way we define the *right-stripped lexicographic square* $`\mathrm{\pounds }^{}`$ and the *right-stripped lexicographic ALN* ($`\mathrm{\pounds }^{}`$ALN), which always outputs $`(p_n,q_n)\mathrm{\pounds }^{}`$; when $`S_n(p,q)=0`$ does not have roots $`(p,q)\mathrm{\pounds }^{}`$ we now set $`p_n:=0`$.
We only consider the case $`domG={}_{}{}^{}\mathrm{\pounds }`$ (the case $`domG=\mathrm{\pounds }^{}`$ is treated analogously); this corresponds to $`f(C_0)=\mathrm{}`$ and $`f(C_1)<\mathrm{}`$. Since $`S_n`$ is continuous, the absence of roots of $`S_n=0`$ in $`{}_{}{}^{}\mathrm{\pounds }`$ in conjunction with (27) means that $`S_n`$ is positive everywhere on $`{}_{}{}^{}\mathrm{\pounds }`$, and so setting $`p_n:=1`$ in this case guarantees that $`{}_{}{}^{}\mathrm{\pounds }`$ALN still ensures (18). It remains to notice that (28) still holds.
## 7 The algorithm
In this short section we extract the decision strategy achieving (9) from our proof of Theorem 1. As we have already noticed (see (25)),
$$\begin{array}{c}S_n(p,q)=\underset{i=1}{\overset{n1}{}}\left(Exp_{\lambda ,G}(p,q)Exp_{\lambda ,G}(p_i,q_i)+𝐊(x_n,x_i)\right)(y_ip_i)\hfill \\ \hfill +\frac{1}{2}\left(Exp_{\lambda ,G}^2(p,q)+𝐊(x_n,x_n)\right)(12p);\end{array}$$
(31)
this immediately leads to the following explicit description.
An algorithm achieving (9)
| Parameters: | game with loss function $`\lambda `$ and canonical choice function $`G`$; |
| --- | --- |
| | kernel $`𝐊`$ on $`𝐗`$ |
FOR $`n=1,2,\mathrm{}`$:
Read $`x_n𝐗`$.
Define $`S_n(p,q)`$ by (31) for all $`(p,q)\mathrm{\pounds }`$ for which $`G(p,q)`$ is defined.
Define $`(p_n,q_n)`$ as any root $`(p,q)`$ of $`S_n(p,q)=0`$;
if there are no roots,
set $`p_n:=(1+signS_n)/2`$ and set $`q_n`$ to any number in $`[0,1]`$.
Set $`\gamma _n:=G(p_n,q_n)`$.
Read $`y_n\{0,1\}`$.
END FOR.
(We saw in the previous section that $`signS_n`$ is well defined and is $`1`$ or $`1`$ in this context.)
The canonical choice functions for the three examples of games given in §3 are as follows: $`G(p,q)=p`$ for the square loss and log loss games, and
$$G(p,q)=\{\begin{array}{cc}0\hfill & \text{if }p<1/2\hfill \\ 1\hfill & \text{if }p>1/2\hfill \\ q\hfill & \text{if }p=1/2\hfill \end{array}$$
(32)
for the absolute loss game.
## 8 Directions of further research
In this section we discuss informally what we consider to be interesting directions of further research.
### Non-convex games
Theorem 1 assumes that the superdecision set is convex. The assumption of convexity is convenient but not indispensable. We will only discuss the simplest non-convex game.
The loss function for the *simple loss game* is the same as for the absolute loss game, $`\lambda (y,\gamma )=|y\gamma |`$, but $`\mathrm{\Gamma }=\{0,1\}`$. Now the approach we have used in this paper does not work: since $`\mathrm{\Gamma }`$ consists of two elements, there is no non-trivial continuous choice function $`G:\mathrm{\pounds }\mathrm{\Gamma }`$ (every continuous image of $`\mathrm{\pounds }`$ is connected: , Theorem 6.1.4).
A natural idea () is to allow Decision Maker to use randomization. The expected loss of a strategy making decision $`1`$ with probability $`\gamma `$ and $`0`$ with probability $`1\gamma `$ is $`|y\gamma |`$, where $`y`$ is the actual observation; therefore, for the simple loss game a randomized decision strategy can guarantee the following analogue of (13):
$$\underset{n=1}{\overset{N}{}}𝔼|y_n\gamma _n|\underset{n=1}{\overset{N}{}}\left|y_nD(x_n)\right|+\frac{\sqrt{6}}{4}\left(2D1_𝒮+1\right)\sqrt{N},$$
(33)
where $`𝔼`$ refers to the strategy’s internal randomization (the decision rules $`D`$ can be allowed to take values in $`[0,1]`$).
The disadvantage of (33) is that typically we are interested in the strategy’s actual rather than expected loss. Our derivation of (33) shows the role of randomization: with our choice function (32) no randomization is required unless $`p=1/2`$. Typically, we rarely find ourselves in a situation of complete uncertainty, $`p_n=1/2`$; therefore, only a little bit of randomization is needed, essentially for tie breaking. The actual loss will be very close to the expected loss. It would be interesting to derive formal statements along these lines.
### Non-binary observations
It would also be interesting to extend this paper’s results to more general observation spaces (first of all, to carry them over to least-squares regression and multi-class classification). The two apparent obstacles to such extensions are that the fundamental equality (23) looks tailored to the binary case $`y\{0,1\}`$ and that Lemma 1 ceases to be obvious outside the binary case. However, (23) only states, in the terminology of , that $`\lambda (p,\gamma )`$ is the game-theoretic expected value of $`\lambda (y,\gamma )`$ (and that reproducing $`\lambda (y,\gamma )`$ given $`\lambda (p,\gamma )`$ can be accomplished by buying $`\lambda (1,\gamma )\lambda (0,\gamma )`$ tickets paying $`y`$ and costing $`p`$ each). Similar equalities hold for many other forecasting protocols. And an analogue of Lemma 1 for a wide class of forecasting protocols is proved in .
### Optimality
An important problem is to investigate the optimality of our algorithm, described in §7: is the bound (9) tight? (The tightness of the bounds in Theorem 2 and Equation (20) is established in .)
### Acknowledgments
This work was partially supported by MRC (grant S505/65) and Royal Society. |
warning/0506/hep-ph0506023.html | ar5iv | text | # 1 Introduction
## 1 Introduction
It is widely believed that conformal field theory is dynamically realized in a large class of non-abelian gauge theories with a certain number of matter multiplets (see Ref.). Conformal gauge theory is very attractive in the phenomenological point of view, since if it includes a SUSY-breaking sector, conformal sequestering of the SUSY breaking may occur, providing a solution to the flavor-changing neutral current (FCNC) problem in the supersymmetric standard model (SSM).<sup>1</sup><sup>1</sup>1See Ref. for some other phenomenological applications of superconformal dynamics. It is tempting to consider vector-like gauge theories for the SUSY breaking, since they are naturally incorporated into vector-like superconformal gauge theories, which are relatively well understood.
In this paper we extend vector-like gauge theories for the SUSY breaking by adding massive hyperquarks to turn the full high-energy theory above the mass threshold into conformal gauge theory. We find, however, that this simple extension does not achieve the conformal sequestering due to the presence of an unwanted global $`U(1)`$ symmetry. To eliminate the unwanted global symmetry we introduce non-abelian gauge interactions acting on the additional massive hyperquarks. We find various examples realizing the sequestering.
We first discuss $`SP(3N+1)\times SP(N)^6`$ gauge theories where all gauge coupling constants at the infrared fixed point are small for $`N>1`$ and perturbative calculations are applicable. We show by an explicit one-loop calculation that the theories have non-trivial fixed points and the sequestering of the SUSY-breaking effects indeed occurs. However, the sequestering is too mild to be applied to the phenomenology, since all the couplings are weak. Therefore, we dwell on strongly coupled conformal gauge theory such as an $`SP(3)\times SP(1)^2`$ theory in this paper.<sup>2</sup><sup>2</sup>2We are unable to prove explicitly that such a theory has a non-trivial infrared fixed point and the required sequestering is obtained, since gauge couplings are all strong. We only state, in this paper, why we expect that is the case.
We also propose a Planck-suppressed gauge mediation which circumvents the tachyonic mass problem for sleptons in anomaly mediation.<sup>3</sup><sup>3</sup>3This construction is essentially independent of the above model of conformal SUSY breaking and serves as a generic way to make anomaly mediation phenomenologically viable. Owing to the gravitational nature of this gauge mediation, the size of the gauge-mediated SUSY breaking is at most comparable to the anomaly-mediation effects. For the lowest messenger scale, the total model provides a hybrid scheme of anomaly and gauge mediations of SUSY breaking.
## 2 Conformal SUSY breaking
The IYIT model for SUSY breaking is based on an $`SP(N)`$ gauge theory with $`2(N+1)`$ chiral superfields (hyperquarks), $`Q_\alpha ^i`$, in the fundamental ($`2N`$-dimensional) representation.<sup>4</sup><sup>4</sup>4We adopt the notation where $`SP(1)=SU(2)`$. Here, $`\alpha =1,\mathrm{},2N`$ and $`i=1,\mathrm{},2(N+1)`$. We introduce $`(N+1)(2N+1)`$ gauge singlet chiral superfields, $`S_{ij}(=S_{ji})`$, and impose the flavor $`SU(2N+2)`$ symmetry in the superpotential,
$$W=\lambda S_{ij}Q^iQ^j,$$
(1)
where $`S_{ij}`$ are assumed to transform as an antisymmetric $`(N+1)(2N+1)`$ representation of the flavor $`SU(2N+2)`$ and we omit the color indices for simplicity. The reason why we impose the $`SU(2N+2)`$ symmetry becomes clear in the next section.
The effective low-energy superpotential is given by
$$W_{\mathrm{eff}}=X(\mathrm{Pf}V^{ij}\mathrm{\Lambda }^{2(N+1)})+\lambda S_{ij}V^{ij},$$
(2)
in terms of gauge invariant low-energy degrees of freedom $`V^{ij}Q^iQ^j`$. Here, $`X`$ is an additional chiral superfield and $`\mathrm{\Lambda }`$ denotes a dynamical scale of the $`SP(N)`$ gauge interaction. We see that the superfields $`S^{ij}`$ have non-vanishing $`F`$ terms in the vacuum and the SUSY is spontaneously broken. Notice here that the model possesses a $`U(1)_R`$ symmetry in addition to the flavor $`SU(2N+2)`$.
### 2.1 conformality
Now let us introduce $`2n_F`$ massive hyperquarks, $`Q^k`$, where $`k=1,\mathrm{},2n_F`$. The mass term is written as
$$W_{\mathrm{mass}}=\underset{i}{}mQ^iQ^{i+n_F}.$$
(3)
Here, $`i`$ runs from 1 to $`n_F`$. Above this mass scale, the high-energy theory is an $`SP(N)`$ gauge theory with $`N_F=2(N+1)+2n_F`$ hyperquarks. The SUSY $`SP(N)`$ gauge theory with $`N_F`$ hyperquarks is expected to be scale-invariant in the infrared for $`3(N+1)<N_F<6(N+1)`$ .
We check, in the following, that the theory with the superpotential Eq.(1) can also be scale-invariant in the infrared. The NSVZ beta function relates the running of the canonical gauge coupling constant to the anomalous dimension factors, $`\gamma _Q`$ and $`\gamma _Q^{}`$, of the hyperquarks, $`Q`$ and $`Q^{}`$, as
$`\mu {\displaystyle \frac{d}{d\mu }}\alpha _g=\alpha _g^2\left[{\displaystyle \frac{3(N+1)(N+1)(1\gamma _Q)n_F(1\gamma _Q^{})}{2\pi (N+1)\alpha _g}}\right],`$ (4)
where $`\alpha _g`$ is defined in terms of the gauge coupling constant $`g`$ of $`SP(N)`$ as $`\alpha _g=g^2/(4\pi )`$ and $`\mu `$ denotes the renormalization scale. Here and hereafter in this section, we neglect the masses of the hyperquarks $`Q^k`$. The beta function of the Yukawa coupling constant in Eq.(1) is also given in terms of the anomalous dimension factors of the hyperquarks, $`\gamma _Q`$, and of the singlet chiral fields, $`\gamma _S`$, by
$`\mu {\displaystyle \frac{d}{d\mu }}\alpha _\lambda =\alpha _\lambda (\gamma _S+2\gamma _Q),`$ (5)
where $`\alpha _\lambda `$ is defined in terms of the Yukawa coupling constant $`\lambda `$ as $`\alpha _\lambda =\lambda ^2/(4\pi )`$.
When the theory is scale-invariant with non-vanishing coupling constants, the beta functions in Eqs.(4) and (5) vanish. That is, we have, at the infrared fixed point,
$`3(N+1)(N+1)(1\gamma _Q)n_F(1\gamma _Q^{})=0,`$ (6)
$`\gamma _S+2\gamma _Q=0.`$ (7)
These conditions determine the anomalous dimensions at the fixed point. The anomalous dimensions at the fixed point are consistent with the unitarity of the theory for
$`1\gamma _Q,1\gamma _Q^{},0\gamma _S,`$ (8)
which comes from the restriction for unitary representation of the superconformal algebra :<sup>5</sup><sup>5</sup>5 Combining Eqs.(7) and (8), we also obtain $`\gamma _Q0`$ and $`\gamma _S2`$. The asymptotic freedom of $`SP(N)`$, namely, $`n_F<2(N+1)`$, results in $`\gamma _Q^{}<\gamma _Q/21/2`$ from Eqs.(6) and (8). the above anomalous dimensions are consistent with the unitarity conditions for any gauge-singlet chiral multiplets such as $`QQ`$, $`Q^{}Q^{}`$, and $`S`$.
Notice that the vanishing of the NSVZ beta function is consistent with the existence of the anomaly free $`U(1)_R`$ symmetry that enters in the superconformal algebra with the charges of the matter fields given by $`R_i=(2+\gamma _i)/3`$; $`i=Q,Q^{},S`$ . In this simple extension of the IYIT model, the anomalous dimensions cannot be determined uniquely from Eqs.(6) and (7), and hence, the charge assignment of the $`U(1)_R`$ is not determined only with this information.
Now, we show by a perturbative calculation that the fixed point is infrared stable. We first see that the gauge and Yukawa coupling constants at the infrared fixed point are small if $`N_F`$ is just below $`6(N+1)`$, as in the case of the Banks-Zaks fixed point . In this case we can obtain the anomalous dimensions at the one-loop level as
$`\gamma _Q`$ $`=`$ $`{\displaystyle \frac{2N+1}{2\pi }}\alpha _\lambda {\displaystyle \frac{2N+1}{4\pi }}\alpha _g,`$ (9)
$`\gamma _Q^{}`$ $`=`$ $`{\displaystyle \frac{2N+1}{4\pi }}\alpha _g,`$ (10)
$`\gamma _S`$ $`=`$ $`{\displaystyle \frac{2N}{2\pi }}\alpha _\lambda .`$ (11)
For $`n_F=2(N+1)\epsilon `$, we determine the coupling constants at the fixed point from Eqs.(6) and (7), as
$`\alpha _g^{}`$ $`=`$ $`{\displaystyle \frac{4\pi \epsilon }{7N^2+9N+2}}\left({\displaystyle \frac{3N+1}{2N+1}}\right)\left(1+𝒪\left({\displaystyle \frac{\epsilon }{N}}\right)\right),`$ (12)
$`\alpha _\lambda ^{}`$ $`=`$ $`{\displaystyle \frac{2\pi \epsilon }{7N^2+9N+2}}\left(1+𝒪\left({\displaystyle \frac{\epsilon }{N}}\right)\right),`$ (13)
and the one-loop approximation is justified a posteriori for small $`\epsilon /N`$.<sup>6</sup><sup>6</sup>6 A non-perturbative determination of the coupling constants through $`a`$-maximization is given in the Appendix A.
We can explicitly examine the infrared stability of the fixed point by considering the renormalizaiton group (RG) evolutions near the fixed point. The RG equations of the small deviations,
$$\mathrm{\Delta }\alpha _g\alpha _g\alpha _g^{},\mathrm{\Delta }\alpha _\lambda \alpha _\lambda \alpha _\lambda ^{},$$
(14)
are given by
$`\mu {\displaystyle \frac{d}{d\mu }}\mathrm{\Delta }\alpha _g`$ $`=`$ $`{\displaystyle \frac{\beta _g}{\alpha _g}}|_{}\mathrm{\Delta }\alpha _g+{\displaystyle \frac{\beta _g}{\alpha _\lambda }}|_{}\mathrm{\Delta }\alpha _\lambda ,`$ (15)
$`\mu {\displaystyle \frac{d}{d\mu }}\mathrm{\Delta }\alpha _\lambda `$ $`=`$ $`{\displaystyle \frac{\beta _\lambda }{\alpha _g}}|_{}\mathrm{\Delta }\alpha _g+{\displaystyle \frac{\beta _\lambda }{\alpha _\lambda }}|_{}\mathrm{\Delta }\alpha _\lambda ,`$ (16)
where $`\beta _G`$ and $`\beta _\lambda `$ denote the beta functions of $`\alpha _g`$ and $`\alpha _\lambda `$ given by Eqs.(4) and (5), respectively, and the values with the subscript “$``$” are evaluated at the fixed point. By using Eqs.(9)-(13), we find that all the eigenvalues of the coefficient matrix $`\{\beta _k/\alpha _l\}`$ are positive at the fixed point in Eqs.(12) and (13), where $`k,l=g,\lambda `$. Therefore, the fixed point in Eqs.(12) and (13) is infrared stable at least against small deviations from the fixed point.<sup>7</sup><sup>7</sup>7 Similar situations of the conformal fixed point with non-trivial Yukawa interactions are discussed in Refs. .
### 2.2 non-sequestering
We are at the point to show that the sequestering of the SUSY breaking does not occur due to an unwanted global $`U(1)`$ symmetry in this simple extension. By following Luty and Sundrum , we consider the RG evolutions of the wave function renormalization factors near the fixed point,
$`{\displaystyle \frac{d}{dt}}\mathrm{\Delta }\mathrm{ln}Z_i=\gamma _i+\gamma _i^{},\mathrm{\Delta }\mathrm{ln}Z_i\mathrm{ln}Z_i+\gamma _i^{}t,t\mathrm{ln}(\mu /M_{}),`$ (17)
where $`i=Q,Q^{},S`$ and $`\gamma _i^{}`$ are the anomalous dimensions at the fixed point given by Eqs.(9)-(11). Here, $`M_{}`$ denotes the scale where the theory enters the conformal regime below the reduced Planck scale $`M_G2.4\times 10^{18}`$ GeV. The deviations from the fixed point can be parameterized by $`\mathrm{\Delta }\alpha _g`$ and $`\mathrm{\Delta }\alpha _\lambda `$ which, in turn, can be expressed as<sup>8</sup><sup>8</sup>8 Without loss of generality, we adopt the convention of the holomorphic gauge coupling in Ref..
$`\mathrm{\Delta }\alpha _g`$ $`=`$ $`{\displaystyle \frac{\alpha _g^2}{2\pi (N+1)\alpha _g}}|_{}((N+1)\mathrm{\Delta }\mathrm{ln}Z_Q+n_F\mathrm{\Delta }\mathrm{ln}Z_Q^{}),`$ (18)
$`\mathrm{\Delta }\alpha _\lambda `$ $`=`$ $`\alpha _\lambda ^{}(2\mathrm{\Delta }\mathrm{ln}Z_Q+\mathrm{\Delta }\mathrm{ln}Z_S).`$ (19)
By using the above expressions, we rewrite the RG equation Eq.(17) as
$`{\displaystyle \frac{d}{dt}}\mathrm{\Delta }\mathrm{ln}Z_i`$ $`=`$ $`\left({\displaystyle \frac{\gamma _i}{\alpha _g}}\right)|_{}\mathrm{\Delta }\alpha _g\left({\displaystyle \frac{\gamma _i}{\alpha _\lambda }}\right)|_{}\mathrm{\Delta }\alpha _\lambda `$ (21)
$`=`$ $`\left(\left({\displaystyle \frac{\gamma _i}{\alpha _g}}\right){\displaystyle \frac{\alpha _g^2}{2\pi (N+1)\alpha _g}}\right)|_{}((N+1)\mathrm{\Delta }\mathrm{ln}Z_Q+n_F\mathrm{\Delta }\mathrm{ln}Z_Q^{})`$
$`+\left(\left({\displaystyle \frac{\gamma _i}{\alpha _\lambda }}\right)\alpha _\lambda \right)|_{}(2\mathrm{\Delta }\mathrm{ln}Z_Q+\mathrm{\Delta }\mathrm{ln}Z_S),`$
and we define the coefficient matrix $`L_{ij}`$ by
$`{\displaystyle \frac{d}{dt}}\mathrm{\Delta }\mathrm{ln}Z_i={\displaystyle \underset{j=Q,Q^{},S}{}}L_{ij}\mathrm{\Delta }\mathrm{ln}Z_j.`$ (22)
When all the eigenvalues of $`L`$ are positive, all $`\mathrm{\Delta }\mathrm{ln}Z_i`$ go to zero as $`t\mathrm{}`$ (the infrared limit) and hence the SUSY breaking is sequestered . Unfortunately, we find that the coefficient matrix $`L`$ has a zero eigenvalue. Thus, one linear combination of $`\mathrm{\Delta }\mathrm{ln}Z_i`$ is constant in the course of the RG evolution and it is not suppressed at the infrared fixed point. We call it as $`\mathrm{\Delta }\mathrm{ln}\overline{Z}`$. Since the vanishing eigenvalue corresponds to the eigenvector $`(\mathrm{\Delta }\mathrm{ln}Z_Q,\mathrm{\Delta }\mathrm{ln}Z_Q^{},\mathrm{\Delta }\mathrm{ln}Z_S)=(1,(N+1)/n_F,2)`$, we find that the solution to the Eq.(22) in the infrared limit is
$`\mathrm{\Delta }\mathrm{ln}Z_Q`$ $``$ $`(\mathrm{\Delta }\mathrm{ln}\overline{Z})_0,`$ (23)
$`\mathrm{\Delta }\mathrm{ln}Z_Q^{}`$ $``$ $`{\displaystyle \frac{N+1}{n_F}}(\mathrm{\Delta }\mathrm{ln}\overline{Z})_0,`$ (24)
$`\mathrm{\Delta }\mathrm{ln}Z_S`$ $``$ $`2(\mathrm{\Delta }\mathrm{ln}\overline{Z})_0,`$ (25)
with an $`𝒪(1)`$ proportionality factor, where $`(\mathrm{\Delta }\mathrm{ln}\overline{Z})_0`$ denotes the value at $`t=0`$. In general, the initial value $`(\mathrm{\Delta }\mathrm{ln}\overline{Z})_0`$ contains visible sector superfields $`q_i`$ as weakly coupled spectators such as
$`(\mathrm{\Delta }\mathrm{ln}\overline{Z})_0{\displaystyle \frac{\kappa _{ab}}{M_G^2}}q_a^{}q_b,`$ (26)
where $`\kappa _{ab}`$ denote $`O(1)`$ coefficients. Therefore, from Eq.(25), we find that the SUSY breaking effects to the visible sector are not sequestered.<sup>9</sup><sup>9</sup>9In Eqs.(23)-(25), we assume that other eigenvalues of $`L`$ are positive. Even if it is not the case, the conclusion is not changed.
The reason of our failure can be traced to the existence of a global $`U(1)`$ symmetry under which the SUSY breaking superfield $`S_{ij}`$ transforms non-trivially. In general, when an anomaly-free (non-$`R`$) $`U(1)`$ symmetry exists, the charge assignment $`\omega _i`$ determines the eigenvector of the coefficient matrix $`L_{ij}`$ for a vanishing eigenvalue:
$`{\displaystyle \underset{j}{}}L_{ij}\omega _j=0.`$ (27)
In the present case, a linear combination $`\mathrm{\Delta }\mathrm{ln}\overline{Z}`$ (of $`\mathrm{\Delta }\mathrm{ln}Z_i`$) remains constant in the infrared limit and the SUSY breaking effects are not sequestered if the SUSY breaking superfields have non-vanishing charges. The eigenvector we have found above corresponds to the charge assignment $`(\omega _Q,\omega _Q^{},\omega _S)=(1,(N+1)/n_F,2)`$ of an anomaly-free $`U(1)`$ symmetry. Thus, in order to realize the sequestering, we should violate the global $`U(1)`$ symmetry under which the SUSY breaking superfields transform non-trivially, provided we do not take $`(\mathrm{\Delta }\mathrm{ln}\overline{Z})_0=0`$ by fine tuning. In the next section, we introduce additional gauge symmetries, where the unwanted $`U(1)`$ symmetry is broken by anomaly due to the new gauge interactions.
## 3 Conformally sequestered extensions
We introduce gauge interactions acting on the massive hyperquarks, $`Q^k`$, where the unwanted global $`U(1)`$ symmetry is broken by anomaly due to the new gauge interactions. We deal, in this section, with $`SP(N)\times SP(N^{})^6`$, $`(N=3N^{}+1)`$ gauge theory, where the former $`SP(N)`$ corresponds to the gauge group for the SUSY breaking and the latter $`SP(N^{})^6`$ gauge group is introduced to break the unwanted $`U(1)`$ symmetry. We list all the matter contents in Table 1. We take such a large gauge group to see explicitly by a perturbative calculation that the conformal sequestering occurs. Indeed all the couplings at the infrared fixed point are weak for $`N^{}>1`$ in the present model.
### 3.1 conformality
Now, we check that the theory with the extended gauge symmetry can be scale-invariant in the infrared. In this model, the beta functions of the $`SP(N)`$ gauge coupling constant $`\alpha _g`$, the $`SP(N^{})`$ gauge coupling constant $`\alpha _g^{}`$, and the Yukawa coupling constant $`\alpha _\lambda `$ in Eq.(1) are given by
$`\mu {\displaystyle \frac{d}{d\mu }}\alpha _g`$ $`=`$ $`\alpha _g^2\left[{\displaystyle \frac{3(N+1)(N+1)(1\gamma _Q)6N^{}(1\gamma _Q^{})}{2\pi (N+1)\alpha _g}}\right],`$ (28)
$`\mu {\displaystyle \frac{d}{d\mu }}\alpha _g^{}`$ $`=`$ $`\alpha _g^{}^2\left[{\displaystyle \frac{3(N^{}+1)N(1\gamma _Q^{})}{2\pi (N^{}+1)\alpha _g^{}}}\right],`$ (29)
$`\mu {\displaystyle \frac{d}{d\mu }}\alpha _\lambda `$ $`=`$ $`\alpha _\lambda (\gamma _S+2\gamma _Q),`$ (30)
where we have assumed that all the $`SP(N^{})`$ sectors are equivalent. Namely, we have imposed an exchange symmetry between any two $`SP(N^{})`$’s in the $`SP(N^{})^6`$ so that the $`SP(N^{})^6`$ has a common gauge coupling constant $`\alpha _g^{}`$. Then, by requiring all the beta functions to vanish, we determine the anomalous dimensions uniquely as
$`\gamma _Q`$ $`=`$ $`{\displaystyle \frac{2(N(N+1)9N^{}(N^{}+1))}{N(N+1)}}={\displaystyle \frac{4}{9N^2+9N^{}+2}},`$ (31)
$`\gamma _Q^{}`$ $`=`$ $`{\displaystyle \frac{N3(N^{}+1)}{N}}={\displaystyle \frac{2}{3N^{}+1}},`$ (32)
$`\gamma _S`$ $`=`$ $`2\gamma _Q.`$ (33)
Here, we have neglected the masses of the hyperquarks $`Q^{}`$.
We also determine the coupling constants at the infrared fixed point by a perturbative calculation. The anomalous dimensions at the one-loop level are given by
$`\gamma _Q`$ $`=`$ $`{\displaystyle \frac{2N+1}{2\pi }}\alpha _\lambda {\displaystyle \frac{2N+1}{4\pi }}\alpha _g,`$ (34)
$`\gamma _Q^{}`$ $`=`$ $`{\displaystyle \frac{2N^{}+1}{4\pi }}\alpha _g^{}{\displaystyle \frac{2N+1}{4\pi }}\alpha _g,`$ (35)
$`\gamma _S`$ $`=`$ $`{\displaystyle \frac{2N}{2\pi }}\alpha _\lambda .`$ (36)
Then, Eqs.(31)-(36) determine the coupling constants at the infrared fixed point by
$`{\displaystyle \frac{N+1}{2\pi }}\alpha _g^{}`$ $`=`$ $`{\displaystyle \frac{8(9N^{}+4)}{3(2N^{}+1)(3N^{}+1)^2}},`$ (37)
$`{\displaystyle \frac{N^{}+1}{2\pi }}\alpha _g^{}^{}`$ $`=`$ $`{\displaystyle \frac{12(N^{}+1)(3N^23N^{}2)}{(2N^{}+1)(3N^{}+2)(3N^{}+1)^2}},`$ (38)
$`{\displaystyle \frac{N}{\pi }}\alpha _\lambda ^{}`$ $`=`$ $`{\displaystyle \frac{8}{(3N^{}+2)(3N^{}+1)}}.`$ (39)
We see that all the coupling constants are small and the perturbative calculation is reliable for $`N^{}>1`$.<sup>10</sup><sup>10</sup>10 For $`N^{}=1`$, although the anomalous dimensions in Eq.(33) satisfy the unitarity bound Eq.(8), the gauge coupling constants of $`SP(N^{})`$ in Eq.(37) turns out to be negative, which implies that the perturbative description is invalid.
The above result enables us to explicitly analyze the infrared stability of the fixed point in the same way as done in the previous section. The RG equations of the small deviations $`\mathrm{\Delta }\alpha _k\alpha _k\alpha _k^{},(k=g,g^{},\lambda )`$ are given by
$`\mu {\displaystyle \frac{d}{d\mu }}\mathrm{\Delta }\alpha _k={\displaystyle \underset{l=g,g^{},\lambda }{}}M_{kl}\mathrm{\Delta }\alpha _l,`$ (40)
where the coefficient matrix $`M`$ is defined by
$`M_{kl}={\displaystyle \frac{\beta _k}{\alpha _l}}|_{}.`$ (41)
Here, “$``$” indicates the values evaluated at the fixed point. If all the eigenvalues of the coefficient matrix are positive, the fixed point in Eqs.(37)-(39) is infrared stable. In Table 2, we show numerical results on the eigenvalues of the matrix $`M`$ for the case of $`N^{}>1`$. ¿From the table, we see that eigenvalues are all positive for $`N^{}>1`$. Therefore, we find that the $`SP(3N^{}+1)\times SP(N^{})^6`$ gauge theory has the stable infrared fixed point in Eqs.(31)-(33).
### 3.2 sequestering
We now discuss the sequestering of the SUSY breaking. The RG equations of the wave function renormalization factors $`Z_i`$ near the fixed point are given by
$`{\displaystyle \frac{d}{dt}}\mathrm{\Delta }\mathrm{ln}Z_i`$ $`=`$ $`{\displaystyle \underset{k=g,g^{},\lambda }{}}\left({\displaystyle \frac{\gamma _i}{\alpha _k}}\right)|_{}\mathrm{\Delta }\alpha _k`$ (42)
$`=`$ $`{\displaystyle \underset{i=Q,Q^{},S}{}}L_{ij}\mathrm{\Delta }\mathrm{ln}Z_j.`$ (43)
Here, the coefficient matrix $`L`$ in the second line is given by using the following relations:
$`\mathrm{\Delta }\alpha _g`$ $`=`$ $`{\displaystyle \frac{\alpha _g^2}{2\pi (N+1)\alpha _g}}|_{}((N+1)\mathrm{\Delta }\mathrm{ln}Z_Q+6N^{}\mathrm{\Delta }\mathrm{ln}Z_Q^{}),`$ (44)
$`\mathrm{\Delta }\alpha _g^{}`$ $`=`$ $`{\displaystyle \frac{\alpha _g^{}^2}{2\pi (N^{}+1)\alpha _g^{}}}|_{}N\mathrm{\Delta }\mathrm{ln}Z_Q^{},`$ (45)
$`\mathrm{\Delta }\alpha _\lambda `$ $`=`$ $`\alpha _\lambda ^{}(2\mathrm{\Delta }\mathrm{ln}Z_Q+\mathrm{\Delta }\mathrm{ln}Z_S).`$ (46)
The sequestering of the SUSY breaking is realized when all the eigenvalues of $`L`$ are positive.
Interestingly, as we show below, the coefficient matrix $`L`$ has the same eigenvalues as the coefficient matrix $`M`$ in Eq.(41). Therefore, the sequestering occurs automatically, if the infrared fixed point determined in Eqs.(31)-(33) is stable. To prove that, we rewrite the conditions for the vanishing beta functions as
$`{\displaystyle \underset{j=Q,Q^{},S}{}}A_{kj}\gamma _j=b_k,`$ (47)
where the coefficient matrix $`A`$ and the vector $`b`$ can be read off from Eqs.(28)-(30) and $`k=g,g^{},\lambda `$. Then, we see the following relations:
$`M_{kl}`$ $`=`$ $`{\displaystyle \underset{j=Q,Q^{},S}{}}A_{kj}\mathrm{\Gamma }_{jl},`$ (48)
$`L_{ij}`$ $`=`$ $`{\displaystyle \underset{k=g,g^{},\lambda }{}}\mathrm{\Gamma }_{ik}A_{kj},`$ (49)
where we have defined
$`\mathrm{\Gamma }_{ik}\left({\displaystyle \frac{\gamma _i}{\alpha _k}}\right)|_{}.`$ (50)
Since the coefficient matrix $`A`$ is invertible, the coefficient matrices $`M`$ and $`L`$ are similar to each other, so that they have the same eigenvalues. Therefore, the sequestering occurs automatically when the anomalous dimensions are uniquely determined by the conditions for the vanishing beta functions (i.e. $`A`$ is invertible) and the fixed point is infrared stable (i.e. all the eigenvalues of $`M`$ are positive). Notice that this is no accident: the conformal sequestering originates from nothing but the attractor structure of the infrared fixed point.
In our $`SP(3N^{}+1)\times SP(N^{})^6`$ model, we have shown that the fixed point is determined from the conditions of vanishing beta functions and the fixed point is infrared stable for $`N^{}>1`$. Thus, we have found that the sequestering is realized in our model. Notice that the relation between the infrared stability and the sequestering holds independently of the perturbative calculation. Therefore, even if perturbative analysis is not applicable, we may argue that the sequestering occurs, if the fixed point is expected to be infrared stable.
It should be noted here that the unwanted global $`U(1)`$ symmetry discussed in the previous section is broken by anomalies of the $`SP(N^{})^6`$ gauge interactions and hence there is no conserved $`U(1)`$ current. This is the reason why the matrix $`M`$ does not have a zero eigenvalue.
In addition to the above global $`U(1)`$ symmetry, there are many unbroken global $`U(1)`$’s acting on the gauge singlet superfields $`S_{ij}`$, which consist of the $`U(1)`$ subgroups of the flavor $`SU(2N+2)`$ of the hyperquarks $`Q^i`$. Thus, there are many linear combinations of the wave function renormalization factors which are not sequestered in the course of the RG evolutions to the infrared fixed point. For example, a linear combination $`\mathrm{\Delta }\mathrm{ln}Z_{S_{12}}\mathrm{\Delta }\mathrm{ln}Z_{S_{34}}`$ is not sequestered, since this corresponds to the global $`U(1)SU(2N+2)`$ symmetry. Fortunately, we can make such non-sequestered combinations vanishing by imposing the flavor $`SU(2N+2)`$ symmetry (or a sufficiently large discrete subgroup thereof) at high energies so that the conformal sequestering of the SUSY breaking is realized. Namely, by assuming that the Kähler potential inducing soft masses for squarks and sleptons is restricted by the flavor $`SU(2N+2)`$ symmetry as
$$\frac{\kappa _{ab}}{M_G^2}\underset{ij}{}S_{ij}^{}S_{ij}q^aq^b,$$
(51)
we can set the linear combinations of $`\mathrm{\Delta }\mathrm{ln}Z`$’s which are not sequestered to be zero. Then, as we have discussed, the remaining combinations of $`\mathrm{\Delta }\mathrm{ln}Z`$’s are sequestered and the squared masses of the sfermions from Eq.(51) are suppressed at the infrared fixed point. This is the reason why we have imposed the flavor $`SU(2N+2)`$ symmetry in the SUSY-breaking sector.
Finally, in the rest of this section, we show that the sequestering is too mild in the present model to solve the FCNC problem. In view of the Table 2, the smallest eigenvalue $`\beta _{}^{}`$ of the coefficient matrix $`L`$ (or equivalently $`M`$) is of the order of $`10^3`$ for $`N^{}2`$. Thus, the linear combination of $`\mathrm{\Delta }\mathrm{ln}Z_i`$ that corresponds to the smallest eigenvalue approaches to the fixed point very slowly, which, in turn, prevents $`\mathrm{\Delta }\mathrm{ln}Z_i`$ from getting up to the fixed point immediately. That is, in the infrared regime ($`t0`$), we find
$`\mathrm{\Delta }\mathrm{ln}Z_S(t)`$ $``$ $`e^{\beta _{}^{}t}\mathrm{\Delta }(\mathrm{ln}\overline{Z})_0,`$ (52)
$`(\mathrm{\Delta }\mathrm{ln}\overline{Z})_0`$ $`=`$ $`c_S(\mathrm{\Delta }\mathrm{ln}Z_S)_0+c_Q(\mathrm{\Delta }\mathrm{ln}Z_Q)_0+c_Q^{}(\mathrm{\Delta }\mathrm{ln}Z_Q^{})_0,`$ (53)
where $`\mathrm{\Delta }\mathrm{ln}\overline{Z}`$ corresponds to the eigenvector for the smallest eigenvalue, the subscript “0” indicates the value at $`t=0`$, and $`c_i`$ denote numerical coefficients. By explicit calculation, we find that the coefficients $`c_i`$ are typically $`𝒪(0.10.01)`$ in our perturbative models. In order to solve the FCNC problem by sequestering, we should require $`\mathrm{\Delta }\mathrm{ln}Z_{S}^{}{}_{}{}^{<}10^7`$ at the SUSY-breaking scale .<sup>11</sup><sup>11</sup>11 Here, we assume that the flavor diagonal masses of the sfermions are of the order of $`1`$ TeV, which are suppressed compared to the gravitino mass of the order of $`100`$ TeV (see discussions in section 5). Thus, without fine tuning among $`\mathrm{ln}Z_i`$, we should require $`e_{}^{\beta _{}^{}t}{}_{}{}^{<}10^7`$. However, since $`\beta _{}^{}`$ is of the order of $`10^3`$, it takes too long to achieve the sufficient sequestering. Therefore, we find that the sufficient sequestering cannot be expected in our perturbative models.
In the perturbative examples, we have seen that the size of the “sequestering speed” $`\beta _{}^{}`$ is not larger than the anomalous dimensions at the fixed point. Thus, in order to realize the sufficient sequestering (i.e. $`\beta _{}^{}=𝒪(1)`$), we should require that the anomalous dimensions at the fixed point are of the order one.<sup>12</sup><sup>12</sup>12It is based on a naive expectation that the speed of the sequestering, $`\beta _{}^{}(\gamma /\alpha )\alpha |_{}`$ or $`(\gamma /\alpha )\alpha ^2|_{}`$, is not so far from $`\gamma _{}`$ even in the strongly coupled case (see Eqs.(43)-(46)). This means that we must consider a strongly coupled conformal gauge theory.<sup>13</sup><sup>13</sup>13Unfortunately, the strongly coupled case $`N^{}=1`$ also seems inadequate since the anomalous dimensions are not sufficiently large. In the next section, we discuss such a strongly coupled theory and present the reason why we consider the sequestering might be also realized there, although the perturbative calculation is not applicable.
## 4 Strongly coupled $`SP(3)\times SP(1)^2`$ model
In this section, we discuss $`SP(3)\times SP(1)^2`$ gauge theory as an example, where $`SP(3)`$ corresponds to the gauge group for the SUSY breaking and the $`SP(1)^2`$ gauge group acts on massive hyperquarks $`Q^k`$. We list the matter contents in Table 3. We assume that such a gauge theory with the Yukawa interaction in Eq.(1) has a non-trivial fixed point. Then, the anomalous dimensions for $`Q,Q^{}`$, and $`S`$ at the fixed point are determined as
$`\gamma _Q=1,\gamma _Q^{}=1,\gamma _S=2,`$ (54)
which sit on a boundary of the unitarity bound Eq.(8).<sup>14</sup><sup>14</sup>14 The reason we take this example is only for simplicity. In the phenomenological point of view, we only require a large size of the anomalous dimensions which satisfy the unitarity bound Eq.(8).
In Table 4, we list some other examples of the $`SP(N)\times SP(N^{})^2`$ gauge theories, which include the cases where the perturbative analysis is marginally applicable. In such examples with gauge symmetry structures similar to $`SP(3)\times SP(1)^2`$, we can explicitly check that the fixed points are infrared stable. Thus, based on these results (i.e. consistency with the unitarity and the presence of similar but calculable examples), we expect that the fixed point with Eq.(54) is infrared stable, although it is hard to check it by an explicit calculation.<sup>15</sup><sup>15</sup>15We know no calculable example that has an infrared unstable (non-trivial) fixed point in the present class of $`SP(N)\times SP(N^{})^n,(n=1,2,\mathrm{})`$ gauge theories.
Now, we discuss the sequestering of the SUSY breaking effects in the strongly coupled $`SP(3)\times SP(1)^2`$ model. As argued in the previous section, if the anomalous dimensions are determined uniquely from the conditions for the vanishing beta functions, then the sequestering is equivalent to the infrared stability of the fixed point. Hence, once we assume that the fixed point with Eq.(54) is infrared stable, the sequestering is guaranteed. By assuming that the “sequestering speed” $`\beta _{}^{}`$ is not so far from the values of $`\gamma _i`$ (see Eqs.(43)-(46)), we expect $`\beta _{}^{}=𝒪(1)`$ in our strongly coupled model. The flavor-changing soft masses are sufficiently suppressed by sequestering<sup>16</sup><sup>16</sup>16 In what follows, we assume that the theory is in the vicinity of the conformal fixed point at $`M_G`$, that is, we assume $`M_{}M_G`$ (see below Eq.(17)). at the energy scale $`\mu `$ as high as
$`\left({\displaystyle \frac{\mu }{M_G}}\right)_{}^<10^{\frac{7}{\beta _{}^{}}}.`$ (55)
When the RG scale $`\mu `$ comes close to the physical mass scale $`m_{\mathrm{phys}}`$ of $`Q^{}`$, the theory ceases to be scale invariant and effectively becomes an asymptotically free $`SP(3)`$ gauge theory of strong coupling with 8 hyperquarks $`Q`$, and finally SUSY is broken dynamically at $`\mu _{}^<m_{\mathrm{phys}}`$. Here, the physical mass is given by
$`m_{\mathrm{phys}}=(mM_G^{\gamma _Q^{}})^{\frac{1}{1\gamma _Q^{}}}=\sqrt{mM_G},`$ (56)
where the last equality results from the unitarity boundary value $`\gamma _Q^{}=1`$ in the present model. Thus, the above condition Eq.(55) for the sufficient sequestering can be rewritten in terms of the mass scale of $`Q^{}`$ or the SUSY breaking scale $`\mathrm{\Lambda }`$ as
$`\left({\displaystyle \frac{\mathrm{\Lambda }}{M_G}}\right)_{}^<\left({\displaystyle \frac{m_{\mathrm{phys}}}{M_G}}\right)_{}^<10^{\frac{7}{\beta _{}^{}}}.`$ (57)
As discussed in the next section, we are interested in the case where the gravitino mass is of the order of $`100`$ TeV, which implies $`\mathrm{\Lambda }10^{1112}`$ GeV. Thus, we claim that the above conditions can be satisfied for $`\beta _{}^{}=𝒪(1)`$, and hence, the FCNC can be suppressed in the present strongly coupled model.
## 5 Circumventing the tachyonic slepton problem
If the sequestering of SUSY breaking occurs sufficiently, the SUSY-breaking masses for the squarks and sleptons become negligibly small at low energies. In this situation we must invoke some mechanism to transmit sizable SUSY-breaking effects to the visible sector of the SSM. The most natural candidate is anomaly mediation . This mechanism is not only theoretically interesting, but also phenomenologically attractive. This is because the gravitino mass is expected at $`𝒪(100)`$ TeV, which provides us with a solution to the gravitino problem . However, the anomaly mediation mechanism suffers from the tachyonic slepton mass problem .
In this section we consider a Planck-suppressed gauge mediation to remedy this phenomenological defect of the anomaly mediation.<sup>17</sup><sup>17</sup>17 In the Appendix B, we also provide a renormalizable setup for such a remedy. Let us introduce a messenger sector which consists of chiral superfields $`\psi `$, $`\overline{\psi }`$, $`\psi ^{}`$, and $`\overline{\psi }^{}`$. Here, $`\psi ,\overline{\psi }`$ and $`\psi ^{},\overline{\psi }^{}`$ transform as vector-like representations under the gauge group of the SSM, and we take them to fit in complete $`SU(5)`$ GUT representations, $`\mathrm{𝟓}+\mathrm{𝟓}^{}`$, for simplicity.
Our additional superpotential terms are given by
$$\delta W=\frac{h}{M_G^2}S_{ij}Q^iQ^j\psi \overline{\psi }+m_m\psi \overline{\psi }^{}+m_m\psi ^{}\overline{\psi },$$
(58)
which let the SUSY breaking intact for a sufficiently large mass parameter $`m_m`$.<sup>18</sup><sup>18</sup>18In fact, $`m_{m}^{}{}_{}{}^{>}m_{3/2}`$ is required, where $`m_{3/2}`$ denotes the gravitino mass. Here, $`h`$ denotes a coupling constant of order one and the combination $`S_{ij}Q^iQ^j`$ stands just for a SUSY-breaking superpotential term which has a non-vanishing $`F`$ component (see Eq.(2)). The SUSY-breaking effects are transmitted to the sfermions and Higgs bosons by the SSM gauge interactions (see Ref.).
In the SUSY-breaking dynamics, we expect $`|S|_{}^<\mathrm{\Lambda }`$ , which yields only Planck-suppressed $`R`$ breaking effects in the gauge mediation. In this case, gauginos do not obtain sizable SUSY-breaking masses via the gauge mediation and the gaugino spectrum is virtually the same as in the purely anomaly-mediated one.<sup>19</sup><sup>19</sup>19 Since the superpotential has the constant term which is required to obtain the flat universe, we may as well introduce an interaction term between $`\psi \overline{\psi }`$ and the constant term. Then, the $`R`$ breaking effects in the gauge mediation possibly become sizable , which may result in the gaugino spectrum different from the one in the pure anomaly mediation. The expression of the gauge mediated mass squared in Eq.(60) may also be altered. On the other hand, the scalar field $`\varphi `$ obtains the mass squared via the gauge mediation for $`m_m<m_{\mathrm{phys}}`$ as
$`m_\varphi ^2`$ $``$ $`2{\displaystyle \underset{a=1,2,3}{}}C_a^\varphi \left({\displaystyle \frac{\alpha _a}{4\pi }}\right)^2{\displaystyle \frac{|hF_S|^2}{m_m^2}}\left({\displaystyle \frac{\mathrm{\Lambda }}{M_G}}\right)^4,`$ (59)
$``$ $`18{\displaystyle \underset{a=1,2,3}{}}C_a^\varphi \left({\displaystyle \frac{\alpha _a}{4\pi }}\right)^2\left({\displaystyle \frac{|h|m_{3/2}}{|\lambda |m_m}}\right)^2m_{3/2}^2,`$ (60)
where $`C_a^\varphi (a=1,2,3)`$ is the quadratic Casimir invariant for each gauge group relevant to the scalar $`\varphi `$.<sup>20</sup><sup>20</sup>20Here, we have neglected RG effects from the MSSM couplings. In the above equation, we have used $`\sqrt{F_S}\sqrt{\lambda }\mathrm{\Lambda }`$ and the gravitino mass $`m_{3/2}`$ given by
$`m_{3/2}{\displaystyle \frac{|F_S|}{\sqrt{3}M_G}}{\displaystyle \frac{|\lambda |\mathrm{\Lambda }^2}{\sqrt{3}M_G}}.`$ (61)
As a result, we find that the gauge-mediated masses squared are comparable to the anomaly mediated ones for $`m_mm_{3/2}`$.<sup>21</sup><sup>21</sup>21 For $`m_{m}^{}{}_{}{}^{<}m_{3/2}`$, even if the total SUSY-breaking were kept intact, messenger scalar particles would become tachyonic. Hence we restrict ourselves to $`m_{m}^{}{}_{}{}^{>}m_{3/2}`$. The choice $`m_mm_{3/2}`$ realizes the lowest-scale model of gauge mediation (see Ref.) for $`m_{3/2}`$ of order $`100`$ TeV. We note that $`m_mm_{3/2}`$ is realized by a relation $`m_mm`$ of the Lagrangian parameters in view of Eq.(56). In particular, the positive contributions to the slepton masses squared in Eq.(60),
$`m_{\stackrel{~}{e}\mathrm{G}.\mathrm{M}.}^2{\displaystyle \frac{3}{5}}18\left({\displaystyle \frac{|h|m_{3/2}}{|\lambda |m_m}}\right)^2\left({\displaystyle \frac{\alpha _1}{4\pi }}\right)^2m_{3/2}^2,`$ (62)
can overwhelm the negative contribution of the anomaly-mediated mass squared,
$`m_{\stackrel{~}{e}\mathrm{A}.\mathrm{M}.}^2{\displaystyle \frac{6}{5}}{\displaystyle \frac{33}{5}}\left({\displaystyle \frac{\alpha _1}{4\pi }}\right)^2m_{3/2}^2.`$ (63)
Therefore, we conclude that the tachyonic slepton problem is resolved in the total model of anomaly and gauge mediation hybrid by tunning scales of these two mediations with each other.
Note that the newly added superpotential term in Eq.(58) does not violate the global $`SU(2N+2)`$ symmetry which is relevant for the conformal sequestering. Hence, we can apply the above mechanism to the conformally sequestered models. However, we should note that in the case of conformally sequestered models, the first term in Eq.(58) is also sequestered in the course of the RG evolution from $`M_{}`$ to $`m_{\mathrm{phys}}`$.<sup>22</sup><sup>22</sup>22The sequestering can be seen through a field redefinition $`\stackrel{~}{S}_{ij}=(1+\lambda ^1h\psi \overline{\psi }/M_G^2)S_{ij}`$, which turns the effects of the superpotential coupling $`h`$ into those of the Kähler couplings appearing as perturbations to the renormalization factors. Thus, in order to realize the sizable gauge mediation effects as in Eq.(62), we need to compensate the sequestering effects by preparing the additional superpotential terms
$$\left(\frac{M_{}}{m_{\mathrm{phys}}}\right)^\beta _{}^{}\frac{h}{M_G^2}S_{ij}Q^iQ^j\psi \overline{\psi }+m_m\psi \overline{\psi }^{}+m_m\psi ^{}\overline{\psi }$$
(64)
at the scale $`M_{}`$ which effectively realize the Eq.(58) after the conformal sequestering. This implies that the higher-dimensional term stems from integrating out an intermediate matter of mass $`(m_{\mathrm{phys}}/M_{})^\beta _{}^{}M_G`$ with Planck-suppressed coupling to $`S_{ij}Q^iQ^j`$.<sup>23</sup><sup>23</sup>23 For example, we may consider a concrete model by introducing extra singlet supermultiplets $`X`$ and $`\overline{X}`$ with a superpotential $`{\displaystyle \frac{h_1}{M_G}}XSQQ+M_XX\overline{X}+h_2\overline{X}\psi \overline{\psi },`$ (65) at the scale $`M_{}`$. Here, $`h_{1,2}`$ denote coupling constants, $`M_X`$ the mass parameter of $`X`$. Then, after integrating out $`X`$ and $`\overline{X}`$, we can effectively obtain the first term in Eq.(58) at the scale $`m_{\mathrm{phys}}`$ for $`M_XM_G(m_{\mathrm{phys}}/M_{})^\beta _{}^{}((h_1h_2)/h)`$. In the above analysis, we have simply used Eq.(58) as a resultant effective superpotential at the scale $`m_{\mathrm{phys}}`$ for the conformally sequestered models.
Finally, we comment on the cosmologcal aspects of this class of models. Since the relic density of the lightest messenger particle is too much to be consistent with the observation, we should require that they decay into the SSM particles at early stage of the universe. This is implemented by introducing small mixings between the messenger particles and the SSM particles. In addition, it should also be noted that there are Goldstone bosons in the SUSY breaking sector, which correspond to the spontaneous breaking of the global $`SU(2N+2)`$ symmetry. However, those massless particles are decoupled from the thermal bath since they only couple with the SSM particles via the Planck-supressed operator, and hence they do not affect the history of the universe.<sup>24</sup><sup>24</sup>24Here, we assume that the inflaton decays dominantly to the SSM particles. Therefore, we find that the present hybrid scheme yields also a consistent scenario from the cosmological point of view.
## Acknowledgements
M.I. and Y.N thank the Japan Society for the Promotion of Science for financial support. The authors acknowledges the referee for useful comments.
## Appendix A Anomalous dimensions from $`a`$-maximization
Recently, Intriligator and Wecht proposed a powerful technique to compute the conformal $`R`$ current in a certain class of conformal field theories in four dimensions and hence the anomalous dimensions thereof . In this appendix we use this so-called $`a`$-maximization method to determine the anomalous dimensions of the fields in the conformally extended IYIT model beyond the Banks-Zaks approximation presented in section 3.
The $`a`$-maximization method simply states that the conformal $`R`$ current appearing in the superconformal algebra maximizes a particular t’Hooft anomaly
$$a=\mathrm{Tr}(3R^3R),$$
(66)
which is related to the conformal anomaly on a curved spacetime
$`{\displaystyle _{S^4}}T_\mu ^\mu .`$ (67)
In our model of the $`SP(N)`$ gauge theory, the candidate of the conformal $`R`$ current contains one free parameter $`x=\gamma _Q`$, from which the corresponding $`R`$ charges are determined by Eqs.(6) and (7) as
$`R_Q={\displaystyle \frac{2}{3}}(1+{\displaystyle \frac{x}{2}}),R_Q^{}={\displaystyle \frac{2}{3}}(1+{\displaystyle \frac{\gamma _Q^{}}{2}}),R_S={\displaystyle \frac{2}{3}}(1+{\displaystyle \frac{2x}{2}})`$ (68)
with
$$\gamma _Q^{}=1\frac{3(N+1)(N+1)(1x)}{2(N+1)\epsilon },$$
(69)
where $`\epsilon =2(N+1)n_F`$.
The claim is that among these one-parameter $`R`$ currents, the conformal one maximizes the anomaly $`a`$, which is obtained as follows:
$`a`$ $`=`$ $`2N(2N+2)\left[3(R_Q1)^3(R_Q1)\right]`$ (72)
$`+2(2N+2\epsilon )2N\left[3(R_Q^{}1)^3(R_Q^{}1)\right]`$
$`+(N+1)(2N+1)\left[3(R_S1)^3(R_S1)\right],`$
where we note that the $`R`$ charges appearing in $`a`$ are those of fermions (i.e. $`R_{\psi _Q}=R_Q1`$) because only fermions contribute to the anomaly. By maximizing $`a`$ with respect to $`x`$, we can determine $`x_{}=\gamma _Q|_{}`$. The unique local maximum is achieved by setting
$`x_{}`$ $`=`$ $`\left(\epsilon ^2(2+3N)4\epsilon (1+N)(2+3N)+(1+N)^2(8+13N)\right)^1A;`$ (73)
$`A`$ $``$ $`44\epsilon +\epsilon ^2+22N16\epsilon N+3\epsilon ^2N+32N^212\epsilon N^2+14N^3+(\epsilon 2(1+N))B,`$ (74)
$`B`$ $``$ $`\sqrt{\epsilon ^2(1+2N)(1+6N)4\epsilon (1+N)(1+2N)(1+6N)+(2+9N+7N^2)^2}.`$ (75)
To compare this rather complicated expression with the perturbative results, we expand Eq.(75) in terms of $`\epsilon `$. Remarkably, the first order approximation is given by
$$x_1=\frac{N}{2+9N+7N^2}\epsilon ,$$
(76)
which completely agrees with our Banks-Zaks-like calculation. Furthermore we can systematically study higher order corrections.<sup>25</sup><sup>25</sup>25For instance, the two loop contribution should be $`x_2={\displaystyle \frac{3N(1+N(7+11N))}{(1+N)^2(2+7N)^3}}\epsilon ^2.`$ It is quite intriguing that the $`a`$-maximization determines all-order loop effects only from the one-loop result Eq.(72).
We can also study $`a`$ of the gauged version of extended IYIT model in section 3, which leads to conformal sequestering. Since the gauging enforces yet another constraint on the anomaly free $`R`$ charge, we obtain a unique $`R`$ charge assignment without using the $`a`$-maximization procedure. It is important to realize, following the general argument of the monotonically decreasing $`a`$, that $`a_{\mathrm{gauged}}`$ is less than $`a_{\mathrm{ungauged}}`$. This is obvious when $`x_{}^{\mathrm{gauged}}`$ is sufficiently close to $`x_{}^{\mathrm{ungauged}}`$, since $`x_{}^{\mathrm{ungauged}}`$ yields the local maximum of $`a(x)`$. For example, we can show by a direct computation that $`a`$ of the gauged extended IYIT model presented in section 3 is always less than that of the ungaged version presented in section 2 for a fixed gauge group. This result is consistent with the fact that our conformal fixed point is a stable one. In particular, it is worthwhile to notice that this is even true for $`N=1`$ case, which cannot be treated in the one-loop approximation.
Finally, it would be an interesting but challenging problem to obtain the speed of the conformal sequestering from the interpolating $`a`$-function. In Ref., the off-shell $`a`$-function is proposed as solving $`a`$-maximization condition with a Lagrange multiplier $`\xi `$ that enforces the constraint on the $`R`$ charge:
$$a(R(\xi ),\xi )=\mathrm{Tr}\left(3R^3R\right)+\xi (\mathrm{constraint}),$$
(77)
where $`R(\xi )`$ is obtained by maximizing $`a`$ with respect to $`R`$ for fixed $`\xi `$, and the constraint is either ABJ anomaly free condition or the requirement that the superpotential be marginal. As was observed in , the first derivative of $`a(\xi )`$ is related to the $`\beta `$ function of the coupling constant.<sup>26</sup><sup>26</sup>26Since the number of the Lagrange multipliers $`\xi `$ agrees with that of marginal deformations, it is conjectured that $`\xi `$ can be regarded as a coupling constant in a certain scheme. Furthermore, the second derivative (Hessian) of $`a(\xi )`$ at a fixed point $`\xi _{}`$ is proportional to the slope of the $`\beta `$ function
$$\frac{^2a(\xi )}{\xi _i\xi _j}|_{}\frac{\beta _i(\xi )}{\xi _j}|_{}.$$
(78)
Consequently there is a chance to read the conformal sequestering matrix without performing the explicit loop calculation even for a strongly coupled theory. Unfortunately, we do not know the proportionality factor (related to the denominator of the NSVZ beta function evaluated at the fixed point) and the transformation matrix $`\{\xi _i/g_j\}`$ non-perturbatively, so we cannot determine the conformal sequestering matrix. Since the conformal sequestering matrix is a physical renormalization invariant quantity while $`a(\xi )`$ is not, we need an off-shell scheme-independent $`a`$-function for our purposes.
## Appendix B Another example of the hybrid scheme
In section 5, we have considered the hybrid model of the anomaly and gauge mediated SUSY breaking, which solves the tachyonic slepton problem in a pure anomaly mediation model. In this appendix, we propose another example of the hybrid model which can be constructed with renormalizable interactions between the SUSY breaking sector and the messenger sector. That is, in addition to the anomaly mediated SUSY breaking, we consider the gauge mediated SUSY breaking discussed in Ref. , where messenger sector consists of $`N_m`$ flavors of chiral superfields $`\psi _i`$ and $`\overline{\psi }_j`$ $`(i,j=1,\mathrm{}N_m)`$ and $`N_m`$ flavors of chiral superfields $`\psi _i^{}`$ and $`\overline{\psi }_j^{}`$ $`(i,j=1,\mathrm{}N_m)`$. Here, $`\psi _i`$, $`\psi _i^{}`$ and $`\overline{\psi }_j`$, $`\overline{\psi }_j^{}`$ transform as $`\mathrm{𝟓}`$ and $`\mathrm{𝟓}^{}`$ of $`SU(5)`$ GUT, respectively. Then, with the superpotential,
$$hS_{ij}\psi _i\overline{\psi }_j+m_m\psi _i\overline{\psi }_i^{}+m_m\psi _i^{}\overline{\psi }_i,$$
(79)
the SUSY-breaking effects are transmitted to the sfermions and Higgs bosons by the gauge interactions. Here, $`h`$ denotes the coupling constant, $`m_m`$ the mass parameter, and we assume that $`h=h_0`$ of order one at $`M_{}^{}{}_{}{}^{<}M_G`$. As discussed in section 3, we impose the global $`SU(2N+2)`$ symmetry to the SUSY breaking sector, and in order for the interaction in Eq.(79) to respect this symmetry, we assume that $`N_m=2(N+1)`$ and $`\psi `$, $`\overline{\psi }`$ and $`\psi ^{}`$, $`\overline{\psi }^{}`$ transform as $`\overline{\mathrm{𝟐}𝐍+\mathrm{𝟐}}`$ and $`\mathrm{𝟐}𝐍+\mathrm{𝟐}`$ representations, respectively, of the $`SU(2N+2)`$.<sup>27</sup><sup>27</sup>27 Here, we assume $`S=0`$, that is, $`R`$ symmetry is not broken. In this case, gauginos do not obtain the SUSY-breaking masses via the gauge mediation and the gaugino spectrum is the same as in the pure anomaly mediation.
In this case, the scalar field $`\varphi `$ obtains the mass squared via the gauge mediation, and at the messenger scale, it is given by,
$`m_\varphi ^2`$ $``$ $`2{\displaystyle \underset{a=1,2,3}{}}C_a^\varphi \left({\displaystyle \frac{\alpha _a}{4\pi }}\right)^2{\displaystyle \frac{|hF_S|^2}{m_m^2}},`$ (80)
$``$ $`6{\displaystyle \underset{a=1,2,3}{}}C_a^\varphi \left({\displaystyle \frac{\alpha _a}{4\pi }}\right)^2\left({\displaystyle \frac{m_{\mathrm{phys}}}{M_{}}}\right)^{\gamma _S}\left({\displaystyle \frac{|h_0|M_G}{m_m}}\right)^2m_{3/2}^2.`$ (81)
Here, we have used RG evolution of $`hF_S`$,
$`hF_S|_{m_m}=hF_S|_{m_{\mathrm{phys}}}=\left({\displaystyle \frac{m_{\mathrm{phys}}}{M_{}}}\right)^{\frac{\gamma _S}{2}}h_0F_S|_{m_{\mathrm{phys}}}\left({\displaystyle \frac{m_{\mathrm{phys}}}{M_{}}}\right)^{\frac{\gamma _S}{2}}h_0m_{3/2}M_G,`$ (82)
where subscripts $`(m_m,m_{\mathrm{phys}})`$ denote the RG scale. Furthermore, in the course of RG running, the scalar mass squared is suppressed by a factor of $`\eta `$ which denotes the sequestering effects between the scales $`m_m`$ and $`m_{\mathrm{phys}}`$,
$`\eta =\{\begin{array}{ccc}(m_{\mathrm{phys}}/m_m)^\beta _{}^{}& \mathrm{for}& m_m>m_{\mathrm{phys}},\\ 1& \mathrm{for}& m_mm_{\mathrm{phys}}.\end{array}`$ (85)
For example, for the $`SP(3)\times SP(1)^2`$ model, the positive contribution to the slepton masses squared in Eq.(81) is given by
$`m_{\stackrel{~}{e}\mathrm{G}.\mathrm{M}.}^2{\displaystyle \frac{3}{5}}\mathrm{\hspace{0.17em}\hspace{0.17em}6}\left({\displaystyle \frac{m_{\mathrm{phys}}}{m_m}}\right)^2\left({\displaystyle \frac{M_G}{M_{}}}\right)^2\left({\displaystyle \frac{\alpha _1}{4\pi }}\right)^2m_{3/2}^2\eta ,`$ (86)
where we are using $`\gamma _S=2`$ and are assuming $`h_0=1`$.<sup>28</sup><sup>28</sup>28Here, we have neglected RG effects from the MSSM couplings. Thus, when $`m_m`$ satisfies
$`m_{m}^{}{}_{}{}^{<}\{\begin{array}{ccc}m_{\mathrm{phys}}(M_G/M_{})^{2/(2+\beta _{}^{})}& \mathrm{for}& m_m>m_{\mathrm{phys}},\\ m_{\mathrm{phys}}(M_G/M_{})& \mathrm{for}& m_mm_{\mathrm{phys}},\end{array}`$ (89)
this contribution overcomes the negative contribution from the anomaly mediated mass squared in Eq.(63).<sup>29</sup><sup>29</sup>29In this model, the gauge coupling constants in the SSM remain perturbative up to Grand Unification scale $`M_{\mathrm{GUT}}2\times 10^{16}`$ GeV for $`m_{m}^{}{}_{}{}^{>}10^{13}`$ GeV. Therefore, we find that this hybrid model also provides a solution to the tachyonic slepton problem with an appropriate choice of the mass scale of messengers.<sup>30</sup><sup>30</sup>30Since the messenger particles are very heavy $`m_{m}^{}{}_{}{}^{>}10^{13}`$ GeV, they are not produced thermally for the reheating temperature of the universe around $`T_R10^{10}`$ GeV, which is very advantageous for the thermal leptogenesis . |
warning/0506/hep-ph0506171.html | ar5iv | text | # Contents
## 1 Introduction
Electroproduction processes in the Bjorken regime probe the partonic structure of the nucleon. Inclusive deep inelastic scattering (DIS) provides extensive information about the sum of quark and antiquark distributions in the nucleon, and allows one to determine the gluon distribution from the observed scaling violations. More detailed information can be obtained from scattering experiments in which one or more hadrons are observed in the final state. There are two basic classes of such experiments. The first one is semi-inclusive deep inelastic scattering, in which one observes a single hadron, carrying a fraction $`z`$ of the photon energy in the target rest frame, in a final state of large average multiplicity. A QCD factorization theorem states that the semi-inclusive cross section in the Bjorken limit can be expressed in terms of distribution functions for quarks, antiquarks and gluons in the target and of the corresponding fragmentation functions into the observed hadron. This allows one to tag the active parton via its fragmentation properties and has recently been used for a flavor decomposition of polarized quark and antiquark distributions in the semi-inclusive production of pions and kaons . In addition, measurements of azimuthal asymmetries in semi-inclusive pion and kaon production, such as the Collins and Sivers asymmetries for a transversely polarized target , provide interesting information about the distribution of the spin and transverse momentum carried by quarks and antiquarks in the nucleon.
The second class are exclusive scattering processes, in which the final state is a recoiling baryon $`B`$ together with a single meson or a few-meson system carrying nearly the full photon energy in the target rest frame. A QCD factorization theorem states that in the Bjorken limit, and for longitudinal photon polarization, the amplitudes of such processes are calculable in terms of the light-cone distribution amplitude of the produced meson and generalized parton distributions (GPDs) for the $`pB`$ transition. Generalized parton distributions provide a wealth of information on the parton structure of the nucleon, in particular about the spatial distribution of partons in the transverse plane and about quark orbital angular momentum, see e.g. the reviews .
The Bjorken limit for semi-inclusive electroproduction implies a large average multiplicity of the produced hadronic system. In practice, however, there can be situations in which individual exclusive channels play an important role. In fixed-target experiments the limited photon energy restricts the phase space for quark fragmentation, in particular at large $`z`$, and for relatively low photon virtuality $`Q^2`$ the suppression of individual exclusive channels due to the faster drop of the exclusive cross sections with $`Q^2`$ may not yet be effective. In particular, phenomenological studies suggest that a large contribution to $`\pi ^\pm `$ production comes from exclusive $`\rho ^0`$ production, with subsequent decay $`\rho ^0\pi ^+\pi ^{}`$ . Fortunately, the cross sections for exclusive $`\rho ^0`$ electroproduction has been measured by the HERMES and CLAS experiments, including the cross section ratio for longitudinal and transverse photons . It is natural to ask whether the strange vector mesons $`\varphi `$ and $`K^{}`$ play an equally important role in semi-inclusive kaon production and whether other exclusive channels may be important, too. A quantitative answer to these questions will help to delineate the limits of the kinematic region where semi-inclusive data can be analyzed using the factorization theorem.
In this paper we investigate the role of exclusive channels in the semi-inclusive production of pions and kaons on the basis of the factorization theorem for hard exclusive processes. Our investigation consists of two parts. Firstly, we evaluate the longitudinal cross section for the exclusive production of pseudoscalar and vector mesons in the leading-twist approximation and at leading order in the strong coupling, focusing on $`\pi `$, $`K`$ and $`\rho `$, $`\varphi `$, $`K^{}`$. We explore uncertainties in the obtained cross sections, in particular those due to the generalized parton distributions, which are still largely unknown and need to be modeled. Such uncertainties will persist if higher-order and higher-twist corrections are included. Seen from a different perspective, they indicate to which extent exclusive meson production is sensitive to GPDs and thus interesting in its own right. It is known that exclusive meson production cross sections at moderate $`Q^2`$ are affected by substantial power corrections. For the production of $`\rho ^0`$, $`\varphi `$ and $`\pi ^+`$ there is data or preliminary data, which we will use to assess the quantitative validity of our calculated cross sections. Secondly, we evaluate the contribution of these exclusive channels to semi-inclusive production of $`\pi `$ and $`K`$, and compare them with the results obtained from leading-twist quark fragmentation. For the exclusive meson production cross sections we rely on experimental data where possible, and only use our leading-twist calculation to estimate the *ratio* of cross sections for measured and unmeasured channels.
The paper is organized as follows. In Sects. 2 and 3 we summarize the leading-twist description of exclusive meson production and describe the models for the GPDs used in our investigation. An analysis of pseudoscalar and vector meson production channels is then given in Sects. 4 and 5. In Sect. 6 we briefly discuss the limitations of our leading-order results and compare them with experimental data. The contribution of exclusive channels to semi-inclusive meson production is compared with leading-twist quark fragmentation in Sect. 7, and in Sect. 8 we summarize our main results. Some important technical details are given in two appendices.
## 2 Exclusive meson production in the leading-twist approximation
Let us consider the exclusive electroproduction process
$$e(k)+p(p)e(k^{})+M(q^{})+B(p^{}),$$
(1)
where $`M`$ is a meson and $`B`$ a baryon, and where four-momenta are indicated in parentheses. Throughout this work we assume beam and target to be unpolarized. We write $`q=kk^{}`$, $`\mathrm{\Delta }=p^{}p`$, and use the standard kinematic variables
$$t=\mathrm{\Delta }^2,Q^2=q^2,W^2=(p+q)^2,x_B=Q^2/(2pq),y=(pq)/(pk).$$
(2)
We respectively write $`m_p`$, $`m_B`$, $`m_M`$ for the masses of the proton, the baryon $`B`$, and the meson $`M`$. With Hand’s convention for the virtual photon flux, the electroproduction cross section is given by
$$\frac{d\sigma (epeMB)}{dQ^2dx_Bdt}=\frac{\alpha _{em}}{2\pi }\frac{y^2}{1ϵ}\frac{1x_B}{x_B}\frac{1}{Q^2}\left[\frac{d\sigma _T}{dt}+ϵ\frac{d\sigma _L}{dt}\right]$$
(3)
in terms of the cross sections $`d\sigma _T/dt`$ and $`d\sigma _L/dt`$ of the $`\gamma ^{}pMB`$ subprocess for transverse and longitudinal $`\gamma ^{}`$, where
$$ϵ=\frac{1y(yx_Bm_p/Q)^2}{1y+y^2/2+(yx_Bm_p/Q)^2}$$
(4)
is the ratio of longitudinal to transverse photon flux.
In the generalized Bjorken limit of large $`Q^2`$ at fixed $`x_B`$ and fixed $`t`$, the process amplitude factorizes into a hard-scattering kernel convoluted with generalized parton distributions for the $`pB`$ transition and with the distribution amplitude of the meson . Example graphs are shown in Fig. 1. In this limit the longitudinal cross section can be written as
$$\frac{d\sigma _L}{dt}=\frac{\alpha _{em}}{Q^6}\frac{x_B^2}{1x_B}\left\{(1\xi ^2)||^2\left[\frac{2\xi (m_B^2m_p^2)+t}{(m_B+m_p)^2}+\xi ^2\right]||^2\left[\xi +\frac{m_Bm_p}{m_B+m_p}\right]\mathrm{\hspace{0.17em}2}\xi \mathrm{Re}(^{})\right\}$$
(5)
for vector mesons, and as
$$\frac{d\sigma _L}{dt}=\frac{\alpha _{em}}{Q^6}\frac{x_B^2}{1x_B}\left\{(1\xi ^2)|\stackrel{~}{}|^2+\frac{(m_Bm_p)^2t}{(m_B+m_p)^2}\xi ^2|\stackrel{~}{}|^2\left[\xi +\frac{m_Bm_p}{m_B+m_p}\right]\mathrm{\hspace{0.17em}2}\xi \mathrm{Re}(\stackrel{~}{}^{}\stackrel{~}{})\right\}$$
(6)
for pseudoscalar mesons $`M`$. The transverse cross section $`d\sigma _T/dt`$ is power suppressed by $`1/Q^2`$ compared with $`d\sigma _L/dt`$. Here we have in addition used the skewness variable
$$\xi =\frac{(pp^{})(q+q^{})}{(p+p^{})(q+q^{})}\frac{x_B}{2x_B},$$
(7)
where the approximation holds in the generalized Bjorken limit. Note that the prefactor in (5) and (6) can be rewritten as $`x_B^2/(1x_B)=4\xi ^2/(1\xi ^2)`$. The quantities $``$, $``$ and $`\stackrel{~}{}`$, $`\stackrel{~}{}`$ are specific for each channel. Throughout this work we will take their leading order approximations in $`\alpha _s`$. To be specific, we have
$`_{K^+\mathrm{\Lambda }}(\xi ,t)`$ $`=`$ $`{\displaystyle \frac{4\pi \alpha _s}{27}}f_K^{}[{\displaystyle _0^1}dz{\displaystyle \frac{1}{z(1z)}}\varphi _{K^+}(z){\displaystyle _1^1}dx{\displaystyle \frac{2H_{p\mathrm{\Lambda }}(x,\xi ,t)+H_{p\mathrm{\Lambda }}(x,\xi ,t)}{\xi xi\epsilon }}`$
$`{\displaystyle _0^1}dz{\displaystyle \frac{2z1}{z(1z)}}\varphi _{K^+}(z){\displaystyle _1^1}dx{\displaystyle \frac{2H_{p\mathrm{\Lambda }}(x,\xi ,t)H_{p\mathrm{\Lambda }}(x,\xi ,t)}{\xi xi\epsilon }}],`$
$`\stackrel{~}{}_{K^+\mathrm{\Lambda }}(\xi ,t)`$ $`=`$ $`{\displaystyle \frac{4\pi \alpha _s}{27}}f_K[{\displaystyle _0^1}dz{\displaystyle \frac{1}{z(1z)}}\varphi _{K^+}(z){\displaystyle _1^1}dx{\displaystyle \frac{2\stackrel{~}{H}_{p\mathrm{\Lambda }}(x,\xi ,t)+\stackrel{~}{H}_{p\mathrm{\Lambda }}(x,\xi ,t)}{\xi xi\epsilon }}`$ (8)
$`{\displaystyle _0^1}dz{\displaystyle \frac{2z1}{z(1z)}}\varphi _{K^+}(z){\displaystyle _1^1}dx{\displaystyle \frac{2\stackrel{~}{H}_{p\mathrm{\Lambda }}(x,\xi ,t)\stackrel{~}{H}_{p\mathrm{\Lambda }}(x,\xi ,t)}{\xi xi\epsilon }}]`$
for $`\gamma ^{}pK^+\mathrm{\Lambda }`$ and $`\gamma ^{}pK^+\mathrm{\Lambda }`$, respectively, with analogous expressions for $`_{K^+\mathrm{\Lambda }}`$ and $`\stackrel{~}{}_{K^+\mathrm{\Lambda }}`$. The GPDs for the $`p\mathrm{\Lambda }`$ transition are defined as
$`{\displaystyle \frac{1}{2}}{\displaystyle \frac{dz^{}}{2\pi }e^{ixP^+z^{}}\mathrm{\Lambda }|\overline{s}(\frac{1}{2}z)\gamma ^+u(\frac{1}{2}z)|p}|_{z^+=0,z_T=0}`$
$`=`$ $`{\displaystyle \frac{1}{2P^+}}\left[H_{p\mathrm{\Lambda }}(x,\xi ,t)\overline{u}\gamma ^+u+E_{p\mathrm{\Lambda }}(x,\xi ,t)\overline{u}{\displaystyle \frac{i\sigma ^{+\alpha }\mathrm{\Delta }_\alpha }{m_\mathrm{\Lambda }+m_p}}u\right],`$
$`{\displaystyle \frac{1}{2}}{\displaystyle \frac{dz^{}}{2\pi }e^{ixP^+z^{}}\mathrm{\Lambda }|\overline{s}(\frac{1}{2}z)\gamma ^+\gamma _5u(\frac{1}{2}z)|p}|_{z^+=0,z_T=0}`$ (9)
$`=`$ $`{\displaystyle \frac{1}{2P^+}}\left[\stackrel{~}{H}_{p\mathrm{\Lambda }}(x,\xi ,t)\overline{u}\gamma ^+\gamma _5u+\stackrel{~}{E}_{p\mathrm{\Lambda }}(x,\xi ,t)\overline{u}{\displaystyle \frac{\gamma _5\mathrm{\Delta }^+}{m_\mathrm{\Lambda }+m_p}}u\right],`$
where we use light-cone coordinates $`v^\pm =(v^0\pm v^3)/\sqrt{2}`$ and $`v_T=(v^1,v^2)`$ for a four-vector $`v`$ and assume light-cone gauge $`A^+=0`$. For brevity we have not displayed the momentum and polarization dependence of the baryon spinors on the right-hand sides. GPDs for other transitions are defined in full analogy. The integrals over meson distribution amplitudes in (2) can be expressed as
$$_0^1𝑑z\frac{1}{z(1z)}\varphi (z)=6\left[1+\underset{n=1}{\overset{\mathrm{}}{}}a_{2n}\right],_0^1𝑑z\frac{2z1}{z(1z)}\varphi (z)=6\underset{n=1}{\overset{\mathrm{}}{}}a_{2n1}$$
(10)
through their coefficients in the expansion
$$\varphi (z)=6z(1z)\left[1+\underset{n=1}{\overset{\mathrm{}}{}}a_nC_n^{3/2}(2z1)\right]$$
(11)
on Gegenbauer polynomials, where $`z`$ is the light-cone momentum fraction carried by the quark in the meson. Note that odd Gegenbauer coefficients $`a_{2n1}`$ describe an asymmetry in the momentum distribution of the quark and antiquark in the meson. They can be nonzero for $`K`$ and $`K^{}`$ due to flavor SU(3) breaking. In (2) to (11) we have not displayed the logarithmic dependence on the renormalization scale in $`\alpha _s`$ and on the factorization scale in the GPDs and the distribution amplitudes.
Using flavor SU(3) symmetry one can relate the transition GPDs from the proton to a hyperon to the distributions $`H^q(x,\xi ,t)`$ for quark flavor $`q`$ in the proton . This gives in particular $`H_{p\mathrm{\Lambda }}=[2H^uH^dH^s]/\sqrt{6}`$ and an analogous relation for $`\stackrel{~}{H}_{p\mathrm{\Lambda }}`$. We use these relations throughout, except for $`\stackrel{~}{E}`$, where there are large effects of SU(3) breaking as we shall see below. Results analogous to (2) hold for all meson channels we consider, see e.g. , and we have collected the relevant combinations of GPDs in Table 1. There we have introduced $`H^{\overline{q}}(x,\xi ,t)=H^q(x,\xi ,t)`$ and $`\stackrel{~}{H}^{\overline{q}}(x,\xi ,t)=\stackrel{~}{H}^q(x,\xi ,t)`$, so that for $`x>0`$ we have simple forward limits
$$H^q(x,0,0)=q(x),H^{\overline{q}}(x,0,0)=\overline{q}(x),\stackrel{~}{H}^q(x,0,0)=\mathrm{\Delta }q(x),\stackrel{~}{H}^{\overline{q}}(x,0,0)=\mathrm{\Delta }\overline{q}(x)$$
(12)
in terms of the unpolarized and polarized quark and antiquark densities in the proton. For gluons we have $`H^g(x,0,0)=xg(x)`$, which is the origin of the additional factors $`x^1`$ in the entries for $`\rho ^0`$, $`\omega `$, $`\varphi `$. In addition to the replacements in Table 1 one has of course to take the appropriate meson distribution amplitude and meson decay constants in (2). For the latter we will take $`f_\pi =131\text{MeV}`$, $`f_K=160\text{MeV}`$, and
$$f_\rho =209\text{MeV},f_\omega =187\text{MeV},f_\varphi =221\text{MeV},f_K^{}=218\text{MeV}$$
(13)
from .
For $`\alpha _s`$ in (2) we will take the one-loop running coupling at the scale $`Q^2`$, with three active quark flavors and $`\mathrm{\Lambda }_{\mathrm{QCD}}=200\text{MeV}`$. This gives $`\alpha _s=0.34`$ at $`Q^2=2.5\text{GeV}^2`$, where we will show most of our numerical results. We will not attempt more refined choices of renormalization scale, as were for instance explored in , since our principal use of the leading-order calculation will be to describe the *relative* size of cross sections for different exclusive channels.
## 3 Modeling the generalized parton distributions
For the calculation of exclusive cross sections we use simple models of GPDs. They have been developed in and been used in most phenomenological analyses so far. Our aim here is not to improve on these models, but instead to see by how much predictions can vary *within* the given framework. We take a factorizing $`t`$ dependence for $`H`$ and $`\stackrel{~}{H}`$,
$`H^q(x,\xi ,t)`$ $`=`$ $`H^q(x,\xi )F_1^p(t),H^g(x,\xi ,t)=H^g(x,\xi )F_1^p(t),`$
$`\stackrel{~}{H}^q(x,\xi ,t)`$ $`=`$ $`\stackrel{~}{H}^q(x,\xi )G_A(t)/G_A(0),`$ (14)
where $`F_1^p(t)`$ is the electromagnetic Dirac form factor of the proton and $`G_A(t)`$ the isovector axial form factor of the nucleon. A more refined version of the model would take different combinations of the proton and neutron form factors for $`H^u`$ and $`H^d`$, but for the low values of $`t`$ dominating integrated cross sections, $`F_1^n(t)`$ is much smaller than $`F_1^p(t)`$ and we simply neglect it. In this sense (3) is consistent with the sum rule for the first moment $`𝑑xH^q(x,\xi ,t)`$. The ansatz for $`\stackrel{~}{H}^q`$ is consistent with the sum rule for $`𝑑x\stackrel{~}{H}^q(x,\xi ,t)`$ to the extent that the (unknown) isoscalar axial form factor has the same $`t`$ dependence as the isovector one. In our numerical evaluations we take the familiar parameterizations
$$F_1^p(t)=\frac{4m_p^22.8t}{4m_p^2t}\frac{1}{[1t/(0.71\text{GeV}^2)]\text{}^2},\frac{G_A(t)}{G_A(0)}=\frac{1}{[1t/(1.05\text{GeV}^2)]\text{}^2}.$$
(15)
We note that for the gluon distribution $`H^g`$ there is no reason a priori to take the electromagnetic form factor $`F_1^p(t)`$ in the ansatz (3). It turns out, however, that $`F_1^p(t)`$ is well approximated by a dipole form $`F_1^p(t)=[1t/(0.98\text{GeV}^2)]^2`$ for $`t`$ up to about $`3\text{GeV}^2`$ and thus close to the two-gluon form factor advocated in .
It is rather certain that the ansatz (3) is too simple and can at best reflect the correct $`t`$ dependence in a limited range of $`x`$ and $`\xi `$ . For $`x`$ and $`\xi `$ in the valence region, say above $`0.2`$, the decrease of GPDs with $`t`$ is most likely less steep than the one of $`F_1^p(t)`$ and $`G_A(t)`$. Whereas there are phenomenological constraints of the $`t`$ behavior of valence quark GPDs and for gluons at small $`x`$ , the behavior for sea quarks is largely unknown, and sea quarks are important for the $`x_B`$ region around $`0.1`$ we will be mostly concerned with. Furthermore, the $`t`$ dependence of meson production at moderate $`Q^2`$ is strongly affected by power corrections, as is for instance seen in the $`Q^2`$ dependence of the logarithmic slope $`B=(/t)\mathrm{log}(d\sigma /dt)|_{t=0}`$ for $`\rho ^0`$ production at very high energies . We note that our ansatz (3) gives a slope parameter $`B4\text{GeV}^2`$, which may be quite realistic for $`x_B`$ around $`0.1`$. Furthermore, cross section ratios should be less affected by the insufficiency of our ansatz, since they are sensitive only to the relative $`t`$ dependence of the contributions from different quark flavors and from gluons.
For the $`t`$ independent functions in (3) we use the double distribution based ansatz of , whose ingredients are the usual parton densities at a given factorization scale $`\mu `$ and a so-called profile parameter $`b`$, where $`\mu `$ and $`b`$ are to be regarded as free parameters of the model. Explicit expressions are given in App. A. We will not take into account the evolution of GPDs, which should not be too problematic since our numerical applications will stay within a rather narrow range of $`Q^2`$.
The modeling of the nucleon helicity-flip distributions $`E^q`$ and $`E^g`$ is still at an early stage of development, with the most advanced considerations focused on the valence quark domain . Fortunately, contributions from $`E`$ enter the unpolarized meson production cross section (5) with prefactors that are quite small in the kinematics we are most interested in. Following the argumentation of that $`E`$ is not much larger than $`H`$ for a given parton species, we hence neglect $`E`$ altogether in our cross section estimates.
The distributions $`\stackrel{~}{E}`$ cannot be neglected since they receive contributions proportional to $`\xi ^1`$ that compensate the kinematic prefactors in the cross section (6). We model them as
$`\stackrel{~}{E}_{pn}(x,\xi ,t)`$ $`=`$ $`{\displaystyle \frac{\theta (|x|\xi )}{2\xi }}\varphi _\pi \left({\displaystyle \frac{x+\xi }{2\xi }}\right){\displaystyle \frac{2m_pf_\pi g_{\pi NN}}{m_\pi ^2t}}{\displaystyle \frac{\mathrm{\Lambda }^2m_\pi ^2}{\mathrm{\Lambda }^2t}},`$
$`\stackrel{~}{E}_{p\mathrm{\Lambda }}(x,\xi ,t)`$ $`=`$ $`{\displaystyle \frac{\theta (|x|\xi )}{2\xi }}\varphi _K\left({\displaystyle \frac{x+\xi }{2\xi }}\right){\displaystyle \frac{(m_p+m_\mathrm{\Lambda })f_Kg_{KN\mathrm{\Lambda }}}{m_K^2t}}{\displaystyle \frac{\mathrm{\Lambda }^2m_K^2}{\mathrm{\Lambda }^2t}},`$ (16)
where the distribution amplitudes $`\varphi `$ are the same as those introduced above. For the coupling constants we take the value $`g_{\pi NN}=2m_pG_A(0)/f_\pi 14.7`$ from the Goldberger-Treiman relation and $`g_{KN\mathrm{\Lambda }}13.3`$ from . Continued to the points $`t=m_\pi ^2`$ or $`t=m_K^2`$ in the unphysical region, the expressions (3) become the well-known results from pion or kaon exchange . These can only be expected to be good approximations for $`t`$ close to the squared meson masses, and for $`t`$ of several $`0.1\text{GeV}^2`$ are to be regarded as extrapolations. In (3) we have included form factors that cut off the $`1/(m_M^2t)`$ behavior of the pure pole terms when $`t`$ becomes large. As default value for the cut-off mass we will take $`\mathrm{\Lambda }=0.8\text{GeV}`$ and study the sensitivity of results to the precise value of this parameter. We note that for $`\mathrm{\Lambda }=0.6\text{GeV}`$ and $`t1\text{GeV}^2`$ the above form of $`\stackrel{~}{E}_{pn}`$ differs by less than $`15\%`$ from the corresponding term calculated in the chiral quark-soliton model as given in Eq. (4.39) of .
With this model for $`\stackrel{~}{E}`$ the longitudinal cross section for $`\gamma ^{}p\pi ^+n`$ takes the form
$`{\displaystyle \frac{d\sigma _L}{dt}}`$ $`=`$ $`{\displaystyle \frac{\alpha _{em}}{Q^6}}{\displaystyle \frac{x_B^2}{1x_B}}\{(1\xi ^2)|\stackrel{~}{}(\xi ,t)|^22m_p\xi \mathrm{Re}\stackrel{~}{}(\xi ,t)Q^2F_\pi (Q^2){\displaystyle \frac{g_{\pi NN}}{m_\pi ^2t}}{\displaystyle \frac{\mathrm{\Lambda }^2m_\pi ^2}{\mathrm{\Lambda }^2t}}`$ (17)
$`{\displaystyle \frac{t}{4}}\left[Q^2F_\pi (Q^2){\displaystyle \frac{g_{\pi NN}}{m_\pi ^2t}}{\displaystyle \frac{\mathrm{\Lambda }^2m_\pi ^2}{\mathrm{\Lambda }^2t}}\right]^2\},`$
where
$$F_\pi (Q^2)=\frac{2\pi \alpha _s}{9}\frac{f_\pi ^2}{Q^2}\left[_0^1𝑑z\frac{1}{z(1z)}\varphi _\pi (z)\right]^2$$
(18)
is the electromagnetic pion form factor to leading order in $`\alpha _s`$ and $`1/Q^2`$. An similar expression involving $`F_{K^+}(Q^2)`$ is obtained for $`\gamma ^{}pK^+\mathrm{\Lambda }`$ according to (6) and (3). Note that the the $`|\stackrel{~}{}|^2`$ term in $`d\sigma _L/dt`$ has no $`x_B`$ dependence other than from the explicit factor $`x_B^2/(1x_B)`$. Within our model the $`|\stackrel{~}{}|^2`$ term reflects the behavior of the polarized parton distributions at momentum fractions of order $`\xi `$, and its contribution to $`d\sigma _L/dt`$ can very roughly be represented by $`[\xi \mathrm{\Delta }q(\xi )]^2`$.
## 4 Exclusive pseudoscalar meson production
In this and the next section we present numerical results for cross sections of exclusive meson production. Our main focus is to compare the rates for different production channels and to investigate model uncertainties. Comparison with data in Sect. 6 will allow us to estimate the shortcomings of the leading approximation in $`1/Q^2`$ and in $`\alpha _s`$, on which our calculations are based.
The factorization theorem for exclusive meson production requires $`t`$ to be much smaller than $`Q^2`$. For definiteness we will give all meson cross sections in this paper integrated over $`t`$ from its smallest kinematically allowed value $`t_0`$ to an upper limit of $`1\text{GeV}^2`$. In generalized Bjorken kinematics we have
$$t_0\frac{2\xi ^2(m_B^2+m_p^2)+2\xi (m_B^2m_p^2)}{1\xi ^2}$$
(19)
with $`\xi `$ defined in (7). For low enough $`x_B`$ most of the cross section should be accumulated in this $`t`$ region, whereas for large $`x_B`$ our cross sections decrease to the extent that $`t_0`$ approaches $`1\text{GeV}^2`$.
To begin with, let us investigate the relative importance of the contributions from $`\stackrel{~}{}`$ and $`\stackrel{~}{}`$ to $`\pi ^+`$ and $`K^+`$ production with our model assumptions. As is seen in Fig. 2, exclusive $`\pi ^+`$ production receives comparable contributions from both the $`|\stackrel{~}{}|^2`$ and the $`|\stackrel{~}{}|^2`$ term in (6), as well as from the interference term proportional to $`\mathrm{Re}(\stackrel{~}{}^{}\stackrel{~}{})`$. Note that the relative weight of the contributions is different for $`d\sigma _L/dt`$, where it strongly depends on $`t`$ given the characteristic $`t`$ dependence of the pion pole term (3). In $`K^+`$ production the influence of $`\stackrel{~}{}`$ is less prominent, since the pole factor $`(m_K^2t)^1`$ gives much less enhancement at small $`t`$ than $`(m_\pi ^2t)^1`$.
To obtain the curves in Fig. 2 we have made a number of choices in the nonperturbative input to the cross section, which we now discuss in turn. In Fig. 3 we show how the cross section changes when we vary the parameter $`\mathrm{\Lambda }`$ in our model for $`\stackrel{~}{E}`$, where $`\mathrm{\Lambda }=1.3\text{GeV}`$ represents an upper limit of the values discussed in the phenomenological study , and $`\mathrm{\Lambda }=0.6\text{GeV}`$ approximates the form factor dependence obtained for the pion pole contribution in , as discussed at the end of Sect. 2. Omitting the form factor altogether (tantamount to setting $`\mathrm{\Lambda }\mathrm{}`$) the $`\pi ^+`$ cross section would increase by more than a factor 1.4 and the $`K^+`$ cross section by more than a factor 1.7 compared with our default choice $`\mathrm{\Lambda }=0.8\text{GeV}`$. Also, the cross sections would considerably increase when raising the upper cutoff in the $`t`$ integration above $`1\text{GeV}^2`$. In other words, the cross section would then receive substantial contributions from values of $`t`$ far away from the region where the pion or kaon pole term can be regarded as a reasonable approximation of $`\stackrel{~}{E}`$.
For the pion distribution amplitude we have taken the asymptotic form $`\varphi _\pi (z)=6z(1z)`$ under scale evolution, which is close to what can be extracted from data on the $`\gamma `$$`\pi `$ transition form factor, see e.g. . The study in quotes limits on $`a_2+a_4`$ at scale $`\mu =1\text{GeV}`$ of about $`\pm 0.3`$ if all other Gegenbauer coefficients are set to zero. This corresponds to a change of the $`\gamma ^{}p\pi ^+n`$ cross section by a factor $`(1+a_2+a_4)^2`$ between $`0.5`$ and $`1.7`$. For the $`K^+`$ distribution amplitude we have taken the asymptotic form as well. Figure 4 shows how the $`K^+`$ cross section changes if instead one takes $`a_1=0.05`$ and $`a_2=0.27`$ at $`\mu =1\text{GeV}`$ from the QCD sum rule calculation . This value of $`a_1`$ is compatible with the findings of .
For our model of $`\stackrel{~}{H}`$ we have taken a double distribution ansatz with a profile parameter $`b=1`$ (see Sect. 2 and App. A). Taking $`b=2`$ instead would decrease the $`K^+`$ cross section by a factor of approximately 0.6. The pion cross section changes less, because of the relative weight of $`\stackrel{~}{}`$ and $`\stackrel{~}{}`$. A more important source of uncertainty is however due to the polarized quark densities used as input to model $`\stackrel{~}{H}`$. As a default we have used the LO parameterization from at a scale $`\mu =1\text{GeV}`$. Using instead the LO parameterization in scenario 1 of at the same scale, the $`K^+`$ cross section changes as shown in Fig. 4. Note that any uncertainty on parton distributions is amplified in the meson production cross section, where GPDs appear squared.
Let us now comment on other pseudoscalar channels. The cross sections for $`\gamma ^{}pK^+\mathrm{\Sigma }^0`$ is about an order of magnitude smaller than the one for $`\gamma ^{}pK^+\mathrm{\Lambda }`$, as is seen in the numerical study of . For the contribution from $`\stackrel{~}{}`$ this can be understood from the flavor structure in Table 1, where for a rough estimate one may concentrate on the dominant terms $`\stackrel{~}{H}^u`$ and $`\stackrel{~}{H}^d`$. For current parameterizations of polarized parton densities the combination $`[2\stackrel{~}{H}^u\stackrel{~}{H}^d]/\sqrt{3}`$ for $`\mathrm{\Lambda }`$ production is clearly larger than $`\stackrel{~}{H}^d`$ in the analogous expression for the $`\mathrm{\Sigma }^0`$ channel. Concerning the contribution from $`\stackrel{~}{}`$, the coupling $`g_{KN\mathrm{\Sigma }^0}`$ is about three times smaller than $`g_{KN\mathrm{\Lambda }}`$ appearing in (3), see . Along the same lines one can see that the cross section for $`\gamma ^{}pK^0\mathrm{\Sigma }^+`$ is of similar size as the one for $`\gamma ^{}pK^+\mathrm{\Sigma }^0`$.
The channel $`\gamma ^{}p\pi ^0p`$ does not receive contributions from the pion pole term in (3) because of charge conjugation invariance, so that in our model it is entirely given by the contribution from $`\stackrel{~}{}`$. In Table 1 we see that the combination $`2\stackrel{~}{H}^u+\stackrel{~}{H}^d`$ for $`\pi ^0`$ production is to be compared with $`\sqrt{2}[\stackrel{~}{H}^u\stackrel{~}{H}^d]`$ for $`\gamma ^{}p\pi ^+n`$, which is of comparable size. One thus expects the $`\pi ^0`$ cross section to be similar to the $`|\stackrel{~}{}|^2`$ part of the $`\pi ^+`$ cross section.
The exclusive channels $`\gamma ^{}p\eta p`$ and $`\gamma ^{}p\eta ^{}p`$ involve the combination $`2\stackrel{~}{H}^u\stackrel{~}{H}^d`$ instead of $`2\stackrel{~}{H}^u+\stackrel{~}{H}^d`$ in the $`\pi ^0`$ case, which is somewhat larger because the polarized distributions $`\mathrm{\Delta }u`$ and $`\mathrm{\Delta }d`$ have opposite sign. The strange quark contribution to these channels involves $`\stackrel{~}{H}^s\stackrel{~}{H}^{\overline{s}}`$, which vanishes in our model with polarized parton distributions satisfying $`\mathrm{\Delta }s(x)=\mathrm{\Delta }\overline{s}(x)`$. A quantitative analysis requires the appropriate decay constants for the $`\eta `$ and $`\eta ^{}`$, see for instance , but one can in general expect comparable cross sections for the $`\pi ^0`$, $`\eta `$ and $`\eta ^{}`$ channels.
## 5 Exclusive vector meson production
Within our model for the GPDs, the cross section for vector meson production is sensitive to unpolarized quark and gluon densities. To obtain an indication of uncertainties we have compared results with the LO distributions from CTEQ6 and the LO distributions from MRST2001 . For consistency we need LO rather than NLO parton densities, which unfortunately are not available in several of the most recent parton fits. We have checked that the NLO distributions from MRST2001 are in good agreement with those in the MRST2002 and MRST2004 analyses for quark and antiquark densities down to about $`x10^2`$ and for the gluon density down to about $`x10^1`$. Comparing the LO distributions of CTEQ6 and MRST2001 at a scale $`\mu ^2=1.2\text{GeV}^2`$ (which is the lowest value accepted by the code for the MRST2001 parameterization) we find that the CTEQ6 gluon is larger for $`x\stackrel{<}{}10^1`$ and smaller for $`x\stackrel{>}{}10^1`$. The $`u`$ quark distribution is quite similar in the two sets for $`x\stackrel{>}{}10^2`$, whereas the $`s`$ quark is significantly smaller for CTEQ6. The distributions for $`d`$, $`\overline{u}`$, $`\overline{d}`$ are quite similar for $`x\stackrel{>}{}10^1`$ and larger for CTEQ6 at smaller $`x`$. The LO parameterization of Alekhin has significantly larger $`u`$ and $`\overline{u}`$ distributions and a smaller gluon than the two other sets. At $`\mu ^2=1.2\text{GeV}^2`$ it has however almost no strange quarks in the proton, which we do not consider physically very plausible and which is in clear contrast with the results of CTEQ6, where a dedicated analysis of data constraining the strangeness distribution was performed. Since our study is crucially dependent on the flavor structure of parton distributions, we have therefore not used . Comparing the different parton sets at the higher scale $`\mu ^2=2.5\text{GeV}^2`$ we find a very similar picture.
In the double distribution model of GPDs we take the profile parameter $`b=2`$ for both quark and gluon distributions. For all mesons we take the asymptotic shape of the distribution amplitude, given that no direct experimental information is available for them. Theoretical estimates do not give stronger deviations from the asymptotic form than those we discussed for pseudoscalar mesons, see e.g. the compilation in . In Fig. 5 we show the individual contributions from quark and gluon distributions to the $`\rho ^0`$ cross section as well as their coherent sum. The clear difference between the CTEQ6 and MRST2001 result reflects the current uncertainty on the usual parton densities at low scales in the relevant range of $`x`$. A striking feature is the clear dominance of the gluon distribution up to quite high values of $`x_B`$. Note that with our model of GPDs the convolutions $`^q(\xi ,t)`$ and $`^g(\xi ,t)`$ are sensitive to the forward parton distributions in a certain range of momentum fractions around $`\xi `$ (see App. A). The strong dominance of gluon over quark distributions at small momentum fractions thus still shows its effect at $`x_B`$ values above $`10^1`$. Note that we have taken the same $`t`$ dependence for quark and gluon GPDs in our model (3), lacking phenomenological evidence to the contrary. Comparison of the $`t`$ dependence e.g. for $`\rho ^0`$ and $`\rho ^+`$ production in equal kinematics could be of help here.
In our numerical calculations we have calculated the integrals (A) and (A) for the meson production amplitude with a lower cutoff at momentum fractions $`x=10^4`$. The cross section for $`x_B=0.05`$ changes by less than 2% if we cut off at $`10^5`$ or at $`10^3`$. It is diminished by about 10% with a cutoff at $`0.005`$, which gives an indication of the relevance of momentum fractions which are an order of magnitude smaller than the $`x_B`$ of the process in this model. Similar changes are observed for the individual quark and gluon distributions. We note that if we take a profile parameter $`b=1`$ for quarks, the quark contribution to the cross section at $`x_B=0.05`$ decreases by 10% when moving the cutoff on $`x`$ from $`x=10^4`$ to $`x=10^3`$ and by 35% when moving it from $`x=10^4`$ to $`0.005`$. Such a strong dependence on momentum fractions well below $`x_B`$ seems difficult to understand in physical terms. We note that in the sea quark sector there are no strong theoretical arguments for taking $`b=1`$, see Sect. 4.4 of .
In Fig. 5 we also observe a significant change of the gluon contribution to the cross section when changing the scale of the parton distributions in the double distribution ansatz (A). In contrast, the quark contribution changes by at most a factor of 1.3, reflecting the relatively weak scale evolution of quark and antiquark distributions compared with gluons in the relevant kinematic region. Changing the scale of the forward distributions in the double distribution model (A) gives a rough indication of how the actual GPDs evolve with $`\mu ^2`$ , so that the strong increase with $`\mu ^2`$ seen in Fig. 5 reflects a strong scale uncertainty of the leading-order approximation in $`\alpha _s`$ for channels involving gluon exchange. A full NLO analysis of meson production is possible using the results of but beyond the scope of this work. We will use the smaller scale $`\mu ^2=1.2\text{GeV}^2`$ in our further studies, because the internal virtualities in the hard-scattering graphs of Fig. 1 are typically smaller than $`Q^2`$ (see also the study of relevant scales in the small-$`x`$ limit). Furthermore, the MRST2001 set gives a better description for the ratio of $`\varphi `$ and $`\rho ^0`$ production with our model (see below) and will hence be our default choice in the following.
In Fig. 6 we show the production cross sections for $`\rho ^+`$, $`\rho ^0`$ and $`\omega `$. The $`x_B`$ behavior of the $`\rho ^+`$ cross section roughly follows the one of $`\xi ^2[u(\xi )d(\xi )+\overline{u}(\xi )\overline{d}(\xi )]^2`$, which is a flavor nonsinglet combination and hence does not display the strong rise of sea quarks or gluons at small $`x`$. The clear suppression of $`\omega `$ production compared with the $`\rho ^0`$ is a consequence of the relative factor in the gluon contribution (see Table 1) and at large $`x_B`$ of the relative size of the flavor combination $`2H^uH^d`$ compared with $`2H^u+H^d`$. We remark that the exclusive channel $`\gamma ^{}pf_2p`$ also contributes to semi-inclusive production of $`\pi ^+`$, $`\pi ^{}`$ and $`\pi ^0`$. It involves the combination $`2H^u+H^d[2H^{\overline{u}}+H^{\overline{d}}]`$, where sea quarks drop out, so that one may expect a cross section of similar size as for $`\rho ^+`$ production. A numerical estimate would however require knowledge of quark and gluon distribution amplitudes of the $`f_2`$, see , and is beyond the scope of this work.
Figure 7 shows our results for $`K^+`$ and $`K^0`$ production. In contrast to $`K^+`$ production, the cross section for the $`\mathrm{\Lambda }`$ channel is not much larger than for the $`\mathrm{\Sigma }^0`$ channel. Consulting Table 1 we see that this is because the contributions from $`u`$ and $`d`$ quarks partially cancel in $`2H^uH^d`$ whereas they add in $`2\stackrel{~}{H}^u\stackrel{~}{H}^d`$. We remark that the results for $`K^+`$ and $`K^0`$ production decrease by less than 25% when instead of MRST2001 we take the CTEQ6 parameterization. The uncertainty due to knowledge of the parton distributions is hence much less than for the gluon dominated channels.
Results for $`\varphi `$ production are shown in Fig. 8. The dominance of the gluon over the strange quark contribution is clearly seen, although strange quarks are not entirely negligible with the MRST2001 parameterization.<sup>1</sup><sup>1</sup>1In the study strange quarks were neglected based on inspection of the CTEQ6 parameterization. Since gluons dominate for most $`x_B`$, we see the same trend concerning differences between the parameterizations and the choice of scale as for $`\rho ^0`$ production. The ratio of $`\sigma _L`$ for $`\varphi `$ and $`\rho ^0`$ production is shown in Fig. 9, where the dependence on $`\mu ^2`$ is seen to be much milder since it partially cancels in the ratio. The difference between CTEQ6 and MRST2001 is still significant and mainly due to the difference in the gluon distributions.
Preliminary data from HERMES give a ratio of about 0.08 for the cross sections of $`\varphi `$ and $`\rho ^0`$ production for $`x_B=0.09`$ and $`Q^2=2.46\text{GeV}^2`$ and for $`x_B=0.13`$ and $`Q^2=3.5\text{GeV}^2`$. These data contain a significant contribution from $`\sigma _T`$, and preliminary HERMES data suggest that $`R=\sigma _L/\sigma _T`$ may be slightly smaller for $`\varphi `$ than for $`\rho ^0`$ production at the same $`Q^2`$. The $`\varphi `$ to $`\rho ^0`$ ratio for $`\sigma _L`$ would then be somewhat larger than 0.08. In addition, one can expect that a narrower shape of the distribution amplitude and power corrections due to the strange quark mass would decrease the estimates in Fig. 9 .
A complete representation of GPDs includes in addition to the double distribution the so-called $`D`$-term . It vanishes in the forward limit $`\xi =0`$ and does not affect the double distribution ansatz we are using. Its contribution to the GPDs can be expanded in Gegenbauer polynomials as
$`H_D^q(x,\xi ,t)`$ $`=`$ $`\theta (|x|\xi )\left(1{\displaystyle \frac{x^2}{\xi ^2}}\right){\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}d_{2n+1}^q(t)C_{2n+1}^{3/2}\left({\displaystyle \frac{x}{\xi }}\right),`$
$`H_D^g(x,\xi ,t)`$ $`=`$ $`\theta (|x|\xi ){\displaystyle \frac{3\xi }{2}}\left(1{\displaystyle \frac{x^2}{\xi ^2}}\right)^2{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}d_{2n+1}^g(t)C_{2n}^{5/3}\left({\displaystyle \frac{x}{\xi }}\right)`$ (20)
for $`\xi >0`$. Such terms contribute to the real part of the convolution integrals needed in the meson production amplitudes as
$`I_D^q`$ $`=`$ $`{\displaystyle _1^1}𝑑x{\displaystyle \frac{H_D^q(x,\xi ,t)}{\xi xiϵ}}={\displaystyle _1^1}𝑑x{\displaystyle \frac{H_D^{\overline{q}}(x,\xi ,t)}{\xi xiϵ}}=\mathrm{\hspace{0.25em}2}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}d_{2n+1}^q(t),`$
$`I_D^g`$ $`=`$ $`{\displaystyle _1^1}𝑑x{\displaystyle \frac{H_D^g(x,\xi ,t)}{x}}{\displaystyle \frac{1}{\xi xiϵ}}=\mathrm{\hspace{0.25em}2}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}d_{2n+1}^g(t).`$ (21)
These terms give a $`\xi `$ independent contribution to $`(\xi ,t)`$, in contrast to the contributions from the double distribution part, which very roughly follow the behavior of $`\xi q(\xi )`$, $`\xi \overline{q}(\xi )`$ or $`\xi g(\xi )`$ and hence grow as $`\xi `$ becomes smaller. In the first three coefficients in the quark $`D`$-term at $`t=0`$ have been extracted from a calculation in the chiral quark-soliton model of the nucleon, giving $`d_1^u(0)2.0`$, $`d_3^u(0)0.6`$, $`d_5^u(0)0.2`$ and equal values for $`d`$ quarks, referring to a scale $`\mu `$ of a few GeV . The gluon $`D`$-term is parametrically subleading at the low scale intrinsic to the model, but evolution to $`\mu `$ of a few GeV can give values similar in size to those for quarks. The values $`I_D^u=I_D^d=5.6`$ turn out to be similar in size and opposite in sign to the real parts of the corresponding integrals from the double distribution part in our model. The effect of such a $`D`$-term is however much weaker on the square $`||^2`$ appearing the cross section, which is dominated by the imaginary parts of the integrals in a large region of $`x_B`$. Taking the above values for the quark $`D`$-term and $`I_D^g=11.2`$ as an order-of-magnitude estimate, we find that the change of the various vector meson cross sections is at the 10% level for $`x_B=0.1`$ and not more than a factor 1.5 in either direction for $`x_B`$ below 0.3, where for definiteness we have taken the MRST2001 distributions at $`\mu ^2=1.2\text{GeV}^2`$.
## 6 Comparison with data and discussion of power corrections
Our calculation of meson production cross sections is based on the leading-twist approximation. It is known that corrections in $`1/Q^2`$ to leading-twist meson cross sections can be substantial for $`Q^2`$ of a few $`\text{GeV}^2`$. A systematic treatment of such power corrections remains an unsolved problem. There is however a number of approaches that allow one to model particular sources of power corrections, see e.g. for a discussion and references. For vector meson production, a considerable suppression of the leading-twist result at moderate $`Q^2`$ is found when including in the hard-scattering kernel the transverse momentum of the quarks in the meson . This means that the transverse resolution power of the virtual photon cannot be neglected compared with the transverse size of the meson. Similarly, the finite transverse momentum of the partons coming from the proton gives rise to power corrections, when it is included in the hard-scattering kernel. Estimating both effects in a calculation considering only quark GPDs , a suppression of the leading-twist cross section for $`\rho ^0`$ production by factors of $`3.3`$, $`4.9`$ and $`9.2`$ was found at $`Q^2=5\text{GeV}^2`$ for $`x_B=0.3`$, $`0.45`$ and $`0.6`$, respectively. The recent study for small $`x_B`$, where only gluon GPDs were retained and only the transverse quark momentum in the meson was taken into account, found corresponding suppression factors of $`4.6`$ and $`6.6`$ for respective values of $`x_B=2.95\times 10^3`$ and $`10^4`$ at $`Q^2=4.8\text{GeV}^2`$. For $`Q^2=10.9\text{GeV}^2`$ and $`x_B=4.3\times 10^3`$ this factor decreases to $`1.9`$. The discrepancy of the calculation including power suppression and experimental data is less than 35% in all three cases.<sup>2</sup><sup>2</sup>2The leading-order formula (90) in with which we obtained the suppression factors just quoted contains in addition an approximation for small $`x_B`$, which should however not be the dominant effect comparing to the full calculation.
In Table 2 we compare our leading-oder results for $`\rho ^0`$ and $`\varphi `$ production with data from HERMES . The discrepancy between our calculation and the $`\rho ^0`$ data is well in the range of what can be explained by suppression from quark transverse momentum (considering in addition the uncertainties of our model for the GPDs). The stronger discrepancy with the $`\varphi `$ data corresponds to our overestimating the $`\varphi `$ to $`\rho ^0`$ production ratio, discussed in the previous subsection.
The CLAS collaboration has published results for $`\rho ^0`$ production at $`x_B=0.31`$ and at $`x_B=0.38`$, with $`Q^2`$ values between $`1.5`$ and $`2.3\text{GeV}^2`$ , and for $`\varphi `$ production at $`x_B=0.29`$ and $`Q^2=1.7\text{GeV}^2`$ . We consider that this kinematics, where the hadronic invariant mass $`W`$ is below $`2.3\text{GeV}`$, is too close to threshold for comparison with a leading-twist calculation. A recent CLAS measurement of $`\omega `$ production at higher energies, with $`Q^2`$ up to $`5.1\text{GeV}^2`$ and $`W`$ up to $`2.8\text{GeV}`$, found that helicity conservation between the $`\gamma ^{}`$ and the $`\omega `$ is strongly violated, in contrast with the predicted behavior in the large-$`Q^2`$ limit. This prevented the extraction of $`\sigma _L`$ and was ascribed to a strong contribution from $`\pi ^0`$ exchange (which is absent in the $`\rho ^0`$ and $`\varphi `$ channels).
Let us now turn to $`\pi ^+`$ production, where the situation is quite different. For the contribution from $`\stackrel{~}{}`$ one expects a similar suppression from quark transverse momentum as in the case of $``$ in vector meson production, which was indeed found in the numerical study . The pion exchange contribution from $`\stackrel{~}{}`$ is described in terms of the pion form factor according to (17), and this relation persists beyond the leading approximation in $`1/Q^2`$ to the extent that the pion emitted from the nucleon is not too far from off-shell. The power corrections for $`\stackrel{~}{}`$ are then the same as those for the pion form factor. The leading-twist expression (18) for $`F_\pi (Q^2)`$, with our choice of $`\alpha _s`$ specified after (13), undershoots the data of by a factor 0.53 at $`Q^2=1\text{GeV}^2`$ and a factor $`0.41`$ at $`Q^2=1.6\text{GeV}^2`$. For $`Q^2=3.3\text{GeV}^2`$ we find this factor to be between 0.34 and 0.77 within the large error bars in . We will not attempt here to summarize the detailed theoretical and phenomenological work on the pion form factor, but remark that in addition to the leading-twist perturbative contribution there is a contribution from the Feynman mechanism, where the photon hits a quark carrying almost all of the pion momentum. This leads to a considerable enhancement over the leading-twist approximation. The calculations of $`F_\pi (Q^2)`$ in , which take this effect into account using different methods, give for instance results larger than our leading-twist value by factors between 2 and 4, even at $`Q^2=10\text{GeV}^2`$. Note that these factors are to be squared in the contribution of $`|\stackrel{~}{}|^2`$ to the $`\pi ^+`$ production cross section. For the production of the neutral pseudoscalars $`\pi ^0`$, $`\eta `$, $`\eta ^{}`$, where there is no pion exchange contribution, one expects that power corrections will decrease the cross section, similarly to the case of vector meson production.
We have further compared our leading-twist calculation of $`ϵ\sigma _L`$ for $`\gamma ^{}p\pi ^+n`$ with preliminary data on $`\sigma _T+ϵ\sigma _L`$ from HERMES . The HERMES data are presented for three different bins in $`x_B`$, with the average values of $`Q^2`$ and $`x_B`$ for individual data points ranging from $`1.5\text{GeV}^2`$ and $`0.1`$ to $`4.2\text{GeV}^2`$ and $`0.17`$ in the first bin, from $`2.5\text{GeV}^2`$ and $`0.21`$ to $`6.3\text{GeV}^2`$ and $`0.25`$ in the second bin, and from $`4.5\text{GeV}^2`$ and $`0.34`$ to $`10.5\text{GeV}^2`$ and $`0.45`$ in the third bin . Averaging the ratio between theoretical and experimental cross sections for the data points in each bin, we find ratios of $`0.42`$, $`0.19`$ and $`0.12`$ in the first, second and third bin, respectively. The large discrepancy at large $`Q^2`$ and large $`x_B`$ (which are strongly correlated in the data) is striking, but not too surprising given the size of corrections just estimated for the pion form factor. The much better agreement at smaller $`Q^2`$ and $`x_B`$ might be accidental, given that we expect comparable contributions from $`\stackrel{~}{}`$ and $`\stackrel{~}{}`$, for with the power corrections go in different directions. Help in clarifying this issue could come from the spin asymmetry for transverse target polarization, which gives access to the relative size of $`\stackrel{~}{}`$ and $`\stackrel{~}{}`$ .
The case of $`\pi ^+`$ production (and also the findings in the $`\omega `$ channel mentioned above) show that there are specific power corrections which will not cancel in cross section ratios for different processes. The situation is however better for channels that are sufficiently similar, as the example of $`\rho ^0`$ and $`\varphi `$ production shows. Corrections due to quark transverse momentum (as well as the overall normalization uncertainty from the scale of $`\alpha _s`$ in our leading-order calculation) will tend to cancel in that case. We hence expect that the overall pattern of differences between various meson cross sections we estimated at leading order will not be overturned in a more realistic treatment, given that these differences are largely controlled by the relevant combinations of quark and gluon distributions.
## 7 Exclusive channels in semi-inclusive pion and kaon production
In semi-inclusive hadron production one considers processes of the type
$$e(k)+p(p)e(k^{})+h(q_h)+X,$$
(22)
where $`h`$ is a specified hadron and $`X`$ an unspecified inclusive final state. A basic observable is the distribution of the produced hadron over the variable
$$z=\frac{q_hp}{qp},$$
(23)
which measures the fraction of the virtual photon energy carried by the produced hadron in the target rest frame. In the Bjorken limit of large $`Q^2`$ at fixed $`x_B`$ and fixed $`z`$, semi-inclusive hadron production can be treated within a QCD factorization approach. The differential cross section factorizes into the distribution of partons of type $`i`$ in the target, the cross section for the virtual photon scattering off this parton, and the fragmentation function $`D^{ih}(z)`$ describing the fragmentation of the struck parton into the hadron $`h`$, which carries a fraction $`z`$ of its longitudinal momentum. To leading order in $`\alpha _s`$ one has
$$\frac{d\sigma }{dQ^2dx_Bdz}=2\pi \alpha _{em}\frac{y^2}{1ϵ}\frac{1}{x_BQ^4}\underset{q}{}e_q^2\left[q(x_B)D^{qh}(z)+\overline{q}(x_B)D^{\overline{q}h}(z)\right],$$
(24)
where the sum is over quark flavors. Note that (24) has the same $`Q^2`$-dependence as the inclusive DIS cross section in the Bjorken regime. The fragmentation functions are process-independent and describe not only semi-inclusive DIS but also $`e^+e^{}`$ annihilation into hadrons and the distribution of leading hadrons in high-$`p_T`$ jets. Their scale evolution is governed by evolution equations analogous to the DGLAP equations for the parton distribution functions. Various parameterizations of the fragmentation functions have been presented in the literature, which fit $`e^+e^{}`$ annihilation and semi-inclusive DIS data at higher scales.
The derivation of semi-inclusive factorization relies on the fact that in the Bjorken limit the inclusive final state $`X`$ has a large invariant mass
$$m_X^2=Q^2\frac{1z}{x_B}+m_p^2(qq_h)^2,$$
(25)
and thus a large average multiplicity (note that the squared momentum transfer $`(qq_h)^2`$ is always negative). The semi-inclusive cross section is thus obtained by summing over many individual channels. In practice $`m_X^2`$ is not very large for $`Q^2`$ values of a few $`\text{GeV}^2`$ typical of fixed-target experiments, for instance at Jefferson Lab or HERMES, especially at high $`z`$. At the same time, for moderate $`Q^2`$ the suppression of the cross sections for exclusive channels relative to the semi-inclusive cross section (see below) may not yet be effective. One thus can reach a situation in which the cross sections of individual exclusive channels becomes comparable to the semi-inclusive one. It is interesting to compare the semi-inclusive cross section (24) with the cross sections of exclusive channels contributing to semi-inclusive production. In the following we investigate the role of exclusive channels in semi-inclusive $`\pi `$ and $`K`$ production on a proton target. We study two types of exclusive channels (see Fig. 10) at a quantitative level:
1. direct exclusive production of pseudoscalar mesons, $`\gamma ^{}p\pi ^+n`$ and $`\gamma ^{}pK^+\mathrm{\Lambda }`$,
2. exclusive production of neutral or charged vector mesons $`\rho ,\varphi ,K^{}`$ with subsequent decay into pseudoscalars.
In the direct production of pseudoscalar mesons (Fig. 10a), the energy fraction $`z`$ carried by the produced meson is related to the invariant momentum transfer to the nucleon $`t`$ by
$$1z=x_B\frac{m_B^2m_p^2t}{Q^2}$$
(26)
according to (25). Exclusive production in the limit of large $`Q^2`$ at fixed $`x_B`$ and fixed $`t`$ thus corresponds to values of $`z`$ very close to $`1`$. For example, $`\pi ^+`$ and $`K^+`$ production corresponds to $`z>0.94`$ in typical HERMES kinematics of $`x_B=0.1`$ and $`Q^2=2.5\text{GeV}^2`$ with a maximum momentum transfer $`|t|=1\text{GeV}^2`$ (since the cross section drops rapidly with $`t`$, most events have $`z`$ values yet closer to unity). Such exclusive contributions can usually be separated experimentally from the semi-inclusive events at smaller values of $`z`$.
The situation is different for the contribution to semi-inclusive production resulting from the decay of exclusively produced vector mesons (Fig. 10b). Since the decay products share the energy of the vector meson, such contributions result in an extended $`z`$ distribution for the pion or kaon, even in the Bjorken limit. With the approximations described in Appendix B, the $`z`$ spectrum of the pseudoscalar meson $`P_1`$ from the decay $`VP_1P_2`$ can be written as
$$\frac{d\sigma (epP_1+P_2B)}{dQ^2dx_Bdz}=\frac{\alpha _{em}}{2\pi }\frac{y^2}{1ϵ}\frac{1x_B}{x_BQ^2}\left[ϵ\sigma _L(\gamma ^{}pVB)D_L(z)+\sigma _T(\gamma ^{}pVB)D_T(z)\right]$$
(27)
with
$$D_L(z)=\frac{3}{2\zeta ^3}(zz_0)^2,D_T(z)=\frac{3}{4\zeta ^3}(zz_1)(z_2z).$$
(28)
Here $`z_1zz_2`$ with
$$z_0=\frac{E_{P1}}{m_V},z_1=z_0\zeta ,z_2=z_0+\zeta ,\zeta =\frac{|𝒒_{P1}|}{m_V}$$
(29)
up to corrections of order $`x_Bm_p^2/Q^2`$. For brevity we have not explicitly indicated the dependence of $`D_L`$ and $`D_T`$ on $`x_B`$ and $`Q^2`$ due to these corrections. The energy and three-momentum of $`P_1`$ in the rest frame of the vector meson
$`E_{P1}`$ $`=`$ $`{\displaystyle \frac{m_V^2+m_{P1}^2m_{P2}^2}{2m_V}},`$
$`|𝒒_{P1}|`$ $`=`$ $`{\displaystyle \frac{\sqrt{m_V^4+m_{P1}^4+m_{P2}^42(m_V^2m_{P1}^2+m_V^2m_{P2}^2+m_{P1}^2m_{P2}^2)}}{2m_V}}`$ (30)
depend only on the meson masses, and so do $`z_0`$, $`z_1`$ and $`z_2`$ in the limit of large $`Q^2`$. In particular, the smallest and largest possible values of $`z`$ for pions from the decay $`\rho \pi \pi `$ are $`z_1=0.04`$ and $`z_2=0.96`$ in that limit. The corresponding values for kaons from $`\varphi KK`$ are $`z_1=0.37`$ and $`z_2=0.63`$. For the kaon from $`K^{}`$ decay one has $`z_1=0.32`$ and $`z_2=0.96`$, and for the pion from $`K^{}`$ decay one has $`z_1=0.04`$ and $`z_2=0.68`$.
According to (5), (6) and (27), the contribution of exclusive vector meson production to the cross section $`d\sigma /(dQ^2dx_Bdz)`$ asymptotically scales as $`1/Q^8`$ at fixed $`x_B`$ and $`z`$ and is thus suppressed by $`1/Q^4`$ compared with the leading behavior (24) of the semi-inclusive cross section. Notice that (24) corresponds to transverse photon polarization, with contributions to the longitudinal cross arising at the level of $`\alpha _s`$ and of $`1/Q^2`$ corrections, just as in the familiar case of inclusive DIS. The situation is opposite for hard exclusive meson production, where $`\sigma _L`$ dominates over $`\sigma _T`$ in the large-$`Q^2`$ limit. Measurements show however that at $`Q^2`$ in the few $`\text{GeV}^2`$ region the ratio $`R=\sigma _L/\sigma _T`$ is still of order 1 in $`\rho ^0`$ and $`\varphi `$ production .
We first consider the semi-inclusive production of pions. Depending on the pion charge, exclusive channels contributing here are direct production $`epe\pi ^+n`$ and the production and decay of $`\rho `$ and $`K^{}(892)`$. The $`\rho `$ decays to almost 100% as $`\rho ^0\pi ^+\pi ^{}`$ and $`\rho ^+\pi ^+\pi ^0`$, and the $`K^{}(892)`$ decays to almost 100% into $`K\pi `$, with branching fractions
$`B(K^+K^+\pi ^0)`$ $`=`$ $`\frac{1}{3},B(K^+K^0\pi ^+)=\frac{2}{3},`$
$`B(K^0K^+\pi ^{})`$ $`=`$ $`\frac{2}{3},B(K^0K^0\pi ^0)=\frac{1}{3}`$ (31)
following from isospin symmetry. Note that in quark fragmentation one has $`\sigma (\pi ^0)=\frac{1}{2}[\sigma (\pi ^+)+\sigma (\pi ^{})]`$, which follows directly from the isospin relations between the pion fragmentation functions. This relation also holds for the contributions from $`K^{}`$ decay, but it is strongly violated for $`\rho `$ decay. For the $`\rho ^0`$ this effect was investigated in Ref. in connection with the separation of $`\overline{u}`$ and $`\overline{d}`$ distributions in the proton using semi-inclusive DIS.
In Fig. 11 we show the result of our leading-twist calculation from Sects. 4 and 5 for the $`ep`$ cross section of $`\pi ^+`$ production. We give all $`ep`$ cross section for a lepton beam energy of $`27.5\text{GeV}`$ in the target rest frame, corresponding to the HERMES experiment, and recall that all our exclusive cross sections are calculated with an upper cutoff of $`1\text{GeV}^2`$ on $`|t|`$. We see that the $`\rho ^0`$ channel is clearly dominating. The $`\omega `$, which decays to almost 100% into $`\pi ^+\pi ^{}\pi ^0`$, is much less prominent. According to our discussion in Sects. 4 and 6 one expects that the contribution from $`epe\pi ^0p`$ to $`\pi ^0`$ production is smaller than in the case of direct $`\pi ^+`$ production. The same holds for the production and decay of $`\eta `$ and $`\eta ^{}`$, which have several three-body decays contributing to all three pion channels. As we argued in Sect. 5 the production of $`f_2(1270)`$, which predominantly decays into $`\pi ^+\pi ^{}`$ and $`\pi ^0\pi ^0`$, may contribute at a similar or lower level as $`\rho ^+`$ decays. In Fig. 12 we show the $`z`$ spectrum arising from different vector meson decays. Whereas $`\rho `$ decays contribute in almost the entire $`z`$ range, pions from $`K^{}`$ decays are limited to $`z`$ values below $`0.7`$. Note that due to charge conjugation invariance the $`z`$ spectrum from $`\rho ^0`$ decays is identical for the $`\pi ^+`$ and $`\pi ^{}`$, and by isospin invariance the same holds for the $`\pi ^+`$ and $`\pi ^0`$ spectra from the decay of $`\rho ^+`$. To illustrate the dependence of the $`z`$ spectrum on the ratio $`R`$ of longitudinal and transverse meson production cross sections we have taken values which correspond to the range measured in $`\rho ^0`$ and $`\varphi `$ production at $`Q^2=2.5\text{GeV}^2`$ .
We now compare the contribution from exclusive channels with the semi-inclusive cross section obtained from the leading-order expression (24) for quark fragmentation. We use the LO parton densities from MRST2001 and the LO fragmentation functions of Kretzer , both at a scale $`Q^2=2.5\text{GeV}^2`$. Let us first take a look at the high-$`z`$ tail of the spectrum, where direct exclusive production contributes. Integrating the semi-inclusive cross section for $`\pi ^+`$ production for $`z>0.9`$ at $`Q^2=2.5\text{GeV}^2`$ and $`x_B=0.1`$, we obtain $`d\sigma /(dQ^2dx_B)=0.19\text{nb}\text{GeV}^2`$ from (24). This number should be understood as a naive extrapolation: the fragmentation functions are not well known for $`z`$ close to 1, and in the above kinematics $`z>0.9`$ corresponds to an invariant mass $`m_X<1.84\text{GeV}`$ according to (25), where leading-twist fragmentation can be just marginally valid. Our leading-twist result for $`epe\pi ^+n`$ gives $`d\sigma /(dQ^2dx_B)=0.045\text{nb}\text{GeV}^2`$, which according to our comparison with preliminary HERMES data (Sect. 6) undershoots the actual cross section by a factor of about 0.4. We thus find that for $`z>0.9`$ direct exclusive pion production may be a substantial part of the semi-inclusive cross section, but cannot be more quantitative given the uncertainties just discussed.
For a realistic estimate of exclusive vector meson production we divide our leading-twist results for $`\sigma _L`$ by a factor 7, except for $`\varphi `$ production. According to Table 2 this brings us close to the HERMES measurement for $`\rho ^0`$ production at $`Q^2=2.3\text{GeV}^2`$ and $`x_B=0.1`$, and according to our arguments in Sect. 6 it should give a reasonable estimate for the other channels. In other words, we assume that the ratio of vector meson channels is sufficiently well described by our leading-twist calculation. Only for $`\varphi `$ production do we divide our leading-twist results by a different factor, namely by 15, following our comparison in Table 2 with preliminary HERMES data in this channel. One might argue that for the production of a $`K^{}`$, which has one light and one strange quark, power corrections are between those for the $`\rho ^0`$ and for the $`\varphi `$, but we refrain from such refinements here. Possible changes by a factor of 2 would in fact not change our conclusions regarding the role of $`K^{}`$ decays. For a prediction of $`\sigma _T`$ we divide $`\sigma _L`$ obtained as just described by the value $`R=\sigma _L/\sigma _T=1.2`$ from preliminary HERMES data for $`\rho ^0`$ production in the relevant kinematics , except for the $`\varphi `$ channel, where instead we take $`R=0.8`$ from the parameterization of preliminary HERMES data in . Variation of $`R`$ as shown in Fig. 12 would not affect the conclusions we shall draw.
Integrating the $`\rho ^0`$ decay contribution to $`\pi ^+`$ production for $`z>0.9`$ at $`Q^2=2.5\text{GeV}^2`$ and $`x_B=0.1`$, we find $`d\sigma /(dQ^2dx_B)=0.13\text{nb}\text{GeV}^2`$, which is surprisingly close to the extrapolation of the fragmentation result given above. The fragmentation formula for $`\pi ^{}`$ production gives $`d\sigma /(dQ^2dx_B)=0.07\text{nb}\text{GeV}^2`$ when integrated over $`z>0.9`$, so that in this case the $`\rho ^0`$ contribution slightly overshoots the naive fragmentation result. In Fig. 13 we show the $`z`$ spectrum of semi-inclusive $`\pi ^+`$ and $`\pi ^{}`$ production, comparing the fragmentation result with the contributions from vector meson decays. Following our above discussion we do not show the cross section from fragmentation for $`z`$ above 0.9. We see that $`\rho ^0`$ production gives a sizable contribution to semi-inclusive production for $`z`$ greater than about 0.8. According to our estimate, $`\rho ^+`$ production is suppressed relative to $`\rho ^0`$ by two orders of magnitude and cannot compete with the cross section from quark fragmentation even at large $`z`$. The $`K^{}`$ decay contribution is somewhat larger in size but limited to $`z`$ below 0.7. The fragmentation result for $`\pi ^0`$ production is just the average of $`\pi ^+`$ and $`\pi ^{}`$ because of isospin invariance. With $`\rho ^0`$ decay being absent and the contributions from other vector mesons being comparatively small, we find no exclusive channel that is prominent in semi-inclusive $`\pi ^0`$ production for the kinematics discussed here. We expect direct exclusive production $`epe\pi ^0p`$ to be much less important at high $`z`$ than in the case of $`\pi ^+`$ production, to the extent that power corrections enhance the $`\pi ^+`$ but suppress the $`\pi ^0`$ compared with the leading approximation in $`1/Q^2`$. Given the relative size of cross sections in Fig. 11 it is clear that pions from $`\omega `$ production are significantly smaller than the fragmentation result for all $`z`$, and we shall not analyze the kinematics of the corresponding three-body decay here.
Turning to semi-inclusive $`K^+`$ and $`K^{}`$ production, we show in Fig. 14 the contributions of the relevant exclusive channels to the $`ep`$ cross section, obtained from our leading-twist calculation in Sect. 5. The production of $`\varphi `$, which decays to approximately 50% into $`K^+K^{}`$, is clearly dominant for $`K^+`$ production, and it is the only channel contributing to $`K^{}`$ production. As is seen in the $`z`$-spectra of Fig. 15, it is however only the $`K^+`$ from $`K^{}`$ decays that extends to $`z`$ values above 0.65.
Integrating the leading-order fragmentation formula (24) for $`K^+`$ production for $`z>0.9`$ at $`Q^2=2.5\text{GeV}^2`$ and $`x_B=0.1`$, we obtain $`d\sigma /(dQ^2dx_B)=0.048\text{nb}\text{GeV}^2`$, which is to be understood as an extrapolation as in the pion case discussed above. Our leading-twist estimate for direct exclusive $`K^+`$ production in Sect. 4 gives $`d\sigma /(dQ^2dx_B)=0.016\text{nb}/\text{GeV}^2`$. Following our discussion in Sect. 6 one expects that power corrections will lead to weaker enhancement than in the case of $`epe\pi ^+n`$ (or possibly even to a suppression), because $`\stackrel{~}{}`$ is more important in the leading-twist cross section for $`K^+`$ production than $`\stackrel{~}{}`$. Nevertheless, the above numbers suggest that direct exclusive $`K^+`$ production may be of some significance at the high-$`z`$ end of the spectrum.
If we integrate the $`K^+`$ spectrum from $`K^+`$ and $`K^0`$ decays for $`z>0.9`$, dividing our leading-twist result by 7 and accounting for the transverse cross section as described above, we obtain $`d\sigma /(dQ^2dx_B)=0.0021\text{nb}\text{GeV}^2`$, which is well below our extrapolation from leading-twist fragmentation. In Fig. 16 we compare the fragmentation result for semi-inclusive $`K^+`$ and $`K^{}`$ production with the individual contributions from $`K^{}`$ and $`\varphi `$ decays. We conclude that, even within the uncertainties of our estimates, contributions from $`K^{}`$ production are only a fraction of the fragmentation result at any $`z`$, and that $`\varphi `$ production, despite its larger cross section, is always well below the semi-inclusive cross section. Our finding concerning the $`\varphi `$ contribution agrees with a recent study of measured kaon multiplicities in . On one hand, kaon production by quark fragmentation is less suppressed compared with pion production than is exclusive $`\varphi `$ production compared with production of $`\rho ^0`$, and on the other hand $`\varphi `$ decays only contribute in a $`z`$-range where the fragmentation functions are still large. Apart possibly from direct $`K^+`$ production at $`z`$ close to 1, we thus find no exclusive channel dominating $`K^+`$ or $`K^{}`$ production in typical HERMES kinematics.
So far we have compared exclusive channels with quark fragmentation at $`x_B=0.1`$ and $`Q^2=2.5\text{GeV}^2`$. We have also performed the comparison of Figs. 13 and 16 for $`x_B=0.3`$ at the same $`Q^2`$. For the vector meson cross sections we used the same values of $`R`$ and the same correction factors as for $`x_B=0.1`$, dividing our leading-twist cross sections by 7 for all vector mesons except the $`\varphi `$, where we divide by 15. In doing this, we assume that the leading-twist approximation gives a realistic description of the $`x_B`$ dependence in this region. We find that our qualitative conclusions do not change when going to the larger value of $`x_B`$.
A comment is in order concerning the treatment of exclusive channels in the analysis of semi-inclusive DIS data when extracting quark fragmentation or distribution functions. It is by no means clear that by subtracting contributions of individual exclusive channels from the total yield one obtains an observable more suitable for comparison with the quark fragmentation formulae. In fact, the derivation of factorization theorems relies on the sum over all channels $`X`$ in (22) to be complete. At sufficiently large $`Q^2`$, each individual exclusive channel by itself is a power correction which may or may not be included in the leading-twist analysis. This situation is similar to the one with the contribution from individual nucleon resonances to the cross section of inclusive DIS. A way to think about the relation of exclusive channels to the leading-twist cross section is quark-hadron duality. It remains a challenge to formulate the problem of quark-hadron duality for semi-inclusive DIS in a quantitative fashion.
Symmetry properties like $`\sigma (\pi ^0)=\frac{1}{2}[\sigma (\pi ^+)+\sigma (\pi ^{})]`$ can emerge from summing over many channels which do not fulfill this relation individually. If however a single channel like $`\rho ^0`$ production dominates the semi-inclusive cross section, it clearly becomes more and more difficult for the remaining channels to compensate the missing symmetry between charged and neutral pion production. In such a situation, parton-hadron duality must cease to work. The outcome of our study is that indeed the only channel whose contribution can become dangerously large in HERMES kinematics, is the $`\rho ^0`$ contribution to $`\pi ^+`$ and $`\pi ^{}`$ production.
## 8 Summary
We have evaluated the cross section for a variety of exclusive meson production channels for moderate to large $`x_B`$ at leading order in $`1/Q^2`$ and in $`\alpha _s`$. Cross sections change significantly when varying the nonperturbative input, generalized parton distributions and meson distribution amplitudes, within plausible limits of current model building. On one hand this implies an uncertainty in predicting these cross sections, but on the other hand it implies that their measurement can ultimately help to constrain the nonperturbative functions, provided theoretical control over corrections to the leading-order formulae. We find the largest cross section uncertainties for $`\rho ^0`$, $`\omega `$ and $`\varphi `$ production, which is sensitive to the generalized gluon distribution over a large range of $`x_B`$, reflecting the current uncertainty of the unpolarized gluon density at low scales. A strong dependence on the factorization scale in these channels underlines the need for analysis at next-to-leading order in $`\alpha _s`$. Comparing our leading-twist cross section with experimental data, we confirm that for $`Q^2`$ of a few $`\text{GeV}^2`$ power corrections are substantial. In particular, the suppression of vector meson cross sections we find is consistent with what has been estimated in the literature from the effects of parton transverse momentum in the hard-scattering subprocess. A consistent description of such effects together with next-to-leading order corrections in $`\alpha _s`$ remains a challenge for theory. As is seen for the ratio of $`\varphi `$ and $`\rho ^0`$ production, the most serious theoretical uncertainties cancel however in cross section ratios for sufficiently similar channels (the main distinction being between channels with and without $`t`$-channel pion exchange).
Rescaling our leading-twist cross sections such as to be consistent with experimental data for $`\rho ^0`$ and $`\varphi `$ production, we have compared their contribution to semi-inclusive pion or kaon production with the result of leading-twist quark fragmentation, focusing on the typical kinematics of HERMES measurements, where $`Q^22.5\text{GeV}^2`$ and $`x_B0.1`$. Within large uncertainties, direct exclusive production of $`\pi ^+`$ and possibly $`K^+`$ appears to be comparable with the fragmentation result extrapolated to the bin $`0.9<z<1`$. Through their decays, exclusively produced $`\rho `$, $`\varphi `$ and $`K^{}`$ contribute in a wide range of $`z`$. Pions from $`K^{}`$ decay and kaons from $`\varphi `$ decay are however limited to $`z`$ below 0.7. With this and the relative size of cross sections, our estimates indicate that in typical HERMES kinematics the only exclusive channel whose cross section can compete with quark fragmentation is the $`\rho ^0`$. The $`\rho ^0`$ saturates the quark fragmentation result for semi-inclusive $`\pi ^+`$ and $`\pi ^{}`$ production at large $`z`$. Since the $`\rho ^0`$ does not contribute to $`\pi ^0`$ and to kaon production, there is no corresponding “dangerous” vector channel in these cases.
## Acknowledgments
We are indebted to E.-C. Aschenauer, H. Avakian, A. Borissov, C. Hadjidakis, D. Hasch, A. Hillenbrand, M. Strikman, M. Vanderhaeghen and A. Vinnikov for valuable discussions and information. This work is supported by the Helmholtz Association, contract number VH-NG-004. This work is supported by U.S. Department of Energy Contract DE-AC05-84ER40150, under which the Southeastern Universities Research Association (SURA) operates the Thomas Jefferson National Accelerator Facility.
## Appendix A Integrals over GPDs within the double distribution model
The $`t`$ independent functions in the ansatz (3) for quark and gluon GPDs are modeled as
$`H^q(x,\xi )`$ $`=`$ $`{\displaystyle _1^1}𝑑\beta {\displaystyle _{1+|\beta |}^{1|\beta |}}𝑑\alpha \delta (x\beta \xi \alpha )h(\beta ,\alpha )\left[\theta (\beta )q(\beta )\theta (\beta )\overline{q}(\beta )\right],`$
$`H^g(x,\xi )`$ $`=`$ $`{\displaystyle _1^1}𝑑\beta {\displaystyle _{1+|\beta |}^{1|\beta |}}𝑑\alpha \delta (x\beta \xi \alpha )h(\beta ,\alpha )\beta \left[\theta (\beta )g(\beta )\theta (\beta )g(\beta )\right],`$
$`\stackrel{~}{H}^q(x,\xi )`$ $`=`$ $`{\displaystyle _1^1}𝑑\beta {\displaystyle _{1+|\beta |}^{1|\beta |}}𝑑\alpha \delta (x\beta \xi \alpha )h(\beta ,\alpha )\left[\theta (\beta )\mathrm{\Delta }q(\beta )+\theta (\beta )\mathrm{\Delta }\overline{q}(\beta )\right],`$ (32)
where $`\theta `$ denotes the usual step function, $`q`$, $`\overline{q}`$, $`\mathrm{\Delta }q`$, $`\mathrm{\Delta }\overline{q}`$ the unpolarized and polarized quark and antiquark distributions, and $`g`$ the unpolarized gluon distribution. The profile function
$$h(\beta ,\alpha )=\frac{\mathrm{\Gamma }(2b+2)}{2^{2b+1}\mathrm{\Gamma }^2(b+1)}\frac{[(1|\beta |)^2\alpha ^2]^b}{(1|\beta |)^{2b+1}}$$
(33)
depends on a parameter $`b`$, which we chose to be either $`b=1`$ or $`b=2`$ in this work.
For meson production amplitudes we need the integrals
$`I^q(\xi )`$ $`=`$ $`{\displaystyle _1^1}𝑑x{\displaystyle \frac{H^q(x,\xi )}{\xi xiϵ}},I^{\overline{q}}(\xi )={\displaystyle _1^1}𝑑x{\displaystyle \frac{H^{\overline{q}}(x,\xi )}{\xi xiϵ}}=[I^q(\xi )]^{},`$
$`\stackrel{~}{I}^q(\xi )`$ $`=`$ $`{\displaystyle _1^1}𝑑x{\displaystyle \frac{\stackrel{~}{H}^q(x,\xi )}{\xi xiϵ}},\stackrel{~}{I}^{\overline{q}}(\xi )={\displaystyle _1^1}𝑑x{\displaystyle \frac{\stackrel{~}{H}^{\overline{q}}(x,\xi )}{\xi xiϵ}}=[\stackrel{~}{I}^q(\xi )]^{},`$ (34)
where we used the definitions $`H^{\overline{q}}(x,\xi )=H^q(x,\xi )`$ and $`\stackrel{~}{H}^{\overline{q}}(x,\xi )=\stackrel{~}{H}^q(x,\xi )`$ together with the fact that these functions are even in $`\xi `$. The required integral for gluons can be brought into the same form as $`I^q(\xi )`$ by rewriting
$$I^g(\xi )=_1^1𝑑x\frac{H^g(x,\xi )}{x}\frac{1}{\xi xiϵ}=\frac{1}{\xi }_1^1𝑑x\frac{H^g(x,\xi )}{\xi xiϵ},$$
(35)
where we used that $`H^g(x,\xi )`$ is even in $`x`$. The imaginary parts of these integrals are readily converted into integrals over $`\beta `$, with
$`\mathrm{Im}I^q(\xi )`$ $`=`$ $`{\displaystyle _0^{\frac{2\xi }{1+\xi }}}𝑑\beta I(\beta ,\xi )q(\beta ),`$
$`I(\beta ,\xi )`$ $`=`$ $`{\displaystyle \frac{\pi \mathrm{\Gamma }(2b+2)}{2^{2b+1}\mathrm{\Gamma }^2(b+1)}}{\displaystyle \frac{(1\xi ^2)^b}{\xi ^{2b+1}}}{\displaystyle \frac{1}{(1\beta )^{2b+1}}}\left({\displaystyle \frac{2\xi }{1+\xi }}\beta \right)^b\beta ^b`$ (36)
for $`\xi >0`$. The function $`I(\beta ,\xi )`$ vanishes at the endpoints of the integration region, which in particular ensures the convergence of the integral at $`\beta =0`$ for common parameterizations of quark densities. To obtain the analogous expressions for $`I^{\overline{q}}`$, $`\stackrel{~}{I}^q`$, $`\stackrel{~}{I}^{\overline{q}}`$ and $`I^g`$ one has to replace $`q(\beta )`$ with $`\overline{q}(\beta )`$, $`\mathrm{\Delta }q(\beta )`$, $`\mathrm{\Delta }\overline{q}(\beta )`$ and $`\beta g(\beta )`$, respectively.
The real parts of the amplitudes involve principal value integrals, whose numerical evaluation requires some care, especially for small $`\xi `$. For our choices of profile parameters $`b=1`$ and $`b=2`$ one can explicitly perform the $`\alpha `$ integral after inserting (A) into (A) and (35). The result is
$`\mathrm{Re}I^q(\xi )`$ $`=`$ $`{\displaystyle _0^1}𝑑\beta \left[R(\beta ,\xi )q(\beta )+R(\beta ,\xi )\overline{q}(\beta )\right]`$
$`=`$ $`{\displaystyle _0^1}𝑑\beta \left\{R(\beta ,\xi )\left[q(\beta )\overline{q}(\beta )\right]+\left[R(\beta ,\xi )+R(\beta ,\xi )\right]\overline{q}(\beta )\right\},`$
$`\mathrm{Re}I^g(\xi )`$ $`=`$ $`{\displaystyle _0^1}𝑑\beta \left[R(\beta ,\xi )+R(\beta ,\xi )\right]\beta g(\beta ),`$
$`\mathrm{Re}\stackrel{~}{I}^q(\xi )`$ $`=`$ $`{\displaystyle _0^1}𝑑\beta \left[R(\beta ,\xi )\mathrm{\Delta }q(\beta )R(\beta ,\xi )\mathrm{\Delta }\overline{q}(\beta )\right],`$ (37)
with
$`R(\beta ,\xi )`$ $`\stackrel{b=1}{=}`$ $`{\displaystyle \frac{3}{4\xi ^3(1\beta )^3}}(\mathrm{\hspace{0.17em}2}\xi (1\beta )(\beta \xi )`$
$`+\beta (1\xi )[\beta (1+\xi )2\xi ]\mathrm{log}{\displaystyle \frac{|\beta (1+\xi )2\xi |}{\beta (1\xi )}}),`$
$`R(\beta ,\xi )`$ $`\stackrel{b=2}{=}`$ $`{\displaystyle \frac{5}{16\xi ^5(1\beta )^5}}(\mathrm{\hspace{0.17em}2}\xi (1\beta )(\beta \xi )[3(\beta \xi )^25\xi ^2(1\beta )^2]`$ (38)
$`+3\beta ^2(1\xi )^2[\beta (1+\xi )2\xi ]^2\mathrm{log}{\displaystyle \frac{|\beta (1+\xi )2\xi |}{\beta (1\xi )}}).`$
For both $`b=1`$ and $`b=2`$, the function $`R(\beta ,\xi )`$ is continuous in the full interval of integration, with finite limits at $`\beta =0`$ and $`\beta =1`$. If $`\xi >0`$ it is positive for $`\beta <\xi `$ and negative for $`\beta >\xi `$, and if $`\xi <0`$ it is negative in the entire interval. Convergence of the integral for polarized quark distributions requires that $`\mathrm{\Delta }q(\beta )`$ and $`\mathrm{\Delta }\overline{q}(\beta )`$ have integrable singularities at $`\beta =0`$, which is the case for the parton densities we use in this study. The unpolarized quark distributions have a steeper behavior at small $`\beta `$, but since $`R(\beta ,\xi )+R(\beta ,\xi )\beta `$ for $`\beta 0`$ it is sufficient to have integrable singularities for $`q(\beta )\overline{q}(\beta )`$ and for $`\beta \overline{q}(\beta )`$.
In Fig. 17 we illustrate the behavior of the functions multiplying the parton distributions in the integrals (A) and (A). The imaginary part of the amplitude involves momentum fractions in the parton densities between $`0`$ and $`2\xi /(1+\xi )=x_B`$, with a maximum of the shape function $`I(\beta ,\xi )`$ for $`\beta `$ around $`\xi `$. In contrast, the real part is sensitive to higher momentum fractions, with a partial cancellation from values above and below $`\xi `$. One also clearly sees the stronger sensitivity to small $`\beta `$ if $`b=1`$. Note that the functions shown in the figure will be multiplied in the amplitude with functions showing a strong rise towards $`\beta =0`$.
## Appendix B Distribution of pions or kaons from vector meson decay
In this appendix we discuss the decay of a vector meson into two pseudoscalar mesons and derive the $`z`$ distribution given in (27). Consider the contribution of $`epVB`$ with subsequent decay $`VP_1P_2`$ to semi-inclusive production $`epP_1+X`$. A useful set of variables to describe the decay of the vector meson are the polar and azimuthal angles $`\theta `$ and $`\phi `$ of $`P_1`$ in the vector meson center-of-mass, as shown in Fig. 18. The distribution in these angles is connected in a straightforward way with the spin density matrix of the produced vector meson , and the phase space element has a factorized form in the variables $`Q^2`$, $`x_B`$, $`t`$ and $`\theta `$, $`\phi `$. The variable $`z`$ used for semi-inclusive production of $`P_1`$ is then given by
$$z=a+b\mathrm{cos}\theta +c\mathrm{sin}\theta \mathrm{cos}\phi $$
(39)
with
$`a`$ $`=`$ $`{\displaystyle \frac{E_{P1}}{m_V}}{\displaystyle \frac{r_2(1+2x_Bm_p^2/Q^2)+r_3\sqrt{1+4x_B^2m_p^2/Q^2}}{2r_1}}\left[1+O(x_B\mathrm{\Delta }_T^2/Q^2)\right]{\displaystyle \frac{E_{P1}}{m_V}},`$
$`b`$ $`=`$ $`{\displaystyle \frac{|𝒒_{P1}|}{m_V}}{\displaystyle \frac{r_3(1+2x_Bm_p^2/Q^2)+r_2\sqrt{1+4x_B^2m_p^2/Q^2}}{2r_1}}\left[1+O(x_B\mathrm{\Delta }_T^2/Q^2)\right]{\displaystyle \frac{|𝒒_{P1}|}{m_V}},`$
$`c`$ $`=`$ $`{\displaystyle \frac{|𝒒_{P1}|\mathrm{\Delta }_T}{Q^2}}{\displaystyle \frac{2x_B}{1x_B}}{\displaystyle \frac{\sqrt{1+4x_B^2m_p^2/Q^2}}{r_3}},`$ (40)
where we abbreviated
$`r_1`$ $`=`$ $`1+{\displaystyle \frac{x_B}{1x_B}}{\displaystyle \frac{m_p^2}{Q^2}},r_2=\mathrm{\hspace{0.25em}1}+{\displaystyle \frac{x_B}{1x_B}}{\displaystyle \frac{m_V^2m_B^2+m_p^2}{Q^2}},`$
$`r_3`$ $`=`$ $`\left[\left(1{\displaystyle \frac{x_B}{1x_B}}{\displaystyle \frac{m_V^2+m_B^2m_p^2}{Q^2}}\right)^2\left({\displaystyle \frac{x_B}{1x_B}}{\displaystyle \frac{2m_Vm_B}{Q^2}}\right)^2\right]^{1/2}.`$ (41)
The energy and momentum $`E_{P1}`$ and $`|𝒒_{P1}|`$ of $`P_1`$ in the rest frame of $`V`$ have already been given in (7), and $`\mathrm{\Delta }_T`$ is the transverse momentum of the scattered baryon with respect to the initial proton in the $`\gamma ^{}p`$ center-of-mass (see Fig. 18). The approximate expressions in (B) are valid up to relative corrections of order $`x_Bm_p^2/Q^2`$ and $`x_B\mathrm{\Delta }_T^2/Q^2`$, and to the same accuracy one has $`\mathrm{\Delta }_T^2=(1x_B)(t_0t)`$. Changing variables from $`\theta `$ to $`z`$ gives for the cross section
$$\frac{d\sigma (epP_1+P_2B)}{dQ^2dx_Bdtd\phi dz}=\frac{1}{bc\mathrm{cot}\theta \mathrm{cos}\phi }\frac{d\sigma (epP_1+P_2B)}{dQ^2dx_Bdtd\phi d\mathrm{cos}\theta }.$$
(42)
In Bjorken kinematics one has $`cb`$ and can replace the Jacobian in (42) by $`1/b`$ (except in the small region where $`\mathrm{sin}\theta c/b`$, which is not relevant for our purposes). Neglecting $`\mathrm{\Delta }_T`$ we get $`z=a+b\mathrm{cos}\theta `$ with $`a`$ and $`b`$ evaluated at $`\mathrm{\Delta }_T=0`$, and integration over $`t`$ and $`\phi `$ gives
$$\frac{d\sigma (epP_1+P_2B)}{dQ^2dx_Bdz}=\frac{3}{4b^3}\left[2(za)^2\frac{d\sigma (epV_LB)}{dQ^2dx_B}+(za+b)(a+bz)\frac{d\sigma (epV_TB)}{dQ^2dx_B}\right]$$
(43)
in terms of the cross sections for the production of longitudinally or transversely polarized vector mesons. Using $`s`$-channel helicity conservation, which is experimentally seen to hold at the few 10% level in $`\rho ^0`$ and $`\varphi `$ production , these cross sections respectively correspond to the production from longitudinally or transversely polarized virtual photons, and we finally obtain (27). In our numerical applications we have used the exact expressions from (B) and (B) at $`\mathrm{\Delta }_T=0`$, and thus in particular neglected $`c`$. Since the integrated cross sections are dominated by small $`\mathrm{\Delta }_T`$, this should be a very good approximation for the values of $`Q^2`$ and $`x_B`$ we focus on in the present study. The inclusion of finite $`\mathrm{\Delta }_T`$ effects in the kinematics would considerably complicate any analysis. |
warning/0506/cs0506042.html | ar5iv | text | # Tree-Based Construction of LDPC Codes
## I Introduction
Low Density Parity Check (LDPC) codes are widely acknowledged to be good codes due to their near Shannon-limit performance when decoded iteratively. However, many structure-based constructions of LDPC codes fail to achieve this level of performance, and are often outperformed by random constructions. (Exceptions include the finite-geometry-based LDPC codes (FG-LDPC) of , which were later generalized in .) Moreover, there are discrepancies between iterative and maximum likelihood (ML) decoding performance of short to moderate blocklength LDPC codes. This behavior has recently been attributed to the presence of so-called pseudocodewords of the LDPC constraint graphs, which are valid solutions of the iterative decoder which may or may not be optimal . Analogous to the role of minimum Hamming distance, $`d_{\mathrm{min}}`$, in ML-decoding, the minimal pseudocodeword weight, $`w_{\mathrm{min}}`$, has been shown to be a leading predictor of performance in iterative decoding. Furthermore, the error floor performance of iterative decoding is dominated by minimal weight pseudocodewords. Although there exist pseudocodewords with weight larger than $`d_{\mathrm{min}}`$ that have adverse affects on decoding, pseudocodewords with weight $`w_{\mathrm{min}}<d_{\mathrm{min}}`$ are especially problematic .
The Type I-A construction and certain cases of the Type II construction presented in this paper are designed so that the resulting codes have minimal pseudocodeword weight equal to the minimum distance of the code, and consequently, these problematic low-weight pseudocodewords are avoided. The resulting codes have minimum distance which meets the lower tree bound originally presented in , and since $`w_{\mathrm{min}}`$ shares the same lower bound , and is upper bounded by $`d_{\mathrm{min}}`$, the proposed constructions have $`w_{\mathrm{min}}=d_{\mathrm{min}}`$. It is worth noting that this property is also a characteristic of some of the FG -LDPC codes , and indeed, the projective-geometry-based codes of arise as special cases of our Type II construction. Furthermore, the Type I-B construction presented herein is a modification of the Type I-A construction, and it yields a family of codes with a wide range of rates and blocklengths that are comparable to those obtained from finite geometries.
We now present the tree bound on $`w_{\mathrm{min}}`$ derived in .
###### Theorem I.1
Let $`G`$ be a bipartite LDPC constraint graph with smallest left (variable node) degree $`d`$ and girth $`g`$. Then the minimal pseudocodeword weight $`w_{\mathrm{min}}`$ (for the AWGN/BSC channels) is lower bounded by
$$w_{\mathrm{min}}\{\begin{array}{cc}1+d+d\left(d1\right)+d\left(d1\right)^2+\mathrm{}+d\left(d1\right)^{\frac{g6}{4}},& \frac{g}{2}\text{ odd }\\ 1+d+d\left(d1\right)+\mathrm{}+d\left(d1\right)^{\frac{g8}{4}}+\left(d1\right)^{\frac{g4}{4}},& \frac{g}{2}\text{ even }\end{array}$$
This bound is also the tree bound on the minimum distance established by Tanner in . And since the set of pseudocodewords includes all codewords, we have $`w_{\mathrm{min}}d_{\mathrm{min}}`$. In the following sections we present two construction techniques of LDPC codes wherein for certain cases, $`w_{\mathrm{min}}=d_{\mathrm{min}}`$.
## II preliminaries
The connection algorithms for the tree constructions Type I-B and Type II are based on mutually orthogonal Latin squares. A well-known construction of a family of mutually orthogonal Latin squares of order $`p^s`$, a power of a prime, may be found in . Let $`M^{(1)},M^{(2)},\mathrm{},M^{(p^s1)}`$ denote $`p^s1`$ mutually orthogonal Latin squares (MOLS) of order $`p^s`$. Let the rows (and columns) of each square be indexed by the integers $`0,1,2,\mathrm{},p^s1`$. Without loss of generality, assume that the first column of each of the Latin squares is $`[0,1,2,\mathrm{},p^s1]^T`$. In addition, define a new square of size $`p^s\times p^s`$, denoted $`M^{(0)}`$, where each column of $`M^{(0)}`$ is $`[0,1,2,\mathrm{},p^s1]^T`$.
## III Tree-based Construction: Type I
In the Type I construction, first a $`d`$-regular tree of alternating variable and constraint node layers is enumerated from a root variable node (layer $`L_0`$) for $`\mathrm{}`$ layers. If $`\mathrm{}`$ is odd (respectively, even), the final layer $`L_\mathrm{}1`$ is composed of variable nodes (respectively, constraint nodes). Call this tree $`T`$. The tree $`T`$ is then reflected across an imaginary horizontal axis to yield another tree, $`T^{}`$, and the variable and constraint nodes are reversed. That is, if layer $`L_i`$ in $`T`$ is composed of variable nodes, then the reflection of $`L_i`$, call it $`L_i^{}`$, is composed of constraint nodes in the reflected tree, $`T^{}`$. The union of these two trees, along with edges connecting the nodes in layers $`L_\mathrm{}1`$ and $`L_\mathrm{}1^{}`$ according to a connection algorithm that is described next, comprise the graph representing a Type I LDPC code. We now present two connection schemes that can be used in this Type I model, and discuss the resulting codes.
### III-A Type I-A
For $`d=3`$, the Type I-A construction yields a $`d`$-regular LDPC constraint graph having $`1+d+d(d1)+\mathrm{}+d(d1)^{\frac{g4}{2}}`$ variable and constraint nodes, and girth $`g`$. The tree $`T`$ has $`\frac{g}{2}`$ layers. To connect the nodes in $`L_{\frac{g}{2}1}`$ to $`L_{\frac{g}{2}1}^{}`$, first label the variable (resp., constraint) nodes in $`L_{\frac{g}{2}1}`$ (resp., $`L_{\frac{g}{2}1}^{}`$) when $`\frac{g}{2}`$ is odd, as $`v_0,v_1,\mathrm{},v_{2^{\frac{g}{2}2}1}`$, $`v_{2^{\frac{g}{2}2}},\mathrm{},v_{22^{\frac{g}{2}2}1},v_{22^{\frac{g}{2}2}},\mathrm{},v_{32^{\frac{g}{2}2}1}`$ (resp., $`c_0,c_1,\mathrm{},c_{32^{\frac{g}{2}2}1}`$). The nodes $`v_0,v_1,\mathrm{},v_{2^{\frac{g}{2}2}1}`$ form the $`0^{th}`$ class, the nodes $`v_{2^{\frac{g}{2}2}},\mathrm{},v_{22^{\frac{g}{2}2}1}`$ form the $`1^{st}`$ class, and the nodes $`v_{22^{\frac{g}{2}2}},\mathrm{},v_{32^{\frac{g}{2}2}1}`$ form the $`2^{nd}`$ class; classify the constraint nodes in a similar manner. In addition, define three permutations $`\pi (),\tau (),\tau ^{}()`$ of the set $`\{0,1,\mathrm{},2^{\frac{g}{2}2}1\}`$ as follows. The nodes in $`L_{\frac{g}{2}1}`$ and $`L_{\frac{g}{2}1}^{}`$ are connected as follows:
1. For $`i=0,1`$, and $`j=0,1,\mathrm{},2^{\frac{g}{2}2}1`$, the variable node $`v_{j+i2^{\frac{g}{2}2}}`$ is connected to nodes $`c_{\pi (j)+i2^{\frac{g}{2}2}}`$ and $`c_{\tau (j)+(i+1)2^{\frac{g}{2}2}}`$.
2. For $`i=2`$, and $`j=0,1,\mathrm{},2^{\frac{g}{2}2}1`$, the variable node $`v_{j+i2^{\frac{g}{2}2}}`$ is connected to nodes $`c_{\pi (j)+22^{\frac{g}{2}2}}`$ and $`c_{\tau ^{}(j)}`$.
The permutations for the cases $`g=6,8,10,12`$ are given below. The above construction can be extended for higher $`g`$ in a natural way and we are working on an explicit closed form expression for the permutations $`\pi ,\tau ,\tau ^{}`$ for higher $`g`$.
$$g=6,\pi =\tau =\tau ^{}=\left(0\right)\left(1\right),\text{the identity permutation.}$$
$$g=8,\pi =\left(0\right)\left(2\right)(1,3),\tau =\left(0\right)\left(2\right)(1,3),\tau ^{}=(0,2)\left(1\right)\left(3\right).$$
$$g=10,\pi =\left(0\right)\left(2\right)\left(4\right)\left(6\right)(1,5)(3,7),\tau =\left(0\right)\left(2\right)\left(4\right)\left(6\right)(1,7)(3,5),$$
$$\tau ^{}=(0,4)(2,6)(1,3)(5,7).$$
$$g=12,\pi =\left(0\right)\left(4\right)\left(8\right)\left(12\right)(2,6)(10,14)(1,9)(3,15)(5,13)(7,11),$$
$$\tau =\left(0\right)(4,12)\left(8\right)(2,6,10,14)(1,15,13,11)(3,9,7,5),$$
$$\tau ^{}=(0,8)(4,12)(2,14)(6,10)(1,3,5,7)(9,11,13,15).$$
When $`\frac{g}{2}`$ is odd, the minimum distance of the resulting code meets the tree bound, and hence, $`d_{\mathrm{min}}=w_{\mathrm{min}}`$. When $`\frac{g}{2}`$ is even, $`d_{\mathrm{min}}`$ is strictly larger than the tree bound; we believe however, that $`w_{\mathrm{min}}`$ is equal to $`d_{\mathrm{min}}`$ in this case as well. Figure 1 illustrates the general construction procedure and Figure 2 shows a girth 10 Type I-A LDPC constraint graph.
### III-B Type I-B
For $`d=p^s,p`$ a prime, the Type I-B construction yields a $`d`$-regular LDPC constraint graph having $`1+d+d(d1)`$ variable and constraint nodes, and girth $`6`$. The tree $`T`$ has 3 layers $`L_0,L_1,`$ and $`L_2`$. $`L_2`$ (resp., $`L_2^{}`$) is composed of $`p^s`$ sets $`\{S_i\}_{i=0}^{p^s1}`$ of $`p^s1`$ variable (resp., constraint) nodes in each set; the set $`S_i`$ corresponds to the children of branch $`i`$ of the root node. Let $`S_i`$ (resp., $`S_i^{}`$) contain the variable (resp., constraint) nodes $`v_{i,1},v_{i,2},\mathrm{},v_{i,p^s1}`$ (resp., $`c_{i,1},c_{i,2},\mathrm{},c_{i,p^s1}`$). To use MOLS of order $`p^s`$ in the connection algorithm, an imaginary node, $`v_{i,0}`$ (resp., $`c_{i,0}`$) is temporarily introduced into each set $`S_i`$ (resp, $`S_i^{}`$). The connection algorithm proceeds as follows:
1. Let $`x_t(i,j)`$ denote the $`(j,t)^{th}`$ entry of the square $`M^{(i)}`$ defined in Section II. For $`i=0,\mathrm{},p^s1`$ and $`j=0,\mathrm{},p^s1`$, connect variable node $`v_{i,j}`$ to constraint nodes $`c_{0,x_0(i,j)},c_{1,x_1(i,j)},\mathrm{},c_{p^s1,x_{p^s1}(i,j)}`$.
2. Delete all imaginary nodes $`\{v_{i,0},c_{i,0}\}_{i=0}^{p^s1}`$ and the edges incident on them.
3. For $`i=1,\mathrm{},p^s1,`$ delete the edge connecting $`v_{0,i}`$ to $`c_{0,i}`$.
The resulting $`d`$-regular constraint graph represents the Type I-B LDPC code. Figure 3 illustrates the construction procedure and Figure 4 provides a specific example of a Type I-B LDPC constraint graph with $`d=4`$; the squares used for constructing this graph are
$$\left[\begin{array}{cccc}0& 0& 0& 0\\ 1& 1& 1& 1\\ 2& 2& 2& 2\\ 3& 3& 3& 3\end{array}\right],\left[\begin{array}{cccc}0& 1& 2& 3\\ 1& 0& 3& 2\\ 2& 3& 0& 1\\ 3& 2& 1& 0\end{array}\right],\left[\begin{array}{cccc}0& 2& 3& 1\\ 1& 3& 2& 0\\ 2& 0& 1& 3\\ 3& 1& 0& 2\end{array}\right],\left[\begin{array}{cccc}0& 3& 1& 2\\ 1& 2& 0& 3\\ 2& 1& 3& 0\\ 3& 0& 2& 1\end{array}\right].$$
The Type I-B algorithm yields LDPC codes having a wide range of rates and blocklengths that are comparable to, but different from, the two-dimensional LDPC codes from finite Euclidean geometries . The Type I-B LDPC codes are $`p^s`$-regular with girth six, blocklength $`N=p^{2s}+1`$, and distance $`d_{\mathrm{min}}p^s+1`$. For degrees of the form $`d=2^s`$, the resulting Type I-B codes have very good rates, above 0.5, and perform well with iterative decoding.
## IV Tree-based Construction: Type II
In the Type II construction, first a $`d`$-regular tree $`T`$ of alternating variable and constraint node layers is enumerated from a root variable node (layer $`L_0`$) for $`\mathrm{}`$ layers, as in Type I. The tree $`T`$ is not reflected; rather, a single layer of $`(d1)^\mathrm{}1`$ nodes is added to form layer $`L_{\mathrm{}}`$. If $`\mathrm{}`$ is odd (resp., even), this layer is composed of constraint (resp., variable) nodes. The union of $`T`$ and $`L_{\mathrm{}}`$, along with edges connecting the nodes in layers $`L_\mathrm{}1`$ and $`L_{\mathrm{}}`$ according to a connection algorithm that is described next, comprise the graph representing a Type II LDPC code. We now present the connection scheme that is used for this Type II model, and discuss the resulting codes.
The connection algorithm for $`\mathrm{}=3`$ and $`\mathrm{}=4`$ proceeds as follows.
### IV-A $`\mathrm{}=3`$
The $`d`$ constraint nodes in $`L_1`$ are labeled $`B_0,B_1,\mathrm{},B_{p^s}`$ to represent the $`d`$ branches stemming from the root node, and the $`d(d1)`$ variable nodes in the third layer $`L_2`$ are labeled as $`B_{0,0},B_{0,1},\mathrm{},B_{0,p^s1}`$, $`B_{1,0},\mathrm{},B_{1,p^s1}`$, $`\mathrm{}`$, $`B_{p^s,0},\mathrm{},B_{p^s,p^s1}`$. The $`p^{2s}`$ constraint nodes in the final layer $`L_{\mathrm{}}=L_3`$ are labeled $`A_{0,0},A_{0,1},\mathrm{},A_{0,p^s1}`$, $`A_{1,0},A_{1,1},\mathrm{},A_{1,p^s1}`$, $`\mathrm{}`$, $`A_{p^s1,0},A_{p^s1,1},\mathrm{},A_{p^s1,p^s1}`$.
1. The constraint nodes in $`L_3`$ are grouped into $`d1=p^s`$ classes of $`d1=p^s`$ nodes in each class. Similarly, the variable nodes in $`L_2`$ are grouped into $`d=p^s+1`$ classes of $`d1=p^s`$ nodes in each class. Those nodes descending from $`B_0`$ form the $`0^{th}`$ class, those descending from $`B_1`$ form the first class, and so on.
2. Each of the variable nodes descending from $`B_0`$ is connected to all the constraint nodes of one class. That is, $`B_{0,0}`$ is connected to $`A_{0,0},A_{0,1},\mathrm{},A_{0,p^s1}`$, $`B_{0,1}`$ is connected to $`A_{1,0},A_{1,1},\mathrm{},A_{1,p^s1}`$, and in general, $`B_{0,k}`$ is connected to $`A_{k,0},A_{k,1},\mathrm{},A_{k,p^s1}`$ which correspond to the constraint nodes in the $`k^{th}`$ class.
3. Let $`x_t(i,j)`$ denote the $`(j,t)^{th}`$ entry of $`M^{(i1)}`$.
4. Then connect the variable node $`B_{i,j}`$ to the constraint nodes
$$A_{0,x_0(i,j)},A_{1,x_1(i,j)},A_{2,x_2(i,j)},\mathrm{},A_{p^s1,x_{p^s1}(i,j)}.$$
Figure 5 illustrates the construction procedure and Figure 6 provides an example of a Type II LDPC constraint graph with degree $`d=4`$ and girth $`g=6`$; the squares used for constructing this example are
$$M^{\left(0\right)}=\left[\begin{array}{ccc}0& 0& 0\\ 1& 1& 1\\ 2& 2& 2\end{array}\right],M^{\left(1\right)}=\left[\begin{array}{ccc}0& 1& 2\\ 1& 2& 0\\ 2& 0& 1\end{array}\right],M^{\left(2\right)}=\left[\begin{array}{ccc}0& 2& 1\\ 1& 0& 2\\ 2& 1& 0\end{array}\right].$$
The ratio of minimum distance to blocklength of the codes is at least $`\frac{2+p^s}{1+p^s+p^{2s}}`$, and the girth is six. For degrees $`d`$ of the form $`d=2^s+1`$, the tree bound on minimum distance and minimum pseudocodeword weight is met, i.e., $`d_{\mathrm{min}}=w_{\mathrm{min}}=2+2^s`$, for the Type II, $`\mathrm{}=3`$, LDPC codes.
### IV-B Relation to finite geometry codes
The codes that result from this $`\mathrm{}=3`$ construction correspond to the two-dimensional projective-geometry-based LDPC (PG LDPC) codes of . With a little modification of the Type II construction, we can also obtain the two-dimensional Euclidean-geometry-based LDPC codes of .
Since a two-dimensional Euclidean geometry may be obtained by deleting certain points and line(s) of a two-dimensional projective geometry, the graph of a two-dimensional EG-LDPC code may be obtained by performing the following operations on the Type II, $`\mathrm{}=3`$, graph:
1. In the tree $`T`$, the root node along with its neighbors, i.e., the constraint nodes in layer $`L_1`$, are deleted.
2. Consequently, the edges from the constraint nodes $`B_0,\mathrm{},B_{p^s}`$ to layer $`L_2`$ are also deleted.
3. At this stage, the remaining variable nodes have degree $`p^s`$, and the remaining constraint nodes have degree $`p^s+1`$. Now, a constraint node from layer $`L_3`$ is chosen, say, constraint node $`A_{0,0}`$. This node and its neighboring variable nodes and the edges incident on them are deleted. Doing so removes exactly one variable node from each class of $`L_2`$, and the degrees of the remaining constraint nodes in $`L_3`$ are lessened by one. Thus, the resulting graph is now $`p^s`$-regular with a girth of six, has $`p^{2s}1`$ constraint and variable nodes , and corresponds to the two-dimensional Euclidean-geometry-based LDPC code $`EG(2,p^s)`$ of .
### IV-C $`\mathrm{}=4`$
1. The tree $`T`$ is now enumerated for four layers, with the nodes in $`L_0,L_1,`$ and $`L_2`$ labeled as in the $`\mathrm{}=3`$ case. For $`i=0,\mathrm{},p^s`$, the constraint nodes in the $`i`$th class of $`L_3`$ are labeled $`B_{i,0,0},B_{i,0,1},\mathrm{},B_{i,0,p^s1},B_{i,1,0},B_{i,1,1},\mathrm{},B_{i,1,p^s1}`$, $`\mathrm{}`$, $`B_{i,p^s1,0},\mathrm{},B_{i,p^s1,p^s1}`$.
2. The $`p^{3s}`$ variable nodes in the final layer $`L_4`$ are labeled $`A_{0,0,0},A_{0,0,1},\mathrm{},A_{0,0,p^s1},A_{0,1,0},A_{0,1,1},\mathrm{},A_{0,1,p^s1}`$, $`\mathrm{}A_{p^s1,0,0},A_{p^s1,0,1},\mathrm{},A_{p^s1,0,p^s1}`$, $`\mathrm{},A_{p^s1,p^s1,0},A_{p^s1,p^s1,1},\mathrm{},A_{p^s1,p^s1,p^s1}`$.
3. For $`0ip^s1`$, $`0jp^s1`$, connect the variable node $`B_{0,i,j}`$, that is in the $`0^{th}`$ class of $`L_3`$, to the constraint nodes $`A_{i,j,0},A_{i,j,1},\mathrm{},A_{i,j,p^s1}.`$
4. Let $`x_t(i,k)=M^{(i1)}(k,t)`$, the $`(k,t)^{th}`$ entry of $`M^{(i1)}`$, and let $`y_t(i,j)=M^{(i)}(j,t)`$, the $`(j,t)^{th}`$ entry of $`M^{(i)}`$, where $`i=i\text{ mod }p^s`$.
5. Then, for $`1ip^s`$, $`0j,kp^s1`$, connect the variable node $`B_{i,j,k}`$ to the constraint nodes
$$A_{0,x_0(i,k),y_0(j,k)},A_{1,x_1(i,k),y_1(j,k)},\mathrm{},A_{p^s1,x_{p^s1}(i,k),y_{p^s1}(j,k)}.$$
The Type II, $`\mathrm{}=4`$, LDPC codes have girth eight, minimum distance $`d_{\mathrm{min}}2(p^s+1)`$, and blocklength $`N=1+p^s+p^{2s}+p^{3s}`$. (We believe that the tree bound on the minimum distance is actually met for all the Type II, $`\mathrm{}=4`$, codes, i.e. $`d_{\mathrm{min}}=w_{\mathrm{min}}=2(p^s+1)`$.) Figure 7 illustrates the general construction procedure. For $`d=3`$, the Type II, $`\mathrm{}=4`$, LDPC constraint graph as shown in Figure 8 corresponds to the $`(2,2)`$-Finite-Generalized-Quadrangles-based LDPC (FGQ LDPC) code of ; the squares used for constructing this code are
$$M^{\left(0\right)}=\left[\begin{array}{cc}0& 0\\ 1& 1\end{array}\right],M^{\left(1\right)}=\left[\begin{array}{cc}0& 1\\ 1& 0\end{array}\right].$$
We believe that the Type II, $`\mathrm{}=4`$, construction results in other FGQ LDPC codes for other choices of $`d`$. The Type II construction algorithm can be modified for larger $`\mathrm{}`$ by involving more iterations of the MOLS in the connection scheme, as will be discussed in a forthcoming paper.
## V Simulation Results
Figures 9, 10, 11, 12 show the bit-error-rate performance of Type I-A, Type I-B, Type II girth six, and Type II girth eight LDPC codes, respectively over a binary input additive white Gaussian noise channel with min-sum iterative decoding. The performance of regular or semi-regular randomly constructed LDPC codes of comparable rates and blocklengths are also shown. (All of the random LDPC codes compared in this paper have a variable node degree of three and are constructed from the online LDPC software available at
http://www.cs.toronto.edu/$`\stackrel{~}{}`$ radford/ldpc.software.html.)
Figure 9 shows that the Type I-A LDPC codes perform substantially better than their random counterparts. Figure 10 reveals that the Type I-B LDPC codes perform better than comparable random LDPC codes at short blocklengths; but as the blocklengths increase, the random LDPC codes tend to perform better in the waterfall region. Eventually however, as the SNR increases, the Type I-B LDPC codes outperform the random ones since, unlike the random codes, they do not have a prominent error floor. Figure 11 reveals that the performance of Type II girth-six LDPC codes is also significantly better than comparable random codes; these codes correspond to the two dimensional PG LDPC codes of . Figure 12 indicates the performance of Type II girth-eight LDPC codes; these codes perform comparably to random codes at short blocklengths, but at large blocklengths, the random codes perform better as they have larger relative minimum distances compared to the Type II girth-eight LDPC codes.
As a general observation, min-sum iterative decoding of most of the tree-based LDPC codes (particularly, Type I-A, Type II, and some Type I-B) presented here did not typically reveal detected errors, i.e., errors caused due to the decoder failing to converge to any valid codeword within the maximum specified number of iterations. Detected errors are caused primarily due to the presence of pseudocodewords, especially those of minimal weight. We think that the lack of detected errors with iterative decoding of these LDPC codes is primarily due to their good minimum pseudocodeword weight $`w_{\mathrm{min}}`$.
## VI Conclusions
The Type I construction yields a family of LDPC codes that, to the best of our knowledge, do not correspond to any of the LDPC codes obtained from finite geometries or other geometrical objects. The two tree-based constructions presented in this paper yield a wide range of codes that perform well when decoded iteratively, largely due to the maximized minimal pseudocodeword weight. However, the overall minimum distance of the code is relatively small. Constructing codes with larger minimum distance, while still maintaining $`d_{\mathrm{min}}=w_{\mathrm{min}}`$, remains an open problem. |
warning/0506/hep-ph0506307.html | ar5iv | text | # NEUTRINO MASS AND MIXING PARAMETERS: A SHORT REVIEW
## 1 Introduction
There is compelling experimental evidence $`^\mathrm{?}`$ that the three known neutrino states with definite flavor $`\nu _\alpha `$ ($`\alpha =e,\mu ,\tau `$) are linear combinations of states with definite mass $`\nu _i`$ ($`i=1,2,3`$), and that the Hamiltonian of neutrino propagation in vacuum $`^\mathrm{?}`$ and matter $`^\mathrm{?}`$ does not commute with flavor. The evidence for flavor nonconservation (i.e., “neutrino oscillations”) comes from a series of experiments performed during about four decades of research with very different neutrino beams and detection techniques: the solar neutrino $`^\mathrm{?}`$ experiments Homestake $`^\mathrm{?}`$, Kamiokande $`^\mathrm{?}`$, SAGE $`^\mathrm{?}`$, GALLEX-GNO $`^{\mathrm{?},\mathrm{?}}`$, Super-Kamiokande (SK) $`^\mathrm{?}`$ and the Sudbury Neutrino Observatory (SNO) $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$; the long-baseline reactor neutrino experiment KamLAND $`^{\mathrm{?},\mathrm{?}}`$; the atmospheric neutrino experiments Super-Kamiokande $`^{\mathrm{?},\mathrm{?}}`$, MACRO $`^\mathrm{?}`$, and Soudan-2 $`^\mathrm{?}`$; and the long-baseline accelerator neutrino experiment KEK-to-Kamioka (K2K) $`^{\mathrm{?},\mathrm{?}}`$.
Together with the null results from the CHOOZ $`^\mathrm{?}`$ short-baseline reactor experiment, the above oscillation data provide stringent constraints on the neutrino mixing matrix, on the splittings between squared neutrino masses, and on matter effects. The absolute neutrino masses are being probed by different, non-oscillation searches: beta decay experiments $`^\mathrm{?}`$, neutrinoless double beta decay searches ($`0\nu 2\beta `$) $`^{\mathrm{?},\mathrm{?}}`$, and precision cosmology $`^{\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?}}`$. Current non-oscillation data provide only upper limits on neutrino masses, with the only exception of the Heidelberg-Moscow $`0\nu 2\beta `$ experiment $`^\mathrm{?}`$, whose claimed signal implies a lower bound on neutrino masses.
Basically all the data are consistent with the simplest extension of the standard electroweak model needed to accommodate nonzero neutrino masses and mixings, namely, with a scenario where the three known flavor states $`\nu _{e,\mu ,\tau }`$ are mixed with only three mass states $`\nu _{1,2,3}`$, no other states or new neutrino interactions being needed. This “standard three-neutrino framework” (as recently reviewed, e.g., in $`^\mathrm{?}`$, where further references can be found) appears thus as a new paradigm of particle and astroparticle physics, which will be tested, refined, and possibly challenged by a series of new, more sensitive experiments planned for the next few years or even for the next decades $`^\mathrm{?}`$. The first challenge might actually come very soon from the running MiniBooNE experiment $`^\mathrm{?}`$, which is probing the only piece of data at variance with the standard three-neutrino framework, namely, the controversial result of the Liquid Scintillator Neutrino Experiment (LSND) $`^\mathrm{?}`$.
In this short review we focus on the current status of the standard three-neutrino framework and on the neutrino mass and mixing parameters which characterize it. The parameters are defined as follows. The unitary mixing matrix $`U`$, in terms of one-particle neutrino states $`|\nu `$, is defined as (see, e.g., $`^\mathrm{?}`$),
$$|\nu _\alpha =\underset{i=1}{\overset{3}{}}U_{\alpha i}^{}|\nu _i.$$
(1)
A common parameterization for the matrix $`U`$ is:
$$U=O_{23}\mathrm{\Gamma }_\delta O_{13}\mathrm{\Gamma }_\delta ^{}O_{12},$$
(2)
where the $`O_{ij}`$’s are real Euler rotations with angles $`\theta _{ij}[0,\pi /2]`$, while $`\mathrm{\Gamma }_\delta `$ embeds a CP-violating phase $`\delta [0,2\pi ]`$,
$$\mathrm{\Gamma }_\delta =\mathrm{diag}(1,1,e^{+i\delta }).$$
(3)
By considering $`\mathrm{\Gamma }_\delta O_{13}\mathrm{\Gamma }_\delta ^{}`$ as a single (complex) rotation, this parametrization coincides with the one recommended \[together with Eq. (1)\] in the Review of Particle Properties $`^\mathrm{?}`$.
For the sake of simplicity, the phase $`\delta `$ will not be considered in full generality hereafter. Numerical examples will refer only to the two inequivalent CP-conserving cases, namely, $`e^{i\delta }=\pm 1`$. In these two cases, the mixing matrix takes a real form $`U_{\mathrm{CP}}`$,
$$U_{\mathrm{CP}}=\left(\begin{array}{ccc}c_{13}c_{12}& s_{12}c_{13}& \pm s_{13}\\ s_{12}c_{23}s_{23}s_{13}c_{12}& c_{23}c_{12}s_{23}s_{13}s_{12}& s_{23}c_{13}\\ s_{23}s_{12}s_{13}c_{23}c_{12}& s_{23}c_{12}s_{13}s_{12}c_{23}& c_{23}c_{13}\end{array}\right),$$
(4)
where the upper (lower) sign refers to $`\delta =0`$ ($`\delta =\pi `$). The two cases are formally related by the replacement $`s_{13}s_{13}`$.
The three-neutrino mass spectrum $`\{m_i\}_{i=1,2,3}`$ is formed by a “doublet” of relatively close states and by a third “lone” neutrino state, which may be either heavier than the doublet (“normal hierarchy,” NH) or lighter (“inverted hierarchy,” IH). Here, the lightest (heaviest) neutrino in the doublet is called $`\nu _1`$ ($`\nu _2`$), so that their squared mass difference is
$$\delta m^2=m_2^2m_1^2>0$$
(5)
by convention. The lone state is then labeled as $`\nu _3`$, and the physical sign of $`m_3^2m_{1,2}^2`$ distinguishes NH from IH.
Concerning the second independent squared mass difference $`\mathrm{\Delta }m^2`$, we define it as $`^{\mathrm{?},\mathrm{?}}`$
$$\mathrm{\Delta }m^2=\left|m_3^2\frac{m_1^2+m_2^2}{2}\right|,$$
(6)
so that the two hierarchies (NH and IH) are simply related by the transformation $`+\mathrm{\Delta }m^2\mathrm{\Delta }m^2`$. The largest and next-to-largest squared mass gaps are given $`\mathrm{\Delta }m^2\pm \delta m^2/2`$ in both cases. More precisely, the squared mass matrix $`M^2`$ reads, in such convention,
$$M^2=\mathrm{diag}(m_1^2,m_2^2,m_3^2)=\frac{m_2^2+m_1^2}{2}\mathbf{\hspace{0.17em}1}+\mathrm{diag}(\frac{\delta m^2}{2},+\frac{\delta m^2}{2},\pm \mathrm{\Delta }m^2),$$
(7)
where the upper (lower) sign refers to normal (inverted) hierarchy.
In the previous equation, the term proportional to the unit matrix $`\mathrm{𝟏}`$ is irrelevant in neutrino oscillations, while it matters in observables sensitive to the absolute neutrino mass scale, such as in $`\beta `$-decay and precision cosmology. In particular, $`\beta `$-decay experiments are sensitive to the so-called effective electron neutrino mass $`m_\beta `$,
$$m_\beta =\left[\underset{i}{}|U_{ei}|^2m_i^2\right]^{\frac{1}{2}}=\left[c_{13}^2c_{12}^2m_1^2+c_{13}^2s_{12}^2m_2^2+s_{13}^2m_3^2\right]^{\frac{1}{2}},$$
(8)
as far as the single $`\nu _i`$ mass states are not experimentally resolvable. On the other hand, precision cosmology is sensitive, to a good approximation (up to small hierarchy-dependent effects which may become important in next-generation precision measurements $`^\mathrm{?}`$) to the sum of neutrino masses $`\mathrm{\Sigma }`$ $`^\mathrm{?}`$,
$$\mathrm{\Sigma }=m_1+m_2+m_3.$$
(9)
Finally, if neutrinos are indistinguishable from their antiparticles (i.e., if they are Majorana rather than Dirac neutrinos), the mixing matrix $`U`$ acquires a (diagonal) extra factor
$$UUU_M,$$
(10)
containing Majorana phases $`\varphi _i`$, which are irrelevant in oscillations but not in neutrinoless double beta decay ($`0\nu 2\beta `$). Using the parametrization
$$U_M=\mathrm{diag}(1,\mathrm{e}^{\frac{\mathrm{i}}{2}\varphi _2},\mathrm{e}^{\frac{\mathrm{i}}{2}(\varphi _3+2\delta )}),$$
(11)
the expression of the effective Majorana mass $`m_{\beta \beta }`$ probed in $`0\nu 2\beta `$ experiments $`^\mathrm{?}`$ takes the form:
$$m_{\beta \beta }=\left|\underset{i}{}U_{ei}^2m_i\right|=\left|c_{13}^2c_{12}^2m_1+c_{13}^2s_{12}^2m_2e^{i\varphi _2}+s_{13}^2m_3e^{i\varphi _3}\right|.$$
(12)
Finally, we remark that the constraints on the neutrino oscillation parameters shown hereafter have been obtained by fitting accurate theoretical predictions to a large set of experimental data, through either least-square or maximum-likelihood methods. In both cases, parameter estimations reduce to finding the minimum of a $`\chi ^2`$ function and to tracing iso-$`\mathrm{\Delta }\chi ^2`$ contours around it. We adopt the convention used in $`^\mathrm{?}`$ and call “region allowed at $`n\sigma `$” the subset of the parameter space obeying the inequality
$$\mathrm{\Delta }\chi ^2n^2.$$
(13)
The projection of such allowed region onto each single parameter provides the $`n\sigma `$ bound on such parameter. In particular, we shall also directly use the relation $`\sqrt{\mathrm{\Delta }\chi ^2}=n`$ to derive allowed parameter ranges at $`n`$ standard deviations. Numerical results and figures are taken from the recent review $`^\mathrm{?}`$ (to which the reader is referred for details), which also makes use of results from Refs. $`^{\mathrm{?},\mathrm{?}}`$. For previous reviews on neutrino parameters see, e.g., Refs. $`^{\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?}}`$.
## 2 Constraints on $`(\delta m^2,\mathrm{sin}^2\theta _{12},\mathrm{sin}^2\theta _{13})`$ from solar+KamLAND data
It is well known that the angle $`\theta _{13}`$ is relatively small and possibly zero. For $`\theta _{13}=0`$, both solar and (long-baseline) reactor neutrino oscillations depend solely on the parameters $`(\delta m^2,\theta _{12})`$. Figure 1 shows the current constraints on such parameters from a global analysis of all the available solar neutrino data $`^\mathrm{?}`$ and of KamLAND data $`^\mathrm{?}`$, both separately and in combination. Although (at $`3\sigma )`$ multiple solutions can explain KamLAND data, the combination with solar data provides a well-defined and unique solution at large mixing angle (LMA) in the mass-mixing parameter space. The identification of such solution represents one of the most impressive recent advances in neutrino physics.
Further progress can be expected in narrowing the parameter space in Fig. 1. The $`\delta m^2`$ uncertainty is currently dominated by the KamLAND observation of half-period of oscillations $`^\mathrm{?}`$ and can be improved with higher statistics $`^\mathrm{?}`$. The $`\mathrm{sin}^2\theta _{12}`$ uncertainty is instead dominated by the SNO ratio of charged-to-neutral current (CC/NC) event rates, which can also be improved with future data $`^\mathrm{?}`$.
The current solar LMA solution, as compared with results prior to the complete SNO-II data set $`^\mathrm{?}`$, is slightly shifted toward larger values of $`\mathrm{sin}^2\theta _{12}`$ and allows higher values of $`\delta m^2`$. \[Our current best-fit point for solar data only is at $`\delta m^2=6.3\times 10^5\mathrm{eV}^2`$ and $`\mathrm{sin}^2\theta _{12}=0.314.`$\] This trend is substantially due to the larger value of the CC/NC ratio measured in the complete SNO II phase (0.34 $`^{\mathrm{?},\mathrm{?}}`$) with respect to the previous central value (0.31 $`^\mathrm{?}`$). We also find that the SNO-II charged-current spectral data $`^\mathrm{?}`$ contribute to allow slightly higher values of $`\delta m^2`$ with respect to older results.
For $`\theta _{13}>0`$, the solar and KamLAND $`\nu `$ parameter space is spanned by $`(\delta m^2,\mathrm{sin}^2\theta _{12},\mathrm{sin}^2\theta _{13})`$. Figure 2 shows the current $`2\sigma `$ bounds in such space (both separately and in combination) in each of the three coordinate planes. Remarkably, both solar and KamLAND data are consistent with $`\theta _{13}`$ being small ($`0`$ at best-fit), in agreement with independent atmospheric, accelerator and short-baseline reactor data (see the next section). The combined upper bound on $`\mathrm{sin}^2\theta _{13}`$ in Fig. 2 is at the interesting level of $`5\%`$.
## 3 Constraints on $`(\mathrm{\Delta }m^2,\mathrm{sin}^2\theta _{23},\mathrm{sin}^2\theta _{13})`$ from SK<sub>ATM</sub>+K2K+CHOOZ data
In the limit $`\delta m^2/\mathrm{\Delta }m^21`$ (one-dominant-mass-scale approximation), the leading parameters in atmospheric and long-baseline accelerator searches are $`(\mathrm{\Delta }m^2,\mathrm{sin}^2\theta _{23},\mathrm{sin}^2\theta _{13})`$. Subleading effects induced by $`\delta m^20`$ (i.e., LMA effects in terrestrial neutrino oscillations $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$) are present, however, even for $`\theta _{13}0`$ $`^\mathrm{?}`$. In accurate calculations, it is worthwhile to include such effects numerically, e.g., by fixing ($`\delta m^2,\mathrm{sin}^2\theta _{12}`$) at their best-fit value in Fig. 1 in a full three-flavor analysis of atmospheric and K2K data, as done in the following two figures.
Figure 3 shows, for $`\theta _{13}=0`$, the results of our analysis of SK (atmospheric) and K2K data, both separately and in combination. The K2K constraints are octant-symmetric and relatively weak in $`\mathrm{sin}^2\theta _{23}`$, while they contribute appreciably to reduce the overall $`\mathrm{\Delta }m^2`$ uncertainty. The SK atmospheric neutrino contraints are instead strong on both mass and mixing parameters, and also slightly asymmetrical $`^\mathrm{?}`$ in $`\mathrm{sin}^2\theta _{23}`$. Unfortunately, current data are not accurate enough to promote this slight asymmetry to a real $`\theta _{23}`$-octant discrimination. However, it is not excluded that future, high-statistics atmospheric neutrino data might be able to do so, if $`\theta _{23}`$ is not too close to $`\pi /4`$ $`^\mathrm{?}`$. Such possible $`\theta _{23}`$ octant asymmetry, together with a measurement of $`\theta _{13}`$, is crucial for model building $`^{\mathrm{?},\mathrm{?}}`$.
For $`\theta _{13}>0`$, SK+K2K data are also sensitive, in principle, to the neutrino mass hierarchy \[$`\mathrm{sign}(\pm \mathrm{\Delta }m^2)=\pm 1`$\] and to the CP parity \[$`\mathrm{cos}\delta =\pm 1`$\]. However, the dependence is very small within the CHOOZ bounds on $`\theta _{13}`$ (see, e.g., Ref. $`^\mathrm{?}`$ and references therein), and thus it makes sense to marginalize the SK+K2K+CHOOZ $`\chi ^2`$ function with respect to hierarchy and CP parity. The results are shown in Fig. 4, in terms of the projections of the $`(\mathrm{\Delta }m^2,\mathrm{sin}^2\theta _{23},\mathrm{sin}^2\theta _{13})`$ region allowed at 1, 2, and $`3\sigma `$ onto each of the coordinate planes (with LMA effects included). The best fit is reached for nonzero $`\theta _{13}`$ (mainly due to a slight preference of low-energy atmospheric data for $`\nu _e`$ event appearance), but $`\theta _{13}=0`$ is allowed within less than $`1\sigma `$. The preferred value of $`\mathrm{sin}^2\theta _{23}`$ remains slightly below maximal mixing. The best-fit value of $`\mathrm{\Delta }m^2`$ is $`2.4\times 10^3`$ eV<sup>2</sup>. Notice that the correlations among the three parameters in Fig. 4 are very weak.
## 4 Global constraints on oscillation parameters
The results of the global analysis of solar+KamLAND data (Sec. 2) and of SK+K2K+CHOOZ data (Sec. 3) can now be merged to provide our best estimates of the five neutrino oscillation parameters $`(\delta m^2,\mathrm{\Delta }m^2,\theta _{12},\theta _{13},\theta _{23})`$, marginalized over the $`2\times 2`$ cases with different mass hierarchies and CP parities (which are physically different but phenomenologically indistinguishable at present). The bounds will be directly shown in terms of the “number of sigmas”, corresponding to the function $`(\mathrm{\Delta }\chi ^2)^{1/2}`$ for each parameter.
Figure 5 shows our global bounds on $`\mathrm{sin}^2\theta _{13}`$, as coming from all data (solid line) and from the following partial data sets: KamLAND (dotted), solar (dot-dashed), solar+KamLAND (short-dashed) and SK+K2K+CHOOZ (long-dashed). Only the latter set, as observed before, gives a weak indication for nonzero $`\theta _{13}`$. Interestingly, solar+KamLAND data are now sufficiently accurate to provide bounds which are not much weaker than the dominant SK+K2K+CHOOZ ones, also because the latter slightly prefer $`\theta _{13}>0`$ as best fit, while the former do not.
Figure 6 shows our global bounds on the four mass-mixing parameters which present both upper and lower limits with high statistical significance. Notice that the accuracy of the parameter estimate is already good enough to lead to almost “linear” errors, especially for $`\delta m^2`$ and $`\mathrm{sin}^2\theta _{12}`$. For $`\mathrm{\Delta }m^2`$ and $`\mathrm{sin}^2\theta _{23}`$, such “linearity” is somewhat worse in the region close to the best fit (say, within $`\pm 1\sigma `$), and thus $`2\sigma `$ (or $`3\sigma `$) errors should be taken as reference.
## 5 Non-oscillation data and their interplay with oscillation constraints
Since oscillation data fix the mass splittings $`\delta m^2`$ and $`\mathrm{\Delta }m^2`$, the observables sensitive to absolute neutrino masses $`(m_\beta ,m_{\beta \beta },\mathrm{\Sigma })`$ are higly correlated with each other, both in normal and in inverted hierarchy: typically, when one increases the other also increases. Therefore, upper bounds on any of them translate into upper bounds on the others. However, the upper bounds on $`m_\beta `$ are currently weak (a few eV) $`^\mathrm{?}`$, and the relevant discussion can be limited to $`(m_{\beta \beta },\mathrm{\Sigma })`$ at present.
Figure 7 shows the impact of all the available non-oscillation data, taken at face value, in the parameter space $`(m_{\beta \beta },\mathrm{\Sigma })`$, at the $`2\sigma `$ level. The horizontal band is allowed by the positive $`0\nu 2\beta `$ experimental claim $`^\mathrm{?}`$ equipped with the nuclear uncertainties of $`^\mathrm{?}`$ as described in $`^\mathrm{?}`$. The slanted bands (for normal and inverted hierarchy) are allowed by all other neutrino data, i.e., by the combination of neutrino oscillation constraints (from Figs. 5 and 6) and of astrophysical and cosmological constraints from Cosmic Microwave Background (CMB) $`^\mathrm{?}`$, large scale structures from galaxy surveys (2dF) $`^\mathrm{?}`$, and small scale structures from Lyman $`\alpha `$ forest data $`^\mathrm{?}`$, as described in $`^\mathrm{?}`$. The tight cosmological upper bound on $`\mathrm{\Sigma }`$ prevents the overlap between the slanted and horizontal bands at $`2\sigma `$, indicating that no global combination of oscillation and non-oscillation data is possible in the sub-eV range. The “discrepancy” is now even stronger than it was found in Ref. $`^\mathrm{?}`$, due to the adoption of smaller $`0\nu 2\beta `$ nuclear uncertainties $`^\mathrm{?}`$. It is premature, however, to derive any definite conclusion as to which piece of the data (or of the $`3\nu `$ scenario) is “wrong” in this conflicting picture. E.g., by relaxing either the $`0\nu 2\beta `$ lower bound or the Ly$`\alpha `$ data, global combinations are possible $`^\mathrm{?}`$. Further experimental and theoretical research is needed to clarify the interplay of absolute neutrino observables in the sub-eV range.
## 6 Summary and Conclusions
There is compelling evidence for neutrino flavor change driven by nonzero masses and mixing angles. Basically all oscillation data (with the only exception of LSND) are consistent within a three-neutrino framework. Within such framework, the global constraints from oscillation data can be summarized (see also Figs. 5 and 6 and Ref. $`^\mathrm{?}`$) through the following $`\pm 2\sigma `$ ranges (95% C.L.) for each parameter:
$`\mathrm{sin}^2\theta _{13}`$ $`=`$ $`0.9_{0.9}^{+2.3}\times 10^2,`$ (14)
$`\delta m^2`$ $`=`$ $`7.92(1\pm 0.09)\times 10^5\mathrm{eV}^2,`$ (15)
$`\mathrm{sin}^2\theta _{12}`$ $`=`$ $`0.314(1_{0.15}^{+0.18}),`$ (16)
$`\mathrm{\Delta }m^2`$ $`=`$ $`2.4(1_{0.26}^{+0.21})\times 10^3\mathrm{eV}^2,`$ (17)
$`\mathrm{sin}^2\theta _{23}`$ $`=`$ $`0.44(1_{0.22}^{+0.41}).`$ (18)
Such ranges are marginalized over the four inequivalent cases $`[\mathrm{sign}(\pm \mathrm{\Delta }m^2)][\mathrm{cos}\delta =\pm 1]`$, i.e., over the two possible hierarchies and the two possible CP-conserving cases, which are currently undistinguishable. \[Notice that the lower error on $`\mathrm{sin}^2\theta _{13}`$ is purely formal, and corresponds to the positivity constraints $`\mathrm{sin}^2\theta _{13}>0`$.\]
Concerning the observables sensitive to absolute masses ($`m_\beta `$, $`m_{\beta \beta }`$ and $`\mathrm{\Sigma }`$), the situation is still unclear. Current constraints at the eV/sub-eV level are dominated by either upper bounds on $`\mathrm{\Sigma }`$ from cosmology or by the $`0\nu 2\beta `$ claim on $`m_{\beta \beta }`$, whose combination is not possible, however, at face value. Further studies and data are need to go beyond the general statement that neutrino masses should be smaller than $`1`$ eV, and to really explore the sub-eV range.
Within the three-neutrino scenario, it appears that the most important unsolved problems require probing $`\theta _{13}`$, $`\delta `$, the hierarchy, and the absolute neutrino masses. Needless to say, further experimental results or theoretical insights might also reserve big surprises and force us to go beyond such scenario, either by adding new neutrino states, or new interactions, or both.
## Acknowledgments
This work is supported by the Italian Ministero dell’Istruzione, Università e Ricerca (MIUR) and Istituto Nazionale di Fisica Nucleare (INFN) through the “Astroparticle Physics” project.
## References |
warning/0506/math0506162.html | ar5iv | text | # Filters and subgroups associated with Hartman measurable functions
## 1 Introduction
### 1.1 Motivation:
In the investigation of finitely additive measures in number theoretic context led to the concept of a Hartman measurable subset $`HG`$ of a discrete abelian group $`G`$. By definition, $`H`$ is Hartman measurable if it is the preimage $`H=\iota _X^1(M)`$ of a continuity set $`MX`$ in a group compactification $`(\iota _X,X)`$ of $`G`$. This, more explicitly, means that $`\iota _X:GC`$ is a group homomorphism with $`\iota (G)`$ dense in the compact group $`X`$ and that $`\mu _X(M)=0`$. Here $`M`$ denotes the topological boundary of $`M`$ and $`\mu _X`$ the normalized Haar measure on $`X`$. By putting $`m_G(H)=\mu _X(M)`$ one can define a finitely additive measure on the Boolean set algebra of all Hartman measurable sets in $`G`$. For the special case $`G=`$ a Hartman set $`H`$, by identification with its characteristic function, can be considered to be a two-sided infinite binary sequence, called a Hartman sequence. Certain number theoretic, ergodic and combinatorial aspects of Hartman sequences have been studied in and , while presents a method to reconstruct the group compactification $`(\iota _X,X)`$ for given $`H`$.
In order to benefit from powerful tools from functional and harmonic analysis it is desirable to study appropriate generalizations of Hartman measurable sets by replacing their characteristic functions by complex valued functions not only taking the values 0 and 1. The natural definition of a Hartman measurable function $`\phi :G`$ is the requirement $`\phi ^{}\iota _X`$, where $`(\iota _X,X)`$ is a group compactification of $`G`$ and $`\phi ^{}`$ is integrable in the Riemann sense, i.e. its points of discontinuity form a null set with respect to the Haar measure on $`X`$. This definition is equivalent to the one of R-almost periodicity, introduced in by S. Hartman. The investigation of the space $`(G)`$ of all Hartman measurable functions on $`G`$ is the content of . Here we are going to transfer ideas from into this context. Thus our main topic is to describe $`(\iota _X,X)`$ only in terms of $`\phi `$. In particular we establish further connections to Fourier analysis. The natural framework for our investigation is that of LCA (locally compact abelian) groups.
### 1.2 Content of the paper
After the introduction we collect in section 2 the necessary preliminary definitions and facts about Hartman measurable sets and functions.
Section 3 treats the following situation: Given a Hartman measurable function $`\phi :GX`$ on an LCA group $`G`$, we know by the very definition of Hartman measurability that there is some group compactification $`(\iota _X,X)`$ of $`G`$ such that $`\phi =\phi ^{}\iota _X`$ for some Riemann integrable function $`\phi ^{}:X`$. We say that $`\phi `$ can be realized in $`(\iota _X,X)`$ resp. by $`\phi ^{}`$. It is easy to see that in this case $`\phi `$ can be realized as well on any ”bigger” compactification $`(\iota _{\stackrel{~}{X}},\stackrel{~}{X})`$. The notion of ”bigger” and ”smaller” is made more precise in the next section.
In particular every Hartman measurable function can be realized in the maximal group compactification of $`G`$, the Bohr compactification $`(\iota _b,bG)`$. The question arises if there is a realization of $`\phi `$ in a group compactification that is as ”small” as possible. If a Hartman measurable function $`\phi `$ possesses a so called aperiodic realization then the group compactification on which this aperiodic realization can be obtained is minimal (Theorem 1). This approach works for arbitrary Hartman measurable $`\phi `$ if one allows ”almost realizations”, i.e. if one demands $`\phi =\phi ^{}\iota _X`$ almost everywhere with respect to the finitely additive Hartman measure $`m_G`$ on $`G`$ rather than $`\phi =\phi ^{}\iota _X`$ everywhere on $`G`$ (Theorem 2). Whenever $`\phi `$ is even almost periodic one can guarantee $`\phi =\phi ^{}\iota _X`$ everywhere on $`G`$. The group compactification on which the minimal realization of $`\phi `$ occurs is unique up to equivalence of group compactifications. It can be obtained by a method involving filters on $`G`$ similar to that presented in .
The content of section 4 is motivated by the following reasoning: Every group compactification of the LCA group $`G`$ corresponds to a (discrete) subgroup $`\mathrm{\Gamma }`$ of the dual $`\widehat{G}`$ in such a way that it is equivalent to the group compactification $`(\iota _\mathrm{\Gamma },C_\mathrm{\Gamma })`$ defined by $`\iota _\mathrm{\Gamma }:g(\chi (g))_{\chi \mathrm{\Gamma }}`$, $`C_\mathrm{\Gamma }:=\overline{\iota _\mathrm{\Gamma }(G)}𝕋^\mathrm{\Gamma }`$. If $`(\iota _X,X)`$ is a group compactification admitting an aperiodic almost realization of the Hartman measurable function $`\phi `$, the corresponding subgroup $`\mathrm{\Gamma }\widehat{G}`$ contains all characters $`\chi `$ such that the corresponding Fourier coefficient $`m_G(\phi \overline{\chi })`$ does not vanish. If $`\phi `$ is almost periodic or if $`\phi `$ can be realized on a finite dimensional compactification this result is sharp in the sense that the subgroup $`\mathrm{\Gamma }`$ is minimal with the above property (Theorem 3). For general Hartman measurable functions the situation is more difficult. This is discussed and illustrated by an example.
Section 5 summarizes the main results and includes an illustrating diagram.
## 2 Preliminaries and Notation
Throughout this paper $`G`$ denotes always an LCA (locally compact abelian) group. For group compactifications of $`G`$ let us write $`(\iota _{X_1},X_1)(\iota _{X_2},X_2)`$ iff there is a continuous group homomorphism $`\pi :X_2X_1`$ such that the diagram
commutes. In this situation we say that $`(\iota _{X_1},X_1)`$ is covered by $`(\iota _{X_2},X_2)`$. If $`(\iota _{X_1},X_1)`$ is covered by $`(\iota _{X_2},X_2)`$ (via $`\pi _1`$) and $`(\iota _{X_2},X_2)`$ is covered by $`(\iota _{X_1},X_1)`$ (via $`\pi _2`$) then $`(\iota _{X_1},X_1)`$ and $`(\iota _{X_2},X_2)`$ are called equivalent. In this case compactness of $`X_1`$ and $`X_2`$ implies that $`\pi _1`$ and $`\pi _2`$ are both topological and algebraic isomorphisms between $`X_1`$ and $`X_2`$. $`\mathrm{"}\mathrm{"}`$ is a partial order on the class of group compactifications modulo equivalence. The maximal element with respect to this order is $`(\iota _b,bG)`$, the Bohr compactification of the topological group $`G`$. Recall that $`AP(G)`$, the set of almost periodic functions on $`G`$, is isometrically isomorphic to $`C(bG)`$, the set of continuous functions on $`bG`$. The mapping $`\iota _b^{}:C(bG)AP(G)`$, defined via $`ff\iota _b`$, is an isometry. Note that this is just a different way to characterize those continuous functions on $`G`$, which can be extended to continuous functions on $`bG`$. This definition (which is best suited for our purposes) is equivalent to the notion of almost periodicity established by Bohr resp. Bochner.
For a locally compact abelian (LCA) group $`G`$ let us denote by $`\widehat{G}`$ the set of all continuous homomorphisms $`\chi :G𝕋/`$. We will occasionally identify $`𝕋`$ with the unimodular group $`\{z:|z|=1\}`$. This will cause no confusion.
$`\widehat{G}`$ endowed with the compact-open topology is an LCA group in its own right. $`\widehat{G}`$ is the Pontryagin dual of $`G`$. Every subset $`\mathrm{\Gamma }\widehat{G}`$ induces a group compactification $`(\iota _\mathrm{\Gamma },C_\mathrm{\Gamma })`$ of $`G`$ via $`\iota _\mathrm{\Gamma }(g):=(\chi (g))_{\chi \mathrm{\Gamma }}𝕋^\mathrm{\Gamma }`$ and $`C_\mathrm{\Gamma }:=\overline{\iota _\mathrm{\Gamma }(G)}𝕋^\mathrm{\Gamma }`$
One can show that every group compactification $`(\iota _X,X)`$ of an LCA group $`G`$ is equivalent to the group compactification induced by the subgroup $`\{\iota _X\eta :\eta \widehat{X}\}\widehat{G}`$ (Theorem 26.13 in ). Thus group compactifications of LCA groups can be described by subgroups of the dual and vice versa.
The system $`\mathrm{\Sigma }(G)𝔓(G)`$ of all Hartman measurable sets on $`G`$, i.e. the system of all preimages $`\iota _b^1(M)`$ of $`\mu _b`$-continuity sets in the Bohr compactification $`(\iota _b,bG)`$ of $`G`$, is a Boolean set algebra and enjoys the property that there exists a unique translation invariant finitely additive probability measure $`m_G`$ on $`\mathrm{\Sigma }(G)`$: $`m_G(\iota _X^1(M^{}))=\mu _X(M^{})`$. For details we refer to .
Let us denote by $`\mathrm{\Delta }`$ the symmetric difference of sets and by $`\tau _g`$ the translation operator on an abelian group defined by $`\tau _g(h):=g+h`$. We introduce two mappings:
* for a Hartman measurable set $`M`$ denote by $`d_M:G[0,1]`$ the mapping $`gm_G(M\mathrm{\Delta }\tau _gM)`$,
* for a $`\mu _X`$-continuity set $`M^{}`$ on some group compactification $`(\iota _X,X)`$ denote by $`d_M^{}:X[0,1]`$ the mapping $`g\mu _X(M^{}\mathrm{\Delta }\tau _gM^{})`$.
Note that the mapping $`d_M^{}`$ (and similarly the mapping $`d_M`$) can be used to define a translation invariant pseudometric by letting $`\rho _M^{}(g,h):=d_M^{}(gh)`$. The set of zeros $`\{g:d_M^{}(g)=0\}`$ is always a closed subgroup. We will denote this subgroup by $`\mathrm{ker}d_M^{}`$.
Now consider sets of the form $`F(M,\epsilon ):=\{gG:d_M(g)<\epsilon \}`$ and denote by $`(M)`$ the filter on $`G`$ generated by $`\{F(M,\epsilon ):\epsilon >0\}`$, i.e. the set of all $`FG`$ such that there exists an $`\epsilon >0`$ with $`F(M,\epsilon )F`$. When we have a realization $`M^{}`$ of $`M`$ on some group compactification $`(\iota _X,X)`$ we can transfer the topological data encoded in the neighborhood filter of the unit $`0_X`$ in $`X`$ to $`G`$ by considering the pullback induced by $`\iota _X`$.
To be precise: Let $`(\iota _X,X)`$ be a group compactification and $`𝔘(X,0_X)`$ the filter of all neighborhoods of the unit $`0_X`$ in $`X`$. By $`𝔘_{(\iota _X,X)}`$ we denote the filter on $`G`$ generated by $`\iota _X^1\left(𝔘(X,0_X)\right)`$. Note that if the mapping $`\iota _X`$ is one-one, $`\iota _X^1\left(𝔘(X,0_X)\right)`$ is already a filter.
For $``$, the group of integers, Theorem 2 in states that for any Hartman set $`M`$ there is a group compactification $`(\iota _X,X)`$ such that $`(M)`$ and $`𝔘(X,0_X)`$ coincide. Hence the filter $`(M)`$ on $``$ contains much information about the group compactification $`(\iota _X,X)`$: If $`M`$ is a Hartman measurable set and $`(\iota _X,X)`$ is a group compactification of the integers such that $`M`$ can be realized on $`X`$ via the continuity set $`M^{}`$ then $`H=\mathrm{ker}d_M^{}`$ is a closed subgroup of $`X`$ and $`(M)=𝔘_{(\pi _H\iota _X,X/H)}`$, for $`\pi _H:XX/H`$ the canonical quotient mapping.
In what follows we need to generalize this result to arbitrary (LCA) groups. This poses no problem since the proof given in for $`G=`$ applies verbatim to an arbitrary topological group.
Recall that for a filter $`𝔓(X)`$ on some set $`X`$ and a function $`\phi :X`$ the filter-limit $`lim_{xX}\phi (x)`$ is defined to be the unique $`\lambda `$ such that $`\epsilon >0`$ we have $`\{xX:|\phi (x)\lambda |<\epsilon \}`$. In the filter $`=(M)`$ is also used to define the subgroup Sub$`(M)`$ of $`𝕋`$ consisting of all those elements $`\alpha `$ such that $`lim_nn\alpha =0`$ (or equivalently: $`lim_ne^{2\pi in\alpha }=1`$).
All three objects - filter, compactification and subgroup - carry the same information regarding a fixed Hartman set $`M`$. It is interesting to note that any subgroup of a compact abelian group $`G`$ can be written as $`\{gG:lim_{\chi \widehat{G}}\chi (g)=1\}`$ for some filter $``$ on $`\widehat{G}`$ (cf. ).
We transfer these concepts into our more general context. To that cause we need the following definitions. Recall that a bounded function $`f`$ on a group compactification $`(\iota _X,X)`$ is called Riemann integrable iff the set $`\text{disc}(f)`$ of points of discontinuity is a $`\mu _X`$-null set, for $`\mu _X`$ the normalized Haar measure on $`X`$. Let us denote the set of all such functions by $`R_{\mu _X}(X)`$ or, simply, $`R(X)`$. We use the following characterization, a proof of which can be found in .
###### Proposition 1.
Let $`X`$ be a compact space and $`\mu _X`$ a finite positive Borel measure on $`X`$. For a bounded real-valued $`\mu _X`$-measurable function $`f`$ the following assertions are equivalent:
1. $`f`$ is Riemann integrable.
2. For every $`\epsilon >0`$ there exist continuous functions $`g_\epsilon `$ and $`h_\epsilon `$ such that $`g_\epsilon fh_\epsilon `$ and $`_X(h_\epsilon g_\epsilon )𝑑\mu _X<\epsilon `$.
Let $`\phi `$ be a function defined on a topological group $`G`$ and $`(\iota _b,bG)`$ the Bohr compactification of $`G`$. We call a function $`\phi _b`$ defined on $`bG`$ an extension (or realization) of $`\phi `$ iff $`\phi =\phi _b\iota _b`$. For example: The set of almost periodic functions on $`G`$ coincides with the set of those functions that can be extended to continuous functions on the Bohr compactification.
###### Definition 2.
Let $`(\iota _b,bG)`$ be the Bohr compactification of the topological group $`G`$. We call a bounded function $`\phi `$ on $`G`$ Hartman measurable iff $`\phi `$ can be extended to a Riemann integrable function $`\phi _b`$ on $`bG`$. The set of Hartman measurable functions $`\{\phi ^{}\iota _b:\phi ^{}R_{\mu _b}(bG)\}`$ is denoted by $`(G)`$.
Given a Hartman measurable function $`\phi `$, we say that *$`\phi ^{}`$ realizes $`\phi `$* if $`\phi ^{}`$ is a Riemann integrable function defined on some group compactification $`(\iota _X,X)`$ such that $`\phi =\phi ^{}\iota _X`$, cf. the diagram below:
In this situation we also say that $`\phi `$ can be realized on $`(\iota _X,X)`$. Most of this paper is devoted to the task of finding a minimal group compactification on which a given $`\phi (G)`$ can be realized. Note that $`\phi ^{}R(X)`$ implies $`\phi _b=\phi ^{}\pi R(bG)`$ (cf. ).
## 3 Filters associated with Hartman measurable functions
By definition every $`\phi (G)`$ has a realization on $`bG`$ by a Riemann integrable function $`\phi ^{}R_{\mu _b}(bG)`$. The mapping
$$d_\phi ^{}:x\phi ^{}\tau _x\phi ^{}_1:=_{bG}|\phi ^{}\tau _x\phi ^{}|𝑑\mu $$
is continuous (cf. , Corollary 2.32). This implies that $`d_\phi :=d_\phi ^{}\iota _b`$ is an almost periodic function on $`G`$.
The finitely additive invariant measure $`m_G`$ can be extended to an invariant mean on $`(G)`$, i.e. to an invariant and non-negative normalized linear functional on $`(G)`$. It will cause no confusion if we denote this invariant mean again by $`m_G`$ (cf. ). Thus we can also write $`d_\phi (g)=m_G(|\phi \tau _g\phi |)`$. It is then obvious to define $`F(\phi ,\epsilon ):=\{gG:d_\phi (\tau _g\phi )<\epsilon \}`$ and denote by $`(\phi )`$ the filter on $`G`$ generated by $`\{F(\phi ,\epsilon ):\epsilon >0\}`$.
In the LCA setting, we can apply the tools developed in to conclude a functional analogue of Theorem 2 in .
###### Definition 3.
Let $`\phi (G)`$ be realized by $`\phi ^{}`$ on the group compactification $`(\iota _X,X)`$. $`\phi ^{}`$ is called an aperiodic realization iff $`\mathrm{ker}d_\phi ^{}:=\{xX:\phi ^{}\tau _x\phi ^{}_1=0\}=\{0_X\}`$.
###### Theorem 1.
Let $`\phi (G)`$ be realized by $`\phi ^{}`$ on the group compactification $`(\iota _X,X)`$. Then $`(\phi )𝔘_{(\iota _X,X)}`$. Furthermore $`(\phi )=𝔘_{(\iota _X,X)}`$ if $`\phi ^{}`$ is an aperiodic realization.
###### Proof.
Suppose $`\phi =\phi ^{}\iota _X`$ with $`\phi ^{}R_{\mu _X}(X)`$ for a group compactification $`(\iota _X,X)`$. For any set $`A(\phi )`$ there exists $`\epsilon >0`$ such that $`d_\phi (x)<\epsilon `$ implies $`xA`$. Using almost periodicity of $`d_\phi `$, i.e. continuity of $`d_\phi ^{}`$, we find a neighborhood $`U𝔘(X,0_X)`$ such that $`d_\phi ^{}(U)[0,\epsilon )`$. For every $`x\iota _X^1(U)`$ we have $`d_\phi (x)<\epsilon `$. Consequently $`\iota _X^1(U)A𝔘_{(\iota _X,X)}`$ and hence $`(\phi )𝔘_{(\iota _X,X)}`$.
Suppose that $`\phi ^{}R_{\mu _X}(X)`$ is aperiodic, i.e. $`d_\phi ^{}(x)=0`$ iff $`x=0_X`$, the unit in $`X`$. Let $`A𝔘_{(\iota _X,X)}`$ be arbitrary; w.l.o.g. we can assume $`A\iota _X^1(U)`$ for an open neighborhood $`U𝔘(X,0_X)`$. Due to the continuity of $`d_\phi ^{}`$ and compactness of $`X`$ we have $`d_\phi ^{}(x)\epsilon >0`$ for $`xXU^{}`$. This implies $`\iota _X(\{gG:d_\phi (g)<\epsilon \})U`$ and hence $`\{gG:d_\phi (g)<\epsilon \}\iota _X^1(U)A(\phi )`$. Thus $`𝔘_{(\iota _X,X)}(\phi )`$ and consequently $`𝔘_{(\iota _X,X)}=(\phi )`$. ∎
###### Definition 4.
Let $`\phi (G)`$ and let $`(\iota _X,X)`$ be a group compactification of $`G`$. A function $`\psi ^{}R_{\mu _X}(X)`$ is called an almost realization of $`\phi `$ iff $`m_G(|\phi \psi |)=0`$ for $`\psi :=\psi ^{}\iota _X`$ and $`m_G`$ the unique invariant mean on $`(G)`$.
###### Theorem 2.
Every $`\phi (G)`$ has an aperiodic almost realization on some group compactification $`(\iota _X,X)`$. If $`\phi ^{}:X`$ is any aperiodic almost realization of $`\phi `$ then $`(\phi )=𝔘_{(\iota _X,X)}`$.
###### Proof.
We only have to prove that an aperiodic almost realization exists, the rest follows from Theorem 1. Let $`\phi ^{}`$ be a realization of $`\phi `$ on $`X`$. The reader will easily check that $`H:=\mathrm{ker}d_\phi ^{}=\{xX:d_\phi ^{}(x)=0\}`$ is a closed subgroup of the compact abelian group $`X`$.
Weil’s formula for continuous functions on quotients (Theorem 3.22 in ) states that there exists $`\alpha >0`$ such that for every $`fC(X)`$
$`{\displaystyle _{X/H}}\left(\underset{={}_{}{}^{\mathrm{}}f(s)}{\underset{}{{\displaystyle _H}f(s+t)𝑑\mu _H(t)}}\right)𝑑\mu _{X/H}(s)`$ $`=`$ $`\alpha {\displaystyle _X}f(u)𝑑\mu _X(u)`$ (1)
holds. This implies that the canonical mapping $`{}_{}{}^{\mathrm{}}:C(X)C(X/H)`$, $`f{}_{}{}^{\mathrm{}}f`$ defined by $`{}_{}{}^{\mathrm{}}f(s+H)=_Hf(s+t)𝑑\mu _H(t)`$ satisfies $`{}_{}{}^{\mathrm{}}f_1\alpha f_1`$. We rescale the Haar measure on $`H`$ such that $`\alpha =1`$. Thus we can extend $`^{\mathrm{}}`$ to a continuous linear operator $`L^1(X)L^1(X/H)`$. Furthermore positivity of $`^{\mathrm{}}`$ enables us to extend $`^{\mathrm{}}`$ to a mapping defined on $`R_{\mu _X}(X)`$ in the following way:
According to Proposition 1 $`fR_{\mu _X}(X)`$ implies that there are $`g_n,h_nC(X)`$ such that $`g_nfh_n`$ and $`h_ng_n_10`$ as $`n\mathrm{}`$. Thus every function $`\stackrel{~}{f}`$ on $`X/H`$ satisfying
$$f_{}:=\underset{n0}{sup}{}_{}{}^{\mathrm{}}g_{n}^{}\stackrel{~}{f}\underset{n0}{inf}{}_{}{}^{\mathrm{}}h_{n}^{}=:f^{}$$
is in $`R_{\mu _{X/H}}(X/H)`$. Note that $`f_{}`$ and $`f^{}`$ are $`\mu _H`$-measurable and coincide $`\mu _H`$-a.e.; to define $`{}_{}{}^{\mathrm{}}f`$ we pick any function $`\stackrel{~}{f}`$ satisfying $`f_{}\stackrel{~}{f}f^{}`$. Then Weil’s formula (1) will still be valid, regardless of the particular choice of the $`g_n,h_n`$ and $`{}_{}{}^{\mathrm{}}f`$.
Since $`\phi ^{}`$ is Riemann integrable on $`X`$, there exist functions $`\phi _nC(X)`$ such that $`\phi ^{}\phi _n_10`$. Note that $`d_{\phi _n}d_\phi ^{}`$ even uniformly on $`X`$:
$$|d_{\phi _n}(s)d_\phi ^{}(s)|=\left|\tau _s\phi _n\phi _n_1\tau _s\phi ^{}\phi ^{}_1\right|2\phi _n\phi ^{}_10.$$
Using the continuity of $`^{\mathrm{}}`$ as a mapping on $`L^1(X)`$ the same argument also shows that $`|d_{}_{}{}^{\mathrm{}}\phi _{}^{}(s+H)d_{{}_{}{}^{\mathrm{}}\phi _{n}^{}}(s+H)|2\phi _n\phi ^{}_10`$ uniformly on $`X/H`$. Suppose $`d_{}_{}{}^{\mathrm{}}\phi ^{}(s+H)=0`$. Then
$$d_\phi ^{}(s)=\underset{n\mathrm{}}{lim}d_{\phi _n}(s)=\underset{n\mathrm{}}{lim}d_{{}_{}{}^{\mathrm{}}\phi _{n}^{}}(s+H)=0$$
implies $`sH`$, i.e. $`s+H=0_X+HX/H`$. So $`{}_{}{}^{\mathrm{}}\phi _{}^{}`$ is aperiodic.
We show that $`\phi ^{}`$ being a realization of $`\phi `$ implies that $`{}_{}{}^{\mathrm{}}\phi _{}^{}`$ is an almost realization of $`\phi `$. By definition $`tH`$ iff $`A_t:=\{sX:\phi ^{}(s+t)=\phi ^{}(s)\}`$ has $`\mu _X`$-measure $`1`$. Applying Weil’s formula (1) to the function $`f=\text{1}\text{I}_{A_t}L^1(X)`$ gives
$$_{X/H}{}_{}{}^{\mathrm{}}f𝑑\mu _{X/H}=_{X/H}{}_{}{}^{\mathrm{}}\text{1}\text{I}_{A_t}^{}(s+H)𝑑\mu _{X/H}(s+H)=_Xf𝑑\mu _X=1.$$
(2)
Plugging the definition of $`^{\mathrm{}}`$ into (2) we get $`\mu _{X/H}`$-a.e. the identity
$${}_{}{}^{\mathrm{}}\text{1}\text{I}_{A_t}^{}(s+H)=_H\text{1}\text{I}_{A_t}(s+u)𝑑\mu _H(u)=1.$$
So for every $`tH`$ and $`\mu _{X/H}`$-a.e. $`s+H`$ we know that the set $`\{uH:\phi ^{}(s+t+u)\phi ^{}(s+u)\}`$ is a $`\mu _H`$-null set. This means
$$\tau _t(\tau _s\phi _{}^{}{}_{|H}{}^{})=\tau _s\phi _{}^{}{}_{|H}{}^{}\mu _H\text{-a.e.}$$
Thus $`\tau _s\phi ^{}`$ is constant $`\mu _H`$-a.e. on $`H`$ and for $`\mu _{X/H}`$ almost all $`s+H`$ we have
$`{}_{}{}^{\mathrm{}}\phi _{}^{}(s+H)`$ $`=`$ $`{\displaystyle _H}\tau _s\phi ^{}(t)𝑑\mu _H(t)={\displaystyle _H}\phi (s)^{}𝑑\mu _H(t)=\phi ^{}(s).`$
Let $`\pi _H:XX/H`$ be the quotient mapping onto the group compactification $`(\iota _{X/H},X/H)`$. Let $`\psi ^{}:={}_{}{}^{\mathrm{}}\phi _{}^{}\pi _H`$. Since $`{}_{}{}^{\mathrm{}}\phi _{}^{}`$ is Riemann integrable on $`X/H`$ it is an elementary fact that $`\psi ^{}`$ is Riemann integrable on $`X`$ (cf. ). Once again, Weil’s formula (1) together with the fact that the Haar measure on the quotient $`X/H`$ is given by $`\mu _{X/H}=\pi _H^1(\mu _X)`$ implies $`\psi ^{}=\phi ^{}`$ $`\mu _X`$-a.e. Thus the function $`\psi `$ defined by
$$\psi :=\psi ^{}\iota _X={}_{}{}^{\mathrm{}}\phi _{}^{}\iota _{X/H}$$
satisfies $`m_G(|\phi \psi |)=\phi ^{}\psi ^{}_1=0`$ for the unique invariant mean $`m_G`$. Thus $`\psi ^{}`$ is the required almost realization of $`\phi `$. ∎
###### Corollary 5.
Every $`\phi AP(G)`$ has an aperiodic realization on some group compactification $`(\iota _X,X)`$.
###### Proof.
We use the notation from Theorem 2. If $`\phi `$ is almost periodic then $`\phi ^{}`$ is continuous. Consequently $`{}_{}{}^{\mathrm{}}\phi _{}^{}`$ and $`\psi ^{}:={}_{}{}^{\mathrm{}}\phi _{}^{}\pi `$ are also continuous. Since these functions coincide $`\mu _X`$-a.e. they coincide everywhere on $`X`$. This implies that $`\phi ^{}`$ is constant on $`H`$-cosets and $`{}_{}{}^{\mathrm{}}\phi _{}^{}(s+H)=\phi ^{}(s)`$ for all $`s+HX/H`$. So $`\phi ^{}`$ is truly a realization of $`\phi `$, not only an almost realization. ∎
This Corollary is a special case of Følner’s ”Main Theorem for Almost Periodic Functions”, for a detailed treatment cf. .
*Remark:* Note that for any given realization of a Hartman measurable function $`\phi (G)`$ on a group compactification $`(\iota _X,X)`$ we can always assume that there exists an aperiodic almost realization of $`\phi `$ on a group compactification $`(\stackrel{~}{\iota _X},\stackrel{~}{X})`$ with $`(\stackrel{~}{\iota _X},\stackrel{~}{X})(\iota _X,X)`$. Since in it is shown that every Hartman measurable function on an LCA group with separable dual has a realization on a metrizable group compactification, every Hartman measurable function on such a group has an aperiodic almost realization on a metrizable group compactification.
###### Lemma 6.
Let $`G`$ be an LCA group and let $`(\iota _X,X)`$ be a group compactification. Then there exists a unique subgroup $`\mathrm{\Gamma }\widehat{G}`$ such that $`(\iota _\mathrm{\Gamma },C_\mathrm{\Gamma })`$ and $`(\iota _X,X)`$ are equivalent. Furthermore $`(\iota _X,X)`$ is the supremum of all group compactifications $`(\iota _\gamma ,C_\gamma )`$ such that $`(\iota _\gamma ,C_\gamma )(\iota _X,X)`$ (writing in short $`(\iota _\gamma ,C_\gamma )`$ for $`(\iota _\gamma ,C_\gamma )`$).
The mapping $`(\iota _X,X)C_\mathrm{\Gamma }`$ is a bijection between equivalence classes of group compactifications of $`G`$ and subgroups of $`\widehat{G}`$.
###### Proof.
See Theorem 26.13 in . ∎
###### Corollary 7.
Let $`\phi (G)`$. Any two group compactifications $`(\iota _{X_1},X_1)`$ and $`(\iota _{X_2},X_2)`$ on which $`\phi `$ has an aperiodic almost realization are equivalent.
###### Proof.
By Theorem 1 we have $`𝔘_{(\iota _{X_1},X_1)}=(\phi )=𝔘_{(\iota _{X_2},X_2)}`$. A straight forward generalization of Theorem 1 in implies that the mapping
$$\mathrm{\Phi }:\widehat{G}\mathrm{\Gamma }(\iota _\mathrm{\Gamma },C_\mathrm{\Gamma })$$
coincides with the composition of the mappings
$`\mathrm{\Sigma }:(\iota _b,bG)(\iota _X,X)`$ $``$ $`𝔘_{(\iota _X,X)},`$
$`\mathrm{\Psi }:𝔓(G)`$ $``$ $`\{\chi \widehat{G}:\underset{gG}{lim}\chi (g)=0\}.`$
Since Lemma 6 states that $`\mathrm{\Phi }=\mathrm{\Psi }\mathrm{\Sigma }`$ is invertible, $`\mathrm{\Sigma }`$ must be one-one. In particular $`𝔘_{(\iota _{X_1},X_1)}=𝔘_{(\iota _{X_2},X_2)}`$ implies that $`(\iota _{X_1},X_1)`$ and $`(\iota _{X_2},X_2)`$ are equivalent group compactifications. ∎
For the rest of this section assume that $`G`$ is an LCA group with separable dual.
###### Corollary 8.
Every filter $`(\phi )`$ with $`\phi (G)`$ coincides with a filter $`𝔘_{(\iota _X,X)}`$ for a metrizable group compactification $`(\iota _X,X)`$. If $`\phi ^{}`$ is an arbitrary realization of $`\phi `$, say on the Bohr compactification $`bG`$, we can take $`XbG/\mathrm{ker}d_\phi ^{}`$.
###### Corollary 9.
Hartman measurable functions induce exactly the filters coming from metrizable group compactifications.
###### Proof.
In Theorem 3 in for every metrizable group compactification $`(\iota _X,X)`$ of the integers $`G=`$, an aperiodic Hartman periodic function of the form $`f=\text{1}\text{I}_A`$ is constructed. The same construction can be done in an arbitrary LCA group $`G`$ as long as the dual $`\widehat{G}`$ contains a countable and dense subset. This shows that any $`𝔘_{(\iota _X,X)}`$ with metrizable $`X`$ can be obtained already by Hartman measurable *sets*, i.e. by a filter $`(\phi )`$ with $`\phi =\text{1}\text{I}_A`$. Since any Hartman measurable function on $`G`$ can be realized on a metrizable group compactification (cf. ). Thus Theorem 1 implies that no filter $`(\phi )`$ can coincide with $`𝔘_{(\iota _X,X)}`$ for a non metrizable group compactification $`(\iota ,C)`$. ∎
## 4 Subgroups associated with Hartman measurable functions
For Hartman measurable $`\phi `$ let us denote by $`\mathrm{\Gamma }(\phi )`$ the (countable) subgroup of $`\widehat{G}`$ generated by the set
$$\text{spec}\phi :=\{\chi \widehat{G}:m_G(\phi \overline{\chi })0\}$$
of all characters with non vanishing Fourier coefficients. We will prove that $`\mathrm{\Gamma }=\mathrm{\Gamma }(\phi )`$ determines a group compactification $`(\iota _\mathrm{\Gamma },C_\mathrm{\Gamma })`$ such that $`\phi `$ can be realized aperiodically on $`C_\mathrm{\Gamma }`$. First we deal with almost periodic functions:
###### Proposition 10.
Let $`\phi AP(G)`$ and $`(\iota _X,X)`$ a group compactification such that every character $`\chi \mathrm{\Gamma }(\phi )`$ has a representation $`\chi =\eta \iota _X`$ with a continuous character $`\eta \widehat{X}`$. Then every function $`f\overline{\text{span}}\mathrm{\Gamma }(\phi )AP(G)`$ has a realization on $`(\iota _X,X)`$.
###### Proof.
This is essentially a reformulation of Theorem 5.7 in . In fact the Stone-Weierstrass Theorem implies that $`\overline{\text{span}}\mathrm{\Gamma }(\phi )=\iota _\mathrm{\Gamma }^{}C(X)`$. Furthermore $`\phi \overline{\text{span}}\mathrm{\Gamma }(\phi )`$, i.e. $`\phi `$ can be realized by some continuous $`\phi ^{}:X`$. ∎
###### Proposition 11.
Let $`\phi AP(G)`$ and $`(\iota _\mathrm{\Gamma },C_\mathrm{\Gamma })`$ the group compactification of $`G`$ induced by the subgroup $`\mathrm{\Gamma }=\mathrm{\Gamma }(\phi )\widehat{G}`$. Then for every continuous character $`\psi \mathrm{\Gamma }(\phi )`$ there exists a continuous $`\psi ^{}:C_\mathrm{\Gamma }`$ such that $`\psi =\psi ^{}\iota _\mathrm{\Gamma }`$.
###### Proof.
Given the group compactification $`(\iota _\mathrm{\Gamma },C_\mathrm{\Gamma })`$, then the compact group $`C_\mathrm{\Gamma }`$ is by definition topologically isomorphic to $`\overline{\{(\chi (g))_{\chi \mathrm{\Gamma }}:gG\}}𝕋^\mathrm{\Gamma }`$.
The restriction of each projection
$$\pi _{\chi _0}:C_\mathrm{\Gamma }𝕋^\mathrm{\Gamma }𝕋,(x_\chi )_{\chi \mathrm{\Gamma }}x_{\chi _0}$$
is a bounded character of $`C_\mathrm{\Gamma }`$ for each $`\chi _0\mathrm{\Gamma }(\phi )`$. I.e. $`\pi _{\chi _0}`$ is an element of $`\widehat{C_\mathrm{\Gamma }}`$. Thus $`\chi _0=\pi _{\chi _0}\iota _\mathrm{\Gamma }`$ for each $`\chi _0\mathrm{\Gamma }(\phi )`$ and we may apply Proposition 10 to obtain the assertion. ∎
###### Proposition 12.
Let $`\phi AP(G)`$ and let $`(\iota _X,X)`$ be a group compactification of $`G`$ such that $`\phi `$ can be realized by a continuous function $`\phi ^{}:X`$. Then each continuous character $`\chi \mathrm{\Gamma }(\phi )`$ has a representation $`\chi =\eta \iota _\mathrm{\Gamma }`$ with $`\eta \widehat{X}`$.
###### Proof.
Obviously it is enough to prove the assertion for a generating subset of $`\mathrm{\Gamma }(\phi )`$. Let $`\chi \widehat{G}`$ be such that $`m_G(\phi \overline{\chi })0`$. Define a linear functional $`m_\chi :C(X)`$ via $`\psi m_\chi (\psi )=m_G((\psi \iota _\mathrm{\Gamma })\overline{\chi })`$. It is routine to check that $`m_\chi `$ is bounded and $`m_\chi =1`$. Since $`X`$ is compact the complex-valued mapping $`\stackrel{~}{\eta }:Xm_\chi (\tau _x\phi ^{})`$ is continuous on $`X`$ (the mapping $`x\tau _x\phi ^{}`$ is continuous). For $`gG`$ we compute
$`\stackrel{~}{\eta }\iota _X(g)`$ $`=`$ $`m_G((\tau _{\iota _X(g)}\phi ^{}\iota _X)\overline{\chi })=m_G(\tau _g(\phi ^{}\iota _X)\overline{\chi })`$
$`=`$ $`m_G((\phi ^{}\iota _X)\tau _g\overline{\chi })=m_G((\phi ^{}\iota _X)\chi (g)\overline{\chi })`$
$`=`$ $`\chi (g)m_\chi (\phi ^{})=\chi (g)\stackrel{~}{\eta }(0).`$
Since $`\stackrel{~}{\eta }(0)=m_\chi (\phi ^{})=m_G(\phi \overline{\chi })0`$ we can define $`\eta :=\stackrel{~}{\eta }(0)^1\stackrel{~}{\eta }`$. The mapping $`\eta :X𝕋`$ is continuous and satisfies the functional equation
$$\eta (\iota _X(g)+\iota _X(h))=\stackrel{~}{\eta }(0)^1\stackrel{~}{\eta }(\iota _X(g)+\iota _X(h))=\chi (g)\chi (h)=\eta (\iota _X(g))\eta (\iota _X(h))$$
on the dense set $`\iota _X(G)`$. Hence $`\eta `$ is a bounded character on $`X`$ and $`\eta \iota _X=\chi `$. ∎
###### Corollary 13.
Let $`\phi (G)`$ be realized by $`\phi ^{}`$ on the group compactification $`(\iota _X,X)`$. Then each $`\chi \mathrm{\Gamma }(\phi )`$ has a representation $`\chi =\eta \iota _X`$ with $`\eta \widehat{X}`$.
###### Proof.
For every $`\chi \widehat{G}`$ with $`m_G(\phi \overline{\chi })=\alpha 0`$ we can pick a continuous function $`\psi ^{}:X`$ such that $`\psi ^{}\phi _1<|\alpha |/2`$. Then $`\psi :=\psi ^{}\iota _X`$ satisfies
$$|m_G(\phi \overline{\chi })m_G(\psi \overline{\chi })|m_G(|\phi \psi |)\psi ^{}\phi ^{}_1<|\alpha |/2.$$
In particular $`m_G(\psi \overline{\chi })0`$. Applying Proposition 12 to the function $`\psi AP(G)`$ yields that the character $`\chi `$ can be realized on $`X`$.∎
Thus for almost periodic functions $`\phi `$ the subgroup $`\mathrm{\Gamma }(\phi )`$ contains all the relevant information to reconstruct $`\phi `$ from its Fourier-data in a minimal way. It is not obvious how to obtain similar results for Hartman measurable functions that are not almost periodic. The following example illustrates how a straight forward approach may fail.
###### Example 14.
Let $`\phi _n(k):=_{j=1}^n\mathrm{cos}^2\left(2\pi \frac{k}{3^j}\right)`$ on $`G=`$. Each $`\phi _n`$ is a finite product of periodic (and hence almost periodic) functions. Since $`AP()`$ is an algebra, $`\phi `$ is almost periodic. In it is shown that $`\phi (k):=lim_n\mathrm{}\phi _n(k)`$ exists and defines a non negative Hartman measurable function with $`m_{}(\phi )=0`$. Since $`\mathrm{\Gamma }(\phi _n)/3^n`$ we have (using obvious notation):
$$\underset{n\mathrm{}}{lim}\mathrm{\Gamma }(\phi _n)=\underset{n=1}{\overset{\mathrm{}}{}}\mathrm{\Gamma }(\phi _n)_3^{\mathrm{}},$$
the Prüfer 3-group (i.e. the subgroup of all complex $`3^n`$-th roots of unity for $`n`$), but
$$\mathrm{\Gamma }(\underset{n\mathrm{}}{lim}\phi _n)=\{0\}.$$
###### Proposition 15.
Let $`\{K_n\}_{n=1}^{\mathrm{}}`$ denote the family of Fejér kernels on $`𝕋^k`$
$$K_n(\mathrm{exp}(it_1),\mathrm{},\mathrm{exp}(it_k))=\frac{1}{k}\underset{j=1}{\overset{k}{}}\left(\frac{\mathrm{sin}(\frac{1}{2}nt_j)}{\mathrm{sin}(\frac{1}{2}t_j)}\right)^2.$$
The linear convolution operators on $`L^1(𝕋^k)`$ defined by
$$\sigma _n:\phi K_n\phi $$
are non negative, their norm is uniformly bounded by $`\sigma _n=1`$ and $`\sigma _n\phi (x)\phi (x)`$ a.e. for every $`\phi L^1(𝕋^k)`$. Furthermore $`\sigma _n\phi \text{span }\mathrm{\Gamma }(\phi )`$ for every $`n`$.
###### Proof.
This is a reformulation of the results in section 44.51 in .∎
Let $`f`$ be Riemann integrable on $`X=𝕋^k`$, w.l.o.g. real-valued, and $`\phi _i,\psi _iC(X)`$ such that $`\phi _if\psi _i`$ and $`\phi _i\psi _i_1<\epsilon _i`$ for a sequence $`\{\epsilon _i\}_{i=1}^{\mathrm{}}`$ of positive real numbers, tending monotonically to $`0`$. We know that $`\sigma _nf(x)f(x)`$ for a.e. $`xX`$. Thus we have
$$\phi _n^{}:=\sigma _n\phi _n\sigma _nf\sigma _n\psi _n=\psi _n^{}$$
and
$$\phi _n^{}\psi _n^{}_1\sigma _n(\phi _n^{}\psi _n^{})_1\sigma _n\phi _n\psi _n_1\epsilon _n.$$
Let $`\phi ^{}:=inf_n\phi _n`$ and $`\psi ^{}:=sup_n\psi _n`$. If we assume w.l.o.g. $`\psi _n`$ to increase and $`\phi _n`$ to decrease as $`n\mathrm{}`$, the same will hold for $`\psi _n^{}`$ and $`\phi _n^{}`$. This implies that in the inequality
$$\phi ^{}(x)=\underset{n\mathrm{}}{lim}\phi _n^{}(x)\underset{n\mathrm{}}{lim\; sup}\sigma _nf\underset{n\mathrm{}}{lim\; inf}\sigma _nf\underset{n\mathrm{}}{lim}\psi _n^{}(x)=\psi ^{}(x)$$
actually equality holds $`\mu _X`$-a.e. on $`X`$. Thus we can apply Proposition 1 and conclude that any function $`f^{}`$ with $`\phi ^{}f^{}\psi ^{}`$ is Riemann integrable (and coincides $`\mu _X`$-a.e. with $`f`$). In particular $`f^{}:=lim\; sup_n\mathrm{}\sigma _nf`$ and $`f_{}:=lim\; inf_n\mathrm{}\sigma _nf`$ are (lower resp. upper semicontinuous) Riemann integrable functions that coincide $`\mu _X`$-a.e. with $`f`$.
Let us call a group compactification $`(\iota _X,X)`$ finite dimensional iff $`X`$ is topologically isomorphic to a closed subgroup of $`𝕋^n`$ for some $`n`$. Note that if $`(\iota _X,X)`$ is finite dimensional, then every group compactification covered by $`(\iota _X,X)`$ is finite dimensional as well. A Hartman measurable function $`\phi (G)`$ can be realized finite dimensionally iff there exists a realization of $`\phi `$ on some finite dimensional group compactification.
###### Proposition 16.
For a compact LCA group $`C`$ the following assertions are equivalent:
1. $`C`$ is finite dimensional,
2. $`\widehat{C}`$ is finitely generated,
3. $`C`$ is topological isomorphic to $`𝕋^k\times F`$ for $`k`$ and a finite group $`F`$ of the form
$$F=\underset{i=1}{\overset{N}{}}(/n_i)^{p_i},p_i\text{ prime}.$$
###### Proof.
Folklore. ∎
###### Proposition 17.
Let $`\phi (G)`$. If $`\phi `$ can be realized finite dimensionally, then there is an almost realization of $`\phi `$ on the (finite dimensional) compactification induced by $`\mathrm{\Gamma }:=\mathrm{\Gamma }(\phi )`$.
###### Proof.
Let $`\phi `$ be realized finite dimensionally on some group compactification $`(\iota _X,X)`$. Since there exists a group compactification covered by $`(\iota _X,X)`$, on which $`\phi `$ can be almost realized aperiodically (cf. Theorem 2), we can assume w.l.o.g. that $`\phi `$ can be almost realized aperiodically already on $`(\iota _X,X)`$. We have to show that $`(\iota _X,X)`$ and $`(\iota _\mathrm{\Gamma },C_\mathrm{\Gamma })`$ are equivalent.
Let $`\psi ^{}`$ be an aperiodic realization of $`\phi `$ on $`C_\mathrm{\Gamma }𝕋^k\times F`$ with $`k`$ and $`F`$ finite. Let us denote the elements of $`𝕋^k\times F`$ by tuples $`(\stackrel{}{\alpha },x)`$. For every fixed $`\stackrel{}{\alpha }𝕋^k`$ define a mapping $`\psi _\stackrel{}{\alpha }:F`$ via
$$\psi _\stackrel{}{\alpha }^{}(x):=\psi ^{}(\stackrel{}{\alpha },x).$$
For each $`\overline{\chi }\widehat{F}`$, the dual of the finite group $`F`$, define the $`F`$-Fourier coefficient of $`\psi _\stackrel{}{\alpha }^{}`$ as
$$c_{\overline{\chi }}(\stackrel{}{\alpha }):=_F\psi _\stackrel{}{\alpha }^{}(x)\overline{\chi }(x)𝑑x=\frac{1}{\mathrm{\#}F}\underset{xF}{}\psi ^{}(\stackrel{}{\alpha };x)\overline{\chi }(x).$$
We want to show that $`c_{\overline{\chi }}:𝕋^k`$ is a Riemann integrable function: The mapping $`\gamma _x:𝕋^kX`$ defined via $`\stackrel{}{\alpha }(\stackrel{}{\alpha };x)`$ is continuous and measure-preserving for every $`xF`$. $`\psi ^{}`$ is by definition Riemann integrable. Thus the mapping $`\psi ^{}\gamma _x:𝕋^k`$ is Riemann integrable for each $`xF`$. Note that
$$c_{\overline{\chi }}(\stackrel{}{\alpha })=\underset{xF}{}(\psi ^{}\gamma _x)(\stackrel{}{\alpha })\overline{\chi }(x).$$
Hence, for each fixed character $`\chi \widehat{F}`$, the mapping $`c_{\overline{\chi }}:𝕋^k`$ defined via $`\stackrel{}{\alpha }_{xF}(\psi ^{}\gamma _x)(\stackrel{}{\alpha })\overline{\chi }(x)`$ is Riemann integrable on $`𝕋^k`$.
Thus Proposition 15 implies $`\sigma _nc_\chi (\stackrel{}{\alpha })c_\chi (\stackrel{}{\alpha })`$ a.e. on $`𝕋^k`$. Taking into account that the Haar measure on $`F`$ is the normalized counting measure, we get
$$\psi _n^{}(\stackrel{}{\alpha };x):=\underset{\overline{\chi }\widehat{F}}{}\left(\sigma _nc_{\overline{\chi }}(\stackrel{}{\alpha })\right)\overline{\chi }(x)\underset{\overline{\chi }\widehat{F}}{}c_{\overline{\chi }}(\stackrel{}{\alpha })\overline{\chi }(x)=\psi _\stackrel{}{\alpha }^{}(x)=\psi ^{}(\stackrel{}{\alpha };x)$$
(3)
for almost every $`\stackrel{}{\alpha }𝕋^k`$ and every $`xF`$, as $`n\mathrm{}`$. Since Haar measure $`\mu _C`$ on $`C`$ is the product measure of the Haar measures on the groups $`𝕋^k`$ and $`F`$, the relation (3) holds $`\mu _C`$-a.e. on $`C`$. We conclude that any function majorizing $`lim\; inf_n\mathrm{}\psi _n^{}`$ and minorizing $`lim\; sup_n\mathrm{}\psi _n^{}`$ is an almost realization of $`\phi `$. Note that according to the properties of the Fejér kernels on $`𝕋^k`$ (see 44.51 in ) for each character $`(\eta \times \chi )(\stackrel{}{\alpha };x):=\eta (\stackrel{}{\alpha })\chi (x)`$, $`\eta \widehat{𝕋^k}`$ and $`\chi \widehat{F}`$, there exists an $`n_0N`$ such that for $`nn_0`$ in the Fourier expansion of $`\psi _n^{}`$ the Fourier coefficient (computed in $`C`$) associated with the character does not vanish iff the $`𝕋^k`$-Fourier coefficient of $`c_\chi `$
$$c_\eta (c_\chi )=_{𝕋^k}c_\chi (\stackrel{}{\alpha })\overline{\eta }(\stackrel{}{\alpha })𝑑\stackrel{}{\alpha }$$
does not vanish. A simple computation shows that the Fourier coefficients of $`\psi ^{}`$ are given by
$`c_{\eta \times \chi }(\psi ^{})`$ $`=`$ $`{\displaystyle _{𝕋^k}}{\displaystyle _F}\psi ^{}(\stackrel{}{\alpha },x)\overline{\eta }(\stackrel{}{\alpha })\overline{\chi }(x)𝑑\stackrel{}{\alpha }𝑑x`$
$`=`$ $`{\displaystyle _{𝕋^k}}c_\chi (\stackrel{}{\alpha })\overline{\eta }(\stackrel{}{\alpha })𝑑\stackrel{}{\alpha }=c_\eta (c_\chi )`$
So the character $`\eta \times \chi `$ contributes to the Fourier expansion of $`\psi ^{}`$ if and only if $`c_{\eta \times \chi }(\psi )0`$. Thus $`\psi _n^{}\text{span }\mathrm{\Gamma }(\phi )`$ for every $`n`$, implying that there exist almost realizations of $`\phi `$ on the group compactification induced by $`\mathrm{\Gamma }(\phi )`$, e.g. $`lim\; inf_n\mathrm{}\psi _n^{}`$ or $`lim\; sup_n\mathrm{}\psi _n^{}`$. ∎
Combining this result with the results of the previous section we obtain
###### Theorem 3.
Let $`\phi (G)`$ and $`\mathrm{\Gamma }=\mathrm{\Gamma }(\phi )\widehat{G}`$. The following assertions hold:
1. $`(\iota _\mathrm{\Gamma },C_\mathrm{\Gamma })(\iota _X,X)`$ for every compactification $`(\iota _X,X)`$ on which $`\phi `$ can be realized. In particular $`(\phi )𝔘_{(\iota _\mathrm{\Gamma },C_\mathrm{\Gamma })}`$.
2. Assume that $`\phi AP(G)`$ or that $`\phi `$ can be realized finite dimensionally. Then $`\phi `$ can be realized aperiodically on $`C_\mathrm{\Gamma }`$. In particular $`(\phi )=𝔘_{(\iota _\mathrm{\Gamma },C_\mathrm{\Gamma })}`$.
We strongly conjecture that the second assertion in Theorem 3 holds for *any* Hartman measurable function, at least on LCA groups $`G`$ with separable dual $`\widehat{G}`$. A proof of this might utilize more general summation methods (in the flavour of Theorems 44.43 and 44.47 in ) than the Fejér summation presented here.
In it is shown that for any Hartman measurable set $`MG=`$ and the induced filter $`=(M)`$ there is an aperiodic realization of $`\phi _M=\text{1}\text{I}_M`$ on the compactification determined by the subgroup Sub$`(M)=\{\alpha :lim_nn\alpha =0\}`$ or, equivalently, Sub$`(M)=\{\alpha :lim_ne^{2\pi in\alpha }=1\}`$.
Together with Theorem 3 this implies that for Hartman sets $`M`$ with finite dimensional realization both the group compactifications of $``$ induced by the subgroups $`\mathrm{\Gamma }(\phi _M)`$ and Sub$`(M)`$ admit aperiodic realizations of $`\phi _M`$. Hence uniqueness of the minimal compactification with aperiodic realization (Corollary 7) implies that in this special case $`\mathrm{\Gamma }(\phi _M)=\text{Sub}(M)`$. In the general situation we can prove up to now far only the following
###### Proposition 18.
For a Hartman measurable function $`\phi (G)`$ let $`=(\phi )`$, $`\mathrm{\Gamma }=\mathrm{\Gamma }(\phi )`$ and $`\text{Sub}(\phi )=\{\chi \widehat{G}:lim_{gG}\chi (g)=1_{}\}`$. Then $`\mathrm{\Gamma }(\phi )\text{Sub}(\phi )`$.
###### Proof.
Suppose $`\chi \mathrm{\Gamma }(\phi )`$. To prove $`lim_{gG}\chi (g)=1_{}`$ (unit element of the multiplicative group of complex numbers) we have to show that for every $`\epsilon >0`$ the set $`\{gG:|1\chi (g)|<\epsilon \}`$ belongs to the filter $`(\phi )`$, i.e. that there exists $`\delta =\delta (\epsilon )>0`$ such that
$$\{gG:m_G(|\tau _g\phi \phi |)<\delta \}\{gG:|1\chi (g)|<\epsilon \}(\phi ).$$
(4)
Using the fact that $`m_G`$ is an invariant mean and that $`\chi `$ is a homomorphism, we have
$$\chi (g)m_G(\phi \overline{\chi })=m_G(\tau _g\phi \overline{\chi })=m_G((\tau _g\phi \phi )\overline{\chi })+m_G(\phi \overline{\chi }).$$
Using $`\chi _{\mathrm{}}=1`$ this implies
$$|1\chi (g)||m_G(\phi \overline{\chi })|=|m_G((\tau _g\phi \phi )\overline{\chi })|m_G(|\tau _g\phi \phi |).$$
Since $`m_G(\phi \overline{\chi })0`$ we can define $`\delta :=\epsilon \frac{m_G(|\tau _g\phi \phi |)}{|m_G(\phi \overline{\chi })|}>0`$. With this choice of $`\delta `$ indeed $`m_G(|\tau _g\phi \phi |)<\delta `$ implies $`|1\chi (g)|<\epsilon `$, i.e. the inclusion (4) holds. ∎
## 5 Summary
The content of the present paper essentially deals with the definition and properties of the objects occurring in the diagram below. Abusing the terminus technicus of *commutative diagrams* in a kind of sloppy way, the theorems of this paper circle around the question under which assumptions this diagram is commutative:
Section 3 deals with the left half of this diagram: to every Hartman measurable function $`\phi `$ a filter $`(\phi )`$ is associated and to every group compactification $`(\iota _X,X)`$ on which $`\phi `$ can be realized a filter $`𝔘(X,0_X)`$ is associated. In general $`𝔘(X,0_X)(\phi )`$, and there always exists compactifications such that equality holds (indicated by $``$).
Section 4 deals with the right half of the diagram: to every Hartman measurable function $`\phi `$ a subgroup $`\mathrm{\Gamma }(\phi )`$ of the dual is associated, which in turn induces a group compactification $`(\iota _\mathrm{\Gamma },C_\mathrm{\Gamma })`$. In general $`(\iota _\mathrm{\Gamma },C_\mathrm{\Gamma })(\iota _X,X)`$ for every group compactification $`(\iota _X,X)`$ on which $`\phi `$ can be realized. If $`\phi `$ is either almost periodic or can be realized finite dimensionally then $`(\iota _A,X_A)`$ is itself a group compactification on which $`\phi `$ can be realized (indicated by $`\genfrac{}{}{0pt}{}{}{}`$) and the filter $`𝔘(X_A,0_{X_A})`$ associated with this particular compactification coincides with $`(\phi )`$. The filter $`(\phi )`$ in turn defines a subgroup Sub($`\phi `$) of the dual $`\widehat{G}`$. While it can be shown that in general $`\mathrm{\Gamma }(\phi )\text{Sub}(\phi )`$ it is and open problem whether this inclusion can be reversed.
Author’s address:
Gabriel Maresch
Technical University Vienna
Institute of Discrete Mathematics and Geometry
Wiedener Hauptstraße 8-10
1040 Vienna, Austria
email: gabriel.maresch@tuwien.ac.at
web: http://www.dmg.tuwien.ac.at/maresch |
warning/0506/math0506106.html | ar5iv | text | # Contents
## 1 Introduction
In Griffiths introduced the Moduli of Polarized Hodge Structures/the period domain $`D`$ and described a dream to enlarge $`D`$ to a moduli space of degenerating polarized Hodge structures. Since in general $`D`$ is not a Hermitian symmetric domain, he asked for the existence of a certain automorphic cohomology theory for $`D`$, generalizing the usual notion of automorphic forms on symmetric Hermitian domains. Since then there have been many efforts in the first part of Griffiths’s dream (see and references there) but the second part still lives in darkness.
I was looking for some analytic spaces over $`D`$ for which one may state Baily-Borel theorem on the unique algebraic structure of quotients of symmetric Hermitian domains by discrete arithmetic groups (see ). I realized that even in the simplest case of Hodge structures, namely $`h^{01}=h^{10}=1`$, such spaces are not well studied. This led me to the definition of a new class of holomorphic functions on the Poincaré upper half plane which generalize the classical modular forms. Since a differential operator acts on them we call them differential modular forms. These new functions are no longer interpreted as holomorphic sections of a positive line bundle on some compactified moduli curve. Nevertheless, they appear in a natural way as coefficients in families of elliptic curves, analogus to Eisenstein series in the Weierstrass Uniformization Theorem. Due to this, I realized that there are very natural holomorphic foliations in the coefficient space of families of varieties whose dynamics is in close relation with the abelian integrals (resp. Hodge structures) of the family. In the case of a three parametric elliptic curve, the mentioned foliation is called Ramanujan foliation because he was the first who obtained a particular leaf of this foliation using Eisenstein series (see Chapter X and (8) bellow). This article will touch some aspects of number theory, holomorphic foliations and Hodge theory which I am going to explain below:
Number theory: Recall the Eisenstein series
(1)
$$g_k(z)=a_k\left(1+(1)^k\frac{4k}{B_k}\underset{n1}{}\sigma _{2k1}(n)e^{2\pi izn}\right),k=1,2,3,z$$
where $`B_k`$ is the $`k`$-th Bernoulli number ($`B_1=\frac{1}{6},B_2=\frac{1}{30},B_3=\frac{1}{42},\mathrm{}`$), $`\sigma _i(n):=_{dn}d^i`$,
(2)
$$a_1=2\zeta (2)\frac{1}{2\pi i},a_2=2\zeta (4)\frac{60}{(2\pi i)^2},a_3=2\zeta (6)\frac{140}{(2\pi i)^3}$$
and $`:=\{x+iyy>0\}`$ is the Poincaré upper half plane. The most well-known differential modular form, which is not a differential of a modular form, is the Eisenstein series $`g_1`$. The idea of differentiating modular forms and getting new modular forms is old and goes back to Ramanujan. Nesterenko’s method (see ) for the proof of transcendency properties of certain numbers is also based on differential equations satisfied by modular forms. However, the precise definition of differential modular forms has been given recently in . In the present article we give another slightly different definition of differential modular forms (see §2.1) over a modular subgroup $`\mathrm{\Gamma }\mathrm{SL}(2,)`$. It is based on a canonical behavior of holomorphic functions on the Poincaré upper half plane under the action of $`\mathrm{SL}(2,)`$. This approach has the advantage that it can be generalized to any modular subgroup of $`\mathrm{SL}(2,)`$ but the one in works only in the case of full modular group $`\mathrm{SL}(2,)`$. The set of differential modular forms in the present article is a bigraded algebra $`M=_{n_0,m}M_m^n`$, $`M_m^0`$ being the set of classical modular forms of weight $`m`$, in which the differential operator $`\frac{d}{dz}`$ maps $`M_m^n`$ to $`M_{m+2}^{n+1}`$. We have $`g_1M_2^1,g_2M_4^0,g_3M_6^0`$ and we prove:
###### Theorem 1.
The functions $`g_1,g_2,g_3`$ are algebraically independent and $`M`$ is freely generated by $`g_1,g_2`$ and $`g_3`$ as a $``$-algebra. In particular each $`M_m^n`$ is a finite dimensional $``$-vector space.
This theorem generalizes the first theorem in each modular forms book that the algebra of modular forms is generated freely by the Eisenstein series $`g_2`$ and $`g_3`$. Our proof gives us also the Ramanujan relations between the $`g_i`$’s. We define the action of Hecke operators on $`M_m^n`$ and it turns out that this is similar to the case of modular forms:
(3)
$$T_pf(z)=p^{mn1}\underset{dp,0bd1}{}d^mf\left(\frac{nz+bd}{d^2}\right),p,fM_m^n.$$
Hecke operators of this type appear in particular in the study of the transfer operator from statistical mechanics which plays an important role in the theory of dynamical zeta functions (see ). The differential operator commutes with Hecke operators (§2.2) and so it induces a map from the set of new differential forms (see §2.7) to itself. Another result which we prove in this article and which could be interesting from the number theory point of view is the following: Let
$$g:=(g_1,g_2,g_3):^3$$
and
$$T:=^3\backslash \{(t_1,t_2,t_3)^327t_3^2t_2^3=0\}.$$
###### Theorem 2.
There are unique analytic functions
$$B_1,B_2:T,B_3:T$$
such that $`B_1`$ does not depend on the variable $`t_1`$ and
(4)
$$B_1g(z)=\mathrm{Im}(z),B_1(t_1,t_2k^4,t_3k^6)=B_1(t)|k|^2$$
(5)
$$B_2g=0,B_2(t_1k^2+k^{}k^1,t_2k^4,t_3k^6)=B_1(t)|k^{}|^2+B_2(t)|k^1|^2+\mathrm{Im}(B_3(t)k^{}\overline{k^1})$$
(6)
$$B_3g=1,B_3(t_1k^2+k^{}k^1,t_2k^4,t_3k^6)=B_3(t)k\overline{k^1}+2\sqrt{1}k\overline{k^{}}B_1(t)$$
for all $`k^{}`$ and $`k^{}`$. Moreover, $`|B_3|`$ restricted to the zero locus of $`B_2`$ is identically one.
Holomorphic singular foliations: Nowadays the theory of holomorphic singular foliations is getting a part of Algebraic Geometry. For a general background in this theory see (). In this direction this article has two novelties which together exhibit a connection between the area of holomorphic foliations and Arithmetic Algebraic Geometry. The first one is as follows: After calculating the Gauss-Manin connection of the following family of elliptic curves
(7)
$$_t:y^24t_0(xt_1)^3+t_2(xt_1)+t_3,t^4$$
and considering its relation with the inverse of the period map, we get the following ordinary differential equation:
(8)
$$\mathrm{Ra}:\{\begin{array}{c}\dot{t}_1=t_1^2\frac{1}{12}t_2\hfill \\ \dot{t}_2=4t_1t_26t_3\hfill \\ \dot{t}_3=6t_1t_3\frac{1}{3}t_2^2\hfill \end{array}$$
which is called the Ramanujan relations/differential equation/foliation, because he had observed that $`g`$ is a solution of (8) (one gets the classical relations by changing the coordinates $`(t_1,t_2,t_3)(\frac{1}{12}t_1,\frac{1}{12}t_2,\frac{2}{3(12)^2}t_3)`$). We denote by $`(\mathrm{Ra})`$ the singular holomorphic foliation induced by (8) in $`^3`$. Its singularities
$$\mathrm{Sing}(\mathrm{Ra})=\{(t_1,12t_1^2,8t_1^3)t_1\}$$
form a one-dimensional curve in $`^3`$. Consider the family (7) with $`t_0=1`$ and define
$$K:=\{tT_\delta \frac{xdx}{y}=0,\text{ for some }\delta H_1(_t,)\}$$
and
$$M_r:=\{tTB_2(t)=r\},r.$$
The last part of Theorem 2 says that $`|B_3|`$ restricted to $`M_0`$ is identically $`1`$. We also define
$$N_w:=\{tM_0B_3(t)=w\},|w|=1,w.$$
For $`t^3\backslash \mathrm{Sing}((\mathrm{Ra}))`$ we denote by $`L_t`$ the leaf of $`(\mathrm{Ra})`$ through $`t`$.
###### Theorem 3.
The following is true:
1. The real analytic varieties $`M_r,r,N_w,|w|=1`$, and the set $`K`$ are $`(\mathrm{Ra})`$-invariant.
2. The set $`K`$ is a dense subset of $`M_0`$ with the following property: For all $`tK`$ the leaf $`L_t`$ intersects $`\mathrm{Sing}(\mathrm{Ra})`$ transversally at some point $`p`$.
3. For all $`tT`$ the leaf $`L_t`$ has an accumulation point at $`T`$ if and only if $`tM_0`$.
The item $`2`$ says that there is a transverse disk to $`\mathrm{Sing}((\mathrm{Ra}))`$ at some point $`p`$ such that $`D\backslash \{p\}`$ is a part of the leaf $`L_t`$. For the proof of the above theorem and a precise description of the foliation $`(\mathrm{Ra})`$ see §5. The proof of this theorem is based on the fact that the foliation $`(\mathrm{Ra})`$ restricted to $`T`$ is uniformized by the inverse of the period map (see for instance for similar topics).
The second novelty in the area of holomorphic foliations is as follows:
###### Theorem 4.
There is no elliptic curve $`E`$ and a differential form of the second type $`\omega `$ on $`E`$, both defined over $`\overline{}`$, such that
(9)
$$_\delta \omega =0$$
for some non-zero topological cycle $`\delta H_1(E,)`$.
This theorem uses Nesterenko’s Theorem (see ) on transcendence properties of the values of Eisenstein series. The above theorem for the case in which $`\omega `$ is of the first kind, is well-known. In this case we can even state it for the field $``$. However, it is trivially false when $`\omega `$ is a differential form of the second kind and we allow transcendental coefficients in $`\omega `$ or the elliptic curve. Generalizations of such theorems to arbitrary curves can be derived from the Abelian Subvariety Theorem (see , ). Abelian integrals of the type (9) appear in deformations of holomorphic foliations with a first integral in complex surfaces (see ). The above result shows that the zeros of arithmetic abelian integrals are quite different from complex ones.
Hodge theory: This article stimulates the hope to realize the second part of Griffiths’ dream with a different formulation. Differential modular forms can also be introduced for a complex manifold $`𝒫`$ over Griffiths period domain $`D`$ with an action of an algebraic group $`G_0`$ from the right. Since they are no longer interpreted as sections of positive line bundles over moduli spaces, the question of the existence of a kind of Baily-Borel Theorem for $`𝒫`$ arises. In the case of Hodge structures with $`h^{01}=h^{10}=1`$ we have
(10)
$$𝒫:=\{\left(\begin{array}{cc}x_1& x_2\\ x_3& x_4\end{array}\right)\mathrm{GL}(2,)\mathrm{Im}(x_1\overline{x_3})>0\},D=$$
(11)
$$G_0=\{\left(\begin{array}{cc}k_1& k_3\\ 0& k_2\end{array}\right)k_3,k_1,k_2^{}\}$$
and we show in this article that $`\mathrm{SL}(2,)\backslash 𝒫`$ has a canonical structure of an algebraic quasi-affine variety such that the action of $`G_0`$ from the right is algebraic. More precisely, we prove that $`\mathrm{SL}(2,)\backslash 𝒫`$ is biholomorphic to $`^4\backslash \{t=(t_0,t_1,t_2,t_3)^427t_0t_3^2t_2^3=0\}`$ and under this biholomorphism the action of $`G_0`$ is given by:
$$tg:=(t_0k_1^1k_2^1,t_1k_1^1k_2+k_3k_1^1,t_2k_1^3k_2,t_3k_1^4k_2^2)$$
(12)
$$t=(t_0,t_1,t_2,t_3)^4,g=\left(\begin{array}{cc}k_1& k_3\\ 0& k_2\end{array}\right)G_0$$
The mentioned biholomorphism is given by the period map (see §4).
Singularity theory: The Brieskorn module/lattice and its Gauss-Manin connection in singularity theory play the role of de Rham cohomology of fibered varieties, and it is a useful object when one wants to calculate the Gauss-Manin connection. In §3 we introduce the associated Brieskorn $`[t,\frac{1}{t_0}]`$-module $`H`$ of the family (7). We prove that $`H`$ is freely generated and we calculate the Gauss-Manin connection in a canonical basis of $`H`$. We generalize a classical Weierstrass Theorem and see that for the inverse of the period map, $`t_i`$ appears as the Eisenstein series $`g_i,i=1,2,3`$. The novelty is the appearance of $`t_1`$ as the Eisenstein series of weight $`2`$. This makes us think about a generalization of K. Saito’s primitive form theory (see ), to the case in which the deformation of a singularity is bigger than the versal deformation.
Final note: An elliptic curve (beside $`^1`$) is one of the most well-studied objects in (Arithmetic) Algebraic Geometry. This makes me feel that some parts of this work has connections to the works of many authors that I have not mentioned here. Here I express my sincere apologizes. My aim of writing this article was just to show, by some simple examples, how the dynamics of foliations and algebraic geometry of fibered varieties can be connected. One of the points in this article which I enjoyed very much, was the calculation of the Gauss-Manin connection and the Ramanujan relations by Singular, . The corresponding algorithms were developed in the articles and the corresponding library in Singular is called brho.lib. Using this library for families of algebraic varieties one can obtain certain holomorphic foliations, whose dynamics has to do with the Hodge structure of Algebraic varieties and in particular with abelian integrals. This opens a new connection between (Arithmetic) Algebraic Geometry and the theory of holomorphic foliations.
Let us now explain the structure of this article. §2 is devoted to the definition of differential modular forms and the action of Hecke operators on them. This section is independent of the other sections. We shall only use the property (19) of the Eisenstein series $`g_1`$. The reader who is interested only in §2, may lose our proof for the main theorem of this section (Theorem 1). §3 is devoted to calculation of the Gauss-Manin connection of the family (7). This section is based on the machinery introduced in . The heart of the present paper is §4 in which we prove that the period map $`\mathrm{𝗉𝗆}`$ is a biholomorphism and then we take the inverse of $`d\mathrm{𝗉𝗆}`$ to obtaine the Ramanujan relations. In §5 we have described the dynamics of the Ramanujan foliation and in particular its uniformization by the inverse of the period map in some quasi-affine variety. Finally, §6 is devoted to the proofs announced in the Introduction.
Acknowledgement: The main ideas of this paper took place in my mind when I was visiting Prof. Sampei Usui at Osaka University. Here I would like to thank him for encouraging me to study Hodge theory and for his help to understand it. I would like to thank Prof. Karl-Hermann Neeb for his interest and carefull reading of the present article.
## 2 $`M_m^n`$-functions
In this section we use the notations $`A=\left(\begin{array}{cc}a_A& b_A\\ c_A& d_A\end{array}\right)\mathrm{SL}(2,)`$ and
$$T=\left(\begin{array}{cc}1& 1\\ 0& 1\end{array}\right),Q=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),x=\left(\begin{array}{cc}x_1& x_2\\ x_3& x_4\end{array}\right),g=\left(\begin{array}{cc}k_1& k_2\\ 0& k_3\end{array}\right),x,g\mathrm{GL}(2,).$$
When there is no confusion we will simply write $`A=\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)`$. We denote by $``$ the Poincaré upper half plane and
$$\mathrm{j}(A,z):=c_Az+d_A$$
For $`A\mathrm{SL}(2,)`$ and $`m`$ we use the slash operator
$$f|_mA=(detA)^{m1}\mathrm{j}(A,z)^mf(Az)$$
For a ring $`R`$ we denote by $`\mathrm{Mat}_p(2,R)`$ the set of $`2\times 2`$-matrices in $`R`$ with the determinant $`p`$.
### 2.1 Definitions
In this section we define the notion of an $`M_m^n`$-function. For $`n=0`$ an $`M_m^0`$-function is a classical modular form of weight $`m`$ on $``$ (see bellow). For arbitrary $`n`$ we define it by induction: A holomorphic function $`f`$ on $``$ is called $`M_m^n`$ if the following three conditions are satisfied:
1. (13)
$$f|_mAf=\underset{i=1}{\overset{n}{}}\left(\genfrac{}{}{0pt}{}{n}{i}\right)c_A^i\mathrm{j}(A,z)^if_i,A\mathrm{SL}(2,)$$
where $`f_i`$ is an $`M_{m2i}^{ni}`$-function.
2. By induction we can write
(14)
$$f_i|_{m2i}Af_i=\underset{j=1}{\overset{ni}{}}\left(\genfrac{}{}{0pt}{}{ni}{j}\right)c_A^j\mathrm{j}(A,z)^jf_{ij}$$
where $`f_{ij}`$ is an $`M_{m2i2j}^{nij}`$-function. We assume that $`f_{ij}=f_{i1,j+1}=\mathrm{}=f_{1,i+j1}=f_{i+j}`$.
3. $`f`$ has finite growth when $`\mathrm{Im}(z)`$ tends to $`+\mathrm{}`$, i.e.
$$\underset{\mathrm{Im}(z)+\mathrm{}}{lim}f(z)=a_{\mathrm{}}<\mathrm{},a_{\mathrm{}}$$
The above definition can be made using a subgroup $`\mathrm{\Gamma }\mathrm{SL}(2,)`$. In this article we mainly deal with full differential modular forms, i.e. the case $`\mathrm{\Gamma }=\mathrm{SL}(2,)`$. We will also denote by $`M_m^n`$ the set of $`M_m^n`$-functions and we set
$$M:=_{n,m,n0}M_m^n$$
A classical modular form satisfies 1 and 2, where instead of the right hand side of (13) we write $`0`$. Note that for an $`M_m^n`$-function $`f`$ the associated functions $`f_i`$ are unique (consider the right hand side of (13) as a polynomial in $`c_A\mathrm{j}(A,z)^1`$ with coefficients $`\left(\genfrac{}{}{0pt}{}{n}{i}\right)f_i`$). The second condition guarantees that the consequent use of (13) for two matrices $`A,B\mathrm{SL}(2,)`$ leads to the same result. To see this, it is useful to define
(15)
$$f||_mA:=(detA)^{mn1}\underset{i=0}{\overset{n}{}}\left(\genfrac{}{}{0pt}{}{n}{i}\right)c_{A^1}^i\mathrm{j}(A,z)^{im}f_i(Az),A\mathrm{GL}(2,),fM_m^n$$
The factor $`detA`$ is introduced because of Hecke operators (see §2.3). This is not an action of $`\mathrm{GL}(2,)`$ on $`M_m^n`$ from the right. The equalities (13) and (14) are written in the form
(16)
$$f_i=f_i||_{m2i}A:=$$
$$(detA)^{mni1}\underset{j=0}{\overset{ni}{}}\left(\genfrac{}{}{0pt}{}{ni}{j}\right)c_{A^1}^j\mathrm{j}(A,z)^{j+2im}f_{i+j}(Az),i=0,1,2,\mathrm{},n,f_0:=f$$
for all $`A\mathrm{SL}(2,)`$ (We have substituted $`A^1z`$ for $`z`$ and then $`A^1`$ for $`A`$).
###### Lemma 1.
We have
$$f||_mA=f||_m(BA),A\mathrm{GL}(2,),B\mathrm{SL}(2,)$$
###### Proof.
The term $`(detA)^{nm+1}f||_mA(z)`$ is equal to:
$`=`$ $`{\displaystyle \underset{i=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n}{i}}\right)c_{A^1}^i\mathrm{j}(A,z)^{im}f_i(B^1BAz)`$
$`=`$ $`{\displaystyle \underset{i=0}{\overset{n}{}}}{\displaystyle \underset{j=0}{\overset{ni}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n}{i}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{ni}{j}}\right)c_{A^1}^ic_{B^1}^j\mathrm{j}(A,z)^{im}\mathrm{j}(B^1,BAz)^{m2ij}f_{i+j}(BAz)`$
$`=`$ $`{\displaystyle \underset{r=0}{\overset{n}{}}}{\displaystyle \underset{s=0}{\overset{r}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n}{s}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{ns}{rs}}\right)c_{A^1}^sc_{B^1}^{rs}\mathrm{j}(A,z)^{sm}\mathrm{j}(B^1,BAz)^{mrs}f_r(BAz)`$
$`=`$ $`{\displaystyle \underset{r=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n}{r}}\right)\mathrm{j}(BAz,z)^{rm}f_r(BAz)\mathrm{j}(A,z)^r({\displaystyle \underset{s=0}{\overset{r}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{r}{s}}\right)\mathrm{j}(BA,z)^sc_{A^1}^sc_{B^1}^{rs})`$
$`=`$ $`{\displaystyle \underset{r=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n}{r}}\right)\mathrm{j}(BA,z)^{rm}f_r(BAz)\mathrm{j}(A,z)^r(\mathrm{j}(BA,z)c_{A^1}+c_{B^1})^r`$
$`=`$ $`{\displaystyle \underset{r=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n}{r}}\right)\mathrm{j}(BA,z)^{rm}c_{(BA)^1}^rf_r(BAz)=(detA)^{nm+1}(f||_mBA)(z)`$
In the second line we have used (13). In the third line we have have changed the counting parameters: $`i+j=r,i=s,0sr`$. In the other lines we have used
$$\mathrm{j}(AB,z)=\mathrm{j}(A,Bz)\mathrm{j}(B,z)$$
and
$$\mathrm{j}(BA,z)c_A+det(A)c_B=c_{BA}\mathrm{j}(A,z),A,B\mathrm{GL}(2,)$$
Since $`f|_mT=f`$ we can write the Fourier expansion of $`f`$ at infinity
$$f=\underset{n=N}{\overset{+\mathrm{}}{}}a_nq^n,a_n,N=0,1,2,\mathrm{},\mathrm{},q=e^{2\pi iz}.$$
The third condition on $`f`$ implies that $`N=0`$.
### 2.2 Algebra of $`M_m^n`$-functions
In this section we use for $`A\mathrm{SL}(2,)`$ the simplification $`c=c_A`$. Recall the Eisenstein series (1) and
(17)
$$\mathrm{\Delta }(z):=(27g_3^2(z)g_2^3(z))=(\frac{2\pi i}{12})^3q\underset{n=1}{\overset{\mathrm{}}{}}(1q^n)^{24}=q24q^2+252q^3+\mathrm{}$$
$$j(z):=\frac{g_2^3(z)}{\mathrm{\Delta }(z)}=q^1+744+196884q+\mathrm{}$$
Note that $`\zeta (2)=\frac{\pi ^2}{6},\zeta (4)=\frac{\pi ^4}{90},\zeta (6)=\frac{\pi ^6}{945}`$ and so
(18)
$$p_{\mathrm{}}:=(a_1,a_2,a_3)=(\frac{2\pi i}{12},12(\frac{2\pi i}{12})^2,8(\frac{2\pi i}{12})^3)$$
where $`a_i`$’s are defined in (2). For $`k2`$ one can write
$$g_k(z)=s_k\underset{0(m,n)^2}{}\frac{1}{(n+mz)^k}M_k^0$$
where $`s_2=\frac{60}{(2\pi i)^2}`$ and $`s_3=\frac{140}{(2\pi i)^3}`$. Now $`g_1`$ satisfies
(19)
$$g_1_2Ag_1=c\mathrm{j}(A,z)^1,A\mathrm{SL}(2,)$$
and so $`g_1M_2^1`$ (see for instance p. 69). The following proposition describes the algebraic structure of $`M_m^n`$:
###### Proposition 1.
The followings are true:
1. For an $`fM_m^1`$ the function $`z(f(\frac{1}{z})f(z))`$ is in $`M_{m2}^0`$, i.e. it is a modular form of weight $`m2`$.
2. $`M_2^1`$ is a one dimensional $``$-vector space generated by $`g_1`$.
3. If $`nn^{}`$ then $`M_m^nM_m^{}^n^{}`$ and
$$M_m^nM_m^{}^n^{}M_{m+m^{}}^{n+n^{}},M_m^n+M_m^n^{}=M_m^n^{}$$
4. For a modular form $`f`$ of weight $`m`$ we have $`f(g_1)^nM_{2n+m}^n`$.
###### Proof.
The first item is a direct consequence of the definition applied to $`A=Q`$. The only modular forms of weight $`0`$ are constant functions. This and 1 imply that $`M_2^1`$ is one dimensional. Item 3 is derived from the definition. Item 4 is a consequence of Item 3. Item 4 was the main idea behind the definition of $`M_m^n`$. ∎
The following proposition shows that $`M`$ is in fact a differential algebra.
###### Proposition 2.
For $`fM_m^n`$ we have $`\frac{df}{dz}M_{m+2}^{n+1}`$ and
(20)
$$\frac{d(f||_mA)}{dz}=\frac{df}{dz}||_{m+2}A,A\mathrm{GL}(2,)$$
###### Proof.
For $`A\mathrm{Mat}_p(2,)`$ the term $`\frac{d(f||_mA)}{dz}`$ is equal to:
$`=`$ $`p^{mn}\left({\displaystyle \underset{i=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n}{i}}\right)c_{A^1}^i((mi)c_{A^1}\mathrm{j}(A,z)^{i1m}f_i(Az)+\mathrm{j}(A,z)^{im2}{\displaystyle \frac{df_i}{dz}}(Az))\right)`$
$`=`$ $`p^{mn}({\displaystyle \underset{i=1}{\overset{n+1}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n}{i1}}\right)c_{A^1}^i\mathrm{j}(A,z)^{i2m}(mi+1)f_{i1}(Az)`$
$`+`$ $`{\displaystyle \underset{i=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n}{i}}\right)c_{A^1}^i\mathrm{j}(A,z)^{i2m}{\displaystyle \frac{df_i}{dz}}(Az))`$
$`=`$ $`p^{mn}\left({\displaystyle \underset{i=0}{\overset{n+1}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n+1}{i}}\right)c_{A^1}^i\mathrm{j}(A,z)^{i2m}\stackrel{~}{f}_i(Az)\right)`$
where
$$\stackrel{~}{f}_i=\frac{i(mi+1)}{n+1}f_{i1}+\frac{n+1i}{n+1}\frac{df_i}{dz},i=0,1,\mathrm{},n+1,f_1=f_{n+1}:=0$$
Now $`\frac{df}{dz}`$ is a an $`M_{m+2}^{n+1}`$-function with the associated $`M_{m2i}^{ni}`$ function $`\stackrel{~}{f}_i`$ for $`i=0,1,2,\mathrm{},n+1`$. The first item of definition has been checked above, using the fact that $`A\mathrm{SL}(2,)`$ and $`f||_mA=f`$. We calculate $`\stackrel{~}{f}_{ij}`$ as above and then we get
$`\stackrel{~}{f}_{ij}`$ $`=`$ $`{\displaystyle \frac{i(mi+1)}{n+1}}f_{i1+j}+{\displaystyle \frac{n+1i}{n+1}}({\displaystyle \frac{j(m2ij+1)}{ni+1}}f_{i+j1}+{\displaystyle \frac{ni+1j}{ni+1}}{\displaystyle \frac{df_{i+j}}{dz}})`$
$`=`$ $`{\displaystyle \frac{(i+j)(m(i+j)+1)}{n+1}}f_{i+j1}+{\displaystyle \frac{n+1ij}{n+1}}{\displaystyle \frac{df_{i+j}}{dz}}=\stackrel{~}{f}_{i+j}`$
This proves Item 2 of the definition of an $`M_{m+2}^{n+1}`$-function. The third item can be checked using
$$\frac{df}{dz}=2\pi iq\frac{df}{dq}$$
We consider $`M`$ as a graded algebra with $`\mathrm{deg}(f)=m,fM_m^n`$. The nontrivial statement about $`M_m^n`$-functions is Theorem 1 in the Introduction. It can be interpreted also as the equality of graded algebras $`M=[g_1,g_2,g_3]`$. It implies that each $`fM_m^n`$ can be written as
$$f=\underset{i=0}{\overset{n}{}}f_ig_1^i,f_iM_{m2i}^0$$
We prove Theorem 1 in §4.2.
The relations between the $`g_i,i=1,2,3`$ and their derivatives are given by Ramanujan’s equalities:
(21)
$$\frac{dg_1}{dz}=g_1^2\frac{1}{12}g_2,\frac{dg_2}{dz}=4g_1g_26g_3,\frac{dg_3}{dz}=6g_1g_3\frac{1}{3}g_2^2$$
(see for instance ). The proof of Theorem 1 will contain a new proof of these equalities.
### 2.3 Hecke operators
For $`p`$ let $`\mathrm{SL}(2,)\backslash \mathrm{Mat}_p(2,)=\{[A_1],[A_2],\mathrm{},[A_s]\}`$. We define the $`p`$-th Hecke operator in the following way
$$T_pf=\underset{k=1}{\overset{s}{}}f||_mA_k,fM_m^n$$
Lemma 1 implies that the above definition does not depend on the choice of $`A_k`$ in the class $`[A_k]`$. Form Proposition 2 one can deduce that the differential operator $`\frac{d}{dz}`$ commutes with the Hecke operator $`T_p`$.
###### Proposition 3.
$`T_p`$ defines a map from $`M_m^n`$ to itself.
This will be proved in §2.5. One can take
$$\stackrel{~}{T}_p:=\underset{dp,0bd1}{}\left(\begin{array}{cc}\frac{p}{d}& b\\ 0& d\end{array}\right)[\mathrm{Mat}_p(2,)]$$
and since for matrices $`\left(\begin{array}{cc}a& b\\ 0& d\end{array}\right)`$ the slash operator $`|_m`$ is $`p^n`$ times $`||_m`$ we have $`T_pf=p^nf|_m\stackrel{~}{T}_p`$ and we get the expression (3) in the Introduction. In a similar way to the case of modular forms (see §6) one can check that
$$T_pT_q=\underset{d(p,q)}{}d^{mn1}T_{\frac{pq}{d^2}}$$
### 2.4 The period domain
The group $`\mathrm{SL}(2,)`$ acts from left on the period domain $`𝒫`$ defined in (10) and $`G_0`$ in (11) acts from right. We consider holomorphic functions on
$$:=\mathrm{SL}(2,)\backslash 𝒫$$
as holomorphic functions
$$f:𝒫,\text{ holomorphic satisfying },f(Az)=f(z),A\mathrm{SL}(2,),z𝒫$$
The determinant function is such a function. The Poincaré upper half plane $``$ is embedded in $`𝒫`$ in the following way:
$$z\stackrel{~}{z}=\left(\begin{array}{cc}z& 1\\ 1& 0\end{array}\right)$$
We denote by $`\stackrel{~}{}`$ the image of $``$ under this map. For a function $`f`$ on $``$ we denote by $`\stackrel{~}{f}`$ the corresponding function on $`\stackrel{~}{}`$.
###### Proposition 4.
There is a unique map
$$\varphi :M𝒪(𝒫),f\varphi (f)=F$$
of the algebra of $`M`$-functions into the algebra of holomorphic functions on $`𝒫`$ such that
1. For all $`fM`$ the restriction of $`F`$ to $`\stackrel{~}{}`$ is equal to $`\stackrel{~}{f}`$.
2. For all $`fM`$ the holomorphic function $`F`$ is $`\mathrm{SL}(2,)`$ invariant.
3. We have
(22)
$$F(xg)=k_2^nk_1^{nm}\underset{i=0}{\overset{n}{}}\left(\genfrac{}{}{0pt}{}{n}{i}\right)k_3^ik_2^iF_i(x),x𝒫,gG_0,$$
where $`F_i=\varphi (f_i)`$.
Conversely, every holomorphic function $`F`$ on $`𝒫`$ which satisfies 2,3 and whose restriction to $`\stackrel{~}{}`$ has finite growth at infinity is of the form $`F=\varphi (f)`$ for some $`fM_m^n`$.
We denote by $`\stackrel{ˇ}{M_m^n}`$ the set of holomorphic functions on $`𝒫`$ which restricted to $`\stackrel{~}{}`$ have finite growth at infinity, are $`\mathrm{SL}(2,)`$ invariant and satisfy (22). For a classical modular form $`f:`$ of weight $`m`$ the associated $`F=\varphi (f)`$ is
$$F(x)=x_3^mf(\frac{x_1}{x_3})\stackrel{ˇ}{M_m^0}$$
We also have
$$det\stackrel{ˇ}{M_0^1}$$
###### Proof.
We have
(23)
$$\left(\begin{array}{cc}x_1& x_2\\ x_3& x_4\end{array}\right)=\left(\begin{array}{cc}\frac{x_1}{x_3}& 1\\ 1& 0\end{array}\right)\left(\begin{array}{cc}x_3& x_4\\ 0& \frac{det(x)}{x_3}\end{array}\right)$$
So we expect $`F`$ to be defined by
$$F(x)=F\left(\left(\begin{array}{cc}\frac{x_1}{x_3}& 1\\ 1& 0\end{array}\right)\left(\begin{array}{cc}x_3& x_4\\ 0& \frac{det(x)}{x_3}\end{array}\right)\right):=x_3^mdet(x)^n\underset{i=0}{\overset{n}{}}\left(\genfrac{}{}{0pt}{}{n}{i}\right)x_4^ix_3^idet(x)^if_i(\frac{x_1}{x_3})$$
This $`F`$ restricted to $`\stackrel{~}{}`$ is $`f`$. By definition of $`F`$ one can rewrite (13) in the form
$$f(A\frac{x_1}{x_3})=(cx_1+dx_3)^{mn}F\left(\begin{array}{cc}x_1& d\\ x_3& c\end{array}\right)$$
where $`A=\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)`$. We check item 3: Let
$$g^{}=\left(\begin{array}{cc}k_1^{}& k_3^{}\\ 0& k_2^{}\end{array}\right):=\left(\begin{array}{cc}x_3& x_4\\ 0& \frac{det(x)}{x_3}\end{array}\right)$$
RHS of (22) $`=`$ $`k_2^nk_1^{nm}{\displaystyle \underset{i=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n}{i}}\right)k_3^ik_2^iF_i(x)`$
$`=`$ $`(k_2k_2^{})^n(k_1k_1^{})^{nm}{\displaystyle \underset{i=0}{\overset{n}{}}}{\displaystyle \underset{j=0}{\overset{ni}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n}{i}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{ni}{j}}\right)k_3^ik_2^ik_2^ik_1^ik_3^jk_2^jf_{ij}({\displaystyle \frac{x_1}{x_3}})`$
$`=`$ $`(k_2k_2^{})^n(k_1k_1^{})^{nm}{\displaystyle \underset{r=0}{\overset{n}{}}}{\displaystyle \underset{s=0}{\overset{r}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n}{s}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{ns}{rs}}\right)k_3^sk_2^sk_2^sk_1^sk_3^{rs}k_2^{r+s}f_{s,rs}({\displaystyle \frac{x_1}{x_3}})`$
$`=`$ $`(k_2k_2^{})^n(k_1k_1^{})^{nm}{\displaystyle \underset{r=0}{\overset{n}{}}}{\displaystyle \underset{s=0}{\overset{r}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n}{r}}\right)(k_2k_3^{}+k_3k_1^{})^r(k_2k_2^{})^rf_r({\displaystyle \frac{x_1}{x_3}})`$
$`=`$ $`F(\left(\begin{array}{cc}z& 1\\ 1& 0\end{array}\right)g^{}g)=F(xg)`$
For all the equalities above we have used the same reasoning as in the proof of Lemma 1. We check that $`F(Ax)=F(x),A\mathrm{SL}(2,)`$: The term $`F(Ax)`$ is equal to
$`=`$ $`(cx_1+dx_3)^mdet(x)^n{\displaystyle \underset{i=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n}{i}}\right)(cx_2+dx_4)^i(cx_1+dx_3)^idet(x)^if_i(A{\displaystyle \frac{x_1}{x_3}})`$
$`=`$ $`(cx_1+dx_3)^mdet(x)^n`$
. $`{\displaystyle \underset{i=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n}{i}}\right)(cx_2+dx_4)^i(cx_1+dx_3)^idet(x)^i(cx_1+dx_3)^{m2i(ni)}F_i\left(\begin{array}{cc}x_1& d\\ x_3& c\end{array}\right)`$
$`=`$ $`F\left(\left(\begin{array}{cc}x_1& d\\ x_3& c\end{array}\right)\left(\begin{array}{cc}1& \frac{cx_2+dx_4}{cx_1+dx_3}\\ 0& \frac{det(x)}{cx_1+dx_3}\end{array}\right)\right)=F(x)`$
Now let $`F`$ satisfy 2,3 and its restriction to $`\stackrel{~}{}`$ has finite growth at infinity put $`f=F_\stackrel{~}{}`$. First we note that
$$\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)\left(\begin{array}{cc}z& 1\\ 1& 0\end{array}\right)=\left(\begin{array}{cc}Az& 1\\ 1& 0\end{array}\right)\left(\begin{array}{cc}\mathrm{j}(A,z)& c\\ 0& \mathrm{j}(A,z)^1det(A)\end{array}\right),A\mathrm{GL}(2,)$$
Now
$`f(Az)`$ $`=`$ $`F\left(\begin{array}{cc}Az& 1\\ 1& 0\end{array}\right)=F(\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)\left(\begin{array}{cc}z& 1\\ 1& 0\end{array}\right)\left(\begin{array}{cc}det(A)\mathrm{j}(A,z)^1& c\\ 0& \mathrm{j}(A,z)\end{array}\right))`$
$`=`$ $`\mathrm{j}(A,z)^n\mathrm{j}(A,z)^{mn}{\displaystyle \underset{i=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n}{i}}\right)c^i\mathrm{j}(A,z)^if_i(x)={\displaystyle \underset{i=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n}{i}}\right)c^i\mathrm{j}(A,z)^{mi}f_i(x)`$
This finishes the proof of our proposition. ∎
### 2.5 Proof of Proposition 3
We define
$`T_p:\stackrel{ˇ}{M_m^n}\stackrel{ˇ}{M_m^n},\stackrel{ˇ}{T}_pF(x)=p^{m2n1}{\displaystyle \underset{k=1}{\overset{s}{}}}F(A_ix).`$
This function has trivially its image in $`\stackrel{ˇ}{M_m^n}`$. We calculate the corresponding function in $`M_m^n`$: The term $`T_pf(z)`$ equals
$`=`$ $`p^{m2n1}{\displaystyle \underset{k=1}{\overset{s}{}}}F(A_k\left(\begin{array}{cc}z& 1\\ 1& 0\end{array}\right))`$
$`=`$ $`p^{m2n1}{\displaystyle \underset{k=1}{\overset{s}{}}}F\left(\left(\begin{array}{cc}A_kz& 1\\ 1& 0\end{array}\right)\left(\begin{array}{cc}\mathrm{j}(A_k,z)& c\\ 0& p.\mathrm{j}(A_k^1,A_kz)\end{array}\right)\right)`$
$`=`$ $`p^{m2n1}{\displaystyle \underset{k=1}{\overset{s}{}}}(p.\mathrm{j}(A_k^1,A_kz))^n(\mathrm{j}(A_k,z))^{nm}{\displaystyle \underset{i=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n}{i}}\right)(cp^1)^i\mathrm{j}(A_k^1,A_kz)^if_i(A_kz)`$
$`=`$ $`{\displaystyle \underset{k=1}{\overset{s}{}}}f||_mA_k`$
where $`A_k=\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)`$. This proves Proposition 3.
### 2.6 Some non-holomorphic functions on the period domain
On the complex manifold $`𝒫`$ we have the following $`\mathrm{SL}(2,)`$ invariant analytic functions:
$$B_1:=\mathrm{Im}(x_1\overline{x_3}),B_2:=\mathrm{Im}(x_2\overline{x_4}),B_3:=x_1\overline{x_4}x_2\overline{x_3}$$
They define analytic functions on $``$ which we denote them by the same letter. They satisfy
$$B_1_\stackrel{~}{}(z)=\mathrm{Im}(z),B_1(xg)=B_1(x)|k_1|^2$$
$$B_2_\stackrel{~}{}(z)=0,B_2(xg)=B_1(x)|k_3|^2+B_2(x)|k_2|^2+\mathrm{Im}(B_3(x)k_3\overline{k_2})$$
$$B_3_\stackrel{~}{}(z)=1,B_3(xg)=B_3(x)k_1\overline{k_2}+2\sqrt{1}k_1\overline{k_3}B_1(x).$$
By equality (23) one can easily see that every point in $`𝒫`$ can be mapped to a point of $`\stackrel{~}{}`$ by an action of a unique element of $`G_0`$. This implies that $`\mathrm{SL}(2,)`$ invariant $`B_i,i=1,2,3`$, with the above properties are unique.
### 2.7 Other topics
The literature of modular forms and its applications in number theory is huge. The first question which naturally arises at this point is as follows: Which part of the theory of modular forms can be generalized to the context of differential modular forms and which arithmetic properties can one expect to find? Since I am not expert in this area, I just mention some subjects which could fit well into this section.
The first of these, is the Eichler-Manin-Shimuara theory of periods for cusp forms (see and its references). Note that the notion “period” in this theory, as far as I know, has nothing to do with the notion of a period in this article. The notion of period appears there because classical modular forms can be intrepreted as sections of a tensor product of the cotangent bundle of a moduli curve and hence a differential multi form, which can be integrated over some path in the moduli curve (see ). The differential modular forms are no longer interpreted as sections of line bundles and this makes the situation more difficult. Lewis type equations attached to differential modular forms will be also of intereset (see ).
Another theory which could be developed for differential modular forms is Atkin-Lehner theory of old and new modular forms (see ). This seems to me to be a quite accessible theory. The $`L`$-functions attached to differential modular forms througth their Fourier expansion may also be of interest (see ).
## 3 Families of elliptic curves and the Gauss-Manin connection
First, let us fix some notation. For a ring $`R`$ we denote by $`R[t]`$ the polynomial ring with coefficients in $`R`$ and the variable $`t:=(t_1,t_2,\mathrm{},t_s)`$. We also define the affine space $`𝔸_R^s=\mathrm{Spec}(R[t])`$ defined over $`R`$. For simplicity we write $`𝔸^s=𝔸_{}^s`$. The set of relative differential $`i`$-forms in $`𝔸_R^s`$ is denoted by $`\mathrm{\Omega }_{𝔸_R^s/\mathrm{Spec}(R)}^i`$.
### 3.1 The family
We consider the following family of elliptic curves
(24)
$$_t:f=0$$
$$f:=y^24T_0(t)x^3+T_1(t)x^2+T_2(t)x+T_3(t)+T_4(t)y+T_5(t)xy,t=(t_1,t_2,\mathrm{},t_s)T$$
where
$$T=𝔸^s\backslash \{\mathrm{\Delta }=0\},T_i[t],i=0,1,\mathrm{},5$$
and
$$\mathrm{\Delta }=(6912T_3^23456T_3T_4^2+432T_4^4)T_0^4+(1152T_1T_2T_3288T_1T_2T_4^2576T_1T_3T_4T_5+$$
$$144T_1T_4^3T_5256T_2^3+384T_2^2T_4T_5288T_2T_3T_5^2120T_2T_4^2T_5^2+144T_3T_4T_5^34T_4^3T_5^3)T_0^3+(64T_1^3T_316T_1^3T_4^216T_1^2T_2^2+16T_1^2T_2T_4T_5$$
$$48T_1^2T_3T_5^2+8T_1^2T_4^2T_5^2+8T_1T_2^2T_5^28T_1T_2T_4T_5^3+12T_1T_3T_5^4T_1T_4^2T_5^4T_2^2T_5^4+T_2T_4T_5^5T_3T_5^6)T_0^2$$
The polynomial $`\mathrm{\Delta }[t]`$ is the locus of parameters in which $`_t`$ is singular. The reader is referred to for the algorithms which calculate $`\mathrm{\Delta }`$ (see also the proof of Proposition 6). To make our notation simpler, we define
$$R:=[t,\frac{1}{T_0}],𝕌_0:=\mathrm{Spec}(R),𝕌_1:=𝔸^2\times (𝔸^s\backslash \{T_0=0\})=\mathrm{Spec}([x,y,t,\frac{1}{T_0}])$$
### 3.2 De Rham cohomology
Let $`\pi :𝕌_1𝕌_0`$ be the projection on the last $`s`$-coordinates. The following quotient
$$H=\frac{\mathrm{\Omega }_{𝕌_1}^1}{f\mathrm{\Omega }_{𝕌_1}^1+df\mathrm{\Omega }_{𝕌_1}^0+\pi ^1\mathrm{\Omega }_{𝕌_0}^1\mathrm{\Omega }_{𝕌_1}^0+d\mathrm{\Omega }_{𝕌_1}^0}\frac{\mathrm{\Omega }_{𝕌_1/𝕌_0}^1}{f\mathrm{\Omega }_{𝕌_1/𝕌_0}^1+df\mathrm{\Omega }_{𝕌_1/𝕌_0}^0+d\mathrm{\Omega }_{𝕌_1/𝕌_0}^0}$$
is an $`R`$-module and will play the role of de Rham cohomology for us. One may call $`H`$ the Brieskorn module associated to the family $`_t`$ in analogy to the local modules introduced by Brieskorn in 1970. The restriction of $`H`$ to the elliptic curve $`_t`$ gives us $`H_{\mathrm{dR}}^1(_t)`$. Set
(25)
$$\omega _1:=\frac{2}{5}(2xdy3ydx),\text{ and }\omega _2:=\frac{2}{7}x(2xdy3ydx).$$
###### Proposition 5.
The $`R`$-module $`H`$ is free and $`\omega _1`$ and $`\omega _2`$ form a basis of $`H`$.
###### Proof.
We consider the classical Brieskorn module
$$\stackrel{~}{H}=\frac{\mathrm{\Omega }_{𝕌_1/𝕌_0}^1}{df\mathrm{\Omega }_{𝕌_1/𝕌_0}^0+d\mathrm{\Omega }_{𝕌_1/𝕌_0}^0}$$
It is a $`R[f]`$-module. It is proved in §6 Proposition 1 that $`\stackrel{~}{H}`$ is freely generated by $`\omega _1,\omega _2`$ as $`R[f]`$-module. Considering the canonical map $`\stackrel{~}{H}H`$ we obtain the assertion of our proposition. ∎
### 3.3 Gauss-Manin connection
Each element of the $`R`$-module $`H`$ can be interpreted as a global section of the first cohomology bundle of the family $`_t`$ over $`T`$. Since the Gauss-Manin connection is a connection in the cohomology bundle, it is natural to find an algebraic definition for it using the $`R`$\- module structure of $`H`$. In this section we do this. Define
$$V:=\frac{\mathrm{\Omega }_{𝕌_1}^2}{df\mathrm{\Omega }_{𝕌_1}^1+f\mathrm{\Omega }_{𝕌_1}^2+\pi ^1\mathrm{\Omega }_{𝕌_0}^1\mathrm{\Omega }_{𝕌_1}^1}\frac{\mathrm{\Omega }_{𝕌_1/𝕌_0}^2}{df\mathrm{\Omega }_{𝕌_1/𝕌_0}^1+f\mathrm{\Omega }_{𝕌_1/𝕌_0}^2}$$
###### Proposition 6.
The polynomial $`\mathrm{\Delta }`$ is a zero divisor of the $`R`$-module $`V`$, i.e.
$$\mathrm{\Delta }V=0$$
###### Proof.
We define
$$\stackrel{~}{V}:=\frac{\mathrm{\Omega }_{𝕌_1/𝕌_0}^2}{df\mathrm{\Omega }_{𝕌_1/𝕌_0}^1}$$
and consider it as an $`R[f]`$-module. It is proved in Lemma 4 that $`B=\{dxdy,xdxdy\}`$ form a basis of $`\stackrel{~}{V}`$. Let $`A`$ be the matrix of multiplication by $`f`$ $`R`$-linear map in $`\stackrel{~}{V}`$ in the basis $`B`$. Then $`S(t,f):=det(Af.I_2)`$ has the property $`S(t,f)\stackrel{~}{V}=0`$, where $`I_2`$ is the identity two by two matrix. This implies that $`S(t,0)V=0`$. Now $`S(0,t)=6912.\mathrm{\Delta }`$ (In fact this is the way we have calculated $`\mathrm{\Delta }`$). ∎
We have a well-defined differential map
$$d:HV$$
and we define the Gauss-Manin connection $`H`$ as follows:
$$:H\mathrm{\Omega }_T^1_RH$$
$$\omega =\frac{1}{\mathrm{\Delta }}\underset{i}{}\alpha _i\beta _i,\text{ where }\mathrm{\Delta }d\omega =\underset{i}{}\alpha _i\beta _i,\alpha _i\mathrm{\Omega }_{𝔸^s}^1,\beta _i\mathrm{\Omega }_{𝕌_1}^1$$
Let $`U`$ be an small open set in $`U`$ and $`\{\delta _t\}_{tU},\delta _tH_1(_t,)`$ be a continuous family of topological one dimensional cycles. The main property of the Gauss-Manin connection is
(26)
$$d(_{\delta _t}\omega )=\alpha _i_{\delta _t}\beta _i,\omega =\underset{i}{}\alpha _i\beta _i,\alpha _i\mathrm{\Omega }_T^1,\beta _iH.$$
Let $`\omega _1,\omega _2`$ be a basis of $`H`$ and define $`\omega =(\omega _1,\omega _2)^t`$. The Gauss-Manin connection in this basis can be written in the following way:
(27)
$$\omega =A\omega ,A=\frac{1}{\mathrm{\Delta }}(\underset{i=1}{\overset{s}{}}A_idt_i)\mathrm{Mat}(2,\mathrm{\Omega }_T^1),A_i\mathrm{Mat}(2,[t]).$$
###### Proposition 7.
Let $`\stackrel{~}{\omega }=S\omega `$ be another basis of $`H`$ and $`\omega =A\omega `$. Then
$$(\stackrel{~}{\omega })=S(S^1dS+A)S^1\stackrel{~}{\omega }$$
###### Proof.
We have
$`(\stackrel{~}{\omega })`$ $`=`$ $`(S\omega )=dS\omega +S\omega =dS.S^1\stackrel{~}{\omega }+SAS^1\stackrel{~}{\omega }`$
$`=`$ $`(dS.S^1+SAS^1)\stackrel{~}{\omega }`$
### 3.4 Classical differential forms
Recall the canonical basis (25) of $`H`$. Assume that in the family (24) $`T_3`$ is of the form $`\stackrel{~}{T}_3+t_3`$, where $`\stackrel{~}{T}_3`$ and all other $`T_i`$’s do not depend on $`t_3`$. Let $`A`$ be the matrix of the Gauss-Manin connection in the basis $`\omega `$ as in (27). We are particularly interested in $`A_3`$ in this case; it satisfies
(28)
$$det(A_3)=321052999680T_0^4\mathrm{\Delta }$$
Define
$$\stackrel{~}{\omega }=\frac{1}{\mathrm{\Delta }}A_3\omega $$
The classical way of calculating Gauss-Manin connection leads to the equalities
(29)
$$\stackrel{~}{\omega }_1=\frac{d\omega _1}{df}=\frac{dx}{y},\stackrel{~}{\omega }_2=\frac{d\omega _2}{df}=\frac{xdx}{y}$$
restricted to the elliptic curves $`_t`$ and up to exact differential forms. Equality (28) implies that $`\stackrel{~}{\omega }`$ form a basis of $`H_\mathrm{\Delta }`$, the localization of $`H`$ over $`\{1,\mathrm{\Delta },\mathrm{\Delta }^2,\mathrm{}\}`$.
### 3.5 An example
The following example plays a basic role in the proof of Theorem 1:
(30)
$$_t:f:=y^24t_0(xt_1)^3+t_2(xt_1)+t_3=0$$
In this example
$$\mathrm{\Delta }=t_0(27t_0t_3^2t_2^3)$$
The calculation of the Gauss-Manin connection with respect to the canonical basis (25) leads to:
$$A_0=\left(\begin{array}{cc}21/2t_0t_1t_2t_39t_0t_3^2+3/4t_2^3& 21/2t_0t_2t_3\\ 21/2t_0t_1^2t_2t_3+9t_0t_1t_3^21/2t_1t_2^35/8t_2^2t_3& 21/2t_0t_1t_2t_318t_0t_3^2+5/4t_2^3\end{array}\right)$$
$$A_1=\left(\begin{array}{cc}0& 0\\ 27t_0^2t_3^2t_0t_2^3& 0\end{array}\right)$$
$$A_2=\left(\begin{array}{cc}63/2t_0^2t_1t_35/4t_0t_2^2& 63/2t_0^2t_3\\ 63/2t_0^2t_1^2t_3+1/2t_0t_1t_2^2+15/8t_0t_2t_3& 63/2t_0^2t_1t_37/4t_0t_2^2\end{array}\right)$$
$$A_3=\left(\begin{array}{cc}21t_0^2t_1t_2+45/2t_0^2t_3& 21t_0^2t_2\\ 21t_0^2t_1^2t_29t_0^2t_1t_35/4t_0t_2^2& 21t_0^2t_1t_2+63/2t_0^2t_3\end{array}\right)$$
Now the same with respect to the classical basis $`\stackrel{~}{\omega }`$ of $`H_\mathrm{\Delta }`$ is given by:
(31)
$$A_0=\left(\begin{array}{cc}3/2t_0t_1t_2t_39t_0t_3^2+1/4t_2^3& 3/2t_0t_2t_3\\ 3/2t_0t_1^2t_2t_3+9t_0t_1t_3^21/2t_1t_2^3+1/8t_2^2t_3& 3/2t_0t_1t_2t_318t_0t_3^2+3/4t_2^3\end{array}\right)$$
$$A_1=\left(\begin{array}{cc}0& 0\\ 27t_0^2t_3^2t_0t_2^3& 0\end{array}\right)$$
$$A_2=\left(\begin{array}{cc}9/2t_0^2t_1t_3+1/4t_0t_2^2& 9/2t_0^2t_3\\ 9/2t_0^2t_1^2t_3+1/2t_0t_1t_2^23/8t_0t_2t_3& 9/2t_0^2t_1t_31/4t_0t_2^2\end{array}\right)$$
$$A_3=\left(\begin{array}{cc}3t_0^2t_1t_29/2t_0^2t_3& 3t_0^2t_2\\ 3t_0^2t_1^2t_29t_0^2t_1t_3+1/4t_0t_2^2& 3t_0^2t_1t_2+9/2t_0^2t_3\end{array}\right)$$
See for the procedures which calculate all matrices above. In the library brho.lib the command gaussmaninp calculate the above matrix. Note that the canonical basis of the Brieskorn module in this library is $`\left(\begin{array}{cc}0& \frac{12}{5}\\ \frac{12}{7}& 0\end{array}\right)\omega `$.
## 4 The period map
### 4.1 Derivation of the period map
Let $`\omega =(\omega _1,\omega _2)^𝗍`$ be a basis of $`H_\mathrm{\Delta }`$. The period map associated to the basis $`\omega `$ is given by:
$$\mathrm{𝗉𝗆}:T\mathrm{SL}(2,)\backslash \mathrm{GL}(2,),t\left[\frac{1}{\sqrt{2\pi i}}\left(\begin{array}{cc}_{\delta _1}\omega _1& _{\delta _1}\omega _2\\ _{\delta _2}\omega _1& _{\delta _2}\omega _2\end{array}\right)\right]$$
It is well-defined and holomorphic. Here $`\sqrt{i}=e^{\frac{2\pi i}{4}}`$ and $`(\delta _1,\delta _2)`$ is a basis of the $``$-module $`H_1(_t,)`$ such that the intersection matrix in this basis is $`\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$.
Let $`\stackrel{~}{\omega }=S\omega ,S\mathrm{Mat}(2,[t,\frac{1}{T_0},\frac{1}{\mathrm{\Delta }}])`$ be another basis of $`H`$ and $`\stackrel{~}{\mathrm{𝗉𝗆}}`$ be the associated period map. Then it is easy to see that
(32)
$$\stackrel{~}{𝗉}𝗆=\mathrm{𝗉𝗆}S^𝗍$$
###### Proposition 8.
Let $`\omega `$ be a basis of $`H_\mathrm{\Delta }`$ with
(33)
$$\omega =A\omega ,A\mathrm{Mat}(2,\mathrm{\Omega }_T^1)$$
Then
(34)
$$d(\mathrm{𝗉𝗆})(t)=\mathrm{𝗉𝗆}(t)A^𝗍,tT,$$
where $`d`$ is the differential map.
###### Proof.
Let $`_i:=_{\frac{}{t_i}}`$ denote the Gauss-Manin connection with respect to the parameter $`t_i`$. Then $`_i\omega =A_i\omega `$ and according to (26) and (32)
$$\frac{}{t_i}\mathrm{𝗉𝗆}(t)=\left(\begin{array}{cc}_{\delta _1}_i\stackrel{~}{\omega }_1& _{\delta _1}_i\stackrel{~}{\omega }_2\\ _{\delta _2}_i\stackrel{~}{\omega }_1& _{\delta _2}_i\stackrel{~}{\omega }_2\end{array}\right)=\mathrm{𝗉𝗆}(t)A_i^𝗍$$
This proves the desired statement. ∎
For the basis $`\omega =(\omega _1,\omega _2)^𝗍`$ if $`\omega _1`$ is such that its restriction to each elliptic curve $`_t,tT`$ is of the first kind then the period map is defined from $`T`$ into $`:=\mathrm{SL}(2,)\backslash 𝒫`$. For instance, the period map associated to the classical basis (29) of $`H_\mathrm{\Delta }`$ has this property.
### 4.2 The Action of an algebraic group
We consider the family of elliptic curves (30). It can be checked easily that (1) is an action of $`G_0`$ on $`𝔸^4`$. (This can be also verified from the proof of the proposition bellow). It is also easy to verify that $`𝔸^4/G_0`$ is isomorphic to $`^1`$ through the map
(35)
$$s:𝔸^4/G_0^1,t[t_2^3:27t_0t_3^2t_2^3]$$
and so
(36)
$$j(t):=\frac{t_2^3}{27t_0t_3^2t_2^3}$$
is $`G_0`$-invariant and gives an isomorphy between $`T/G_0`$ and $`𝔸`$. Recall the basis $`\stackrel{~}{\omega }`$ of $`H_\mathrm{\Delta }`$ in (29).
###### Proposition 9.
The period $`\mathrm{𝗉𝗆}`$ associated to the basis $`\stackrel{~}{\omega }`$ is a biholomorphism and
(37)
$$\mathrm{𝗉𝗆}(tg)=\mathrm{𝗉𝗆}(t)g,t𝔸^4,gG_0$$
###### Proof.
We first prove (37). Let
$$\alpha :𝔸^2𝔸^2,(x,y)(k_2^1k_1xk_3k_2^1,k_2^1k_1^2y)$$
Then
$$k_2^2k_1^4\alpha ^1(f)=y^24t_0k_2^2k_1^4(k_2^1k_1xk_3k_2^1t_1)^3+t_2k_2^2k_1^4(k_2^1k_1xk_3k_2^1t_1)+t_3k_2^2k_1^4$$
$$y^24t_0k_1^1k_2^1(x(t_1k_2k_1^1+k_3k_1^1))^3+t_2k_1^3k_2(x(t_1k_2k_1^1+k_3k_1^1))+t_3k_1^4k_2^2$$
This implies that $`\alpha `$ induces an isomorphism of elliptic curves
$$\alpha :_{tg}_t$$
Now
$$\alpha ^1\stackrel{~}{\omega }=\left(\begin{array}{cc}k_1^1& 0\\ k_3k_2^1k_1^1& k_2^1\end{array}\right)\stackrel{~}{\omega }=\left(\begin{array}{cc}k_1& 0\\ k_3& k_2\end{array}\right)^1\stackrel{~}{\omega }$$
By the equality (32) we have
$$\mathrm{𝗉𝗆}(t)=\mathrm{𝗉𝗆}(tg).g^1$$
which proves (37).
The matrix of the Gauss-Manin connection in the basis $`\stackrel{~}{\omega }`$ for the family (30) is calculated in §3.5. Let $`B`$ be a $`4\times 4`$ matrix and the $`i`$-th row of $`B`$ constitutes of the first and second rows of $`A_i`$. Then
$$det(B)=\frac{3}{4}t_0\mathrm{\Delta }^3$$
shows that the period map $`\mathrm{𝗉𝗆}`$ is regular at each point $`tT`$ and hence it is locally a biholomorphism.
The period map $`\mathrm{𝗉𝗆}`$ induces a local biholomorphic map $`\overline{𝗉}𝗆:T/G_0\mathrm{SL}(2,)\backslash `$ and so we have the local biholomorphism $`\overline{𝗉}𝗆j^1:𝔸𝔸`$. One can compactify $`\mathrm{SL}(2,)\backslash `$ by adding the cusp $`\mathrm{SL}(2,)/=\{c\}`$ (see ) and the map $`\overline{𝗉}𝗆j^1`$ is continuous at $`v`$ sending $`v`$ to $`c`$, where $`v`$ is the point induced by $`t_027t_3^2t_2^3=0`$ in $`𝔸^4/G_0`$. Using Picard’s Great Theorem we conclude that $`j^1\overline{𝗉}𝗆`$ is a biholomorphism and so $`\mathrm{𝗉𝗆}`$ is a biholomorphism. ∎
### 4.3 The inverse of the period map
We denote by
$$F=(F_0,F_1,F_2,F_3):𝒫T$$
the inverse of the period map.
###### Proposition 10.
The following is true:
1. $`F_0(x)=det(x)^1`$.
2. For $`i=2,3`$
$$F_i=det(x)^{1i}\stackrel{ˇ}{g_i}\stackrel{ˇ}{M_{2i}^0}$$
where $`g_i`$ is the Eisenstein series (1).
3. $`F_1=\stackrel{ˇ}{g_1}\stackrel{ˇ}{M_2^1}`$.
###### Proof.
Taking $`F`$ of (37) we have
$$F_0(xg)=F_0(x)k_1^1k_2^1,$$
(38)
$$F_1(xg)=F_1(x)k_1^1k_2+k_3k_1^1,$$
$$F_2(xg)=F_2(x)k_1^3k_2,F_3(xg)=F_3(x)k_1^4k_2^2,x,gG_0$$
By Legendre’s Theorem $`det(x)`$ is equal to one on $`V:=\mathrm{𝗉𝗆}(1\times 0\times 𝔸\times 𝔸)`$ and so the same is true for $`F_0det(x)`$. But the last function is invariant under the action of $`G_0`$ and so it is the constant function $`1`$. This proves the first item. Let $`G_i=F_idet(x)^{i1},i=1,2,3`$. The equalities (38) imply that $`G_i,i=2,3`$ do not depend on $`x_2,x_4`$. Now the map $`(t_2,t_3)\pi \mathrm{𝗉𝗆}(1,0,t_2,t_3)`$, where $`\pi `$ is the projection on the $`x_1,x_3`$ coordinates, is the classical period map (see for instance see and and the appendix of ) and this implies that $`G_i=\stackrel{ˇ}{g_i},i=2,3`$. Note that in our definition of the period map the factor $`\frac{1}{\sqrt{2\pi i}}`$ appears. In particular $`F_i,i=2,3`$ have finite growth at infinity. The fact that $`F_1`$ has finite growth at infinity follows form the Ramanujan relations (21) and the equality $`\frac{d}{dz}=2\pi iq\frac{d}{dq}`$. Since $`G_1\stackrel{ˇ}{M_2^1}`$, $`\stackrel{ˇ}{M_2^1}`$ is a one dimensional space, both $`g_1,G_1`$ satisfy (38) and $`M_2^0=\mathrm{}`$, we have $`G_1=g_1`$. ∎
### 4.4 Ramanujan relations
We proved in Lemma 9 that the period map $`\mathrm{𝗉𝗆}`$ associated to $`\stackrel{~}{\omega }`$ is a biholomorphism. According to (34), the inverse $`F`$ of $`\mathrm{𝗉𝗆}`$ satisfies the differential equation
$$x.A(F(x))^𝗍=I$$
We consider $`\mathrm{𝗉𝗆}`$ as a map sending the vector $`(t_0,t_1,t_2,t_3)`$ to $`(x_1,x_2,x_3,x_4)`$. Its derivative at $`t`$ is a $`4\times 4`$ matrix whose $`i`$-th column constitutes of the first and second row of $`\frac{1}{\mathrm{\Delta }}xA_i^𝗍`$. We use (31) to derive the equality
$$\left(dF\right)_x=\left(d\mathrm{𝗉𝗆}\right)_t^1=$$
$$det\left(x\right)^1\left(\begin{array}{cccc}F_0x_4& F_0x_3& F_0x_2& F_0x_1\\ \frac{1}{12F_0}\left(12F_0F_1^2x_312F_0F_1x_4F_2x_3\right)& F_1x_3+x_4& \frac{1}{12F_0}\left(12F_0F_1^2x_1+12F_0F_1x_2+F_2x_1\right)& F_1x_1x_2\\ 4F_1F_2x_33F_2x_46F_3x_3& F_2x_3& 4F_1F_2x_1+3F_2x_2+6F_3x_1& F_2x_1\\ \frac{1}{3F_0}\left(18F_0F_1F_3x_312F_0F_3x_4F_2^2x_3\right)& 2F_3x_3& \frac{1}{3F_0}\left(18F_0F_1F_3x_1+12F_0F_3x_2+F_2^2x_1\right)& 2F_3x_1\end{array}\right).$$
For $`g_i:=F_i_\stackrel{~}{}`$ we obtain the Ramanujan relations (21).
### 4.5 Other topics
We have seen that the period map satisfies the differential equation
(39)
$$d(\mathrm{𝗉𝗆}^𝗍)=A\mathrm{𝗉𝗆}^𝗍$$
This can be interpreted as a multi-variable Fuchsian system/Picard-Fuchs equation. Restricting (39) to a line in $`𝔸^4`$, one gets a usual one variable Picard-Fuchs equation. For many of them the solution and hence the associated abelian integral can be given explicitly by hypergeometric functions (see ). It seems to me that with the notion of period map in this article it is possible to classify all the algebraic values of $`G`$-functions obtained by the methods in .
Another interesting subject could be the construction of a logarithmic structure on $`T`$ in the context of algebraic geometry. Such constructions for the moduli of polarized Hodge structures is done in and in the context of analytic geometry .
## 5 The Ramanujan foliation
The whole theory developed in the previous sections could be done by setting $`t_0=1`$. In this section we assume that $`t_0=1`$ and use the same notations for $`\mathrm{𝗉𝗆},𝒫,G_0,T,\mathrm{\Delta }`$ and so on. For instance redefine
$$𝒫:=\{x=\left(\begin{array}{cc}x_1& x_2\\ x_3& x_4\end{array}\right)\mathrm{GL}(2,)\mathrm{Im}(x_1\overline{x_3})>0,det(x)=1\}$$
and
$$G_0=\{\left(\begin{array}{cc}k& k^{}\\ 0& k^1\end{array}\right)k^{},k^{}\}$$
The action of $`G_0`$ on $`𝔸^3`$ is given by
$$tg:=(t_1k^2+k^{}k^1,t_2k^4,t_3k^6),t=(t_1,t_2,t_3)𝔸^3,g=\left(\begin{array}{cc}k& k^{}\\ 0& k^1\end{array}\right)G_0$$
We also define
$$g=(g_1,g_2,g_3):T𝔸^3$$
### 5.1 The Ramanujan foliation
We write (8) in the vector filed form
$$\mathrm{Ra}:=(t_1^2\frac{1}{12}t_2)\frac{}{t_1}+(4t_1t_26t_3)\frac{}{t_2}+(6t_1t_3\frac{1}{3}t_2^2)\frac{}{t_3}$$
It is also useful to define the differential forms
$$\eta _1:=(t_1^2\frac{1}{12}t_2)dt_2(4t_1t_26t_3)dt_1$$
$$\eta _2:=(4t_1t_26t_3)dt_3(6t_1t_3\frac{1}{3}t_2^2)dt_2,$$
$$\eta _3:=(t_1^2\frac{1}{12}t_2)dt_3(6t_1t_3\frac{1}{3}t_2^2)dt_1$$
and say the the foliation $`(\mathrm{dR})`$ is induced by $`\eta _i,i=1,2,3`$. The singularities of (8) are the points $`t𝔸^3`$ such that $`\mathrm{Ra}(t)=0`$. It turns out that
$$\mathrm{Sing}(\mathrm{Ra})=\{(t_1,12t_1^2,8t_1^3)t_1𝔸\},$$
which is a one dimensional curve and lies in $`\{\mathrm{\Delta }=0\}`$. We have
$`d\mathrm{\Delta }(\mathrm{Ra})`$ $`=`$ $`(2.27t_3dt_33t_2^2dt_2)(\mathrm{Ra})`$
$`=`$ $`2.27t_3(6t_1t_3{\displaystyle \frac{1}{3}}t_2^2)3t_2^2(4t_1t_26t_3)`$
$`=`$ $`12t_1\mathrm{\Delta }`$
This implies that the variety $`\mathrm{\Delta }_0:=\{\mathrm{\Delta }=0\}`$ is invariant by the foliation $`(\mathrm{Ra})`$. Inside $`\mathrm{\Delta }_0`$ we have the algebraic leaf $`\{(t_1,0,0)𝔸^3\}`$ of $`(\mathrm{Ra})`$. On $`\mathrm{Sing}(\mathrm{Ra})`$ we have a special point $`p_{\mathrm{}}`$ given by (18). It is the limit of $`g(z)`$ when $`\mathrm{Im}(z)`$ tends to $`+\mathrm{}`$.
###### Proposition 11.
The following is a uniformization of the foliation $`(\mathrm{Ra})`$ restricted to $`T`$:
$$u:\times (𝔸^2\backslash \{(0,0)\})T$$
(40)
$$(z,c_2,c_4)g(z)\left(\begin{array}{cc}(c_4zc_2)^1& c_4\\ 0& c_4zc_2\end{array}\right)=$$
$$(g_1(z)(c_4zc_2)^2+(c_4zc_2),g_2(z)(c_4zc_2)^4,g_3(z)(c_4zc_2)^6).$$
###### Proof.
One may check directly that for fixed $`c_2,c_4`$ the map induced by $`u`$ is tangent to (8). In general it is not a solution of (8). This is the main reason for naming “foliation”. We give another proof which uses the period map: From (31) we have
$$d(\mathrm{𝗉𝗆})(t)=\mathrm{𝗉𝗆}(t)\left(\begin{array}{cc}\frac{3}{4}\eta _2& \frac{3}{2}(3t_3dt_22t_2dt_3)\\ \frac{9}{2}t_3\eta _13t_2\eta _3+\frac{3}{2}t_1\eta _2& \frac{3}{4}\eta _2\end{array}\right)^𝗍.$$
Therefore
$$d(\mathrm{𝗉𝗆}(t))(\mathrm{Ra}(t))=\mathrm{𝗉𝗆}(t)\left(\begin{array}{cc}0& 0\\ & 0\end{array}\right)=\left(\begin{array}{cc}& 0\\ & 0\end{array}\right).$$
This implies that the $`x_2`$ and $`x_4`$ coordinates of the pull forward of the vector field $`\mathrm{Ra}`$ by $`\mathrm{𝗉𝗆}`$ are zero. Therefore, the leaves of $`(\mathrm{Ra})`$ in the period domain are of the form
$$\left(\begin{array}{cc}z(c_4zc_2)^1& c_2\\ (c_4zc_2)^1& c_4\end{array}\right)=\left(\begin{array}{cc}z& 1\\ 1& 0\end{array}\right)\left(\begin{array}{cc}(c_4zc_2)^1& c_4\\ 0& c_4zc_2\end{array}\right).$$
Due to Proposition 11, our foliation $`(\mathrm{Ra})`$ can be considered as a kind of Hilbert modular foliation (see ).
### 5.2 The family $`y^24x^3+t_1x^2+t_2x+t_3`$
The family (7) can be rewritten in the form
$$y^24t_0x^3+12t_0t_1x^2+(12t_0t_1^2+t_2)x+(4t_0t_1^3t_2t_1+t_3)=0$$
The mapping
$$\alpha :𝔸^4𝔸^4,t(t_0,12t_0t_1,12t_0t_1^2+t_2,4t_0t_1^3t_2t_1+t_3)$$
is an isomorphism and so we can re state Proposition 9 for the family
$$_t:y^24t_0x^3t_1x^2t_2xt_3=0.$$
The inverse of the period map in this case is given by $`G=(G_0,G_1,G_2,G_3)`$ with
$$G_0=F_0\stackrel{ˇ}{M_0^1},G_1=12F_0F_1\stackrel{ˇ}{M_2^0},$$
$$G_2=12F_0F_1^2+F_2\stackrel{ˇ}{M_4^1},G_3=4F_0F_1^3F_2F_1+F_3\stackrel{ˇ}{M_6^2}$$
In this case the singular fibers are parameterized by the zeros of
$$\mathrm{\Delta }:=t_0(432t_0^2t_3^2+72t_0t_1t_2t_316t_0t_2^3+4t_1^3t_3t_1^2t_2^2).$$
The Ramanujan relations take the simpler form:
(41)
$$\{\begin{array}{c}\dot{t_1}=t_2\hfill \\ \dot{t_2}=6t_3\hfill \\ \dot{t_3}=t_1t_3\frac{1}{4}t_2^2\hfill \end{array}$$
They have the solution
$$g:=(12g_1,12g_1^2+g_2,4g_1^3g_2g_1+g_3)$$
## 6 Proofs
Now we are in a position to prove the theorems announced in the Introduction.
### 6.1 Proof of Theorem 1
We prove that $`\stackrel{ˇ}{M}`$ as a $``$-algebra is freely generated by $`\frac{1}{F_0},F_i,i=0,1,2,3`$. Let $`\stackrel{~}{F}\stackrel{ˇ}{M_m^n}`$ and $`\stackrel{~}{F}_i\stackrel{ˇ}{M_{m2i}^{ni}}`$ be its associated functions. Since the period map $`\mathrm{𝗉𝗆}:T`$ is a biholomorphism, there exist holomorphic functions $`p_ii=0,1,\mathrm{},n,p_0:=p`$ defined on $`T`$ such that $`\stackrel{~}{F}_i=p_i(F_0,F_1,F_2,F_3)`$. The property (22) of $`\stackrel{~}{F}`$ implies that:
(42)
$$p(tg)=k_2^nk_1^{nm}\underset{i=0}{\overset{n}{}}\left(\genfrac{}{}{0pt}{}{n}{i}\right)k_3^ik_2^ip_i(t),gG_0,tT.$$
Take $`g=\left(\begin{array}{cc}1& t_1\\ 0& 1\end{array}\right)`$ and $`t=(t_0,0,t_1,t_3)`$. Then
$$p(t_0,t_1,t_2,t_3)=\underset{i=0}{\overset{n}{}}\left(\genfrac{}{}{0pt}{}{n}{i}\right)t_1^ip_i(t_0,0,t_2,t_3).$$
This implies that $`p`$ is a polynomial of degree at most $`n`$ in the variable $`t_1`$. Let us re write $`p(t)=_{i=0}^nt_1^iq_i`$, where $`q_i`$’s are holomorphic functions on $`𝔸^3\backslash \{t𝔸^3\mathrm{\Delta }=0\}`$. We apply (42) to $`g=\left(\begin{array}{cc}k& 0\\ 0& t_0^1k^1\end{array}\right)`$ and consider the coefficients of $`t_1^i,i=1,2\mathrm{},n`$. We get
(43)
$$q_i(1,t_2t_0^1k^4,t_3t_0^1k^6)=t_0^nk^{m+2i}q_i(t),i=1,2,\mathrm{},n.$$
Take $`t_0=1`$. The growth condition on $`\stackrel{~}{F}`$ is translated through the period map into the following fact: $`p`$ restricted to a transversal disk to $`\mathrm{\Delta }=0`$ at $`p_{\mathrm{}}`$ has a holomorphic extension to $`p_{\mathrm{}}`$. This will also imply the similar growth conditions for $`q_i`$’s. The classical fact that the set of modular forms is generated by the Eisenstein series $`g_2`$ and $`g_3`$ and (43) with $`t_0=1`$ imply that $`q_i(1,t_2,t_3)`$ is a homogeneous polynomial of degree $`m2i`$ in the graded ring $`[t_2,t_3],\mathrm{deg}(t_2)=4,\mathrm{deg}(t_3)=6`$. We conclude that $`p`$ is of the form
$$p=t_0^n\underset{i=0}{\overset{n}{}}t_1^iq_i(1,t_2t_0^1,t_3t_0^1).$$
In other words, $`p`$ is homogeneous of degree $`m`$ in the variables $`t_1,t_2,t_3`$ with $`\mathrm{deg}(t_i)=2i,i=1,2,3`$.
### 6.2 Proof of Theorem 2
In §2.6 we described some analytic functions $`B_i,i=1,2,3`$, on $``$ which had some compatibility properties with the action of $`G_0`$ on $``$. We use Proposition 9 and transfer them to the world of coefficients. This will prove the existence and uniqueness of the functions $`B_1,B_2,B_3`$. For the sake of simplicity we have used the same letters to name these functions.
By the properties that $`B_1`$ has we can say more about it. In (4) we put $`k=1`$ and we conclude that $`B_1`$ is independent of the variable $`t_1`$. The function $`B_2|\mathrm{\Delta }|^{\frac{1}{6}}`$ is $`G_0`$ invariant and so there is an analytic function $`b_2:𝔸`$ such that
$$B_2(t)=\frac{b_2(j(t))}{|\mathrm{\Delta }(t)|^{\frac{1}{6}}}.$$
Translating this to $``$, we have
$$\mathrm{Im}(z)=\frac{b_2(j(z))}{|\mathrm{\Delta }(z)|^{\frac{1}{6}}}$$
where the above $`j`$ and $`\mathrm{\Delta }`$ are the ones on §2.2.
The proof of the last part of the theorem is as follows: On $`M_0`$ an $`x𝒫`$ can be written in the form $`\left(\begin{array}{cc}x_1& x_4r\\ x_3& x_4\end{array}\right),r,x_4(x_1rx_3)=1`$. Then
(44)
$$B_3(x)=\overline{x_4}(x_1rx_3)=\frac{\overline{x_4}}{x_4}.$$
### 6.3 Proof of Theorem 3
We follow the notations introduced in §5. In particular we work with the family (7) with $`t_0=1`$. The leaves of the pull-forward of the foliation $`(\mathrm{Ra})`$ by the period map $`\mathrm{𝗉𝗆}`$ have constant $`x_2`$ and $`x_4`$ coordinates. By definition of $`B_2:=\mathrm{Im}(x_2\overline{x_4})`$ in the period domain, we conclude that $`M_r`$’s are $`(\mathrm{Ra})`$-invariant. The equality (44) implies that $`N_w`$’s are $`(\mathrm{Ra})`$-invariants.
Let us now prove item 2: Take $`tK`$ and a cycle $`\delta H_1(_t,)`$ such that $`_\delta \stackrel{~}{\omega }_2=0`$ and $`\delta `$ is not of the form $`n\delta ^{}`$ for some $`2n`$ and $`\delta ^{}H_1(_t,)`$. We choose another $`\delta ^{}H_1(_t,)`$ such that $`(\delta ^{},\delta )`$ is a basis of $`H_1(_t,)`$ and the intersection matrix in this basis is $`\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$. Now $`\mathrm{𝗉𝗆}(t)`$ has zero $`x_4`$-coordinates and so its $`B_2`$ is zero. This implies that $`KM_0`$. It is dense because an element $`\left(\begin{array}{cc}x_1& x_4r\\ x_3& x_4\end{array}\right)M_0`$ can be approximated by the elements in $`M_0`$ with $`r`$.
The image of the map $`g`$ is the locus of the points $`t_0`$ in $`T`$ such that $`\mathrm{𝗉𝗆}(t_0)`$ is of the form $`\left(\begin{array}{cc}z& 1\\ 1& 0\end{array}\right)`$ in a basis of $`H_1(_t,)`$. We look $`g`$ as a function of $`q=e^{2\pi iz}`$ and we have
$$g(0)=p_{\mathrm{}},\frac{g}{q}(0)=(24a_1,240a_2,504a_3)$$
where $`p_{\mathrm{}}=(a_1,a_2,a_3)`$ is the one in (18). This implies that the image of $`g`$ intersects $`\mathrm{Sing}(\mathrm{Ra})`$ transversally. For a $`tK`$ the $`x_4`$-coordinate of $`\mathrm{𝗉𝗆}`$ is zero and the leaf through $`t`$, namely $`L_t`$, has constant $`x_2`$-coordinate, namely $`c_2`$. By (40) $`L_t`$ is uniformized by
$$u(z)=(c_2^2g_1(z),c_2^4g_2(z),c_2^6g_3(z)),z.$$
This implies that $`L_t`$ intersects $`\mathrm{Sing}(\mathrm{Ra})`$ transversally at $`(c_2^2a_1,c_2^4a_2,c_2^6a_3)`$.
We prove item 3: Let $`tT`$ and the leaf $`L_t`$ through $`t`$ have an accumulation point at $`t_0T`$. We use the period map $`\mathrm{𝗉𝗆}`$ and look $`(\mathrm{Ra})`$ in the period domain. For $`(c_2,c_4)𝔸^2\backslash \{0\}`$ the set $`S=\{A(c_2,c_4)^𝗍A\mathrm{SL}(2,)\}`$ has an accumulation point in $`𝔸^2`$ if and only if $`\frac{c_2}{c_4}\mathrm{}`$ or equivalently $`B_2(t)=0`$.
### 6.4 Proof of Theorem 4
Let $`k`$ be an algebraically closed field of charachteristic $`0`$ , for instance take $`k=\overline{}`$.
###### Proposition 12.
The quasi affine variety
$$T=\mathrm{Spec}(k[t,\frac{1}{\mathrm{\Delta }}])$$
is the moduli of $`(E,[\omega _1],[\omega _2])`$’s, where $`E`$ is an elliptic curve defined over $`k`$, $`\omega _1`$ is a differential form of the first kind on $`E`$ and $`([\omega _1],[\omega _2])`$ is a basis of $`H_{\mathrm{dR}}^1(E)`$.
###### Proof.
For simplicity we do not write more $`[.]`$ for differential forms. The $`j`$ invariant (36) classifies the ellipric curves over $`k`$ (see Theorem 4.1). Therefore, for a given elliptic curve $`E/k`$ we can find parameter $`t𝔸_k^4`$ such that $`E_t`$ over $`k`$. Under this isomorphism we write
$$\left(\begin{array}{c}\omega _1\\ \omega _2\end{array}\right)=g^𝗍\left(\begin{array}{c}\frac{dx}{y}\\ \frac{xdx}{y}\end{array}\right),\text{ in }H_{\mathrm{dR}}^1(_t)$$
for some $`gG_0`$, where $`\omega _1,\omega _2`$ are as in the proposition. Now, the triple $`(E,\omega _1,\omega _2)`$ is isomorphic to $`(_{tg},\frac{dx}{y},\frac{xdx}{y})`$. Since $`j:𝔸^4/G_0𝔸`$ is an isomorphism, every triple $`(E,\omega _1,\omega _2)`$ is represented exactly by one parameter $`tT`$. ∎
By Proposition 12 the hypothesis of Theorem 4 gives us a parameter $`tT`$ such that $`_\delta \frac{xdx}{y}=0`$, for some $`\delta H_1(_t,)`$. We can assume that $`\delta `$ is not a multiple of another cycle in $`H_1(_t,)`$. The corresponding period matrix of $`t`$ in a basis $`(\delta ^{},\delta )`$ of $`H_1(_t,)`$ has zero $`x_4`$-coordinate and so the numbers
$$t_0=det(x)^1,t_i=det(x)^{1i}x_3^{2i}g_i(\frac{x_1}{x_3}),i=2,3,t_1=F_1(\left(\begin{array}{cc}x_1& x_2\\ x_3& 0\end{array}\right))=det(x)x_3^2g_1(\frac{x_1}{x_3})$$
all are in $`\overline{}`$. This implies that for $`z=\frac{x_1}{x_3}`$ we have
$$\frac{g_3}{g_1^3}(z),\frac{g_4}{g_1^2}(z),\frac{g_3^2}{g_2^3}(z)\overline{}.$$
This is in contradiction with
Theorem (Nesterenko 1996, ) For any $`z`$, the set
$$e^{2\pi iz},\frac{g_1(z)}{a_1},\frac{g_2(z)}{a_2},\frac{g_3(z)}{a_3}$$
contains at least three algebraically independent numbers over $``$.
A direct corollary of Theorem 4 is that the multivalued function
$$I(t)=\frac{_{\delta _t}\frac{xdx}{y}}{_{\delta _t}\frac{dx}{y}}$$
defined in $`T`$ never takes algebraic values for algebraic $`t`$.
### 6.5 Other topics
As a person who has started his mathematical career by studying holomorphic foliations in complex manifolds, I would be interested to describe completely the dynamics of the foliation $`(\mathrm{Ra})`$ and in particular the behavior of the leaves near $`\mathrm{\Delta }=0`$. The leaves of $`(\mathrm{Ra})`$ in $`\mathrm{\Delta }=0`$ parameterize degenerated elliptic curves. Can one describe their behavior by abelian integrals?
Our proof of Theorem 1 is completely based on the existence of the nice family (7) for which the period map is a biholomorphism. To prove Theorem 1 for certain subgroups of $`\mathrm{SL}(2,)`$ by the methods of this article, one must find a four parameter family of elliptic curves such that the period map is an etale covering, i.e. it is a local biholomorphism of finite degree. The inverse of the period map is a multi valued function whose restriction to the simply connected domain $`\stackrel{~}{}`$ gives rise to a finite number of holomorphic functions on $``$. These new functions are differential modular with respect to a subgroup of $`\mathrm{SL}(2,)`$ which can be calculated by some topological data, such as the homotopy group of $`T`$, attached to the period map.
The Moduli space $`T`$ is not a Shimura variety (see ). In this point it would be too much speculation to say that spaces like $`T`$ can be constructed for arbitrary moduli of polarized Hodge structures. Nevertheless, more constructions of such spaces may result in generalizations of Shimura varieties. |
warning/0506/astro-ph0506026.html | ar5iv | text | # The all-sky distribution of 511 keV electron-positron annihilation emission Based on observations with INTEGRAL, an ESA project with instruments and science data centre funded by ESA member states (especially the PI countries: Denmark, France, Germany, Italy, Switzerland, Spain), Czech Republic and Poland, and with the participation of Russia and the USA.
## 1 Introduction
Since the first detection (Johnson & Haymes johnson73 (1973)) and the subsequent firm identification (Leventhal et al. leventhal78 (1978)) of the galactic 511 keV annihilation line, the origin of galactic positrons has been a lively topic of scientific debate. Among the proposed candidates for sources of positrons figure cosmic-ray interactions with the interstellar medium (Ramaty et al. ramaty70 (1970)), pulsars (Sturrock sturrock71 (1971)), compact objects housing either neutron stars or black holes (Ramaty & Lingenfelter ramaty79 (1979)), gamma-ray bursts (Lingenfelter & Hueter lingenfelter84 (1984)), (light) dark matter (Rudaz & Stecker rudaz88 (1988); Boehm et al. boehm04 (2004)), and stars expelling radioactive nuclei produced by nucleosynthesis, such as supernovae (Clayton clayton73 (1973)), hypernovae (Cassé et al. casse04 (2004)), novae (Clayton & Hoyle clayton74 (1974)), red giants (Norgaard norgaard80 (1980)), and Wolf-Rayet stars (Dearborn & Blake dearborn85 (1985)). It seems difficult to disentangle the primary galactic positron source based only on theoretical grounds, mainly due to the (highly) uncertain positron yields, but also due to the uncertain distribution and duty cycle of the source populations.
Help is expected from a detailed study of the 511 keV line emission morphology. The celestial 511 keV intensity distribution should be tied to the spatial source distribution, although positron diffusion and effects associated with the annihilation physics may to some extent blur this link. First estimations of the 511 keV emission morphology were obtained by the Oriented Scintillation Spectrometer Experiment (OSSE) on-board the Compton Gamma-Ray Observatory (CGRO) satellite (Purcell et al. purcell94 (1994); Cheng et al. cheng97 (1997); Purcell et al. purcell97 (1997); Milne et al. milne00 (2000); Milne et al. milne01 (2001)), but observations were restricted to the inner Galaxy, giving only a limited view of the 511 keV emission distribution. With the launch of ESA’s INTEGRAL satellite in October 2002, a new gamma-ray observatory is available that allows a detailed study of positron annihilation signatures. In particular, the imaging spectrometer SPI (Vedrenne et al. vedrenne03 (2003)), one of the two prime instruments on-board INTEGRAL, has been optimised for the study of line radiation, combining high-resolution spectroscopy (R $`250`$ at 511 keV) with modest angular resolution ($`3^{}`$ FWHM).
We present in this work an all-sky map of 511 keV gamma-ray line emission, with the goals of determining the morphology of the emission in the Galaxy and of searching for previously unknown sources of 511 keV emission anywhere in the sky. The present public data archive does not yet cover the entire celestial sphere, but the unexposed regions are limited to a few areas at high galactic latitudes, comprising less than $`5\%`$ of the sky.The resulting point-source sensitivity is better than $`2\times 10^4`$ ph cm<sup>-2</sup>s<sup>-1</sup> for many regions along the galactic plane, allowing for the first time the extraction of information about the distribution of positron annihilation all over the Galaxy. We do not address the distribution of positronium continuum emission in this paper, since the subtraction of the diffuse galactic continuum emission is a distinct data analysis challenge. A map of positronium continuum emission will be presented elsewhere (Weidenspointner et al., in preparation).
Earlier results on the 511 keV line emission morphology as observed by SPI have been presented by Jean et al. (2003a ), Knödlseder et al. (2004a ), and Weidenspointner et al. (weidenspointner04 (2004)), and were based on observations performed during the galactic centre deep exposure (GCDE) of 2003. Using a “light bucket” approach which neglects the coding properties of the SPI mask, Teegarden et al. (teegarden05 (2005)) derived upper limits on electron-positron annihilation radiation from the galactic disk using core-programme data combined with open-programme observations at low galactic latitudes ($`|b|20^{}`$). In the present paper we provide for the first time an all-sky analysis using all public data of the first INTEGRAL mission year.
Spectroscopic characteristics of the 511 keV line based on SPI data have been published by Jean et al. (2003a ), Lonjou et al. (lonjou04 (2004)), and Churazov et al. (churazov05 (2005)). We will present the 511 keV line profile that we obtain from the all-sky dataset elsewhere (Jean et al., in preparation).
This paper is organised as follows. Section 2 describes the observations and the data preparation. Section 3 explains the treatment of the instrumental background. In section 4, we present the first all-sky map of 511 keV gamma-ray line radiation and determine the morphology of the emission. Section 4 also describe searches for correlations with tracers of galactic source populations in order to shed light on the origin of the positrons. In section 5 we discuss the implications of the observations for the galactic origin of positrons, and we conclude in section 6.
## 2 Observations and Data Preparation
The data that were analysed in this work consist of those included in the December 10, 2004 public INTEGRAL data release (i.e. orbital revolutions 19-76, 79-80, 89-122) plus the INTEGRAL Science Working Team data of the Vela region observed during revolutions 81-88. The data span the IJD epoch $`1073.3941383.573`$, where IJD is the Julian Date minus 2 451 544.5 days.
We screened the data for anomalously high counting rates (typically occurring at the beginning and the end of an orbital revolution due to the exit and entry of the radiation belts) and for periods of solar activity (as monitored by the SPI anticoincidence system) and excluded these periods from the data. This data screening has turned out to be crucial for reducing the systematic uncertainties in the data analysis related to instrumental background variations. After data screening, the dataset consists of 6821 pointed observations, with a total exposure time of 15.3 Ms. Typical exposure times per pointing are $`12003400`$ seconds, but a few long staring observations of up to 113 ks exposure time are also included.
Figure 1 shows a map of the resulting effective SPI exposure at 511 keV. The maximum exposure of $`2.1\times 10^8`$ cm<sup>2</sup> s occurs towards the galactic centre region thanks to data obtained during a long dedicated observation of this region.<sup>1</sup><sup>1</sup>1 To obtain the effective exposure time, the exposure has to be divided by the effective area at 511 keV of about 75 cm<sup>2</sup>. A relatively uniform exposure of $`3\times 10^7`$ cm<sup>2</sup> s has been achieved for galactic longitudes $`|l|50^{}`$ and latitudes $`|b|15^{}`$. Regions of peculiarly high exposure ($`5\times 10^7`$ cm<sup>2</sup> s) are found in Cygnus, Vela and towards the Large Magellanic Cloud. In addition, particularly well exposed sources ($`2\times 10^7`$ cm<sup>2</sup> s) are the Crab nebula, 3C 273, NGC 4151, M 94, NGC 936 (during the SN2003 gs outburst) and the Coma cluster. Unexposed regions are found mostly at intermediate galactic latitudes ($`|b|30^{}60^{}`$), and towards the south galactic pole.
A map of the resulting narrow-line $`3\sigma `$ point-source sensitivity of SPI at 511 keV is shown in Fig. 2. To evaluate the sensitivity, an energy band of 7 keV centred at 511 keV has been used. The choice of such a relatively wide band eliminates any bias due to the germanium detector degradation and annealing cycles, as well as any bias/effect due to gain calibration uncertainties. It also takes into account moderate 511 keV line broadening, as reported by Jean et al. (2003a ).
Over large regions of the sky, and in particular in the galactic plane, a sensitivity better than $`2\times 10^4`$ ph cm<sup>-2</sup>s<sup>-1</sup> is reached. A best point-source sensitivity of $`5\times 10^5`$ ph cm<sup>-2</sup>s<sup>-1</sup> is achieved towards the galactic centre direction. The sensitivity to extended diffuse emission becomes slightly worse with increasing emission size, and depends on the exposure pattern in the region of interest. For example, for a 2d angular Gaussian surface brightness distribution centred on the galactic centre, the 511 keV line sensitivity worsens from $`5\times 10^5`$ ph cm<sup>-2</sup>s<sup>-1</sup> for a galactic centre point-source to $`7\times 10^5`$ ph cm<sup>-2</sup>s<sup>-1</sup> for an extended source of $`8^{}`$ FWHM.
Only single-detector event data have been analysed in this work (multiple-detector event data do not contribute significantly to the SPI sensitivity at an energy of 511 keV; c.f. Roques et al. roques03 (2003)). Energy calibration was performed orbit-wise, resulting in a relative (orbit-to-orbit) calibration precision of $`0.01`$ keV and an absolute accuracy of $`0.05`$ keV (Lonjou et al. lonjou04 (2004)).
The data have been analysed by sorting the events in a 3-dimensional data-space, spanned by the (calibrated) event energy, the detector number, and the SPI pointing number. An energy binning of 0.5 keV has been chosen, well below the instrumental energy resolution of $`2.12`$ keV at 511 keV.
## 3 Background modelling
The most crucial step in SPI data analysis consists of the precise modelling of the time variability of the instrumental background. In the region of the 511 keV line, the instrumental background consists of a nearly flat continuum and a (broadened) instrumental 511 keV line originating from positron annihilation within the telescope (Teegarden et al. teegarden04 (2004)). Since the time variation of the continuum component differs from that of the line component we model them independently. The background model for a given data-space bin, indexed in the following by the pointing number $`p`$, the detector number $`d`$ and the energy bin $`e`$, is then given by
$$b_{p,d,e}=b_{p,d,e}^{\mathrm{cont}}+b_{p,d,e}^{\mathrm{line}}$$
(1)
(note that for the analysis presented in this work a single energy bin has been used, covering the energy interval $`507.5514.5`$ keV; however, for clarity and reference in future works we give here the complete energy-dependent formalism).
The time variation of the continuum component is extrapolated from that observed in an continuum energy band adjacent to the 511 keV line. We used the energy band $`E_{\mathrm{adj}}=523545`$ keV, situated above the 511 keV line, in order to exclude any bias due to positronium continuum emission that appears below 511 keV. To reduce the statistical uncertainty that arises from the limited counting statistics, we smoothed the time variation by locally adjusting the rate of saturated events in the germanium detectors (GEDSAT) to the adjacent counting rate (GEDSAT turned out to provide a good first order tracer of the background variation in SPI; c.f. Jean et al. 2003b ). The predicted number of continuum background counts in data-space bin $`(p,d,e)`$ is then given by
$`b_{p,d,e}^{\mathrm{cont}}`$ $`=`$ $`g_{p,d}\times T_{p,d}\times `$ (2)
$`{\displaystyle \frac{\mathrm{\Delta }_e}{_{e^{}E_{\mathrm{adj}}}\mathrm{\Delta }_e^{}}}\times {\displaystyle \frac{_{p^{}=p\mathrm{\Delta }_p}^{p+\mathrm{\Delta }_p}_{e^{}E_{\mathrm{adj}}}n_{p^{},d,e^{}}}{_{p^{}=p\mathrm{\Delta }_p}^{p+\mathrm{\Delta }_p}g_{p^{},d}\times T_{p^{},d}}},`$
where
* $`g_{p,d}`$ is the GEDSAT rate for detector $`d`$, averaged over the time period spanned by pointing $`p`$, given in units of counts s<sup>-1</sup>,
* $`T_{p^{},d}`$ is the lifetime for detector $`d`$ during pointing $`p^{}`$, given in units of seconds,
* $`\mathrm{\Delta }_e`$ is the energy bin size for spectral bin $`e`$, given in units of keV (here $`\mathrm{\Delta }_e=0.5`$ keV), and
* $`n_{p^{},d,e^{}}`$ is the number of observed counts for pointing $`p^{}`$, detector $`d`$, and energy bin $`e^{}`$, given in units of counts.
The number of pointings used for smoothing, given by $`2\mathrm{\Delta }_p+1`$, is determined for each pointing $`p`$ and detector $`d`$ by satisfying the constraint
$$\underset{\mathrm{\Delta }_p0}{\mathrm{min}}\left(\underset{p^{}=p\mathrm{\Delta }_p}{\overset{p+\mathrm{\Delta }_p}{}}T_{p^{},d}T_{\mathrm{min}}\right).$$
(3)
An accumulated lifetime of $`T_{\mathrm{min}}=20`$ hours has shown to provide an optimum compromise between reducing the statistical uncertainty (due to the limited number of events in the adjacent energy band) and reducing the systematic uncertainty (due to the fact that the GEDSAT rate does not predict the background to infinite precision). In other words, continuum background variations shorter than $`20`$ hours are modelled by the GEDSAT rate while variations on longer time scales are modelled by the observed event rate in the $`523545`$ keV band.
The time variation of the line component was modelled for each detector $`d`$ and energy bin $`e`$ separately using a multi-component template of the form
$$b_{p,d,e}^{\mathrm{line}}=\beta _{d,e}^{(1)}+\beta _{d,e}^{(2)}\times g_{p,d}+\beta _{d,e}^{(3)}_{t_0}^tg_d(t^{})e^{(t^{}t)/\tau }dt^{}.$$
(4)
This template consists of a constant term $`\beta _{d,e}^{(1)}`$ plus the GEDSAT rate $`g_{p,d}`$ scaled by $`\beta _{d,e}^{(2)}`$ plus the GEDSAT rate $`g_d(t^{})`$ convolved with an exponential decay law, scaled by $`\beta _{d,e}^{(3)}`$ (the convolution integral is taken from the start of the INTEGRAL mission $`t_0`$ up to date $`t`$). The coefficients $`\beta _{d,e}^{(i)}`$ of the template are adjusted during the analysis for each SPI detector $`d`$ and energy bin $`e`$ using a maximum likelihood fitting procedure (again, in the analysis presented in this paper only a single energy bin is used). The constant term has been introduced to provide for non-linearities between the background variation and the GEDSAT rate. In fact, it turns out that $`\beta _{d,e}^{(1)}`$ are negative. An equally good background predictor is obtained if the GEDSAT rate raised to a power of $`1.1`$ is taken, but using a constant instead of a powerlaw has the advantage of having the background variation template decomposed into a linear combination of terms. The third component makes provision for a long term build-up that is seen in the intensity of the 511 keV background line, and that is tentatively attributed to production of the isotope <sup>65</sup>Zn which has a decay time of $`\tau =352`$ days. The precise value of $`\tau `$ is in fact weakly constrained by the present data, and a linear slope provides an equally good fit of the instrumental 511 keV line background.
Although the background model defined by Eqs. (1)–(4), which hereafter is called model DETE, predicts the instrumental background to good accuracy, significant residuals remain after subtracting off the background model and a model of the sky intensity distribution from the data (c.f. Fig. 3). We found that these residuals can lead to systematic biases in the study of the morphology of the 511 keV emission, in particular for the determination of the longitude profile of the emission. These biases can be explained by the telescope pointing strategy that has been adopted for a large fraction of the galactic centre deep exposure (GCDE) of the INTEGRAL core program: a slow ($`5^{}`$/day) scan of the galactic plane from negative towards positive longitudes combined with rapid ($`3^{}`$/hour) excursions in galactic latitude. As result, the longitude profile of the 511 keV line emission is encoded in count rate variations on timescales of days while the latitude profile is encoded in count rate variations on timescales of hours. Within a few hours the SPI instrumental background is sufficiently stable to be accurately predicted by our model, hence the latitude profile is rather well determined. However, on timescales of days the background variations are more difficult to predict to sufficient accuracy, potentially leading to systematic trends in the determination of the longitude profile.<sup>2</sup><sup>2</sup>2 In a preliminary analysis in which we treated a much smaller dataset, systematic background uncertainties suggested a significant elongation of the galactic centre bulge emission along the galactic plane. This elongation was artificial and had been produced by background variations that were not fully explained by our model. Removing the short period of data with the strongest background variations removed also the apparent elongation of the bulge emission.
In order to improve the background model on long timescales, we studied also a class of models where we adjust the longterm variations during model fitting. For this purpose we adjust the model parameters $`\beta _{d,e}^{(2)}`$ not only for all SPI detectors but also for time intervals of fixed duration $`T`$. In this way, systematic uncertainties in the background model on timescales longer than $`T`$ are removed. Fitting the background for each orbital revolution ($`T3`$ days) is adequate to reduce systematic trends well below the statistical uncertainties (c.f. bottom panel of Fig. 3). This method is similar to the method that we applied in our earlier works (Jean et al. 2003a ; Knödlseder et al. 2004a ; Weidenspointner et al. weidenspointner04 (2004)), with the difference that we now also fit the background model for each of the SPI detectors separately, and that we included in addition a constant term and a build-up term in the model (see Eq. 4). Hereafter this second background model is called ORBIT-DETE.
The introduction of additional parameters in ORBIT-DETE with respect to DETE leads to a substantial loss in sensitivity. The detection significance of galactic centre 511 keV line emission drops from $`50\sigma `$ for DETE to $`22\sigma `$ for ORBIT-DETE. However, it was found that the statistical accuracy of the morphology determination, which is driven by the count rate contrast in the data-space rather than the count rate level, is not degraded by the introduction of additional parameters, as long as $`T2`$ days. Consequently, using the ORBIT-DETE model for the morphological characterisation of the 511 keV line emission is the optimum choice that keeps a high statistical accuracy while reducing the systematic uncertainties in the analysis.
On the other hand, despite the systematic uncertainties, DETE is accurate enough to allow for a precise determination of 511 keV line flux levels. This is related to the fact that flux measurements require an average determination of the count rate level and are not sensitive to the count rate contrast. Apparently, the count rate residuals approximately average to zero (c.f. Fig. 3).
We therefore opted for a two step approach where we first determine the morphology using ORBIT-DETE, and then, using the optimum morphology parameters, determine the 511 keV flux using DETE. In this way we recover the good sensitivity of SPI for 511 keV flux measurements that was reduced by a factor of $`2`$ by the usage of ORBIT-DETE. The comparison of the flux levels determined using DETE and ORBIT-DETE provides us with a measure of the systematic uncertainty in the flux determination, which in general is smaller than the statistical uncertainty obtained with DETE. We add the systematic to the statistical uncertainty in quadrature and quote the result as total error on the flux measurement. In cases where uncertainties in the morphology (such as the size of the emission region) introduce some uncertainty on the flux, we have also added this uncertainty to the total error in quadrature.
## 4 Results
### 4.1 Imaging
To determine a model independent map of the 511 keV gamma-ray line intensity distribution over the sky, we employed the Richardson-Lucy algorithm (Richardson richardson72 (1972); Lucy lucy74 (1974)). This type of algorithm is widely used for image deconvolution, and has in particular been successfully employed for the analysis of gamma-ray data of CGRO (Knödlseder et al. 1999a ; Milne et al. milne00 (2000)).
We implemented the accelerated version ML-LINB-1 of Kaufman (kaufman87 (1987)) of the Richardson-Lucy algorithm for our analysis, which iteratively updates the sky intensity distribution $`f_j^kf_j^{k+1}`$ using the relation
$$f_j^{k+1}=f_j^k+\lambda ^kw_jf_j^k\left(\frac{_{i=1}^N\left(\frac{n_i}{e_i^k}1\right)R_{ij}}{_{i=1}^NR_{ij}}\right)$$
(5)
where $`R_{ij}`$ is the instrumental response matrix (linking the data space, indexed by $`i`$, to the image space, indexed by $`j`$), $`n_i`$ is the number of counts measured in data space bin $`i`$, $`e_i^k=_{j=1}^MR_{ij}f_j^k+b_i`$ is the predicted number of counts in data space bin $`i`$ after iteration $`k`$ ($`b_i`$ being the predicted number of instrumental background counts for bin $`i`$), $`N`$ and $`M`$ are the dimensions of the data and image space, respectively, and $`\lambda ^k`$ is an acceleration factor that is obtained by constrained maximum likelihood fitting (with the constraint that the resulting sky intensities remain positive).
To avoid noise artefacts in the weakly exposed regions of the sky, we weighted the image increment with a quantity that is related to the sensitivity of the instrument, given by $`w_j=(_{i=1}^NR_{ij})^{1/2}`$. We verified that introducing this weighting had no impact on the image reconstruction in the well exposed regions of the sky. In addition, we smoothed the iterative corrections on the right hand side of Eq. 5 using a $`5^{}\times 5^{}`$ boxcar average. In this way the effective number of free parameters in the reconstruction is reduced and image noise is damped to an acceptable level. The application of more sophisticated image reconstruction methods involving wavelet based multi-resolution algorithms aiming at a complete suppression of image noise (Knödlseder et al. 1999a ) will be presented elsewhere.
The resulting all-sky image of the 511 keV line emission is shown in Fig. 4, longitude and latitude profiles of the emission are shown in Fig. 5. We have chosen to stop the iterative procedure after iteration 17 since at this point the recovered flux and the fit quality correspond approximately to the values that we achieve by fitting astrophysical models to the data (c.f. section 4.2). In this way we make sure that we are not in the regime of overfitting, which is characterised by substantial image noise and artificial image structures. On the other hand, simulations showed that faint diffuse emission, as expected for example for a galactic disk component, would not be recovered at this point.
Figure 4 reveals that the 511 keV sky is dominated by prominent emission from the bulge region of the Galaxy. Beyond the galactic bulge, no additional 511 keV emission is seen all over the sky, despite the good exposure in some regions (e.g. Cygnus, Vela, LMC, anticentre, north galactic pole region). The 511 keV emission appears symmetric and centred on the galactic centre, with indications for a slight latitude flattening. The latitude flattening could be either due to an inherent asymmetry of the bulge component or due to the presence of an underlying faint galactic disk component. Indeed, if the Richardson-Lucy iterations are continued, a faint disk-like structure emerges (c.f. Fig. 6). Yet the image starts to become polluted by noise and we cannot exclude the possibility that the apparent disk emission is artificially created by the exposure pattern that follows the galactic plane. Therefore we employ more quantitative methods in the next section to assess the significance of the possible disk emission.
By fitting Gaussian functions to the longitude and latitude profiles of the image (c.f. Fig. 5) we estimate the extent of the emission to $`13^{}\times 10^{}`$ (FWHM). Figure 5 indicates, however, that the emission profiles are not well represented by Gaussian functions. The emission is better described by a compact (FWHM $`5^{}`$) core and a more extended halo (FWHM $`10^{}20^{}`$). We want to emphasise, however, that this qualitative analysis should not be pushed too far, since image deconvolution is a non-linear process which is easily affected by image noise and exposure biases.
### 4.2 Morphological characterisation
#### 4.2.1 Method
To make a quantitative assessment of the morphology of the 511 keV line emission we use a maximum likelihood multi-component model fitting algorithm. Assuming Poisson noise for the measured number $`n_i`$ of events in each of the $`N`$ data-space bins, the algorithm maximises the log likelihood
$$\mathrm{ln}L=\underset{i=1}{\overset{N}{}}n_i\mathrm{ln}e_ie_i\mathrm{ln}n_i!,$$
(6)
where $`e_i=_k\alpha _ks_i^k+b_i(𝜷)`$ is the predicted number of (source plus background) counts in data space bin $`i`$, $`s_i^k=_{j=1}^Mf_j^kR_{ij}`$ is the sky intensity model $`f_j^k`$ folded into the data space ($`R_{ij}`$ being the instrumental response matrix), $`b_i(𝜷)`$ is the background model (c.f. Fig. 1), and $`\alpha _k`$ and $`𝜷`$ are scaling factors for the sky intensity and the background model, respectively, that are adjusted by the fit.
Detection significances (and parameter errors) are estimated using the maximum likelihood ratio test (Cash cash79 (1979)). We calculate the maximum log likelihood-ratio $`\text{MLR}=2(\mathrm{ln}L_0\mathrm{ln}L_1)`$ between two models (hypotheses), where for the first one we constrain a number $`p`$ of the parameters to specific values (resulting in $`L_0`$) while for the second one all parameters are left free (resulting in $`L_1`$). In the case that $`L_1`$ provides a satisfactory fit of the data, MLR is then distributed like a $`\chi _p^2`$ distribution with $`p`$ degrees of freedom. Statistical parameter errors were estimated using the formalism of Strong (strong85 (1985)). Throughout this paper the error bars quoted are $`1\sigma `$.
We call the maximum log likelihood-ratio (MLR) of a model the difference between the log likelihood obtained by fitting all model parameters and the log likelihood obtained by fitting only the background model to the data (i.e. for $`L_1`$ all parameters $`\alpha _k`$ and $`𝜷`$ vary freely while for $`L_0`$ all $`\alpha _k`$ are constrained to zero and only the $`𝜷`$ are allowed to vary). To compare models with different numbers of free parameters, we quote the reduced maximum log likelihood-ratio, $`\text{RMLR}=\text{MLR}\text{DOF}`$, with DOF being the number of free parameters $`\alpha _k`$ of the sky intensity model.
#### 4.2.2 2d surface brightness distribution
As a first step we characterise the apparent morphology of the 511 keV line emission on the sky using a 2d angular Gaussian surface brightness distribution for which we determined the centroid, $`l_0,b_0`$, the longitude and latitude extent, $`\mathrm{\Delta }l,\mathrm{\Delta }b`$, and the 511 keV line flux. The results of this analysis are summarised in Table 1, the best fitting model intensity distribution is shown in Fig. 8.
The analysis confirms our earlier findings (Jean et al. 2003a ; Knödlseder et al. 2004a ; Weidenspointner et al. weidenspointner04 (2004)) of a compact and symmetric 511 keV line emission distribution towards the galactic centre. The centroid of the emission appears slightly offset from the galactic centre direction, at the statistical $`2\sigma `$ level, but we do not claim that this offset is significant. From our earlier analyses we learned that the centroid can be shifted by this amount simply from the combined effect of statistical and systematic biases in the modelling of the instrumental background.
Within the statistical uncertainties, the emission appears fully symmetric, with an extension of $`8^{}`$ (FWHM). Formally, we determine a marginal emission flattening of $`\mathrm{\Delta }b/\mathrm{\Delta }l=0.89\pm 0.14`$. The total 511 keV flux is $`(1.09\pm 0.04)\times 10^3`$ ph cm<sup>-2</sup>s<sup>-1</sup>, where the quoted error includes the uncertainty in the extent of the emission and the statistical and systematic measurement errors (c.f. section 2).
The RMLR of $`462.2`$ that has been obtained using the ORBIT-DETE background model converts into a formal detection significance of $`22\sigma `$. Using the DETE background model and including the systematic uncertainties results in a substantially higher detection significance of $`34\sigma `$. Neglecting systematic uncertainties would even boost the detection significance towards $`49\sigma `$.
#### 4.2.3 Galactic models
To determine the galactic positron-electron annihilation rate requires modelling the spatial distribution of the positron-electron annihilation. The 511 keV photon luminosity L<sub>511</sub> is related to the positron luminosity L<sub>p</sub> through $`\text{L}\text{511}=(21.5f_\mathrm{p})\times \text{L}\text{p}`$ where $`f_\mathrm{p}`$ is the positronium (Ps) fraction, defined as the fraction of positrons that decay via positronium formation (Brown, & Leventhal brown87 (1987)). Using $`f_\mathrm{p}=0.93\pm 0.04`$ that has been determined from OSSE observations (Kinzer et al. kinzer01 (2001)) results in a conversion from 511 keV photon luminosity to a positron-electron annihilation rate of $`\text{L}\text{p}=(1.64\pm 0.06)\times \text{L}\text{511}`$.
We here compare models of bulge, disk, and halo components with the data. Based on galactic model density distributions $`\rho (x,y,z)`$ we calculate the expected all-sky 511 keV intensity $`f(l,b)`$ towards direction $`(l,b)`$ by integrating the volume emissivity $`\rho (x,y,z)`$ along the line of sight $`s`$:
$$f(l,b)=\frac{1}{4\pi }\rho (x,y,z)ds$$
(7)
(the galactic centre has been assumed to be at a distance of $`R_{\mathrm{}}=8.5`$ kpc). Galactic 511 keV photon luminosities are calculated by integrating $`\rho (x,y,z)`$ over the galactic volume,
$$\text{L}\text{p}=\rho (x,y,z)s^2dsd\mathrm{\Omega }$$
(8)
assuming an outer Galaxy radius of $`R_{\mathrm{max}}=15`$ kpc.
Since the 511 keV line emission is primarily arising from the galactic centre region we fitted in a first step models of the galactic stellar bulge to the data. To account for uncertainties in our knowledge about the morphology of this component (which are related to our location in the galactic plane amid the obscuration by interstellar dust) we compared a variety of proposed bulge models to the data. The models were gathered from Dwek et al. (dwek95 (1995)) and Freudenreich (freudenreich98 (1998)), who modelled the distribution of K and M giant stars using DIRBE near-infrared skymaps, and from Picaud & Robin (picaud04 (2004)) who analysed data from the DENIS near-infrared survey. There is an accumulating body of evidence that the stellar distribution in the bulge is bar-shaped, and except for models G0 and E0, all employed bulge models have triaxial morphologies that differ in the orientation angles, the scale lengths, and the radial density profiles. Details of the models are given in Appendix A, the results of the analysis are summarised in Table 2, and best fitting 511 keV intensity distributions are shown in Fig. 8.
The best fitting bulge models are E3, G3, S<sub>F</sub>, and E<sub>F</sub>. Reasonably good fits are also obtained for P<sub>F</sub> and S<sub>PR</sub>, while only moderate fits are achieved for the remaining models. Our ranking is similar to that established from the analysis of the DIRBE and DENIS near-infrared data (Dwek et al. dwek95 (1995); Freudenreich freudenreich98 (1998); Picaud & Robin picaud04 (2004)). The best fitting bulge models fit the data as well as the adjusted 2d angular Gaussian surface brightness distribution. This means that models of the galactic stellar bulge are able to explain satisfactorily the morphology of the 511 keV bulge emission.
In a second step we fitted the 511 keV emission using parametric models of the galactic bulge and halo morphology in order to determine the scale of the emission. For the bulge models G0’ and E0’ we adjust the radial scale length ($`R_0`$) and vertical scale height ($`z_0`$), while for the galactic halo model H’ we determine the density slope powerlaw index ($`n`$), the inner cutoff radius ($`a_\mathrm{c}`$), and the axis ratio ($`ϵ`$). In addition, we employed a model composed of a set of galactocentric nested shells of constant density (model ‘Shells’) to determine the radial density profile of the 511 keV emission. We varied the radii of the shells and the number of shells in order to maximise the MLR, whilst limiting the number of shells to the minimum required to satisfactorily describe the data.
The data suggest a symmetric bulge emission profile, with scale lengths between $`300`$ and $`600`$ pc. The RMLRs are comparable to the best fitting bulge models that we tested before. The data are equally well fitted by a model of the galactic halo, with a density powerlaw index of $`n=3.0\pm 0.3`$, an inner cutoff radius of $`a_\mathrm{c}=0.39\pm 0.08`$ kpc, and a flattening of $`ϵ=0.81\pm 0.12`$. Most studies of the stellar halo population suggest power indices between $`2.4`$ and $`3.5`$ and flattenings in the range $`0.6`$ to $`1.0`$, while the inner cutoff radius is basically undetermined (Robin et al. robin00 (2000) and references therein). Our values are compatible with those of the stellar halo population, but the large uncertainties in the stellar halo morphology do not allow firm conclusions to be drawn.
The nested shell model provides the best fit to the data thanks to its flexibility in adjusting the radial density profile of the emission. A satisfactory fit is achieved by using two shells with radii $`00.5`$ and $`0.51.5`$ kpc; splitting up these shells in a finer binning, moving the shell interface radius or adding more shells does not significantly improve the fit. In particular, we detect no significant 511 keV bulge emission from galactocentric distances $`1.5`$ kpc. The radial dependence of the 511 keV volume emissivity is plotted in Fig. 7. For illustration we added the result of a third shell to the figure that covers radial distances of $`1.53.0`$ kpc and for which the flux is consistent with zero. Our fit reveals a drop in the annihilation emissivity by one order of magnitude between the inner $`00.5`$ kpc and the outer $`0.51.5`$ kpc shell, confirming the existence of a narrow core plus an extended halo of 511 keV emission that has already been suggested by the imaging analysis (c.f. section 4.1).
In a third step we added galactic disk components to the bulge and halo models. For the galactic disk we tested models of young (model D0) and old (model D1) stellar populations (Robin et al. robin03 (2003)). With both models we find clear evidence for 511 keV line emission from the galactic disk. Adding disk models D0 and D1 to bulge or halo models consistently improves the fit leading to a detection of the disk emission at the $`34\sigma `$ level.<sup>3</sup><sup>3</sup>3 Using the quoted RMLRs, the formal significance of the disk emission amounts only to $`23\sigma `$. However, the tabulated RMLRs have been obtained using the ORBIT-DETE background model which is less sensitive to 511 keV line emission than the DETE model. Using the procedure outlined in section 3 we reduce the flux uncertainties and increase the detection significance to $`34\sigma `$. Formally, D1 provides a better fit than D0, but the difference is marginal. The signal from the disk is still too faint in our present dataset to deduce anything about its morphology.
The flux, luminosity and annihilation rate in the bulge, halo and disk components are summarised in Table 3. Recall that the bulge and halo components are alternatives and their contribution should not be added to derive the total galactic values. Either component provides an almost equally good fit to the data. Due to their degeneracy fitting both simultaneously is not meaningful.
The halo model leads to a considerably larger flux, luminosity and annihilation rate than the bulge model due to the presence of a flat and extended tail in this distribution (c.f. Fig. 8). Currently, our data do not allow to detect this tail, and thus, they do not allow to discriminate between bulge and halo models. Future deep observations at intermediate galactic latitudes that are scheduled for the INTEGRAL AO-3 observing period aim in measuring this emission tail, promising to provide constraints that will allow in the future to disentangle between the different emission morphologies.
The data suggest bulge-to-disk 511 keV flux ratios in the range $`13`$, where the lower boundary is obtained for the short scale-length old stellar disk model D1 which suggests larger disk flux values than the young stellar disk model D0. Halo-to-disk 511 keV flux ratios are even larger, in the range $`24`$, owing to the larger flux in the halo component. The large uncertainty in these ratios arises from the low intensity of the galactic disk component, which for the analysed dataset is just above the SPI detection limit.
We also note that the bulge-to-disk 511 keV photon luminosity ratio is much higher than the bulge-to-disk flux ratio and lies in the range $`39`$. This difference is explained by the fact that the average squared distance $`\overline{s}^2=\rho (s)s^2dsd\mathrm{\Omega }/\rho (s)dsd\mathrm{\Omega }`$, which defines the distance at which a source of luminosity L<sub>p</sub> produces the observed 511 keV line flux, is smaller for the galactic disk than for the galactic bulge. In other words, to produce the same 511 keV flux at Earth, the intrinsic luminosity of the bulge has to be larger than that of the disk.<sup>4</sup><sup>4</sup>4 The difference between bulge-to-disk flux and luminosity ratio is only important for our home Galaxy and is related to the fact that the Sun is located within the galactic radius. For external galaxies this difference disappears since their bulge and disk appear at the same distance to us. It is therefore important to quote explicitly the quantity for which we discuss the bulge-to-disk ratio. The same rational also holds for the halo-to-disk 511 keV photon luminosity ratio, which is larger than the corresponding flux ratio, and which is comprised between $`413`$.
### 4.3 Correlation with tracer maps
To gain insight into the nature of the galactic positron sources, we searched for correlations between the 511 keV line emission morphology and all-sky intensity distributions observed at other wavelengths. This work was inspired by a similar study that Knödlseder et al. (1999b ) performed to understand the morphology of the 1.8 MeV gamma-ray line emission (arising from radioactive decay of <sup>26</sup>Al) observed by the COMPTEL telescope aboard CGRO. Through their analysis, the authors could establish a tight correlation between the morphology of galactic microwave free-free emission and that of 1.8 MeV line emission, hinting towards a massive star origin of <sup>26</sup>Al. The tracer maps used for the comparison are those listed in Knödlseder et al. (1999b ). For a detailed description of the maps the reader is referred to that work.
Figure 9 summarises the result of the correlation study for our 511 keV dataset, where we show the RMLR as a function of the tracer map (ordered by increasing photon energy or decreasing wavelength). None of the tracer maps is consistent with the data. The maximum RMLR that is reached ($`353.7`$ for the DIRBE 2.2 $`\mu `$m map) is more than $`100`$ units smaller than the values obtained for the parametric models of the previous sections. Apparently, the 511 keV emission morphology is unique and cannot be represented by any known celestial intensity distribution.
Nevertheless, Fig. 9 shows a clear trend, where the data favour maps in the near-infrared domain (DIRBE 1.25 $`\mu `$m \- 4.9 $`\mu `$m) and the hard X-ray band (HEAO-1) over maps observed at longer wavelengths. In particular, the worst fits are obtained in the microwave and far-infrared domain where the skymaps trace the young stellar population, either through their ionising radiation (DMR maps at $`\nu 53`$ GHz), or through their related molecular gas (CO) and cold dust emission (DIRBE 100 $`\mu `$m \- 240 $`\mu `$m). From this it is clear that the bulk of the 511 keV emission is not related to the young massive stellar population of the Galaxy.
On the contrary, all best fitting tracers maps show the characteristic features of an old stellar population: a strong bulge component combined with a short scale radius disk component. Apparently, the 511 keV data tend to favour such morphologies. This is illustrated in Fig. 10 where we plot the RMLR as function of the bulge-to-disk flux ratio of the tracer map, defined as the flux contained within a circular region of $`6^{}`$ in radius around the galactic centre, divided by the flux within galactic latitudes $`b=\pm 20^{}`$ outside the circular bulge region (note that the precise B/D value depends of course on the exact definitions chosen for the two regions, but our purpose here is to illustrate a trend). Clearly, there is a strong correlation between the B/D flux ratio and the RMLR, in the sense that the larger the B/D flux ratio, the larger the RMLR. In particular, the DIRBE near-infrared maps and the HEAO-1 map show the largest B/D flux ratios of all tracer maps ($`0.2`$). Thus, finding a population of objects which show a large B/D ratio could provide the key for finding the galactic source of positrons.
### 4.4 Point-source search
The modest angular resolution of SPI of about $`3^{}`$ (Vedrenne et al. vedrenne03 (2003)) makes it difficult to distinguish between point source, point-like, and small-scale diffuse emission. So in principle we cannot exclude the possibility that the 511 keV gamma-ray line emission that is seen towards the galactic bulge region is made of a limited number of point sources that blend to simulate diffuse emission. In due course the question of any point source contribution to the flux will be best addressed using data from the imager IBIS on INTEGRAL, which is relatively insensitive to diffuse emission, in conjunction with that from SPI. Such work is underway and will be reported on separately. We here limit ourselves to the constraints which can be placed from SPI data alone on such a contribution.
We have therefore used the SPIROS algorithm (Skinner & Connell skinner03 (2003)) to search our dataset for the positions and fluxes of point sources that are compatible with our data. SPIROS searches for the most probable position for a point source and fits a source at that position before repeating iteratively the search using the residuals after sources already found are taken into account. At each iteration the positions and fluxes of all sources that have been found are optimised by maximising a goodness of fit parameter (the $`\chi ^2`$ statistic was used here).
If the 511 keV emission is intrinsically diffuse then application of this algorithm will lead to sources being placed at selected positions (regions of high flux and local noise peaks) until a distribution of emission is found that is consistent (given limited statistics) with the data. Such a description is unlikely, however, to be unique and if most or all of the flux is from diffuse emission then the particular source positions found will have no astrophysical significance. We therefore do not present here the detailed results of this blind source search and we restrict ourselves to discussion of the general conclusions which can be drawn from the analysis.
For models with 7 point sources based on the SPIROS solution after iteration 7 we obtain $`\text{RMLR}=462.5`$, slightly inferior to our best fitting diffuse models. Subsequent iterations suggest point sources at the edge of the exposed regions which are obviously spurious. Only after iteration 13 is another point source found in the galactic bulge region. Fitting this source together with the 7 sources found earlier leads to $`\text{RMLR}=471.9`$, comparable to the best fitting diffuse models. The total flux attributed to the 8 point sources is $`1.1\times 10^3`$ ph cm<sup>-2</sup>s<sup>-1</sup>, comparable with the values obtained for the bulge component using our best fitting diffuse models (c.f. Table 2). We therefore simply conclude that at least 8 point sources would be needed to satisfactorily describe the SPI data.<sup>5</sup><sup>5</sup>5 Given the angular resolution of SPI, compact sources with an extent $`1^{}2^{}`$ would be considered as point sources in this context.
In addition to the blind search for point sources, we also looked for evidence of 511 keV gamma-ray line emission from a list of potential candidate objects. Our list comprises compact objects, pulsars, supernova remnants, star forming regions, globular clusters, nearby (active) galaxies, and galaxy clusters. Depending on the expected source extent, we searched for either point source emission or extended emission, modelled by a 2d angular Gaussian surface brightness distribution for which we specify the centroid and the FWHM extension.
The results of the analysis are summarised in Table 4. None of the sources we searched for showed a significant 511 keV flux, hence we only quote ($`3\sigma `$) upper limits in Table 4. Since the emission of the Galaxy may interfere with the emission from the specific sources (due to the large field of view of the SPI instrument), we also included models for the diffuse galactic 511 keV emission for the source search. We have used combinations Shells+D0, Shells+D1, H’+D0 and H’+D1 to cover the range of plausible best fitting diffuse models (c.f. Fig. 8) and quote always the most conservative flux limit. The (less sensitive) ORBIT-DETE background model has been used to ensure that systematic uncertainties are negligible.
## 5 Discussion
### 5.1 Comparison with earlier measurements
Observations of the 511 keV line emission have been made by a large number of balloon and satellite borne telescopes, yet only a few of them provided constraining information on the emission morphology. The OSSE instrument that flew during 1991-2000 on-board CGRO accumulated so far the largest database for studying the 511 keV line intensity distribution. The observations of the Gamma-Ray Spectrometer on-board the Solar Maximum Mission (SMM) (1980-1988) and of the Transient Gamma-Ray Spectrometer (TGRS) on-board the WIND mission (1995-1997) have also been used to estimate the overall 511 keV line flux and maximum emission size (Purcell et al. purcell94 (1994); Kinzer et al. kinzer96 (1996); Tueller et al. tueller96 (1996); Cheng et al. cheng97 (1997); Purcell et al. purcell97 (1997); Harris et al. harris98 (1998); Milne et al. milne00 (2000); Kinzer et al. kinzer01 (2001)).
The picture that emerged prior to the INTEGRAL launch was the following. From the general trend that instruments with larger fields-of-view show larger fluxes it was inferred that the 511 keV emission is extended. OSSE observations strongly exclude a single point-source located at the GC (Purcell et al. purcell94 (1994)). The OSSE observations suggest at least two emission components, one being a spheroidal bulge and the other being a galactic disk component. A third component, named the Positive Latitude Enhancement (PLE), situated about $`9^{}12^{}`$ above the GC has been reported (Cheng et al. cheng97 (1997); Purcell et al. purcell97 (1997)), but the morphology and intensity of this component was in fact only poorly determined by the data (Milne et al. milne00 (2000)).
The emission, which showed no significant offset from the GC, was well fitted by either a model comprising a narrow ($`56^{}`$ FWHM) Gaussian bulge plus $`35^{}`$ FWHM Gaussian and CO-like disk components, or by a centre-truncated R<sup>1/4</sup> spheroid plus exponential disk model (Purcell et al. purcell97 (1997); Milne et al. milne00 (2000); Kinzer et al. kinzer01 (2001)). The total 511 keV gamma-ray line flux was estimated to be $`(13)\times 10^3`$ ph cm<sup>-2</sup>s<sup>-1</sup>. The distribution of flux between the bulge and disk components was only weakly constrained by the observations, and depended sensitively on the assumed bulge shape. In particular, estimates for the bulge-to-disk (B/D) flux ratio varied from $`0.23.3`$ depending upon whether the bulge component features a halo (which leads to a large B/D ratio) or not (Milne et al. milne00 (2000)).
Our analysis basically confirms the pre-INTEGRAL observations (c.f. Table 5). One difference is that the bulge appears slightly larger in our analysis when compared to the OSSE result. We note that OSSE performed differential measurements using its $`4^{}\times 11^{}`$ collimator which may bias the results towards small values (Kinzer et al. kinzer01 (2001)), but in any case, the discrepancy is not very significant and is not surprising in view of possible systematic uncertainties.
Another difference with respect to OSSE is that we find no evidence for a feature resembling the PLE. Fitting a model of the PLE (2d Gaussian of $`5^{}`$ FWHM located at $`l=2^{}`$ and $`b=8^{}`$) on top of the bulge results in a $`3\sigma `$ upper flux limit of $`1.5\times 10^4`$ ph cm<sup>-2</sup>s<sup>-1</sup> for the PLE. The OSSE team has gradually reduced their estimates of the flux and significance attributed to this emission feature from $`5\times 10^4`$ ph cm<sup>-2</sup>s<sup>-1</sup> (Purcell et al. purcell97 (1997)) down to an upper limit of $`1\times 10^4`$ ph cm<sup>-2</sup>s<sup>-1</sup> (Milne et al. milne01 (2001)). Recently it has been suggested that data analysis problems linked with variable continuum emission may account for the reported PLE (Milne milne04 (2004)), so perhaps our non-detection of a PLE feature is not surprising.
Until now, there have been very few published upper limits on 511 keV gamma-ray line emission from point sources which take account of diffuse emission. Examples are the $`3\sigma `$ limits of $`1.6\times 10^4`$ ph cm<sup>-2</sup>s<sup>-1</sup> for 1E 1740.7-2942 (Purcell et al. purcell94 (1994)) and $`1.4\times 10^4`$ ph cm<sup>-2</sup>s<sup>-1</sup> for inner Galaxy sources (Milne et al. milne01 (2001)). Our upper limits, summarised in Table 4, are somewhat more stringent.
Finally, we want to mention that the method of analysis used in this work assumes that the 511 keV line emission is not time variable. From our analysis of SPI data alone we have no indications for time variability. Furthermore, OSSE and TGRS measurements revealed no significant time variability (Purcell et al. purcell97 (1997); Harris et al. harris98 (1998)) and in addition, our 511 keV line flux measurements are consistent with those of OSSE and TGRS. Thus we believe that our assumption of non-variable 511 keV line emission is reasonable. We have, however, not yet performed a thorough analysis on all relevant timescales.
### 5.2 General considerations
The most distinctive morphological feature of the 511 keV emission is the large B/D luminosity ratio of $`39`$. Unless there is a mechanism that strongly suppresses positron annihilation in the galactic disk, or that somehow transports positrons from the disk into the galactic bulge or halo where they annihilate, the positron source population we are seeking for should also exhibit such a high B/D ratio.
The B/D luminosity ratio of $`39`$ is considerably larger than the B/D mass ratio of $`0.31.0`$ of our Galaxy (e.g. Caldwell & Ostriker caldwell81 (1981); Freudenreich freudenreich98 (1998); Bissantz & Gerhard bissantz02 (2002); Robin et al. robin03 (2003)). The uncertainty in the galactic B/D mass ratio is partly due to differences in the modelling of the disk component, where disk profiles exhibiting a central hole or depletion lead to B/D ratios at the high end, while double exponential profiles without hole favour B/D ratios at the low end. Since we employed in our analysis disk models with central holes from Robin et al. (robin03 (2003)), we should for consistency compare our 511 keV B/D luminosity ratio to their (large) B/D ratio of $`1`$.<sup>6</sup><sup>6</sup>6 We note that, independently, large B/D ratios are also favoured by microlensing surveys towards the galactic bulge region (Binney & Evans binney01 (2001)). But even with such large B/D mass ratios, the source population we are seeking for should still be at least 3 times more abundant in the bulge than in the disk of the Galaxy. We therefore conclude that the primary positron source of the Galaxy is clearly associated with the galactic bulge. It therefore should belong to the old stellar population.
Furthermore, the fact that the 511 keV emission matches well the morphology of the stellar bulge suggests that positron diffusion probably plays only a minor role. Were positron diffusion to be important we would expect to find substantial 511 keV emission in gas-rich regions adjacent to the rather gas-poor galactic bulge, such as the molecular ring structure at galactocentric distances of $`4`$ kpc. However, we do not find any evidence for 511 keV emission correlated with this structure. We therefore conclude that positron diffusion is negligible at galactic scales (i.e. kpc scales).
### 5.3 Constraints on the disk source
One certain source of positrons in the disk of the Galaxy is the radioisotope <sup>26</sup>Al. It decays with a lifetime of $`\tau 10^6`$ yr with emission of a 1809 keV gamma-ray photon; $`85\%`$ of the decays are also accompanied by the emission of a positron. The galactic distribution of <sup>26</sup>Al is well known thanks to observations of the COMPTEL telescope aboard CGRO, and follows that of the young stellar population. Thus, under the assumption that the positrons annihilate close to their production site, 511 keV line emission along the galactic plane is expected, showing the morphological characteristics of a young stellar population.
The expected 511 keV line flux $`F_{511}`$ due to <sup>26</sup>Al decay is related to the 1809 keV line flux $`F_{1809}`$ through $`F_{511}=0.85\times (21.5f_\mathrm{p})\times F_{1809}`$. Using the COMPTEL measurement of the 1809 keV flux along the galactic plane, $`F_{1809}=9\times 10^4`$ ph cm<sup>-2</sup>s<sup>-1</sup> (c.f. Table 4.3 in Knödlseder knoedl97 (1997)), and assuming $`f_\mathrm{p}=0.93`$ leads to an expected 511 keV line flux of $`F_{511}=5\times 10^4`$ ph cm<sup>-2</sup>s<sup>-1</sup>.
Fitting our model of the young stellar population (model D0), together with bulge models to the data suggests a disk flux in the range $`(46)\times 10^4`$ ph cm<sup>-2</sup>s<sup>-1</sup>. To explore the sensitivity of the disk flux on the assumed disk model we also paired the shell model with tracers of 1809 keV line emission, such as the DMR free-free and the DIRBE 240 $`\mu `$m emission maps (Knödlseder et al. 1999b ). This resulted in slightly larger (and more significant) disk fluxes of $`(8.3\pm 2.3)\times 10^4`$ ph cm<sup>-2</sup>s<sup>-1</sup>. Comparing these values to what we expect from <sup>26</sup>Al suggests that $`60100\%`$ of the galactic plane emission may be attributed to $`\beta ^+`$-decay of <sup>26</sup>Al.
If this contribution is subtracted a 511 keV disk flux of at most $`3\times 10^4`$ ph cm<sup>-2</sup>s<sup>-1</sup> remains to be explained by other positron sources. A possible candidate is the radioisotope <sup>44</sup>Ti whose daughter isotope <sup>44</sup>Sc decays via $`\beta ^+`$-decay into stable <sup>44</sup>Ca ($`\tau 87`$ yr). In contrast to <sup>26</sup>Al there is no firm measurement of the present day galactic <sup>44</sup>Ti mass (nor of its spatial distribution), but simple chemical evolution arguments lead to the expectation that about $`4\times 10^6`$ $`M_{\mathrm{}}`$ of <sup>44</sup>Ti are produced per year (Leising & Share leising94 (1994)). Under the assumption that all positrons escape the production site this yield translates into an annihilation rate of $`3\times 10^{42}`$ s<sup>-1</sup>. Assuming further that <sup>44</sup>Ti is distributed following model D0, a 511 keV disk flux of $`8\times 10^4`$ ph cm<sup>-2</sup>s<sup>-1</sup> is expected. In view of the approximate estimation and in particular in view of the uncertainty about the spatial distribution this value seems in reasonable agreement with the remaining 511 keV flux of $`3\times 10^4`$ ph cm<sup>-2</sup>s<sup>-1</sup>.
It is intriguing that the galactic disk flux could be entirely explained by the radioactive decay of <sup>26</sup>Al and <sup>44</sup>Ti. This would suggest that once the <sup>26</sup>Al and <sup>44</sup>Ti contributions have been subtracted only the bright bulge component of 511 keV emission remains, which would then demand a specific source population that is only confined to the inner Galaxy. However, we cannot immediately draw this conclusion. Fitting an old stellar population disk model, with larger scale height and smaller radial scale than the young one, increases the estimate of the disk flux by about a factor of two, leaving room for a weak disk component not associated with <sup>26</sup>Al or <sup>44</sup>Ti. This would suggest lower limits on the B/D flux (luminosity) ratio of $`2`$ ($`6`$) for the source population that gives rise to the galactic bulge emission.
### 5.4 Constraints on the bulge source
#### 5.4.1 Massive stars, core collapse supernovae, pulsars
The bulge dominance of the 511 keV emission immediately excludes scenarios in which the bulk of galactic positron production is related to massive stars. Such scenarios include the production of $`\beta ^+`$-decay radioisotopes produced by Wolf-Rayet stars and all types of core collapse supernovae (including hypernovae) and the pair production in the strong magnetic fields of pulsars. Massive stars may well explain the faint disk component of 511 keV emission via the radioactive decay of <sup>26</sup>Al and <sup>44</sup>Ti (c.f. section 5.3). They cannot, however, explain the majority of the emission, which would in that case resemble the 1809 keV line emission.
#### 5.4.2 Hypernovae
Cassé et al. (casse04 (2004)) proposed that a recent hypernova at the galactic centre could be responsible for the observed positron emission, but there is no observational evidence that such an explosion indeed took place. Hypernovae are believed to be related to Wolf-Rayet stars, which are distributed in the galactic disk following the spiral arm pattern, hence it would be much more likely to find a recent hypernova at an arbitrary position along the galactic plane (or at the tangent points of the spiral arms) rather than at the position of the galactic centre. And even in the rare event of a hypernova exploding right at the galactic centre, it would be difficult to explain why the resulting 511 keV annihilation radiation (arising from the $`\beta ^+`$-decay of freshly synthesised <sup>56</sup>Co) should reflect the stellar morphology of the galactic bulge. We therefore conclude that it is unlikely that galactic centre hypernovae are the source of the bulge positrons.
#### 5.4.3 Cosmic-ray interactions
Interactions of cosmic-ray particles with the ambient interstellar medium may produce positrons, primarily via the $`\mathrm{N}+\mathrm{p}\pi ^+\mathrm{e}^+`$ reaction channel (N stands for nucleus). The effect of the cosmic-ray interaction is best seen in the GeV gamma-ray domain, and has been comprehensively mapped by the EGRET satellite aboard CGRO. The EGRET all-sky map shows dominant emission from the galactic plane, which follows a linear combination of various gas and dust tracers in the galaxy. Thus cosmic-ray interactions should also lead to disk dominated 511 keV emission, which is at odds with our observations.
#### 5.4.4 X-ray binaries
Positron production in X-ray binaries may occur either as a result of $`\gamma \gamma `$ pair production in the luminous compact region around the compact object or due to nuclear interactions that may form excited nuclei that subsequently decay through the emission of positrons. Galactic black holes and microquasars, where the positrons are ejected in a relativistic jet with Lorentz factors of a few (Dermer & Murphy dermer01 (2001)), are the two leading candidates.
X-ray binaries are separated into two classes, depending on whether the donor is a high-mass star (HMXB) or a low-mass star (LMXB). The two classes show clearly different spatial distributions (Grimm et al. grimm02 (2002)). HMXB are associated with the young stellar population and are primarily found in the galactic disk. Consequently they can immediately be excluded as the source of the bulge positrons. LMXB, in contrast, are strongly concentrated towards the galactic bulge, and are more promising source candidates. Among the 150 LMXBs listed in the catalogue of Liu et al. (liu01 (2001)), more than $`50\%`$ are observed towards the galactic bulge. Correcting for completeness, Grimm et al. (grimm02 (2002)) find a B/D ratio of $`0.9`$ and a vertical scale height of 410 pc for the LMXB distribution.
Formally, the LMXB B/D ratio is considerably below the value required by our 511 keV data, yet the large vertical scale height of LMXB could lead to a scenario where a substantial fraction of positrons from disk LMXB may escape into the galactic halo. This scenario works as follows. Since the scale height of LMXB (410 pc) considerably exceeds the scale height of the dense interstellar gas layer of our Galaxy ($`100`$ pc), positrons from disk LMXB are ejected into rather low-density regions, typically a factor of $`10100`$ less dense than regions found near the galactic plane (Ferrière ferriere98 (1998)). Before positrons can annihilate they have to slow down considerably, mostly through Coulomb interactions, with a characteristic timescale of $`\tau _{\mathrm{SD}}10^5n^1`$ yr, $`n`$ being the ISM density in units of cm<sup>-3</sup> (Forman et al. forman86 (1986)). Consequently, positrons live $`10100`$ times longer at large scale heights than near the galactic plane, allowing for substantial diffusion before annihilation takes place. The typical diffusion length depends much on the magnetic field configuration and the amount of ISM turbulence at large scale heights, but qualitatively it seems plausible that a considerable fraction of the disk positrons may annihilate in the galactic halo.
The resulting broad diffuse component of 511 keV emission would be difficult to detect with SPI. In particular, the present dataset, for which good exposure is restricted to a band $`b\pm 15^{}`$ along the galactic plane (c.f. Fig. 1), makes it virtually impossible to measure a disk component with a broad-latitude distribution of 511 keV emission. Therefore it would be sufficient that $`2/3`$ of the positrons produced by disk LMXB escape into the galactic halo to reconcile the LMXB distribution with a B/D ratio of 3. If we require a more extreme B/D ratio of 6, as expected after subtraction of the <sup>26</sup>Al component from the disk (c.f. section 5.3), a positron escape fraction of $`80\%`$ would be needed for disk LMXB. It remains to be seen whether such large escape fractions are feasible.
An alternative way to test the LMXB scenario is to search for 511 keV line emission from individual bright and/or nearby objects. So far no emission is seen towards the interesting candidates Sco X-1 (the brightest LMXB) and Cen X-4 and A0620-00 (probably the most nearby LMXB at $`1.2`$ kpc), but we plan for deep observations of these objects in the near future to search for their annihilation signatures. The detection of positron annihilation signatures from nearby objects could however be hampered by (even modest) positron diffusion away from the sources, which would lead to extended 511 keV emission halos around the objects. So even for the modest angular resolution of SPI of $`3^{}`$ nearby individual LMXB could appear as extended sources, and their low surface brightness could make their detection more difficult.
#### 5.4.5 Classical novae
Among all proposed positron candidate sources, classical novae, i.e. thermonuclear runaways on white dwarfs in accreting binary systems, are the sources for which the largest B/D ratios of $`34`$ have been suggested (Della Valle & Duerbeck dellavalle93 (1993); Della Valle & Livio 1994a ). Interstellar extinction, in particular towards the galactic bulge region, makes it virtually impossible to derive their spatial distribution in the Milky Way directly, but novae are readily observed in nearby external galaxies which may serve as templates (e.g. Shafter et al. shafter00 (2000)). Due to its proximity and due to its similarity to the Milky Way, M31 is the primary source of information, and modern investigations indicate that novae reside primarily in the bulge region of M31 (Ciardullo et al. ciardullo87 (1987); Capaccioli et al. capaccioli89 (1989)). Although it had been suggested that selection effects may have “faked” such a finding (Hatano et al. hatano97 (1997)) the recent study of Shafter & Irby (shafter01 (2001)) demonstrates that such biases, if they exist, must be small.
Novae produce positrons via the $`\beta ^+`$-decay of radioactive isotopes synthesised during the thermonuclear runaway, mainly of <sup>13</sup>N, <sup>18</sup>F, and <sup>22</sup>Na (lifetimes $`\tau =14`$ min, $`2.6`$ hr, and $`3.75`$ yr, respectively). <sup>22</sup>Na yields of $`6\times 10^9`$ $`M_{\mathrm{}}`$, as suggested by theoretical nucleosynthesis calculations for ONe novae (Hernanz et al. hernanz02 (2002)), would require nova rates of $`1600`$ yr<sup>-1</sup> to maintain positron production and annihilation in an equilibrium state, a rate which is considerably above estimates of $`35\pm 11`$ yr<sup>-1</sup> for all types of novae in the entire Galaxy (Shafter shafter97 (1997)). <sup>13</sup>N yields of $`2\times 10^7`$ $`M_{\mathrm{}}`$ for low-mass CO novae are more promising (Hernanz et al. hernanz02 (2002)) since they would require nova rates of only $`26`$ yr<sup>-1</sup> if all positrons could indeed escape from the nova envelope into the ISM. However, with a <sup>13</sup>N lifetime of $`14`$ min it seems unlikely that this would be possible.
It is probable that large fractions of the <sup>13</sup>N positrons annihilate within the dense nova envelope, leading to prompt annihilation that could give rise to transient annihilation signatures (Leising & Clayton leising87 (1987); Gómez-Gomar et al. gomez98 (1998)). This signature has been sought using various gamma-ray telescopes, but has so far eluded detection (see Hernanz & José hernanz04 (2004) and references therein). Detection of the transient signature may help to shed light on the positron escape fraction, and could show whether novae contribute to the galactic bulge positron budget or not.
#### 5.4.6 Thermonuclear supernovae
In view of their potential to produce large numbers of positrons, thermonuclear Type Ia supernovae (SN Ia) are often considered as the most plausible source of positrons in the Milky Way (Dermer & Murphy dermer01 (2001)). SN Ia produce positrons via the $`\beta ^+`$-decay of radioactive <sup>56</sup>Co ($`\tau =111`$ days). Expected <sup>56</sup>Co yields of $`0.6`$ $`M_{\mathrm{}}`$ provide $`2.5\times 10^{54}`$ positrons per event, although, as with novae, prompt annihilation in the supernova envelope probably prevents large fractions of the positrons from escaping into the ISM. From the analysis of late light curves of SN Ia Milne et al. (milne99 (1999)) derive a mean escaped positron yield of $`8\times 10^{52}`$ positrons per SN Ia, corresponding to a positron escape fraction of $`f0.03`$. A recent study of SN 2000cx even suggests $`f0`$, but SN 2000cx was an unusual event that may not represent the average SN in the bulge of our Galaxy (Sollerman et al. sollerman04 (2004)).
Assuming therefore $`f=0.03`$ a bulge SN Ia rate of $`0.6`$ per century is required to maintain the observed 511 keV luminosity in a steady state. Unfortunately the rate and distribution of SN Ia in our Galaxy are only poorly known. The galactic SN Ia rate is generally inferred from rates observed in external galaxies which are then scaled to the mass and the type of the Milky Way. In that way rates of $`0.31.1`$ SN Ia per century are derived (Tammann et al. tammann94 (1994); Cappellaro et al. cappellaro97 (1997); Mannucci et al. mannucci05 (2005)), sufficient to maintain the galactic 511 keV luminosity. In contrast, when we follow the suggestion of Prantzos (prantzos04 (2004)) and derive the bulge SN Ia rate by scaling the SN Ia rate observed in early-type galaxies to the mass of the galactic bulge, a bulge SN Ia rate of $`0.08`$ SN Ia per century is obtained. This value is much too low to explain the observed bulge 511 keV luminosity. It is difficult to judge if the galactic bulge can indeed be considered as a downsized version of an elliptical galaxy, in particular in view of the differences in the evolution of the galactic bulge and an elliptical galaxy. Furthermore, it is suggested that different SN Ia explosion mechanisms exist in different types of galaxy (e.g. Della Valle & Livio 1994b ; Howell howell01 (2001); Mannucci et al. mannucci05 (2005)) making the proposed extrapolation even more uncertain.
Observations of external galaxies indicate that SN Ia distributions are strongly peaked towards galactic centres (Bartunov et al. bartunov92 (1992)), yet reliable determinations of B/D ratios are difficult in view of observational biases and selection effects (Wang et al. wang97 (1997); Hatano et al. hatano98 (1998); Howell et al. howell00 (2000)). If the SN Ia distributions follows that of novae (both populations are believed to arise from accreting white dwarfs; see van den Bergh vandenbergh88 (1988)) one can expect B/D ratios of $`34`$. Even higher B/D ratios can be achieved if part of the positrons produced by disk SN Ia, which have a vertical scale height of $`330`$ pc (Chen et al. chen01 (2001)), escape into the halo (c.f. section 5.4.4). Thus SN Ia could indeed present the required characteristics that are needed to explain the positron distribution and annihilation rate in the Galaxy.
Alternatively, instead of explaining the bulge emission globally we may search for 511 keV emission from nearby Type Ia supernovae remnants, such as SN 1006, the Lupus Loop, or the Tycho SNR. Assuming a mean escaped positron yield of $`8\times 10^{52}`$ positrons per SN Ia (Milne et al. milne99 (1999)) and a positronium fraction of $`f_\mathrm{p}=0.93`$, the mean expected 511 keV line flux from an individual SN Ia is estimated to
$$F_{511}=1.3\times 10^4\left(\frac{1\mathrm{kpc}}{D}\right)^2\left(\frac{10^5\mathrm{yr}}{\tau }\right)\text{ph }\text{cm}\text{-2}\text{s}\text{-1}$$
(9)
where $`D`$ is the distance to the supernova remnant in units of kpc, and $`\tau `$ is the mean lifetime of the positrons in years, between $`10^310^7`$ yr, depending on the density, temperature, and ionisation state of the annihilating medium (Guessoum et al. guessoum91 (1991)). Taking distances of $`D2`$ kpc for SN 1006 (Laming et al. laming96 (1996)), of $`D1`$ kpc for the Lupus Loop (Leahy et al. leahy91 (1991)), and of $`D2.3`$ kpc for Tycho (Hughes hughes00 (2000)) and assuming a positron lifetime of $`10^5`$ yr results in predicted 511 keV line fluxes of $`0.3\times 10^4`$ ph cm<sup>-2</sup>s<sup>-1</sup>, $`1.3\times 10^4`$ ph cm<sup>-2</sup>s<sup>-1</sup>, and $`0.2\times 10^4`$ ph cm<sup>-2</sup>s<sup>-1</sup>, respectively. At least for the Lupus Loop, the predicted flux is close to our upper 511 keV flux limit, indicating that dedicated deep observations of nearby supernova remnants can help to answer the question about the galactic positron source. Such dedicated observations with INTEGRAL are already scheduled.
#### 5.4.7 Light dark matter annihilation
Light dark matter (1-100 MeV) annihilation, as suggested recently by Boehm et al. (boehm04 (2004)), is probably the most exotic but also the most exciting candidate source of galactic positrons. Unfortunately, the spatial distribution of dark matter in general, and light dark matter in particular, is only poorly constrained by observational data, at least for the inner Galaxy. The debate of whether the dark matter profile shows a cusp towards the galactic centre is still not settled, but it seems clear now that, dynamically, dark matter plays only a minor role in the inner 3 kpc of the Galaxy. In this region the stellar mass dominates (Binney & Evans binney01 (2001); Klypin et al. klypin02 (2002)).
Due to these uncertainties it is difficult to judge whether the observed 511 keV emission could be explained by dark matter annihilation. Maybe more promising is the idea to search for signatures of dark matter annihilation in nearby, external galaxies. Hooper et al. (hooper04 (2004)) suggested that nearby dwarf spheroidal galaxies may provide prominent sources of 511 keV line emission due to the high densities of dark matter that are known to be present. They proposed the nearby Sagittarius dwarf galaxy (Sgr dwarf) as most promising candidate and estimate 511 keV line fluxes in the range of $`(17)\times 10^4`$ ph cm<sup>-2</sup>s<sup>-1</sup>, depending on the assumed dark matter halo profile. Cordier et al. (cordier04 (2004)) searched for emission from this Galaxy using SPI and obtained a $`(3\sigma )`$ upper limit of $`3.8\times 10^4`$ ph cm<sup>-2</sup>s<sup>-1</sup>. Our upper limit of $`1.7\times 10^4`$ ph cm<sup>-2</sup>s<sup>-1</sup> is substantially lower, and excludes almost all types of halo models for this galaxy, in particular those with a central cusp.
Standard cold dark matter cosmology predicts cuspy dark matter distributions (Klypin et al. klypin02 (2002)), so in principle Sgr dwarf should have been detected by SPI if dark matter annihilation were a viable scenario. Maybe dark matter halos are less cuspy than theory predicts? This possibility is indeed indicated by observations of our own Galaxy (Binney & Evans binney01 (2001)) and dwarf galaxies (Blais-Ouellette et al. blais99 (1999); Kleyna et al. kleyna03 (2003)). But in this case dark matter annihilation should not lead to a compact but to a rather extended 511 keV emission feature – in contradiction to what SPI observations of the inner Galaxy suggest. From the arguments given one may question the dark matter scenario. However, it is certainly premature to reject them totally because of the uncertainties in the dark halo profiles and the annihilation conditions.
## 6 Conclusions
Our first mapping of 511 keV gamma-ray line emission over a large fraction of the celestial sphere leads us to the following observations:
1. 511 keV emission is significantly ($`50\sigma `$) detected towards the galactic bulge region, and, at a very low level ($`4\sigma `$), from the galactic disk
2. there is no evidence for a point-like source in addition to the diffuse emission, down to a typical flux limit of $`10^4`$ ph cm<sup>-2</sup>s<sup>-1</sup>
3. there is no evidence for the positive latitude enhancement that has been reported from OSSE measurements; the $`3\sigma `$ upper flux limit for this feature is $`1.5\times 10^4`$ ph cm<sup>-2</sup>s<sup>-1</sup>
4. the bulge emission is spherically symmetric and is centred on the galactic centre with an extension of $`8^{}`$ (FWHM); it is equally well described by models that represent the stellar bulge or the halo populations
5. the bulge annihilation rate is $`(1.5\pm 0.1)\times 10^{43}`$ s<sup>-1</sup>, the disk annihilation rate is $`(0.3\pm 0.2)\times 10^{43}`$ s<sup>-1</sup>
6. the bulge-to-disk luminosity ratio lies in the range $`39`$
The bulge dominated 511 keV line emission morphology suggests an old stellar population as the main galactic positron source. In contrast, the faint disk emission is well explained by the release of positrons during the radioactive decay of <sup>26</sup>Al that originated from massive stars, with a possible contribution from <sup>44</sup>Ti synthesised during supernova explosions.
The extreme bulge-to-disk ratio that is observed in the 511 keV luminosity imposes severe constraints on the principal galactic positron source. Type Ia supernovae, low-mass X-ray binaries or dark matter annihilation may possibly satisfy these constraints, but uncertainties in the knowledge about the spatial distribution of these objects and the positron escape processes prevents us from drawing firm conclusions. Novae could probably most easily explain the large B/D ratios, yet an implausibly large positron escape fraction from <sup>13</sup>N decay would be required to accommodate the observed annihilation rate. SN Ia could explain the annihilation rate for a modest positron escape fraction, but it is questionable if they have the required large B/D ratio. LMXB could reproduce the observed B/D ratio provided that a substantial fraction of positrons ejected by disk LMXB escape into the halo. Light dark matter is an exciting option, but it remains to be seen if the observed 511 keV emission distribution is compatible with the profile of the galactic dark matter halo.
Future deep observations of individual nearby candidate sources may provide the means to identify the galactic positron source. As such we will soon observe the X-ray binaries Sco X-1 and Cen X-4 with INTEGRAL, and observations of the nearby supernova remnant SN 1006 are already scheduled. We cannot be sure that any of these observations will allow the detection of a 511 keV signal, but were such a signal detected we would gain important new insights in the primary source of positrons in our Galaxy.
## Appendix A Summary of galactic density profiles
##### Bulge models:
We model the galactic bulge as a triaxial stellar bar for which the parameters are summarised in Table 6. The apparent intensity distribution on the sky is computed in the galactic frame which we define by a right-handed cartesian coordinate system where the Sun is located on the negative y-axis, at $`R_{\mathrm{}}=8.5`$ kpc. The transformation from the galactic frame into the bar frame is performed by two consecutive rotations: the first, represented by the matrix $`R_X(\alpha )`$, is a counterclockwise rotation by an angle $`\alpha `$ around the z-axis; the second, represented by the matrix $`R_Y(\beta )`$, is a clockwise rotation by an angle $`\beta `$ around the new y-axis, i.e.
$$r^{}=R_Y(\beta )R_X(\alpha )r$$
(10)
The effective bar radius is defined by
$$R_\mathrm{s}=\left(\left[\left(\frac{|x^{}|}{a_\mathrm{x}}\right)^C_{}+\left(\frac{|y^{}|}{a_\mathrm{y}}\right)^C_{}\right]^{\frac{C_{}}{C_{}}}+\left(\frac{|z^{}|}{a_\mathrm{z}}\right)^C_{}\right)^{\frac{1}{C_{}}}$$
(11)
where $`a_\mathrm{x}`$, $`a_\mathrm{y}`$, and $`a_\mathrm{z}`$ are the scale lengths and $`C_{}`$ and $`C_{}`$ are the face-on and the edge-on shape parameters. The radial dependencies of the bar density are given by various density profiles that we designate by labels:
$`\mathrm{G}:`$ $`\rho _\mathrm{G}=`$ $`\rho _0\times f_{\mathrm{max}}(R_{\mathrm{xy}})\times \mathrm{exp}(0.5R_\mathrm{s}^2)`$ (12)
$`\mathrm{G3}:`$ $`\rho _{\mathrm{G3}}=`$ $`\rho _0\times f_{\mathrm{max}}(R_{\mathrm{xy}})\times R_\mathrm{s}^{1.8}\mathrm{exp}(R_\mathrm{s}^3)`$ (13)
$`\mathrm{E}:`$ $`\rho _\mathrm{E}=`$ $`\rho _0\times f_{\mathrm{max}}(R_{\mathrm{xy}})\times \mathrm{exp}(R_\mathrm{s}^n)`$ (14)
$`\mathrm{E3}:`$ $`\rho _{\mathrm{E3}}=`$ $`\rho _0\times f_{\mathrm{max}}(R_{\mathrm{xy}})\times K_0(R_\mathrm{s})`$ (15)
$`\mathrm{S}:`$ $`\rho _\mathrm{S}=`$ $`\rho _0\times f_{\mathrm{max}}(R_{\mathrm{xy}})\times \mathrm{sech}^2(R_\mathrm{s})`$ (16)
$`\mathrm{P}:`$ $`\rho _\mathrm{P}=`$ $`\rho _0\times f_{\mathrm{max}}(R_{\mathrm{xy}})\times [1+(R_s/R_c)]^1`$ (17)
where
$`f_{\mathrm{max}}(R_{\mathrm{xy}})=\{\begin{array}{ccc}& 1.0\hfill & \text{for }R_{\mathrm{xy}}R_{\mathrm{max}}\hfill \\ & \mathrm{exp}\left(\frac{1}{2}\left(\frac{R_{\mathrm{xy}}R_{\mathrm{max}}}{a_{\mathrm{max}}}\right)^2\right)\hfill & \text{for }R_{\mathrm{xy}}>R_{\mathrm{max}}\hfill \end{array}`$ (20)
is a cutoff function and
$$R_{\mathrm{xy}}=\sqrt{x^2+y^2}$$
(21)
is the distance from the galactic centre in the xy-plane.
##### Model D0:
For the young disk population we use the model proposed by Robin et al. (robin03 (2003)) to describe the young (age $`<0.15`$ Gyr) stellar disk population of the Galaxy:
$$\rho (x,y,z)=\rho _0(\mathrm{exp}((a/R_0)^2)\mathrm{exp}((a/R_i)^2))$$
(22)
where
$$a^2=x^2+y^2+z^2/ϵ^2.$$
(23)
This model presents a truncated exponential disk profile with a fixed disk axis ratio of $`ϵ=0.014`$, a fixed disk scale radius of $`R_0=5`$ kpc, and a fixed inner disk truncation radius of $`R_i=3`$ kpc. The vertical exponential scale height of the disk is $`z_0=70`$ pc.
##### Model D1:
For the old disk population we use the model proposed by Robin et al. (robin03 (2003)) to describe the old (age $`710`$ Gyr) stellar disk population of the Galaxy:
$`\rho (x,y,z)`$ $`=`$ $`\rho _0(\mathrm{exp}((0.25+a^2/R_0^2)^{1/2})`$ (24)
$`\mathrm{exp}((0.25+a^2/R_i^2)^{1/2}))`$
This model presents a truncated exponential disk profile with a fixed disk axis ratio of $`ϵ=0.0791`$, a fixed disk scale radius of $`R_0=2.53`$ kpc, and a fixed inner disk truncation radius of $`R_i=1.32`$ kpc. The vertical exponential scale height of the disk is $`z_0=200`$ pc.
##### Model H:
To model the stellar halo we use the general model proposed by Robin et al. (robin03 (2003)):
$$\rho (R,z)=\rho _0(a/R_{\mathrm{}})^n$$
(25)
where $`a`$ is defined by Eq. 23 (the flatness of the model is set by the value of the axis ratio $`ϵ`$ in Eq. 23), and $`aa_c`$, avoiding a singularity at the galactic centre. $`n`$ determines the slope of the density profile.
##### Model S:
A set of galactocentric nested spherical shells of constant density. The radii of the shells and their number has been varied in order to maximise the MLR, while using the minimum number of shells required to satisfactorily describe the data.
###### Acknowledgements.
The SPI project has been completed under the responsibility and leadership of CNES. We are grateful to ASI, CEA, CNES, DLR, ESA, INTA, NASA and OSTC for support. |
warning/0506/quant-ph0506225.html | ar5iv | text | # Optimal Bell tests do not require maximally entangled states
## Abstract
Any Bell test consists of a sequence of measurements on a quantum state in space-like separated regions. Thus, a state is better than others for a Bell test when, for the optimal measurements and the same number of trials, the probability of existence of a local model for the observed outcomes is smaller. The maximization over states and measurements defines the optimal nonlocality proof. Numerical results show that the required optimal state does not have to be maximally entangled.
As first shown by Bell Bell in 1964, the correlations among the measurement outcomes of space-like separated parties on some quantum states cannot be reproduced by a local theory. This fact is often referred to as quantum nonlocality and has been recognized as the most intriguing quantum feature. The fundamental importance of the work by Bell was that it provided conditions for experimentally testing Quantum Mechanics (QM) versus the whole set of local models, the so-called Bell inequalities. The experimental demonstration exp , up to some loopholes, of a Bell inequality violation definitely closed the Einstein-Podolsky-Rosen program EPR for the existence of a local theory alternative to QM.
The interest on quantum correlations, or entanglement, has dramatically increased during the last two decades due to the emerging field of Quantum Information Science (QIS) book . It has been realized that quantum states provide new ways of information processing and communication without analog in Classical Information. The essential resource for most of these applications are entangled states. This has motivated a strong effort devoted to the characterization and quantification of the entanglement of quantum states. Although many questions remain unanswered, the problem is completely solved for the case of pure states in bipartite systems. For a state $`|\mathrm{\Psi }_A_B`$, its amount of entanglement is specified by the so-called entropy of entanglement BBPS , $`E(\mathrm{\Psi })=S(\rho _A)`$, where $`S`$ is the usual von Neumann entropy and $`\rho _A=\text{tr}_B(|\mathrm{\Psi }\mathrm{\Psi }|)`$. In particular, this means that the maximally entangled state in a bipartite system of dimension $`d\times d`$, reads
$$|\mathrm{\Phi }_d=\frac{1}{\sqrt{d}}\underset{i=1}{\overset{d}{}}|ii,$$
(1)
where $`\{i\}`$ define orthonormal bases in $`_A`$ and $`_B`$.
Apart from their importance for quantum information applications, entangled states provide the only known way of establishing nonlocal correlations among space-like separated parties. It is meant here by nonlocal those correlations that (i) cannot be explained by a local model but (ii) do not allow any faster-than-light communication, that is, they are consistent with the no-signaling condition. Indeed, it is a well-established result that a quantum pure state violates a Bell inequality if and only if it is entangled Gisin . However, it is also well known that there exist nonlocal correlations that are not achievable by measuring quantum states PR . In a similar line of thought, it has very often been assumed that $`|\mathrm{\Phi }_d`$ represents the most nonlocal quantum state too. However, no precise demonstration of this fact has ever been given and, indeed, it is one of the scopes of this work to raise some doubts about this statement.
In what follows, entangled states constitute a resource for constructing nonlocality proofs. The strength of a Bell experiment has to be computed by means of statistical tools: a Bell test is better than another when, for the same number of trials, the probability that a local model explains the observed outcomes is smaller. Recall that statistical fluctuations on finite samples allow a local theory to predict the possibility of data violating a Bell inequality. The goal is then to identify those states needed in the construction of optimal Bell tests. The importance of constructing optimal nonlocality proofs is two-fold. First, from an experimental point of view, they allow improving present Bell experiments, especially in terms of the needed resources. Second, Bell tests also represent an important tool for QIS useful . In particular, they are useful for testing the quantumness of devices. This is a hardly explored problem in QIS that, for instance, can be especially relevant in cryptographic applications MY : given some observed correlations among several parties, how can its quantum origin be certified? Could these correlations have alternatively been established by classical means, i.e. shared randomness? Bell inequalities provide an answer to the previous questions.
The scenario: We consider the standard scenario for any Bell test. Two space-like separated parties, called Alice and Bob, share copies of a pure quantum state $`|\mathrm{\Psi }_A_B`$, of dimension $`d\times d`$. They can choose among $`m`$ possible measurements, each of $`n`$ outcomes, this being denoted by $`m\times n`$. $`A_j^i`$ denotes the positive operator corresponding to the outcome $`j`$ of the measurement $`i`$ for Alice, so $`_jA_j^i=\text{1}\text{1}`$. Similarly, Bob’s measurement operators are denoted by $`B_j^i`$. The probability that Alice and Bob obtain the outcomes $`j_A`$ and $`j_B`$ after applying the measurement $`i_A`$ and $`i_B`$, where $`j_A,j_B=1,\mathrm{},n`$ and $`i_A,i_B=1,\mathrm{},m`$, on $`|\mathrm{\Psi }`$ is
$$p_Q(j_A,j_B|i_A,i_B)=\text{tr}\left(A_{j_A}^{i_A}B_{j_B}^{i_B}|\mathrm{\Psi }\mathrm{\Psi }|\right).$$
(2)
In a Bell experiment, a quantum state is prepared and sent to the parties who measure it. After $`N`$ repetitions of the experiment, the frequencies of the results define a $`(m^2n^2)`$-dimensional vector whose components tend to $`p_Q(j_A,j_B|i_A,i_B)`$ when $`N\mathrm{}`$. A vector of probabilities is achievable using QM when there exist a state $`|\mathrm{\Psi }`$ and measurements $`A_{j_A}^{i_A}`$ and $`B_{j_B}^{i_B}`$ satisfying (2).
On the other hand, in a local model any observed correlation between measurement results in space-like separated regions should come from initially shared random data, denoted by $`\lambda `$. QM is nonlocal because some of the vectors (2) do not allow a local description, i.e. they cannot be written as
$$p_L(j_A,j_B|i_A,i_B)=\underset{\lambda }{}p(\lambda )p(j_A|i_A,\lambda )p(j_B|i_B,\lambda ).$$
(3)
Therefore, shared quantum states can be used to establish nonlocal correlations.
The goal of any Bell experiment is to test the hypothesis $`𝒬`$, “the observed outcomes are governed by a quantum probability distribution (2)”, against the composite hypothesis $``$, “there exists a local model (3) reproducing the data” vDGG . The statistical tool that quantifies the average amount of support in favor of $`𝒬`$ against $``$ per trial when the data are generated by $`𝒬`$ is the so-called relative entropy or Kullback-Leibler (KL) Divergence CT , $`D`$. For two probability distributions, $`\stackrel{}{p}_1`$ and $`\stackrel{}{p}_2`$ associated to the event $`\{z\}`$, it reads
$$D(\stackrel{}{p}_1||\stackrel{}{p}_2)=\underset{z}{}p_1(z)\mathrm{log}(p_1(z)/p_2(z)).$$
(4)
We denote by $`\stackrel{}{q}=(q(1,1,1,1),\mathrm{},q(m,m,n,n))`$ the quantum probabilities for the measurement settings $`i_A`$ and $`i_B`$, and outcomes $`j_A`$ and $`j_B`$. Using (2), one has
$$q(i_A,i_B,j_A,j_B)=\text{tr}\left(A_{j_A}^{i_A}B_{j_B}^{i_B}|\mathrm{\Psi }\mathrm{\Psi }|\right)p_M(i_A,i_B),$$
(5)
where $`p_M(i_A,i_B)`$ characterizes the choice of measurements by Alice and Bob. Now, the support in favor of $`𝒬`$ against $``$ provided by these quantum data is vDGG
$$D(𝒬||)=\underset{\stackrel{}{p}}{\mathrm{min}}D(\stackrel{}{q}||\stackrel{}{p}),$$
(6)
where the minimization runs over all alternative local models, $`\stackrel{}{p}`$. The vector $`\stackrel{}{p}`$ is defined analogously to (5), replacing the quantum term $`p_Q`$ (2) by a local model $`p_L`$ (3). This quantity gives the statistical figure of merit to be maximized in any nonlocality test vDGG . It is worth mentioning here that the KL Divergence (i) is an asymptotic measure and (ii) appears as the measure of statistical support for the two most commonly used methods for hypothesis testing, frequentist and Bayesian (see vDGG for more details). Moreover, and despite not being symmetric, it can be seen as a measure of statistical distance between probability distributions.
It is often convenient to interpret a Bell test as a game between a quantum and a local player vDGG ; Peres . The quantum player has to design an experimental situation for which the local player is unable to provide a model. Thus, the quantum player looks for the experiment that gives him the victory with the minimal number of repetitions, i.e. his task consists of designing the Bell test maximizing (6). In order to do that, he can choose the state to be prepared, the measurements, and the probability governing the choice of measurements, $`p_M(i_A,i_B)`$. Notice that we do not impose $`p_M(i_A,i_B)`$ to be product. Indeed, one could think of a configuration where an external referee is sending the choice of measurements to the parties. On the other hand, the local player only assumes the existence of a local model. In particular, he is allowed to change his description according to the observed data.
Results: In what follows, the optimality of Bell tests is analyzed according to the KL Divergence. The optimization of (6) in full generality is a very hard problem. Here, we mainly consider the standard situation where Alice and Bob apply two projective measurements, i.e. $`m=2`$, $`n=d`$ and $`A_i^j`$ and $`B_i^j`$ are mutually orthogonal one-dimensional projectors. For a fixed number of measurements, it is possible to search numerically the state and measurements defining the optimal Bell test. In the qubit case, $`d=2`$, the best nonlocality proof is given by the maximally entangled state $`|\mathrm{\Phi }_2`$ (1) and the measurements maximizing the violation of the CHSH inequality CHSH , as expected. The KL Divergence turns out to be equal to 0.046 bits vDGG and the optimal choice of settings is completely random, $`p_M(i_A,i_B)=1/4`$. Actually, the optimal choice of settings turns out to be random for all the situations considered in this work.
Moving to higher dimension, the optimal measurements for the maximally entangled state of two three-dimensional systems are the ones maximizing its violation of the CGLMP inequality CGLMP . The statistical strength is of 0.058 bits, reflecting the fact that quantum nonlocality increases with the dimension KGZMZ . However, it is known that the largest violation of the CGLMP inequality is given by a nonmaximally entangled state ADGL
$$|\mathrm{\Psi }_3^{mv}=\gamma (|00+|11)+\sqrt{12\gamma ^2}|22,$$
(7)
where $`\gamma 0.617`$. The measurements maximizing its statistical strength are again those maximizing its Bell violation (which are the same as for $`|\mathrm{\Phi }_3`$) and give 0.072 bits, larger than the value obtained for the maximally entangled state. The maximization now over the space of measurements, choices of settings and states gives the same measurements as above but for a different state,
$$|\mathrm{\Psi }_3\delta (|00+|11)+\sqrt{12\delta ^2}|22,$$
(8)
where $`\delta 0.642`$. Therefore, the $`3\times 3`$ state producing the optimal nonlocality test with two projective measurements per site does not have maximal entanglement. Actually, this state has even less entanglement than $`|\mathrm{\Psi }_3^{mv}`$. All these results are summarized in Table I. It is worth mentioning here that the optimal measurements are the same for all three states. Similar results are obtained for $`d=4`$: (i) the optimal measurements are those maximizing the Bell violation for $`|\mathrm{\Phi }_4`$ CGLMP ; ADGL but (ii) the optimal state, $`|\mathrm{\Psi }_4`$, is not maximally entangled. The corresponding KL Divergence is of 0.098 bits. The problem in full generality becomes intractable for larger $`d`$, so the following simplifications are considered.
First, the $`2\times d`$ measurements are taken equal to those maximizing the Bell violation for $`|\mathrm{\Phi }_d`$: the parties apply a unitary operation with only nonzero terms in the diagonal, $`e^{i\varphi _a(j)}`$ for Alice and $`e^{i\phi _b(j)}`$ for Bob, with $`j=0,\mathrm{},d1`$ and $`a,b=1,2`$. These phases read CGLMP
$$\varphi _1(j)=0\varphi _2(j)=\frac{\pi }{d}j\phi _1(j)=\frac{\pi }{2d}j\phi _2(j)=\frac{\pi }{2d}j.$$
(9)
Then, Alice carries out a discrete Fourier transform, $`U_{FT}`$, and Bob applies $`U_{FT}^{}`$, and they measure in the computational basis. Thus, it is assumed in what follows that these measurements define the optimal $`2\times d`$ Bell test. This is known to be the case for qubits, and our numerical results indicate that this also happens for $`d=3,4`$.
Once the settings are fixed, the problem is cast in a formulation very similar to a standard Bell inequality. The goal is now to obtain the $`d\times d`$ state maximizing (6) for the given settings. Let $`\stackrel{}{q}^s`$ and $`\stackrel{}{p}^s`$ denote a pair forming a solution to this problem, i.e.
$$\underset{\stackrel{}{q}}{\mathrm{max}}\underset{\stackrel{}{p}}{\mathrm{min}}D(\stackrel{}{q}||\stackrel{}{p})=D(\stackrel{}{q}^s||\stackrel{}{p}^s)=D^s.$$
(10)
For small deviations from this solution one has $`D(\stackrel{}{q}^s+\delta \stackrel{}{q}||\stackrel{}{p}^s)D^s`$ and $`D(\stackrel{}{q}^s||\stackrel{}{p}^s+\delta \stackrel{}{p})D^s`$. Therefore, all vectors of quantum and local probabilities close to the previous solution satisfy
$`{\displaystyle \underset{i}{}}\mathrm{log}\left({\displaystyle \frac{q_i^s}{p_i^s}}\right)q_i`$ $``$ $`D^s`$ (11)
$`{\displaystyle \underset{i}{}}{\displaystyle \frac{q_i^s}{p_i^s}}p_i`$ $``$ $`1.`$ (12)
The values $`D^s`$ and 1 are found by substituting $`q_i=q_i^s`$ and $`p_i=p_i^s`$, respectively. Actually this has to be true for all $`\stackrel{}{q}`$ and $`\stackrel{}{p}`$. If this was not the case, using convexity arguments one could construct a point arbitrarily close to $`\stackrel{}{p}^s`$ or $`\stackrel{}{q}^s`$ violating these conditions. Indeed, assume there exists a vector of quantum probabilities $`\stackrel{}{q}^{}`$ violating (11). Then, $`(1ϵ)\stackrel{}{q}^s+ϵ\stackrel{}{q}^{}`$ would also violate the same condition for arbitrarily small $`ϵ`$. Note that the quantity on the left hand side of (11) can be seen as the mean value of a Bell operator, while (12) defines a Bell inequality. Then, $`\stackrel{}{q}^s`$ maximizes (11) over all $`\stackrel{}{q}`$, while $`\stackrel{}{p}^s`$ does it for (12).
After inspection, one can see that the Bell inequality (12) corresponding to the optimal solution for $`d=2,3,4`$ is of the CGLMP form, up to taking a linear combination with the normalization condition $`_ip_i=1`$. Actually, Eq. (12) can be rewritten in these three cases as
$$[A_1B_1]+[B_1A_2]+[A_2B_2]+[B_2A_11]d1,$$
(13)
where $`[X]`$ stands for $`X`$ modulo $`d`$ and $`X=p(X=1)+2p(X=2)+\mathrm{}+(d1)p(X=d1)`$. This inequality easily follows from the identity
$$[A_1B_1+B_1A_2+A_2B_2+B_2A_11]=d1,$$
(14)
and the fact that $`[X]+[Y][X+Y]`$. One can see that Eq. (13) represents an extremely compact way of writing all CGLMP inequalities for arbitrary dimension. Then, it is assumed that the inequality (12), derived from Eq. (10), has the CGLMP form (13), up to linear combination with $`_ip_i=1`$, also for $`d>4`$. Thus the $`q_i^s/p_i^s`$ terms in (11) and (12) are known functions of one parameter.
The problem has now been hugely simplified. Under the mentioned assumptions, the state for an optimal $`2\times d`$ Bell test is given by the eigenvector of largest eigenvalue of the Bell operator (11), where the measurements are fixed as before, and where the coefficients of the Bell operator are determined (up to one unknown parameter, over which we also optimize) by the CGLMP inequality. The associated eigenvalue gives the optimal KL Divergence. This computation can be done up to very large dimension, the results can be found in Fig. 1.
Discussion: Figure 1 shows several interesting features. First of all, one can see that for an optimal $`2\times d`$ Bell test, there is no need for systems of very large dimension. Actually, the simplest CHSH scenario for the singlet state already constitutes a reasonably good test for ruling out local models. However, beyond this simple case, none of the optimal Bell tests requires a maximally entangled state. In all the studied situations, the Schmidt basis for the optimal state was the computational one. Assuming this is always the case, we can compute the conjectured optimal state for large $`d`$, say $`d=1000`$, finding that $`E(\mathrm{\Psi }_d)\mathrm{ln}d0.69\mathrm{log}d`$ bits.
It also follows from Fig. 1 that two measurements per site may not be optimal for large $`d`$. For instance, when $`d=16`$ the conjectured optimal $`2\times 16`$ test is worse than the $`4\times 16`$ test consisting of two independent realizations of the optimal $`2\times 4`$ Bell test for two copies of $`|\mathrm{\Psi }_4`$. Indeed, it is always possible to interpret two independent realizations of this Bell test as a “new” $`2^2\times 4^2`$ Bell test for $`|\mathrm{\Psi }^2`$. Using that (i) the KL Divergence is additive, $`D(\stackrel{}{q}^2||\stackrel{}{p}^2)=2D(\stackrel{}{q}||\stackrel{}{p})`$, and (ii) the closest local model to two independent realizations of the same Bell test corresponds to two independent realizations of the best local model for the single-copy case, the KL Divergence for this test is twice the initial one.
A priori, one would have expected the maximally entangled state to be the optimal state for any Bell test. A thorough numerical search of Bell tests for the maximally entangled state $`|\mathrm{\Phi }_3`$ using more settings per site and general measurements has been performed. No improvement over the optimal $`2\times 3`$ case was obtained. Actually, it is remarkable that Bell tests with two projective measurements per site are so good for low dimensional systems. Therefore, all the previous numerical results show that beyond qubits and for the same amount of resources (system dimension and number and type of settings) the optimal state for a Bell test is not maximally entangled.
Conclusions: Non-local correlations constitute an information theoretic resource per se BLMPPR , that can be distributed by means of quantum states. It is known that there are nonlocal correlations that cannot be established by measuring quantum states PR . Moreover, the nonlocal correlations obtained from the maximally entangled state $`|\mathrm{\Phi }_3`$ seem to be less robust against noise than those from $`|\mathrm{\Psi }_3^{mv}`$ ADGL . Actually, the communication cost of simulating the nonlocal correlations for $`|\mathrm{\Psi }_3^{mv}`$ seems to be higher than for $`|\mathrm{\Phi }_3`$ Pironio . More recently, it has been shown that the so-called nonlocal machine PR ; BLMPPR is sufficient for the simulation of the correlations in a singlet state CGMP , but it fails for some nonmaximally entangled states of two qubits BGS . All these result suggest that, despite the fact that all pure entangled states contain nonlocal correlations Gisin , the relation between entanglement and nonlocality is subtler than firstly expected, since they may represent different information resources.
In this work, entangled states are analyzed as a tool for the construction of Bell tests. For all the studied scenarios and beyond the qubit case, the states needed for an optimal Bell test are not maximally entangled.
This work is supported by the ESF, an MCYT “Ramón y Cajal” grant, the Generalitat de Catalunya, the Swiss NCCR “Quantum Photonics” and OFES within the EU project RESQ (IST-2001-37559). |
warning/0506/hep-th0506021.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The gauge theory/string theory correspondence furnishes a powerful set of tools for understanding gauge theories at strong coupling by performing computations in a dual string theory at weak coupling. However, the correspondence is only well-understood in systems where the string background is highly symmetric and nearly flat, but we expect that the duals to many interesting gauge theories (such as large-$`N`$ QCD or SQCD) will not have these properties. It is therefore an interesting challenge to study less symmetric string backgrounds, and in particular to study backgrounds with reduced supersymmetry.
One interesting class of models arises from compactifications of string theory on noncompact Calabi-Yau manifolds with D3-branes at conical singularities, which generically give rise to $`𝒩=1`$ gauge theories with product gauge groups and bifundamental matter. These models are attractive for several reasons. They possess minimal supersymmetry and are therefore closer to realistic gauge theories than the well studied $`𝒩=4`$ case; also, they lead to conformal field theories where the quantum conformal invariance is not obvious by inspection of the field theory (but where the supergravity dual makes conformal invariance manifest.) Perhaps the most striking feature of these theories is that one can break conformal invariance in a controlled way by adding fluxes through cycles of the Calabi-Yau geometry, which induce RG flow and confinement at low energies.
However, one missing element of these models is fundamental matter. Aside from being experimentally important, fundamentals give rise to many interesting things such as the phase structure of super-Yang-Mills theory in the infrared. In confining theories the fundamentals of course do not appear as asymptotic states but are instead confined in mesons and baryons.
In this note we study the mesonic fluctuations of a particular set of mesons in the conifold theory of Klebanov and Witten . This theory is interesting for its relative simplicity and also because its non-conformal version flows to a theory very similar to $`𝒩=1`$ pure glue theory in the infrared. Moreover all metrics for the corresponding supergravity solutions are known, allowing explicit computations. The mesons which we study arise as fluctuations on D7-branes which are embedded in the string background. The fundamental fields come from strings connecting the stack of D3-branes to the D7-branes. In the usual decoupling limit, the 3-7 strings and 3-3 strings, which describe the gauge theory, have a dual description in terms of the closed strings and 7-7 strings. The closed strings are the usual glueballs of the strongly coupled field theory while the open 7-7 strings are naturally identified with the mesons.
We will compute the spectrum of operator dimensions, which, as we will see, can be done exactly for a large portion of the states, and we will study the effect of giving masses to the quarks (which requires numerical work).
The paper is organized as follows. In section 2 we review the geometry of the conifold. In section 3 we discuss adding flavor to the Klebanov-Witten field theory by the addition of probe D7-branes. In section 4 we compute the spectrum for scalar mesons. In the case of massive quarks, we compute the mass spectrum numerically, but in the massless case (corresponding to the UV limit of the gauge theory) we obtain the spectrum analytically. In section 5 we discuss our results.
## 2 Review of the Conifold
In this section we briefly review the geometry of the conifold in order to fix notation. Useful references are .
The conifold is a non-compact Calabi-Yau 3-fold, defined by the equation
$`z_1z_2z_3z_4=0`$ (1)
in $`𝐂^4`$. Because Eqn.(1) is invariant under an overall real rescaling of the coordinates, this space is a cone, whose base is the Einstein space $`T^{1,1}`$ . The metric on the conifold may be cast in the form
$`ds_6^2=dr^2+r^2ds_{T^{1,1}}^2,`$ (2)
where
$$ds_{T^{1,1}}^2=\frac{1}{9}\left(d\psi +\underset{i=1}{\overset{2}{}}\mathrm{cos}\theta _id\varphi _i\right)^2+\frac{1}{6}\underset{i=1}{\overset{2}{}}\left(d\theta _i^2+\mathrm{sin}^2\theta _id\varphi _i^2\right)$$
(3)
is the metric on $`T^{1,1}`$. Here $`\psi `$ is an angular coordinate which ranges from $`0`$ to $`4\pi `$, while $`(\theta _1,\varphi _1)`$ and $`(\theta _2,\varphi _2)`$ parametrize two $`𝐒^2`$s in the standard way. This form of the metric shows that $`T^{1,1}`$ is a $`U(1)`$ bundle over $`𝐒^2\times 𝐒^2`$.
These angular coordinates are related to the $`z_i`$ variables by
$`z_1`$ $`=`$ $`r^{3/2}e^{i/2(\psi \varphi _1\varphi _2)}\mathrm{sin}{\displaystyle \frac{\theta _1}{2}}\mathrm{sin}{\displaystyle \frac{\theta _2}{2}},`$
$`z_2`$ $`=`$ $`r^{3/2}e^{i/2(\psi +\varphi _1+\varphi _2)}\mathrm{cos}{\displaystyle \frac{\theta _1}{2}}\mathrm{cos}{\displaystyle \frac{\theta _2}{2}},`$ (4)
$`z_3`$ $`=`$ $`r^{3/2}e^{i/2(\psi +\varphi _1\varphi _2)}\mathrm{cos}{\displaystyle \frac{\theta _1}{2}}\mathrm{sin}{\displaystyle \frac{\theta _2}{2}},`$
$`z_4`$ $`=`$ $`r^{3/2}e^{i/2(\psi \varphi _1+\varphi _2)}\mathrm{sin}{\displaystyle \frac{\theta _1}{2}}\mathrm{cos}{\displaystyle \frac{\theta _2}{2}}.`$
It is also sometimes helpful to consider a set of “homogeneous” coordinates $`A_a,B_b`$ where $`a,b=1,2`$, in terms of which the $`z_i`$ are
$`z_1`$ $`=`$ $`A_1B_1,z_2=A_2B_2,`$ (5)
$`z_3`$ $`=`$ $`A_1B_2,z_4=A_2B_1.`$ (6)
With this parameterization the $`z_i`$ obviously solve the defining equation of the conifold.
We may also parameterize the conifold in terms of an alternative set of complex variables $`w_i`$, given by
$`\begin{array}{cc}z_1=w_1+iw_2,\hfill & z_2=w_1iw_2,\hfill \\ z_3=w_3+iw_4,\hfill & z_4=w_3+iw_4.\hfill \end{array}`$ (9)
The conifold equation may now be written as
$`{\displaystyle w_i^2}=0`$ (10)
and we identify the $`T^{1,1}`$ base of the cone as the intersection of the conifold with the surface
$`{\displaystyle |w_i|^2}=r^3.`$ (11)
$`T^{1,1}`$ described in this way is explicitly invariant under $`SO(4)SU(2)\times SU(2)`$ rotations of the $`w_i`$ coordinates and under an overall phase rotation. Thus the symmetry group of $`T^{1,1}`$ is $`SU(2)\times SU(2)\times U(1)`$.
An important fact about $`T^{1,1}`$ is that it has Betti numbers $`b_2,b_3=1`$. The corresponding two-cycle and three-cycle may be expressed in terms of harmonic differential forms:
$`\omega _2`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\mathrm{\Omega }_{11}\mathrm{\Omega }_{22}\right),`$ (12)
$`\omega _3`$ $`=`$ $`\zeta \omega _2.`$ (13)
In this paper we will consider D7-branes in the model of Klebanov and Witten . This model is a particularly simple $`𝒩=1`$ gauge/gravity dual, obtained by placing a stack of $`N`$ D3-branes near a conifold singularity. The branes source the RR 5-form flux and warp the geometry:
$`ds_{10}^2`$ $`=`$ $`h(r)^{1/2}dx_\mu dx^\mu +h(r)^{1/2}(dr^2+r^2ds_{T^{1,1}}^2)`$ (14)
$`h(r)`$ $`=`$ $`1+{\displaystyle \frac{L^4}{r^4}}`$ (15)
$`g_sF_5`$ $`=`$ $`d^4xdh^1+(d^4xdh^1)`$ (16)
$`L^4`$ $`=`$ $`{\displaystyle \frac{27}{4}}\pi g_sN\alpha ^2.`$ (17)
Hereafter, we specialize to the near-horizon limit $`r/L1`$, and set $`L=1`$ for convenience. It may be easily restored by dimensional analysis at any point.
The dual field theory has gauge group $`SU(N)\times SU(N)`$ and matter fields $`A_{1,2},B_{1,2}`$ which transform in the bifundamental color representations $`(𝐍,\overline{𝐍})_c`$ and $`(\overline{𝐍},𝐍)_c`$. The theory also has a superpotential
$`W=\lambda Tr(A_iB_jA_kB_l)ϵ^{ik}ϵ^{jl}.`$ (18)
By solving the F-term equations for this superpotential, we obtain supersymmetric vacua for arbitrary diagonal $`A_{1,2}`$ and $`B_{1,2}`$, so that the moduli space of the field theory is precisely that of $`N`$ D3-branes placed on the conifold.
## 3 Adding flavor
In this section we review the procedure of adding flavor branes to AdS/CFT in general and make several useful comments on adding flavor to the Klebanov-Witten field theory both in terms of the bulk geometry and the dual field theory. This general procedure was first pointed out in and was exploited in the $`AdS_5\times S^5`$ case in . Some other examples of flavored theories with probe branes have been studied in .
One way to add flavor to AdS/CFT is to take a system of D3-branes and then to add D7-branes which fill the four $`x^\mu `$ directions and four of the six transverse dimensions . In flat space such a configuration of branes is clearly supersymmetric. As usual there is an $`𝒩=4`$ $`SU(N)`$ SYM theory living on the D3-branes. Strings with one end on a D3-brane and one end on a D7-brane couple to the fields of the D3-brane gauge theory as quarks.
For AdS/CFT purposes we can now take the supergravity approximation in which D3-branes are replaced by an $`AdS_5\times S^5`$ geometry with Ramond-Ramond flux, while we retain the D7-branes as probes which fill the five AdS directions and which wrap a topologically trivial 3-cycle of the internal 5-manifold (for example an $`S^3`$ submanifold of the $`S^5`$ of $`AdS_5\times S^5`$). The triviality of the 3-cycle guarantees that the brane carries no net charge and will not introduce any tadpoles. On the other hand, topological triviality also suggests that one might be able to shrink the $`S^3`$ and slip it off of the $`S^5`$, naively in contradiction with the flat space picture of D3 and D7-branes. It turns out that subtleties of the $`AdS`$ geometry play a key role in ensuring stability. The mass eigenvalues of modes controlling the D-brane slipping off the 3-cycle are negative, but are above the Breitenlohner-Freedman bound , so that the 7-brane embedding is stable.
In the flat space picture, if the D3-branes and D7-branes intersect then the quarks are massless, and if the D3-branes and D7-branes are separated then the quarks are massive. This translates nicely into the AdS picture in the following way. A D7-brane which intersects the D3-branes in flat space gets mapped to a D7-brane which fills the whole AdS space and wraps a three-sphere of constant size in the $`S^5`$. On the other hand, a D7-brane separated from the stack of D3-branes maps to a D7-brane which wraps an $`S^3`$ with some asymptotic size at large AdS radius, but this $`S^3`$ shrinks to zero size at some finite radius (which is possible because of the topological triviality). In the 5-dimensional AdS space the D7-brane appears to fill out the radial direction up to some minimal radius where it “ends.”
It is interesting of course to consider theories with branes in spaces which are not flat. The basic picture of D3 and D7 branes contributing gauge fields and quarks will not change, but many details are different. For simplicity throughout this paper we specialize to the case of a single D7-brane. If the number of D3-branes is large then the D7 backreaction can be systematically neglected and it is appropriate to treat the D7-brane as a probe, which we do throughout this paper. Inclusion of backreaction effects in other geometries has been explored in .
Let us consider D7-branes embedded in the geometry of the conifold by the equation $`z_1=\mu `$. In terms of the standard coordinate system,
$`z_1=r^{3/2}e^{i/2(\psi \varphi _1\varphi _2)}\mathrm{sin}{\displaystyle \frac{\theta _1}{2}}\mathrm{sin}{\displaystyle \frac{\theta _2}{2}}`$
so the embedding equation gives two conditions, one on the magnitude of $`z_1`$ and one on the phase:
$`r_0`$ $`=`$ $`\left({\displaystyle \frac{|\mu |}{\mathrm{sin}\frac{\theta _1}{2}\mathrm{sin}\frac{\theta _2}{2}}}\right)^{2/3},`$ (19)
$`\psi _0`$ $`=`$ $`\varphi _1+\varphi _2+\mathrm{const}.`$ (20)
This embedding can be explicitly shown to be supersymmetric by considering the $`\kappa `$-symmetry on the worldvolume of the brane . A slightly different embedding equation was studied in the warped deformed conifold by .
It was proposed in that the embedding $`z_1=\mu `$ leads to fields, summarized in Table 1 and a superpotential of the form
$`W`$ $`=`$ $`W_{flavors}+W_{masses},`$ (21)
$`W_{flavors}`$ $`=`$ $`h\stackrel{~}{q}A_1Q+gqB_1\stackrel{~}{Q},W_{masses}=\mu _1q\stackrel{~}{q}+\mu _2Q\stackrel{~}{Q}.`$ (22)
To relate this superpotential to the D7-brane geometry, let us probe the space with a single D3-brane, which corresponds to giving some expectation values to $`A_1`$ and $`B_1`$. One then finds that the theory on this probe has a massless flavor when $`A_1B_1=\mu _1\mu _2/(gh)`$, which is exactly of the form of the embedding equation $`z_1=\mu `$. Part of the motivation for this superpotential was a comparison with a type IIA brane construction where a D6-brane splits on an NS5-brane, contributing two flavor branes and correspondingly two sets of flavors. For the type IIB picture, in the massless limit of the field theory, this corresponds nicely to the presence of two solution branches of $`z_1=0`$, namely $`\theta _1=0`$ and $`\theta _2=0`$. If the quarks are massless there is an $`SU(K)\times SU(K)`$ flavor symmetry, where $`K`$ is the number of probe D7-branes. If the quarks are massive then the two branches of the D7-branes connect and the flavor symmetry is broken down to the diagonal $`SU(K)`$.
An alternative perspective is to suppose that one of the masses $`\mu _i`$ is larger than the other and then integrate out the associated flavors. Then one obtains a quartic superpotential of the form
$`W=q(A_1B_1\mu )\stackrel{~}{q}`$ (23)
which again produces the appropriate massless locus for a D3-brane probe. Our probe calculations will show that this quartic superpotential is consistent with adding D7-branes with massive flavors (and then with the limit where we take the masses to zero.) Of course, because we believe the quarks can be massive the consistency was virtually guaranteed. However, taking the limit $`\mu 0`$ and setting $`\mu =0`$ are different things, and it is unclear whether the theory corresponding to the cubic superpotential can be realized or not.
## 4 Scalar mesons
In this section, we compute the dimension and mass spectra of the scalar mesons. As discussed in the introduction, in the probe and decoupling limits the 7-7 strings are identified with the mesons in the dual field theory. We will thus be able to extract the mass spectrum of the spin=0 mesons and their conformal dimension in the UV limit by studying the 7-7 strings.
The semiclassical dynamics of this D7-brane are captured by the Dirac-Born-Infeld action
$`S_{DBI}=\tau _7{\displaystyle d^8\xi \sqrt{\underset{ij}{det}\left(g_{MN}+F_{MN}\right)\frac{y^M}{\xi ^i}\frac{y^N}{\xi ^j}}}+{\displaystyle \frac{g_s\tau _8}{2}}{\displaystyle C_4}F_2F_2.`$ (24)
where $`\xi ^i`$ are coordinates on the D7-brane. We will compute the spectrum of fluctuations for the D7-branes using this action.
Let us consider the fluctuations of scalar modes alone, with all D7-brane gauge fields turned off. Then the DBI action is simply the worldvolume of the 7-brane. Let us choose as coordinates on the brane eight of the spacetime coordinates: $`(x^\mu ,\theta _1,\theta _2,\varphi _1,\varphi _2).`$ The fluctuations can be described by setting
$`r`$ $`=`$ $`r_0(\theta _i)(1+\chi (x^\mu ,\theta _i,\varphi _j)),`$ (25)
$`\psi `$ $`=`$ $`\psi _0(\varphi _i)+3\eta (x^\mu ,\theta _i,\varphi _j).`$ (26)
The unperturbed induced metric takes the form
$`g_{MN}`$ $`=`$ $`\left(\begin{array}{ccc}r_0^2\eta _{\mu \nu }& 0& 0\\ 0& g_{\theta _i\theta _j}& 0\\ 0& 0& g_{\varphi _i\varphi _j}\end{array}\right)`$ (30)
$`g_{\theta _i\theta _j}`$ $`=`$ $`\left(\begin{array}{cc}\frac{1}{6}+\frac{1}{9}\mathrm{cot}^2\frac{\theta _1}{2}& \frac{1}{9}\mathrm{cot}\frac{\theta _1}{2}\mathrm{cot}\frac{\theta _2}{2}\\ \frac{1}{9}\mathrm{cot}\frac{\theta _1}{2}\mathrm{cot}\frac{\theta _2}{2}& \frac{1}{6}+\frac{1}{9}\mathrm{cot}^2\frac{\theta _2}{2}\end{array}\right)`$ (33)
$`g_{\varphi _i\varphi _j}`$ $`=`$ $`\left(\begin{array}{cc}\frac{1}{6}\mathrm{sin}^2\theta _1+\frac{1}{9}(1+\mathrm{cos}\theta _1)^2& \frac{1}{9}(1+\mathrm{cos}\theta _1)(1+\mathrm{cos}\theta _2)\\ \frac{1}{9}(1+\mathrm{cos}\theta _1)(1+\mathrm{cos}\theta _2)& \frac{1}{6}\mathrm{sin}^2\theta _2+\frac{1}{9}(1+\mathrm{cos}\theta _2)^2\end{array}\right)`$ (36)
One expands about this metric via the matrix identity
$`\sqrt{detA+\delta A}=\sqrt{detA}(1+{\displaystyle \frac{1}{2}}\mathrm{Tr}A^1\delta A+{\displaystyle \frac{1}{8}}(\mathrm{Tr}A^1\delta A)^2{\displaystyle \frac{1}{4}}\mathrm{Tr}A^1\delta AA^1\delta A+\mathrm{}`$ (37)
The terms first order in the fluctuations $`\chi `$ and $`\eta `$ turn out to be total derivatives, as is necessary for our embedding to be a solution of the equations of motion. The quadratic order fluctuations lead to an action of the form
$`S`$ $`=`$ $`\tau _7{\displaystyle }d^4xd\theta _1d\theta _2d\varphi _1d\varphi _2\sqrt{detg_0}{\displaystyle \frac{1}{C}}[{\displaystyle \frac{1}{2}}g_0^{ab}_a\chi _b\chi +{\displaystyle \frac{1}{2}}g_0^{ab}_a\eta _b\eta `$
$`+`$ $`{\displaystyle \frac{4}{\mathrm{sin}^2\frac{\theta _2}{2}}}\chi _{\varphi _2}\eta {\displaystyle \frac{2}{C\mathrm{sin}^2\frac{\theta _2}{2}}}\left(\mathrm{cot}{\displaystyle \frac{\theta _1}{2}}_{\theta _1}+\mathrm{cot}{\displaystyle \frac{\theta _2}{2}}_{\theta _2}\right)\chi _{\varphi _2}\eta `$
$`+`$ $`{\displaystyle \frac{4}{\mathrm{sin}^2\frac{\theta _1}{2}}}\chi _{\varphi _1}\eta {\displaystyle \frac{2}{C\mathrm{sin}^2\frac{\theta _1}{2}}}(\mathrm{cot}{\displaystyle \frac{\theta _1}{2}}_{\theta _1}+\mathrm{cot}{\displaystyle \frac{\theta _2}{2}}_{\theta _2})\chi _{\varphi _1}\eta ]`$
with
$`C=1+{\displaystyle \frac{2}{3}}\mathrm{cot}^2{\displaystyle \frac{\theta _1}{2}}+{\displaystyle \frac{2}{3}}\mathrm{cot}^2{\displaystyle \frac{\theta _2}{2}}.`$ (39)
### 4.1 The UV/massless limit
Even though the quarks have mass, if we flow to the UV, they effectively become massless and conformal symmetry is restored, at least at the classical level. It is interesting to inquire what the dimensions of the operators are in the UV field theory. In the dual, this corresponds to computing near the boundary of $`AdS`$.
Examining (19) we see that there are two ways to go near the boundary: $`\theta _10`$ or $`\theta _20`$. Which one we choose will determine which side of the conifold we are on near the boundary. In the current setup, the physics is symmetric between exchange of $`\theta _1`$ and $`\theta _2`$ so we will simply choose the $`\theta _10`$ limit. We can compute the dimensions of the operators in the field theory by examining the scaling of the 7-7 strings near the boundary.
We define $`r_0=\mu ^{2/3}e^{\beta /3}`$ (now a good coordinate because of the conformal invariance), and $`\mathrm{cos}^2\frac{\theta _2}{2}=x`$. Defining the linear combination of fields
$`\mathrm{\Phi }^\pm =\chi \pm i\eta `$ (40)
we find that the equations for $`\mathrm{\Phi }^\pm `$ are two fully decoupled partial differential equations. This equation is solved by a separation of variables ansatz
$`\mathrm{\Phi }^\pm =\rho ^\pm (x)e^{k\beta /3}e^{im_1\varphi _1+im_2\varphi _2}.`$ (41)
The equations of motion for the scalar reduces to ordinary differential equations for the functions $`\rho (x)^\pm `$ which take the form
$`((1x){\displaystyle \frac{}{x}}x{\displaystyle \frac{}{x}}+{\displaystyle \frac{1}{6}}k(k1){\displaystyle \frac{m_1^2}{8}}{\displaystyle \frac{3x}{1x}}+{\displaystyle \frac{m_1m_2}{2(1x)}}`$
$`{\displaystyle \frac{1}{4}}{\displaystyle \frac{m_2^2}{x(1x)}}{\displaystyle \frac{m_1}{4}}{\displaystyle \frac{1+x}{1x}}\pm {\displaystyle \frac{m_2}{2(1x)}})\rho ^\pm (x)`$ $`=`$ $`0.`$ (42)
This equation has singularities only at $`x=0,1,\mathrm{}`$ and is therefore of hypergeometric type. To see this explicitly we define new functions by rescaling the $`\rho ^\pm `$ by factors of $`x`$ and $`(1x)`$, which allow us to remove the terms in (42) proportional to $`\frac{1}{x}`$ and $`\frac{1}{1x}`$. Explicitly, we write
$`\rho ^\pm (x)`$ $`=`$ $`x^p(1x)^qf^\pm (x)`$ (43)
for which the equation of motion becomes
$`[x(1x){\displaystyle \frac{^2}{x^2}}+(1+2px(1+2p+2q)){\displaystyle \frac{}{x}}+{\displaystyle \frac{(q^2q+\frac{1}{4}\frac{1}{4}(m_1m_2\pm 1)^2)}{1x}}`$
$`+{\displaystyle \frac{(p^2\frac{m_2^2}{4})}{x}}+{\displaystyle \frac{1}{6}}k(k1){\displaystyle \frac{1}{4}}({\displaystyle \frac{m_1^2}{2}}m_1)(p+q)^2]f^\pm (x)=0.`$ (44)
Appropriate choice of the parameters $`p`$ and $`q`$ eliminates the terms proportional to $`1/x`$ and $`1/(1x)`$. The wavefunctions are given by
$`f^\pm (x)=_2F_1(\alpha ,\alpha +2(p+q),1+2p;x)`$ (45)
which are regular when $`\alpha `$ is a non-negative integer; it turns out that the original $`\rho `$ are also regular with this condition over the range $`0x1`$ ($`0\theta _2\pi `$) which encompasses our domain. We also find that there are two possible values of $`k`$:
$`k_1`$ $`=`$ $`{\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{2}}\sqrt{1+3m_1^26m_1+24(\alpha +p+q)^2}`$ (46)
$`k_2`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{2}}\sqrt{1+3m_1^26m_1+24(\alpha +p+q)^2}.`$ (47)
To be painfully explicit, we exhibit the solutions for $`f^+`$ (the $`f^{}`$ are straightforwardly related.) It is clear that there are always two choices of $`p`$ and $`q`$ which do the trick; to make regularity transparent we will always choose $`p`$ and $`q`$ to be positive. We then have four cases:
* $`m_20,m_1m_2`$: We choose $`p=m_2/2`$ and $`q=1+\frac{1}{2}(m_1m_2)`$, so that $`f^+=_2F_1(\alpha ,2+m_1+\alpha ,1+m_2;x)`$. The wavefunction is regular if $`\alpha `$ is a non-negative integer (negative $`\alpha `$ gives irregular or redundant solutions) and so the two values of $`k`$ are quantized to be
$`k={\displaystyle \frac{1}{2}}\pm {\displaystyle \frac{1}{2}}\sqrt{1+3m_1^26m_1+24(\alpha +{\displaystyle \frac{m_1}{2}}+1)^2}.`$ (48)
* $`m_20,m_1<m_2`$: Again we choose $`p=m_2/2`$, but now to make regularity obvious we take $`q=(m_2m_1)/2`$, such that $`q>0`$. Now $`f^+=_2F_1(\alpha ,2m_2m_1+\alpha ,1+m_2;x)`$, again with $`\alpha `$ a non-negative integer. The quantized values of $`k`$ are
$`k={\displaystyle \frac{1}{2}}\pm {\displaystyle \frac{1}{2}}\sqrt{1+3m_1^26m_1+24(\alpha +m_2{\displaystyle \frac{m_1}{2}})^2}.`$ (49)
* $`m_2<0,m_1m_2`$: Now we choose $`p=m_2/2`$ and $`q=1+\frac{1}{2}(m_1m_2)`$, finding that $`f^+=_2F_1(\alpha ,2+m_12m_2+\alpha ,1m_2;x)`$, with $`\alpha `$ a non-negative integer. The allowed values of $`k`$ are
$`k={\displaystyle \frac{1}{2}}\pm {\displaystyle \frac{1}{2}}\sqrt{1+3m_1^26m_1+24(\alpha +{\displaystyle \frac{m_1}{2}}m_2+1)^2}.`$ (50)
* $`m_2<0,m_1<m_2`$: Now we choose $`p=m_2/2`$ and $`q=(m_2m_1)/2`$, finding that $`f^+=_2F_1(\alpha ,\alpha m_1,1m_2;x)`$, with $`\alpha `$ a non-negative integer. The allowed values of $`k`$ are
$`k={\displaystyle \frac{1}{2}}\pm {\displaystyle \frac{1}{2}}\sqrt{1+3m_1^26m_1+24(\alpha {\displaystyle \frac{m_1}{2}})^2}.`$ (51)
To find the dimensions of the operators, we recall that $`e^{k\beta /3}r^k`$. In the AdS/CFT correspondence a minimal massless scalar field dual to an operator of dimension $`\mathrm{\Delta }`$ scales as $`r^\mathrm{\Delta }`$ for its normalizable part and $`r^{\mathrm{\Delta }4}`$ for its non-normalizable part. However, by examining (4) we see that the kinetic terms for these scalars are not canonically normalized, which means that the possible scalings at infinity are modified to $`r^{\mathrm{\Delta }+p}`$ and $`r^{\mathrm{\Delta }4+p}`$ for some $`p`$. Using the values for $`k_1`$ and $`k_2`$ we have $`2\mathrm{\Delta }4=k_1k_2`$, which one can compute straightforwardly.
The dimensions are mostly complicated irrational numbers (reminiscent of the closed string spectrum on $`T^{1,1}`$ ) but a few features of the spectrum stand out. The lowest mode has $`m_1=m_2=\alpha =0`$, and is simply a constant; it can be assigned dimension 5/2 or 3/2. From the earlier discussion of the massive field theory, it is natural to choose dimension 3/2 and associate this mode with the operator $`q\stackrel{~}{q}`$. Note also that the mode with $`m_1=m_2=1`$ and $`\alpha =0`$ has dimension 3, appropriate for a superpotential term – we identify this mode with the operator $`qA_2B_2\stackrel{~}{q}`$. If added to the superpotential, this operator would change the D7-brane embedding from $`z_1=\mu `$ to $`z_1+ϵz_2=\mu `$.
For large $`m_1`$, all the dimensions scale as $`\mathrm{\Delta }3m_1/2`$. This is consistent with identification of the corresponding gauge theory operators as
$`q(AB)(AB)\mathrm{}(AB)\stackrel{~}{q}`$ (52)
where, ignoring the $`q`$ fields, each insertion of $`(AB)`$ should increase the dimension by 3/2 and the relevant $`SU(2)`$ charge (associated with $`m_1`$) by one unit. Unlike the case of baryonic operators on the conifold, where one finds an exact scaling $`\mathrm{\Delta }3N/4`$ , we see that the mesons only exhibit a simple scaling with the charge in a large-charge limit. For small charges there are boundary effects due to the quarks which, at least in the large-$`N`$ limit, are completely calculable here. This behavior should also be contrasted with the case of flavors added to $`𝒩=4`$ super-YM theory, where the meson dimensions were pure integers.
If we take instead the limit of large $`\alpha `$, we see that the dimensions scale as $`\mathrm{\Delta }\sqrt{6}\alpha `$. It would be interesting to find an explanation for this curious scaling in the field theory.
### 4.2 The mass spectra
Having obtained the conformal dimensions for the mesons in the conformal limit, we would like to compute their spectrum by solving the full differential equation without taking any simplifying limits. Unfortunately, we will find that the equation is not amenable to analytic solution and so we will have to appeal to numerical methods. We will display selected results from several cases that are illustrative of the general behavior.
We will again find that the linear combination of fields (40) decouples the equations of motion. Rewriting the action (4) in terms of $`\mathrm{\Phi }^\pm `$ and varying gives
$`{\displaystyle \frac{1}{\sqrt{g_0}}}_a\left({\displaystyle \frac{1}{C}}\sqrt{g_0}g_0^{ab}_b\mathrm{\Phi }^\pm \right)\pm 3i\left[{\displaystyle \frac{1}{\sqrt{g_0}}}_{\theta _i}(\sqrt{g_0}\gamma ^{\theta _i\varphi _j})\gamma ^{\varphi _j}\right]_{\varphi _j}\mathrm{\Phi }^\pm =0,`$ (53)
where the $`a,b`$ indices run over the $`x^\mu `$ and the $`\theta _{1,2},\varphi _{1,2}`$, and
$`\gamma ^{\theta _i\varphi _j}`$ $`=`$ $`{\displaystyle \frac{4(1+\mathrm{cos}\theta _j)}{C^2\mathrm{sin}^2\theta _j}}_{\theta _i}\mathrm{ln}r_0(\theta _i,\theta _j),`$ (54)
$`\gamma ^{\varphi _j}`$ $`=`$ $`{\displaystyle \frac{8(1+\mathrm{cos}\theta _j)}{3C\mathrm{sin}^2\theta _j}}.`$ (55)
The inverse components of the metric are straightforward to find from the form given in (30-36). Since $`_{x^\mu }`$ and $`_{\varphi _i}`$ are Killing vectors we can write
$`\mathrm{\Phi }^\pm =\psi ^\pm (\theta _1,\theta _2)e^{ikx}e^{im_1\varphi _1+im_2\varphi _2}.`$ (56)
We find that (53) becomes
$``$ $`_{\theta _i}\left({\displaystyle \frac{1}{C}}\sqrt{g_0}g_0^{\theta _i\theta _j}_{\theta _j}\psi ^\pm \right)+{\displaystyle \frac{1}{C}}\sqrt{g_0}g_0^{\varphi _i\varphi _j}m_im_j\psi ^\pm \pm 3\left[_{\theta _i}(\sqrt{g_0}\gamma ^{\theta _i\varphi _j})\sqrt{g_0}\gamma ^{\varphi _j}\right]m_j\psi ^\pm `$ (57)
$`=`$ $`{\displaystyle \frac{\sqrt{g_0}}{Cr_0^2(\theta _1,\theta _2)}}k^2\psi ^\pm .`$
Note that for massive modes $`k^2=k_\mu k^\mu <0`$ since it must be timelike. Using simple Kaluza-Klein arguments we see that the mass of the mesons in the dual field theory is $`M^2=k^2`$. It is evident from the form of this equation that the only difference between the equation for $`\psi ^+`$ and for $`\psi ^{}`$ is in a term proportional to $`m_j`$. Thus, we can choose to solve for $`\psi ^+`$ without loss of generality. This equation cannot be solved analytically. In addition, we have found no simple way to separate the equation in the $`\theta _1`$,$`\theta _2`$ directions and so we must use a numerical approach to solving the partial differential equation.
#### 4.2.1 The numerical approach
Because we are unable to separate the partial differential equations into ordinary differential equations, we must use a technique slightly more involved than the regular finite difference scheme and shooting technique. We will make use of the finite element method and the Arnoldi algorithm via Matlab to solve for the mass eigenvalues, $`k^2`$. We will use a mesh with 2779 nodes and 5392 triangles for all problems involved. The $`\theta _i`$s are part of two different $`S^2`$s and thus range from $`0`$ to $`\pi `$. Since we already know $`\theta _i0`$ corresponds to going near the boundary, and we want normalizable modes, we place Dirichlet boundary conditions at $`\theta _{1,2}=0`$. We must also demand regularity at $`\theta _{1,2}=\pi `$, which corresponds to placing Neumann boundary conditions at $`\theta _{1,2}=\pi `$.
We will first examine the simplest case, when $`m_1=m_2=0`$. Setting $`\mu ^{4/3}=.02`$ we solve (57) for the first 50 eigenvalues. The eigenvalues break up into different series corresponding to the number of nodes in the $`(\theta _1,\theta _2)`$ plane. In figures 1 and 2 we display the first two such series. Higher series have similar behavior. The $`+`$ signs denote actual mass eigenvalues, while the solid lines are best fit lines.
We will also find similar behavior for modes with $`m_{1,2}0`$. In figure 3 we display the zero node modes for the case $`m_1=1`$, $`m_2=2`$ with $`\mu ^{4/3}=2`$ (note: changing $`\mu `$ just changes the eigenvalues by an overall scaling, as expected since it merely scales the mass gap for the quarks). We find similar behavior for other values of $`m_{1,2}`$ and different number of nodes.
For all the cases we see that in the large $`n`$ limit we have $`Mn\mu ^{2/3}`$ as we would expect. Restoring $`L`$ by dimensional analysis, and using $`m_q\frac{\mu ^{2/3}}{2\pi \alpha ^{}}`$ we find the the mass gap for the lightest meson is
$`M_{gap}{\displaystyle \frac{m_q}{\sqrt{g_sN}}}`$ (58)
Therefore, in the supergravity regime where $`g_sN1`$ we find the meson mass is much smaller than the quark mass. At large t’Hooft coupling we find that the binding energy of the mesons almost completely cancels the rest energy of the quarks. This is similar to the situation in $`AdS_5\times S^5`$ .
## 5 Discussion
In this note we have computed the spectrum of mesons in an $`𝒩=1`$ field theory corresponding to fluctuations in the position of a holomorphically embedded D7-brane. In the limit of nearly massless quarks, the field theory is classically conformal, and also conformal at large-$`N`$, and the spectrum turns out to be computable exactly, where the dimensions in general are complicated irrationals.
There are a few operators for which the exact results are simple. Among these are the lowest mode, corresponding to a mass term, with dimension 3/2, and a mode corresponding to a BPS fluctuation of the D7-brane, with dimension 3. The existence of these operators suggests that a consistent superpotential for our flavored theory is
$`W=qA_1B_1\stackrel{~}{q}.`$ (59)
It would be interesting to study the Klebanov-Strassler theory obtained at the end of the duality cascade with the addition of 3-form flux with this superpotential.
In the strictly massless limit, $`z_1=0`$, it is possible to relax our embedding condition slightly. Specifically, with a nonzero mass we imposed a relation between the azimuthal coordinates, $`\psi \varphi _1\varphi _2=0`$. However, when the mass is zero this condition need not apply; it would be nice to see what relaxing this condition would mean for the field theory (in particular, whether it is possible to realize the cubic superpotential discussed in section 3.)
The appearance of irrational dimensions is not surprising, in light of similar results for the glueball spectrum of the conifold . However, this feature of the meson spectrum differs from the $`𝒩=4`$ case, where the meson dimensions were pure integers. In particular, we do not find a tower of states with spacing 3/2, except in the large R-charge limit; more precisely, in this limit the spectrum is of the form $`3k/2+O(1/J)`$. It might be possible to compute these $`1/J`$ corrections in a plane-wave limit, or perhaps in some other formalism. It would be interesting if such a comparison with our exact results were possible.
We have also numerically computed the spectrum for the case of massive quarks. In the large $`g_sN`$ limit the meson mass gap is significantly smaller than the quark masses. We have uncovered a relatively simple quadratic scaling behavior for the meson masses. It would be nice to find, either with analytical or more numerical work, the exact functional dependence on $`n,m_1,m_2`$ etc.
All of our calculations have been in the probe limit and further studies of the backreaction would be interesting, especially for the Klebanov-Strassler deformed conifold theory. However, it may still be possible to learn things from further study of probe theories. In particular, it would be interesting to study the dynamics of nontrivial classical field configurations in the D7-brane worldvolume. Such fields would correspond to dissolved D3-branes or anti-D3-branes. The anti-brane case is particularly interesting, as it would break supersymmetry along the lines of the KKLT scenario<sup>1</sup><sup>1</sup>1We thank S. Trivedi for this suggestion., but with the possibility for some moduli to be fixed by the D7-brane. We leave these suggestions for the future.
Acknowledgments
We thank A. Maharana for discussions and collaboration during the early stages of our work. We thank H. Elvang, E. Katz, M. Strassler, S. Trivedi, and J. Wacker for discussions and correspondence. We are especially grateful to Joel Giedt and Leo Pando-Zayas for bringing some typos in a previous version to our attention. TSL thanks the Kavli Institute for Theoretical Physics for warm hospitality during much of the preparation of this work. TSL was supported in part by the KITP under National Science Foundation grant PHY99-07949, the National Science Foundation under grants PHY-0331728 and OISE-0443607, and the Department of Energy under grant DE-FG02-95ER40893. The work of P. O. is supported in part by the DOE under grant DOE91-ER-40618 and by the NSF under grant PHY00-98395. |
warning/0506/hep-th0506098.html | ar5iv | text | # BRST-antifield Quantization: a Short Review
## 1 Introduction
Gauge symmetries are omnipresent in theoretical physics, especially in particle physics. Well-known examples of gauge theories are QED and QCD. A common feature of gauge theories is the appearance of unphysical degrees of freedom in the Lagrangian. Because of this, the naive path integral for gauge theories is meaningless since integrating over gauge directions in the measure would make it infinite-valued:
$$DA_\mu e^{iS}=\mathrm{}$$
(1.1)
The redundant gauge variables must be removed from the theory by considering gauge-fixing conditions. When this is done, gauge invariance is of course lost and it is not clear how to control the physics. The modern approach to cope with these problems in the case of general gauge theories was developed by Batalin and Vilkovisky , building on earlier work by Zinn-Justin , Kallosh and de Wit and van Holten . It goes under the name of BV or antifield formalism and is the method explained in these lectures. The BRST symmetry is central to it . A complete coverage of the topic and further references can be found in .
## 2 Structure of gauge symmetries
Let us give a brief overview of the different types of gauge theories that one may encounter.
### 2.1 Yang-Mills type
We are all familiar with Yang-Mills gauge theories. In the absence of matter, the action is given by
$`S_0[A_\mu ^a]`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle F_{\mu \nu }^aF_a^{\mu \nu }d^nx}`$ (2.1)
$`F_{\mu \nu }^a`$ $`=`$ $`_\mu A_\nu ^a_\nu A_\mu ^a+f_{bc}^aA_\mu ^bA_\nu ^c.`$ (2.2)
$`F_{\mu \nu }^a`$ is the field strength tensor, $`A_\mu ^a`$ is the gauge field and $`f_{bc}^a`$ are the structure constants of the associated gauge group. The dimensionality of space-time is $`n`$. The gauge transformation takes infinitesimally the form
$`A_\mu ^a(x)^{}`$ $`=`$ $`A_\mu ^a(x)+\delta _ϵA_\mu ^a`$ (2.3)
$`\delta _ϵA_\mu ^a(x)`$ $`=`$ $`D_\mu ϵ^a(x)=_\mu ϵ^a+f_{bc}^aA_\mu ^bϵ^c(x),`$ (2.4)
where $`ϵ^c(x)`$ is a set of arbitrary functions, the gauge parameters. In these theories, the commutator of infinitesimal gauge transformations reads
$$[\delta _\epsilon ,\delta _\eta ]X=\delta _\xi X,\xi ^a=f_{bc}^a\epsilon ^b\eta ^c$$
(2.5)
with $`\epsilon ,\eta ,\xi `$ gauge parameters and where $`X`$ can be any field. It is clear from this formula that the algebra of the gauge transformations closes off-shell as the commutator of the gauge transformations is again a gauge transformation of the same type, without using the equations of motion.
### 2.2 Closure only on-shell
Off-shell closure holds for Yang-Mills gauge theories but is not a general feature of gauge systems. The gauge transformations might close only when the equations of motion hold. Notable examples where this is the case are extended supergravity theories. Rather than discussing the gauge structure of supergravities, which is rather intricate, we shall illustrate “closure only on-shell” in the case of a much simpler (but of no direct physical interest) system. Consider the following Lagrangian:
$$=\frac{1}{2}(\dot{q}^1\dot{q}^2\dot{q}^3)^2=\frac{1}{2}\dot{y}^2,y=q^1q^2q^3$$
(2.6)
for a model with three coordinates $`q^1`$, $`q^2`$ and $`q^3`$. $``$ is invariant under two different sets of gauge transformations, which can be taken to be:
$$\delta _\epsilon q^1=\epsilon +\epsilon q^2\ddot{y},\delta _\epsilon q^2=\epsilon ,\delta _\epsilon q^3=\epsilon q^2\ddot{y}$$
(2.7)
and
$$\delta _\eta q^1=0,\delta _\eta q^2=\eta ,\delta _\eta q^3=\eta $$
(2.8)
We can easily calculate the commutators of the gauge transformations on the fields:
$$[\delta _\epsilon ,\delta _\eta ]q^1=\epsilon \eta \ddot{y},[\delta _\epsilon ,\delta _\eta ]q^2=0,[\delta _\epsilon ,\delta _\eta ]q^3=\epsilon \eta \ddot{y}.$$
(2.9)
From (2.6), the equation of motion (eom) for $`y`$ is $`\ddot{y}=0`$. We see that the algebra of the gauge transformations (2.7) and (2.8) is closed (in fact, abelian) only up to equations of motion, i.e., only on-shell.
### 2.3 Reducible gauge theories
The gauge transformations might also be “reducible”, i.e., dependent. Consider the theory of an abelian $`2`$-form $`B_{\mu \nu }=B_{\nu \mu }`$. The field strength is given by $`H_{\mu \nu \rho }=_\mu B_{\nu \rho }+_\nu B_{\rho \mu }+_\rho B_{\mu \nu }`$. The Lagrangian reads:
$$=\frac{1}{12}H_{\mu \nu \rho }H^{\mu \nu \rho }$$
(2.10)
and is invariant under gauge transformations
$$\delta _\mathrm{\Lambda }B_{\mu \nu }=_\mu \mathrm{\Lambda }_\nu _\nu \mathrm{\Lambda }_\mu $$
(2.11)
where $`\mathrm{\Lambda }`$ is the gauge parameter. These transformations vanish for a class of parameters $`\mathrm{\Lambda }_\mu =_\mu ϵ`$, meaning that the gauge parameters are not all independent. Such gauge transformations are called reducible, and the corresponding gauge theory is said to be reducible. Two-forms define a natural generalization of electromagnetism, $`A_\mu B_{\mu \nu }`$ and occur in many models of unification; the main difference with electromagnetism being the irreducibility of the latter.
### 2.4 Reducibility on-shell
The last feature that we want to illustrate is the possibility that the reducibility of the gauge transformations holds only on-shell. One can reformulate the previous free $`2`$-form model by introducing an auxiliary field $`A_\mu `$. The Lagrangian is then given by:
$$=\frac{1}{12}A_\mu \epsilon ^{\mu \nu \rho \sigma }H_{\nu \rho \sigma }\frac{1}{8}A_\mu A^\mu .$$
(2.12)
This Lagrangian reduces to (2.10) by inserting the equation of motion for $`A_\mu `$:
$$A_\mu =\frac{1}{3}\epsilon ^{\mu \nu \rho \sigma }H_{\nu \rho \sigma }.$$
(2.13)
The gauge transformations are of the form:
$`\delta _\mathrm{\Lambda }B_{\mu \nu }`$ $`=`$ $`_\mu \mathrm{\Lambda }_\nu _\nu \mathrm{\Lambda }_\mu `$ (2.14)
$`\delta _\mathrm{\Lambda }A_\mu `$ $`=`$ $`0.`$ (2.15)
We can introduce interactions by considering Lie-algebra-valued fields $`A_\mu =A_\mu ^aT_a`$, $`B_{\mu \nu }=B_{\mu \nu }^aT_a`$ ($`T_a`$: generators of the gauge group) and covariant derivatives instead of partial derivatives, $`_\mu D_\mu `$. This is the so-called Freedman-Townsend model .
For the Freedman-Townsend model, the gauge transformations vanish for parameters $`\mathrm{\Lambda }_\mu =D_\mu ϵ`$. But now, this vanishing occurs only if the equations of motion are satisfied. This is due to $`[D_\mu ,D_\nu ]F_{\mu \nu }`$, and $`F_{\mu \nu }=0`$ is the eom for $`B`$. The theory is said in that case to be reducible on-shell.
The antifield-BRST formalism is capable of handling all the gauge structures described here, while the original methods were devised only for off-shell closed, irreducible gauge algebras. This wide range of application of the antifield formalism is one of its main virtues.
Remark: recent considerations on the structure of gauge symmetries, including reducible ones, may be found in .
## 3 Algebraic tools
BRST theory uses crucially cohomological ideas and tools. In the following, some definitions are collected and a useful technique for the computation of cohomologies is illustrated.
### 3.1 Cohomology
Let us consider a nilpotent linear operator $`D`$ of order $`2`$: $`D^2=0`$. Because of this property, the image of $`D`$ is contained in the kernel of $`D`$, Im $`D`$ $``$ Ker $`D`$. The cohomology of the operator $`D`$ is defined as the following quotient space:
$$H(D)\frac{\text{Ker}D}{\text{Im}D}.$$
(3.1)
### 3.2 De Rham $`d`$
As a familiar example, let us discuss the de Rham $`d`$-operator. This will enable us to introduce further tools.
In a coordinate patch, a $`p`$-form is an object of the form
$$w=\frac{1}{p!}w_{i_1\mathrm{}i_p}dx^{i_1}\mathrm{}dx^{i_p}$$
(3.2)
where the coefficients $`w_{i_1\mathrm{}i_p}`$ are totally antisymmetric functions of the coordinates and $``$ refers to the exterior product. \[Appropriate transition conditions should hold in the overlap of two patches, but these will not be discussed here.\] The vector space of $`p`$-forms on $`M`$ is denoted by $`\mathrm{\Omega }^p(M)`$. The direct sum of $`\mathrm{\Omega }^p(M)`$, $`p=0,\mathrm{},m\text{dim}M`$ defines the space of all forms on $`M`$:
$$\mathrm{\Omega }^{}(M)\mathrm{\Omega }^0(M)\mathrm{}\mathrm{\Omega }^m(M)$$
(3.3)
The exterior derivative $`\text{d}_p`$ is a map $`\mathrm{\Omega }^p(M)\mathrm{\Omega }^{p+1}(M)`$ whose action on a $`p`$-form $`w`$ is defined by<sup>1</sup><sup>1</sup>1The subscript $`p`$ is often not written; the exterior derivative is then referred as d.
$$\text{d}_pw=\frac{1}{p!}\frac{w_{i_1\mathrm{}i_p}}{x^j}dx^jdx^{i_1}\mathrm{}dx^{i_p}.$$
(3.4)
Important properties of d are its nilpotency of order two
$$\text{d}^2=0(\text{i.e.}\text{d}_{p+1}d_p=0)$$
(3.5)
and the fact that it is an odd derivative:
$$\text{d}(w\eta )=\text{d}w\eta +(1)^pw\text{d}\eta .$$
(3.6)
The nilpotency can easily be proved by direct computation:
$$\text{d}^2w=\frac{1}{p!}\frac{^2w_{i_1\mathrm{}i_p}}{x^kx^j}dx^kdx^jdx^{i_1}\mathrm{}dx^{i_p}.$$
(3.7)
This expression clearly vanishes since the coefficients are symmetric in $`k,j`$ while $`dx^kdx^j`$ is antisymmetric. Therefore, Im $`\text{d}_p`$ $``$ Ker $`\text{d}_{p+1}`$.
An element of Ker $`\text{d}_p`$ is said to be “closed” or “a cocycle” (of d), d$`\alpha =0`$. An “exact form” or “coboundary” (of d) lives in Im $`\text{d}_{p1}`$; it is thus such that $`\alpha =\text{d}\beta `$, for some $`(p1)`$-form $`\beta `$. The $`p`$th de Rham cohomology group is defined as:
$$H^p(\text{d})=\frac{\text{Ker}\text{d}_p}{\text{Im}\text{d}_{p1}}.$$
(3.8)
Another important operation is the interior product $`\text{i}_X`$: $`\mathrm{\Omega }^p(M)\mathrm{\Omega }^{p1}(M)`$, where $`X=X^j/x^j`$. The action of $`\text{i}_X`$ on a $`p`$-form $`w`$ reads:
$$\text{i}_Xw=\frac{1}{(p1)!}X^jw_{ji_2\mathrm{}i_p}dx^{i_2}\mathrm{}dx^{i_p}.$$
(3.9)
The interior product also satisfy (3.5) and (3.6) as well.
### 3.3 Poincaré lemma
A useful tool for computing cohomologies is given by contracting homotopies. We illustrate the techniques by computing the cohomology of $`d`$ in a special case.
Let $`M`$ be $`^n`$, and $`w`$ a closed $`p`$-form on M, $`p>0`$. The Poincaré lemma states that all closed forms in degree $`>0`$ are exact<sup>2</sup><sup>2</sup>2For $`M^n`$, closed forms might not be globally exact (although they are locally so). The Poincaré lemma fails in such a case.. More precisely:
$`H^p(\text{d})`$ $`=`$ $`0p>0`$ (3.10)
$`H^p(\text{d})`$ $``$ $`p=0.`$
Proof:
Assume the coefficients $`w_{i_1\mathrm{}i_p}`$ to be polynomial in $`x^i`$ (this restriction is not necessary and is made only to simplify the discussion). Let:
$$\text{d}=dx^i\frac{}{x^i},i_x=x^i\frac{^L}{(dx^i)}(i_X,X^j=x^j).$$
(3.11)
Define the counting operator $`N`$ as the combination<sup>3</sup><sup>3</sup>3It is the Lie derivative of a form along $`X`$, $`Nw=_Xw`$. $`N=\text{d}i_X+i_X\text{d}`$. It is such that $`Nx^i=x^i`$, $`Ndx^i=dx^i`$. For a general form $`\alpha `$ we have $`N\alpha =k\alpha `$ with $`k`$ the total polynomial degree (in $`x^i`$ and $`dx^i`$). It then follows that for:
* $`k0`$:
$$\alpha =\frac{k}{k}\alpha =N\left(\frac{1}{k}\alpha \right)=(\text{d}i_x+i_x\text{d})\left(\frac{1}{k}\alpha \right).$$
Thus, if $`\text{d}\alpha =0`$ then $`\alpha =\text{d}(i_x\left(\frac{1}{k}\alpha \right))`$.
* $`k=0`$: $`\text{Im}\text{d}_1`$ has no meaning, there is no such a thing as a $`(1)`$-form. We can say that $`\mathrm{\Omega }^1(M)`$ is empty and $`H^0(\text{d})=\text{Ker}\text{d}_0`$. So constants are the only members of the cohomology.
This proves the Poincaré lemma.
### 3.4 Local functions
A local function $`f`$ is a smooth function of the spacetime coordinates, the field variables and their respective derivatives up to a finite order, $`f=f(x,[\phi ])=f(x^\mu ,\phi ^i,_\mu \phi ^i,\mathrm{},_{\mu _1\mathrm{}\mu _k}\phi ^i)`$. In field theory local functions are usually polynomial in the derivatives. The following discussion is however more general than that, and it remains valid in the case of arbitrary smooth local functions.
The Euler-Lagrange derivative $`\frac{\delta }{\delta \phi ^i}`$ of a local function $`f`$ is defined by
$$\frac{\delta f}{\delta \phi ^i}=\underset{k0}{}()^k_{\mu _1}\mathrm{}_{\mu _k}\frac{f}{(_{\mu _1}\mathrm{}_{\mu _k}\phi ^i)}$$
(3.12)
(with $`_{\mu _1}\mathrm{}_{\mu _k}\phi ^i)`$ the last derivative of $`\phi `$ occurring in $`f`$).
Theorem : A local function is a total derivative iff it has vanishing Euler-Lagrange derivatives with respect to all fields:
$$f=_\mu j^\mu \frac{\delta f}{\delta \phi ^i}=0\phi ^i.$$
(3.13)
A proof of this theorem will not be given here but can be found for instance in (see also references given in ).
### 3.5 Local differential forms
Local $`p`$-forms are differential forms whose coefficients are local functions:
$$w=\frac{1}{p!}w_{i_1\mathrm{}i_p}(x,[\phi ])dx^{i_1}\mathrm{}dx^{i_p}$$
(3.14)
Consider a local $`n`$-form $`w=fd^nx`$ in $`^n`$; $`w`$ is trivially closed. It is further exact iff the function $`f`$ is a total derivative:
$$w=\text{d}\alpha f=_\mu j^\mu \frac{\delta f}{\delta \phi ^i}=0\phi ^i$$
(3.15)
### 3.6 Algebraic Poincaré lemma
The algebraic Poincaré lemma gives the cohomology of d in the algebra of local forms:
$`p=n:H^p(\text{d})`$ $``$ $`0\text{and characterized above}`$
$`0<p<n:H^p(\text{d})`$ $`=`$ $`0`$ (3.16)
$`p=0:H^p(\text{d})`$ $``$ $`.`$
The distinguishing feature compared with the Poincaré lemma for ordinary exterior forms not depending on local fields is the appearance of a non-vanishing cohomology at $`p=n`$; non trivial local $`n`$-forms $`w=fd^nx`$ are such that at least one of the derivatives $`\delta f/\delta \phi ^i`$ is not identically zero. A proof of the algebraic Poincaré lemma in the case of polynomial dependence on derivatives can be found in ; for a more general proof see .
Example. Consider local $`1`$-forms in $`^1`$, $`w=(t,q,\dot{q},\ddot{q},\mathrm{})dt`$. For example:
$$w_1=_1dt=\frac{1}{2}\dot{q}^2dt,w_2=_2dt=\dot{q}qdt.$$
(3.17)
These are obviously closed, but are they exact? We work out the Euler-Lagrange derivatives of $`_{1,2}`$:
$$\frac{\delta _1}{\delta q}=\ddot{q},\frac{\delta _2}{\delta q}=0.$$
(3.18)
Therefore, from (3.15), $`w_2`$ is an exact form (and indeed, $`w_2=\text{d}\left(\frac{1}{2}q^2\right)`$), while $`w_1`$ cannot be written as the exterior derivative of a local function. $`H^1(\text{d})`$ is clearly non trivial.
## 4 BRST construction
We stated in the introduction that gauge invariance is lost after the necessary gauge fixing. The central idea of the BRST construction is to replace the original gauge symmetry by a rigid symmetry, the BRST symmetry $`s`$, which is still present even after one has fixed the gauge. This is achieved by introducing extra fields in the theory: the ghost fields and the conjugate antifields. The operator $`s`$ acts on the enlarged space of fields, ghosts and antifields. An extended action involving all these variables can be constructed in such a way that it is BRST invariant. The operator $`s`$ is called the BRST differential and it is nilpotent: $`s^2=0`$. Therefore, cohomological groups $`H^k(s)`$ can be constructed. The BRST differential fulfills:
$$H^0(s)=\text{Gauge invariant functions (“Observables”)}$$
(4.1)
In this way we recover the gauge symmetry. This is BRST theory in a nutshell. These important statements will now be discussed in more detail.
### 4.1 Master equation
Consider a gauge theory of fields $`\phi ^i`$ described by a classical action $`S_0(\phi ^i)`$ on a manifold $`M`$. The equations of motion constrain the fields to a submanifold, denoted $`\mathrm{\Sigma }`$. The action is invariant under gauge transformations<sup>4</sup><sup>4</sup>4We use De Witt’s condensed notation; see Appendix I.
$$\delta _\epsilon \phi ^i=R_\alpha ^i\epsilon ^\alpha .$$
(4.2)
Assume the theory to be (on-shell) reducible, with no reducibility on the reducibility functions. In such a case there are relations among the gauge parameters but no relations among the relations. The relations among the gauge parameters can be written as:
$$Z_\mathrm{\Delta }^\alpha R_\alpha ^i=C_\mathrm{\Delta }^{ij}\frac{\delta S_0}{\delta \phi ^j}.$$
(4.3)
We proceed as follows. For each commuting (anticommuting) gauge parameter $`\epsilon ^\alpha `$ one introduces a fermionic (bosonic) ghost variable $`c^\alpha `$. We also introduce ghosts of ghosts $`c^\mathrm{\Delta }`$, one for each (independent) reducibility identity of the theory. The set of original fields, ghosts and ghosts of ghosts are collectively denoted as $`\mathrm{\Phi }^A`$. We double now the configuration space by considering conjugate fields w.r.t. each of the $`\mathrm{\Phi }`$’ s: the anti-fields $`\mathrm{\Phi }_A^{}`$. They are postulated to have opposite (Grassmann) parity. Gradings are assigned to the various fields as displayed in table 1.
We define the antibracket<sup>5</sup><sup>5</sup>5Properties of anti-brackets are listed in Appendix II. of two functionals $`F(\mathrm{\Phi }^A,\mathrm{\Phi }_A^{})`$, $`G(\mathrm{\Phi }^A,\mathrm{\Phi }_A^{})`$ by:
$$(F,G)=\frac{\delta ^RF}{\delta \mathrm{\Phi }^A}\frac{\delta ^LG}{\delta \mathrm{\Phi }_A^{}}\frac{\delta ^RF}{\delta \mathrm{\Phi }_A^{}}\frac{\delta ^LG}{\delta \mathrm{\Phi }^A}.$$
(4.4)
Here, $`L`$ (respectively $`R`$) refers to the standard left (respectively right) derivative. These are related as
$$\frac{\delta ^RF}{\delta (\text{field})}(1)^{\epsilon _{\text{field}}(\epsilon _F+1)}\frac{\delta ^LF}{\delta (\text{field})},$$
(4.5)
where $`\epsilon `$ denotes the Grassmann-parity. The BRST transformation of any functional $`F(\mathrm{\Phi }^A,\mathrm{\Phi }_A^{})`$ can be written in terms of antibrackets:
$$sF=(S,F).$$
(4.6)
The generating function $`S(\mathrm{\Phi }^A,\mathrm{\Phi }_A^{})`$ of the BRST transformation is sometimes called the generalized action. The BRST transformation is nilpotent of order two; this is reflected in the (classical) master equation:
(S,S)=0 (4.7)
The solution to the master equation is unique up to canonical transformations. It can be constructed in a sequential form<sup>6</sup><sup>6</sup>6There is however no guarantee for this sequence to be finite!:
$`S`$ $`=`$ $`S_0+S_1+S_2+\mathrm{}`$
$`S_0\text{classical action},S_1`$ $`=`$ $`\phi _i^{}R_\alpha ^ic^\alpha ,S_2=c_\alpha ^{}Z_\mathrm{\Delta }^\alpha c^\mathrm{\Delta }+\mathrm{}`$
The proof of this statement (including the reducible case) can be found in and references therein. Locality of $`S`$ under general conditions is established in .
The solution $`S`$ of the master equation is key to the BRST-antifield formalism. It can be written down explicitly for the Yang-Mills theory and the abelian 2-form model introduced earlier:
Yang-Mills.
$$S=\frac{1}{4}d^nxF_{\mu \nu }^aF_a^{\mu \nu }+d^nxA_a^\mu D_\mu c^a+\frac{1}{2}d^nxc_a^{}f_{bc}^ac^bc^c$$
(4.9)
Abelian $`\mathrm{𝟐}`$-form.
$$S=\frac{1}{12}d^nxH_{\mu \nu \rho }H^{\mu \nu \rho }+d^nxB^{\mu \nu }(_\mu c_\nu _\nu c_\mu )+d^nxc^\mu _\mu c$$
(4.10)
In those cases, the solution $`S`$ of the master equation is linear in the antifields. For gauge systems with an “open algebra” (i.e., for which the gauge transformations close only on-shell), or for on-shell reducible gauge theories, the solution of the master equation is more complicated. It contains terms that are indeed non linear in the antifields. These terms are essential for getting the correct gauge fixed action below. Without them, one would not derive the correct Feynman rules leading to gauge-independent amplitudes.
Exercises.
1.) Consider a nilpotent operator $`\mathrm{\Omega }`$ of order $`N`$ (i.e. $`\mathrm{\Omega }^N0`$).
(i) prove that the only eigenvalue of $`\mathrm{\Omega }`$ is zero.
(ii) for $`N=2`$, analyze the cohomology of $`\mathrm{\Omega }`$ in terms of its Jordan decomposition.
(iii) for $`N=3`$, Im $`\mathrm{\Omega }^2`$ $``$ Ker $`\mathrm{\Omega }`$. The corresponding cohomologies are defined as
$$H_{(1)}(\mathrm{\Omega })\frac{\text{Ker}\mathrm{\Omega }}{\text{Im}\mathrm{\Omega }^2},H_{(2)}(\mathrm{\Omega })\frac{\text{Ker}\mathrm{\Omega }^2}{\text{Im}\mathrm{\Omega }}$$
Calculate these.
Hint: The Jordan decomposition of a matrix is a block-diagonal form. Each such block, called Jordan block, has on its diagonal always the same eigenvalue and 1 in the upper secondary diagonal.
2.) Prove that P$`(\phi ^i,_\mu \phi ^i,\mathrm{},_{\mu _1\mathrm{}\mu _k}\phi ^i)d^nx`$ is exact iff $`\frac{\delta P}{\delta \phi ^i}=0`$, where $`\frac{\delta }{\delta \phi ^i}`$ is the Euler-Lagrange derivative.
Hint: Relate $`N=\phi ^i\frac{}{\phi ^i}+(_\mu \phi ^i)\frac{}{(_\mu \phi ^i)}+\mathrm{}`$ to $`\frac{\delta P}{\delta \phi ^i}`$
3.) Write explicitly $`R_\alpha ^i`$, $`Z_\mathrm{\Delta }^\alpha `$, $`C_\mathrm{\Delta }^{ij}`$ for the Freedman-Townsend model.
4.) Write Noether’s identities (see Appendix I) for Yang-Mills, gravity and an abelian $`2`$-form gauge theory.
5.) Check the properties of anti-brackets given in Appendix II.
6.) Consider an irreducible gauge theory with gauge transformations closing off-shell and forming a group. The solution of the master equation is given by
$$S=S_0+\phi _i^{}R_\alpha ^ic^\alpha +\frac{1}{2}c_\alpha ^{}f_{\beta \gamma }^\alpha c^\beta c^\gamma $$
(4.11)
where $`f_{\beta \gamma }^\alpha `$ are the structure constants of the gauge group. Verify that $`S`$ satisfies the master equation $`(S,S)=0`$.
7.) Define the operator $`\mathrm{\Delta }`$ as
$$\mathrm{\Delta }F=(1)^{\epsilon _A}\frac{\delta ^L}{\delta \varphi ^A}\frac{\delta ^LF}{\delta \varphi _A^{}}.$$
(4.12)
Prove the following statements:
(i) $`\epsilon (\mathrm{\Delta })=1`$
(ii) $`\mathrm{\Delta }^2=0`$
(iii) gh $`\mathrm{\Delta }=1`$
(iv) $`\mathrm{\Delta }(\alpha ,\beta )=(\mathrm{\Delta }\alpha ,\beta )(\alpha ,\mathrm{\Delta }\beta )(1)^{\epsilon _\alpha }`$
(v) $`\mathrm{\Delta }(\alpha \beta )=(\mathrm{\Delta }\alpha )\beta +(1)^{\epsilon _\alpha }\alpha (\mathrm{\Delta }\beta )+(1)^{\epsilon _\alpha }(\alpha ,\beta )`$
Verify that the superjacobian for the change of variables $`\varphi ^A\varphi ^{}_{}{}^{}A+(\mu S,\varphi ^A)`$ is $`1(\mathrm{\Delta }S)\mu `$, where $`\mu `$ is a fermionic constant.
## 5 Observables
It is now time to substantiate the claim
$$H^0(s)\mathrm{Observables}$$
(5.1)
where, as we have just seen, the BRST differential is given by $`sF=(S,F)`$. In particular,
$`s\mathrm{\Phi }^A`$ $`=`$ $`(S,\mathrm{\Phi }^A)={\displaystyle \frac{\delta ^RS}{\delta \mathrm{\Phi }_A^{}}}={\displaystyle \underset{k}{}}{\displaystyle \frac{\delta ^RS_k}{\delta \mathrm{\Phi }_A^{}}}`$ (5.2)
$`s\mathrm{\Phi }_A^{}`$ $`=`$ $`(S,\mathrm{\Phi }_A^{})={\displaystyle \frac{\delta ^RS}{\delta \mathrm{\Phi }^A}}={\displaystyle \underset{k}{}}{\displaystyle \frac{\delta ^RS_k}{\delta \mathrm{\Phi }^A}}.`$ (5.3)
Note that the ghost number of $`S`$ is 0, the ghost number of the BRST transformation is 1 as well as the one of the antibracket, so that the gradings of both sides of the equation $`sF=(S,F)`$ match. We shall actually not provide the detailed proof of (5.1) here, but instead, we give only the key ingredients that underlie it, referring again to for more information.
To that end, we expand the BRST transformations of all the variables according to the antifield number, as in . So one has,
$`S`$ $`=`$ $`{\displaystyle \underset{k0}{}}S_k`$ (5.4)
$`s`$ $`=`$ $`\delta +\gamma +{\displaystyle \underset{i>0}{}}s_i.`$ (5.5)
The first term $`\delta `$ has antifield number -1 and is called the Koszul-Tate differential, the second term $`\gamma `$ has antifield number 0 and is called the longitudinal differential, and the next terms $`s_i`$ have antifield number $`i`$. Although the expansion stops at $`\delta +\gamma `$ for Yang-Mills (as it follows from the solution of the master equation given above), higher order terms are present for gauge theories with an open algebra, or on-shell reducible theories<sup>7</sup><sup>7</sup>7The expansion of $`s`$ is connected to spectral sequences, which will not be discussed in detail here ..
Explicitly, one finds for the Koszul-Tate differential $`\delta \phi ^i=0`$, as there is no field operator of anti-field number -1 and $`\delta \phi _i^{}=\delta S_0/\delta \phi ^i`$ . For the longitudinal differential it follows in the same way that
$$\gamma \phi ^i=R_\alpha ^ic^\alpha .$$
(5.6)
Observe also that
$$0=s^2=\delta ^2+\{\delta ,\gamma \}+(\gamma ^2+\{\delta ,s_1\})+\mathrm{}$$
(5.7)
In this equation, each term has to vanish separately, as each term is of different antifield number.
Let $`A`$ be a BRST-closed function(al), $`sA=0`$. We must compute the equivalence class
$$AA+sB.$$
(5.8)
Since we are dealing with observables, the only relevant operators are of of ghost number 0, thus $`\mathrm{gh}A=0`$ and $`\mathrm{gh}B=1`$. The latter can only be satisfied, if $`B`$ contains at least one anti-field. Expanding $`A`$ and $`B`$ in antifield number yields
$`A={\displaystyle \underset{k0}{}}A_k=A_0+A_1+A_2+\mathrm{}`$ (5.9)
$`B={\displaystyle \underset{k1}{}}B_k=B_1+B_2+B_3+\mathrm{}`$ (5.10)
Acting with $`s`$ on $`A`$ using the expansion (5.5) then gives
$$(\delta +\gamma +\mathrm{})(A_0+A_1+\mathrm{})=(\gamma A_0+\delta A_1)+\mathrm{},$$
(5.11)
where the term in parentheses collects all antifield number zero contributions. The condition $`sA=0`$ implies that this term must vanish on its own and thus $`\gamma A_0=\delta A_1`$. Furthermore,one finds
$$A+sB=A+(\delta +\gamma +\mathrm{})B=(A_0+\delta B_1)+\mathrm{}$$
(5.12)
where the last term in parentheses is again the antifield zero contribution. Using (5.3),
$$\delta B_1=\frac{\delta B_1}{\delta \mathrm{\Phi }^{}}\frac{\delta S_0}{\delta \mathrm{\Phi }^i},$$
(5.13)
we see that the second term of the antifield zero contribution in $`A+sB`$ vanishes when the equations of motions are fulfilled, i.e. on-shell. A similar property holds for $`\delta A_1`$. There is therefore a clear connection of $`\delta `$ to the dynamics and the equations of motions. Note that the “on-shell functions” can be viewed as the equivalence classes of functions on $`M`$ identified when they coincide on $`\mathrm{\Sigma }`$, i.e., $`C^{\mathrm{}}(\mathrm{\Sigma })=C^{\mathrm{}}(M)/𝒩`$, where the ideal $`𝒩`$ contains all the functions that vanish on-shell.
From (5.6) and the fact that $`\gamma `$ is a derivation, one gets for<sup>8</sup><sup>8</sup>8$`A_0`$ can only depend on $`\varphi _i`$, as it has pure ghost and antifield number 0.$`\gamma A_0`$
$$\gamma A_0=\frac{\delta A_0}{\delta \phi ^i}R_\alpha ^ic^\alpha .$$
(5.14)
As this is a gauge transformation (with gauge parameters replaced by the ghosts), $`\gamma A_0`$ vanishes if $`A_0`$ is gauge-invariant. The longitudinal differential is associated with gauge transformations. We thus see that a necessary condition for $`A`$ to be BRST-closed is that its first term $`A_0`$ be gauge-invariant on-shell. And furthermore, two such $`A_0`$’s are equivalent when they coincide on-shell.
It turns out that the condition on $`A_0`$ is also sufficient for $`A`$ to be BRST-closed, in the sense that given an $`A_0`$ that is gauge-invariant on-shell, one can complete it by terms $`A_1`$, $`A_2`$ … of higher antifield number so that $`sA=0`$.
To summarize: the term that determines the cohomological class of a BRST cocycle is the first term $`A_0`$. This term must be an observable, in that it must be gauge invariant on-shell. We can therefore conclude that $`H^0(s)`$ captures indeed the concept of observables. The differential $`\delta `$ reduces from the manifold $`M`$ to the on-shell manifold $`\mathrm{\Sigma }`$ and $`\gamma `$ further to $`\mathrm{\Sigma }/G`$, the set of all gauge-invariant functions, where $`G`$ is the set of all gauge orbits.
## 6 Path Integral and Gauge-fixing
We first consider the Yang-Mills case. To perform actual path-integral calculations, it is necessary to gauge-fix the theory. To perform this task, it is convenient to add additional fields, the anti-ghost $`\overline{c}_a`$ and auxiliary fields, the Nakanishi-Lautrup fields $`b_a`$. They transform as $`s\overline{c}_ab_a`$ and $`sb_a=0`$. We take $`\overline{c}_a`$ and $`b_a`$ to have ghost number -1 and 0, respectively. The corresponding antifields $`\overline{c}^a`$ and $`b^a`$ have thus ghost number 0 and -1, respectively. Furthermore a contracting homotopy argument similar to the one given above for the Poincaré lemma shows that the counting operator of $`\overline{c}_a`$, $`b_a`$ and their conjugate antifields is BRST exact. Hence the cohomology is not altered by the introduction of these new variables. In particular, the set of observables is not affected. The solution of the master equation with the new variables included reads, for Yang-Mills theory
$`S`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle d^nxF_{\mu \nu }^aF_a^{\mu \nu }}+{\displaystyle d^nxA_a^\mu D_\mu c^a}`$ (6.1)
$`+{\displaystyle \frac{1}{2}}{\displaystyle d^nxc_a^{}f_{bc}^ac^bc^c}i{\displaystyle d^nx\overline{c}^ab_a}.`$
The last term is called the non-minimal part.
Theorem
The generating functional
$$𝒵=𝒟\mathrm{\Phi }^A\mathrm{exp}\left(\frac{i}{\mathrm{}}S_\psi [\mathrm{\Phi }^A]\right),$$
(6.2)
does not depend on the choice of $`\psi `$. Here, $`\psi `$ is called the gauge-fixing fermion, and has Grassmann-parity 1 (hence its name) and ghost number -1. In (6.2), the notation
$`\mathrm{\Phi }^a`$ $`=`$ $`(A_\mu ^a,c^a,\overline{c}^a,b_a)`$ (6.3)
$`\mathrm{\Phi }_a^{}`$ $`=`$ $`(A_\mu ^a,c^a,\overline{c}^a,b^a),`$ (6.4)
has been used and the “gauge-fixed action” $`S_\psi [\mathrm{\Phi }^A]`$ is given by
$$S_\psi [\mathrm{\Phi }^A]=S\left[\mathrm{\Phi }^A,\mathrm{\Phi }_A^{}=\frac{\delta \psi }{\delta \mathrm{\Phi }^A}\right].$$
(6.5)
This theorem is proved in section 8 below.
Before turning to the proof, we want to illustrate formula (6.2) by showing how one can choose the gauge-fixing fermion $`\psi `$ to reproduce familiar expressions for the path integral of the Yang-Mills field. A possible choice, which leads to non-degenerate propagators for all fields and ghosts, is given by
$$\psi =id^nx\overline{c}^a\left(^a+\frac{\alpha }{2}b^a\right),$$
(6.6)
where $`^a`$ is the gauge condition, e.g. $`^a=^\mu A_\mu ^a`$ for covariant gauges and $`\alpha `$ is the gauge parameter. This leads to
$`\overline{c}^a`$ $`={\displaystyle \frac{\delta \psi }{\delta \overline{c}_a}}=`$ $`i\left(^a+{\displaystyle \frac{\alpha }{2}}b^a\right)`$ (6.7)
$`A_{a\mu }^{}`$ $`={\displaystyle \frac{\delta \psi }{\delta A_\mu ^a}}=`$ $`i_\mu \overline{c}_a`$ (6.8)
$`b^a`$ $`={\displaystyle \frac{\delta \psi }{\delta b_a}}=`$ $`i{\displaystyle \frac{\alpha }{2}}\overline{c}^a`$ (6.9)
$`c_a^{}`$ $`={\displaystyle \frac{\delta \psi }{\delta c^a}}=`$ $`0.`$ (6.10)
One then gets the familiar gauge-fixed Yang-Mills action
$`S_\psi [A_\mu ^a,c^a,\overline{c}_a,b_a]`$
$`={\displaystyle d^nx\left(\frac{1}{4}F_{\mu \nu }^aF_a^{\mu \nu }i^\mu \overline{c}_aD_\mu c^a+\left(^a+\frac{\alpha }{2}b^a\right)b_a\right)},`$ (6.11)
which is usually obtained by the Fadeev-Popov procedure . The conventional Landau gauge is recovered by setting $`\alpha =0`$. As the resulting path integral can be written as
$$𝒟[A_\mu ^a\overline{c}^ac^a]\delta (^\mu A_\mu ^a)e^{\left(\frac{i}{\mathrm{}}S_{gf}[A_\mu ^a,\overline{c}^a,c^a]\right)},$$
(6.12)
the transversality of the gauge boson is directly implemented. Landau gauge is thus called a strict gauge. An example of a non-strict gauge is Feynman gauge, $`\alpha =1`$, in which case the equations of motion do not imply $`^a=0`$ but yield instead $`b^a^a`$.
The choice of the gauge-fixing fermion is not unique. One can add to $`\psi `$ the term $`\psi \overline{c}_a\overline{c}_bc^c`$, which yields quartic ghost couplings. Quartic ghost renormalizations may even be needed without such explicit terms, e.g. when using $`^a=^\mu A_\mu ^a+d_{bc}^aA_\mu ^bA^{\mu c}`$, where $`d_{bc}^a`$ is a symmetric tensor in color space (see ).
By appropriately choosing the gauge fixing fermion, one can reduce the path integral to an expression that involves only the physical (transverse) degrees of freedom and which is manifestly unitary in the physical subspace (equal to the reduced phase space path integral). Independence on the choice of $`\psi `$ (still to be proved) guarantees then that the expression (6.2) is correct. We shall not demonstrate here the equivalence of (6.2) with the reduced phase space path integral. The reader may find a discussion of that point in .
As a final point, we note that in order for (6.2) to be indeed independent on the choice of $`\psi `$, it is necessary that the measure be BRST invariant. This can be investigated using the operator $`\mathrm{\Delta }`$, already defined in the exercises in (4.12) as
$$\mathrm{\Delta }=(1)^{ϵ_A}\frac{\delta ^L}{\delta \mathrm{\Phi }^A}\frac{\delta ^L}{\delta \mathrm{\Phi }_A^{}}.$$
(6.13)
The BRST transformation can be written as
$$\mathrm{\Phi }^A\mathrm{\Phi }^A^{}=\mathrm{\Phi }^A+(\mu S,\mathrm{\Phi }^A)=\mathrm{\Phi }^A(\mathrm{\Phi }^A,S)\mu =\mathrm{\Phi }^A\frac{\delta ^LS}{\delta \mathrm{\Phi }_A^{}}\mu ,$$
(6.14)
where $`\mu `$ is a constant, anti-commuting parameter. The Jacobian of this transformation is given by
$$J_{AB}=\frac{\delta ^L\mathrm{\Phi }^A^{}}{\delta \mathrm{\Phi }^B}=\delta ^{AB}\frac{\delta ^L}{\delta \mathrm{\Phi }^B}\frac{\delta ^LS}{\delta \mathrm{\Phi }_A^{}}\mu .$$
(6.15)
\[As the Jacobian involves commuting and anti-commuting fields, the Jacobian “determinant” is actually a super-determinant.\] Therefore the measure transforms as
$$𝒟\mathrm{\Phi }^A\mathrm{sdet}J𝒟\mathrm{\Phi }^A^{}.$$
(6.16)
For an infinitesimal transformation, the super-determinant can be approximated by the super-trace
$$\mathrm{sdet}J1+\mathrm{str}\left(\frac{\delta ^L\delta ^LS}{\delta \mathrm{\Phi }_B\delta \mathrm{\Phi }^A}\mu \right)=1(1)^{ϵ_A}\frac{\delta ^L\delta ^LS}{\delta \mathrm{\Phi }_A\delta \mathrm{\Phi }^A}\mu =1+\mathrm{\Delta }S.$$
(6.17)
It follows that the measure is BRST-invariant iff $`\mathrm{\Delta }S=0`$. The property $`\mathrm{\Delta }S=0`$ can be shown by explicit calculation for pure Yang-Mills theory . The more general case will be treated in the last section. Further interesting properties of $`\mathrm{\Delta }`$ and of the formalism are developed in .
Exercises.
8.) For Yang-Mills theory in 4 dimensions, compute the dimensionality of all fields. Is there any freedom? What is the most general gauge fixing fermion $`\psi `$ of mass dimension 3? What are the restrictions on the gauge-fixing fermion, if the action is required to be invariant under the transformation $`\overline{c}^a\overline{c}^a+ϵ^a`$?
9.) (a) Show that for a functional $`W`$ with $`\mathrm{gh}W=0`$ and Grassmann-parity 0
$$\mathrm{\Delta }e^{\frac{i}{\mathrm{}}W}=0\frac{1}{2}(W,W)i\mathrm{}\mathrm{\Delta }W=0,$$
(6.18)
where $`\mathrm{\Delta }`$ is the operator defined by (6.13).
(b) Define the operator $`\sigma `$ as
$$\sigma \alpha (W,\alpha )i\mathrm{}\alpha ,$$
(6.19)
where $`W`$ satisfies (6.18). Then show that
$$\sigma \alpha =0\mathrm{\Delta }(\alpha e^{\frac{i}{\mathrm{}}W})=0.$$
(6.20)
(c) Show that $`\sigma `$ is nilpotent, $`\sigma ^2=0`$.
(d) Show that if $`\alpha =\sigma \beta `$, then $`\alpha \mathrm{exp}iW/\mathrm{}=\mathrm{\Delta }`$(something).
10.) Write explicitly the action of $`s`$, $`\delta `$, and $`\gamma `$ on all fields and antifields for Yang-Mills theory.
11.) For Yang-Mills theory, consider $`H(\gamma )`$ in the space of polynomials in the ghost fields $`c^a`$, i.e. $`a_0+a_ac^a+a_{ab}c^ac^b+\mathrm{}`$. Compute $`H^0(\gamma )`$ and $`H^1(\gamma )`$. Show in particular that $`H^2(\gamma )`$ parameterizes the non-trivial central extensions $`h_{ab}`$, i.e., non trivial modifications of the algebra of the form $`[X_a,X_b]=f_{bc}^aX_a+h_{ab}1`$.
## 7 Beyond Yang-Mills
The results of the previous section generalize straightforwardly to gauge theories other than Yang-Mills. The solution of the master equation $`S`$ is of the form
$$S=S_0+\varphi _i^{}R_\alpha ^ic^\alpha +\mathrm{},$$
(7.1)
where $`S_0`$ is the classical action. The second term is uniquely determined by the gauge transformation (4.2), and all further terms depend on the specific theory. While the expansion of the solution of the master equation stops at antifield number one in the Yang-Mills case, one gets higher order terms in the case of open gauge systems. It is again often convenient to extend to the non-minimal sector by introducing $`\overline{c}_\alpha `$ and $`b_\alpha `$ and the corresponding antifields in a similar manner to what has been done in the case of Yang-Mills theory.
Assuming again $`\mathrm{\Delta }S=0`$, the quantized theory follows from a path integral
$$𝒵=𝒟\mathrm{\Phi }^A\mathrm{exp}\left(\frac{i}{\mathrm{}}S_\psi [\mathrm{\Phi }^A]\right).$$
(7.2)
$`S_\psi `$ is the solution $`S`$ of the master equation $`(S,S)=0`$, in which the antifields have been eliminated by use of the gauge-fixing fermion $`\psi `$ as before,
$$S_\psi [\mathrm{\Phi }^A]=S_\psi \left[\mathrm{\Phi }^A,\mathrm{\Phi }_A^{}=\frac{\delta \psi }{\delta \mathrm{\Phi }^A}\right].$$
(7.3)
The gauge-fixing fermion has again odd Grassmann-parity and ghost number -1. It is in general given by a local expression
$$\psi =d^nx\chi (\mathrm{\Phi }^A,_\mu \mathrm{\Phi }^A,\mathrm{},_{\mu _1}\mathrm{}_{\mu _k}\mathrm{\Phi }^A).$$
(7.4)
For theories with an open algebra, the terms quadratic in the antifields will lead to quartic (or higher) ghost-antighost vertices in the gauge-fixed action. While these terms are a gauge-dependent option in the Yang-Mills case, they have an unavoidable character (in relativistic gauges) for open gauge algebras. These terms, which follow directly from the general construction of the gauge-fixed action, cannot be obtained through the exponentiation of a determinant, since this procedure always produces an expression which is quadratic in the ghosts.
A useful concept is that of gauge-fixed BRST transformation, which is what $`s`$ becomes after gauge-fixing. It is denoted by $`s_\psi `$ and defined as
$$s_\psi \mathrm{\Phi }^A=(s\mathrm{\Phi }^A)|_{\mathrm{\Phi }_A^{}=\delta \psi /\delta \mathrm{\Phi }^A}.$$
(7.5)
Note that $`s_\psi ^2=0`$ is in general only valid on (gauge-fixed) shell, i.e. for field configurations satisfying $`\delta S_\psi /\delta \mathrm{\Phi }^A=0`$.
If $`S`$ is linear in the antifields, i.e. if the gauge algebra closes off-shell, $`s_\psi \mathrm{\Phi }^A=s\mathrm{\Phi }^A`$. This follows directly from
$$s\mathrm{\Phi }^A=(S,\mathrm{\Phi }^A)=\frac{\delta ^RS}{\delta \mathrm{\Phi }_A^{}},$$
(7.6)
which is independent of the antifields, if $`S`$ depends only linearly on the antifields. In that case $`s_\psi ^2=0`$ even off-shell provided one keeps all the variables. This is the case in Yang-Mills theory. \[In Yang-Mills theory, one often eliminates the auxiliary fields $`b_a`$ by means of their own equations of motion. One then loses off-shell nilpotency on the antighosts, for which $`s\overline{c}_ab_a`$, even though $`s_\psi ^2=0`$ is true off-shell beforehand.\]
If $`S`$ is linear in the antifields, one may in fact write it as
$$S=S_0(s\mathrm{\Phi }^A)\mathrm{\Phi }_A^{},$$
(7.7)
by virtue of (7.6). The gauge fixed version is then
$$S_\psi =S_0(s_\psi \mathrm{\Phi }^A)\frac{\delta \psi }{\delta \mathrm{\Phi }^A}=S_0s_\psi \psi .$$
(7.8)
Therefore $`s_\psi S_\psi =0`$, as the first term is BRST invariant, and the second term is annihilated by $`s_\psi `$ by virtue of $`s_\psi ^2=0`$, which holds off-shell when $`S`$ is linear in the antifields as we have just pointed out. The property $`s_\psi S_\psi =0`$ is actually quite general and holds even when $`S`$ is not linear in the antifields. It can be proved directly as follows,
$$s_\psi S_\psi =(s_\psi \mathrm{\Phi }^A)\frac{\delta ^LS_\psi }{\delta \mathrm{\Phi }^A}=\frac{\delta ^RS}{\delta \mathrm{\Phi }_A^{}}\frac{\delta ^LS_\psi }{\delta \mathrm{\Phi }^A}.$$
(7.9)
The left-derivative in this expression is a total derivative, as $`S_\psi `$ depends on $`\mathrm{\Phi }_A`$ directly and though the gauge-fixing fermion. Using the chain rule, this yields by virtue of (7.3)
$$\frac{\delta ^RS}{\delta \mathrm{\Phi }_A^{}}\left(\frac{\delta S}{\delta \mathrm{\Phi }^A}+\frac{\delta ^2\psi }{\delta \mathrm{\Phi }^A\delta \mathrm{\Phi }^B}\frac{\delta S}{\delta \mathrm{\Phi }_B^{}}\right)=0.$$
(7.10)
The second term vanishes because the product of the functional derivatives of $`S`$ have a symmetry in ($`A`$, $`B`$) opposite to that of the second functional derivative of $`\psi `$. The first term vanishes by the master equation, thus proving the claim.
Further information on the gauge-fixed action and the gauge-fixed cohomology can be found in .
## 8 Quantum Master Equation
In order to prove gauge independence of the expressions given above, it is necessary to discuss two important features of the path integral.
* Assume that a theory of fields $`\chi ^\alpha `$ is given, governed by the action $`S[\chi ^\alpha ]`$, with no gauge invariance (this could be the gauge-fixed action). Expectation values are, after proper normalization, calculated as
$$<F>=𝒟\chi F\mathrm{exp}\left(\frac{i}{\mathrm{}}S[\chi ]\right).$$
(8.1)
The Dyson-Schwinger equations can be directly derived from the vanishing of the path integral of a total derivative,
$$𝒟\chi \frac{\delta }{\delta \chi ^\alpha }\left(Fe^{\frac{i}{\mathrm{}}S}\right)=0$$
(8.2)
(which is itself a consequence of translation invariance of the measure). This leads to
$$\frac{\delta F}{\delta \chi ^\alpha }+\frac{i}{\mathrm{}}F\frac{\delta S}{\delta \chi ^\alpha }=0,$$
(8.3)
which is equivalent to
$$F\frac{\delta S}{\delta \chi ^\alpha }=i\mathrm{}\frac{\delta F}{\delta \chi ^\alpha }.$$
(8.4)
This expression contains in the l.h.s. the expectation values of the classical equations of motions. In the classical limit $`\mathrm{}0`$, the r.h.s. vanishes and the classical equations of motion hold.
* There is another aspect that we shall have to take into account. If $`\chi ^\alpha `$ is changed under a transformation, $`\chi ^\alpha \chi ^\alpha +ϵ^\alpha `$, where $`ϵ^\alpha `$ depends on the fields, the expectation value (8.1) is in general not invariant. Furthermore, invariance of the classical action is not sufficient to guarantee that the path integral is invariant. One needs also invariance of the measure.
We shall be concerned with BRST invariance of the path integral constructed above. We have seen that the gauge-fixed action is BRST invariant. But the measure might not be. If it is not, one may, in some cases, restore invariance by taking a different measure. \[The measure is in fact dictated by unitarity and may indeed not be equal to the trivial measure $`𝒟\mathrm{\Phi }`$.\] Invariance is quite crucial in the case of BRST symmetry, since it is BRST symmetry that guarantees gauge-independence of the results. The non-trivial measure terms can be exponentiated in the action. Since there is an overall $`(1/\mathrm{})`$ in front of $`S`$, the measure terms appear as quantum corrections to $`S`$. So, one replaces the classical action by a “quantum action”
$$W=S+\mathrm{}M_1+\mathrm{}^2M_2+\mathrm{},$$
(8.5)
where the functionals $`M_i`$ stem from non-trivial measure factors. The theorem proved below states that quantum averages are gauge-independent if the master equation is replaced by the “quantum master equation”
$$\frac{1}{2}(W,W)=i\mathrm{}\mathrm{\Delta }W,$$
(8.6)
where $`\mathrm{\Delta }`$ is defined in (6.13). Note that if $`\mathrm{\Delta }S=0`$, the Jacobian is unity for the BRST transformation and $`W`$ might then be taken equal to $`S`$. The quantum master equation reduces to the classical master equation considered above, which is solved by $`S`$. While there is always a solution to the classical master equation, the solution to the quantum master equation might get obstructed. We shall investigate this question below. For the moment, we assume that there is no obstruction.
We can now state the correct, general rules, for computing expectation values of observables (including $`1`$): these are the quantum averages, weighted by $`\mathrm{exp}(\frac{i}{\mathrm{}}W)`$, of the BRST observables corrected by the addition of appropriate $`\mathrm{}`$ (and possibly also higher) order terms. Namely, consider a classical observable $`A_0`$. Construct its (in fact, one of its) BRST-invariant extension $`A=A_0+\text{ghost terms}`$, so that $`(S,A)=0`$. The BRST cocycle $`A`$ has to be augmented as
$$A\alpha =A_0+\mathrm{}B_1+\mathrm{}^2B_2+\mathrm{}$$
(8.7)
where the terms of order $`\mathrm{}`$ and higher must be such that $`\sigma \alpha =0`$, where $`\sigma `$ was defined in the exercises as
$$\sigma \alpha (W,\alpha )i\mathrm{}\mathrm{\Delta }\alpha ,$$
(8.8)
with $`W`$ the solution of the quantum master equation. (Note that these $`B`$-terms come over and above the ghost terms needed classically to fulfill $`(S,A)=0`$.) The operator $`\sigma `$ is the quantum generalization of $`s`$. The $`\psi `$-independent expectation value $`<A_0>`$ of the observable $`A_0[\varphi ^i]`$ is computed from $`\alpha `$ as
$$<A_0>=𝒟\mathrm{\Phi }^A\alpha \left(\mathrm{\Phi }^A,\mathrm{\Phi }_A^{}=\frac{\delta \psi }{\delta \mathrm{\Phi }^A}\right)\mathrm{exp}\left(\frac{i}{\mathrm{}}W\left[\mathrm{\Phi }^A,\mathrm{\Phi }_A^{}=\frac{\delta \psi }{\delta \mathrm{\Phi }^A}\right]\right).$$
(8.9)
The claim is that this expectation value does not depend on the choice of the gauge-fixing fermion
$`\psi ^{}=\psi +\delta \psi \psi +\mu `$ (8.10)
$`<A_0>_\psi ^{}=<A_0>_\psi ,`$ (8.11)
where $`\mu `$ is an arbitrary modification of $`\psi `$.
To prove the claim, we denote the argument of the integral by $`V`$ for convenience in the following. It has been already shown in the exercises that
$$\mathrm{\Delta }V=0\sigma \alpha =0.$$
(8.12)
Now, perform the variation of the gauge fixing functional (8.10). The variation of the quantum average (8.9) is equal to
$$𝒟\mathrm{\Phi }\frac{\delta ^L\mu }{\delta \mathrm{\Phi }^A}\frac{\delta ^LV}{\delta \mathrm{\Phi }_A^{}}.$$
(8.13)
To evaluate this expression, note that
$$\frac{\delta ^L}{\delta \mathrm{\Phi }^A}\left(\mu \frac{\delta ^LV}{\delta \mathrm{\Phi }_A^{}}\right)=\frac{\delta ^L\mu }{\delta \mathrm{\Phi }^A}\frac{\delta ^LV}{\delta \mathrm{\Phi }_A^{}}+(1)^{ϵ_A}\frac{\delta ^L\delta ^LV}{\delta \mathrm{\Phi }^A\delta \mathrm{\Phi }_A^{}}.$$
(8.14)
The derivatives are total ones. Denoting partial derivatives by a prime ’, the last term can be rewritten as
$$\frac{\delta ^L\delta ^LV}{\delta \mathrm{\Phi }^A\delta \mathrm{\Phi }_A^{}}=\frac{\delta ^{}_{}{}^{}L\delta ^LV}{\delta ^{}\mathrm{\Phi }^A\delta \mathrm{\Phi }_A^{}}+\frac{\delta ^2\psi }{\delta \mathrm{\Phi }^A\delta \mathrm{\Phi }^B}\frac{\delta ^L\delta ^LV}{\delta \mathrm{\Phi }_B^{}\delta \mathrm{\Phi }_A^{}}=\mathrm{\Delta }V=0.$$
(8.15)
As in the case of equation (7.10), the second term vanishes by parity arguments. Thus the integral (8.13) can be rewritten as a total derivative in field space, which vanishes in view of translation invariance of the standard measure. Therefore the path integral does not get modified if one changes the gauge-fixing fermion, as claimed.
Given $`A_0`$, its BRST extension is determined up to a BRST exact term $`sB`$, see (5.8). This ambiguity can be extended to higher orders in $`\mathrm{}`$ as
$$\alpha \alpha +\sigma \beta .$$
(8.16)
It has been shown in the exercises that
$$(\sigma \beta )\mathrm{exp}\left(\frac{i}{\mathrm{}}W\right)\mathrm{\Delta }\left(\beta \mathrm{exp}\left(\frac{i}{\mathrm{}}W\right)\right).$$
(8.17)
As the r.h.s is a total derivative by its definition, (6.13), the path-integral over the l.h.s. vanishes
$$𝒟\mathrm{\Phi }\sigma \beta \mathrm{exp}\left(\frac{i}{\mathrm{}}W\right)=0.$$
(8.18)
Therefore adding any element in the image of $`\sigma `$ does not alter the quantum averages. The path integral associates a unique answer to a given cohomological class of $`\sigma `$, i.e., does not depend on the choice of representative.
Note that the ambiguity in $`\alpha `$, given $`A_0`$, is more than just adding a $`\sigma `$-trivial term to $`\alpha `$. At each order in $`\mathrm{}`$ one may add a non trivial new observable since this does not modify the classical limit. This addition is relevant, in the sense that it changes the expectation value by terms of order $`\mathrm{}`$ or higher. This is an unavoidable quantum ambiguity. A similar ambiguity exists for the quantum measure (i.e., the $`M_1`$, $`M_2`$ etc. terms in $`W`$). These terms do not spoil BRST invariance and must be determined by other criteria, e.g., by comparison with the Hamiltonian formalism.
## 9 Anomalies
We close this brief survey by analyzing the possible obstructions to the existence of a solution $`W`$ to the quantum master equation. This leads to the important concept of anomalies. The fact that anomalies in the Batalin-Vilkovisky formalism appear as an incurable violation of the BRST invariance of the measure was first investigated in .
Analyzing the obstructions to the existence of a solution to the quantum master equation can be done most easily from a direct $`\mathrm{}`$ expansion. To order $`\mathrm{}^0`$ one gets from the quantum master equation
$$\frac{1}{2}(S,S)=0,$$
(9.1)
which is the classical master equation. This equation is certainly fulfilled, since there is no obstruction to the existence of $`S`$.
To the next order $`\mathrm{}`$, the quantum master equation yields
$$sM_1=(S,M_1)=i\mathrm{\Delta }S.$$
(9.2)
Given the $`\mathrm{}^0`$-term $`S`$, this equation has a solution for $`M_1`$ if $`s\mathrm{\Delta }S=(S,\mathrm{\Delta }S)=0`$. This condition is necessary but in general not sufficient (see below). To prove that the condition $`(S,\mathrm{\Delta }S)=0`$ holds, we note that
$$\mathrm{\Delta }(\alpha ,\beta )=(\mathrm{\Delta }\alpha ,\beta )(1)^{ϵ_\alpha }(\alpha ,\mathrm{\Delta }\beta ),$$
(9.3)
as was proven in the exercises. This property uses $`\mathrm{\Delta }^2=0`$ and the generalized Leibniz rule
$$\mathrm{\Delta }(\alpha \beta )=(\mathrm{\Delta }\alpha )\beta +(1)^{ϵ_\alpha }\alpha \mathrm{\Delta }\beta +(1)^{ϵ_\alpha }(\alpha ,\beta ).$$
(9.4)
For $`\alpha =\beta =S`$, the l.h.s of (9.3) vanishes by $`(S,S)=0`$. In view of the gradings of $`S`$, the r.h.s yields $`2(\mathrm{\Delta }S,S)`$. Thus, $`(\mathrm{\Delta }S,S)=0`$, that is, $`\mathrm{\Delta }S`$ is closed, as requested.
This does not imply that $`\mathrm{\Delta }S`$ is exact, however, unless the cohomological group $`H^1(s)`$ vanishes at ghost number one (recall that $`\mathrm{\Delta }S`$ has ghost number one). But (9.2) requires $`\mathrm{\Delta }S`$ to be exact. If $`\mathrm{\Delta }S`$ is not exact, there is no $`M_1`$ and therefore, no way to define a BRST invariant measure such that the quantum averages do not depend on the gauge fixing fermion. This presumably signals a serious pathology of the theory. If $`\mathrm{\Delta }S`$ is exact, $`M_1`$ exists and one can investigate the problem of existence of the next term $`M_2`$. One easily verifies that it is again $`H^1(s)`$ that measures the potential obstructions to the existence of this next term $`M_2`$, as well as the existence of the subsequent terms $`M_3`$ etc.
In particular, if one can show that $`H^1(s)`$ vanishes, one is guaranteed that a solution of the quantum master equation exists. If $`H^1(s)0`$, further work is required since one must check that one does not hit an obstruction. Note that $`W`$ should be a local functional (with possibly infinite coupling constants), so that the relevant space in which to compute the cohomology is that of local functionals.
The computation of the local cohomology of the BRST operator for Yang-Mills gauge theory has been carried out in , following earlier work without antifields . See for a review.
## Acknowledgements
We are grateful to the organizers of the $`10^{th}`$ Saalburg Summer School on Modern Theoretical Methods for their kind invitation and hospitality. We are also grateful to Jim Stasheff for a careful reading of the manuscript. The work of MH is supported in part by the “Interuniversity Attraction Poles Programme – Belgian Science Policy ”, by IISN-Belgium (convention 4.4505.86) and by the European Commission FP6 programme MRTN-CT-2004-005104, in which he is associated to the V.U.Brussel (Belgium).
## Appendix A Appendix I: De Witt’s notation.
We review the De Witt’s condensed notation. This notation makes it possible to write gauge transformations in a more compact form:
$$\delta _\epsilon \phi ^i=R_\alpha ^i\epsilon ^\alpha \delta _\epsilon \phi ^i(x)=d^nyR_\alpha ^i(y,x)\epsilon ^\alpha (y).$$
(A.1)
For example, the transformation of the Yang-Mills gauge field
$$\delta _\epsilon A_\mu ^a=D_\mu \epsilon ^a=_\mu \epsilon ^a+f_{cb}^aA_\mu ^b\epsilon ^c$$
(A.2)
can be written as
$$\delta _\epsilon A_\mu ^a=R_{\mu b}^a\epsilon ^b,R_{\mu b}^a(x,y)=_\mu \delta (xy)\delta _b^a+f_{bc}^aA_\mu ^b\delta (xy).$$
(A.3)
Noether’s identities have the simple form:
$$\frac{\delta S_0}{\delta \phi ^i}R_\alpha ^i=0.$$
(A.4)
## Appendix B Appendix II: Properties of anti-brackets.
(i) $`(F,G)=(1)^{(\epsilon _F+1)(\epsilon _G+1)}(G,F)`$, where $`\epsilon _F=0(1)`$ for $`F`$ bosonic (fermionic); the anti-bracket is symmetric if both $`F`$ and $`G`$ are bosonic, and antisymmetric otherwise.
(ii) Jacobi identitiy:
$`(1)^{(\epsilon _F+1)(\epsilon _H+1)}(F,(G,H))+\text{cyclic permutations}=0`$
(iii) $`(FG,H)=F(G,H)+(F,H)G(1)^{\epsilon _G(\epsilon _H+1)}`$;
$`(F,GH)=(F,G)H+G(F,H)(1)^{\epsilon _G(\epsilon _F+1)}`$
(iv) $`\mathrm{gh}\left((F,G)\right)=\mathrm{gh}F+\mathrm{gh}G+1`$ |
warning/0506/math0506390.html | ar5iv | text | # Knots and words
## 1. Introduction
C. F. Gauss \[Ga\] introduced a method allowing to encode closed planar curves by words of a certain type called now Gauss words. This method extends to planar knot diagrams and their isotopy. This gives a description of the set of isotopy types of classical knots in terms of words and their transformations.
To state our results we need to generalize both knots and words. An appropriate generalization of knots is provided by long virtual knots, see \[Kau2\], \[GPV\]. We shall use an equivalent formulation in terms of stable equivalence classes of (pointed oriented) knot diagrams on surfaces, see \[KK\], \[CKS\]. The set of these equivalence classes $`𝒦`$ contains the set $`𝒮`$ of isotopy classes of oriented knots in $`S^3`$. On the combinatorial side, the Gauss words generalize to so-called nanowords, see \[Tu2\]. Our main result is a bijection between $`𝒦`$ and the set of (appropriately defined) homotopy classes of nanowords in the alphabet consisting of 4 letters. Note that the image of $`𝒮𝒦`$ under this bijection can be described via the Rosenstiehl theorem \[Ro\] giving necessary and sufficient conditions for the curve corresponding to a given Gauss word to be planar. Thus, isotopy classification of knots in $`S^3`$ is a special instance of homotopy classification of words. Similarly, classical and virtual links can be interpreted as nanophrases. This gives a broader perspective to knot theory. A number of methods of knot theory including the Kauffman bracket polynomial, the Jones polynomial, the knot quandle etc. can be extended to the more general setting of words and phrases in arbitrary alphabets.
The theory of knots and links being very reach, it is natural to study related simpler objects. One simplification of links in cylinders over surfaces is obtained by projecting them to the surfaces, i.e., by forgetting the over/under-crossing information in link diagrams. The stable equivalence classes of curves on surfaces were studied in \[Kad\], \[Tu1\]. We describe these classes in terms of nanowords and nanophrases over an alphabet consisting of only two letters. We also suggest two different simplifications of knot theory in terms of nanophrases over a 2-letter alphabet. The resulting combinatorial objects are called pseudo-links and quasi-links. Every oriented link in a cylinder over an oriented surface projects to a (multi-component) curve on the surface, to a pseudo-link and to a quasi-link. Their distinctive features are contained in the following facts: the Jones polynomial of a link depends only on the underlying pseudo-link; the fundamental group of the 2-fold branched covering of a classical link depends only on the underlying quasi-link.
The plan of the paper is as follows. In Sect. 2 we recall nanowords and their homotopy. In Sect. 3 we identify the stable equivalence classes of pointed curves on surfaces with homotopy classes of nanowords in a 2-letter alphabet. In Sect. 4 we identify the stable equivalence classes of pointed knot diagrams on surfaces with homotopy classes of nanowords in a 4-letter alphabet. In Sect. 5 we show how to get rid of the base points. In Sect. 6 we extend these results to links. In Sect. 7 we introduce pseudo-links and quasi-links. In Sect. 8 we discuss the bracket polynomial and the Jones polynomial. In Sect. 9 and 10 we discuss the keis of nanophrases.
Conventions. Throughout the paper, all surfaces, curves, knots, links, and knot and link diagrams are oriented, unless explicitly stated to the contrary.
## 2. Nanowords and homotopy
### 2.1. Words and nanowords
An alphabet is a set and letters are its elements. A word of length $`n1`$ in an alphabet $`𝒜`$ is a mapping $`w:\widehat{n}𝒜`$ where $`\widehat{n}=\{1,2,\mathrm{},n\}`$. A word $`w:\widehat{n}𝒜`$ is usually encoded by the sequence of letters $`w(1)w(2)\mathrm{}w(n)`$. A word $`w:\widehat{n}𝒜`$ is a Gauss word if each element of $`𝒜`$ is the image of precisely two elements of $`\widehat{n}`$.
For a set $`\alpha `$, an $`\alpha `$-alphabet is a set $`𝒜`$ endowed with a mapping $`𝒜\alpha `$ called projection. The image of $`A𝒜`$ under this mapping is denoted $`|A|`$. A nanoword over $`\alpha `$ is a pair (an $`\alpha `$-alphabet $`𝒜`$, a Gauss word in the alphabet $`𝒜`$). For example, any Gauss word $`w`$ in the alphabet $`\alpha `$ yields a nanoword $`(𝒜=\alpha ,w)`$, see \[Tu2\] for more examples and further details. By definition, there is a unique empty nanoword $`\mathrm{}`$ of length 0.
An isomorphism of $`\alpha `$-alphabets $`𝒜_1`$, $`𝒜_2`$ is a bijection $`f:𝒜_1𝒜_2`$ such that $`|A|=|f(A)|`$ for all $`A𝒜_1`$. Two nanowords ($`𝒜_1`$, $`w_1`$) and ($`𝒜_2`$, $`w_2`$) over $`\alpha `$ are isomorphic if there is an isomorphism of $`\alpha `$-alphabets $`f:𝒜_1𝒜_2`$ such that $`w_2=fw_1`$.
### 2.2. Homotopy of nanowords
A homotopy data consists of a set $`\alpha `$ with involution $`\tau :\alpha \alpha `$ and a set $`S\alpha ^3=\alpha \times \alpha \times \alpha `$. The following three transformations of nanowords over $`\alpha `$ are called $`S`$-homotopy moves or simply homotopy moves.
(1). The first move applies to any nanoword of the form $`(𝒜,xAAy)`$ where $`A𝒜`$ and $`x,y`$ are words in the alphabet $`𝒜^{}=𝒜\{A\}`$. It transforms $`(𝒜,xAAy)`$ into the nanoword $`(𝒜^{},xy)`$ where the structure of an $`\alpha `$-alphabet in $`𝒜^{}`$ is obtained by restricting the one in $`𝒜`$. Note that $`xy`$ is a Gauss word in the alphabet $`𝒜^{}`$.
The inverse move $`(𝒜^{},xy)(𝒜^{}\{A\},xAAy)`$ adds a new letter $`A`$ with arbitrary $`|A|\alpha `$ and replaces the Gauss word $`xy`$ in the $`\alpha `$-alphabet $`𝒜^{}`$ with $`xAAy`$.
(2). The second move applies to a nanoword of the form $`(𝒜,xAByBAz)`$ where $`A,B𝒜`$ with $`|B|=\tau (|A|)`$ and $`x,y,z`$ are words in the alphabet $`𝒜^{}=𝒜\{A,B\}`$. This nanoword is transformed into $`(𝒜^{},xyz)`$ where the structure of an $`\alpha `$-alphabet in $`𝒜^{}`$ is obtained by restricting the one in $`𝒜`$.
The inverse move $`(𝒜^{},xyz)(𝒜^{}\{A,B\},xAByBAz)`$ adds two new letters $`A,B`$ with arbitrary $`|A|\alpha `$ and $`|B|=\tau (|A|)`$ and replaces the Gauss word $`xyz`$ in the alphabet $`𝒜^{}`$ with $`xAByBAz`$.
(3) The third move applies to a nanoword of the form $`(𝒜,xAByACzBCt)`$ where $`A,B,C𝒜`$ are distinct letters such that $`(|A|,|B|,|C|)S`$ and $`x,y,z,t`$ are words in the alphabet $`𝒜\{A,B,C\}`$. The move transforms $`(𝒜,xAByACzBCt)`$ into $`(𝒜,xBAyCAzCBt)`$. The inverse move applies if $`(|A|,|B|,|C|)S`$ and transforms a nanoword $`(𝒜,xBAyCAzCBt)`$ into $`(𝒜,xAByACzBCt)`$.
Two nanowords over $`\alpha `$ are $`S`$-homotopic if they can be obtained from each other by a finite sequence of homotopy moves (1) – (3), the inverse moves, and isomorphisms. The relation of $`S`$-homotopy is denoted $`_S`$. The set of $`S`$-homotopy classes of nanowords over $`\alpha `$ is denoted $`𝒩(\alpha ,S)`$. (The involution $`\tau `$ is omitted in this notation for shortness.)
Recall the following two lemmas from \[Tu2\], Sect. 3.2.
###### Lemma 2.2.1.
Let $`A,B,C`$ be distinct letters in an $`\alpha `$-alphabet $`𝒜`$ and let $`x,y,z,t`$ be words in the alphabet $`𝒜\{A,B,C\}`$ such that $`xyzt`$ is a Gauss word in this alphabet. Then
(i) $`(𝒜,xAByCAzBCt)_S(𝒜,xBAyACzCBt)`$ for $`(|A|,\tau (|B|),|C|)S`$;
(ii) $`(𝒜,xAByCAzCBt)_S(𝒜,xBAyACzBCt)`$ for $`(\tau (|A|),\tau (|B|),|C|)S`$;
(iii) $`(𝒜,xAByACzCBt)_S(𝒜,xBAyCAzBCt)`$ for $`(\tau (|A|),|B|,|C|)S`$.
A homotopy data $`(\alpha ,S)`$ is admissible if $`S(\alpha \times a\times a)\mathrm{}`$ for all $`a\alpha `$. For instance, if $`S`$ contains the diagonal $`\{(a,a,a)\}_{a\alpha }`$ of $`\alpha ^3`$, then $`(\alpha ,S)`$ is admissible.
###### Lemma 2.2.2.
Let $`(𝒜,xAByABz)`$ be a nanoword over $`\alpha `$ where $`A,B𝒜`$ with $`|B|=\tau (|A|)`$ and $`x,y,z`$ are words in the alphabet $`𝒜\{A,B\}`$. If $`(\alpha ,S)`$ is admissible, then $`(𝒜,xAByABz)_S(𝒜\{A,B\},xyz)`$.
A morphism $`(\alpha ,S)(\alpha ^{},S^{})`$ between two homotopy data is an equivariant mapping $`f:\alpha \alpha ^{}`$ such that $`(f\times f\times f)(S)S^{}`$. Given $`f`$, we can transform a nanoword $`(𝒜,w)`$ over $`\alpha `$ into a nanoword $`(𝒜^{},w^{})`$ over $`\alpha ^{}`$ where $`𝒜^{}=𝒜`$ as sets, $`w^{}=w`$, and the projection $`𝒜^{}\alpha ^{}`$ is the composition of the projection $`𝒜^{}=𝒜\alpha `$ with $`f`$. This transformation is compatible with homotopy and induces a monoid homomorphism $`𝒩(\alpha ,S)𝒩(\alpha ^{},S^{})`$.
## 3. Curves versus words
### 3.1. Curves
By a curve, we mean the image of a generic immersion of an oriented circle into an oriented surface. The word “generic” means that the curve has only a finite set of self-intersections which are all double and transversal. A curve is pointed if it is endowed with a base point (the origin) which is not a self-intersection. Two pointed curves are stably homeomorphic if there is a homeomorphism of their regular neighborhoods in the ambient surfaces mapping the first curve onto the second one and preserving the origin of the curve and the orientations of the curve and the surface. In particular, attaching a 1-handle to the ambient surface away from a curve or removing such a handle does not change the stable homeomorphism type of the curve.
Following \[KK\], \[CKS\], we call two pointed curves stably equivalent if they can be related by a finite sequence of the following transformations: (i) replacing the curve with a stably homeomorphic one; (ii) homotopy of the curve in its ambient surface away from the origin. Note that such a homotopy may push a branch of the curve across another branch or a double point but not across the origin of the curve.
Denote $`𝒞`$ the set of stable equivalence classes of pointed curves. This set is a monoid with multiplication defined by connected sum at the origin. We are far from understanding the algebraic structure of $`𝒞`$. Several stable equivalence invariants of pointed curves were introduced in \[Tu1\], \[SW\] where curves are studied in terms of virtual strings.
We show now that the study of $`𝒞`$ is an instance of homotopy theory of words.
### 3.2. Homotopy data $`(\alpha _0,S_0)`$
Consider the homotopy data $`(\alpha _0,S_0)`$ where $`\alpha _0`$ is the set $`\{a,b\}`$ with involution $`\tau :\alpha _0\alpha _0`$ permuting $`a,b`$ and $`S_0=\{(a,a,a)\}_{a\alpha _0}`$ is the diagonal. This homotopy data is admissible in the sense of Sect. 2.2.
###### Theorem 3.2.1.
There is a canonical bijection $`𝒞=𝒩(\alpha _0,S_0)`$.
###### Proof.
We associate with any pointed curve $`f`$ a nanoword $`w(f)`$ over $`\alpha _0`$. Let us label the double points of $`f`$ by (distinct) letters $`A_1,\mathrm{},A_n`$ where $`n`$ is the number of double points. Starting at the origin of $`f`$ and following along $`f`$ in the positive direction we write down the labels of all double points until the return to the origin. Since every double point is traversed twice, this gives a Gauss word $`w(f)`$ in the alphabet $`𝒜=\{A_1,\mathrm{},A_n\}`$. Let $`t_i^1`$ (resp. $`t_i^2`$) be the tangent vector to $`f`$ at the crossing point labeled by $`A_i`$ appearing at the first (resp. second) passage through this crossing. Set $`|A_i|=a`$ if the pair $`(t_i^1,t_i^2)`$ is positively oriented and $`|A_i|=b`$ otherwise. This makes $`𝒜`$ into an $`\alpha _0`$-alphabet and makes $`w`$ into a nanoword over $`\alpha _0`$. This nanoword is well defined up to isomorphism.
We claim that stably equivalent pointed curves give rise to $`S_0`$-homotopic nanowords. Stable homeomorphisms of curves preserve the nanoword up to isomorphism. We need to show that a homotopy of a curve $`f`$ in its ambient surface away from the origin does not change the $`S_0`$-homotopy class of $`w(f)`$. Such a homotopy can be obtained by an ambient isotopy and a finite sequence of local deformations shown in Figure 1 and the inverse deformations. It is understood that all deformations in Figure 1 are effected away from the origin of $`f`$.
An ambient isotopy does not change $`w(f)`$. A local deformation of the first type changes $`w(f)`$ via the first homotopy move. Depending on the orientations of the two branches of $`f`$, a local deformation of the second type changes $`w(f)`$ via one of the moves $`xAByBAzxyz`$ or $`xAByABzxyz`$ where $`|B|=\tau (|A|)`$. In both cases the $`S_0`$-homotopy class of $`w(f)`$ is preserved. Consider the local deformation of the third type. It suffices to consider the case where all three branches are oriented upwards. (The deformations involving other orientations of the branches can be obtained as compositions of this one with ambient isotopy and local deformations of the second type.) There are 6 cases to consider depending on the order in which one traverses the three branches involved. Let $`I`$ (resp. $`II,III`$) be the branch connecting the leftmost (resp. intermediate, rightmost) bottom point to the rightmost (resp. intermediate, leftmost) top point. If one traverses these branches in the order $`I,II,III`$ (resp. $`III,II,I`$), then $`|A|=|B|=|C|=a`$ (resp. $`|A|=|B|=|C|=b`$) and the deformation changes $`w(f)`$ via the third homotopy move (resp. its inverse). If one traverses these branches in the order $`II,III,I`$ (resp. $`I,III,II`$), then the deformation changes $`w(f)`$ via the second homotopy from Lemma 2.2.1 (resp. its inverse). Finally, if one traverses these three branches in the order $`II,I,III`$ (resp. $`III,I,II`$), then the deformation changes $`w(f)`$ via the third homotopy from Lemma 2.2.1 (resp. its inverse).
We can conclude that the formula $`fw(f)`$ defines a mapping $`W:𝒞𝒩(\alpha _0,S_0)`$. We claim that it is bijective. It is easy to see that for every nanoword $`w`$ over $`\alpha _0`$ there is a unique (up to stable homeomorphism) pointed curve $`f`$ such that $`w(f)=w`$. Indeed, knowing $`w`$ we can uniquely recover an oriented regular neighborhood of such a curve in the ambient surface (this well known construction is described in detail in \[Tu1\], Sect. 4.1 in terms of virtual strings. Note that the notion of an open virtual string is equivalent to the one of an isomorphism class of a nanoword over $`\alpha _0`$.) This implies that $`W`$ is surjective.
The injectivity of $`W`$ follows from the fact that if two nanowords are related by the $`i`$-th homotopy move with $`i=1,2,3`$ then they can be represented by pointed curves related by the $`i`$-th deformation in Figure 1 (effected away from the origin). Here for $`i=2`$ the two branches in Figure 1 are oriented in opposite directions and for $`i=3`$ all the branches are oriented upwards and traversed in the order $`I,II,III`$. ∎
## 4. Knots versus words
### 4.1. Knot diagrams
By a knot diagram, we mean a (generic oriented) curve on an (oriented) surface such that at each crossing point of the curve one of the two branches is distinguished. The distinguished branch is the over-crossing and the second branch is the under-crossing. A knot diagram is pointed if it is endowed with a base point (the origin) distinct from the crossing points. Two pointed knot diagrams are stably homeomorphic if there is a homeomorphism of their regular neighborhoods in the ambient surfaces mapping the first diagram onto the second one and preserving the origin, the over/undercrossings, and the orientations of the surface and the curve.
Following \[KK\], \[CKS\], we call two pointed knot diagrams stably equivalent if they can be related by a finite sequence of the following transformations: (i) replacing a knot diagram with a stably homeomorphic one; (ii) the usual Reidemeister moves on a knot diagram in its ambient surface away from the origin. The latter moves may push a branch of the diagram above or below a double point or another branch but not across the origin. It should be stressed that removing a closed subset from the ambient surface away from a knot diagram or attaching a 1-handle away from the diagram does not change the stable equivalence type of the diagram.
Denote $`𝒦`$ the set of stable equivalence classes of pointed knot diagrams. The elements of $`𝒦`$ bijectively correspond to long virtual knots in the sense of \[Kau2\], \[GPV\]. Every (oriented) knot $`KS^3`$ determines an element of $`𝒦`$ obtained by presenting $`K`$ by a diagram on $`S^2`$ and picking an arbitrary base point. This yields a well defined mapping from the set of isotopy classes of classical knots into $`𝒦`$. This mapping is essentially injective, see \[Kau2\], \[GPV\].
Forgetting the over/under-crossing information we obtain a natural projection $`𝒦𝒞`$. We now interpret $`𝒦`$ in terms of words.
### 4.2. Homotopy data $`(\alpha _{},S_{})`$
Consider the homotopy data $`(\alpha _{},S_{})`$ where $`\alpha _{}=\{a_+,a_{},b_+,b_{}\}`$ with involution $`\tau :\alpha _{}\alpha _{}`$ defined by $`\tau (a_\pm )=b_{},\tau (b_\pm )=a_{}`$ and $`S_{}\alpha _{}\times \alpha _{}\times \alpha _{}`$ consists of the following 12 triples:
$$(a_\pm ,a_\pm ,a_\pm ),(a_\pm ,a_\pm ,a_{}),(a_{},a_\pm ,a_\pm ),(b_\pm ,b_\pm ,b_\pm ),(b_\pm ,b_\pm ,b_{}),(b_{},b_\pm ,b_\pm ).$$
This homotopy data is admissible in the sense of Sect. 2.2.
Forgetting the signs, we obtain a projection $`\alpha _{}\alpha _0`$. Applying it, we can transform a nanoword over $`\alpha _{}`$ into a nanoword over $`\alpha _0`$. This induces a monoid homomorphism $`𝒩(\alpha _{},S_{})𝒩(\alpha _0,S_0)`$.
###### Theorem 4.2.1.
There is a canonical bijection $`𝒦=𝒩(\alpha _{},S_{})`$. Under this bijection, the monoid homomorphism $`𝒩(\alpha _{},S_{})𝒩(\alpha _0,S_0)`$ corresponds to the natural projection $`𝒦𝒞`$.
###### Proof.
The proof reproduces the proof of Theorem 3.2.1 with a few changes. We begin by associating with any pointed knot diagram $`F`$ a nanoword $`w=w(F)`$ over $`\alpha _{}`$. As usual, each crossing of $`F`$ gives rise to a sign $`\pm `$. It is $`+`$ if the over-going branch crosses the under-going branch from left to right and $``$ otherwise. To define $`w`$, label the double points of $`F`$ by (distinct) letters $`A_1,\mathrm{},A_n`$ where $`n`$ is the number of double points. Starting at the origin of $`F`$ and following along $`F`$ we write down the labels of all double points until the return to the origin. This gives a Gauss word $`w`$ in the alphabet $`𝒜=\{A_1,\mathrm{},A_n\}`$. Let $`t_i^1`$ (resp. $`t_i^2`$) be the tangent vector to $`F`$ at the crossing labeled $`A_i`$ appearing at the first (resp. second) passage through this crossing. Let $`\epsilon (i)=\pm `$ be the sign of this crossing. Set $`|A_i|=a_{\epsilon (i)}`$ if the pair $`(t_i^1,t_i^2)`$ is positively oriented and $`|A_i|=b_{\epsilon (i)}`$ otherwise. This makes $`𝒜`$ into an $`\alpha _{}`$-alphabet and makes $`w=w(F)`$ into a nanoword over $`\alpha _{}`$. This nanoword is well defined up to isomorphism.
We need to verify that stably equivalent pointed knot diagrams give rise to $`S_{}`$-homotopic nanowords. Stable homeomorphisms preserve the nanoword up to isomorphism. We need to show that the $`S_{}`$-homotopy class of $`w(F)`$ is preserved under the Reidemeister moves on $`F`$ away from the origin. The first Reidemeister move changes $`w(F)`$ via the first homotopy move. Depending on the orientations of the two branches of $`F`$ involved in the second Reidemeister move, the nanoword $`w(F)`$ changes via one of the moves $`(𝒜,xAByBAz)(𝒜\{A,B\},xyz)`$ or $`(𝒜,xAByABz)(𝒜\{A,B\},xyz)`$ where $`A,B𝒜`$ with $`|B|=\tau (|A|)`$. In both cases the $`S_{}`$-homotopy class of $`w(F)`$ is preserved. Consider the third Reidemeister move. It suffices to consider the case where all three branches are oriented in the same direction, say upwards, and the signs of all crossings are $`+`$. (This is the classical “braid move” $`\sigma _1\sigma _2\sigma _1\sigma _2\sigma _1\sigma _2`$; the moves involving other orientations of the branches and/or other signs of crossings can be obtained as compositions of this move with second Reidemeister moves.) There are 6 cases to consider depending on the order in which one traverses the three branches involved. Let $`I`$ (resp. $`II`$, $`III`$) be the branch connecting the leftmost (resp. intermediate, rightmost) bottom point to the rightmost (resp. intermediate, leftmost) top point. If one traverses these branches in the order $`I,II,III`$ (resp. $`III,II,I`$), then $`|A|=|B|=|C|=a_+`$ (resp. $`|A|=|B|=|C|=b_+`$) and the deformation changes $`w(F)`$ via the third homotopy move (resp. its inverse) where we use that $`(a_+,a_+,a_+)S_{}`$ (resp. that $`(b_+,b_+,b_+)S_{}`$). If one traverses these three branches in the order $`II,III,I`$ (resp. $`I,III,II`$), then the deformation changes $`w(F)`$ via the second homotopy from Lemma 2.2.1 (resp. its inverse) where we use that $`(a_{},a_{},a_+)S_{}`$ (resp. that $`(b_{},b_{},b_+)S_{}`$). Finally, if one traverses these branches in the order $`(II,I,III)`$ (resp. $`III,I,II`$), then the deformation changes $`w(F)`$ via the third homotopy from Lemma 2.2.1 (resp. its inverse) where we use that $`(b_+,b_{},b_{})S_{}`$ (resp. that $`(a_+,a_{},a_{})S_{}`$).
Thus the formula $`Fw(F)`$ defines a mapping $`W:𝒦𝒩(\alpha _{},S_{})`$. We claim that $`W`$ is a bijection. As in the case of curves, for every nanoword $`w`$ over $`\alpha _{}`$ there is a unique (up to stable homeomorphism) pointed knot diagram $`F`$ such that $`w(F)=w`$. Indeed it suffices to realize the underlying nanoword over $`\alpha _0`$ by a pointed curve and then to choose the over/under-crossings to ensure the right signs at all crossings. This implies that $`W`$ is surjective.
To prove the injectivity of $`W`$ it suffices to observe that if two nanowords over $`\alpha _{}`$ are related by the $`i`$-th homotopy move with $`i=1,2,3`$ then they can be represented by pointed knot diagrams related by the $`i`$-th Reidemeister move (effected away from the origin). The cases $`i=1,2`$ are straightforward. For $`i=3`$ the 12 elements of the set $`S_{}`$ lead to all 12 possible choices of over/under-crossings in the third move on Figure 1 leading to admissible Reidemeister moves (for this argument we can assume that the branches $`I,II,III`$ are oriented upwards and traversed either in the order $`I,II,III`$ or $`III,II,I`$).
The last claim of the theorem follows from the definitions. ∎
### 4.3. Examples
The nanoword $`ABCABC`$ with $`|A|=|C|=a_+,|B|=b_+`$ represents a pointed trefoil, see Figure 2 where the thick point is the origin of the diagram. The nanoword $`ABCADCBD`$ with $`|A|=|D|=b_+,|B|=b_{},|C|=a_{}`$ represents a pointed figure eight knot, see Figure 2.
## 5. Eliminating the origin
### 5.1. Shifts
Fix an involution $`\nu `$ in a set $`\alpha `$ called the shift involution. The $`\nu `$-shift of a nanoword $`(𝒜,w:\widehat{n}𝒜)`$ over $`\alpha `$ is the nanoword $`(𝒜^{},w^{}:\widehat{n}𝒜^{})`$ obtained by moving the first letter $`A=w(1)`$ of $`w`$ to the end and applying $`\nu `$ to $`|A|\alpha `$. More precisely, $`𝒜^{}=(𝒜\{A\})\{A_\nu \}`$ where $`A_\nu `$ is a “new” letter not belonging to $`𝒜`$. The projection $`𝒜^{}\alpha `$ extends the given projection $`𝒜\{A\}\alpha `$ by $`|A_\nu |=\nu (|A|)`$. The word $`w^{}`$ in the alphabet $`𝒜^{}`$ is defined by $`w^{}=xA_\nu yA_\nu `$ for $`w=AxAy`$.
Given a homotopy data $`(\alpha ,S)`$ and a shift involution $`\nu `$ in $`\alpha `$, we can quotient the set of nanowords over $`\alpha `$ by the equivalence relation generated by $`S`$-homotopy and $`\nu `$-shifts. The resulting set is denoted $`𝒩(\alpha ,S,\nu )`$. There is a natural projection $`𝒩(\alpha ,S)𝒩(\alpha ,S,\nu )`$ but there is no natural multiplication in $`𝒩(\alpha ,S,\nu )`$.
### 5.2. Non-pointed knots
Stable equivalence can be defined for (non-pointed) knot diagrams through repeating the definition in the pointed case but omitting all references to base points. Denote $`\widehat{𝒦}`$ the set of stable equivalence classes of knot diagrams. Each knot in the cylinder over a surface represents an element in $`\widehat{𝒦}`$ depending only on the isotopy type of the knot.
As we know, a pointed knot diagram gives rise to a nanoword in the alphabet $`\alpha _{}=\{a_+,a_{},b_+,b_{}\}`$. This nanoword is preserved when the origin is pushed along the generic part of the diagram. When the origin jumps over a double point, the nanoword is modified by the $`\nu `$-shift where $`\nu :\alpha _{}\alpha _{}`$ is the involution sending $`a_\pm `$ to $`b_\pm `$. Theorem 4.2.1 implies that $`\widehat{𝒦}=𝒩(\alpha _{},S_{},\nu )`$.
Non-oriented knots can be treated similarly, we do it in the next section in a more general setting of links.
## 6. Nanophrases and links
### 6.1. Nanophrases
A nanophrase of length $`k0`$ over a set $`\alpha `$ is a tuple consisting of an $`\alpha `$-alphabet $`𝒜`$ and a sequence of $`k`$ words $`w_1,\mathrm{},w_k`$ in the alphabet $`𝒜`$ such that their concatenation $`w_1w_2\mathrm{}w_k`$ is a Gauss word in this alphabet. We denote this nanophrase by $`(𝒜,(w_1|w_2|\mathrm{}|w_k))`$ or shorter by $`(w_1|w_2|\mathrm{}|w_k)`$. Note that some of the words $`w_1,\mathrm{},w_k`$ may be empty.
By definition, there is a unique empty nanophrase of length 0 (the corresponding $`\alpha `$-alphabet $`𝒜`$ is void).
Any nanoword $`w`$ over $`\alpha `$ yields a nanophrase $`(w)`$ of length $`1`$. In the sequel we make no difference between nanowords and nanophrases of length 1.
Isomorphism of two nanophrases is an isomorphism of $`\alpha `$-alphabets transforming the first sequence of words into the second one. Given a homotopy data $`\alpha ,\tau ,S`$, we define homotopy moves on nanophrases as in Sect. 2.2 with the only difference that the 2-letter sub-words $`AA,AB,BA,AC,BC`$ etc. modified by these moves may belong to different words of the phrase. Isomorphisms and homotopy moves generate an equivalence relation $`_S`$ of $`S`$-homotopy on the class of nanophrases over $`\alpha `$. Examples:
$$(AB|AC|BC)_S(BA|CA|CB),(AB|ADDCBC)_S(BA|CACB)$$
provided $`(|A|,|B|,|C|)S`$. The length of a nanophrase is preserved under $`S`$-homotopy.
Lemmas 2.2.1 and 2.2.2 extend to nanophrases with the only change that the 2-letter sub-words $`AB,BA,CA`$ etc. may belong to different words of the phrase.
### 6.2. Operations on nanophrases
Fix a homotopy data $`(\alpha ,\tau ,S)`$ and a shift involution $`\nu `$ in $`\alpha `$. We define $`\nu `$-shifts, $`\nu `$-inversions, and $`\nu `$-permutations of words in a nanophrase $`P=(𝒜,(w_1|w_2|\mathrm{}|w_k))`$ over $`\alpha `$.
We can $`\nu `$-shift the $`i`$-th word $`w_i`$ in $`P`$ through moving the first letter, say $`A`$, of $`w_i`$ to the end of $`w_i`$ keeping $`|A|\alpha `$ if $`A`$ appears in $`w_i`$ only once and applying $`\nu `$ to $`|A|`$ if $`A`$ appears in $`w_i`$ twice. All other words in $`P`$ are preserved.
To define inversions, we need more notation. For a word $`w`$ in $`𝒜`$, denote by $`𝒜_w`$ the same alphabet $`𝒜`$ with new projection $`|\mathrm{}|_w`$ to $`\alpha `$ defined as follows: for $`A𝒜`$ set $`|A|_w=\tau (|A|)`$ if $`A`$ occurs in $`w`$ once, $`|A|_w=\nu (|A|)`$ if $`A`$ occurs in $`w`$ twice, and $`|A|_w=|A|`$ otherwise. The $`\nu `$-inversion of the $`i`$-th word in $`P`$ replaces $`w_i`$ with the opposite word $`(w_i)^{}`$ obtained by reading $`w_i`$ from right to left and replaces the $`\alpha `$-alphabet $`𝒜`$ with $`𝒜_{w_i}`$. All other words in $`P`$ are preserved.
The words in $`P`$ can be permuted in an arbitrary way, producing thus new nanophrases over $`\alpha `$. We will need more sophisticated permutations of words depending on $`\nu `$. We begin with notation. For two words $`u,v`$ in the alphabet $`𝒜`$, consider the mapping $`𝒜\alpha `$ sending $`A𝒜`$ to $`\nu (|A|)\alpha `$ if $`A`$ appears both in $`u`$ and $`v`$ and sending $`A`$ to $`|A|`$ otherwise. This mapping makes the set $`𝒜`$ into an $`\alpha `$-alphabet denoted $`𝒜_{uv}`$. For $`i=1,\mathrm{},k1`$, the $`\nu `$-permutation of the $`i`$-th and $`(i+1)`$-st words transforms $`P=(𝒜,(w_1|w_2|\mathrm{}|w_k))`$ into the nanophrase
$$(𝒜_{w_iw_{i+1}},(w_1|w_2|\mathrm{}|w_{i1}|w_{i+1}|w_i|w_{i+2}|\mathrm{}w_k)).$$
This operation is involutive. The $`\nu `$-permutations define an action of the symmetric group $`S_k`$ on the set of nanophrases of length $`k`$.
Denote $`𝒫(\alpha ,S,\nu )`$ the set of nanophrases over $`\alpha `$ quotiented by the equivalence relation generated by $`S`$-homotopy, $`\nu `$-permutations and $`\nu `$-shifts on words. Denote $`𝒫_u(\alpha ,S,\nu )`$ the set of nanophrases over $`\alpha `$ quotiented by the equivalence relation generated by the same operations and the $`\nu `$-inversions.
### 6.3. Link diagrams
Link diagrams on (oriented) surfaces are defined in the same way as knot diagrams with the difference that they may be formed by several (transversal generic oriented closed) curves rather than only one curve. These curves are components of the diagram. A link diagram is pointed if each component is endowed with a base point (the origin) distinct from the crossing points of the diagram. A link diagram is ordered if its components are numerated by $`1,2,\mathrm{},k`$ where $`k`$ is the number of the components. Two ordered pointed link diagrams are stably homeomorphic if there is an orientation preserving homeomorphism of their regular neighborhoods in the ambient surfaces mapping the first diagram onto the second one and preserving the over/undercrossings and the order, the origins and the orientations of the components.
The stable equivalence of ordered pointed link diagrams is generated by the same transformations as in the case of knots. These transformations should preserve the order and the origins of the components; the Reidemeister moves are allowed only away from the origins.
Denote $``$ the set of stable equivalence classes of ordered pointed link diagrams. Recall the homotopy data $`\alpha _{},\tau ,S_{}`$ defined in Sect. 4.2.
###### Theorem 6.3.1.
There is a canonical bijection $`=𝒫(\alpha _{},S_{})`$.
The proof of this theorem is analogous to the proof of Theorem 4.2.1. To write down the nanophrase associated with an ordered pointed link diagram one goes along the first component starting at its origin, then along the second component, etc.
Forgetting the order and the origins of link components, we obtain a notion of stable equivalence for (non-ordered non-pointed) link diagrams. Denote $`\widehat{}`$ the set of equivalence classes of link diagrams. As in the case of knots, each link in the cylinder over a surface represents an element in $`\widehat{}`$ depending only on the isotopy class of this link. Theorem 6.3.1 implies that $`\widehat{}=𝒫(\alpha _{},S_{},\nu )`$ where $`\nu :\alpha _{}\alpha _{}`$ is the involution sending $`a_\pm `$ to $`b_\pm `$.
Additionally forgetting link orientations, we obtain a notion of stable equivalence for unoriented link diagrams (on oriented surfaces). Denote $`\widehat{}_u`$ the set of equivalence classes of unoriented link diagrams. Theorem 6.3.1 implies that $`\widehat{}_u=𝒫_u(\alpha _{},S_{},\nu )`$.
### 6.4. Remarks
1. Theorem 6.3.1 can be extended to framed links. Consider the involution $`\nu \tau =\tau \nu :\alpha _{}\alpha _{}`$ sending $`a_\pm ,b_\pm `$ to $`a_{},b_{}`$, respectively. A framed homotopy $``$ of nanophrases over $`\alpha _{}`$ is defined as the $`S_{}`$-homotopy with the first homotopy move replaced by the following “framed homotopy move” on a nanoword in a nanophrase: $`xAAyBBzxyz`$ provided $`|A|=\nu \tau (|B|)`$. Framed homotopic nanophrases are $`S_{}`$-homotopic; the converse is in general not true. Lemmas 2.2.1 and 2.2.2 for $`\alpha =\alpha _{}`$ extend to this setting by replacing $`_S`$ with $``$. The proof of Lemma 2.2.1 in \[Tu2\] does not use the first homotopy move. The proof of Lemma 2.2.2 in \[Tu2\] uses the first homotopy move but can be easily modified to use the framed move instead.
Let us show that a nanoword $`xAABBy`$ is framed homotopic to $`xy`$ provided $`|A|=\tau (|B|)`$. Pick letters $`C,D`$ not appearing in $`x,y`$ with $`|C|=|B|,|D|=|A|`$. Then
$$xAABByxACDABBCDyxCADBACBDyxDBBDyxy.$$
Here we insert $`CD\mathrm{}CD`$, apply homotopy (iii) of Lemma 2.2.1, delete $`CA\mathrm{}AC`$, and finally delete $`DBBD`$. A similar argument shows that a transformation $`xAAByxBAAy`$ preserves the framed homotopy class. These observations easily imply that framed homotopy classes of nanophrases over $`\alpha _{}`$ bijectively correspond to stable equivalence classes of framed ordered pointed link diagrams.
2. The results of Sect. 5 and 6 have an obvious version for systems of transversal curves on surfaces; it suffices to replace $`\alpha _{}`$ by $`\alpha _0=\{a,b\}`$.
## 7. Pseudo-links and quasi-links
### 7.1. Pseudo-links
Set $`\alpha _1=\{1,1\}`$ with involution $`\tau `$ permuting $`1`$ and $`1`$ and let $`S_1\alpha _1\times \alpha _1\times \alpha _1`$ consist of the following 6 triples:
$$(1,1,1),(1,1,1),(1,1,1),(1,1,1),(1,1,1),(1,1,1).$$
This homotopy data is admissible in the sense of Sect. 2.2. As a shift involution in $`\alpha _1`$, we take the identity mapping $`\mathrm{id}:\alpha _1\alpha _1`$. The corresponding permutations and shifts of words in nanophrases over $`\alpha _1`$ are the ordinary permutations and cyclic shifts of words (involving no modification of the underlying $`\alpha _1`$-alphabets). Nanophrases over $`\alpha _1`$ considered up to permutations and cyclic shifts of words are called pseudo-links.
This terminology is justified by the following connections to knot theory. Consider the projection $`\alpha _{}\alpha _1`$ sending $`a_+,b_+`$ to $`1`$ and $`a_{},b_{}`$ to $`1`$. This projection transforms $`S_{}(\alpha _{})^3`$ into $`S_1`$. It commutes with $`\tau `$ and with the shift involutions in $`\alpha _{},\alpha _1`$ (the shift involution in $`\alpha _{}`$ is defined by $`\nu (a_\pm )=b_\pm `$). Applying the projection $`\alpha _{}\alpha _1`$, we can transform any nanophrase over $`\alpha _{}`$ into a nanophrase over $`\alpha _1`$. Clearly, $`S_{}`$-homotopic nanophrases over $`\alpha _{}`$ yield $`S_1`$-homotopic nanophrases over $`\alpha _1`$. This induces a mapping $`𝒫(\alpha _{},S_{})𝒫(\alpha _1,S_1)`$. Quotienting by permutations and shifts of words we obtain a mapping from $`𝒫(\alpha _{},S_{},\nu )`$ to $`𝒫(\alpha _1,S_1,\mathrm{id})`$. Further quotienting by inversions of words we obtain a mapping $`𝒫_u(\alpha _{},S_{},\nu )𝒫_u(\alpha _1,S_1,\mathrm{id})`$. All these mappings are surjective.
By Sect. 6, a link diagram $`D`$ on a surface yields a nanophrase over $`\alpha _{}`$. Projecting to $`\alpha _1`$, we obtain a nanophrase $`p(D)`$ over $`\alpha _1`$. If $`D`$ is pointed and ordered, then $`p(D)`$ is well-defined, otherwise $`p(D)`$ is defined only up to permutations and shifts of words. The class of $`p(D)`$ in $`𝒫(\alpha _1,S_1,\mathrm{id})`$ is an invariant of stable equivalence of $`D`$. We call $`p(D)`$ the underlying pseudo-link of $`D`$. Further projecting to $`𝒫_u(\alpha _1,S_1,\mathrm{id})`$ we obtain an invariant independent of the orientation of $`D`$.
In the next section we explain that the Jones polynomial of a link depends only on the underlying pseudo-link. This shows that pseudo-links are highly non-trivial objects retaining important features of links.
### 7.2. Quasi-links
Set $`\alpha _2=\{c,d\}`$ with the identity involution $`\tau =\mathrm{id}:\alpha _2\alpha _2`$ and let $`S_2\alpha _2\times \alpha _2\times \alpha _2`$ consist of the following 6 triples:
$$(c,c,c),(c,c,d),(d,c,c),(d,d,d),(d,d,c),(c,d,d).$$
This homotopy data essentially differs from $`(\alpha _1,S_1)`$ by the choice of $`\tau `$. The homotopy data $`(\alpha _2,S_2)`$ is admissible. As a shift involution $`\nu _2`$ in $`\alpha _2`$, we take the permutation of $`c`$ and $`d`$. Nanophrases over $`\alpha _2`$ considered up to $`\nu _2`$-permutations and $`\nu _2`$-shifts of words are called quasi-links.
Connections to knot theory go as follows. Consider the projection $`\alpha _{}\alpha _2`$ sending $`a_+,b_{}`$ to $`c`$ and $`a_{},b_+`$ to $`d`$. This projection transforms $`S_{}(\alpha _{})^3`$ into $`S_2`$, commutes with $`\tau `$ and with the shift. Applying this projection, we can transform any nanophrase over $`\alpha _{}`$ into a nanophrase over $`\alpha _2`$. This induces a mapping $`𝒫(\alpha _{},S_{})𝒫(\alpha _2,S_2)`$. Quotienting by permutations and shifts of words (and eventually by inversions of words) we obtain projections $`𝒫(\alpha _{},S_{},\nu )𝒫(\alpha _2,S_2,\nu _2)`$ and $`𝒫_u(\alpha _{},S_{},\nu )𝒫_u(\alpha _2,S_2,\nu _2)`$.
A link diagram $`D`$ yields a nanophrase over $`\alpha _{}`$ whose projection to $`\alpha _2`$ is a nanophrase over $`\alpha _2`$ denoted $`q(D)`$. If $`D`$ is pointed and ordered, then $`q(D)`$ is well defined, otherwise $`q(D)`$ is defined only up to $`\nu _2`$-permutations and $`\nu _2`$-shifts of words. The class of $`q(D)`$ in $`𝒫(\alpha _2,S_2,\nu _2)`$ is an invariant of stable equivalence of $`D`$. We call $`q(D)`$ the underlying quasi-link of $`D`$. Further projecting to $`𝒫_u(\alpha _2,S_2,\nu _2)`$ we obtain an invariant independent of the orientation of $`D`$.
Quasi-links will be further discussed in Sect. 9.
### 7.3. Remarks
1. As explained above, there are three natural projections from the set of link diagrams on surfaces to simpler objects. They map a link diagram to the underlying family of curves, the underlying pseudo-link and the underlying quasi-link. The length of the resulting nanophrases is equal to the number of link components.
2. The homotopy data $`(\alpha _0,S_0)`$ and $`(\alpha _1,S_1)`$ are closely related. Consider the bijection from $`\alpha _0=\{a,b\}`$ to $`\alpha _1=\{1,1\}`$ sending $`a`$ to $`1`$ and $`b`$ to $`1`$. This bijection commutes with the involution $`\tau `$ in $`\alpha _0,\alpha _1`$ and transforms $`S_0(\alpha _0)^3`$ into a subset of $`S_1(\alpha _1)^3`$. In this way any nanophrase over $`\alpha _0`$ determines a pseudo-link and homotopic nanophrases yield homotopic pseudo-links. Thus, the homotopy theory of pseudo-links is a quotient of the homotopy theory of curves. However, there is no way to recover the pseudo-link $`p(D)`$ underlying a link diagram $`D`$ from the system of curves underlying $`D`$. Note also that the shift involutions in $`\alpha _0`$ and $`\alpha _1`$ do not match: the first one permutes $`a`$ and $`b`$ while the second one is the identity.
## 8. The bracket polynomial
The aim of this section is to construct a polynomial invariant of pseudo-links whose value on the underlying pseudo-link of a link diagram is equal to the Jones polynomial of the link. We begin by recalling Kauffman’s bracket polynomial.
### 8.1. Bracket polynomial of links
L. Kauffman \[Kau1\] associated with every non-empty link diagram $`D`$ on a surface a 1-variable Laurent polynomial $`D`$ called the bracket polynomial of $`D`$. This polynomial is defined by expanding each crossing of $`D`$ as a linear combination of two uncrossings with coefficients $`t`$ and $`t^1`$, see Figure 3. This expands $`D`$ as a linear combination of diagrams with no crossings. Each $`d`$-component diagram with no crossings is then replaced with $`(t^2+t^2)^{d1}`$. The bracket polynomial depends neither on the orientation of $`D`$, nor on an order of its components, nor on a choice of base points. The bracket polynomial is invariant under the second and the third Reidemeister moves and is multiplied by $`t^{\pm 3}`$ under the first Reidemeister move.
We can translate the bracket polynomial to the language of nanophrases over $`\alpha _{}`$. A nanophrase $`P`$ over $`\alpha _{}`$ gives rise to a pointed ordered link diagram on a surface. Let $`P[t,t^1]`$ be the bracket polynomial of this diagram. This polynomial is invariant under $`\nu `$-shifts, $`\nu `$-inversions, and $`\nu `$-permutations on words in $`P`$ since they are translated to diagrams as change of base points, orientation reversal, and change of order of components. The polynomial $`P`$ is invariant under the second and third $`S_{}`$-homotopy moves on $`P`$ since they are translated to diagrams as the second and third Reidemeister moves. Under the move deleting $`AA`$ from a word of $`P`$, the bracket polynomial is multiplied by $`t^{3\epsilon (A)}`$ where $`\epsilon (A)=+1`$ if $`|A|\{a_+,b_+\}`$ and $`\epsilon (A)=1`$ if $`|A|\{a_{},b_{}\}`$. When an empty word $`\mathrm{}`$ is deleted from a nanophrase $`P`$ of length $`2`$ the polynomial $`P`$ is divided by $`(t^2+t^2)`$. Clearly, $`(\mathrm{})=1`$.
To translate the Kauffman crossing expansion to this setting, we need the following notation. Given a phrase $`P`$ in an $`\alpha _{}`$-alphabet $`𝒜`$ and a word $`w`$ in the alphabet $`𝒜`$, denote by $`P_w`$ the same phrase $`P`$ in the $`\alpha _{}`$-alphabet $`𝒜_w`$ defined in Sect. 6.2. Denote by $`w^{}`$ the word in the alphabet $`𝒜`$ obtained by writing the letters of $`w`$ in the opposite order.
The Kauffman crossing expansion applied to a self-crossing of a link component and to a crossing of two different components implies the following two recursive relations for the bracket polynomial of nanophrases:
$$(P_1|AwAz|P_2)=t^{\epsilon (A)}(P_1|w|z|P_2)+t^{\epsilon (A)}(P_1|w^{}z|P_2)_w,$$
$$(P_1|Aw|Az|P_2)=t^{\epsilon (A)}(P_1|wz|P_2)+t^{\epsilon (A)}(P_1|w^{}z|P_2)_w.$$
Here $`w`$ and $`z`$ are words in an $`\alpha _{}`$-alphabet $`𝒜`$, $`A`$ is a letter in $`𝒜`$, and $`P_1,P_2`$ are finite sequences of words in $`𝒜`$ such that every letter of $`𝒜`$ appears in the phrase $`(P_1|AwAz|P_2)`$ twice. In the first formula $`w`$ and $`z`$ are parts of a word $`AwAz`$ while in the second formula $`Aw`$ and $`Az`$ are two consecutive words. These formulas and the properties of $`P`$ listed above allow us to compute this polynomial recursively.
### 8.2. Bracket polynomial of pseudo-links
In the recursive formulas above, the right hand side depends on $`\epsilon (A)`$ rather than on $`A`$. This implies (by induction on the number of letters in a nanophrase) that the bracket of a nanophrase over $`\alpha _{}`$ depends only on the underlying nanophrase over $`\alpha _1`$. Any pseudo-link $`p`$ determines a Laurent polynomial $`p[t^{\pm 1}]`$ by $`p=\stackrel{~}{p}`$ where $`\stackrel{~}{p}`$ is any nanophrase over $`\alpha _{}`$ whose projection to $`\alpha _1`$ equals $`p`$. The polynomial $`p`$ is invariant under shifts, inversions, and permutations of words in $`p`$. It is preserved under the second and third $`S_1`$-homotopy moves and is multiplied by $`t^{3|A|}`$ under the move deleting $`AA`$ from a word of $`p`$ (where $`|A|\alpha _1=\{1,1\}`$). To compute $`p`$ one can use the recursive relations above with $`\epsilon (A)`$ replaced everywhere by $`|A|\{1,1\}`$.
As an illustration, we compute the bracket for the nanoword $`ABCABC`$ over $`\alpha _1`$ where $`|A|=|B|=|C|=1`$. We have
$$ABCABC=t(BC|BC)+t^1(C_\tau B_\tau B_\tau C_\tau )$$
$$=t(t(CC)+t^1(C_\tau C_\tau ))+t^1(t^{3|B_\tau |})(C_\tau C_\tau )$$
$$=t^2(t^{3|C|})t^{3|C_\tau |}+t^1(t^{3|B_\tau |})(t^{3|C_\tau |})=t^5t^3+t^7$$
where $`|B_\tau |=\tau (|B|)=1`$ and $`|C_\tau |=\tau (|C|)=1`$. This is compatible with the usual formula for the bracket of the standard diagram of a trefoil.
### 8.3. The Jones polynomial of pseudo-links
For a pseudo-link $`p`$ we define the writhe $`|p|=_A|A|`$ where $`A`$ runs over all letters occurring in $`p`$. The polynomial $`J(p)=(t)^{3|p|}p`$ is invariant under all $`S_1`$-homotopy moves on $`p`$. For a pseudo-link $`p`$ arising from a link in $`S^3`$, the polynomial $`p`$ is equal to the Jones polynomial of this link up to a re-parametrization.
### 8.4. Polynomials of phrases
An $`\alpha _1`$-alphabet is nothing but a bipartitioned set, that is a set decomposed as a disjoint union of two subsets (the preimages of $`\pm 1\alpha _1`$). Any phrase $`P`$ in an $`\alpha _1`$-alphabet $`𝒜`$ gives rise to a polynomial $`P[t^{\pm 1}]`$ as follows. It is explained in \[Tu2\] that a word $`w`$ in any alphabet determines in a canonical way a nanoword $`w^d`$ over this alphabet. The same procedure applies to phrases and derives from $`P`$ a nanophrase $`P^d`$ over $`𝒜`$. (Each letter $`A𝒜`$ occurring $`m_A`$ times in $`P`$ gives rise to $`m_A(m_A+1)/2`$ distinct letters each occurring in $`P^d`$ twice.) Composing with projection $`𝒜\alpha _1`$ we obtain from $`P^d`$ a pseudo-link $`(P^d)_1`$. Set $`P=(P^d)_1`$. Similarly, we define the Jones polynomial of $`P`$ by $`J(P)=J((P^d)_1)=(t)^{3|P|}P`$ with $`|P|=_{A𝒜}|A|m_A(m_A+1)/2`$ where $`|A|=\pm 1`$ is the image of $`A`$ in $`\alpha _1`$ and $`m_A`$ is the number of entries of $`A`$ in $`P`$. The polynomials $`P,J(P)`$ are interesting invariants of phrases in bipartitioned alphabets. One natural question is to characterize the polynomials that arise from phrases in this way.
## 9. Keis
Keis were introduced by M. Takasaki in 1942 as abstractions of symmetries, see \[Kam\] for a survey of keis and related objects (quandles, racks, etc.). Here we recall from \[Tu2\] the concept of an $`\alpha `$-kei where $`\alpha `$ is a set with involution $`\tau `$. This will be instrumental in the next section where we discuss keis of nanophrases.
### 9.1. $`\alpha `$-keis
Consider a set $`X`$ and suppose that each $`a\alpha `$ gives rise to a bijection $`xax:XX`$ and to a binary operation $`(x,y)x_ay`$ on $`X`$. The set $`X`$ is an $`\alpha `$-kei and the mappings $`xax,(x,y)x_ay`$ are kei operations if the following axioms are satisfied:
(i) $`ax_ax=x`$ for all $`a\alpha ,xX`$;
(ii) $`a(x_ay)=ax_aay`$ for all $`a\alpha ,x,yX`$;
(iii) $`(x_ay)_az=(x_aaz)_a(y_az)`$ for all $`a\alpha ,x,y,zX`$;
(iv) $`a\tau (a)x=x`$ for all $`xX,a\alpha `$ and
(v) $`(x_ay)_{\tau (a)}ay=x`$ for all $`x,yX,a\alpha `$.
A morphism of $`\alpha `$-keis $`XX^{}`$ is a set-theoretic mapping commuting with the kei operations in $`X,X^{}`$. Given an $`\alpha `$-kei $`X`$, we define an $`\alpha `$-kei $`\overline{X}`$ to be the same set $`X`$ with new kei operations $`ax:=\tau (a)x`$, $`x_ay:=x_{\tau (a)}y`$ for $`x,yX,a\alpha `$. Clearly, $`\overline{\overline{X}}=X`$.
In analogy with group theory, one can define presentations of $`\alpha `$-keis by generators and relations. A presentation of an $`\alpha `$-kei $`X`$ by generators and relations yields a presentation of $`\overline{X}`$ by generators and relations by replacing every letter $`a\alpha `$ appearing in the relations by $`\tau (a)`$.
In the simplest case where $`\alpha =\{a\}`$ is a 1-element set, an $`\alpha `$-kei is a set $`X`$ with involution $`x\stackrel{~}{x}=ax`$ and a binary operation $`(x,y)xy=x_ay`$ such that $`\stackrel{~}{x}x=x`$; $`\stackrel{~}{xy}=\stackrel{~}{x}\stackrel{~}{y}`$; $`(xy)z=(x\stackrel{~}{z})(yz)`$, and $`(xy)\stackrel{~}{y}=x`$ for all $`x,y,zX`$. When the involution $`x\stackrel{~}{x}`$ is the identity, these axioms are equivalent to those of a kei, see \[Kam\].
The $`\alpha `$-keis generalize quandles: there is a canonical bijection between quandles and $`\alpha `$-keis $`X`$ such that $`\alpha =\{a,b\}`$ is a 2-element set with involution permuting $`a,b`$ and $`ax=bx=x`$ for all $`xX`$, cf. \[Tu2\], Lemma 14.7.1.
### 9.2. Core $`\alpha `$-keis
The following construction of $`\alpha `$-keis provides a vast set of examples. By a $`\tau `$-compatible action of $`\alpha `$ on a group $`G`$ we mean a set of group automorphisms $`\{GG,gag\}_{a\alpha }`$ such that $`a\tau (a)g=g`$ for all $`a\alpha ,gG`$. It is easy to check that such an action together with kei operations $`g_ah=h(\tau (a)g)^1h`$ make $`G`$ into an $`\alpha `$-kei. It is called the core of $`G`$ and denoted $`\mathrm{core}(G)`$. For $`\alpha `$ consisting of one element that acts on $`G`$ as the identity, this construction is due to D. Joyce (cf. \[FR\], p. 349).
The construction of the core has a natural adjoint associating with an arbitrary $`\alpha `$-kei $`X`$ a group $`\mathrm{\Gamma }_X`$ with generators $`\{[x]\}_{xX}`$ and relations $`[a(x_by)]=[ay][a\tau (b)x]^1[ay]`$ for all $`a,b\alpha ,x,yX`$. We endow $`\mathrm{\Gamma }_X`$ with the $`\tau `$-compatible action of $`\alpha `$ defined on the generators by $`a[x]=[ax]`$ for $`a\alpha ,xX`$. Given a group $`G`$ with $`\tau `$-compatible action of $`\alpha `$ and a kei morphism $`f:X\mathrm{core}(G)`$, there is a unique group homomorphism $`\mathrm{\Gamma }_XG`$ whose composition with the inclusion $`X\mathrm{\Gamma }_X,x[x]`$ is equal to $`f`$. This universal property characterizes $`\mathrm{\Gamma }_X`$ up to isomorphism.
A presentation of $`\mathrm{\Gamma }_X`$ by generators and relations can be read from an arbitrary presentation $`[S:R]`$ of $`X`$ by generators and relations (cf. \[FR\], Lemma 4.3). Namely, $`\mathrm{\Gamma }_X`$ is generated by the symbols $`\{as\}_{a\alpha ,sS}`$ subject to the relations obtained from $`R`$ by replacing all terms of type $`a(x_by)`$ by $`(ay)(a\tau (b)x)^1(ay)`$.
## 10. Keis of nanophrases
A well known construction due to S. Matveev and D. Joyce associates quandles with link diagrams. Since link diagrams are nanophrases over $`\alpha _{}`$, one may attempt to generalize this construction to nanophrases over an arbitrary alphabet $`\alpha `$ with involution $`\tau `$. We do it here starting with certain additional data.
### 10.1. Keis and homotopy
Fix an equivalence relation $``$ on $`\alpha `$ such that $`ab\tau (a)\tau (b)`$ for $`a,b\alpha `$. Let $`\underset{¯}{\alpha }=\alpha /`$ with involution $`\underset{¯}{\tau }`$ induced by $`\tau `$. For $`a\alpha `$, denote its projection to $`\underset{¯}{\alpha }`$ by $`\underset{¯}{a}`$.
Fix a set (possibly empty) $`\beta \alpha `$ such that $`\tau (\beta )=\beta `$. We associate with any nanophrase $`P=(𝒜,(w_1,\mathrm{},w_k))`$ over $`\alpha `$ an $`\underset{¯}{\alpha }`$-kei $`\kappa _\beta (P)`$ as follows. Let $`n_r`$ be the length of the word $`w_r`$ for $`r=1,\mathrm{},k`$. Each letter $`A𝒜`$ appears in $`P`$ twice, say, first time at the $`i_1`$-th position in $`w_{r_1}`$ and second time at the $`i_2`$-th position in $`w_{r_2}`$ where $`1i_1n_{r_1},1i_2n_{r_2}`$, $`r_1r_2`$, and $`r_1=r_2i_1<i_2`$. The $`\underset{¯}{\alpha }`$-kei $`\kappa _\beta (P)`$ is generated by the symbols $`\{x_s^r\}`$ where $`1rk`$ and $`0sn_r`$. Each $`A𝒜`$ gives rise to two defining relation: if $`a=|A|\beta `$, then
$$x_{i_1}^{r_1}=\underset{¯}{a}x_{i_11}^{r_1},x_{i_2}^{r_2}=x_{i_21}^{r_2}_{\underset{¯}{a}}x_{i_11}^{r_1},$$
and if $`a=|A|\alpha \beta `$, then
$$x_{i_1}^{r_1}=x_{i_11}^{r_1}_{\underset{¯}{a}}x_{i_21}^{r_2},x_{i_2}^{r_2}=\underset{¯}{a}x_{i_21}^{r_2}.$$
These generators and relations define the $`\underset{¯}{\alpha }`$-kei $`\kappa _\beta (P)`$. It has two sets of distinguished elements $`x_0^1,x_0^2,\mathrm{},x_0^k`$ (the inputs) and $`x_{n_1}^1,x_{n_2}^2,\mathrm{},x_{n_k}^k`$ (the outputs). Adding the relations $`x_0^r=x_{n_r}^r`$ for $`r=1,\mathrm{},k`$ we obtain a quotient $`\underset{¯}{\alpha }`$-kei $`\widehat{\kappa }_\beta (P)`$.
Note the obvious $`\underset{¯}{\alpha }`$-kei isomorphism $`\kappa _\beta (P^{})\overline{\kappa _{\alpha \beta }(P)}`$ where $`P^{}`$ is $`P`$ read from right to left. This isomorphism transforms the $`r`$-th input (resp. output) into the $`(n+1r)`$-th output (resp. input). Clearly, $`\widehat{\kappa }_\beta (P^{})\overline{\widehat{\kappa }_{\alpha \beta }(P)}`$.
For the next definition, it is convenient to set $`\beta _0=\beta `$ and $`\beta _1=\alpha \beta `$. Let $`S=S(\beta ,)\alpha ^3`$ consist of all triples $`(a,b,c)\alpha ^3`$ such that
\- $`abc`$ and $`a,b,c\beta _i`$ for some $`i\{0,1\}`$;
\- or $`ab\tau (c)`$ and $`a,b\beta _i,c\beta _{1i}`$ for some $`i\{0,1\}`$;
\- or $`\tau (a)bc`$ and $`b,c\beta _i,a\beta _{1i}`$ for some $`i\{0,1\}`$.
The set $`S`$ contains the diagonal of $`\alpha ^3`$ and therefore the homotopy data $`(\alpha ,S)`$ is admissible.
###### Theorem 10.1.1.
For any nanophrase $`P`$ over $`\alpha `$, the $`\underset{¯}{\alpha }`$-kei $`\kappa _\beta (P)`$ is invariant under $`S`$-homotopy moves.
The proof goes by repeating the proof of Lemma 15.1.1 in \[Tu2\].
The next theorem shows that for an appropriate choice of the shift involution $`\nu `$, the kei $`\kappa _\beta (P)`$ is also preserved under $`\nu `$-permutations on $`P`$ and its quotient $`\widehat{\kappa }_\beta (P)`$ is preserved under $`\nu `$-permutations and $`\nu `$-shifts.
###### Theorem 10.1.2.
Let $`\nu :\alpha \alpha `$ be an involution such that $`\nu (\beta )=\alpha \beta `$ and $`a\nu (a)`$ for all $`a\alpha `$. For a nanophrase $`P`$ over $`\alpha `$, the $`\underset{¯}{\alpha }`$-kei $`\kappa _\beta (P)`$ is invariant under $`\nu `$-permutations on the words of $`P`$. The quotient $`\underset{¯}{\alpha }`$-kei $`\widehat{\kappa }_\beta (P)`$ is invariant under $`\nu `$-permutations and $`\nu `$-shifts on the words of $`P`$.
We leave the proof to the reader as an exercise.
### 10.2. Examples
1. Consider the alphabet $`\alpha _{}=\{a_+,a_{},b_+,b_{}\}`$ with involution $`\tau (a_\pm )=b_{}`$, shift involution $`\nu (a_\pm )=b_\pm `$, and distinguished subset $`\beta =\{a_+,b_{}\}`$. Provide $`\alpha _{}`$ with equivalence relation $`a_+b_+,a_{}b_{}`$. This data satisfies all the conditions of Theorems 10.1.1 and 10.1.2. Clearly, $`S(\beta ,)=S_{}(\alpha _{})^3`$ is the set defined in Sect. 4.2. This yields for any nanophrase $`P`$ over $`\alpha _{}`$ an $`\underset{¯}{\alpha }_{}`$-kei $`\kappa _\beta (P)`$ invariant under $`S_{}`$-homotopy. The quotient $`\underset{¯}{\alpha }_{}`$-kei $`\widehat{\kappa }_\beta (P)`$ is also invariant under $`\nu `$-shifts and $`\nu `$-permutations. The set $`\underset{¯}{\alpha }_{}=\{+,\}`$ consists of 2 elements permuted by $`\underset{¯}{\tau }`$. For the nanophrase $`P`$ associated with a link diagram on a surface, the $`\underset{¯}{\alpha }_{}`$-kei $`\widehat{\kappa }_\beta (P)`$ is invariant under stable equivalence and independent of the choice of the order and the base points of the link components. The quotient of $`\widehat{\kappa }_\beta (P)`$ by $`ax=x`$ for all $`a\alpha _{}=\{+,\},x\widehat{\kappa }_\beta (P)`$ with binary operation $`_+`$ is the standard link quandle (see \[FR\], \[Kam\], \[Kau2\]).
2. Consider the alphabet $`\alpha _0=\{a,b\}`$ with involution $`\tau (a)=b`$ and distinguished subset $`\beta _0=\alpha `$. As the equivalence relation $``$ in $`\alpha _0`$ we take the equality $`=`$. This data satisfies the conditions of Theorem 10.1.1 where $`S(\beta _0,)=S_0(\alpha _0)^3`$ is the diagonal. For any nanophrase $`P`$ over $`\alpha _0`$, we obtain an $`\underset{¯}{\alpha }_0`$-kei $`\kappa _{\beta _0}(P)`$ invariant under $`S_0`$-homotopy. This example is contained in the previous one: the mapping $`\alpha _0\alpha _{}`$ defined by $`aa_+,bb_{}`$ transforms $`P`$ into a nanophrase $`P_{}`$ over $`\alpha _{}`$ and $`\kappa _{\beta _0}(P)=\kappa _\beta (P_{})`$.
3. The homotopy data $`(\alpha _1,S_1)`$ from Sect. 7 cannot be obtained by the methods of Sect. 10.1 and does not lead to keis. Pseudo-links have no keis.
4. Consider the alphabet $`\alpha _2=\{c,d\}`$ with trivial involution $`\tau =\mathrm{id}`$, shift involution $`\nu `$ permuting $`c`$ and $`d`$, and distinguished subset $`\beta =\{c\}`$. Provide $`\alpha _2`$ with trivial equivalence relation $``$ (all elements are equivalent). This data satisfies the conditions of Theorems 10.1.1 and 10.1.2. Clearly, $`S(\beta ,)=S_2(\alpha _2)^3`$ is the set defined in Sect. 7.2. This yields for any quasi-link $`P`$ (= a nanophrase over $`\alpha _2`$) an $`\underset{¯}{\alpha }_2`$-kei $`\kappa _\beta (P)`$ invariant under $`S_2`$-homotopy. The quotient $`\underset{¯}{\alpha }_2`$-kei $`\widehat{\kappa }_\beta (P)`$ is also invariant under $`\nu `$-shifts and $`\nu `$-permutations. The set $`\underset{¯}{\alpha }_2`$ consists here of 1 element. A study of this kei should lead to interesting homotopy invariants of quasi-links.
When $`P`$ is obtained from an oriented link $`LS^3`$ by taking the associated nanophrase over $`\alpha _{}`$ and projecting to $`\alpha _2`$, the group $`\mathrm{\Gamma }=\mathrm{\Gamma }_{\widehat{\kappa }_\beta (P)}`$ (defined in Sect. 9.2) is closely related to the fundamental group of the 2-fold branched cover $`M`$ of $`S^3`$ with branching set $`L`$. Namely, the group $`\pi _1(M)`$ is the quotient of $`\mathrm{\Gamma }`$ by the relations $`ag=g`$ for the unique $`a\underset{¯}{\alpha }_2`$ and all $`g\mathrm{\Gamma }`$. This is obtained by comparing the presentation of $`\widehat{\kappa }_\beta (P)`$ as above with the Wada presentation of $`\pi _1(M)`$, both computed from a diagram of $`L`$. This observation easily extends to the nanophrase derived from (a diagram of) a link $`L\mathrm{\Sigma }\times [0,1]`$ where $`\mathrm{\Sigma }`$ is a surface. Here one should use the 2-fold branched cover of $`\mathrm{\Sigma }\times [0,1]/\mathrm{\Sigma }\times 1`$ with branching set $`L`$, cf. the argument in \[KK\], Prop. 5.1 and the proof of Wada’s theorem in \[Pr\], p. 287.
5. Let $`\alpha _{},\tau ,\nu ,\beta ,`$ be as in Example 1. Pick a set $`\gamma `$ and consider the alphabet $`\alpha _\gamma =\alpha _{}\times \gamma `$ with involution $`\tau \times \mathrm{id}`$, shift involution $`\nu _\gamma =\nu \times \mathrm{id}`$, and distinguished subset $`\beta \times \gamma `$. Provide $`\alpha _\gamma `$ with equivalence relation $`_\gamma `$ as follows: two pairs $`(x,c),(y,d)`$ with $`x,y\alpha _{},c,d\gamma `$ are equivalent if $`xy`$ and $`c=d`$. This data satisfies the conditions of Theorems 10.1.1 and 10.1.2. The set $`S_\gamma =S(\beta \times \gamma ,_\gamma )`$ is the product of $`S_{}(\alpha _{})^3`$ and the diagonal of $`\gamma ^3`$. This yields for any nanophrase $`P`$ over $`\alpha _\gamma `$ an $`\underset{¯}{\alpha _\gamma }`$-kei $`\kappa _{\beta \times \gamma }(P)`$ invariant under $`S_\gamma `$-homotopy where $`\underset{¯}{\alpha _\gamma }=\alpha _\gamma /_\gamma =\{+,\}\times \gamma `$. The quotient $`\underset{¯}{\alpha _\gamma }`$-kei $`\widehat{\kappa }_{\beta \times \gamma }(P)`$ is also invariant under $`\nu _\gamma `$-shifts and $`\nu _\gamma `$-permutations.
Nanophrases over $`\alpha _\gamma `$ have a simple geometric interpretation. Let us call an (ordered pointed) link diagram on an (oriented) surface $`\gamma `$-colored if all its crossings are endowed with elements of $`\gamma `$ (the colors). Homeomorphisms of $`\gamma `$-colored link diagrams should preserve the colors of the crossings. Stable equivalence of $`\gamma `$-colored link diagrams is defined as in the non-colored case with the following restrictions on the Reidemeister moves: the second move is allowed only when it involves two crossings of the same color, the third move is allowed only when it involves three crossings of the same color which is kept under the move. The crossings not involved in the moves keep their color. Denote $`_\gamma `$ the set of stable equivalence classes of $`\gamma `$-colored ordered pointed link diagrams. The same arguments as in Theorem 6.3.1 show that $`_\gamma =𝒫(\alpha _\gamma ,S_\gamma )`$. This equality implies a similar equality for non-ordered non-pointed link diagrams. The construction above associates with every $`\gamma `$-colored link diagram an $`\underset{¯}{\alpha _\gamma }`$-kei invariant under stable equivalence.
### 10.3. Remark
Further invariants of nanophrases can be derived from their keis by abelianization \[Tu2\]. Another interesting possibility is to define cohomology of $`\alpha `$-keis and to derive homotopy invariants of nanophrases from cocycles and state sums. |
warning/0506/math-ph0506057.html | ar5iv | text | # 1 A schematic sketch of the structure of the simplest projective Hjelmslev plane, 𝑃𝐻(2,2). Shown are all the 28 of its points (represented by small filled circles), grouped into seven pairwise disjoint sets (neighbour classes), each of cardinality four, as well as 24 of its lines (drawn as solid, dashed, dotted and dot-dashed curves), forming six different neighbour classes. In order to make the sketch more illustrative, different neighbour classes of lines have different colour. Also shown is the associated Fano plane, 𝑃𝐺(2,2), whose points are represented by seven big circles, six of its lines are drawn as pairs of line segments and the remaining line as a pair of concentric circles. Notice the intricate character of pairwise intersection of the lines of 𝑃𝐻(2,2); two lines from distinct neighbour classes have just one point in common, while any two lines within a neighbour class share (𝑞=)2 points, both of the same neighbour class.
Hjelmslev Geometry of Mutually Unbiased Bases
Metod Saniga and Michel Planat
Astronomical Institute, Slovak Academy of Sciences,
05960 Tatranská Lomnica, Slovak Republic
(msaniga@astro.sk)
and
Institut FEMTO-ST, CNRS, Laboratoire de Physique et Métrologie des Oscillateurs,
32 Avenue de l’Observatoire, F-25044 Besançon, France
(planat@lpmo.edu)
Abstract
The basic combinatorial properties of a complete set of mutually unbiased bases (MUBs) of a $`q`$-dimensional Hilbert space $`_q`$, $`q=p^r`$ with $`p`$ being a prime and $`r`$ a positive integer, are shown to be qualitatively mimicked by the configuration of points lying on a proper conic in a projective Hjelmslev plane defined over a Galois ring of characteristic $`p^2`$ and rank $`r`$. The $`q`$ vectors of a basis of $`_q`$ correspond to the $`q`$ points of a (so-called) neighbour class and the $`q+1`$ MUBs answer to the total number of (pairwise disjoint) neighbour classes on the conic.
MSC Codes: 51C05 – 81R99 – 81Q99
PACS Numbers: 02.10.Hh – 02.40.Dr – 03.65.Ca
Keywords: Projective Hjelmslev Planes – Proper Conics – Mutually Unbiased Bases
Two distinct orthonormal bases of a $`q`$-dimensional Hilbert space, $`_q`$, are said to be mutually unbiased if all inner products between any element of the first basis and any element of the second basis are of the same value $`1/\sqrt{q}`$. This concept plays a key role in a search for a rigorous formulation of quantum complementarity and lends itself to numerous applications in quantum information theory. It is a well-known fact (see, e.g., and references therein) that $`_q`$ supports at most $`q+1`$ pairwise mutually unbiased bases (MUBs) and various algebraic geometrical constructions of such $`q+1`$, or complete, sets of MUBs have been found when $`q=p^r`$, with $`p`$ being a prime and $`r`$ a positive integer. In our recent paper we have demonstrated that the bases of such a set can be viewed as points of a proper conic (or, more generally, of an oval) in a projective plane of order $`q`$. In this short note we extend and qualitatively finalize this picture by showing that also individual vectors of every such a basis can be represented by points, although these points are of a different nature and require a more general projective setting, that of a projective Hjelmslev plane .
To this end in view, we shall first introduce the basics of the Galois ring theory (see, e.g., for the symbols, notation and further details). Let, as above, $`p`$ be a prime number and $`r`$ a positive integer, and let $`f(x)𝒵_{p^2}[x]`$ be a monic polynomial of degree $`r`$ whose image in $`𝒵_p[x]`$ is irreducible. Then $`GR(p^2,r)𝒵_{p^2}[x]/(f)`$ is a ring, called a Galois ring, of characteristic $`p^2`$ and rank $`r`$, whose maximal ideal is $`pGR(p^2,r)`$. In this ring there exists a non-zero element $`\zeta `$ of order $`p^r1`$ that is a root of $`f(x)`$ over $`𝒵_{p^2}`$, with $`f(x)`$ dividing $`x^{p^r1}1`$ in $`𝒵_{p^2}[x]`$. Then any element of $`GR(p^2,r)`$ can uniquely be written in the form
$$g=a+pb,$$
(1)
where both $`a`$ and $`b`$ belong to the so-called Teichmüller set $`𝒯_r`$,
$$𝒯_r\{0,1,\zeta ,\zeta ^2,\mathrm{},\zeta ^{p^r2}\},$$
(2)
having
$$q=p^r$$
(3)
elements. From Eq. (1) it is obvious that $`g`$ is a unit (i.e., an invertible element) of $`GR(p^2,r)`$ iff $`a0`$ and a zero-divisor iff $`a=0`$. It then follows that $`GR(p^2,r)`$ has $`\mathrm{\#}_\mathrm{t}=q^2`$ elements in total, out of which there are $`\mathrm{\#}_\mathrm{z}=q`$ zero-divisors and $`\mathrm{\#}_\mathrm{u}=q^2q=q(q1)`$ units. Next, let “$`^\overline{}`$” denote reduction modulo $`p`$; then obviously $`\overline{𝒯}_r=GF(q)`$, the Galois field of order $`q`$, and $`\overline{\zeta }`$ is a primitive element of $`GF(q)`$. Finally, one notes that any two Galois rings of the same characteristic and rank are isomorphic.
Now we have a sufficient algebraic background to introduce the concept of a projective Hjelmslev plane over $`GR(p^2,r)`$, henceforth referred to as $`PH(2,q)`$.<sup>1</sup><sup>1</sup>1This is, of course, a very specific, and rather elementary, kind of projective Hjelmslev plane; its most general, axiomatic definition can be found, for example, in . $`PH(2,q)`$ is an incidence structure whose points are classes of ordered triples $`(\varrho \stackrel{˘}{x}_1,\varrho \stackrel{˘}{x}_2,\varrho \stackrel{˘}{x}_3)`$, where both $`\varrho `$ and at least one $`\stackrel{˘}{x}_i`$ ($`i`$=1,2,3) are units, whose lines are classes of ordered triples $`(\sigma \stackrel{˘}{l}_1,\sigma \stackrel{˘}{l}_2,\sigma \stackrel{˘}{l}_3)`$, where both $`\sigma `$ and at least one $`\stackrel{˘}{l}_i`$ ($`i`$=1,2,3) are units, and the incidence relation is given by
$$\underset{i=1}{\overset{3}{}}\stackrel{˘}{l}_i\stackrel{˘}{x}_i\stackrel{˘}{l}_1\stackrel{˘}{x}_1+\stackrel{˘}{l}_2\stackrel{˘}{x}_2+\stackrel{˘}{l}_3\stackrel{˘}{x}_3=0.$$
(4)
From this definintion it follows that in $`PH(2,q)`$ — as in any ordinary projective plane — there is a perfect duality between points and lines; that is, instead viewing the points of the plane as the fundamental entities, and the lines as ranges (loci) of points, we may equally well take the lines as primary geometric constituents and define points in terms of lines, characterizing a point by the complete set of lines passing through it. It is also straightforward to see that this plane contains
$$\mathrm{\#}_{\mathrm{trip}}=\frac{\mathrm{\#}_\mathrm{t}^3\mathrm{\#}_\mathrm{z}^3}{\mathrm{\#}_\mathrm{u}}=\frac{(q^2)^3q^3}{q(q1)}=\frac{q^3(q^31)}{q(q1)}=q^2\left(q^2+q+1\right)$$
(5)
points/lines and that the number of points/lines incident with a given line/point is, in light of Eq. (4), equal to the number of non-equivalent couples $`(\varrho \stackrel{˘}{x}_1,\varrho \stackrel{˘}{x}_2)`$/$`(\sigma \stackrel{˘}{l}_1,\sigma \stackrel{˘}{l}_2)`$, i.e.
$$\mathrm{\#}_{\mathrm{coup}}=\frac{\mathrm{\#}_\mathrm{t}^2\mathrm{\#}_\mathrm{z}^2}{\mathrm{\#}_\mathrm{u}}=\frac{(q^2)^2q^2}{q(q1)}=\frac{q^2(q^21)}{q(q1)}=q\left(q+1\right).$$
(6)
These figures should be compared with those characterizing ordinary finite planes of order $`q`$, which read $`\overline{\mathrm{\#}}_{\mathrm{trip}}=q^2+q+1`$ and $`\overline{\mathrm{\#}}_{\mathrm{coup}}=q+1`$, respectively (e.g., ).
Any projective Hjelmslev plane, $`PH(2,q)`$ in particular, is endowed with a very important, and of crucial relevance when it comes to MUBs, property that has no analogue in an ordinary projective plane — the so-called neighbour (or, as occasionally referred to, non-remoteness) relation. Namely (see, e.g., ), we say that two points $`A`$ and $`B`$ are neighbour, and write $`AB`$, if either $`A=B`$, or $`AB`$ and there exist two different lines incident with both; otherwise, they are called nonneighbour, or remote. The same symbol and the dual definition is used for neighbour lines. Let us find the cardinality of the set of neighbours of a given point/line of $`PH(2,q)`$. Algebraically speaking, given a point $`\varrho \stackrel{˘}{x}_i`$, $`i`$=1,2,3, the points that are its neighbours must be of the form $`\varrho \left(\stackrel{˘}{x}_i+p\stackrel{˘}{y}_i\right)`$, with $`\stackrel{˘}{y}_i𝒯_r`$; for two points are neighbour iff their corresponding coordinates differ by a zero divisor . Although there are $`q^3`$ different choices for the triple $`(\stackrel{˘}{y}_1,\stackrel{˘}{y}_2,\stackrel{˘}{y}_3)`$, only $`q^3/q=q^2`$ of the classes $`\varrho \left(\stackrel{˘}{x}_i+p\stackrel{˘}{y}_i\right)`$ represent distinct points because $`\varrho \left(\stackrel{˘}{x}_i+p\stackrel{˘}{y}_i\right)`$ and $`\varrho \left(\stackrel{˘}{x}_i+p(\stackrel{˘}{y}_i+\kappa \stackrel{˘}{x}_i)\right)`$ represent one and the same point as $`\kappa `$ runs through all the $`q`$ elements of $`𝒯_r`$. Hence, every point/line of $`PH(2,q)`$ has $`q^2`$ neighbours, the point/line in question inclusive. Following the same line of reasoning, but restricting only to couples of coordinates, we find that given a point $`P`$ and a line $``$, $`P`$ incident with $``$, there exist exactly $`(q^2/q=)q`$ points on $``$ that are neighbour to $`P`$ and, dually, $`q`$ lines through $`P`$ that are neighbour to $``$.
Clearly, as $`AA`$ (reflexivity), $`AB`$ implies $`BA`$ (symmetry) and $`AB`$ and $`BC`$ implies $`AC`$ (transitivity), the neighbour relation is an equivalence relation. Given “$``$” and a point $`P`$/line $``$, we call the subset of all points $`Q`$/lines $`𝒦`$ of $`PH(2,q)`$ satisfying $`PQ`$/$`𝒦`$ the neighbour class of $`P`$/$``$. And since “$``$” is an equivalence relation, the aggregate of neighbour classes partitions the plane, i.e. the plane consists of a disjoint union of neighbour classes of points/lines. The modulo-$`p`$-mapping then “induces” a so-called canonical epimorphism of $`PH(2,q)`$ into $`PG(2,q)`$, the ordinary projective plane defined over $`GF(q)`$, with the neighbour classes being the cosets of this epimorphism . Loosely rephrased, $`PH(2,q)`$ comprises $`q^2+q+1`$ “clusters” of neighbour points/lines, each of cardinality $`q^2`$, such that its restriction modulo the neighbour relation is the ordinary projective plane $`PG(2,q)`$ every single point/line of which encompasses the whole “cluster” of these neighbour points/lines. Analogously, each line of $`PH(2,q)`$ consists of $`q+1`$ neighbour classes, each of cardinality $`q`$, such that its “$`^\overline{}`$” image is the ordinary projective line in the $`PG(2,q)`$ whose points are exactly these neighbour classes.
Let us illustrate these remarkable properties on the simplest possible example that is furnished by $`PH(2,q=2)`$, i.e. the plane defined over $`GR(4,1)`$ whose epimorphic “shadow” is the simplest projective plane $`PG(2,2)`$ — the Fano plane. As partially depicted in Fig. 1, this plane consists of seven classes of quadruples of neighbour points/lines, each point/line featuring three classes of couples of neighbour lines/points. When modulo-two-projected, each quadruple of neighbour points/lines goes into a single point/line of the associated Fano plane.
The most relevant geometrical object for our model is, of course, a conic, that is a curve $`𝒬`$ of $`PH(2,q)`$ whose points obey the equation
$$𝒬:\underset{ij}{}c_{ij}\stackrel{˘}{x}_i\stackrel{˘}{x}_jc_{11}\stackrel{˘}{x}_1^2+c_{22}\stackrel{˘}{x}_2^2+c_{33}\stackrel{˘}{x}_3^2+c_{12}\stackrel{˘}{x}_1\stackrel{˘}{x}_2+c_{13}\stackrel{˘}{x}_1\stackrel{˘}{x}_3+c_{23}\stackrel{˘}{x}_2\stackrel{˘}{x}_3=0,$$
(7)
with at least one of the $`c_{ij}{}_{}{}^{}s`$ being a unit of $`GR(p^2,r)`$. In particular, we are interested in a proper conic, which is a conic whose equation cannot be reduced into a form featuring fewer variables whatever coordinate transformation one would employ. It is known (see, e.g., ) that the equation of a proper conic of $`PH(2,q)`$ can always be brought into a “canonical” form
$$𝒬^{}:\stackrel{˘}{x}_1\stackrel{˘}{x}_3\stackrel{˘}{x}_2^2=0$$
(8)
from which it readily follows that any such conic is endowed, like a line, with $`q^2+q=q(q+1)`$ points; $`q^2`$ of them are of the form
$$\varrho \stackrel{˘}{x}_i=(1,\sigma ,\sigma ^2),$$
(9)
where the parameter $`\sigma `$ runs through all the elements of $`GR(p^2,r)`$, whilst the remaining $`q`$ are represented by
$$\varrho \stackrel{˘}{x}_i=(0,\delta ,1),$$
(10)
with $`\delta `$ running through all the zero-divisors of $`GR(p^2,r)`$. And each point of a proper conic, like that of a line, has $`q`$ neighbours; for the neighbours of a particular point $`\sigma =\sigma _0`$ of (9) are of the form
$$\varrho \stackrel{˘}{x}_i=(1,\sigma _0+p\kappa ,(\sigma _0+p\kappa )^2)=(1,\sigma _0+p\kappa ,\sigma _0^2+p2\kappa )$$
(11)
and there are obviously $`q`$ of them (the point in question inclusive) as $`\kappa `$ runs through $`𝒯_r`$, and all the $`q`$ points of (10) are the neighbours of any of them. All in all, a proper conic, like a line, of $`PH(2,q)`$ features $`q+1`$ pairwise disjoint classes of neighbour points, each having $`q`$ elements, these classes being the single points of its modular image in $`PG(2,q)`$. To illustrate the case, several proper conics in $`PH(2,2)`$ are shown in Fig. 2.
At this point our algebraic geometrical machinery is elaborate enough to generalize and qualitatively complete the geometrical picture of MUBs proposed in where we have argued that a basis of $`_q`$, $`q`$ given by (3), can be regarded as a point of an arc in $`PG(2,q)`$, with a complete set of MUBs corresponding to a proper conic (or, in the case of $`p`$=2, to a more general geometrical object called oval). This model, however, lacks a geometrical interpretation of the individual vectors of a basis, which can be achieved in our extended projective setting à la Hjelmslev only. Namely, taking any complete, i.e. of cardinality $`q`$+1, set of MUBs, its bases are now viewed as the neighbour classes of points of a proper conic of $`PH(2,q)`$ and the vectors of a given basis have their counterpart in the points of the corresponding neighbour class. The property of different vectors of a basis being pairwise orthogonal is then geometrically embodied in the fact that the corresponding points are all neighbour, whilst the property of two different bases being mutually unbiased answers to the fact that the points of any two neighbour classes are remote from each other. It is left to the reader as an easy exercise to check that “rephrasing these statements modulo $`p`$” one recovers all the conic-related properties of MUBs given in , irrespective of the value of $`p`$. The ($`p`$=2) case of “non-conic” MUBs is here, however, much more complex and intricate than that in the ordinary projective planes and will properly be dealt with in a separate paper.
To conclude, it must be stressed that this remarkable analogy between complete sets of MUBs and ovals/conics is worked out at the level of cardinalities only and thus still remains a conjecture. Hence, the next crucial step to be done is to construct an expliciting mapping by associating a MUB to each neighbour class of the points of the conic and a state vector of this MUB to a particular point of the class. This is a much more delicate issue, as there are (at least) two non-isomorphic kinds of projective Hjelmslev planes of order $`q=p^r`$ that have exactly the same “cardinality” properties, viz. the plane defined over the Galois ring $`GR(p^2,r)`$ and the one defined over the ring of “dual” numbers, $`GF(q)[x]/(x^2)GF(q)+eGF(q)`$, where $`e^2=0`$. Even for the simplest case ($`p`$=2 and $`r`$=1) there is an intricate difference in geometry between the two planes, as the former contains ($`q^2+q+1=`$)7-arcs, while the latter not (see, e.g., ). A thorough exploration of the fine structure of these two Hjelmslev geometries, as well as of a number of other finite Hjelmslev and related ring planes, is therefore a principal theoretical task for making further progress in this direction.
Acknowledgement
The first author is grateful to Prof. Olav Arnfinn Laudal for a number of enlightening comments concerning the structure of projective geometries over (Galois) rings and to Prof. Mark Stuckey for being a great help to him in obtaining a copy of . We also thank two of the referees for their constructive comments and Mr. Pavol Bendík for careful drawing of the figures. |
warning/0506/physics0506171.html | ar5iv | text | # How to observe ⁸B solar neutrinos in liquid scintillator detectors
## 1 Introduction
Observations of solar neutrinos have offered the first experimental evidence in favor of non-standard effects, in particular neutrino flavor transitions induced by non-zero neutrino masses and mixings. Solar neutrinos have been detected by radiochemical experiments (i. e., Homestake , Gallex/GNO and SAGE ) which give an energy-integrated information on the solar neutrino fluxes, and by real time water Cherenkov detectors (i. e., Kamiokande, Super-Kamiokande and SNO ) which allow to observe the spectral distribution of solar neutrino events. However, the detection threshold in Cherenkov detectors is limited to about $`5`$ MeV by the radiopurity of the target mass and, as a consequence, only high energy $`{}_{}{}^{8}\mathrm{B}`$ solar neutrinos spectrum has been measured.
In the next future, liquid organic scintillator detectors, such as KamLAND and Borexino , will be operating with the goal of measuring the low energy solar neutrino fluxes, in particular $`{}_{}{}^{7}\mathrm{Be}`$, CNO and pep solar neutrinos. The KamLAND experiment is a 1 kT detector located in the Kamioka mine (Japan), at a depth of 2700 m.w.e. of rock, operating continuously since January 2002, with the main goal of measuring the flux of the $`\overline{\nu }_\mathrm{e}`$’s coming from all the Japanese nuclear power plants. This experiment has spectacularly confirmed the so-called Large Mixing Angle (LMA) solution to the solar neutrino problem (see, e. g., for a recent reanalysis). Borexino is a $`0.3`$ kT liquid scintillator detector which is being commissioned at Gran Sasso (Italy), under 3800 m.w.e. of rock, whose main goal is the measurement of the $`{}_{}{}^{7}\mathrm{Be}`$ solar neutrino flux. Moreover, it has been recently proposed to realize a $`1`$ kT liquid scintillator detector, denominated SNO+, at the SNO site (SNOLab, Canada) under 6000 m.w.e. of rock, after the completion of the SNO detector physics program . It is also under discussion the possibility to realize a gigantic ($`30\mathrm{k}\mathrm{T}`$) liquid scintillator detector, the Low Energy Neutrino Astrophysics (LENA) detector , in the Pyhäsalmi mine (Finland) at a depth of 1450 m ($`4000`$ m.w.e.), although other sites (e. g., underwater in the site of Pylos in Greece) have also been proposed. The observation of solar neutrinos in these detectors, through $`\nu `$e elastic scattering, is not a simple task, since neutrino events cannot be separated from the background, and it can be accomplished only if the detectors contamination will be kept very low . Moreover, only mono-energetic sources such as $`{}_{}{}^{7}\mathrm{Be}`$ or pep neutrinos can be detected, taking advantage of the Compton-like shoulder edge produced in the event spectrum.
In this Letter, we show that organic liquid scintillator detectors can also measure the $`{}_{}{}^{8}\mathrm{B}`$ solar neutrino flux by means of the $`\nu _\mathrm{e}`$ charged current interaction with the $`{}_{}{}^{13}\mathrm{C}`$ nuclei naturally contained in the scintillators. The possibility to use $`{}_{}{}^{13}\mathrm{C}`$ as a target for $`{}_{}{}^{8}\mathrm{B}`$ neutrinos was pointed out in the past by . Here, we propose a technique to tag the solar neutrino events. Namely, we propose to identify the signal by looking at the time and space coincidence with the decay of the produced $`{}_{}{}^{13}\mathrm{N}`$ nuclei. We perform a detailed calculation of the solar neutrino signal and of the background in KamLAND, Borexino and SNO+, showing that these detectors will be able to extract the signal with a reasonable uncertainty in a few years of data taking. It should be stressed that the proposed technique does not involve any modification of the experimental setup, since one expects a background-to-signal ratio of the order of 1 or less even assuming the natural isotopic abundance of $`{}_{}{}^{13}\mathrm{C}`$ ($`1\%`$) and the contamination levels already reached in the KamLAND detector .
The Letter is organized as follows. In the next section we discuss the neutrino interactions with $`{}_{}{}^{13}\mathrm{C}`$. In Sec. III we calculate the solar neutrino event rates. In Sec. IV we analyze the space and time coincidence with the decay of the produced $`{}_{}{}^{13}\mathrm{N}`$ nuclei. Background issues are described in Sec. V. Sec. VI presents the expected sensitivity for KamLAND, Borexino and SNO+ and prospects for LENA. In Sec. VII we draw our conclusions.
## 2 Neutrino interactions on $`{}_{}{}^{13}\mathrm{C}`$
The $`{}_{}{}^{13}\mathrm{C}`$ is a stable isotope of carbon with a natural isotopic abundance $`I=1.07`$%. A small amount of $`{}_{}{}^{13}\mathrm{C}`$ is, thus, naturally present in organic liquid scintillators and can be used as a target for neutrino detection. The relevant detection process in our discussion is the charged current (CC) transition to $`{}_{}{}^{13}\mathrm{N}`$ ground state:
$$\nu _\mathrm{e}+^{13}\mathrm{C}^{13}\mathrm{N}(\mathrm{gnd})+\mathrm{e}^{}.$$
(1)
The reaction threshold is $`Q=2.22`$ MeV and, thus, only $`{}_{}{}^{8}\mathrm{B}`$ solar neutrinos are detectable (with a neglible contribution from hep neutrinos). In liquid scintillators one observes the electron produced in the final state with a visible energy which, neglecting the detector energy resolution, is simply equal to the electron kinetic energy $`T_\mathrm{e}`$. The cross section of reaction (1) is known with great accuracy, since it can be deduced from the decay time of $`{}_{}{}^{13}\mathrm{N}`$. One has :
$`\sigma (E_\nu )`$ $`=`$ $`{\displaystyle \frac{2\pi ^2\mathrm{ln}2}{m_\mathrm{e}^5ft}}p_\mathrm{e}E_\mathrm{e}F(Z,E_\mathrm{e})`$ (2)
$`=`$ $`0.2167\times 10^{43}\mathrm{cm}^2{\displaystyle \frac{p_\mathrm{e}E_\mathrm{e}}{\mathrm{MeV}^2}}F(Z,E_\mathrm{e}),`$
where $`E_\mathrm{e}=E_\nu Q+m_\mathrm{e}`$ is the electron energy,<sup>1</sup><sup>1</sup>1We neglected the small recoil energy of the $`{}_{}{}^{13}\mathrm{N}`$ nucleus (of the order of few keV). In this assumption, one has simply $`T_\mathrm{e}=E_\mathrm{e}m_\mathrm{e}=E_\nu Q`$. $`p_\mathrm{e}`$ is the electron momentum, $`F(Z=7,E_\mathrm{e})`$ is the Fermi factor and the $`ft`$value of $`{}_{}{}^{13}\mathrm{N}`$ decay is experimentally determined as $`\mathrm{log}(ft/\mathrm{s})^{\mathrm{exp}}=3.667\pm 0.001`$ . By averaging the cross section over the $`{}_{}{}^{8}\mathrm{B}`$ neutrino spectrum, one obtains $`\sigma =8.57\times 10^{43}\mathrm{cm}^2`$, which is about one order of magnitude larger than the cross section of $`\nu _\mathrm{e}\mathrm{e}\nu _\mathrm{e}\mathrm{e}`$ scattering.
The peculiarity of process (1) is that it can be monitored by looking for the delayed coincidence with the positron emitted in the $`{}_{}{}^{13}\mathrm{N}`$ decay:
$$^{13}\mathrm{N}^{13}\mathrm{C}+\nu _\mathrm{e}+\mathrm{e}^+,$$
(3)
which occurs with $`99.8`$% branching ratio (0.2% of $`{}_{}{}^{13}\mathrm{N}`$ nuclei undergo electron capture) and a decay time $`\tau =862.6`$ s. In this case, the visible energy is the sum of the positron kinetic energy and the energy released in $`\mathrm{e}^+\mathrm{e}^{}`$ annihilation, so that the delayed events have a continuous energy spectrum in the range $`[1.02,2.22]`$ MeV. Moreover, in the absence of macroscopic motions in the detector, the $`{}_{}{}^{13}\mathrm{N}`$ nucleus essentially does not move from its original position. The expected displacement due to recoil and diffusion<sup>2</sup><sup>2</sup>2Typical value of diffusion coefficients in liquids are $`D10^5`$ cm<sup>2</sup>s<sup>-1</sup>, which correspond to an average displacement in the time $`\tau =862.6`$ s equal to $`l_{\mathrm{diff}}=\sqrt{2D\tau }0.1`$ cm. during the decay time $`\tau `$ is, indeed, smaller than the typical detector spatial resolution, $`\sigma 10`$ cm. This means that the prompt event produced by the reaction (1) and the delayed event produced by the decay (3) have to be observed essentially in the same position. This condition, as we will see in the following sections, is extremely effective in reducing the background.
Other interaction channels of low energy neutrinos with $`{}_{}{}^{13}\mathrm{C}`$ can, in principle, be considered. First, we discuss the CC transition to $`{}_{}{}^{13}\mathrm{N}`$ excited states. For solar neutrinos, only the lowest excited state (at 3.51 MeV) could be of practical importance. The cross section for this process is about 30% of that for the ground state and it is calculated theoretically with an uncertainty at the level of $`3040`$. However, the $`{}_{}{}^{13}\mathrm{N}`$$`{}_{}{}^{}(3.51\mathrm{MeV})`$ decays almost immediately to $`{}_{}{}^{12}\mathrm{C}+`$p with almost 100% branching ratio . As a consequence, it cannot be discriminated by the coincidence with the delayed events (3). The other relevant process is the neutral current (NC) transition:
$$\nu _\mathrm{x}+^{13}\mathrm{C}\nu _\mathrm{x}+^{13}\mathrm{C}^{}.$$
(4)
Here, only the excited state $`{}_{}{}^{13}\mathrm{C}`$$`{}_{}{}^{}(3.68\mathrm{MeV})`$ is relevant and the excitation to other levels has negligible cross section. The cross section for this process, averaged over the $`{}_{}{}^{8}\mathrm{B}`$ neutrino spectrum, is $`\sigma _{\mathrm{NC}}=1.15\times 10^{43}\mathrm{cm}^2`$ and it is affected by about $`3040\%`$ uncertainty. The process is in principle very interesting since it can give a measure of the total $`{}_{}{}^{8}\mathrm{B}`$ flux and can be also tagged by a monochromatic $`\gamma `$-ray emission back to the ground state. However, the great uncertainty and the low cross section make this process hard to be competitive with the NC measurement made by SNO.
In this Letter, we will mainly focus on the information which can be obtained from reaction (1), which, at present, seems more interesting in view of the larger and much better known cross section, and of the delayed detection tagging with the $`{}_{}{}^{13}\mathrm{N}`$ decay.
## 3 The $`{}_{}{}^{8}\mathrm{B}`$ solar neutrino signal
A neutrino of energy $`E_\nu `$ interacting with $`{}_{}{}^{13}\mathrm{C}`$ through reaction (1) produces an electron with kinetic energy $`T_\mathrm{e}=E_\nu Q`$ (neglecting the $`{}_{}{}^{13}\mathrm{N}`$ recoil). The rate $`R_\nu `$ of prompt events (per unit mass) produced by $`{}_{}{}^{8}\mathrm{B}`$ solar neutrinos in the energy window $`[T_{e,1},T_{e,2}]`$ is thus simply given by:
$$R_\nu =n_{{}_{}{}^{13}\mathrm{C}}\mathrm{\Phi }_{T_{e,1}+Q}^{T_{e,2}+Q}dE_\nu \sigma (E_\nu )\lambda (E_\nu )P_{\mathrm{ee}}(E_\nu ),$$
(5)
where $`\mathrm{\Phi }=5.79\times 10^6`$ cm<sup>-2</sup>s<sup>-1</sup> is the boron neutrino flux , $`\lambda (E_\nu )`$ <sup>3</sup><sup>3</sup>3A useful approximation for the $`{}_{}{}^{8}\mathrm{B}`$ spectrum is the following:
$$\lambda (E_\nu )=\frac{1}{E_0}\frac{\mathrm{\Gamma }(\alpha +\beta )}{\mathrm{\Gamma }(\alpha )\mathrm{\Gamma }(\beta )}x^{\alpha 1}(1x)^{\beta 1},$$
with $`x=E/E_0`$, $`\alpha =2.92`$, $`\beta =3.49`$, and $`E_0=14.8`$MeV. is the boron neutrino spectrum , $`\sigma (E_\nu )`$ is the interaction cross section and $`P_{\mathrm{ee}}(E_\nu )`$ is the electron neutrino survival probability. In the previous relation, $`n_{{}_{}{}^{13}\mathrm{C}}`$ is the number of $`{}_{}{}^{13}\mathrm{C}`$ atoms per unit mass, which depends on the scintillator chemical composition:
$$n_{{}_{}{}^{13}\mathrm{C}}=\frac{I}{u}\underset{k}{}f_k\frac{X_k}{\mu _k},$$
(6)
where $`I=1.07\times 10^2`$ is the isotopic abundance of $`{}_{}{}^{13}\mathrm{C}`$, $`u=1.661\times 10^{33}`$ kT is the atomic mass unit, $`f_k`$ is the mass fraction of the $`k`$-th component into the scintillator, $`X_k`$ the stoichiometric coefficient of carbon in the molecule, and $`\mu _k`$ is the molecular mass of the $`k`$-th molecule. KamLAND scintillator is composed by 80% of dodecane (C<sub>12</sub>H<sub>26</sub>, with a molecular mass $`\mu =170.33`$) and 20% of pseudocumene (C<sub>9</sub>H<sub>12</sub>, with a molecular mass $`\mu =120.19`$), which correspond to $`n_{{}_{}{}^{13}\mathrm{C}}=4.60\times 10^{29}`$ kT<sup>-1</sup>. Borexino is, instead, composed by 100% of pseudocumene, corresponding to $`n_{{}_{}{}^{13}\mathrm{C}}=4.82\times 10^{29}`$ kT<sup>-1</sup>. The SNO+ liquid scintillator composition has still to be decided. For simplicity, we will assume, here and in the following, that it will be the same as in Borexino (i. e., 100% pseudocumene).
In Fig. 1 we show the function $`\varrho (E_\nu )P_{\mathrm{ee}}(E_\nu )\lambda (E_\nu )\sigma (E_\nu )`$ (normalized to unity) which gives the relative contribution of neutrinos of different energies to the total signal from reaction (1). The solid line is obtained in the assumption of an undistorted $`{}_{}{}^{8}\mathrm{B}`$ neutrino spectrum (which can be intended as the non oscillatory scenario or a constant suppression of $`\nu _\mathrm{e}`$). The dashed line is obtained in the assumption of $`\nu _\mathrm{e}\nu _{\mu \tau }`$ flavor transitions for the following oscillation parameters which are the current best fit (in $`2\nu `$) for the whole solar and KamLAND data :
$`\delta m^2`$ $`=`$ $`7.92\times 10^5\mathrm{eV}^2,`$
$`\mathrm{sin}^22\theta `$ $`=`$ $`0.86.`$ (7)
The electron neutrino survival probability has been calculated taking into account the Mikheyev-Smirnov-Wolfenstein (MSW) effect in the Sun (for simplicity we have not considered the oscillations in the Earth matter).<sup>4</sup><sup>4</sup>4A simple and accurate approximation to calculate the adiabatic MSW survival probability can be found in . Neglecting the detector energy resolution, the function $`f(T_\mathrm{e})\varrho (T_\mathrm{e}+Q)`$ also gives the spectral distribution of solar neutrino events, since detection reaction kinematics implies a one-to-one relation (i. e., $`T_\mathrm{e}=E_\nu Q`$) between the electron and neutrino energies. The event spectrum is, in principle, extremely sensitive to a possible deformation of parent solar neutrino spectrum. In particular, the differences between the two curves in Fig. 1 directly reflect the behavior of the electron neutrino survival probability in LMA scenarios. Namely, the rise of the LMA spectrum (dashed line) with respect to the standard case (solid line) below $`E_\nu 7`$ MeV ($`T_\mathrm{e}5`$ MeV in terms of electron energy) is due to the transition from vacuum averaged neutrino oscillations at small energies to purely adiabatic transitions at large energies. The observation of this feature would be very important as a final confirmation of matter effect in solar neutrino oscillations. However, it will be extremely hard to observe it in the present detectors, due to the smallness of the expected event rates.
In first two columns of Tab. 1, we show the neutrino event rates (given in counts$``$kTy<sup>-1</sup>) expected in KamLAND and Borexino scintillators in the energy windows $`[T_{\mathrm{e},1},T_{\mathrm{e},2}]=[2.8,16]`$ MeV and $`[2.8,5.5]`$ MeV, assuming the oscillation parameters in Eq. (3). The lower bound ($`2.8`$ MeV) has been chosen to reduce the background from $`\mathrm{U}\mathrm{Th}`$ contamination. The upper bound $`5.5`$ MeV has been chosen to focus on the low energy part of the spectrum, which, as explained above, is particularly interesting and, moreover, is practically unexplored by Super-Kamiokande and SNO. As we can see, the counting rates are of the order of $`1020`$ counts$``$kTy<sup>-1</sup>. Moreover, they will be further reduced by the cuts, essential to reduce the background. However, as we will see in the next sections, the background levels are extremely low so that it will be possible to extract the solar neutrino signal with a reasonable uncertainty. We remark that the proposed measure does not require any modification of the present experimental set-up and, although with a larger uncertainty, can be complementary to those coming from SNO and Super-Kamiokande.
## 4 Tagging the events
In order to reduce the background, we can take advantage of the time and space coincidence of neutrino events with the positron emitted in the $`{}_{}{}^{13}\mathrm{N}`$ decay. A candidate prompt event will be tagged as signal only if followed by a delayed event in the energy window $`[1.02,2.22]`$ MeV, within a time interval $`\mathrm{\Delta }t=𝒯\tau `$ (where $`\tau =862.6`$ s is the $`{}_{}{}^{13}\mathrm{N}`$ decay time), and inside a sphere of radius $`r=\sigma `$ from the prompt event detection point (where $`\sigma 10`$ cm is the typical detector spatial resolution, see, e. g., ). The signal event rate is thus given by:
$$S=R_\nu ϵ(𝒯,),$$
(8)
where the global efficiency of the coincidence, $`ϵ(𝒯,)`$, is determined by the combined efficiency for the cut in space, $`\xi ()`$, and in time, $`\eta (𝒯)`$:
$$ϵ(𝒯,)=\xi ()\eta (𝒯).$$
(9)
The function $`\eta (𝒯)`$ is simply equal to the probability that the $`{}_{}{}^{13}\mathrm{N}`$ nucleus decays within the time $`\mathrm{\Delta }t=𝒯\tau `$:
$$\eta (𝒯)=1\mathrm{exp}(𝒯).$$
(10)
The function $`\xi ()`$ is instead the probability that the prompt event and the delayed event, which are assumed to occur in the same position, are detected at a distance smaller than $`r=\sigma `$. Modelling the detector spatial resolution with a gaussian function, one obtains:<sup>5</sup><sup>5</sup>5If we model the detector spatial resolution with a gaussian function $`f(𝐱,𝐱_0)\mathrm{exp}[(𝐱𝐱_0)^2/2\sigma ^2]`$, where $`𝐱_0`$ is the true position of the event and $`𝐱`$ is the observed position of the event, the probability that the prompt event is observed in the position $`𝐱_\mathrm{p}`$ and the delayed in the position $`𝐱_\mathrm{d}`$ can be cast as $`\mathrm{d}P(𝐱_\mathrm{p},𝐱_\mathrm{d})=f(𝐱_\mathrm{p},𝐱_0)f(𝐱_\mathrm{d},𝐱_0)\mathrm{d}^3𝐱_\mathrm{p}\mathrm{d}^3𝐱_\mathrm{d}\mathrm{exp}\left[\left(𝐫+𝐲\right)^2/4\sigma ^2\right]\mathrm{d}^3𝐫\mathrm{d}^3𝐲`$, where $`𝐫=𝐱_\mathrm{p}𝐱_\mathrm{d}`$ and $`𝐲=𝐱_\mathrm{p}+𝐱_\mathrm{d}2𝐱_0`$. Actually, only the distance $`r=|𝐫|`$ is observable, so that, integrating $`𝐫`$ on a sphere of radius $`\sigma `$ and $`𝐲`$ on $`𝐑^3`$, one obtains Eq. (11).
$$\xi ()=\frac{_0^{}dxx^2\mathrm{exp}(x^2/4)}{_0^{\mathrm{}}dxx^2\mathrm{exp}(x^2/4)}=\mathrm{erf}\left(\frac{}{2}\right)\sqrt{\frac{1}{\pi }}\mathrm{exp}\left(\frac{^2}{4}\right).$$
(11)
We remark, that the above equation is valid in the assumption that the displacement between the point where $`{}_{}{}^{13}\mathrm{N}`$ is created and the point where it decays is small with respect to $`\sigma `$. For that to happen, the macroscopic motions in the liquid scintillator have to be sufficiently slow. This can be achieved, for example, by maintaining a small temperature gradient pointing upward everywhere in the detector. KamLAND data show that the measured average displacement of the diffusive $`{}_{}{}^{222}\mathrm{Rn}`$ over its $`5.5`$ d mean life is less than 1 m . Therefore, the assumption that $`{}_{}{}^{13}\mathrm{N}`$ nuclides displacement over their $`15`$ minutes lifetime can be kept smaller than detector resolution seems justified.
Finally, we consider the possibility that the delayed energy window is reduced with respect to the full energy range ($`[1.02,2.22]`$ MeV) of $`{}_{}{}^{13}\mathrm{N}`$ decay spectrum. In this case the signal event rate is given by:
$$S=R_\nu ϵ(𝒯,)(E_1,E_2)=R_\nu ϵ(𝒯,)_{E_1}^{E_2}dE_\mathrm{d}\chi (E_\mathrm{d}),$$
(12)
where $`(E_1,E_2)`$ is the fraction of decay events in the adopted energy window $`[E_1,E_2]`$ \[$`\chi (E)`$ is the normalized $`{}_{}{}^{13}\mathrm{N}`$ decay spectrum\].
In last two columns of Tab. 1, we show the signal event rates (given in counts$``$kTy<sup>-1</sup>) expected in the KamLAND and Borexino liquid scintillators, after that the efficiency cuts are applied. Here, for illustrative purposes, we consider a time cut at $`𝒯=2`$ and a space cut at $`=3`$, which correspond to a global efficiency $`ϵ(,𝒯)=0.68`$. It is clear, however, that the cuts must be optimized according to each detector’s capabilities and performances. For the KamLAND detector, moreover, we restrict the delayed energy window to $`[E_1,E_2]=[1.3,2.22]`$ MeV in order to reduce the background from $`{}_{}{}^{210}\mathrm{Bi}`$ originated by the decay of $`{}_{}{}^{210}\mathrm{Pb}`$ that can be either produced by build-up due to $`{}_{}{}^{222}\mathrm{Rn}`$ contamination in the liquid scintillator or caused by an intrinsic impurity.<sup>6</sup><sup>6</sup>6At present $`{}_{}{}^{210}\mathrm{Pb}`$ is an important contamination in KamLAND which will be removed by purification to allow solar neutrinos detection. This reduces further the efficiency by a factor $`0.77`$. For Borexino (and SNO+), we consider instead the full range $`[1.02,2.22]`$ MeV, assuming that the $`{}_{}{}^{210}\mathrm{Bi}`$ background contribution will be further reduced, since this is a pre-requisite for the observation of sub-MeV solar neutrinos.<sup>7</sup><sup>7</sup>7If the detector has an intrinsic efficiency $`ϵ_\mathrm{p}`$ ($`ϵ_\mathrm{d}`$) for the prompt (delayed) window, the global efficiency is further reduced by a factor $`ϵ_\mathrm{p}ϵ_\mathrm{d}`$. We assume for simplicity $`ϵ_\mathrm{p}=ϵ_\mathrm{d}=1`$.
As a final result, the expected signal event rates are at the level of $`1020`$ counts$``$kTy<sup>-1</sup>. In order to observe such low counting rates, one clearly needs detectors with sufficiently low background levels (and, of course, efficient background rejection). Present detectors, as we shall see in the next section, already satisfy this requirement.
## 5 The background
There are three main sources of background for the proposed measure. These are:
1. Internal background due to $`\mathrm{U}\mathrm{Th}`$ contamination and to contamination from long lived radon daughters out of secular equilibrium with $`{}_{}{}^{238}\mathrm{U}`$ (in particular $`{}_{}{}^{210}\mathrm{Pb}`$);
2. Cosmogenic background due to muon-induced production of radioactive nuclides, such as $`{}_{}{}^{11}\mathrm{C}`$, $`{}_{}{}^{10}\mathrm{C}`$ , etc.;
3. Elastic $`\nu `$e scattering by solar neutrinos.
These background sources are well known, so that it is possible to perform a detailed analysis of their relevance. This is clearly important, because it allows to make a realistic estimate. We remark, however, that the real background level will be measured directly by the experiments with great accuracy, being the background event rate much larger than the signal event rate both in the prompt and delayed energy window (before space and time cuts are applied).
In Tab. 2, we show the contribution of each background source to the total background rates (per unit mass) in KamLAND, Borexino and SNO+. We have assumed that the $`\mathrm{U}\mathrm{Th}`$ contamination in the detectors is at the $`10^{17}`$ g/g level, which correspond to the present contamination in the KamLAND detector . The internal background due to the elastic scattering of solar neutrinos on electrons has been evaluated assuming $`\nu _\mathrm{e}\nu _{\mu \tau }`$ flavor oscillations with the oscillation parameters given in Eq. (3). Finally, the cosmogenic contribution has been obtained by rescaling the results of , which are relative to Borexino, also to KamLAND and SNO+. This can be done by considering that the muon induced background $`R_\mu `$ in a given experiment scales as $`R_\mu n_{{}_{}{}^{12}\mathrm{C}}\mathrm{\Phi }_\mu E_\mu ^\alpha `$, where $`\mathrm{\Phi }_\mu `$ is the muon flux at the experimental site, $`E_\mu `$ is the average muon energy, $`n_{{}_{}{}^{12}\mathrm{C}}`$ is the number of $`{}_{}{}^{12}\mathrm{C}`$ nuclei per unit mass in the scintillator (the $`{}_{}{}^{12}\mathrm{C}`$ is the most relevant target for muon induced radioactive nuclei production in liquid organic scintillators) and $`\alpha =0.73`$.<sup>8</sup><sup>8</sup>8The cosmogenic production cross section scales with the energy as $`\sigma E_\mu ^\alpha `$ with $`\alpha 0.73`$ . From the data in Tab. 3 (see for details) we calculate that the cosmogenic background is $`7`$ times larger in KamLAND than in Borexino, while is $`94`$ times lower in SNO+.
From Tab. 2 we see that cosmogenic background is the dominant component for KamLAND and Borexino both in the prompt and delayed window, while it give only a minor contribution to the total background in SNO+. It is interesting to give a closer look at the various cosmogenic background sources. The relevant cosmogenic nuclei and their half-lives are reported in Tab. 4. We see that only $`{}_{}{}^{11}\mathrm{C}`$ and $`{}_{}{}^{10}\mathrm{C}`$ nuclei have a long lifetime. The background coming from the other cosmogenics can be efficiently rejected simply extending the muonic veto to few seconds after the muon passage through the detector. However, in order to avoid further losses of efficiency due to the veto dead-time, we rejected only the background events coming from the (fast) decay of $`{}_{}{}^{12}\mathrm{B}`$. In the last four rows of Tab. 4 we show the event rates (given in counts$``$kTy<sup>-1</sup>) expected in Borexino for the four energy windows considered in this work. We see that the main contribution to cosmogenic background in the prompt energy windows come from $`{}_{}{}^{10}\mathrm{C}`$, while the dominant source in delayed windows is provided by $`{}_{}{}^{11}\mathrm{C}`$ nuclides. It was recently shown that $`{}_{}{}^{11}\mathrm{C}`$-induced background can be greatly reduced by a three fold coincidence with the parent muon track and the subsequent neutron capture on protons. However, in order to be extremely conservative, we have not considered this possibility.<sup>9</sup><sup>9</sup>9An additional background not included above is provided by the interaction of cosmogenic protons with $`{}_{}{}^{13}\mathrm{C}`$, according to $`{}_{}{}^{13}\mathrm{C}(`$p,n$`)^{13}\mathrm{N}`$. These events are potentially dangerous because, being followed by the decay of the produced $`{}_{}{}^{13}\mathrm{N}`$, they cannot be discriminated by the coincidence (a rejection is, anyhow, possible by looking at the subsequent neutron capture on protons). We have estimated, by MonteCarlo simulations, that this background component is negligible, being the cosmogenic proton interaction rate with $`{}_{}{}^{13}\mathrm{C}`$ of the order $`10^2`$ kTy<sup>-1</sup> in Borexino.
In the last two columns of Tab. 2, we show the background in the prompt energy windows after the coincidence criteria are applied. This is obtained by considering that the probability to have a background event in the delayed energy window when space and time cuts are applied is simply equal to the average number of delayed background events during the time interval $`\mathrm{\Delta }t=𝒯\tau `$ and inside the spherical volume $`V=(4/3)\pi (\sigma )^3`$, being this number much smaller than one.<sup>10</sup><sup>10</sup>10The probability to have at least one background event inside a certain (spatial and/or temporal and/or energy) window should be calculated by means of a Poissonian distribution. However, since the number of background events is very small, the probability is practically equal to the average number of events in the considered window. The rate of fake coincidences $`B`$ is thus given by:
$$B=(B_\mathrm{p}B_\mathrm{d}\rho )\left[\frac{4}{3}\pi (\sigma )^3𝒯\tau \right],$$
(13)
where $`B_\mathrm{p}`$ and $`B_\mathrm{d}`$ are the prompt and delayed background rates (per unit mass), $`\rho `$ is the liquid scintillator density (equal to $`\rho =0.78\mathrm{g}/\mathrm{cm}^3`$ for the KamLAND scintillator and $`\rho =0.88\mathrm{g}/\mathrm{cm}^3`$ for the Borexino scintillator) and we have taken $`=3`$ and $`𝒯=2`$. We see that, despite of the large number of background events (several thousands per kTy) in the prompt and delayed energy window, the fake coincidences are rare (tenth per kTy in KamLAND, few per kTy in Borexino or almost absent in SNO+), and comparable or lower than the expected signal. For this reason we think that a measure of the $`{}_{}{}^{8}\mathrm{B}`$ solar neutrino flux is feasible.
## 6 Expected sensitivity and future prospects
Since real and fake coincidences are indistinguishable, the number of signal events $`N_S`$ has to be obtained from the difference between the total number of observed events $`N_T`$ and the number of background events $`N_B`$:
$$N_SS=N_TN_B=N_TB,$$
(14)
where $`S`$ is the signal event rate, $`B`$ is the background event rate and $``$ is the total detector exposure. The uncertainty of the number of signal events $`\mathrm{\Delta }N_S`$ is obtained propagating the error in Eq. (14):
$$\mathrm{\Delta }N_S=\sqrt{N_T+(\mathrm{\Delta }B)^2},$$
(15)
where we assumed that the total number of events is affected by a Poissonian uncertainty $`\mathrm{\Delta }N_T=\sqrt{N_T}`$ and we combined in quadrature the errors. Dividing by $`N_S`$ we obtain the fractional uncertainty $`\delta S`$ of the signal rate:
$$\delta S\frac{\mathrm{\Delta }S}{S}=\sqrt{\frac{1+r}{S}+r^2\delta B^2},$$
(16)
where $`r=B/S`$ is background-to-signal ratio and $`\delta B`$ is the fractional uncertainty of the background rate.
As anticipated in the previous section, the background will be directly measured by the experiments. More precisely, one measures the prompt and the delayed background rates and, then, determines the final background rate $`B`$ through Eq. (13). The uncertainty $`\delta B`$ can, thus, be estimated as:
$$\delta B=\sqrt{\delta B_\mathrm{p}^2+\delta B_\mathrm{d}^2}=\frac{1}{^{1/2}}\sqrt{\frac{1}{B_\mathrm{p}}+\frac{1}{B_\mathrm{d}}}$$
(17)
where we considered that the fractional uncertainties of the prompt and delayed background rates are given by $`\delta B_{\mathrm{p},\mathrm{d}}=[B_{\mathrm{p},\mathrm{d}}]^{1/2}`$,as prescribed by Poissonian statistics.By using Eq. (17), Eq. (16) can be cast as:
$$\delta S=\frac{1}{\sqrt{S}}\sqrt{(1+r)+\left[\frac{S}{B_\mathrm{p}}+\frac{S}{B_\mathrm{d}}\right]r^2}.$$
(18)
It is immediately evident that the second term in the square root of the above expression is always negligible, being the ratios $`r^2S/B_{\mathrm{p},\mathrm{d}}`$ of the order of $`10^3`$ or less in the various detectors. This means that the contribution of the background uncertainty to the total error budget is always negligible:
$$\delta S\sqrt{\frac{1+r}{S}}.$$
(19)
This situation is, in principle, favorable. The fractional uncertainty $`\delta S`$ scales as $`^{1/2}`$, and, thus, it is possible to obtain a good sensitivity if the detector exposure is large enough. In the following we assume $`=1`$ kTy.
In the first two columns of Tab. 5 we show the background-to-signal ratio $`r`$ for the three experiments under study in the two energy windows $`[2.8,16]`$ MeV and $`[2.8,5.5]`$ MeV, assuming a space cut $`=3`$ and a time cut $`𝒯=2`$. In the third and fourth columns, we show the corresponding sensitivity $`\delta S`$, calculated according to Eq. (19), assuming a total exposure $`=1`$ kTy. In the last two columns, we give the minimal values for $`\delta S`$ obtained by choosing the optimal values of $``$ and $`𝒯`$ which minimize the quantity $`(1+r)/S`$ in each experiment.<sup>11</sup><sup>11</sup>11This exercise have to be done numerically, since $``$ and $`𝒯`$ enter in Eq. (19) in a non-trivial way. Clearly, the lower is the background level in the detector, the larger is the space-time window which has to be considered. In KamLAND the optimal sensitivity is obtained with $`=2.82`$ and $`𝒯=1.65`$ in the energy window $`[2.8,16]`$ MeV and $`=2.64`$ and $`𝒯=1.44`$ in the energy window $`[2.8,5.5]`$ MeV. Larger cuts are preferable for Borexino and SNO+ (of the order of $`5`$ and $`𝒯5`$). However, the choice of cuts is not crucial these (less noisy) detectors and the sensitivity does not change dramatically if we take tighter cuts.
In KamLAND, due to the large cosmogenic contribution, the background-to-signal ratio is equal to about 2.5 in the energy window $`[2.8,16]`$ MeV (while it is equal to about 7 if we restrict to $`[2.8,5.5]`$ MeV). This corresponds to an expected sensitivity $`\delta S`$ equal to about 50% in one year of data taking (assuming $``$1 kT fiducial mass). We remind that KamLAND has already analyzed data for a total exposure $`=0.766`$ kTy , corresponding to $`18`$ solar neutrino events in the window $`[2.8,16]`$ MeV which can be extracted with about 60% uncertainty.<sup>12</sup><sup>12</sup>12Of course, with the proposed cuts, only about 9 of these events would be tagged as candidate. Despite the large uncertainty, we believe that this measure would, anyhow, represent a milestone, since it would be the first observation of solar neutrinos into liquid scintillator detectors. We hope that the KamLAND collaboration will try to extract this piece of information from their own set of data.
In Borexino and in SNO+, due to the larger depth of the experimental sites, the background-to-signal ratio is much less than one. The sensitivity is thus only limited by the statistical error of the signal events. The low background level allows to explore with sufficient accuracy the energy windows $`[2.8,5.5]`$ MeV for which, at present, we have no direct information. This also indicates that, in these experiments, it will be possible to decrease the lower bound of the energy windows ($`2.8`$ MeV) with only a moderate decrease of the expected sensitivity.
We remark that, even if the background is negligible, the low expected counting rates do not allow to observe a possible distortion of solar neutrino event spectrum, unless the $`{}_{}{}^{13}\mathrm{C}`$ abundance is enriched and/or one considers gigantic detectors. In principle, $`{}_{}{}^{13}\mathrm{C}`$ enrichment is possible.<sup>13</sup><sup>13</sup>13The $`{}_{}{}^{13}\mathrm{C}`$ is mainly used in health diagnostic, since it has a specific NMR signature. However, the current separation techniques probably do not allow a massive production of this isotope. For this reason, we do not consider in detail this possibility. One should note, however, that even a partial enrichment (e. g., corresponding to a $`{}_{}{}^{13}\mathrm{C}`$ abundance of the order $`2030`$%) could allow to obtain important results, like e. g., the high accuracy determination of the solar neutrino spectrum down to energies equal to about $`E_\nu 3`$ MeV (or the observation of hep solar neutrinos).
Finally, we briefly discuss the perspectives for gigantic liquid scintillator detectors planned in the future. In particular a $`30`$ kT detector, the Low Energy Neutrino Astrophysics (LENA), has been proposed . The site proposed for the experiment is Pyhäsalmi mine in Finland at a depth comparable to that of Borexino ($`4000`$ m.w.e.). This means that the cosmogenic background will be sufficiently low for the proposed measure. The scintillator should be composed mainly by PXE (C<sub>16</sub>H<sub>18</sub>, with a molecular mass $`\mu =210.31`$ and a density density $`\rho =0.998\mathrm{g}/\mathrm{cm}^3`$), but, of course, the final composition has yet to be decided. It is clear that, in such a large detector, it will be very hard to keep the internal background low. However, the gain in statistics will probably overcompensate this limitation. To give an example, in one year of data taking with a fiducial mass equal to $`30`$ kT, one obtains a better sensitivity than in SNO+ (with 1 kT fiducial mass), even assuming a background-to-signal ratio of the order of ten. For this reason, we believe that the LENA detector has the capability to perform a precise measure of the $`{}_{}{}^{8}\mathrm{B}`$ flux (comparable to that provided by Super-Kamiokande and/or SNO) in a few years of data taking.
## 7 Summary and conclusions
In this Letter we have discussed the possibility to detect $`{}_{}{}^{8}\mathrm{B}`$ solar neutrinos by using the $`\nu _\mathrm{e}`$ CC-interaction with $`{}_{}{}^{13}\mathrm{C}`$ nuclei naturally contained in organic liquid scintillators. The proposed detection process has a low threshold ($`Q=2.22\mathrm{MeV}`$) and large and well-know cross section. Moreover, one can take advantage of the subsequent decay of the produced $`{}_{}{}^{13}\mathrm{N}`$ nuclei to discriminate neutrino events from the background.
We have calculated the expected event rates (of the order of $`20`$ kTy<sup>-1</sup>) for KamLAND, Borexino and an hypothetical Borexino-like experiment situated at SNOlab (SNO+). Moreover, we have evaluated thoroughly all the possible sources of (external and internal) background in the three considered detectors. We have shown that the background-to-signal ratio is $`2`$ in KamLAND, while is much less than 1 in Borexino and SNO+.
Finally, we have calculated the expected sensitivity for the various experiments. Assuming an exposure equal to $`=1`$ kTy, the solar neutrino signal can be extracted with uncertainty of the order of $`50`$% in KamLAND and $`2025`$% in Borexino and SNO+. The expected sensitivity scales as $`^{1/2}`$, since background is directly measured by the experiments. Gigantic (such as LENA) and/or enriched detectors, having a much larger statistics, will have the possibility to perform a very precise measurement of the $`{}_{}{}^{8}\mathrm{B}`$ neutrino flux.
It should be stressed that the proposed measure does not require any modification of the standard experimental set-up. The KamLAND detector should be able to extract about $`18`$ $`{}_{}{}^{8}\mathrm{B}`$ solar neutrino events from the already collected data (corresponding to a total exposure equal to $`0.766`$ kTy).
## Acknowledgments
The authors are grateful to many participants of the NOW 2004 workshop for interesting discussions, where this work has begun. We thank C. Galbiati and D. Franco for helpful discussions on $`{}_{}{}^{13}\mathrm{C}(`$p,n$`)^{13}\mathrm{N}`$ background, and L. Oberauer and M. Chen for reading the manuscript and useful comments. F. V. also thanks G. Fiorentini for useful discussions and comments. This work is supported by the Italian Ministero dell’Istruzione, Università e Ricerca (MIUR) and Istituto Nazionale di Fisica Nucleare (INFN) through the “Astroparticle Physics” research project. |
warning/0506/hep-ph0506080.html | ar5iv | text | # Analytical and numerical evaluation of the Debye and Meissner masses in dense neutral three-flavor quark matter
## I Introduction
The rich phase structure of matter at high baryon density has attracted much interest over decades theoretically and phenomenologically. At sufficiently high density and low temperature, wherever quarks feel an attractive force as is suggested from the one gluon exchange between quarks that are antisymmetric in color, the Fermi surface of quark matter is unstable against forming a Cooper pair whose condensation leads to color superconductivity reviews . It is well established that the color-flavor locked (CFL) phase is a ground state of three-flavor quark matter in the asymptotically high density region where the strange quark mass $`M_s`$ is negligible as compared with the quark chemical potential $`\mu `$ Alford:1998mk .
If it is realized in a bulk system like the cores of compact stellar objects, quark matter must be neutral in electric and color charges Iida:2000ha ; Alford:2002kj ; Steiner:2002gx . As long as $`M_s/\mu `$ is negligibly small, three-flavor quark matter, which is composed of the equal number of $`u`$, $`d`$, and $`s`$ quarks, satisfies electric and color neutrality on its own regardless of whether it is in the normal or CFL phase. When $`M_s`$ comes to have a substantial effect suppressing $`s`$ quarks, there arise various possibilities in forming Cooper pairs and in achieving neutrality, that brings intricate subtleties into the phase structure especially in the intermediate density region where $`M_s`$ can compete $`\mu `$.
When the system is normal quark matter, a finite electron density is required to neutralize the system electrically because the number of $`s`$ quarks is reduced by $`M_s0`$. In CFL quark matter, on the other hand, neutrality can be fulfilled even without electrons rigidly at zero temperature Rajagopal:2000ff and approximately at low temperatures below an insulator-metal crossover Ruster:2004eg ; Fukushima:2004zq . This is because the BCS ansatz for Cooper pairing between cross-species of quarks enforces equality in the numbers of red, green, blue, and $`u`$, $`d`$, $`s`$ quarks and thus neutrality. At zero temperature, another candidate, that is, the two-flavor color superconducting (2SC) phase, is known to cost a larger free energy than the CFL phase under the neutrality conditions, or, it could exist for strong coupling between quarks in some density windows Alford:2002kj ; Alford:2003fq ; Fukushima:2004zq .
Electronless and neutral CFL quark matter lasts as long as the Fermi energy mismatch when two quarks were not paired, which is given by $`M_s^2/2\mu `$ as explained later, is less than the energy gap $`\mathrm{\Delta }`$. There should appear a new state of matter once the energy gain by releasing the pressure between cross-species of quarks surpasses the gap energy, which occurs when $`M_s^2/2\mu >\mathrm{\Delta }`$. Then, some of quark energy dispersion relations pass across zero and their pairing is disrupted in the corresponding momentum (blocking) region that is from one zero to another zero of the quark excitation energy. The new state is called the gapless CFL (gCFL) phase Alford:2003fq and, if any other phase transitions such as the chiral phase transition and the transition to a crystalline color superconductor loff ; Alford:2000ze lie at lower densities than the gapless onset, the gCFL phase must show up next as the density goes down from the CFL side.
The gapless or breached pairing superconductivity was first discussed in a non-relativistic model and was recognized as an unstable state Sarma and its analogue in three-flavor quark matter was considered in Ref. Alford:1999xc . It was discovered later on that stable gapless superconductivity could be possible under the constraint of particle number fixing in a non-relativistic model Gubankova:2003uj (see also Ref. Forbes:2004cr for more discussions on stability) and of electric and color neutrality in the gapless regime of the 2SC (g2SC) phase Shovkovy:2003uu .
Besides gapless superconductivity, there is another competing possibility, i.e. the mixed phase, to realize neutrality Neumann:2002jm ; Shovkovy:2003ce ; Bedaque:2003hi , and if the mixed phase has a lower free energy, the gapless phase would not come to realization. According to Ref. Reddy:2004my it can be the case in fact for the g2SC phase which should be taken over energetically by the mixed phase where normal and color superconducting phases coexist, and consequently the g2SC phase may be less of a reality. It must deserve further investigation, however, especially with the screening effects in the mixed phase taken into account to declare the existence or non-existence of the (g)2SC phase. In the case of the gCFL phase, on the other hand, the free energy comparison indicates that the gCFL phase should be stable energetically apart from a minor exception of the CFL-2SC mixed phase possibility, which will be presumably ruled out, however, once the surface tension and the Coulomb energy corrections are included Alford:2004nf . Hence, so far, the gCFL phase remains as a likely candidate for the ground state at moderate density and is considered to be relevant to neutron star physics Alford:2004zr .
However, chromomagnetic instability first found in the g2SC phase Huang:2004bg and later confirmed also in the gCFL phase Casalbuoni:2004tb ; Alford:2005qw implies that the gapless phase is not stable against perturbation of transverse (chromomagnetic) gluon fields and the true ground state must need something further. It is negative Meissner masses squared, i.e., imaginary Meissner masses, that the authors of Refs. Huang:2004bg ; Casalbuoni:2004tb ; Alford:2005qw have actually revealed in the gapless phases. Since the Meissner mass is the screening mass for transverse gluons, the instability signifies spontaneous generation of the expectation value of gauge fields. Keeping in mind a simple relation in superconductors between the gauge fields and currents in the London gauge, we can regard chromomagnetic instability as color current generation. (An interpretation of instability as spontaneous baryon current generation has been argued in Ref. Huang:2005pv .) To put it another way, it is possible to interpret it as instability towards a state characterized by diquark condensates with color-dependent oscillation in coordinate space Giannakis:2004pf , because gauge fields in the quark propagator generally reside in the form of the covariant derivative. This state of quark matter is a sort of crystalline color superconductors under the single colored plain wave ansatz. It should be worth emphasizing that the difference is only in intuitive pictures, and in effect, these above interpretations are equivalent, which is mathematically described by the gauge transformation. Still, anyway, it is controversial what the true ground state should be that supersedes the homogeneous gCFL phase.
The purpose of this paper is to compute the Debye and Meissner masses in the (g)CFL phase and to present detailed analyses on the nature of chromomagnetic instability which occurs in connection with the presence of gapless quarks. Our central results are condensed in the unstable regions on the phase diagram presented in Figs. 1, 2, and 3. In both the CFL and gCFL phases there is no mixing with other gluons for $`A_1`$, $`A_2`$, $`A_4`$, $`A_5`$, $`A_6`$, and $`A_7`$, among which two gluons of $`(A_1,A_2)`$, $`(A_4,A_5)`$, and $`(A_6,A_7)`$ respectively have the same screening masses. We shall write $`A_{1,2}`$ for instance to mean whichever of $`A_1`$ or $`A_2`$ equivalently from now on. As for $`A_3`$, $`A_8`$, and the electromagnetic field $`A_\gamma `$, on the other hand, mixing between them arises and thus we have to refer to eigenvalues of the $`3\times 3`$ mass squared matrix to investigate instability. It should be noted that the color indices are labeled according to the conventional Gell-Mann matrices in color space.
Figure 1 shows the unstable region enclosed with thick lines inside which the Meissner mass squared for $`A_{1,2}`$ is negative. The phase diagram underlaid is what has been revealed in Ref. Fukushima:2004zq for the case of weak coupling yielding $`\mathrm{\Delta }_0=25\mathrm{MeV}`$ at $`M_s=T=0`$. The gCFL phase starts appearing at $`M_s^2/\mu =47.1\mathrm{MeV}`$ which is well close to the kinematical estimate $`(M_s^2/\mu )_\mathrm{c}2\mathrm{\Delta }`$. In the vicinity of the phase boundary line on which one of three gaps vanishes, i.e., $`\mathrm{\Delta }_2=0`$, there comes out a rather complicated structure which is not robust but strongly depends on a choice of the coupling strength. We will not present results for other coupling cases as done in Ref. Fukushima:2004zq , for our purpose here is not to disclose the phase structure but to argue characteristic features of chromomagnetic instability and to settle where it becomes relevant on the phase diagram.
We observe that instability grows for $`A_{1,2}`$ as soon as the gCFL state occurs at zero temperature. The instability penetrates from the CFL region where all $`u`$-$`d`$, $`d`$-$`s`$, and $`s`$-$`u`$ quarks make Cooper pairs into the uSC region where $`u`$-$`d`$ and $`s`$-$`u`$ quarks remain to pair, but never enters the 2SC region where only the $`u`$-$`d`$ quark pairing gap is nonvanishing. Actually in the 2SC phase, $`A_{1,2}`$ become irrelevant to the Higgs-Anderson mechanism, and therefore the Meissner mass for $`A_{1,2}`$ must be zero there.
The unstable regions for $`A_{4,5}`$ and $`A_{6,7}`$ are depicted in Fig. 2. In the 2SC phase there is no discrimination between $`A_{4,5}`$ and $`A_{6,7}`$ and the difference is manifested only when quark matter lies in the CFL or uSC phase. Our mapping is consistent with all of the known results both in the 2SC phase Huang:2004bg and in the CFL phase Casalbuoni:2004tb ; no instability takes place near the gCFL onset for $`A_{4,5}`$ nor $`A_{6,7}`$, while the large $`M_s^2/\mu `$ regions exhibit instability which turns out to be related to the g2SC instability. At small temperatures and large $`M_s^2/\mu `$ the Meissner masses squared for $`A_{4,5}`$ and $`A_{6,7}`$ are negative until the system reaches the phase boundary from which the system enters the phase of unpaired quark matter (UQM).
It seems to be somewhat tricky to understand the unstable regions for mixed modes of $`A_3`$, $`A_8`$, and $`A_\gamma `$ shown in Fig. 3. At zero or extremely low temperatures whose energy scale is determined by the electron contribution as discussed later, the whole gCFL phase is unstable and our results qualitatively agree with Ref. Casalbuoni:2004tb . At finite temperature there are two distinct regions where instability remains. The instability near the gCFL onset eventually disappears as the temperature grows, while the instability at larger $`M_s^2/\mu `$ goes further into the 2SC phase. In fact in the 2SC phase, $`A_3`$ decouples from others and its Meissner mass is reduced to zero, and one of the two eigenmodes composed of $`A_8`$ and $`A_\gamma `$ exhibits chromomagnetic instability as is consistent with the findings in the g2SC phase Huang:2004bg . For even stronger coupling when the 2SC (not necessarily g2SC) phase is possible at zero temperature (see Fig. 17 in Ref. Fukushima:2004zq ), we have numerically checked that instability at larger $`M_s^2/\mu `$ starts exactly at the g2SC onset.
The exact correspondence between the gapless and instability onsets is apparent only at zero temperature since the finite-temperature effects allow for thermal quark excitations which are not clearly distinguishable from quarks in the blocking region. An interesting manifestation of this is the unstable region in the uSC phase in which there are no gapless quarks below $`T5\mathrm{MeV}`$ (see the dotted curve in the uSC region shown in Fig. 1 in Ref. Fukushima:2004zq ). In a sense, at finite temperature, we can say that the instability penetrates into gapped sides, as is also observed from Fig. 5 in Ref. Alford:2005qw .
We are explaining how we have come by these instability mappings in a numerical way in later discussions, starting with the following subsections in which we shall make a brief review of the gCFL phase and of the derivation of the Debye and Meissner screening masses of gauge fields. The readers who are already familiar with these basics can skip most of them and jump to Sec. II.
### I.1 Gapless CFL phase
Our strategy to approach the Debye and Meissner masses is based on the thermodynamic potential that has been formulated within an effective model. In the following subsubsections, we discuss the model and approximations to describe the gCFL phase, and then illustrate the quark excitation energies as a function of the momentum (i.e. the dispersion relations).
#### I.1.1 Model and approximations
In this paper we adopt essentially the same model and approximations as used in Refs. Alford:2003fq ; Fukushima:2004zq . The only difference is that the $`M_s`$ effect is incorporated as an effective chemical potential shift. This approximation is necessary to relate the Meissner mass to the potential curvature with respect to gluon source fields in a simple way. The model employed here is the Nambu–Jona-Lasinio (NJL) model with four-fermion interaction. In our model study we assume that the predominant diquark condensate is antisymmetric in Dirac indices, antisymmetric in color, and thus antisymmetric in flavor;
$$\psi _i^aC\gamma _5\psi _j^b\mathrm{\Delta }_1ϵ^{ab1}ϵ_{ij1}+\mathrm{\Delta }_2ϵ^{ab2}ϵ_{ij2}+\mathrm{\Delta }_3ϵ^{ab3}ϵ_{ij3},$$
(1)
where ($`i`$,$`j`$) and ($`a`$,$`b`$) represent the flavor indices ($`u`$,$`d`$,$`s`$) and the color triplet indices (red,green,blue) respectively. The gap parameters $`\mathrm{\Delta }_1`$, $`\mathrm{\Delta }_2`$, and $`\mathrm{\Delta }_3`$ describe a $`9\times 9`$ matrix in color-flavor space that takes the form,
$$𝚫=\left(\begin{array}{ccccccccc}0& \mathrm{\Delta }_3& \mathrm{\Delta }_2& 0& 0& 0& 0& 0& 0\\ \mathrm{\Delta }_3& 0& \mathrm{\Delta }_1& 0& 0& 0& 0& 0& 0\\ \mathrm{\Delta }_2& \mathrm{\Delta }_1& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& \mathrm{\Delta }_3& 0& 0& 0& 0\\ 0& 0& 0& \mathrm{\Delta }_3& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& \mathrm{\Delta }_2& 0& 0\\ 0& 0& 0& 0& 0& \mathrm{\Delta }_2& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& \mathrm{\Delta }_1\\ 0& 0& 0& 0& 0& 0& 0& \mathrm{\Delta }_1& 0\end{array}\right)$$
(2)
in the basis $`(ru,gd,bs,gu,rd,rs,bu,bd,gs)`$. In the same way as in Refs. Alford:2003fq ; Fukushima:2004zq , we ignore diquark condensates which are symmetric in color, though we know that they do not break any new symmetry and consequently can take a finite value. This approximation is motivated by the fact that the QCD interaction is repulsive between quarks in this channel and such diquark condensates have been actually known to be insignificant quantitatively Ruster:2004eg .
In the mean-field approximation where the condensates (1) are relevant, we only have to consider the four-fermion interaction in the diquark-diquark channel that can be generally in hand after appropriate Fierz transformation. That is,
$$=\overline{\psi }(i\text{/}+𝝁\gamma ^0𝑴)\psi +_{\text{int}},$$
(3)
where
$$\begin{array}{cc}\hfill _{\text{int}}& =\frac{G}{4}(\overline{\psi }_i^ai\gamma _5ϵ^{ab\eta }ϵ_{ij\eta }C\overline{\psi }_j^{Tb})\hfill \\ & \times (\psi _i^{}^{Ta^{}}Ci\gamma _5ϵ^{a^{}b^{}\eta }ϵ_{i^{}j^{}\eta }\psi _j^{}^b^{}).\hfill \end{array}$$
(4)
Here the mass matrix $`𝑴`$ is unity in color and $`\text{diag}(0,0,M_s)`$ in flavor ($`u`$,$`d`$,$`s`$) space in our approximation and $`𝝁`$ the matrix of quark chemical potentials in color-flavor space. Bold symbols generally denote matrices in color-flavor space. The charge conjugation matrix is $`C=i\gamma ^2\gamma ^0`$.
The effect of nonzero $`𝑴`$ on the particle dispersion relation becomes apparent in the vicinity of the Fermi surface, where it can be well approximated by a shift in quark chemical potentials by $`𝑴𝑴^{}/2\mu `$. We shall make use of this prescription that actually turns out to be essential to allow us to relate the Meissner mass to the potential curvature. In the present work, as in Refs. Alford:2003fq ; Fukushima:2004zq , $`𝑴`$ is not current but constituent quark mass and is treated as an input parameter. We know that this approximation provides a correct phase structure unless the chiral phase transition cuts deeply into the gCFL region under the choice of strong coupling in the chiral sector Abuki:2004zk .
Chemical potentials are fixed by equilibration and neutrality. Stable quark matter should be a color singlet as a whole and neutral in the electric charge under $`\beta `$-equilibrium. Color singletness is a more stringent condition than neutrality. It has been shown, however, that no energy cost is needed to project a neutral state onto a singlet state in color in the thermodynamic limit Amore:2001uf . Hence, it is sufficient to impose global neutrality with respect to the electric and color charges. To describe that, we shall consider $`\mu _e`$, $`\mu _3`$, and $`\mu _8`$ that are coupled to negative $`Q=\text{diag}(\frac{2}{3},\frac{1}{3},\frac{1}{3})`$ in flavor ($`u`$,$`d`$,$`s`$) space so that a positive $`\mu _e`$ corresponds to the electron density, $`T_3=\text{diag}(\frac{1}{2},\frac{1}{2},0)`$, and $`\frac{2}{\sqrt{3}}T_8=\text{diag}(\frac{1}{3},\frac{1}{3},\frac{2}{3})`$ in color (red, green, blue) space, respectively. Then, 9 diagonal components of the effective chemical potential matrix $`𝝁_{\text{eff}}`$ (including the shift by $`𝑴𝑴^{}/2\mu `$) in color-flavor space are explicitly
$$\begin{array}{ccc}\hfill \mu _{ru}& =& \mu \frac{2}{3}\mu _e+\frac{1}{2}\mu _3+\frac{1}{3}\mu _8,\hfill \\ \hfill \mu _{gd}& =& \mu +\frac{1}{3}\mu _e\frac{1}{2}\mu _3+\frac{1}{3}\mu _8,\hfill \\ \hfill \mu _{bs}& =& \mu +\frac{1}{3}\mu _e\frac{2}{3}\mu _8\frac{M_s^2}{2\mu },\hfill \\ \hfill \mu _{gu}& =& \mu \frac{2}{3}\mu _e\frac{1}{2}\mu _3+\frac{1}{3}\mu _8,\hfill \\ \hfill \mu _{rd}& =& \mu +\frac{1}{3}\mu _e+\frac{1}{2}\mu _3+\frac{1}{3}\mu _8,\hfill \\ \hfill \mu _{rs}& =& \mu +\frac{1}{3}\mu _e+\frac{1}{2}\mu _3+\frac{1}{3}\mu _8\frac{M_s^2}{2\mu },\hfill \\ \hfill \mu _{bu}& =& \mu \frac{2}{3}\mu _e\frac{2}{3}\mu _8,\hfill \\ \hfill \mu _{bd}& =& \mu +\frac{1}{3}\mu _e\frac{2}{3}\mu _8.\hfill \\ \hfill \mu _{gs}& =& \mu +\frac{1}{3}\mu _e\frac{1}{2}\mu _3+\frac{1}{3}\mu _8\frac{M_s^2}{2\mu },\hfill \end{array}$$
(5)
The gap parameters corresponding to (1) in this model are precisely defined as
$$\mathrm{\Delta }_\eta =\frac{1}{2}G\psi _i^{T\alpha }Ci\gamma _5ϵ^{\alpha \beta \eta }ϵ_{ij\eta }\psi _j^\beta .$$
(6)
With all these definitions, if we choose Nambu-Gor’kov basis as $`\mathrm{\Psi }(p)=(\psi (p),C\overline{\psi }^T(p))^T`$, we can express the mean-field (inverse) quark propagator in a simple form,
$$iS^1(p)=\left(\begin{array}{cc}\text{/}p+𝝁\gamma ^0& i\gamma _5𝚫\\ i\gamma _5𝚫& \text{/}p𝝁\gamma ^0\end{array}\right)$$
(7)
after rearranging the Dirac matrices. The quark propagator is a $`72\times 72`$ matrix in color, flavor, spin, and Nambu-Gor’kov space. From zeros of the inverse propagator, we can read 72 energy dispersion relations $`\epsilon _i(p)`$. The thermodynamic potential $`\mathrm{\Omega }`$ is thus written in terms of $`\epsilon _i`$’s as
$$\begin{array}{cc}\hfill \mathrm{\Omega }& =\frac{1}{8\pi ^2}_0^\mathrm{\Lambda }𝑑pp^2\underset{j=1}{\overset{72}{}}\left\{|\epsilon _i|+2T\mathrm{ln}\left(1+e^{|\epsilon _i|/T}\right)\right\}\hfill \\ & +\frac{1}{G}\left(\mathrm{\Delta }_1^2+\mathrm{\Delta }_2^2+\mathrm{\Delta }_3^2\right)\frac{\mu _e^4}{12\pi ^2}\frac{\mu _e^2T^2}{6}\frac{7\pi ^2T^4}{180}\hfill \end{array}$$
(8)
with the electron contribution in the last three terms. In order to fix three gap parameters and three chemical potentials at each $`M_s`$, $`\mu `$, and $`T`$, we will simultaneously solve three gap equations,
$$\frac{\mathrm{\Omega }}{\mathrm{\Delta }_1}=\frac{\mathrm{\Omega }}{\mathrm{\Delta }_2}=\frac{\mathrm{\Omega }}{\mathrm{\Delta }_3}=0,$$
(9)
and three neutrality conditions,
$$\frac{\mathrm{\Omega }}{\mu _e}=\frac{\mathrm{\Omega }}{\mu _3}=\frac{\mathrm{\Omega }}{\mu _8}=0.$$
(10)
Here $`\mathrm{\Lambda }`$ in (8) is the ultraviolet cut-off parameter and we use $`\mathrm{\Lambda }=800\mathrm{MeV}`$ as in Refs. Alford:2003fq ; Fukushima:2004zq . The coupling constant $`G`$ is chosen to yield $`\mathrm{\Delta }_0=25\mathrm{MeV}`$ when both $`M_s`$ and $`T`$ are zero. In the present work we take the quark chemical potential $`\mu =500\mathrm{MeV}`$ that roughly corresponds to 10 times the normal nuclear density in the model.
When $`M_s`$ is not so large as to disrupt any Cooper pair, $`\mu _e`$, $`\mu _3`$, and $`\mu _8`$ satisfying electric and color neutrality have been analyzed at zero temperature in a model-independent way Alford:2002kj . Two of three can be fixed as a function of the third that we will choose $`\mu _e`$ here, then
$`\mu _3`$ $`=\mu _e,`$ (11)
$`\mu _8`$ $`={\displaystyle \frac{M_s^2}{2\mu }}+{\displaystyle \frac{\mu _e}{2}}.`$ (12)
The important feature of the CFL phase we should note is that the thermodynamic potential (8) at $`T=0`$ with the electron contribution disregarded is independent of
$$\mu _{\stackrel{~}{Q}}=\frac{4}{9}\left(\mu _e+\mu _3+\frac{1}{2}\mu _8\right),$$
(13)
that is the chemical potential for $`\stackrel{~}{Q}=QT_3\frac{1}{\sqrt{3}}T_8`$. It is known that $`\stackrel{~}{Q}`$ is a generator of the “rotated electromagnetism” that is never broken by any condensate of the form (1). At zero temperature $`\stackrel{~}{Q}`$-charged excitations in CFL quark matter are all gapped and the CFL phase is a $`\stackrel{~}{Q}`$-insulator Rajagopal:2000ff .
In the presence of the electron contribution to the thermodynamic potential, in reality, there is a gentle curvature on the plateau of the thermodynamic potential provided by the last three terms in (8) which select $`\mu _e=0`$, meaning zero electron density. Although it makes only slight changes in energy, the existence of the electron effects is crucial in understanding the $`A_{1,2}`$ instability.
#### I.1.2 Gapless dispersion relations
Since the gap matrix (2) is block-diagonal, the $`72\times 72`$ quark propagator matrix is block-diagonal as well in color-flavor space and can be divided into four parts; one $`24\times 24`$ part for $`ru`$-$`gd`$-$`bs`$ quark pairing with $`\mathrm{\Delta }_1`$, $`\mathrm{\Delta }_2`$, and $`\mathrm{\Delta }_3`$ and three $`16\times 16`$ parts for quark pairing of $`bd`$-$`gs`$ with $`\mathrm{\Delta }_1`$, $`rs`$-$`bu`$ with $`\mathrm{\Delta }_2`$, and $`gu`$-$`rd`$ with $`\mathrm{\Delta }_3`$. It might be possible but quite hard to handle the $`24\times 24`$ part analytically. Fortunately, as we will see shortly, this $`24\times 24`$ intricate part has little to do with chromomagnetic instability, so we will limit our discussion here to the propagator for $`bd`$-$`gs`$, $`rs`$-$`bu`$, and $`gu`$-$`rd`$ pairings.
Leaving the explicit expression of the quark propagator until Sec. II, let us elucidate here the quark energy dispersion relations obtained from zeros of the inverse propagator. In general for the $`16\times 16`$ part involving two species $`A`$ and $`B`$ quarks the energy dispersion relation takes the form,
$`\epsilon _{A\text{-}B}`$ $`=\sqrt{(p\overline{\mu }_{A\text{-}B})^2+\mathrm{\Delta }_{AB}^2}+\delta \mu _{A\text{-}B},`$ (14)
$`\stackrel{~}{\epsilon }_{A\text{-}B}`$ $`=\sqrt{(p+\overline{\mu }_{A\text{-}B})^2+\mathrm{\Delta }_{AB}^2}+\delta \mu _{A\text{-}B}`$ (15)
for quasi-particle and antiparticle excitations respectively, where $`\overline{\mu }_{A\text{-}B}=\frac{1}{2}(\mu _A+\mu _B)`$ and $`\delta \mu _{A\text{-}B}=\frac{1}{2}(\mu _A\mu _B)`$ and $`\mathrm{\Delta }_{AB}`$ is the energy gap for $`A`$-$`B`$ pairing. In our definition $`\delta \mu _{A\text{-}B}=\delta \mu _{B\text{-}A}`$. Gapless dispersion relations come about once
$$|\delta \mu _{A\text{-}B}|>\mathrm{\Delta }_{AB}.$$
(16)
Then, using effective chemical potentials (5) with the known CFL solutions (11) and (12), we have
$`\delta \mu _{gu\text{-}rd}`$ $`=\frac{1}{2}\mu _e\frac{1}{2}\mu _3=\mu _e,`$ (17)
$`\delta \mu _{rs\text{-}bu}`$ $`=\frac{1}{2}\mu _e+\frac{1}{4}\mu _3+\frac{1}{2}\mu _8\frac{M_s^2}{4\mu }=\mu _e\frac{M_s^2}{2\mu },`$ (18)
$`\delta \mu _{bd\text{-}gs}`$ $`=\frac{1}{4}\mu _3\frac{1}{2}\mu _8+\frac{M_s^2}{4\mu }=\frac{M_s^2}{2\mu },`$ (19)
and
$$\overline{\mu }_{gu\text{-}rd}=\overline{\mu }_{bd\text{-}gs}=\overline{\mu }_{rs\text{-}bu}=\mu \frac{M_s^2}{6\mu }$$
(20)
in the CFL phase. If it were not for the electron terms in the thermodynamic potential, $`\mu _e`$ could lie anywhere as far as it does not disrupt the pairing of $`gu`$-$`rd`$ nor $`rs`$-$`bu`$ quarks which are $`\stackrel{~}{Q}`$-charged. In the region between the dashed and dot-dashed curves shown in Fig. 4, the system remains to be a $`\stackrel{~}{Q}`$-insulator, meaning the bandgap. The CFL solution ceases to exist at the point where these two curves meet, before which a first-order phase transition to unpaired quark matter takes place as indicated by the dotted vertical line (see discussions in Ref. Alford:2003fq for details).
By substituting $`\mu _e=0`$, as is favored in the CFL phase with electrons, for the above expressions, we readily realize that the $`bd`$-$`gs`$ and $`rs`$-$`bu`$ dispersion relations are identical supposing $`\mathrm{\Delta }_1=\mathrm{\Delta }_2=\mathrm{\Delta }`$. In fact, the $`bd`$-$`gs`$ and $`rs`$-$`bu`$ quark pairings are breached at the same time when
$$|\delta \mu _{bd\text{-}gs}|=|\delta \mu _{rs\text{-}bu}|=\frac{M_s^2}{2\mu }>\mathrm{\Delta },$$
(21)
which corresponds to the fact that in Fig. 4 the solid and dashed lines cross at $`\mu _e=0`$. This coincidence of the gapless onset locations for $`bd`$-$`gs`$ and $`rs`$-$`bu`$ quarks is responsible for instability for $`A_{1,2}`$ in the gCFL phase, as only $`A_{1,2}`$ can excite a pair of $`rs`$ and $`gs`$ quarks at once (see Fig. 6). In other words, if the electron terms are absent and $`\mu _e`$ is chosen to be positive, say $`10\mathrm{MeV}`$, the $`A_{1,2}`$ instability would disappear, while there would remain the instability for other gluons.
Once quark matter enters the gCFL phase with larger $`M_s^2/\mu `$ than the onset, the $`bd`$-$`gs`$ and $`rs`$-$`bu`$ quark dispersion relations are both gapless but no longer degenerated. The blocking momentum region for $`bd`$-$`gs`$ quarks becomes wider with increasing $`M_s^2/\mu `$ as shown by the solid curve in Fig. 5 and $`\mathrm{\Delta }_1`$ decreases accordingly because quarks within the blocking region cannot participate in pairing. Since $`rs`$ and $`bu`$ quarks have nonzero $`\stackrel{~}{Q}`$-charge, the blocking momentum region for $`rs`$-$`bu`$ quarks is not allowed to be arbitrarily wider and is determined by the requirement to cancel $`\stackrel{~}{Q}`$-charge brought by electrons whose density is $`\mu _e^3`$. That means, the width of the blocking momentum region for $`rs`$-$`bu`$ quarks is estimated as $`\delta p\mu _e^3/3\overline{\mu }^2\overline{\mu }`$ Alford:2003fq . The actual form of the $`rs`$-$`bu`$ dispersion relation is, as a result, kept to be almost quadratic with a tiny blocking region anywhere in the gCFL phase as seen by the dashed curve in Fig. 5.
Although the $`ru`$-$`gd`$-$`bs`$ part is hard to anticipate a priori, there is no gapless mode involved in this sector. The dotted curve in Fig. 5 represents one of energy dispersion relations for $`ru`$-$`gd`$-$`bs`$ quarks that has the lowest energy among them. It is clear from the figure that $`ru`$-$`gd`$-$`bs`$ quarks are all gapped at $`M_s^2/\mu =80\mathrm{MeV}`$ and it is also the case for larger $`M_s^2/\mu `$. Hence, the three (solid, dashed, and dot-dashed) lines drawn for $`bd`$-$`gs`$, $`rs`$-$`bu`$, and $`gu`$-$`rd`$ quarks respectively in Fig. 4 suffice for judging if gapless quarks are present or not.
### I.2 Debye and Meissner screening masses
In a medium at finite temperature and density the gauge fields are screened by the polarization of charged thermal excitations. The screening effect on the gauge fields differs depending on the longitudinal and transverse directions due to the lack of Lorentz invariance. The screening mass, which is the inverse of the screening length, in the longitudinal direction is called the (chromo)electric or Debye screening mass, and in (color) superconductors in particular the transverse screening mass, the (chromo)magnetic or Meissner screening mass in QED (QCD). In the normal phase any perturbative calculation leads to vanishing magnetic mass in the static limit. In the superconducting phase, on the other hand, a finite magnetic screening mass results from the Higgs-Anderson mechanism and embodies the Meissner effect so that superconductors can exclude the magnetic field.
The Debye and Meissner masses are calculated from the eigenvalues of the mass squared matrix in color defined by the self-energies;
$`m_{D,\alpha \beta }^2`$ $`=\underset{q0}{lim}\mathrm{\Pi }_{\alpha \beta }^{00}(\omega =0,\stackrel{}{q}),`$ (22)
$`m_{M,\alpha \beta }^2`$ $`=\frac{1}{2}\underset{q0}{lim}\left(\delta _{ij}\widehat{q}_i\widehat{q}_j\right)\mathrm{\Pi }_{\alpha \beta }^{ij}(\omega =0,\stackrel{}{q}),`$ (23)
where ($`\alpha `$,$`\beta `$) represent the color octet (or photon) indices and we defined $`\widehat{q}=\stackrel{}{q}/|\stackrel{}{q}|`$, and $`\mathrm{\Pi }_{\alpha \beta }^{\mu \nu }`$ is the self-energy with the Lorentz indices ($`\mu `$,$`\nu `$) for the gluons ($`\alpha ,\beta 8`$) or photon ($`\gamma `$ instead of $`\alpha ,\beta `$) or their mixing. (We always use the subscript $`\gamma `$ for the sake of meaning photon.) At the one-loop level at high density, the quark loop contribution to $`\mathrm{\Pi }_{\alpha \beta }^{\mu \nu }`$, which is proportional to $`\mu ^2`$, is predominant, so we shall consider only the polarization tensor coming from one-loop diagrams of quarks.
In the symmetric CFL state with $`M_s=0`$, the mass squared matrix is diagonal in color and all eight gluons have the same screening mass. The Debye and Meissner masses in a color superconductor have been calculated diagrammatically in Ref. Rischke:2000qz and we shall quote the results here;
$`m_D^2`$ $`={\displaystyle \frac{218\mathrm{ln}2}{6}}m_g^22.576m_g^2,`$ (24)
$`m_M^2`$ $`={\displaystyle \frac{1}{3}}m_D^20.859m_g^2`$ (25)
for all eight gluons. Here $`m_g^2=g^2\mu ^2/6\pi ^2`$ and $`g`$ is the strong coupling constant. Note that our definition of $`m_g^2`$ is different by a factor $`N_f=3`$ from the definition given in Ref. Rischke:2000qz .
Once mixing between gluons and photon, which exists among $`A_8`$ and $`A_\gamma `$ in the CFL phase at $`M_s=0`$, and generally among $`A_3`$, $`A_8`$, and $`A_\gamma `$ for $`M_s0`$, is taken into account, the screening masses are to be modified. They actually have been evaluated analytically only when $`M_s=0`$ Schmitt:2003aa and after diagonalization two eigenvalues composed of $`A_8`$ and $`A_\gamma `$ read
$`\stackrel{~}{m}_{D,88}^2`$ $`=\left(1+{\displaystyle \frac{4}{3}}{\displaystyle \frac{e^2}{g^2}}\right)m_D^2,`$ (26)
$`\stackrel{~}{m}_{M,88}^2`$ $`=\left(1+{\displaystyle \frac{4}{3}}{\displaystyle \frac{e^2}{g^2}}\right)m_M^2,`$ (27)
$`\stackrel{~}{m}_{D,\gamma \gamma }^2`$ $`=\stackrel{~}{m}_{M,\gamma \gamma }^2=0,`$ (28)
where $`e`$ is the electromagnetic coupling constant. The last relation (28) follows from the facts that the condensates (1) leave a $`\mathrm{U}(1)_{\stackrel{~}{Q}}`$ symmetry unbroken (that is; $`\stackrel{~}{m}_{M,\gamma \gamma }^2=0`$), and that the system is a $`\stackrel{~}{Q}`$-insulator in which $`\stackrel{~}{Q}`$-charged excitations are all gapped (that is; $`\stackrel{~}{m}_{D,\gamma \gamma }^2=0`$).
For later convenience we shall carefully look into how these analytical expressions (24) and (25) are structured from distinct contributions of particles and antiparticles. When quarks are massless, the quark propagator can be specifically separated into the particle and antiparticle parts (see Eq. (40)). At the quark one-loop level, the gluon or photon self-energies consist of the particle-particle, particle-antiparticle, and antiparticle-antiparticle excitations. (In the case of the normal phase the particle-particle excitations in the vector channels would be to be articulated as the particle-hole excitations.) The last contribution from only the antiparticle excitations is just negligible so that we shall always omit it in our discussion.
In the $`2\times 2`$ space of Nambu-Gor’kov doubling, the quark propagator has not only the diagonal (normal) components but also the off-diagonal (abnormal) components proportional to the gap $`\mathrm{\Delta }`$ that connect two quasi-particles.
The Debye mass arises only from the particle-particle contribution. It can be further divided into the Nambu-Gor’kov diagonal and off-diagonal parts as follows;
$$m_D^2=[m_D^2]_{\text{diag}}+[m_D^2]_{\text{off}},$$
(29)
where
$`[m_D^2]_{\text{diag}}={\displaystyle \frac{9}{2}}m_g^2,`$ (30)
$`[m_D^2]_{\text{off}}=\left(1+{\displaystyle \frac{4}{3}}\mathrm{ln}2\right)m_g^2.`$ (31)
The sum of these two parts properly reproduces (24) as it should.
In evaluating the Meissner mass, on the other hand, it is essentially important to note that the particle-antiparticle excitation produces a significant contribution comparable to the particle-particle one. The Meissner mass consists of the diagonal and off-diagonal parts,
$$m_M^2=[m_M^2]_{\text{diag}}+[m_M^2]_{\text{off}},$$
(32)
that is further split into the particle-particle (p-p) and particle-antiparticle (p-a) parts,
$$[m_M^2]_{\text{diag}}=[m_M^2]_{\text{diag(p-p)}}+2[m_M^2]_{\text{diag(p-a)}}$$
(33)
with
$`[m_M^2]_{\text{diag(p-p)}}`$ $`={\displaystyle \frac{1}{3}}[m_D^2]_{\text{diag}},`$ (34)
$`[m_M^2]_{\text{diag(p-a)}}`$ $`={\displaystyle \frac{3}{2}}m_g^2,`$ (35)
and
$$[m_M^2]_{\text{off}}=[m_M^2]_{\text{off(p-p)}}=\frac{1}{3}[m_D^2]_{\text{off}},$$
(36)
from which we can easily make sure $`m_M^2=\frac{1}{3}m_D^2`$ as is well-known in the CFL phase Rischke:2000qz ; Son:1999cm . The particle-antiparticle contribution in the off-diagonal part $`[m_M^2]_{\text{off(p-a)}}`$ is vanishingly small, so we dropped it off from the above relations.
It is of great importance to realize that the diagonal p-p contribution to the Meissner mass is negative one third of the Debye mass counterpart, while the p-a excitations provide a positive contribution that is twice larger than the p-p’s. Adding them together we finally acquire the Meissner mass contribution that is positive one third of the Debye mass contribution. In the off-diagonal part, in contrast, the situation is much simpler and the ratio is positive one third itself.
In other words, in the diagonal part, the p-p loops always tend to induce paramagnetism, while diamagnetism originates from the p-a loops. Usually in the superconducting phase, the diamagnetic tendency is greater enough to bring about the Meissner effect. In gapless superconductors, however, antiparticles are never gapless and only the p-p loops are abnormally enhanced due to gapless quarks and their large density of states near the Fermi surface causes the opposite phenomenon to the Meissner effect, i.e., chromomagnetic instability. Also, this can be seen from the relations (34) and (36) which hold in the presence of finite $`M_s`$ or even in the gCFL phase. In the gCFL phase, a large density of states leads to a large Debye mass. Then, through (34) and (36), it should be accompanied by a large negative contribution to the Meissner mass squared, which eventually results in chromomagnetic instability.
As for (35), because the $`M_s`$ dependence of antiparticle excitations is suppressed by $`M_s/\mu `$, the relation is hardly changed for any $`M_s\mu `$.
## II computation of self-energies
In this section we shall diagrammatically compute the calculable parts of one-loop self-energies for gluons and photon and derive the analytical expression for the singular part of the Debye and Meissner masses.
Using the mean-field quark propagator $`S(p)`$, and the vertex matrices, $`\mathrm{\Gamma }_\alpha ^\mu `$, in color, flavor, spin, and Nambu-Gor’kov space, we can write the one-loop self-energies in a general form as
$$i\mathrm{\Pi }_{\alpha \beta }^{\mu \nu }(q)=\frac{1}{2}^T\frac{d^4p}{(2\pi )^4}\text{tr}\mathrm{\Gamma }_\alpha ^\mu S(q+p)\mathrm{\Gamma }_\beta ^\nu S(p).$$
(37)
The vertices $`\mathrm{\Gamma }_\alpha ^\mu `$ take a matrix form,
$$\mathrm{\Gamma }_\alpha ^\mu =\left[\begin{array}{cc}ig\gamma ^\mu T_\alpha \delta _{ij}& 0\\ 0& ig\gamma ^\mu T_\alpha \delta _{ij}\end{array}\right]$$
(38)
for gluons ($`\alpha =1,\mathrm{},8`$) where $`T_\alpha `$ are the $`\mathrm{SU}(3)_{\text{color}}`$ generators defined by the Gell-Mann matrices in color space and normalized as $`\text{tr}T_\alpha T_\beta =\frac{1}{2}\delta _{\alpha \beta }`$. For photon the vertex reads
$$\mathrm{\Gamma }_\gamma ^\mu =\left[\begin{array}{cc}ie\gamma ^\mu Q_{ij}& 0\\ 0& ie\gamma ^\mu Q_{ij}\end{array}\right],$$
(39)
which is unity in color. We see that the flavor is not changed at any vertices but the color can be converted through the off-diagonal components in $`T_1`$, $`T_2`$, $`T_4`$, $`T_5`$, $`T_6`$, and $`T_7`$. We call the gluons corresponding to $`T_3`$ and $`T_8`$ as the color-diagonal gluons.
The quark propagator is divided into the particle and antiparticle parts by the energy projection operators $`\mathrm{\Lambda }_p^\pm =\frac{1}{2}(1\pm \gamma ^0\stackrel{}{\gamma }\widehat{p})`$. The explicit propagator in the $`bd`$-$`gs`$ sector, for instance, can be written down after some rearrangement of the Nambu-Gor’kov (1,2) (where 2 is assigned to the Nambu-Gor’kov doubler in our convention) and color-flavor ($`bd`$,$`gs`$) indices, that is, in the basis ($`bd`$-1,$`gs`$-2,$`gs`$-1,$`bd`$-2), we have
$$\begin{array}{cc}& iS(p)=\left[\begin{array}{cccc}\mathrm{\Lambda }_p^{}& 0& 0& 0\\ 0& \mathrm{\Lambda }_p^+& 0& 0\\ 0& 0& \mathrm{\Lambda }_p^{}& 0\\ 0& 0& 0& \mathrm{\Lambda }_p^+\end{array}\right]\gamma ^0\left[\begin{array}{cc}S_{bd\text{-}gs}^\mathrm{a}(p)& 0\\ 0& S_{gs\text{-}bd}^\mathrm{a}(p)\end{array}\right]\hfill \\ & +\left[\begin{array}{cccc}\mathrm{\Lambda }_p^+& 0& 0& 0\\ 0& \mathrm{\Lambda }_p^{}& 0& 0\\ 0& 0& \mathrm{\Lambda }_p^+& 0\\ 0& 0& 0& \mathrm{\Lambda }_p^{}\end{array}\right]\gamma ^0\left[\begin{array}{cc}S_{bd\text{-}gs}^\mathrm{p}(p)& 0\\ 0& S_{gs\text{-}bd}^\mathrm{p}(p)\end{array}\right],\hfill \end{array}$$
(40)
where the antiparticle part is
$$\begin{array}{cc}\hfill S_{A\text{-}B}^\mathrm{a}(p)& =\frac{1}{\left(p_0+\stackrel{~}{\epsilon }_{A\text{-}B}\right)\left(p_0\stackrel{~}{\epsilon }_{B\text{-}A}\right)}\hfill \\ & \times \left[\begin{array}{cc}p_0(p+\mu _B)& i\mathrm{\Delta }_1\gamma _5\gamma ^0\\ i\mathrm{\Delta }_1\gamma _5\gamma ^0& p_0+(p+\mu _A)\end{array}\right],\hfill \end{array}$$
(41)
and the particle part is
$$\begin{array}{cc}\hfill S_{A\text{-}B}^\mathrm{p}(p)& =\frac{1}{\left(p_0+\epsilon _{A\text{-}B}\right)\left(p_0\epsilon _{B\text{-}A}\right)}\hfill \\ & \times \left[\begin{array}{cc}p_0+(p\mu _B)& i\mathrm{\Delta }_1\gamma _5\gamma ^0\\ i\mathrm{\Delta }_1\gamma _5\gamma ^0& p_0(p\mu _A)\end{array}\right].\hfill \end{array}$$
(42)
In this representation the structure of the Dirac indices becomes much simpler and is given by $`\mathrm{\Lambda }_p^\pm \gamma ^0`$ and $`\gamma _5\gamma ^0`$ attached with $`\mathrm{\Delta }`$. Then, it is important to notice that the $`(\mu ,\nu )`$ dependence in the integrand of (37) comes from the trace over the Dirac indices alone, that is
$`𝒯_{\text{diag(p-p)}}^{\mu \nu }(p;q)`$ $`=\text{tr}\left[\gamma ^\mu \mathrm{\Lambda }_{q+p}^+\gamma ^0\gamma ^\nu \mathrm{\Lambda }_p^+\gamma ^0\right],`$ (43)
$`𝒯_{\text{diag(p-a)}}^{\mu \nu }(p;q)`$ $`=\text{tr}\left[\gamma ^\mu \mathrm{\Lambda }_{q+p}^+\gamma ^0\gamma ^\nu \mathrm{\Lambda }_p^{}\gamma ^0\right]`$ (44)
for p-p and p-a loops made of the Nambu-Gor’kov diagonal components of the propagator (41) and (42). For the Nambu-Gor’kov off-diagonal components we have likewise;
$`𝒯_{\text{off(p-p)}}^{\mu \nu }(p;q)`$ $`=\text{tr}[\gamma ^\mu \mathrm{\Lambda }_{q+p}^+\gamma _5\gamma ^\nu \mathrm{\Lambda }_p^{}\gamma _5],`$ (45)
$`𝒯_{\text{off(p-a)}}^{\mu \nu }(p;q)`$ $`=\text{tr}[\gamma ^\mu \mathrm{\Lambda }_{q+p}^+\gamma _5\gamma ^\nu \mathrm{\Lambda }_p^+\gamma _5].`$ (46)
### II.1 Singularities in $`A_{1,2}`$
We will first discuss $`A_1`$ and $`A_2`$ that are degenerated. The generators in color, $`T_1`$ and $`T_2`$, have a non-vanishing component only between red and green and the flavor is kept unchanged at the vertices, so that the self-energy for $`A_{1,2}`$ stems from four diagrams;
$$\begin{array}{cc}\hfill \mathrm{\Pi }_{11}^{\mu \nu }& =\mathrm{\Pi }_{\text{diag(}rg;s\text{)}}^{\mu \nu }+\mathrm{\Pi }_{\text{diag(}rg;u\text{)}}^{\mu \nu }\hfill \\ & +\mathrm{\Pi }_{\text{diag(}rg;d\text{)}}^{\mu \nu }+\mathrm{\Pi }_{\text{off(}ru\text{-}gd\text{;}gu\text{-}rd\text{)}}^{\mu \nu }.\hfill \end{array}$$
(47)
Here, diag($`rg;s`$) means the quark loop composed of the diagonal component of $`rs`$ and $`gs`$ quarks (see Fig. 6) and diag($`rg;u`$) and diag($`rg;d`$) should be understood in the same manner. The last contribution, off($`ru`$-$`gd`$;$`gu`$-$`rd`$), represents the self-energy contribution from the off-diagonal components of $`ru`$-$`gd`$ and $`gu`$-$`rd`$ propagation, that is, the loop made up with $`ru`$ quarks turning into $`gd`$ and $`rd`$ quarks turning into $`gu`$ as shown in Fig. 7.
In the case of $`A_1`$ (and $`A_2`$ equivalently) only $`\mathrm{\Pi }_{\text{diag(}rg;s\text{)}}^{\mu \nu }`$ may contain gapless quarks in the gCFL phase in both lines of the loop diagram, and is expected to provide a singular contribution. As a matter of fact, the singular part presumably dominates over the screening mass behavior near the gapless onset as a function of $`M_s`$. This expectation is confirmed by the comparison between functional forms of only the $`rs`$-$`gs`$ polarization and the full result derived from the potential curvature we are dealing with in the next section. From Fig. 8, taking in the full results in advance, we can clearly observe that the physics near the gCFL onset is to be described by the $`rs`$-$`gs`$ excitation only.
The Debye screening mass, or $`\mathrm{\Pi }^{00}`$, has only the particle-particle part (and the antiparticle-antiparticle part that is negligible) because $`𝒯_{\text{diag(p-a)}}^{00}(p;0)=𝒯_{\text{off(p-a)}}^{00}(p;0)=0`$ and is written as
$$\begin{array}{cc}& \mathrm{\Pi }_{\text{diag(}rg;s\text{)}}^{00}(q0)\hfill \\ & =\frac{g^2}{8\pi ^2}𝑑pp^2𝒯_{\text{diag(p-p)}}^{00}(p;0)𝒰_{\text{diag(}rg;s\text{)}}(p),\hfill \end{array}$$
(48)
where $`𝒯_{\text{diag(p-h)}}^{00}(p;0)=2`$ and we defined
$$\begin{array}{cc}& 𝒰_{\text{diag(}rg;s\text{)}}=f_{\text{p-p}}[\epsilon _{bu\text{-}rs};bd,gs,bu]\hfill \\ & +f_{\text{p-p}}[\epsilon _{rs\text{-}bu};bd,gs,bu]+f_{\text{p-p}}[\epsilon _{bd\text{-}gs};bu,rs,bd]\hfill \\ & +f_{\text{p-p}}[\epsilon _{gs\text{-}bd};bu,rs,bd].\hfill \end{array}$$
(49)
The p-p integrand $`f_{\text{p-p}}`$ is explicitly given by
$$\begin{array}{cc}& f_{\text{p-p}}[\epsilon _{A\text{-}B};C,D,E]\hfill \\ & =\frac{\left[\epsilon _{A\text{-}B}+(p\mu _E)\right]\left[\epsilon _{A\text{-}B}+(p\mu _C)\right]}{\left[\epsilon _{A\text{-}B}+\epsilon _{B\text{-}A}\right]\left[\epsilon _{A\text{-}B}\epsilon _{C\text{-}D}\right]\left[\epsilon _{A\text{-}B}+\epsilon _{D\text{-}C}\right]}\mathrm{tanh}\left[\frac{\epsilon _{A\text{-}B}}{2T}\right]\hfill \end{array}$$
(50)
with the energy dispersion relation $`\epsilon _{A\text{-}B}`$ defined in (15).
In the CFL phase, as we have already mentioned, the $`bd`$-$`gs`$ and $`rs`$-$`bu`$ quark dispersion relations are identical and there are vanishing combinations in the denominator of (50) like $`\epsilon _{gs\text{-}bd}(q+p)\epsilon _{rs\text{-}bu}(p)0`$ as $`q0`$. The numerator goes to zero in the same limit of $`q0`$ and consequently it amounts to the derivative with respect to $`\epsilon `$ acting onto the distribution function, i.e., $`(/\epsilon _{A\text{-}B})\mathrm{tanh}[\epsilon _{A\text{-}B}/2T]`$. At zero temperature this is proportional to the delta function $`\delta (\epsilon _{A\text{-}B})`$ that takes a finite value only when gapless quarks (that is; $`\epsilon _{A\text{-}B}(p)=0`$ for some $`p`$) are present.
Using the solution (11) and (12) and assuming $`\mathrm{\Delta }_1=\mathrm{\Delta }_2=\mathrm{\Delta }`$ that is known to be a good approximation in the CFL phase, the polarization at zero temperature can simplify as
$$\begin{array}{cc}& \mathrm{\Pi }_{\text{diag(}rg;s\text{)}}^{00}\hfill \\ & =\frac{g^2}{4\pi ^2}𝑑pp^2\left\{\frac{[\epsilon (p\overline{\mu })]^2}{2\epsilon ^2}\delta \left(\epsilon \frac{M_s^2}{2\mu }\right)+\frac{\mathrm{\Delta }^2}{2\epsilon ^3}\right\}\hfill \end{array}$$
(51)
with $`\overline{\mu }=\mu M_s^2/6\mu `$ and $`\epsilon =\sqrt{(p\overline{\mu })^2+\mathrm{\Delta }^2}`$. This expression is valid only up to the onset where the gCFL starts. Actually, right at the onset where $`M_s^2/2\mu =\mathrm{\Delta }`$, the energy dispersion relation for $`gs`$-$`bd`$ and $`rs`$-$`bu`$ quarks is quadratic near the Fermi momentum and approximated as $`\epsilon (p)M_s^2/2\mu (p\overline{\mu })^2/2\mathrm{\Delta }`$. Then, the $`p`$-integration for the first term of (51) picks up the density of states at $`p=\overline{\mu }`$ and ends up with an infrared divergence. We would emphasize that this singular behavior originates from the electron pressure which ensures $`\mu _e=0`$. Otherwise, two quarks running through the loop diagram would not be gapless simultaneously at the gCFL onset.
In the CFL region before the gCFL occurs, it is rather the second term of (51) which is a major contribution, which can be further approximated as
$$\mathrm{\Pi }_{\text{diag(}rg;s\text{)}}^{00}=\frac{g^2\mu ^2}{4\pi ^2}=\frac{3}{2}m_g^2$$
(52)
with any logarithmic divergence neglected. The logarithmic terms are, in fact, given by $`\mathrm{\Delta }^2\mathrm{ln}[\mathrm{\Lambda }/\mathrm{\Delta }]`$ that is much smaller than the leading terms of order $`\overline{\mu }^2`$ and corrections are $``$1% for our choice of $`\mathrm{\Lambda }=800\mathrm{MeV}`$ and $`\mathrm{\Delta }_0=25\mathrm{MeV}`$. We will always drop any logarithmic terms $`\mathrm{ln}[\mathrm{\Lambda }/\mathrm{\Delta }]`$ throughout this paper.
Hence, the contribution to the Debye mass in the CFL side is obtained as
$$[m_D^2]_{\text{diag(}rg;s\text{)}}=\frac{3}{2}m_g^2.$$
(53)
Interestingly enough, this result exactly corresponds to one-flavor contribution (i.e. one-third) of (30) and does not have $`M_s`$-dependence until the gCFL onset.
We can immediately extend our discussion to the Meissner mass once if we analyze the ($`\mu `$,$`\nu `$) structure in the polarization. The p-p contribution to the Meissner mass is
$$\begin{array}{cc}& \frac{1}{2}\left(\delta _{ij}\widehat{q}_i\widehat{q}_j\right)𝒯_{\text{diag(p-p)}}^{ij}(p;0)\hfill \\ \hfill =& \frac{1}{2}\left(\delta _{ij}\widehat{q}_i\widehat{q}_j\right)\times 2\widehat{p}^i\widehat{p}^j=1(\widehat{q}\widehat{p})^2\hfill \\ \hfill & \frac{2}{3}=\frac{1}{3}𝒯_{\text{diag(p-p)}}^{00}(p;0).\hfill \end{array}$$
(54)
In the third line we made use of averaging over the angle integration that leads to one third, that is, $`d^3p(\widehat{q}\widehat{p})^2𝒰(p)=\frac{1}{3}d^3p𝒰(p)`$. In evaluating the Debye and Meissner masses $`𝒰_{\text{diag(}rg;s\text{)}}(p)`$ is common, so that from the above relation we reach the general relation,
$$[m_M^2]_{\text{diag(}rg;s\text{)}}^{\text{(p-p)}}=\frac{1}{3}[m_D^2]_{\text{diag(}rg;s\text{)}},$$
(55)
that is completely consistent with (34).
The p-a part generates ultraviolet divergent terms proportional to $`\mathrm{\Lambda }^2`$, $`\mathrm{\Delta }^2\mathrm{ln}[\mathrm{\Lambda }/\mathrm{\Delta }]`$, and $`(M_s^2/\mu )^2\mathrm{ln}[\mathrm{\Lambda }/\mathrm{\Delta }]`$ Alford . In the present work, as in the treatment of the Debye mass, we ignore any logarithmic divergences. We can rewrite $`f_{\text{p-p}}`$ into $`f_{\text{p-a}}`$ immediately, in view of the difference between (41) and (42), by replacing corresponding $`\epsilon _{A\text{-}B}`$ by $`\stackrel{~}{\epsilon }_{A\text{-}B}`$, and $`p\mu _A`$ by $`(p+\mu _A)`$. A similar analysis to (54) on the Dirac trace part leads to a factor as follows;
$$\begin{array}{cc}& \frac{1}{2}\left(\delta _{ij}\widehat{q}_i\widehat{q}_j\right)𝒯_{\text{diag(p-a)}}^{ij}(p;0)\hfill \\ \hfill =& \frac{1}{2}\left(\delta _{ij}\widehat{q}_i\widehat{q}_j\right)\times 2\left[\delta ^{ij}\widehat{p}^i\widehat{p}^j\right]=1+(\widehat{q}\widehat{p})^2\frac{4}{3}.\hfill \end{array}$$
(56)
Then, after some calculations along the same line as Refs. Huang:2004bg ; Rischke:2000qz with the approximation that $`\mathrm{\Delta }`$ and $`M_s`$ are all neglected, we acquire
$$\begin{array}{cc}& [m_M^2]_{\text{diag(}rg;s\text{)}}^{\text{(p-a)}}\hfill \\ \hfill =& \frac{1}{2}\underset{q0}{lim}(\delta _{ij}\widehat{q}_i\widehat{q}_j)[\mathrm{\Pi }_{\text{diag(}rg;s\text{)}}^{ij}]^{\text{(p-a)}}\hfill \\ \hfill & \frac{1}{2}m_g^2\frac{\mathrm{\Lambda }^2}{12\pi ^2},\hfill \end{array}$$
(57)
that is consistent with (35) divided by the flavor factor $`N_\mathrm{f}=3`$ up to the ultraviolet divergence.
After all, the sum over the p-p, p-a, and a-p parts finally amounts to
$$\begin{array}{cc}& [m_M^2]_{\text{diag(}rg;s\text{)}}=[m_M^2]_{\text{diag(}rg;s\text{)}}^{\text{(p-p)}}\hfill \\ & +2\times [m_M^2]_{\text{diag(}rg;s\text{)}}^{\text{(p-a)}}+\text{(subtraction)}=\frac{1}{2}m_g^2\hfill \end{array}$$
(58)
in the CFL phase, where “subtraction” is a term added by hand to remove the $`\mathrm{\Lambda }^2`$-term which should be renormalized. Eq. (57) indicates
$$\text{(subtraction)}=\frac{\mathrm{\Lambda }^2}{6\pi ^2}$$
(59)
per one polarization diagram of Nambu-Gor’kov diagonal particle-antiparticle excitations that should contain ultraviolet divergences.
### II.2 No singularities in $`A_{4,5}`$ and $`A_{6,7}`$
The analyses in the previous subsection give us information about the origin of singularities around the gCFL onset. When two quarks involved in the loop diagram becomes gapless at the same time, the divergent density of states compels us to face the imaginary Meissner mass. We can easily comprehend that this kind of singularity near the gapless onset is absent for $`A_{4,5}`$ and $`A_{6,7}`$.
The self-energies are diagrammatically decomposed as
$$\begin{array}{cc}\hfill \mathrm{\Pi }_{44}^{\mu \nu }& =\mathrm{\Pi }_{\text{diag(}br;s\text{)}}^{\mu \nu }+\mathrm{\Pi }_{\text{diag(}br;u\text{)}}^{\mu \nu }\hfill \\ & +\mathrm{\Pi }_{\text{diag(}br;d\text{)}}^{\mu \nu }+\mathrm{\Pi }_{\text{off(}bu\text{-}rs\text{;}ru\text{-}bs\text{)}}^{\mu \nu }\hfill \end{array}$$
(60)
and
$$\begin{array}{cc}\hfill \mathrm{\Pi }_{66}^{\mu \nu }& =\mathrm{\Pi }_{\text{diag(}gb;s\text{)}}^{\mu \nu }+\mathrm{\Pi }_{\text{diag(}gb;u\text{)}}^{\mu \nu }\hfill \\ & +\mathrm{\Pi }_{\text{diag(}gb;d\text{)}}^{\mu \nu }+\mathrm{\Pi }_{\text{off(}gd\text{-}bs\text{;}bd\text{-}gs\text{)}}^{\mu \nu }.\hfill \end{array}$$
(61)
None of these above polarizations is composed of would-be gapless quarks only, i.e. $`bd`$-$`gs`$ and $`rs`$-$`bu`$ quarks only. Therefore, we can conclude at this stage of analyses that $`A_{4,5}`$ and $`A_{6,7}`$ do not exhibit chromomagnetic instability at least in the vicinity of the gCFL onset.
### II.3 Singularities in $`A_{3,8,\gamma }`$
For color-diagonal gluons $`A_3`$, $`A_8`$, and photon $`A_\gamma `$, neither color nor flavor is changed at the vertices. This means that 9 diagonal and 6 off-diagonal diagrams are possible. For example, the self-energy for $`A_8`$, for instance, has two types of singularities; one is associated with the $`bd`$-$`gs`$ gapless quark dispersion relation,
$$\begin{array}{cc}\hfill \left[\mathrm{\Pi }_{88}^{\mu \nu }\right]_{bd\text{-}gs}& =\frac{1}{3}\mathrm{\Pi }_{\text{diag(}gg;s\text{)}}^{\mu \nu }+\frac{4}{3}\mathrm{\Pi }_{\text{diag(}bb;d\text{)}}^{\mu \nu }\hfill \\ & \frac{2}{3}\mathrm{\Pi }_{\text{off(}bd\text{-}gs\text{;}bd\text{-}gs\text{)}}^{\mu \nu },\hfill \end{array}$$
(62)
and the other is associated with the $`rs`$-$`bu`$ gapless quark (almost quadratic) dispersion relation,
$$\begin{array}{cc}\hfill \left[\mathrm{\Pi }_{88}^{\mu \nu }\right]_{rs\text{-}bu}& =\frac{1}{3}\mathrm{\Pi }_{\text{diag(}rr;s\text{)}}^{\mu \nu }+\frac{4}{3}\mathrm{\Pi }_{\text{diag(}bb;u\text{)}}^{\mu \nu }\hfill \\ & \frac{2}{3}\mathrm{\Pi }_{\text{off(}rs\text{-}bu\text{;}rs\text{-}bu\text{)}}^{\mu \nu }.\hfill \end{array}$$
(63)
The one-loop diagrams corresponding to $`\mathrm{\Pi }_{\text{diag(}bb;d\text{)}}^{\mu \nu }`$ and $`\mathrm{\Pi }_{\text{off(}bd\text{-}gs\text{;}bd\text{-}gs\text{)}}^{\mu \nu }`$ are shown in Figs. 9 and 10 to take some examples.
For completeness we shall enumerate the rest here. The non-singular contributions to the $`A_8`$ polarization are
$$\begin{array}{cc}& \left[\mathrm{\Pi }_{88}^{\mu \nu }\right]_{\text{non-singular}}=\frac{1}{3}\mathrm{\Pi }_{\text{diag(}rr;u\text{)}}^{\mu \nu }+\frac{1}{3}\mathrm{\Pi }_{\text{diag(}rr;d\text{)}}^{\mu \nu }\hfill \\ & +\frac{1}{3}\mathrm{\Pi }_{\text{diag(}gg;u\text{)}}^{\mu \nu }+\frac{1}{3}\mathrm{\Pi }_{\text{diag(}gg;d\text{)}}^{\mu \nu }+\frac{4}{3}\mathrm{\Pi }_{\text{diag(}bb;s\text{)}}^{\mu \nu }\hfill \\ & +\frac{1}{3}\mathrm{\Pi }_{\text{off(}gu\text{-}rd\text{;}gu\text{-}rd\text{)}}^{\mu \nu }+\frac{1}{3}\mathrm{\Pi }_{\text{off(}ru\text{-}gd\text{;}ru\text{-}gd\text{)}}^{\mu \nu }\hfill \\ & \frac{2}{3}\mathrm{\Pi }_{\text{off(}gd\text{-}bs\text{;}gd\text{-}bs\text{)}}^{\mu \nu }\frac{2}{3}\mathrm{\Pi }_{\text{off(}bs\text{-}ru\text{;}bs\text{-}ru\text{)}}^{\mu \nu },\hfill \end{array}$$
(64)
The diagonal contributions, as shown in Fig. 9, are easily available from the previous results only by replacing the color-flavor indices appropriately and adding $`\frac{1}{2}`$ to adjust the combinatorial factor. It should be noted that instability occurs in this case whenever any of quarks becomes gapless because two virtually excited quarks are identical in the color-diagonal channels.
The off-diagonal expression is also deduced simply by means of proper arrangement of the color-flavor indices and replacement of the numerator of (50) as $`\epsilon _{A\text{-}B}\pm (p\mu _B)i\mathrm{\Delta }_{AB}`$. Because $`\epsilon _{A\text{-}B}\pm (p\mu _B)\mathrm{\Delta }_{AB}`$ in the vicinity of the Fermi surface $`p\mu `$ that is dominating over the $`p`$-integration, the functional form of the off-diagonal contribution is essentially the same as the diagonal contributions. There are two important features to remark here, though.
First, the self-energy contributions coming from the off-diagonal loop are free from ultraviolet quadratic divergence even in p-a excitations. It is because the off-diagonal propagator is proportional to, not $`\epsilon _{A\text{-}B}\pm (p\mu _B)`$, but $`\mathrm{\Delta }_{AB}`$ that is not rising with increasing $`p`$. Thus, the subtraction term proportional to $`\mathrm{\Lambda }^2`$ is no longer necessary.
Second, from the trace over the Dirac indices (45) and (46) we notice that
$`𝒯_{\text{off}}^{00}`$ $`=𝒯_{\text{diag}}^{00},`$ (65)
$`𝒯_{\text{off}}^{ij}`$ $`=𝒯_{\text{diag}}^{ij}`$ (66)
for both p-p and p-a cases, that means the diagonal and off-diagonal diagrams contribute to the Debye and Meissner masses differently. Using (65), (66) and from a similar analysis to (54), we immediately see the relation,
$$[m_M^2]_{\text{off(}bd\text{-}gs\text{;}bd\text{-}gs\text{)}}=\frac{1}{3}[m_D^2]_{\text{off(}bd\text{-}gs\text{;}bd\text{-}gs\text{)}},$$
(67)
that is consistent with (36). The sign difference between (55) and (67) originates from the negative sign in (65). The p-a off-diagonal contribution to the Meissner mass squared turns out to be negligible.
In addition, the negative sign difference between (65) and (66) can explain the results for the “neutral” and “charged” condensates discussed in a two-species model in Ref. Alford:2005qw . As we have mentioned above, the functional form of $`\mathrm{\Pi }_{\text{diag(}gg;s\text{)}}^{\mu \nu }`$, $`\mathrm{\Pi }_{\text{diag(}bb;d\text{)}}^{\mu \nu }`$, and $`\mathrm{\Pi }_{\text{off(}bd\text{-}gs\text{;}bd\text{-}gs\text{)}}^{\mu \nu }`$ are almost identical up to the combinatorial factor $`\frac{1}{2}`$. Let’s consider the “neutral” case fictitiously following the argument in Ref. Alford:2005qw , for which we assume that $`bd`$ and $`gs`$ quarks had the opposite charge corresponding to the two-species model composed of $`bd`$ and $`gs`$. Then, the polarization related to the screening masses would be $`\mathrm{\Pi }^{\mu \nu }=\mathrm{\Pi }_{\text{diag(}gg;s\text{)}}^{\mu \nu }+\mathrm{\Pi }_{\text{diag(}bb;d\text{)}}^{\mu \nu }\mathrm{\Pi }_{\text{off(}bd\text{-}gs\text{;}bd\text{-}gs\text{)}}^{\mu \nu }`$. From (65) we see that the Debye mass positively diverges at the gapless onset, while the singularities in the Meissner mass cancel to vanish. In the “charged” case, on the other hand, we assume that $`bd`$ and $`gs`$ quarks had the same charge. In the same way, the polarization leading to the screening masses would be $`\mathrm{\Pi }^{\mu \nu }=\mathrm{\Pi }_{\text{diag(}gg;s\text{)}}^{\mu \nu }+\mathrm{\Pi }_{\text{diag(}bb;d\text{)}}^{\mu \nu }+\mathrm{\Pi }_{\text{off(}bd\text{-}gs\text{;}bd\text{-}gs\text{)}}^{\mu \nu }`$. Obviously, the Debye mass has no singularity at the gapless onset, but the Meissner mass negatively diverges. This is exactly what was found in Ref. Alford:2005qw . In QCD, however, the diquark condensate is actually a mixture of the “neutral” and “charged” condensates and the Debye and Meissner masses both diverge positively and negatively from the common origin at the gCFL onset.
As a final remark of this section, we shall draw attention to the contrasting difference between the singular parts (62) related to $`bd`$-$`gs`$ quarks and (63) related to $`rs`$-$`bu`$ quarks. In general, the Meissner mass negatively diverges on the $`bd`$-$`gs`$, $`rs`$-$`bu`$, and $`gu`$-$`rd`$ boundary lines drawn in Fig. 4. When $`M_s`$ grows larger, as seen in Fig. 5, the blocking region of the $`bd`$-$`gs`$ dispersion relation becomes wider and the singular character in (62) weakens. In contrast, $`\mu _e`$ is kept to be very close to the $`rs`$-$`bu`$ line in Fig. 4 for all $`M_s`$ and in all the gCFL region (63) remains to provide a negatively huge contribution to the Meissner mass. If it were not for electrons, $`\mu _e`$ would be within the bandgap and (63) would not be singular at all. In this sense, the singularity of (63) is induced by the presence of the electron terms which force $`\mu _e`$ to stay in the singular region slightly below the dashed line in Fig. 4. To evaluate that huge contribution correctly from (63), it is essential to resolve a tiny blocking region $`\delta p\mu _e^3/3\overline{\mu }^2`$ with a great accuracy, which is next to impossible technically.
Although the strength of singularities (62) and (63) at the gCFL onset is exactly equal at $`T=0`$, the temperature dependence is drastically different; as seen from Fig. 11, (63) is no longer singular even at $`T=0.5\mathrm{MeV}`$. This is because at finite temperature $`rs`$-$`bu`$ quarks not from the blocking region but from thermal excitations can cancel the $`\stackrel{~}{Q}`$-charge driven by electrons. There needs not to be a blocking region associated with $`rs`$-$`bu`$ quarks therefore. Actually, the $`rs`$-$`bu`$ quark density is estimated as $`\sqrt{2\pi \mathrm{\Delta }T}\overline{\mu }^2/\pi ^2`$ when its dispersion relation is quadratic, and it can balance out the electron density $`\mu _e^3/3\pi ^2`$ even at $`T=(\mu _e^3/3\overline{\mu })^2/2\pi \mathrm{\Delta }`$ that is as tiny as $``$ eV for $`\mu _e`$ a few MeV. Consequently there is no gapless $`rs`$-$`bu`$ quarks any more already at tiny temperatures.
## III potential curvature
In order to derive the information on the full evaluation of screening masses, let us elaborate a numerical method based on the previous analytical consideration to compute the Debye and Meissner masses here. It is well-known that the Debye mass can be expressed as the potential curvature as follows kapusta ;
$$m_{D,\alpha \beta }^2=\frac{^2\mathrm{\Omega }_\mu }{\mu _\alpha \mu _\beta },$$
(68)
where $`\mu `$’s are the color chemical potential coupled to the generators $`T_\alpha `$. That means, $`\mu `$’s here involve off-diagonal components in color-flavor space besides (5). The thermodynamic potential defined with all $`\mu `$’s is denoted by $`\mathrm{\Omega }_\mu `$ here. The potential curvature (68) is evaluated at vanishing $`\mu `$’s except $`\mu _3`$ and $`\mu _8`$ which are fixed by the color neutrality conditions.In QCD at finite temperature, moreover, pure gluonic loops produce a Debye mass $`gT`$, which is not included in Eq. (68).
All we have to know to evaluate the thermodynamic potential are the energy dispersion relations that are much easier than calculating the loop integral (37) directly. We show the Debye masses as a function of $`M_s^2/\mu `$ inferred from (68) in Figs. 12 and 13. The singular behavior associated with the gapless onset emerges in an upward direction, that is accounted for from the relation (34). In other words, through the relation (34), the positively rising results for the Debye masses squared imply negative Meissner masses squared for $`A_{1,2}`$ and $`A_{3,8,\gamma }`$ around the gCFL onset and for all the gluons for large $`M_s^2/\mu `$. It is worth noting that the $`\stackrel{~}{Q}`$-photon has a finite Debye mass for large $`M_s^2/\mu `$ because the system in the gCFL phase has a nonzero density of electrons that can screen $`\stackrel{~}{Q}`$-charge.
Once we develop a relation connecting the thermodynamic potential to the Meissner mass like the above relation (68) with respect to the Debye mass, we can numerically compute the Meissner mass including not only the singular parts but also all the contributions. As a matter of fact, we have found a relation,
$$(m_{M,\alpha \beta }^2)_{\text{bare}}=\frac{1}{3}\underset{i=1}{\overset{3}{}}\frac{^2\mathrm{\Omega }_A}{A_{\alpha ,i}A_{\beta ,i}}|_{A=0},$$
(69)
where “bare” means the Meissner mass squared with terms diverging like $`\mathrm{\Lambda }^2`$, and $`\mathrm{\Omega }_A`$ is the thermodynamic potential defined in the presence of the gauge fields $`𝑨_i`$ (Note that bold symbols stand for matrices in color-flavor space in our notation.) in the Lagrangian (3) as
$$_A=\overline{\psi }(i\text{/}+𝝁_{\text{eff}}\gamma ^0\stackrel{}{𝑨}\stackrel{}{\gamma })\psi +_{\text{int}}.$$
(70)
Equation (69) needs some more explanation. It should be worth mentioning that (69) is a well-known general relation if we do accomplish the angle integration with respect to the loop momentum $`𝐩`$. However, this is technically difficult, for $`𝒑`$ makes an angle with $`𝑨_i`$ on evaluating $`\mathrm{\Omega }_A`$. The benefit of (69) is that we can reduce the $`p`$-integration into a one-dimensional $`p_3`$-integration as if the rotational symmetry were not affected by $`𝑨_i`$. We shall take advantage of this simplification at the price of treating $`M_s`$ as chemical potential shifts.
This relation can be proved in the following way. When the momentum $`𝒑`$ in the loop-integral is chosen to be in the $`3`$-direction, we can show that,
$$\begin{array}{cc}& 𝒯_{\text{diag(p-p)}}^{11}(p;0)=𝒯_{\text{diag(p-p)}}^{22}(p;0)=0,\hfill \\ & 𝒯_{\text{diag(p-p)}}^{33}(p;0)=2,\hfill \\ & 𝒯_{\text{diag(p-a)}}^{11}(p;0)=𝒯_{\text{diag(p-a)}}^{22}(p;0)=2,\hfill \\ & 𝒯_{\text{diag(p-a)}}^{33}(p;0)=0,\hfill \end{array}$$
(71)
indicating that $`\mathrm{\Pi }^{33}`$ contains only the p-p contribution and $`\mathrm{\Pi }^{11}`$ and $`\mathrm{\Pi }^{22}`$ the p-a contribution. From these expressions in the same way as (54) we eventually arrive at the relation,
$$m_{M,\alpha \beta }^2=\frac{1}{3}\left[\mathrm{\Pi }_{\alpha \beta }^{11}(0,0)+\mathrm{\Pi }_{\alpha \beta }^{22}(0,0)+\mathrm{\Pi }_{\alpha \beta }^{33}(0,0)\right].$$
(72)
We should remark that, even though we do not know an explicit expression for the $`ru`$-$`gd`$-$`bs`$ quark propagator, the ($`\mu ,\nu `$) structure comes out only through the energy projection operators and the above argument is generally valid for any quark propagator as long as $`M_s`$ can be treated as an effective chemical potential shift. The relation (72) is readily rewritten in terms of the thermodynamic potential, which leads us to (69).
In our model study, the four-fermion interaction is non-renormalizable and $`\mathrm{\Omega }_A`$ suffers ultraviolet divergent terms proportional to $`\mathrm{\Lambda }^2`$. Thus the Meissner masses derived directly from the potential curvature should have those divergences to be subtracted properly by hand. As argued in the previous analytical study, we conjecture that the subtraction term is simply inferred by counting divergent diagrams, which results in;
$`m_{M,\alpha \alpha }^2`$ $`=(m_{M,\alpha \alpha }^2)_{\text{bare}}+3\times {\displaystyle \frac{g^2\mathrm{\Lambda }^2}{6\pi ^2}},(1\alpha 8)`$ (73)
$`m_{M,\gamma \gamma }^2`$ $`=(m_{M,\gamma \gamma }^2)_{\text{bare}}+4\times {\displaystyle \frac{e^2\mathrm{\Lambda }^2}{6\pi ^2}},`$ (74)
$`m_{M,\alpha \beta }^2`$ $`=(m_{M,\alpha \beta }^2)_{\text{bare}},(1\alpha \beta 8)`$ (75)
$`m_{M,\alpha \gamma }^2`$ $`=(m_{M,\alpha \gamma }^2)_{\text{bare}}.`$ (76)
These are all determined simply from the diagrammatic counting of the Nambu-Gor’kov diagonal p-a loops and the associated subtraction (59) with proper coefficients. In the case of $`A_8`$ for example, the divergence factor in unit of $`\mathrm{\Lambda }^2/6\pi ^2`$ coming from (62) is $`\frac{1}{2}\times (\frac{1}{3}+\frac{4}{3})=\frac{5}{6}`$ where $`\frac{1}{2}`$ is a combinatorial factor. In the same way (63) and (64) give $`\frac{5}{6}`$ and $`\frac{8}{6}`$ respectively. The sum of these results in $`3\times (\mathrm{\Lambda }^2/6\pi ^2)`$ as in (73). It is interesting to see that the divergent terms cancel each other in the color-mixing and photon-mixing channels. From these above expressions we have obtained the results as presented in Figs. 14, 15, 16, and 17.
Now, before addressing the instability problem, let us compare our numerical results shown in Figs. 14, 15, 16, and 17 to the known analytical estimates. In view of our results at $`M_s=0`$, all the gluons except $`A_8`$ have the identical Meissner mass $`0.892m_g^2`$ which is consistent with (25). We have confirmed that the agreement becomes better for larger $`\mu \mathrm{\Delta }`$ because the analytical estimates require this. The mixing channel with $`e/g=1/2`$ (chosen presumably larger than realistic values to make the mixing effects apparently visible) results in $`1.18m_g^2`$ for $`\stackrel{~}{A}_8`$ and zero for $`\stackrel{~}{A}_\gamma `$, which is also consistent with the analytic estimate $`\stackrel{~}{m}_{M,88}^2=1.15m_g^2`$ from (27). All these suggest that our method works well.
Figure 14 shows the Meissner mass squared for $`A_{1,2}`$ as a function of $`M_s^2/\mu `$. The vertical dotted line corresponds to the gCFL onset and is located near the onset condition given kinematically $`(M_s^2/\mu )_\mathrm{c}2\mathrm{\Delta }`$. (We chose the model parameters to yield $`\mathrm{\Delta }_0=25\mathrm{MeV}`$ for $`M_s=T=0`$.) We see that the Meissner mass squared is negative in the entire gCFL region at zero temperature, that means chromomagnetic instability. The zero temperature behavior for $`A_{1,2}`$ is consistent with the results reported in Ref. Casalbuoni:2004tb except for the asymptotic property that the Meissner mass in Fig. 14 does not approach zero for large $`M_s^2/\mu `$ as it does in Ref. Casalbuoni:2004tb . The singularity around the gCFL onset is drastically smeared by the temperature effect as seen in the curve for $`T=0.5\mathrm{MeV}`$, though chromomagnetic instability still exists for $`M_s^2/\mu >84\mathrm{MeV}`$. At higher temperatures, as shown by the $`T=5\mathrm{MeV}`$ curve, the Meissner screening mass smoothly goes to zero with increasing $`M_s^2/\mu `$ since the system then comes to reach the normal phase without any first-order phase transition. (The $`T=5\mathrm{MeV}`$ results in fact have a tiny unstable region at $`M_s^2/\mu =125\mathrm{MeV}`$, see Fig. 1.)
In Fig. 15 (and in Fig. 16) we plot the results for $`A_{4,5}`$ (and for $`A_{6,7}`$ respectively). The gross features are in good agreement with the results presented in Ref. Casalbuoni:2004tb and all these gluons do not suffer chromomagnetic instability until $`M_s^2/\mu >105\mathrm{MeV}`$ at zero temperature. We remark that the instabilities for $`M_s^2/\mu 100\mathrm{MeV}`$ persist even at $`T=5\mathrm{MeV}`$, but it disappears at $`T=10\mathrm{MeV}`$.
The Meissner masses corresponding to $`A_3`$, $`A_8`$, and $`A_\gamma `$, as we have exposited before, are negatively divergent if any of quark excitation energies take the quadratic form as a function of the momentum. This situation takes place only at the gCFL onset. The results inside the gCFL regime are negatively large but not divergent actually. In fact a finite density of electrons requires an almost but not exactly quadratic dispersion relation anywhere in the gCFL phase. The outputs accordingly have severe dependence on the precise form of the almost quadratic dispersion relation, that is, the deviation from the quadratic form, which is hard to evaluate numerically with an extreme accuracy. We have confirmed the above mentioned behavior quantitatively in our calculations, that is, negatively large Meissner masses squared for $`A_3`$, $`A_8`$, and $`A_\gamma `$ in the entire gCFL region, but we could not reduce tremendous numerical ambiguity. Instead of tackling numerical difficulties that would be milder at finite temperature, we have calculated the Meissner masses at $`T=0.5\mathrm{MeV}`$. In Fig. 17 the dotted curves are the masses without mixing taken into account and the solid curves are the eigenvalues of the $`3\times 3`$ mass squared matrix. We clearly see that instability certainly occurs in one or two of three eigenmodes. Furthermore, we note that any effect of the almost quadratic dispersion relation cannot be perceived from the $`T=0.5\mathrm{MeV}`$ results because it is already smooth enough at that temperature as is obvious also in Fig. 11.
The unstable regions read from these figures are to be considered as slices with $`T`$ fixed on the phase diagram. It is instructive to compile all these slices and map respective unstable regions onto the phase diagram. Figure 18 is one example to locate the phase boundary and the unstable region and to see them on a common basis. In this way we clarified the unstable regions on the phase diagram as we have already discussed at the beginning of this paper.
## IV summary and open questions
In summary, in dense neutral three-flavor quark matter, we figured out the computational procedure to derive the Debye and Meissner screening masses from curvatures of the thermodynamic potential with respect to gluon source fields. We investigated chromomagnetic instability by changing the temperature $`T`$ and the Fermi momentum mismatch $`M_s^2/\mu `$ and explored the unstable regions for respective gluons and photon on the phase diagram involving the (g)CFL, uSC, and (g)2SC phases.
We can say that the instability lines we added on the renewed phase diagram tell us, at least, at what $`\mu `$ or $`M_s`$ the homogeneous (g)CFL, uSC, and (g)2SC phases remain to be stable at a certain $`T`$. From the present work, however, we cannot extract any clue about what is going on inside the unstable regions. One possible solution would be a crystalline-like phase as is suggested in the two-flavor case Giannakis:2004pf . If this is true also in three-flavor quark matter, it is most likely that the instability lines are to be regarded as the boundaries that separate the homogeneous and crystalline superconducting phases. In this sense, the boundaries drawn by thick lines in Figs. 1, 2, and 3 have a definite physical meaning even after a true ground state inside the unstable regions will be identified. Strictly speaking, once chromomagnetic instability occurs in any gluon channel, it would affect the other gluon channels. Consequently instability lines with respect to one gluon going inside the unstable region with respect to another gluon might be substantially different from what we have seen in this work. They are actually the most outer boundaries only that separate the homogeneous and crystalline phases, as shown in Fig. 19.
Our approach here is essentially along the same line as in Ref. Fukushima:2004zq , and so this work leaves the same open questions enumerated in Ref. Fukushima:2004zq ; how gauge field fluctuations affect the critical point and the order of respective phase transitions Matsuura:2003md , what the nature of $`K^0`$-condensation in the gCFL phase is and how it affects instability Kryjevski:2004jw ; Buballa:2004sx , what difference the ’t Hooft (instanton) interaction makes that induces an interaction like $`|\mathrm{\Delta }_3|^2M_s`$, where the chiral phase transition is located and how our ($`M_s^2/\mu `$-$`T`$) phase diagram is mapped onto the ($`\mu `$-$`T`$) plane with $`M_s`$ solved dynamically.
In addition to these above mentioned issues, we should make an improvement how to take account of $`M_s`$ fully. The derivation of our formula (69) needs the energy projection operator for massless quarks and works out as long as the $`M_s`$ effects can be well-approximated by a chemical potential shift. From this reason, we could not augment the phase diagrams for stronger couplings (Figs. 16 and 17 in Ref. Fukushima:2004zq for example) with the instability lines. The improvement could be implemented by the full angle-integration in $`𝒑`$ to acquire $`\mathrm{\Omega }_A`$.
From the theoretical point of view, logarithmic divergences that have been simply neglected in the present work must be taken seriously. In particular, logarithmic terms depending on $`M_s`$ seem to remain even in the normal phase Alford , while in our work the Meissner masses turn out to be zero properly in the normal phase, suggesting some cancellation mechanism between $`M_s`$ dependent divergences. This must be further investigated in a field-theoretical way. Another problem of theoretical interest is how much singularities are inhibited by non-perturbative resummation. However small the coupling constant $`g`$ is, the polarization is divergently large at the one-loop level, and generally, that demands a sort of resummation, for instance, by means of the ladder approximation. The instability might not be cured itself by resummation, but the singularities near the gCFL onset would be possibly smoothened.
Phenomenologically, since not only the gCFL but also the g2SC phase is within our scope especially once we will be able to approach the stronger coupling case, it is inevitable to examine the competition between the mixed and crystalline phases around the g2SC region. The entanglement between these phase possibilities should play an important role to complement the phase diagram for the large $`M_s`$ or moderate density region iida .
###### Acknowledgements.
The author acknowledges helpful conversations with M. Alford, M. Forbes, C. Kouvaris, M. Huang, K. Iida, K. Rajagopal, A. Schmitt, and I. Shovkovy. This work was supported in part by RIKEN BNL Research Center and the U.S. Department of Energy under cooperative research agreement #DE-AC02-98CH10886. |
warning/0506/gr-qc0506095.html | ar5iv | text | # Phantom cosmology with general potentials
## 1 Introduction
There is now substantial agreement that the expansion of the universe is accelerated, with evidence from type Ia supernovae , WMAP data , and large scale structure surveys . Two classes of models aim at explaining the observed cosmic acceleration: one modifies gravity on large scales by introducing corrections to the Einstein-Hilbert or the Palatini action (often there are instabilities or incorrect post–Newtonian limits ), and the other advocates the existence of dark energy with density $`\rho `$ and exotic pressure $`P<\rho /3`$.
Most dark energy candidates proposed thus far are scalar fields; there is marginal evidence for an effective equation of state parameter of the dark energy $`wP/\rho <1`$, which is equivalent to increasing Hubble parameter $`\dot{H}>0`$. If confirmed, this measurement is important because an increasing $`H`$ cannot be explained by Einstein gravity with a canonical scalar field . If the universe really superaccelerates it may end its existence in a finite time in an explosive expansion of the scale factor accompanied by diverging dark energy density, called a Big Rip . Other kinds of singularities at a finite time in the future discussed in the literature are called “sudden future singularities” and are characterized by finite scale factor and Hubble parameter, and diverging $`\dot{H}`$ and dark pressure .
Simple models of a superaccelerated universe employ a scalar field coupled nonminimally to the Ricci curvature or a phantom field, i.e., a minimally coupled scalar field with negative kinetic energy . Phantom fields and fields with non–canonical energy are present in string theories and supergravity , are associated to bulk viscosity due to particle production , or arise in higher order theories . It has also been proposed that early inflation and late time acceleration of the universe can be unified in a single theory based on a phantom field . A phantom field poses several challenges: it is subject to severe instabilities, which may perhaps be avoided by thinking of the phantom as an effective field theory resulting from some fundamental theory with positive energies . A fundamental quantum phantom is very difficult to stabilize .
In general, a form of dark energy with $`w<1`$ is worrisome because it violates the energy conditions cherished by most physicists and opens the door to disturbing exotica. For example, even a small amount of exotic matter violating the weak energy condition may be sufficient to open up a wormhole and make time travel possible . On the other hand, even a simple classical scalar field coupled nonminimally to the Ricci curvature may violate all of the energy conditions , and one should probably not be too rigid in rejecting phantom fields a priori.
The peculiar dynamics arising from a negative kinetic energy density has been studied with the help of a toy model consisting of two mutually coupled oscillators, one with negative kinetic energy representing the phantom field, and the other with positive kinetic energy representing the gravitational field . However, this toy model fails to properly describe the dynamics and it is preferable to study the actual dynamical system. Furthermore, a number of papers in the literature study phantom models with different specific potentials . Here we study the dynamics without toy models or approximations and we deduce the behaviour and the late time state of a phantom–dominated universe without specifying the form of the potential. The assumption of a spatially flat Friedmann–Lemaitre–Robertson–Walker (FLRW) universe plus general assumptions on the potential (such as boundedness, monotonicity, etc.) allow one to derive the asymptotic dynamics. In addition to the need to understand phantom dynamics with general potentials, our study is motivated by the question whether the universe will end in a Big Rip or will expand forever. In fact, a universe that superaccelerates may end its existence in a Big Rip or it may expand forever, depending on the dynamical equations. Often authors working with a specific choice of the phantom potential find a de Sitter regime as the final stage of evolution (late time attractor). We show that certain features of the phantom potential determine whether the final state is an indefinite expansion or a Big Rip, and the analysis needs not be repeated for all the possible potentials. We adopt the notations and conventions of and units such that $`c=G=\mathrm{}=1`$.
## 2 The phase space of phantom cosmology
Based on the recent cosmological observations we adopt the FLRW line element
$$ds^2=dt^2+a^2(t)\left(dx^2+dy^2+dz^2\right)$$
(2.1)
in comoving coordinates $`(t,x,y,z)`$. We consider the situation in which dark energy has already come to dominate the cosmic dynamics, hence a phantom field is the only form of matter in the Einstein equations. The energy density and pressure of the phantom are, respectively,
$$\rho =\frac{1}{2}\dot{\varphi }^2+V(\varphi ),P=\frac{1}{2}\dot{\varphi }^2V(\varphi ).$$
(2.2)
The distinguishing negative kinetic energy is evident in eqs. (2.2) which exhibit the “wrong” sign for the kinetic term $`\dot{\varphi }^2/2`$. The Einstein equations for $`\varphi (t)`$ and the scale factor $`a(t)`$ are
$`H^2={\displaystyle \frac{\kappa }{6}}\left(\dot{\varphi }^2+2V\right),`$ (2.3)
$`\dot{H}+H^2={\displaystyle \frac{\kappa }{3}}\left(\dot{\varphi }^2+V\right),`$ (2.4)
$`\ddot{\varphi }+3H\dot{\varphi }{\displaystyle \frac{dV}{d\varphi }}=0,`$ (2.5)
where $`H\dot{a}/a`$ is the Hubble parameter, an overdot denotes differentiation with respect to the comoving time $`t`$, and $`\kappa 8\pi G`$. The sign of the term $`dV/d\varphi `$ in the Klein–Gordon equation (2.5) is the opposite of the usual one. The potential is required to be positive by eq. (2.3), and this modification of the usual Klein–Gordon equation implies that a phantom falls up in an increasing potential, is repelled by a minimum, and settles in a maximum of $`V`$ . In addition, a monotonically increasing unbounded potential may be expected to generate runaway solutions – the analogy with a canonical scalar field in a negative potential not bounded from below applies (these properties are proven later in this paper).
Eq. (2.3) implies that the phantom energy density is non–negative for any solution of eqs. (2.3)–(2.5). Only two equations in the set (2.3)–(2.5) are independent – when $`\dot{\varphi }0`$ one can derive the Klein–Gordon equation (2.5) from the other two.
The field equations can be obtained from the Lagrangian $`L`$ or the Hamiltonian $``$
$$L=3a\dot{a}^2+\kappa a^3\left[\frac{\dot{\varphi }^2}{2}+V(\varphi )\right],=a^3\left[H^2+\frac{\kappa }{6}\dot{\varphi }^2\frac{\kappa V(\varphi )}{3}\right],$$
(2.6)
the Hamiltonian constraint (2.3) corresponding to $`=0`$. The orbits of the solutions of eqs. (2.3)–(2.5) are constrained to the surface of constant energy $`=0`$ in the phase space $`(a,\dot{a},\varphi ,\dot{\varphi })`$. We choose as dynamical variables the Hubble parameter and the scalar field $`(H,\varphi )`$ which correspond to physical observables<sup>1</sup><sup>1</sup>1Other authors choose as dynamical variables various combinations of $`H`$ and $`\varphi `$ but this obscures the physical interpretation of their results.. The phase space accessible to the orbits of the solutions is the two-dimensional energy surface $`=0`$ in the three-dimensional space $`(H,\varphi ,\dot{\varphi })`$; this surface is in general curved (the appendix provides an example).
Once the values of $`H`$ and $`\varphi `$ are given, one computes the corresponding value(s) of $`\dot{\varphi }`$ by rewriting the Hamiltonian constraint (2.3) as the algebraic equation for $`\dot{\varphi }`$
$$\dot{\varphi }^22V+\frac{6H^2}{\kappa }=0,$$
(2.7)
with solutions
$$\dot{\varphi }_\pm (H,\varphi )=\pm \sqrt{2\left[V(\varphi )\frac{3H^2}{\kappa }\right]}\pm \sqrt{F(H,\varphi )}.$$
(2.8)
Depending the form of $`V\left(\varphi \right)`$, there may be a region $``$ of the phase space forbidden to the dynamics and corresponding to a negative sign of $`F2\left(V3H^2/\kappa \right)`$:
$$\{(H,\varphi ,\dot{\varphi }):V(\varphi )<\frac{3H^2}{\kappa }\}.$$
(2.9)
When $`F(H,\varphi )>0`$ there are two distinct values of $`\dot{\varphi }`$ for each value of $`H`$ and $`\varphi `$, corresponding to the fact that the two-dimensional phase space consists of two sheets joining each other only at the points of the set
$$\{(H,\varphi ,\dot{\varphi }):V(\varphi )=\frac{3H^2}{\kappa },\dot{\varphi }=0\},$$
(2.10)
which constitutes the boundary of the forbidden region $``$ and lies in the $`(H,\varphi )`$ plane. We denote the sheet corresponding to the positive sign of $`\dot{\varphi }`$ “upper sheet”, while the other is called “lower sheet”. This structure of the phase space is general in scalar–tensor cosmology and is well known in the theory of a scalar field with canonical kinetic energy coupled non–minimally to the Ricci curvature . Phantom cosmology can be seen as a special case of scalar–tensor theory for which there is always a forbidden region in the phase space, while this needs not be the case in other scalar–tensor theories. The equilibrium points of the system (2.4)–(2.5) consist of de Sitter spaces with constant scalar field $`(H_0,\varphi _0)`$. If they exist, these fixed points satisfy the constraints imposed by eqs. (2.3)–(2.5)
$$H_0=\pm \sqrt{\frac{\kappa V_0}{3}},V_0^{}=0,$$
(2.11)
where $`V_0V(\varphi _0)`$ and $`V_0^{}\frac{dV}{d\varphi }|_{\varphi _0}`$. The Hamiltonian constraint (2.3) can only be satisfied if $`V(\varphi )0`$, which guarantees that $`H_0`$ is real – then there are de Sitter fixed points provided that $`dV/d\varphi `$ has zeros (the existence of equilibrium points is not guaranteed in all scalar–tensor theories). These fixed points lie on the boundary $``$ of the forbidden region<sup>2</sup><sup>2</sup>2In other scalar–tensor theories the fixed points, if they exist, may be located anywhere on the energy surface $`=0`$..
What happens to an orbit with initial conditions chosen exactly on $``$ ? To answer this question we consider the tangent to the orbits
$$\stackrel{}{T}=(\dot{H},\dot{\varphi },\ddot{\varphi })=(\frac{\kappa }{2}\dot{\varphi }^2,\dot{\varphi },\frac{dV}{d\varphi }3H\dot{\varphi })$$
(2.12)
in the $`(H,\varphi ,\dot{\varphi })`$ space. On $``$ it is $`\dot{\varphi }=0`$ and $`\stackrel{}{T}_{}=(0,0,dV/d\varphi )`$. Hence if $`dV/d\varphi |_{}>0`$, an orbit beginning on $``$ will move into the upper sheet, while if $`dV/d\varphi |_{}<0`$, the orbit will move into the lower sheet. If instead $`dV/d\varphi |_{}=0`$ one has a fixed point satisfying the properties $`\dot{H}=0`$, $`\dot{\varphi }=0`$, $`H^2=\kappa V/3`$, and $`V^{}=0`$.
The following result is true for any form of the potential $`V(\varphi )`$: $`H`$ always increases along the orbit of any solution except at points where $`\dot{\varphi }=0`$, at which also $`\dot{H}`$ vanishes.
This follows by combining eqs. (2.3) and (2.4) to obtain<sup>3</sup><sup>3</sup>3Eq. (2.13) is equivalent to eq. (7) of multiplied by $`\dot{\varphi }`$.
$$\dot{H}=\frac{\kappa }{2}\dot{\varphi }^2,$$
(2.13)
hence $`\dot{H}>0`$ except at the points where $`\dot{\varphi }=0`$, in particular the fixed points $`(H_0,\varphi _0)`$.
Eq. (2.13) is expressed by saying that the universe always superaccelerates (i.e., $`\dot{H}>0`$ as opposed to acceleration defined by $`\ddot{a}=\dot{H}+H^2>0`$), or that the phantom field is a form of superquintessence . Eqs. (2.3) and (2.13) imply that
$$\dot{H}+3H^2=\kappa V.$$
(2.14)
A second result is that there are no limit cycles (periodic orbits).
In fact, $`H(t)`$ is non–decreasing and hence it cannot periodically return to its initial value apart from the trivial case of a fixed point. If $`dV/d\varphi `$ has definite sign this possibility is also excluded, cf. eq. (2.11).
## 3 Bounded potential
If the potential $`V(\varphi )`$ is bounded from above by a (positive) constant $`V_0`$, then the asymptotic solution of eqs. (2.3)–(2.5) at large times is such that $`\dot{H}0`$ and $`\dot{\varphi }0`$. The resulting de Sitter attractor $`(H_0,\varphi _0)`$ is a global attractor.
In fact, according to eq. (2.14), $`\dot{H}=\kappa V3H^2`$; assume that $`\dot{H}`$ does not tend to zero (or that it does not tend to zero faster than $`t^1`$) as $`t+\mathrm{}`$. Then the limit
$$\underset{t+\mathrm{}}{lim}H(t)=\underset{t+\mathrm{}}{lim}_{t_0}^t\dot{H}(\tau )𝑑\tau =+\mathrm{}$$
(3.1)
since $`\dot{H}>0`$. But then $`V(t)>3H^2/\kappa +\mathrm{}`$ as $`t+\mathrm{}`$ and this contradicts the fact that $`V`$ is bounded from above. Hence it must be $`\dot{H}0`$ (faster than $`t^1`$) as $`t+\mathrm{}`$.
Since $`\dot{H}=\kappa \dot{\varphi }^2/2`$, then $`\dot{\varphi }0`$ (faster than $`t^{1/2}`$) as well when $`t+\mathrm{}`$. The asymptotic state for a potential bounded from above is one with $`\dot{H}0`$ and $`\dot{\varphi }0`$. Eq. (2.14) gives the asymptotic relation $`H_0^2=\kappa V/3`$ and the Klein–Gordon equation reduces, in this limit, to $`dV/d\varphi |_{\varphi _0}=0`$ and the asymptotic state is the de Sitter fixed point $`(H_0,\varphi _0)`$. The function
$$L(H,\varphi )=\left(HH_0\right)^2$$
(3.2)
defined on the entire phase space is a Ljapunov function – in fact $`L(H,\varphi )>0`$ $`(H,\varphi )(H_0,\varphi _0)`$, $`L(H_0,\varphi _0)=0`$, and $`dL/dt=\kappa \dot{\varphi }^2\left(HH_0\right)0(H,\varphi )(H_0,\varphi _0)`$ because $`H`$ tends to $`H_0`$ while non–decreasing, hence $`HH_0`$ $`t`$. The system is stable and the attraction basin of the de Sitter attractor is the entire phase space. This excludes the possibility of a Big Rip or a sudden future singularity.
If the potential has maxima, minima, or inflexion points, there will be fixed points of the system (2.3)–(2.5). The character of these points determines the stability of the fixed point. One can assess the stability by considering the perturbed equations (2.4)–(2.5) for homogeneous perturbations that depend only on time, as done in . It is more meaningful to study stability against inhomogeneous perturbations which depend on both space and time. Then the complication arises that a gauge-independent formalism is required. The study of inhomogeneous perturbations is carried out in for generalized gravity theories described by the action
$$S=d^4x\sqrt{g}\left[\frac{1}{2}f(\varphi ,R)\frac{\omega (\varphi )}{2}g^{ab}_a\varphi _b\varphi V(\varphi )\right].$$
(3.3)
The result of is that there is stability if and only if
$$\frac{\frac{f_{\varphi \varphi }(\varphi _0)}{2}V_{\varphi \varphi }(\varphi _0)+\frac{6f_{\varphi R}(\varphi _0)H_0^2}{f_R(\varphi _0)}}{\omega _0\left(1+\frac{3f_{\varphi R}^2(\varphi _0)}{2\omega _0f_R(\varphi _0)}\right)}0.$$
(3.4)
Phantom cosmology is the special case $`f(\varphi ,R)=R`$, $`\omega =1`$, $`f_{\varphi R}=f_{\varphi \varphi }=0`$, $`f_R=1`$ and stability corresponds to $`V_0^{\prime \prime }0`$, where a prime denotes differentiation with respect to $`\varphi `$. The perturbation analysis of establishes that the de Sitter fixed point is stable against inhomogeneous linear perturbations if $`V(\varphi )`$ has a maximum there. Here we recover this result in a different approach that does not rely on a linear perturbation analysis and it is valid to any order, extending and complementing the results of applied to phantom cosmology. Moreover, does not draw conclusions about the size of the attraction basin, while here we establish that this attraction basin is the entire phase space. This agrees with, and extends, the result of Guo et al. that, if a late time attractor exists in the phase space of phantom cosmology with bounded potential, it is unique.
The phantom field settling in a maximum and producing a late time de Sitter attractor is consistent with the literature selecting specific potentials $`V(\varphi )`$ and with general conjectures . Let us see some examples. Singh, Sami and Dadhich consider the bell-shaped potential $`V(\varphi )=V_0/\mathrm{cosh}\left(\kappa \alpha \varphi \right)`$ attaining its global maximum $`V_0`$ at $`\varphi =0`$, and they find a late time de Sitter attractor. The same type of attractor is found by Carroll, Hoffman and Trodden for the Gaussian potential $`V(\varphi )=V_0\text{e}^{\varphi ^2/\sigma }`$.
Hao and Li consider the inverse power-law potential $`V(\varphi )=V_0\sigma \left(\varphi /\varphi _0\right)^2`$ for $`\left|\varphi \right|\varphi _0\sqrt{\frac{V_0}{\sigma }}`$, with maximum $`V_0`$ at $`\varphi =0`$. They find again a de Sitter attractor empty of scalar field ($`\varphi =0`$) and with a truly constant cosmological constant $`\mathrm{\Lambda }=V_0`$.
Another situation in which the potential is bounded from above occurs when $`V(\varphi )`$ has a horizontal asymptote $`V_0`$ approached from below as $`\varphi \pm \mathrm{}`$. Under these conditions the scalar field rolls up the slope of the potential but its (negative) kinetic energy is dissipated by friction – described by the term $`3H\dot{\varphi }`$ in eq. (2.5) – while the force term $`dV/d\varphi `$ tends to zero approaching the asymptote, with $`3H\dot{\varphi }dV/d\varphi `$ and $`\dot{\varphi }0`$. As a result the motion stops ($`\dot{\varphi }0`$) and a de Sitter regime with constant $`\varphi `$ and $`H`$ is approached. An example is the potential studied by Sami and Toporensky $`V(\varphi )=V_0\left[1\text{e}^{c\varphi ^2}\right]`$, which has the shape of an inverted bell with $`V_0`$ as horizontal asymptote. These authors find a late time de Sitter attractor with cosmological constant $`\mathrm{\Lambda }=V_0`$ .
Guo et al. consider the bounded potential $`V(\varphi )=V_0\left[1+\mathrm{cos}\left(\varphi /f\right)\right]`$ and find a slow-climb solution in which $`\varphi `$ settles in the maximum of the potential at $`\varphi =0`$ in a de Sitter regime with $`(H_0,\varphi _0)=(\sqrt{\kappa V_0/3},0)`$.
## 4 Unbounded potential
If the phantom potential is not bounded from above it would seem that the asymptotic state could be either a de Sitter regime, or a very different state, possibly a Big Rip, depending on the shape of $`V(\varphi )`$ and its slope. Physically, the motion of $`\varphi (t)`$ is analogous to the motion of a ball with negative kinetic energy falling up a hill under the force $`dV/d\varphi `$, damped by the term $`3H\dot{\varphi }`$. If this friction becomes negligible in comparison with the force and inertial terms, $`3H\dot{\varphi }<<dV/d\varphi ,\ddot{\varphi }`$, one has a “slow-climb” regime analogous to the slow-roll, potential-dominated, regime of inflation for a canonical scalar field . This happens if the potential becomes sufficiently steep. If instead friction is comparable to the force term, motion could stop ($`\dot{H}0`$, $`\dot{\varphi }0`$) corresponding to a de Sitter regime. If $`V`$ is strictly increasing this situation is forbidden and the expansion is super–exponential: if $`V(\varphi )`$ is not bounded from above and is strictly increasing the solution of eqs. (2.3)–(2.5) cannot approach a de Sitter fixed point at late times, but $`H`$ and $`\varphi `$ go to infinity.
In fact, a de Sitter fixed point has $`V^{}=0`$ and in order to approach a de Sitter fixed point it must be $`V^{}0`$, which is incompatible with our assumptions. Hence $`H`$ and $`\varphi `$ go to infinity.
This result explains the examples available in the literature, in which an unbounded potential with $`dV/d\varphi >0`$ never produces a late time de Sitter solution but rather gives a universe that expands faster than exponentially. Sami and Toporensky and Guo et al. consider the simple potential $`V(\varphi )=m^2\varphi ^2/2`$. Let us consider the region $`\varphi >0`$ where $`dV/d\varphi `$ is positive – there $`\dot{\varphi }`$ and $`\dot{H}`$ approach constant values at late times ,
$$\dot{\varphi }\sqrt{\frac{2}{3\kappa }}m,\dot{H}\frac{\kappa m}{\sqrt{3}}\dot{\varphi }$$
(4.1)
as $`t+\mathrm{}`$, and $`H,\varphi t`$ asymptotically. The motion becomes potential–dominated at large times and the ratio of kinetic to potential energy $`\frac{\dot{\varphi }^2}{2V}\frac{\dot{\varphi }^2}{m^2\dot{\varphi }^2t^2}0`$ so that $`V>>\dot{\varphi }^2/2`$ asymptotically. Both $`H`$ and $`\varphi `$ diverge and the scale factor $`a(t)=a_0\text{e}^{\alpha t^2/2}`$ diverges faster than exponentially. However, it takes an infinite time for $`\varphi `$ and $`H`$ to reach infinity and there is no Big Rip. The effective equation of state parameter
$$w\frac{P}{\rho }=\frac{\dot{\varphi }^2+m^2\varphi ^2}{\dot{\varphi }^2m^2\varphi ^2}1$$
(4.2)
as $`t+\mathrm{}`$. The asymptotic equation of state approximates $`P=\rho `$, which is the exact equation of state for de Sitter space, but the solution is not de Sitter space nor does it approach it. Asymptotically, $`\ddot{\varphi }0`$ and the friction and force terms balance each other, $`3H\dot{\varphi }dV/d\varphi `$. The appendix is devoted to this example and shows that (4.1) is an attractor.
Sami and Toporensky also consider power-law potentials $`V(\varphi )=V_0\varphi ^\alpha `$ (with $`\alpha >0`$) and they find a “slow–climb” regime characterized by
$$\dot{\varphi }\frac{V^{}}{3H},H^2\frac{\kappa V}{3},\frac{\dot{\varphi }^2}{2V}\frac{1}{6\kappa }\left(\frac{V^{}}{V}\right)^2$$
(4.3)
if $`\alpha <4`$. For $`\alpha =4`$ they find exponential growth of both $`\varphi `$ and $`H`$, while if the power-law potential is steeper ($`\alpha >4`$), they find a Big Rip solution with $`\varphi (t)\left(t_0t\right)^{\frac{2}{4\alpha }}`$. This result is confirmed by Guo et al. and agrees with the qualitative argument on the steepness of the potential $`V(\varphi )`$.
Hao and Li consider the exponential potential $`V(\varphi )=V_0\mathrm{exp}\left(c\varphi \right)`$ (with $`c>0`$), which is unbounded and has $`dV/d\varphi >0`$ everywhere. They find an attractor with equation of state parameter $`w<1`$ which makes the Big Rip unavoidable. In the same authors consider again the exponential potential $`V(\varphi )=V_0\mathrm{exp}\left[\sqrt{3}\kappa A\left(\varphi \varphi _0\right)\right]`$ with $`V_0,A>0`$ finding a Big Rip attractor in parametric form $`(a(\varphi ),t(\varphi ))`$. By eliminating the parameter $`\varphi `$ one finds the scale factor
$$a(t)=\frac{a_0}{\left(\frac{\sqrt{6V_0}\kappa A^2}{2\sqrt{A^2+2}}\right)^{\frac{2}{3A^2}}}\left[\frac{2\sqrt{A^2+2}}{\sqrt{6V_0}\kappa A^2}+t_0t\right]^{\frac{2}{3A^2}}.$$
(4.4)
In simple terms, the effective equation of state parameter is constant, $`w=(1+A^2)`$, and the scale factor can be written as $`a\left(t_{}t\right)^{\frac{2}{3(w+1)}}`$. It is well known that a spatially flat FLRW universe filled with a fluid with constant equation of state $`P=w\rho `$ has scale factor $`at^{\frac{2}{3(1+w)}}`$ .
It seems intuitive that a potential $`V(\varphi )`$ unbounded and steeper than the exponential potential will always cause a sudden future singularity. This can be shown for extremely steep potentials: if $`V(\varphi )`$ has a vertical asymptote the universe evolves to a future singularity in a finite time.
In fact, since $`V(\varphi )+\mathrm{}`$ as $`\varphi \varphi _c`$, it is $`dV/d\varphi >0`$. $`V`$ is unbounded and strictly increasing and, due to our previous result, $`\dot{\varphi }`$ cannot tend to zero. The only two possibilities are that $`\dot{\varphi }`$ tends to a finite limit $`\dot{\varphi }_c`$ or that $`\dot{\varphi }`$ diverges. In both cases $`\varphi `$ reaches the critical value $`\varphi _c`$ in a finite time and $`V`$ and $`dV/d\varphi `$ diverge. $`\dot{\varphi }`$ cannot have a finite limit $`\dot{\varphi }_c`$: in fact, if this happens, then $`\dot{H}=\kappa \dot{\varphi }^2/2\dot{H}_c\kappa \dot{\varphi }_c^2/2`$ and $`\ddot{\varphi }0`$. Then the modified Klein–Gordon equation (2.5) yields $`3H\dot{\varphi }dV/d\varphi +\mathrm{}`$ asymptotically and $`H`$ must diverge in a finite time, which is in contradiction with $`\dot{H}`$const. Hence it must be $`\dot{\varphi }+\mathrm{}`$ and $`\dot{H}=\kappa \dot{\varphi }^2/2+\mathrm{}`$. It may happen that both $`H`$ and $`\varphi `$ diverge or that they stay finite with their derivatives diverging at a sudden future singularity. Examples of this kind of singularities have been found in other models of dark energy based on the tachyonic field , the brane world , and cosmology with the Gauss–Bonnet term .
Out last result is consistent with the numerical examples of , who find a sudden future singularity for potentials with a vertical asymptote.
## 5 Discussion and conclusions
The dynamics of phantom cosmology are studied without assuming specific potentials $`V(\varphi )`$, but assuming that the phantom dark energy has already come to dominate the cosmic dynamics. In the literature there are also scenarios in which the late time state of the universe is dominated again by cold (dark and ordinary) matter with zero pressure and the phantom energy decays : such scenarios are a priori excluded from our analysis.
A clear and unified picture of the dynamics is obtained: the phase space is a two–dimensional surface in the $`(H,\varphi ,\dot{\varphi })`$ space and is composed of two sheets. All the equilibrium points, which exist if and only if $`dV/d\varphi `$ has zeros, are de Sitter spaces located on the curve where the two sheets touch each other – they are given by eqs. (2.11). $`H`$ is always increasing except when $`\dot{\varphi }`$ vanishes, and there are no periodic orbits. In the presence of a bounded potential $`V(\varphi )`$ the universe has a late time de Sitter attractor whose attraction basin is the entire phase space. If $`V(\varphi )`$ is unbounded there cannot be a de Sitter attractor and the late time cosmic expansion is either super–exponential with infinite lifetime or else the universe ends in a Big Rip or other sudden future singularity. Unfortunately, in the general situation, we cannot provide a sharp criterion on the potential to discriminate between these two possibilities. Many results for specific potentials that are found in the literature are recovered and predictions are made for new scenarios featuring potentials with the shapes discussed.
It is worth remembering that there are serious doubts on whether phantom fields can be stable and that the classical equations of motion considered have to be superseded by a semiclassical treatment near the Big Rip singularity. Indeed there is some evidence that semiclassical backreaction may avoid the Big Rip .
Finally, there are compelling arguments for a field with canonical kinetic energy to couple nonminimally to the Ricci curvature . In particular, conformal coupling is associated with an infrared fixed point of the renormalization group and avoids causal pathologies . It is not clear whether the same arguments apply to a phantom field and force it to couple nonminimally (a nonminimally coupled phantom is considered in ). The slow–climb regime for a phantom would be altered by the inclusion of non–minimal coupling terms in eqs. (2.3)–(2.5), similarly to what happens in slow-roll inflation . The dynamics of a nonminimally coupled phantom will be studied elsewhere.
## Acknowledgments
The author thanks Werner Israel for pointing out Ref. . This work was supported by the Natural Sciences and Engineering Research Council of Canada (NSERC).
## Appendix: massive phantom field
Here we consider a phantom field with potential $`V(\varphi )=m^2\varphi ^2/2`$ as an example. The phase space is the two–dimensional surface $`\frac{6}{\kappa }H^2+\dot{\varphi }^2=m^2\varphi ^2`$. By using the dimensionless variables
$$x\sqrt{\kappa }H,y\frac{\kappa m}{\sqrt{6}}\varphi ,z\frac{\kappa }{\sqrt{6}}\dot{\varphi },$$
(A.1)
the equation of the surface can be written as $`y=\pm \sqrt{x^2+z^2}`$, which describes a cone with axis along the $`y`$–axis and vertex in the origin. The upper sheet is described by $`y=+\sqrt{x^2+z^2}`$ and the lower sheet by $`y=\sqrt{x^2+z^2}`$; the two sheets join on the plane $`z=0`$ along the lines $`y=\pm x`$. Although we restrict our analysis to spatially flat ($`K=0`$) universes, it can be shown that the cone separates trajectories belonging to an open ($`K=1`$) universe lying inside the cone from orbits belonging to a closed ($`K=+1`$) universe lying outside the cone (cf. ).
The only equilibrium point satisfying eqs. (2.11) is the Minkowski space $`(x,y,z)=(0,0,0)`$. Our results imply that $`H`$ and $`\varphi `$ go to infinity for any solution except the equilibrium point, and hence the equilibrium point is unstable. Since $`V^{\prime \prime }(0)=m^2>0`$ the fixed point (which coincides with the global minimum of the potential) is unstable – see the discussion of Sec. III.
This three-dimensional picture should be compared with fig. 3 of showing a projection of the phase space (a two-dimensional cone) and of some orbits onto the $`(\varphi ,\dot{\varphi })`$ plane, corresponding to our $`(y,z)`$ plane. A similar conical geometry of the phase space structure appears if one considers inflation realized by a massive scalar with the “right” sign of the kinetic energy density .
The trajectories of the solutions exhibit a late time attractor found analytically in (eq. (4.1)). The authors of show numerically the convergence of many solutions to this asymptotic solution but do not prove analytically its attractor nature. The attractor can be obtained by requiring an asymptotic linear relationship between $`H`$ and $`\varphi `$, as suggested by numerical integrations. By setting $`H=\alpha \varphi `$ one obtains, using eq. (2.13), $`\dot{\varphi }=2\alpha /\kappa `$. Eq. (2.3) yields $`H\pm \sqrt{\frac{2}{3\kappa }}\alpha mt`$, while eq. (2.4) is identically satisfied and eq. (2.5) gives $`\alpha =m\sqrt{\kappa /6}`$. Since $`\dot{H}0`$ the negative sign is rejected. The asymptotic solution found in this way coincides with the one of
$$(H_{},\varphi _{})=(\frac{m^2}{3}t,\sqrt{\frac{2}{3\kappa }}mt).$$
(A.2)
The stability with respect to linear homogeneous perturbations is assessed by considering the perturbed Hubble parameter and scalar field $`H=H_{}+\delta H`$, $`\varphi =\varphi _{}+\delta \varphi `$, and inserting these expressions into the evolution equations (2.3)–(2.5). It is convenient to consider the contrasts $`\mathrm{\Delta }_1\delta H/H_{}`$ and $`\mathrm{\Delta }_2\delta \varphi /\varphi _{}`$. One finds $`\mathrm{\Delta }_1=\sqrt{\frac{3\kappa }{2}}\frac{1}{mt}\left(\delta \varphi \frac{\delta \dot{\varphi }}{m^2t}\right)`$ and eq. (2.5) yields the evolution equation for $`\mathrm{\Delta }_2`$
$$\ddot{\mathrm{\Delta }_2}+\left(\frac{2}{t}m+m^2t\right)\dot{\mathrm{\Delta }_2}+\left(m^2\frac{1}{t^2}\right)\mathrm{\Delta }_2=0$$
(A.3)
with asymptotic solution $`\mathrm{\Delta }_2`$const. and $`\mathrm{\Delta }_1=\mathrm{\Delta }_2\left(1\frac{1}{m^2t^2}\right)\frac{\dot{\mathrm{\Delta }_2}}{m}\text{constant}`$. The solution $`(H_{},\varphi _{})`$ is stable and is an attractor, in agreement with the numerical results. |
warning/0506/quant-ph0506174.html | ar5iv | text | # Quantumness of ensemble from no-broadcasting principle
## I Introduction
The main conceptual novelty of quantum information revolution was the notion of quantum information itself. Already before quantum information era there had been attempts to consider information more generally than classically Ingarden (1976); Ohya and Petz (1993). However those concepts did not have operational meaning. The real breakthrough turned out to be the concept of sending quantum states intact, initiated in the seminal work on teleportation Bennett et al. (1983) (for experimental realizations see Bouwmeester et al. (1997); Boschi et al. (1998); A.Furusawa et al. (1998)) and in Schumacher quantum noiseless coding theorem Schumacher (1995).
This concept was truly revolutionary, because previously, the dominating point of view was the Copenhagen interpretation. The latter is extremely epistemic: according to it, the wave function is merely a description of the measuring and preparing apparatuses. The emphasis is put on classical input (parameters of preparation) and classical output (clicks of detectors). From such a point of view, it was rather impossible to think that a message could be something else than a classical message. Indeed any process of sending message is actually some quantum experiment. If the interpretation of an experiment is dominated by the idea of having classical input and classical output, then one imagines that the sender must input a classical message and the receiver must also receive a classical message, not a quantum state.
Another obstacle against arriving at the concept of quantum information was that information was thought as closely related to knowledge: even if communicated and processed by physical means, (classical) information is something that one can get to know. However quantum information is more like a “thing” than knowledge. Knowledge can be shared, quantum information can not: if one passes some quantum information to another person, she cannot keep it at the same time. For example, in the process of teleportation the input state transferred to the receiver is destroyed at the sender site.
The above obstacles have a common denominator: in the pre-quantum-information era the emphasis was put on the subject, not on the object. This passive paradigm Horodecki et al. (2004) has shaped for a long time the understanding of quantum mechanics. The quantum information revolution forces us to put emphasis on the object, because quantum information cannot be expressed in a natural way in terms of preparation and measurement procedures. Well, one can argue that everything can be earlier or later recast in terms of preparation procedure and classical outcomes. An experimenter is even forced to consider an experiment in such terms. However, if we treat such approach as the unique one, we may overlook some phenomena, in a similar way, as looking at the sky, it is rather hard to recognize the rule that governs the planetary motion.
As well known, Asher Peres was one of the authors of the paper on teleportation, which initiated the quantum information revolution. Quite paradoxically, being undoubtedly one of the fathers of the quantum information domain Peres (2004), he nevertheless maintained an epistemic point of view, considering the wave function as a mere tool for prediction of probabilities. His point of view was motivated by the fact that the attempts of ontologization of quantum mechanics lead notoriously to paradoxes Peres (1993). We hope that one can avoid this by some kind of reformulation of the notion of reality (see e.g. d’Espagnat (1995)).
As we have mentioned, the most distinctive feature of quantum information, which makes it a “thing” rather than “knowledge”, is that it cannot be shared (much as a (classical) physical object that can not be both here and there at the same time). For pure states, this is the content of the famous no-cloning theorem Wootters and Zurek (1982); Dieks (1982) (see pioneering attempt by Wigner in Wigner (1967)). Recently, various properties of quantum information in the context of no-cloning and dual no-deleting principle Pati and Braunstein (2000) have been analyzed Jozsa (2004); Horodecki et al. (a); Pati and Sanders . In particular, in Horodecki et al. (a) a new angle of looking at no-cloning principle and no-deleting principles was proposed. Namely, both principles can be viewed as consequences of (i) a principle of conservation of information and (ii) the fact that in the quantum case, two copies can have more information than one copy.
In Horodecki et al. (a) the difference between quantum and classical information was not its conservation (which we have postulated for both), but just the relation between information of two copies and one copy. However, the no-cloning distinguishes between quantum and classical only on the level of pure states. Indeed, classical probability distribution also cannot be cloned. It is the no-broadcasting theorem Barnum et al. (1996) that reports full difference: it says that states can be broadcast if and only if they commute with each other. Broadcasting is a generalization of cloning that allows for correlations between copies. Thus, extending Horodecki et al. (a), we can say that two copies, even realized marginally on separate systems, contain more information than a single copy only if the states do not commute.
In this paper we develop a quantitative approach to this peculiar property of quantum information. To this end we start with axiomatically defined information $`I`$ of an ensemble, and then study how much the information of two copies exceeds the information of one copy (where we allow the copies to be correlated). The difference denoted by $`I_q`$ shows quantumness of the ensemble. For any classical ensembles (ensembles of commuting states) the difference vanishes, because such ensembles can be broadcast. A desired property of $`I`$ is that it always feels that an ensemble cannot be broadcast, i.e. that for such ensemble, the difference $`I_q`$ is always strictly positive. From Barnum et al. (1996) it follows that one can build one such function $`I`$ using fidelity. However fidelity is not the best function to quantify information, as it is hard to obtain from it an extensive quantity. Therefore we turn to entropic quantities. Namely we show that also Holevo quantity
$$\chi (\{p_i,\rho _i\})=S(\underset{i}{}p_i\rho _i)\underset{i}{}p_iS(\rho _i)$$
(1)
is good, in a sense that difference between one and two copies is always nonzero for non-broadcastable ensembles. We also discuss other possibilities.
The function $`I_q`$ can also serve as an axiomatically defined quantum contents of ensemble, and can serve as a recipe to produce a new function quantifying “quantumness of ensemble” (see Fuchs and Sasaki ; Horodecki et al. (b)). By showing that Holevo quantity gives nonzero $`I_q`$ for any noncommuting ensemble, we have obtained a new measure of quantumness, that it nonzero for any nonclassical ensemble.
## II Cloning versus broadcasting
The no-cloning theorem states that there does not exist any process, which turns two distinct nonorthogonal quantum states $`\psi ,\varphi `$ into states $`\psi \psi ,\varphi \varphi `$ respectively. In Horodecki et al. (a) the theorem was connected with a principle of conservation of information. Namely, two copies contain more information than one copy, and therefore it is impossible to produce two copies out of one copy. Even if the information is not conserved, it is at least monotone under operations, and still the main argument holds: cloning is in general impossible, because two copies have more information than one copy.
We actually do not know very well what information is, so it is safer to use “information monotones”, i.e. functions that do not increase under physical operations. In Horodecki et al. (a) we have taken entropy as a measure of informational contents of ensemble of states. Entropy is a monotone, if one restricts to pure states. (In the following, since we deal with mixed states, the Holevo function would be more appropriate.) If one is concerned with only two states, one can take a function of just two arguments, such as fidelity, and this was considered in Alicki (2005).
Why is no-cloning a non-classical feature? The question arises, because one also cannot clone classical probabilities, by the same reason: the information contained in two samples of a probability distribution is more than that contained in single sample. One answer is the following. In the quantum world, pure states cannot be cloned, while in the classical world, pure states (i.e. probability distributions with one probability equal to 1 and all others to 0) can be cloned. However, one may not be satisfied with putting quantum and classical state on the same footing. Indeed, quantum pure states involve probabilities, so it may be more appropriate to compare quantum pure states with all classical states, not only pure ones. If so, then one can conclude that no-cloning principle holds both in the quantum and in the classical world Alicki (2005). However there is still a fundamental difference: in classical theory one can always broadcast information. Namely, from two classical states $`\rho `$ and $`\sigma `$ one can obtain states $`\rho _{AB}`$ and $`\sigma _{AB}`$ such that $`\rho _A=\rho _B=\rho `$ and $`\sigma _A=\sigma _B=\sigma `$. The difference between cloning and broadcasting is that in cloning one requires to obtain independent copies:
$$\rho \rho \rho ;$$
(2)
in broadcasting correlations between the copies are allowed, i.e. we require only
$$\rho \rho _{AB},$$
(3)
where $`\rho _A=\rho _B=\rho `$.
Quantumly, broadcasting is not always possible. This can be inferred already from the no-cloning theorem: the latter says that nonorthogonal pure states cannot be cloned. However for pure states broadcasting is equivalent to cloning. The question of broadcasting quantum states was first considered in Barnum et al. (1996). There it was shown that much more is impossible, than predicted by no-cloning theorem. Namely, the states cannot be broadcast if and only if they do not commute. Thus we see that impossibility of broadcasting is purely quantum feature, as it goes in parallel with noncommutativity - the main feature distinguishing quantum theory from classical one.
## III Information contents of one copy versus two copies
In view of the above remarks, one can apply the concept of information contents to the problem of broadcasting. In quantum theory broadcasting is impossible, because, for noncommuting states, two copies (however correlated) will have more information than single copy. As we have said, we can formalize it without knowing what information really is, but rather assuming that it cannot increase under operations Horodecki et al. (2003a, b). Indeed, whatever information is, in any theory, it cannot increase under operations allowed in the theory. It follows, in particular, that for reversible operations, information is conserved.
In the context of cloning/broadcasting, we talk about sets of states. Thus information would be a monotonous function of set of states (cf. Horodecki et al. (2003a, b) where we analyzed information as a function of states themselves). We require any candidate for information to satisfy the following postulates:
1. $`I0`$,
2. $`I(𝒮)=0`$ if and only if $`𝒮`$ contains one and only one element,
3. $`I(\mathrm{\Lambda }(𝒮))I(𝒮)`$ (monotonicity),
where $`𝒮=\{\rho _i\}_i`$ denotes a set of states, $`\mathrm{\Lambda }`$ is any quantum operation, and $`\mathrm{\Lambda }(𝒮)\{\mathrm{\Lambda }(\rho _i)\}_i`$. An example is Holevo function of the ensemble with equal apriori probabilities. If we consider just two states, then one can also use fidelity as in Alicki (2005) (actually in Barnum et al. (1996) the authors used fidelity to show that noncommuting states cannot be broadcast). In next section we will consider ensembles of states, i.e. sets of states with ascribed probabilities. The postulates are then analogous. Then the example is just Holevo function of the ensemble.
Consider an information function $`I`$ satisfying the above postulates. Then broadcasting of states $`\sigma ,\rho `$ is impossible, if for any states $`\rho _{AB},\sigma _{AB}`$ which have $`\rho `$ and $`\sigma `$ on both subsystems respectively the function is greater than for $`\rho `$ and $`\sigma `$:
$$I(\rho ,\sigma )<I(\rho _{AB},\sigma _{AB}).$$
(4)
Indeed, one can obtain $`\rho ,\sigma `$ from $`\rho _{AB},\sigma _{AB}`$ by partial trace. Thus for sure, $`I(\rho ,\sigma )I(\rho _{AB},\sigma _{AB})`$. If broadcasting is possible, there exists an operation transforming $`\rho `$ into $`\rho _{AB}`$ and $`\sigma `$ into $`\sigma _{AB}`$; then we have to have also $`I(\rho ,\sigma )I(\rho _{AB},\sigma _{AB})`$. Consequently for states that can be broadcast we have
$$I(\rho ,\sigma )=I(\rho _{AB},\sigma _{AB}).$$
(5)
In other words, broadcasting means that the operation of obtaining one copy from two copies is reversible, and thus any information monotone must be conserved. Then if there are states for which it cannot be conserved, then they cannot be broadcast.
This suggests to define a new quantity $`I_q`$ for any information monotone $`I`$. The quantity would report how much the information contents of two copies exceeds the information contents of one copy. Since two copies can be realized in many different ways, we will take the infimum over all realization of two copies. Thus, for any information monotone $`I`$ we define
$$I_q(\rho ,\sigma )=\underset{\rho _{AB},\sigma _{AB}}{inf}I(\rho _{AB},\sigma _{AB})I(\rho ,\sigma )$$
(6)
where infimum is taken over such states $`\rho _{AB}`$, $`\sigma _{AB}`$ that $`\rho _A=\rho _B=\rho `$, $`\sigma _A=\sigma _B=\sigma `$.
Now, it would be good, if $`I_q`$ is nonzero if and only if the states cannot be broadcast. This would mean that the function reports presence of quantum information if only it is indeed present (i.e. if only the states cannot be broadcast). Choosing a particular information monotone, it is not easy to check that for all states giving rise to two copies, the monotone is greater than for single copies. In Barnum et al. (1996) it was shown that for fidelity it is the case. Thus if we take $`I^F=1F(\sigma ,\rho )`$, then $`I_q^F`$ is nonzero if and only if the states cannot be broadcast.
However fidelity cannot give rise to an extensive quantity, which we think is more appropriate to quantify information in whatever context. Moreover, fidelity can be defined just for two states, while it is convenient to extend the quantity to ensembles of more than two states. Therefore we propose to choose $`I`$ based on Holevo quantity $`\chi `$ given by
$$\chi (\{p_i,\rho _i\})=S\left(\underset{i}{}p_i\rho _i\right)\underset{i}{}p_iS(\rho _i);$$
(7)
we will denote also our information monotone by $`\chi `$:
$$\chi (\rho ,\sigma ):=\chi \left(\{(1/2,\rho ),(1/2,\sigma )\}\right).$$
(8)
Let us now argue that $`\chi _q`$ is nonzero if and only if the states cannot be broadcast
###### Proposition 1
The quantity $`\chi _q(\rho ,\sigma )`$ is nonzero if and only if the states $`\rho ,\sigma `$ cannot be broadcast.
Proof. We write the quantity $`\chi _q`$ as a difference of two relative entropies:
$$\chi _q=S(\rho _{ABC}|\rho _{AB}\rho _C)S(\rho _{AC}|\rho _A\rho _C)$$
(9)
where
$$\rho _{ABC}=\frac{1}{2}|0_C0|\rho _{AB}+\frac{1}{2}|1_C1|\sigma _{AB}$$
(10)
with $`|i`$ orthogonal states, and $`\rho _{AB}`$, $`\sigma _{AB}`$ are states that optimize the infimum in the definition of $`\chi _q`$. Now we use a theorem by Petz Petz (1986) (see also Hayden et al. (2004)) which says that if $`S(\rho |\sigma )=S(\mathrm{\Lambda }(\rho )|\mathrm{\Lambda }(\sigma ))`$ where $`\mathrm{\Lambda }`$ is a trace preserving completely positive (CPTP) map, then there exists another CPTP map $`\mathrm{\Gamma }`$ which reverses the map $`\mathrm{\Lambda }`$ on the considered states:
$$\mathrm{\Gamma }(\mathrm{\Lambda }(\rho ))=\rho ,\mathrm{\Gamma }(\mathrm{\Lambda }(\sigma ))=\sigma .$$
(11)
In our case the map $`\mathrm{\Lambda }`$ is the partial trace over system $`B`$. Petz gives explicit form of the map, from which it follows that in our case the map $`\mathrm{\Gamma }`$ must be of the form $`\mathrm{\Gamma }_{AAB}\mathrm{id}_C`$. One can alternatively get it from the following two facts: the second argument of relative entropy in our formulas is a product state; the $`C`$ part is the same in both arguments. From the former it follows that the map must be product, from the latter that the $`C`$ part must be identity. Thus, if $`\chi _q=0`$ then there exists a map that produces states $`\rho _{AB}`$ and $`\sigma _{AB}`$ from $`\rho _A`$ and $`\sigma _B`$, which means that the states can be broadcast. This ends the proof.
## IV Quantumness of ensemble
In Fuchs and Sasaki ; Horodecki et al. (a) measures of quantumness of an ensemble have been proposed. Here we propose a new way of quantifying quantumness of ensemble, by looking at the difference between information contents in one copy and information contents in two or more copies. Thus, for a given information monotone $`I`$ on ensembles, we define $`I_q^{(n)}`$ as follows
$$I_q^{(n)}()=infI(_{A_1A_2\mathrm{}A_n})I()$$
(12)
where infimum is taken over all ensembles $`_{A_1A_2\mathrm{}A_n}`$ which, when partially traced over all subsystems but one, reproduce ensemble $``$. To obtain a particular measure of quantumness, we take $`I`$ to be Holevo quantity. The proposition 1 holds also in this more general case, so that we obtain that $`\chi _q^{(n)}()`$ is nonzero if and only if there exist two states in ensemble that can not be broadcast. However, by Barnum et al. (1996) we know that states can be broadcast if and only if they commute. Thus we obtain that $`\chi _q^{(n)}()`$ is zero if and only if the ensemble is entirely classical, i.e. all states commute with each other.
We now can consider the limit of $`n\mathrm{}`$. Then, for ensembles of pure states we have $`I_q^{(\mathrm{})}()=H(\{p_i\})I()`$, where $`H`$ is Shannon entropy. Thus here quantumness reports just how the states are indistinguishable. It would be interesting to understand what is the result for ensembles of mixed states in the limit of infinite amount of copies.
Another candidate for information monotone that should feel quantumness if only it is present is accessible information $`I_{\mathrm{acc}}`$. Indeed, the original “meaning” of the Holevo quantity is that of being an upper bound to the accessible information; such a bound is achieved exactly when the states forming the ensemble commute.
That $`I_{\mathrm{acc}}`$ is an information monotone comes from its very definition as the maximal mutual information between the (classical) input of a sender and the (classical) output of the receiver; if an operation could increase it, it could be used by the receiver to achieve a better mutual information.
While the Holevo quantity is in principle easily computed as a function of the ensemble alone, the evaluation of $`I_{\mathrm{acc}}`$ requires to find the optimal measurement strategy to achieve the maximal mutual information. Notice that $`\chi _q()0`$ can be considered a consequence of strong subadditivity; $`I_q^{\mathrm{acc}}()0`$, apart coming from monotonicity, can be understood as the fact that, given an optimal POVM $`\{M_i\}`$ for $``$, the POVM $`\{M_i^{\mathrm{AB}}=M_i1\}`$ provides a lower bound to $`I_{\mathrm{acc}}(_{\mathrm{AB}})`$ equal to $`I_{\mathrm{acc}}()`$. On the other hand, having at disposal two broadcast copies intuitively gives the receiver the opportunity to discern better the different states appearing in the ensemble.
However we have not been able to prove that associated $`I_q^{\mathrm{acc}}`$ is nonzero if and only if the states do not commute. Such a proof is not immediate because it involves two parallel optimizations, both for $`\{M_i\}`$ and $`\{M_i^{\mathrm{AB}}\}`$; neither the original proof of Holevo (1973) (Kholevo) nor the proof by Fuchs and Caves Fuchs and Caves (1994) can be directly applied. Let us consider the case of pure states. Then, the limiting case of quantumness based on $`I_{\mathrm{acc}}`$ is simply given by $`(I_{\mathrm{acc}})_q^{\mathrm{}}=H(\{p_i\})S(\rho )`$. We can compare now our quantities with the quantumness proposed by Fuchs Fuchs as $`Q_F=\chi I_{\mathrm{acc}}`$. We thus have for ensembles of pure states
$`\chi _q^{\mathrm{}}()`$ $`=`$ $`H(\{p_i\})S(\rho )`$ (13)
$`(I_{\mathrm{acc}})_q^{\mathrm{}}()`$ $`=`$ $`H(\{p_i\})I_{\mathrm{acc}}()`$ (14)
$`Q_F()`$ $`=`$ $`S(\rho )I_{\mathrm{acc}}()`$ (15)
Thus we obtain that for pure state ensembles Fuchs’ measure is the difference between our two:
$$Q_F()=(I_{acc})_q^{\mathrm{}}()\chi _q^{\mathrm{}}()$$
(16)
It would be interesting to compare these quantities with quantumness of ensemble proposed in Horodecki et al. (b). However it is not easy to get expression of that measure.
### IV.1 Permanence and superbroadcasting
Quantumnes first considered by Jozsa Jozsa (2004), called by him permanence. In particular he asked what is needed to get two copies of an ensemble of pure states, if we already possess one copy. It turns out, that one has essentially to bring in the second copy. More precisely, in Jozsa (2004) it was proved that, given any finite set of states $`\{\psi _i\}`$ containing no orthogonal pairs of states and a set $`\{\rho _i\}`$ of (generally mixed) states indexed by the same labels, there is an operation
$$|\psi _i\rho _i|\psi _i|\psi _i$$
(17)
if and only if there is an operation
$$\rho _i|\psi _i.$$
This indeed means that the original state $`|\psi _i`$ (to be cloned) does not help in the process and the information must be provided completely by means of the ancilla state $`\rho _i`$.
It is interesting how this property can be generalized to mixed states, where the natural paradigm is broadcasting rather than cloning.
One may reformulate the problem of realizing the transformation (17) in a more general context. Suppose A and B share one unknown state out of a set of possible states $`\{\rho _1^{AB},\mathrm{},\rho _n^{AB}`$}. The task is to transform $`\rho _i^B`$ into $`\rho _i^A`$, for all $`i=1,\mathrm{},N`$, so that the reduced states are the same at both sites. B may send his (unknown) subsystem to A through a perfect quantum channel.
If the reduced states on A side are pure the shared states must be product and we go back to the original problem: if such pure states are not orthogonal, they cannot be broadcast (cloned), and the new copy must be brought in, if possible, by B by a local operation. B does everything alone, rather than help.
If instead the reduced states $`\rho _i^A`$ are mixed, sending the B part to A may be useful, as the following example shows.
Consider a pair of orthogonal states
$$|\psi _1^{AB}=|00,|\psi _2^{AB}(a)=\sqrt{a}|11+\sqrt{a}|10+\sqrt{12a}|01,$$
with $`0a1/2`$. The corresponding reduced states are
$$\rho _1^A=\rho _1^B=|00|,\rho _2^A(a)=\left(\begin{array}{cc}12a& \sqrt{a(12a)}\\ \sqrt{a(12a)}& 2a\end{array}\right),\rho _2^B(a)=\left(\begin{array}{cc}1a& a\\ a& a\end{array}\right).$$
We have
$$[\rho _1^A,\rho _2^A(a)]=\left(\begin{array}{cc}0& \sqrt{a(12a)}\\ \sqrt{a(12a)}& 0\end{array}\right),$$
(18)
so that for $`a0,1/2`$ the reduced states on A side do not commute and can not be broadcast. Moreover, for qubits it holds in general Alberti and Uhlmann (1980) that, fixed two pairs of states $`\{\rho _1,\rho _2\}`$ and $`\{\sigma _1,\sigma _2\}`$, there exists a CPTP map $`\mathrm{\Lambda }`$ such that $`\sigma _i=\mathrm{\Lambda }[\rho _i]`$, $`i=1,2`$, if and only if
$$\rho _1t\rho _2_1\sigma _1t\sigma _2_1t^+,$$
(19)
with $`A_1=\mathrm{Tr}\sqrt{AA^{}}`$ the trace norm. In our case the condition for
$$\{\rho _1^B,\rho _2^B(a)\}\{\rho _1^A,\rho _2^A(a)\}$$
(20)
to be realized by an operation acting on one subsystem only can be easily checked to be satisfied only for $`a=0`$.
Therefore for $`0<a<1/2`$ the reduced states $`\rho _1^A,\rho _2^A(a)`$ can neither be broadcast nor be the output of an operation performed on B subsystem only. On the other hand, since the total states are orthogonal, it is possible to implement a global transformation such that for the reduced density matrices (20) holds.
We conclude that, considering (reduced) mixed states in the broadcasting framework, global operations involving the total system whose reduced state is to be cloned may be helpful, contrary to the result (for pure states) of Jozsa (2004). The question of what are the most general conditions on the set of states $`\{\rho _i^{AB}\}`$ to obtain a result similar to the latter remains open. For example, even limiting the problem to initial factorized states, i.e. $`\{\rho _i^{AB}=\rho _i^A\rho _i^B\}`$, and requiring that the $`\rho _i^A`$ states do not commute pairwise, i.e. $`[\rho _i^A,\rho _j^A]0`$ for $`ij`$, so that they cannot be broadcast, B subsystems may help. A trivial case where this happens is the following: if $`\rho _i^A=\rho _i^{A_1}\rho _i^{A_2}`$ and $`\rho _i^B=\rho _i^{B_1}\rho _i^{B_2}`$, with $`\rho _i^{B_2}=\rho _i^{A_2}`$, then the possibility of broadcasting the A<sub>1</sub> parts, i.e. $`[\rho _i^{A_1},\rho _j^{A_1}]=0`$ for all $`i,j`$, is sufficient for broadcasting the whole states $`\rho _i^A`$.
Another hint to how any concept of permanence for mixed states needs a more complex approach is given by *superbroadcasting*. In fact, in Ariano (2005) it was proved that, given $`N`$ independent copies of an arbitrary mixed qubit state $`\rho `$, i.e. given $`\rho ^N`$, it is possible to *broadcast* $`M`$ copies of it, with $`M`$ arbitrarily greater than $`N`$ if $`\rho `$ is sufficiently mixed and $`N6`$. Such result does not contradict neither Barnum et al. (1996) nor the present work, since it is $`\rho `$ which is broadcast, not $`\rho ^N`$, i.e. the task is different or, from another point of view, the initial resource is greater.
It is worth noting that if the “standard” broadcasting, from one to two copies, is possible, then it is reversible, so that the information content of the original ensemble and the broadcast ensemble is the same according to our postulates. That the same holds for superbroadcasting is not evident: after superbroadcasting we are assured to be able to get, locally, only a single copy of $`\rho `$. The information content of the initial independent $`N`$ copies is smeared over $`M`$ qubits and, in general, may be not possible to recover it completely. We will consider this topic elsewhere.
## V Appendix: Quantum information and isomorphism principle
As we have dealt with characterization of quantum information it is natural to ask about its role and status in quantum physics. In particular, our motivation to discuss quantum information in the context of philosophy of physics follows in part from the fact that its impact on interpretative problems is rather little. For instance, in a recent interesting review article on interpretations of Quantum Mechanics the term “quantum information” does not occur even once Schlosshauer (2004). Does it mean that a consistent interpretation can dispense with this notion? We think, that it is unlikely. Instead, it cannot be excluded that a permanent interpretative crisis follows from the fact that the concept of quantum information goes beyond standard axioms of quantum formalism.
In this context a basic question arises: what fundamental condition should any interpretation (philosophy) satisfy to be adequate? We believe that the “minimal” condition is that it cannot ignore the rather profound fact that Nature allows to describe itself. It seems that any philosophy (interpretation) of quantum mechanics should take it seriously. This fact suggests that Nature has an ordered structure, and this order is partially revealed in any successful mathematical description. This can be formulated as follows Horodecki et al. (2004):
Isomorphism principle: Any consistent description of Nature is a sort of isomorphism between the laws of Nature and their mathematical representation.
Such a principle could have been regarded as quite trivial before quantum mechanics was born. Indeed, in classical theories, the observer was passive, hence the structure of theory could be directly assigned to structure of Nature. Physical notions were easily associated with some realities. However, in quantum mechanics there does not exist a passive observer, and the postulate of existence of an isomorphism is nontrivial.
The isomorphism means that the theoretical structure consistent with physical phenomena, although not a real thing, is an isomorphic image of the existing reality. If accepted, it supports the view, that in quantum information era any attempts to understand quantum formalism should take into account the notion of quantum information. Indeed, one can ask: why and for what particular feature does Nature require an abstract mathematical description in terms of Hilbert space? From the point of view of the isomorphism principle, the answer is: the Hilbert space formalism reflects the structure of Nature itself, rather than being only an abstract, descriptive tool. The basic notions associated with Hilbert space formalism should consequently be also taken into account in building a consistent interpretation. One of such a basic notions is undoubtedly quantum information, and in this context, it should be treated as seriously as energy inf . For example, basing on the close relationships between the notion of quantum information and the notion of entanglement Schumacher (1995); Bennett et al. (1997), in Horodecki et al. (2001) principle of conservation of information was formulated in terms of entanglement: Entanglement does not change under local operations in closed system (see also Horodecki and Horodecki (1998) in this context).
According to the isomorphism principle, the quantum information, though not necessary a real thing, reflects some physical reality. In particular this allows to avoid the longstanding dilemma between Scylla ontology and Charybdis epistemology which symbolizes the two opposite (extremal) views Einstein and Bohr on the nature of physical reality in relation to quantum formalism. Indeed, the heart of the Copenhagen interpretation is a Ptolemaic paradigm, taking as its “reference frame” preparation parameters and outcomes of measurements. It involves the “surface” of reality in the sense that the wave function provides only a mathematical representation of our knowledge about the experimental setup. From this point of view “photons are clicks in photon detectors” and “there is no quantum information, there is only a quantum way of handling information” phy . In contrast, Einstein’s ontological concept of physical reality involves the objective state of a system specified by a set of parameters independently of our knowledge of them. The Kochen-Specker theorem shows that extremal version of Einstein’s view cannot be valid any longer.
Remarkably, the isomorphism principle indicates a more suitable approach which lies between the two extremes. It involves a Copernican like active paradigm which takes as a reference frame “what is actually processed in the laboratory”. Three stages are distinguished: preparation, control and measurement. The new feature here is the introduction of the control stage as an autonomous part. The stage contains a quantum system, which may be a compound of subsystems. The latter can be localized in space, controlled individually, and communicated. In particular, the preparation part can be almost completely absorbed into the control part. For example, in quantum computation it is only necessary to prepare the standard input state. Consequently a quantum experiment can be thought of (on the conceptual level) as being mostly control followed by measurement. In accordance with the isomorphism principle the wave function is not only a tool for calculation of probabilities but it is the isomorphic image of what is actually processed in the laboratory. Basing on the information isomorphism we can claim that quantum information is carried by a quantum system and that the wave function is the image of this information. It should be emphasized here, that there is a substantial difference between the quantum information and its image: In contrast to the wave function, the quantum information itself cannot be regarded as a sequence of classical symbols.
In this context it is natural to ask how quantum information suggested by the isomorphism principle manifests itself in reality. It seems that the simplest criterion may by related to the concept of resource. As we know, in quantum case the newly discovered resources are highly nonintuitive and much more subtle than classical ones. One may consider the following resource criterion: If a quantum property associated with quantum system can be used as resource in some nonclassical tasks, we say that it can be related to reality. In this sense, weaker than the commonly used one, such properties as quantumness of ensembles and entanglement reflect some reality. As a matter of fact, quantum information in the form of nonorthogonal quantum states and quantum information in the form of entanglement are related to each other: possibility of communication of nonorthogonal states is equivalent to possibility of sharing entanglement Schumacher (1995); Bennett et al. (1997); Horodecki and Horodecki (1998). However the full meaning of quantum information is still far from being clear Gisin .
There is a practical reason, for which the isomorphism principle seems to be more appealing than the passive Ptolemaic paradigm. Isomorphism asks us to take seriously quantum formalism; the Copenhagen interpretation has not taken it seriously enough, and as a result the discovery that quantum states may be processed, has been done surprisingly late! A very illustrative example of a discovery that was not possible within “Ptolemaic approach” was teleportation. Namely, the measurement in Bell basis in teleportation is used as a control operation: the measurement induces a change in the state of the shared pair, and the outcomes are used to further manipulate such state.
It cannot be excluded that, not taking advantage of this lesson, we might again miss other important elements of Nature for a long time.
We are grateful to Professor Jerzy Janik for arranging a stimulating meeting in the Henryk Niewodniczański Institute of Nuclear Physics. We also thank Robert Alicki, Karol Horodecki, Jonathan Oppenheim, Aditi and Ujjwal Sen’s, Andreas Winter and Dong Yang for stimulating discussions. The work is supported by Polish Ministry of Scientific Research and Information Technology under the (solicited) grant no. PBZ-MIN-008/P03/2003 and by EC grants RESQ, contract no. IST-2001-37559 and QUPRODIS, contract no. IST-2001-38877. M. P. acknowledges support from CNR-NATO. |
warning/0506/physics0506038.html | ar5iv | text | # Study of P,T-Parity Violation Effects in Polar Heavy-Atom Molecules
## Introduction
It is well recognized that polar diatomics containing heavy elements are very promising objects for the experimental search for the break of inversion symmetry (P) and time-reversal invariance (T). Though the search for the P,T-parity nonconservation (PNC) effects in heavy atoms and heavy-atom molecules has produced null results up to now, there are serious reasons to search for them with the presently accessible (expected) level of experimental sensitivity. The observation of non-zero P,T-odd effects at this level would indicate the presence of so-called “new physics”Erler:04 beyond the Standard Model (SM) of electroweak and strong interactions Glashow:61 ; Weinberg:67 ; Salam:68 ; Weinberg:72 that is certainly of fundamental importance. Despite well known drawbacks and unresolved problems of the Standard Model (radiative corrections to the Higgs mass are quadratically divergent; rather artifi ial Higgs mechanism of symmetry breaking is not yet verified in experiment; the problem of CP-violation is not well understood, where “C” is charge conjugation symmetry etc.) there are no experimental data available which would be in direct contradiction with this theory (see section 2 and papers Commins:99 ; Erler:04 for more details and references). In turn, some popular extensions of the Standard Model, which allow one to overcome its disadvantages, are not confirmed experimentally.
A crucial feature of PNC experiments in atoms, molecules, liquids or solids is that for interpretation of measured data in terms of fundamental constants of the P,T-odd interactions, one must calculate those properties of the systems, which establish a connection between the measured data and studied fundamental constants (see section 3). These properties are described by operators heavily concentrated near or on heavy nuclei; they cannot be measured and their theoretical study is not a trivial task. During the last several years the significance of (and requirement for) ab initio calculation of electronic structure providing a high level of reliability and accuracy in accounting for both relativistic and correlation effects has only increased (see sections 2 and 9).
The main goal of the paper is to discuss the present status of relativistic calculations of P,T-odd properties in heavy-atom molecules, the two-step methodology used in these calculations, and the accuracy of this method. The historical background of the PNC study in atoms and molecules, its current status and some general remarks on the PNC experiments are presented in sections 1, 2 and 3, correspondingly. The ab initio relativistic methods and the two-step techniques of calculation designed for studying PNC properties in heavy-atom molecules are discussed in sections 4 and 5. The calculations of PNC properties and hyperfine structure in molecules PbF, HgF, YbF, BaF, TlF and PbO are presented in sections 69. Concluding remarks are outlined in section 10.
## 1 Study of P- and T-parity nonconservation effects in heavy-atom molecules: Historical background
After discovery of the combined charge and space parity violation, or CP-violation, in $`K_L^0`$-meson decay Christenson:64 , the search for the electric dipole moments (EDMs) of elementary particles has become one of the most fundamental problems in physics Commins:99 ; Sapirstein:02aa ; Berger:04 ; Ginges:04 ; Erler:04 . A permanent EDM is induced by the weak interaction that breaks both the space symmetry inversion and time-reversal invariance Landau:57 . Considerable experimental effort has been invested in probing for atomic EDMs induced by EDMs of the proton, neutron and electron, and by P,T-odd interactions between them. The best available restriction for the electron EDM, $`d_e`$, was obtained in the atomic Tl experiment Regan:02 , which established an upper limit of $`|d_e|<1.6\times 10^{27}e`$cm, where $`e`$ is the charge of the electron. The benchmark upper limit on a nuclear EDM is obtained in atomic experiment on <sup>199</sup>Hg Romalis:01 , $`|d_{\mathrm{Hg}}|<2.1\times 10^{28}e`$cm, from which the best restriction on the proton EDM, $`|d_p|<5.4\times 10^{24}e`$cm, was also recently obtained by Dmitriev & Sen’kov Dmitriev:03 (the previous upper limit on the proton EDM was obtained in the TlF experiment, see below).
Since 1967, when Sandars suggested the use of polar heavy-atom molecules in the experimental search for the proton EDM Sandars:67 , molecules have been considered the most promising objects for such experiments. Sandars also noticed earlier Sandars:65 that the P- and P,T-parity nonconservation effects are strongly enhanced in heavy atoms due to relativistic and other effects. For example, in paramagnetic atoms the enhancement factor for an electron EDM, $`d_{\mathrm{atom}}/d_e`$, is roughly proportional to $`\alpha ^2Z^3\alpha _D`$, where $`\alpha 1/137`$ is the fine structure constant, $`Z`$ is the nuclear charge and $`\alpha _D`$ is the atomic polarisability. It can be of order 100 or greater for highly polarizable heavy atoms ($`Z50`$). Furthermore, the effective intramolecular electric field acting on electrons in polar molecules can be five or more orders of magnitude higher than the maximal field accessible in a laboratory. The first molecular EDM experiment was performed on TlF by Sandarset al. Hinds:76 (Oxford, UK); it was interpreted as a search for the proton EDM and other nuclear P,T-odd effects. In 1991, in the last series of the <sup>205</sup>TlF experiments by Hindset al. Cho:91 (Yale, USA), the restriction $`d_p=(4\pm 6)\times 10^{23}e\mathrm{cm}`$ was obtained (this was recalculated in 2002 by Petrovet al. Petrov:02 as $`d_p=(1.7\pm 2.8)\times 10^{23}e\mathrm{cm}`$).
In 1978 the experimental investigation of the electron EDM and other PNC effects was further stimulated by Labzowskyet al. Labzowsky:78 ; Gorshkov:79 and Sushkov & Flambaum Sushkov:78 who clarified the possibilities of additional enhancement of these effects in diatomic radicals like BiS and PbF due to the closeness of levels of opposite parity in $`\mathrm{\Omega }`$-doublets having a $`{}_{}{}^{2}\mathrm{\Pi }_{1/2}^{}`$ ground state. Then Sushkovet al. Sushkov:84 and Flambaum & Khriplovich Flambaum:85b suggested the use of $`\mathrm{\Omega }`$-doubling in diatomic radicals with a $`{}_{}{}^{2}\mathrm{\Sigma }_{1/2}^{}`$ ground state for such experiments and the HgF, HgH and BaF molecules were first studied semiempirically by Kozlov Kozlov:85 . At the same time, the first two-step ab initio calculation of PNC effects in PbF initiated by Labzowsky was finished by Titovet al. Titov:85Dism ; Kozlov:87 . A few years later, Hinds started an experimental search for the electron EDM in the YbF molecule, on which the first result was obtained by his group in 2002 (Sussex, UK) Hudson:02 , $`d_e=(0.2\pm 3.2)\times 10^{26}e\mathrm{cm}`$. Though that restriction is worse than the best current $`d_e`$ datum (from the Tl experiment, see above), nevertheless, it is limited by only counting statistics, as Hindset al. pointed out in Hudson:02 .
A new series of electron EDM experiments on YbF by Hinds’ group (Imperial College, UK) are in progress and a new generation of electron EDM experiments using a vapor cell, on the metastable $`a(1)`$ state of PbO, is being prepared by the group of DeMille (Yale, USA). The unique suitability of PbO for searching for the elusive $`d_e`$ is demonstrated by the very high projected statistical sensitivity of the Yale experiment to the electron EDM. In prospect, it allows one to detect $`d_e`$ of order of $`10^{29}÷10^{31}e`$cm DeMille:00 , two–four orders of magnitude lower than the current limit quoted above. Some other candidates for the EDM experiments, in particular, HgH, HgF, TeO, and HI<sup>+</sup> are being discussed and an experiment on PbF is planned (Oklahoma Univ., USA).
## 2 Present status of the electron EDM search
As is mentioned in the introduction, the observation of a non-zero EDM would point out the presence of so called “new physics” (see Okun:82 ; Erler:04 and references) beyond the Standard Model Glashow:61 ; Weinberg:67 ; Salam:68 ; Weinberg:72 ; Kobayashi:73 or CP violation in the QCD sector of SM, $`SU(3)_C`$. The discovery of a lepton EDM (electron EDM in our case) would have an advantage as compared to the cases of neutron or proton EDMs because the latter are not considered as elementary particles within the SM and its extensions.
In Table 1 some estimates for the electron EDM predicted by different theoretical models are given (e.g., see Commins:99 for more details). One can see from the table that the most conservative estimate is given by the Standard Model. This is explained by severe cancellations and suppressions of the contributions producing the electron EDM within the SM. In turn, the “new physics” (extensions of the Standard Model: supersymmetry (SUSY) Kazakov:00 ; Mohapatra:03 ; Erler:04 multi-Higgs Barr:92a ; Barr:93b ; Ginzburg:04 , left-right symmetry Pati:74 ; Barr:93b ; Mohapatra:03 , lepton flavor-changing Liu:94 ; Masina:04 etc.) is very sensitive to the EDMs of elementary particles. This is especially true for the minimal (“naive”) SUSY model, which predicts an electron EDM already at the level of $`10^{25}e\mathrm{cm}`$. However, the best experimental estimate on the electron EDM, $`1.6\times 10^{27}e\mathrm{cm}`$, obtained in the experiment on the Tl atom Regan:02 , is almost two orders of magnitude smaller. More sophisticated SUSY models (which are extremely popular among theorists because they allow one to overcome serious theoretical drawbacks of SM, explain the “gauge hierarchy problem”, solve the problem of dark matter in astrophysics etc.) still predict the electron EDM at the level of $`10^{27}e\mathrm{cm}`$ or somewhat smaller. Since the Tl experiment is finished now, an intriguing expectation is connected with the ongoing experiment on the $`a(1)`$ state of the PbO molecule, which is expected to be sensitive to the electron EDM at least two orders of magnitude smaller. Thus, the most popular extensions of SM can be severely examined by this experiment, i.e. even the result compatible with zero will dramatically influence their status.
## 3 General remarks on experimental search for EDMs in atoms and molecules
The experiments to search for EDMs in atoms and molecules are carried out using different approaches Khriplovich:97 ; Commins:99 . The experimental technique depends on the properties of the atoms and molecules used in such an experiment. These properties influence the atomic and molecular sources, resonance region and detector. For example, for diatomic radicals like YbF or PbF the experiments on molecular beams are most reasonable, while for molecules with closed electronic shells in the ground state like PbO the EDM measurements can be carried out in optical cells.
Nevertheless, the statistical sensitivity of the experiments to the electron or proton EDM usually depends on some parameters common for all such EDM experiments. The easiest way to see this is to apply the Heisenberg uncertainty principle to evaluate the sensitivity of the EDM measurement. Suppose that the EDM of a molecule is measured in some electric field, $`\stackrel{}{E}`$. (Do not confuse the EDM of a polar molecule with the large conventional dipole moment of the molecule, which averages to zero in the absence of external electric field in the laboratory coordinate system. In contrast to the latter, the (vanishingly small) molecular EDM can exist only due to P,T-odd interactions; it is permanent and its direction depends on the sign of the projection of the total electronic momentum on the molecular axis. See Commins:99 for more details.) Thus the energy of interaction of the molecular EDM, $`\stackrel{}{d}=d\stackrel{}{\sigma }`$ (where $`\stackrel{}{\sigma }`$ is a unit vector along the total angular momentum of the molecule), with the electric field is $`\stackrel{}{d}\stackrel{}{E}`$ (linear Stark effect) and the energy difference between the levels with opposite directions of the total angular momentum (leading to the contributions of different signs) is $`2\stackrel{}{d}\stackrel{}{E}`$. If a measurement is carried out by detecting the energy shift during a time $`T`$, the uncertainty in the $`d`$ determination is $`\delta d=\mathrm{}/(2T\stackrel{}{E}\stackrel{}{\sigma })`$. For such measurement on $`N`$ uncorrelated molecules one, obviously, has
$$\delta d=\mathrm{}/(2T\sqrt{N}\stackrel{}{E}\stackrel{}{\sigma })=\mathrm{}/(2TE_\sigma \sqrt{\tau dN/dt}),$$
where $`dN/dt`$ is the counting rate, $`E_\sigma =\stackrel{}{E}\stackrel{}{\sigma }`$, and $`\tau `$ is the total measurement time (usually $`\tau T`$ and $`T`$ is limited by the coherence time of the considered system). Up to now we deal with the molecular EDM $`\stackrel{}{d}`$. Let us write $`d=Gd_e`$, where $`d_e`$ is the value of the electron EDM (the same is valid, of course, for the proton EDM) and $`G`$ is the proportionality coefficient (usually called the enhancement factor). Thus, the final expression for $`\delta d_e`$ is
$$\delta d_e=\frac{\mathrm{}}{2TGE_\sigma \sqrt{\tau dN/dt}}=\frac{\mathrm{}}{2TW\sqrt{\tau dN/dt}},$$
(1)
where the value $`W=GE_\sigma `$ is the effective electric field in the molecule, which can be interpreted as the field that should be applied along the EDM of a free electron to give the energy shift $`2Wd_e2E_\sigma d`$.
From expression (1), the basic conditions which should be met in any prospective EDM experiment can be derived:
1. The counting rate ($`dN/dt`$) should be made as high as possible. From this point of view the experiments on vapor cells, like that planned for PbO, have a clear advantage as compared to beam experiments because molecular vapor density can be usually made much higher than molecular beam density. Thus, in the experiment on the PbO cell the counting rate is estimated to be of order $`10^{11}`$$`10^{15}`$ Hz DeMille:00 , while in the experiment on the YbF molecular beam the counting rate was of order 10<sup>3</sup>—10<sup>4</sup> Hz Hudson:02 .
2. It is crucial to attain high coherence time $`T`$. In the beam experiments that time is just the time of flight through the region with the electric field. For a gas-dynamic molecular source the typical time of flight is $`110`$ ms. On the other hand, for the PbO experiment in vapor cell $`T`$ is close to the lifetime of the excited (metastable) state $`a(1)`$, $`T0.1`$ ms. So, the beam experiments have advantage in the coherence time.
3. It is also critical to have a high value of the effective electric field $`W`$, acting on the electron. The only way to know that parameter is to perform relativistic calculations. It is notable that the first semiempirical estimates of this kind were performed by Sandars in Sandars:65 ; Sandars:67 for Cs and TlF, correspondingly. In these papers the importance of accounting for relativistic effects and using heavy atoms and heavy-atom molecules in EDM experiments was first understood.
The expected energy difference, $`2\stackrel{}{d}\stackrel{}{E}`$ is extremely small even for completely polarized heavy-atom molecules. Thus, in practice, the EDM experiment is usually carried out in parallel and antiparallel electric and magnetic ($`\stackrel{}{B}`$) fields. Interaction energy of the molecular magnetic moment, $`\stackrel{}{\mu }`$, with the magnetic field is much higher than that of the EDM with the electric field and the energy differences are
$$2\stackrel{}{\mu }\stackrel{}{B}+2\stackrel{}{d}\stackrel{}{E}$$
and
$$2\stackrel{}{\mu }\stackrel{}{B}2\stackrel{}{d}\stackrel{}{E}$$
for parallel and antiparallel fields, respectively (in practice, the atomic or molecular spin precession is usually studied instead of direct measurement of the energy shift, see Khriplovich:97 ). When the electric field is reversed, the energy shift, $`4\stackrel{}{d}\stackrel{}{E}=4d_eW`$, points to the existence of the permanent molecular EDM. The same measurement technique is applicable to studying other P,T-odd interactions in atoms and molecules.
The electronic structure parameters describing the P,T-odd interactions of electrons (sections 6, 7, and 9) and nucleons (section 8) including the interactions with their EDMs should be reliably calculated for interpretation of the experimental data. Moreover, ab initio calculations of some molecular properties are usually required even for the stage of preparation of the experimental setup. Thus, electronic structure calculations suppose a high level of accounting for both correlations and relativistic effects (see below). Modern methods of relativistic ab initio calculations (including very recently developed approaches) allow one to achieve the required high accuracy. These approaches will be outlined and discussed in the following sections.
## 4 Heavy-atom molecules: Computational strategies.
The most straightforward method for electronic structure calculation of heavy-atom molecules is solution of the eigenvalue problem using the Dirac-Coulomb (DC) or Dirac-Coulomb-Breit (DCB) Hamiltonians Schwerdtfeger:02bb ; Schwerdtfeger:04aa ; Hirao:04 when some approximation for the four-component wave function is chosen.
However, even applying the four-component single configuration (SCF) approximation, Dirac-Fock (DF) or Dirac-Fock-Breit (DFB), to calculation of heavy-atom molecules (followed by transformation of two-electron integrals to the basis of molecular spinors is not always an easy task) because a very large set of primitive atomic basis functions can be required for such all-electron four-component SCF calculations (see Mosyagin:05a ). Starting from the Pauli approximation and Foldy–Wouthuysen transformation, many different two-component approaches were developed in which only large components are treated explicitly (e.g., see Wood:78 ; Barthelat:80 ; Lenthe:93 ; Wolf:02 and references). In addition, the approaches with perturbative treatment of relativistic effects Kutzelnigg:90 have been developed in which a nonrelativistic wavefunction is used as reference. During the last few years, good progress was also attained in four-component techniques Dyall:02a ; Visscher:02aa ; Grant:04A ; Schwerdtfeger:02bb which allowed one to reduce efforts in calculation and transformation of two-electron matrix elements with small components of four-component molecular spinors. These developments are applied, in particular, in the dirac DIRAC and bertha Quiney:99 ; BERTHA molecular programs. Thus, accurate DC(B) calculations of relatively simple heavy-atom molecules can be performed on modern computers now.
The greatest computational savings are achieved when the two-component relativistic effective core potential (RECP) approximation suggested originally by Leeet al. Lee:77 is used (e.g., see reviews Ermler:88 ; Schwerdtfeger:03 ; YSLee:04 ). There are several reasons for using RECPs (including model potentials and pseudopotentials) in calculations of complicated heavy-atom molecules, clusters and solids. The RECP approaches allow one to exclude the large number of chemically inactive electrons from molecular calculations and to treat explicitly only valence and outermost core electrons from the beginning. Then, the oscillations of the valence spinors are usually smoothed in heavy-atom cores simultaneously with excluding small components from the explicit treatment. As a result the number of primitive basis functions can be reduced dramatically; this is especially important for calculation and transformation of two-electron integrals when studying many-atomic systems and compounds of very heavy elements including actinides and superheavies. The RECP method is based on a well-developed nonrelativistic technique of calculations; however, an effective spin-orbit interaction and other scalar-relativistic effects are taken into account usually by means of radially-local Ermler:88 ; Schwerdtfeger:03 ; YSLee:04 ; Teichteil:04 , separable Blochl:90 ; Vanderbilt:90 ; Theurich:01 or Huzinaga-type Bonifacic:74 ; Katsuki:88b ; Seijo:04 operators.
Correlation molecular calculations with RECPs are naturally performed in the basis of spin-orbitals (and not of spinors as is in all-electron four-component calculations) even for the cases when quantum electrodynamics (two-electron Breit etc.) effects are taken into account Petrov:04b ; Mosyagin:05a . Note, however, that the DCB technique with the separated spin-free and spin-dependent terms also has been developed Dyall:94 , but it can be efficiently applied only in the cases when spin-dependent effects can be neglected both for valence and for core shells. In the RECP method, the interactions with the excluded inner core shells (spinors!) are described by spin-dependent potentials whereas the explicitly treated valence and outer core shells are usually described by spin-orbitals in molecular calculations. It means that some “soft” way of accounting for the core-valence orthogonality constraints is applied in the latter case Titov:99 (note, meantime, that the strict core-valence orthogonality can be retrieved after the RECP calculation by using the restoration procedures described below). Another merit of the RECP method is in its natural ability to account for correlations with the explicitly excluded inner core electrons Mosyagin:04a (this direction is actively developed during last years). The use of the molecular spin-orbitals and the “correlated” RECPs allows one to reduce dramatically the expenses at the stage of correlation calculation of heavy-atom molecules. These are important advantages when a very high level of accounting for correlations is required even in studying diatomics (e.g., see calculations of PbO described in section 9). Thus, many complications of the DC(B) molecular calculations are avoided when employing RECPs.
The “shape-consistent” (or “norm-conserving”) RECP approaches are most widely employed in calculations of heavy-atom molecules though “energy-adjusted/consistent” pseudopotentials Schwerdtfeger:03 by Stuttgart team are also actively used as well as the Huzinaga-type “ab initio model potentials” Seijo:04 . In plane wave calculations of many-atom systems and in molecular dynamics, the separable pseudopotentials Blochl:90 ; Vanderbilt:90 ; Theurich:01 are more popular now because they provide linear scaling of computational effort with the basis set size in contrast to the radially-local RECPs. The nonrelativistic shape-consistent effective core potential was first proposed by Durand & Barthelat Durand:75 and then a modified scheme of the pseudoorbital construction was suggested by Christiansenet al. Christiansen:79 and by Hamannet al. Hamann:79 .
In a series of papers (see Titov:99 ; Titov:00a ; Titov:02Dism ; Petrov:04b ; Mosyagin:04a and references) a generalized RECP approach was developed that involves both radially-local, separable and Huzinaga-type potentials as its components in particular cases. It allows one to attain very high accuracy of calculation of valence properties and electronic densities in the valence region when treating outermost core shells in calculations explicitly (see section 5 for more details).
Nevertheless, calculation of such properties as spin-dependent electronic densities near nuclei, hyperfine constants, P,T-parity nonconservation effects, chemical shifts etc. with the help of the two-component pseudospinors smoothed in cores is impossible. We should notice, however, that the above core properties (and the majority of other properties of practical interest which are described by the operators heavily concentrated within inner cores or on nuclei) are mainly determined by electronic densities of the valence and outer core shells near to, or on, nuclei. The valence shells can be open or easily perturbed by external fields, chemical bonding etc., whereas outer core shells are noticeably polarized (relaxed) in contrast to the inner core shells. Therefore, accurate calculation of electronic structure in the valence and outer core region is of primary interest for such properties.
For evaluation of the matrix elements of the operators concentrated on (or close to) nuclei, proper shapes of the valence molecular four-component spinors must be restored in atomic core regions after performing the RECP calculation of that molecule. In 1959, a nonrelativistic procedure of restoration of the orbitals from smoothed Phillips–Kleinman pseudoorbitals was proposed Phillips:59 based on the orthogonalization of the latter to the original atomic core orbitals. In 1985, Pacios & Christiansen Pacios:85 suggested a modified orthogonalization scheme in the case of shape-consistent pseudospinors. At the same time, a simple procedure of “nonvariational” one-center restoration (NOCR, see below) employing the idea of generation of equivalent basis sets in four-component Dirac-Fock and two-component RECP/SCF calculations was proposed and first applied in Titov:85Dism to evaluation of the P,T-odd spin-rotational Hamiltonian parameters in the PbF molecule. In 1994, a similar procedure was used by Blöchl inside the augmentation regions Blochl:94 in solids to construct the transformation operator between pseudoorbitals (“PS”) and original orbitals (“AE”) in his projector augmented-wave method.
All the above restoration schemes are called “nonvariational” as compared to the “variational” one-center restoration (VOCR, see below) procedure proposed in Titov:92A ; Titov:96 . Proper behavior of the molecular orbitals (four-component spinors) in atomic cores of molecules can be restored in the scope of a variational procedure if the molecular pseudoorbitals (two-component pseudospinors) match correctly the original orbitals (large components of bispinors) in the valence region after the molecular RECP calculation. As is demonstrated in Titov:99 ; Mosyagin:05a , this condition is rather correct when the shape-consistent RECP is involved to the molecular calculation with explicitly treated outermost core orbitals and, especially, when the generalized RECP operator is used since the above matching condition is implemented at their generation.
At the restoration stage, a one-center expansion in the spherical harmonics with numerical radial parts is most appropriate both for orbitals (spinors) and for the description of “external” interactions with respect to the core regions of a considered molecule. In the scope of the discussed two-step methods for the electronic structure calculation of a molecule, finite nucleus models and quantum electrodynamic terms including, in particular, two-electron Breit interaction may be taken into account without problems Petrov:04b .
One-center expansion was first applied to whole molecules by Desclaux & Pyykkö in relativistic and nonrelativistic Hartree-Fock calculations for the series CH<sub>4</sub> to PbH<sub>4</sub> Desclaux:74b and then in the Dirac-Fock calculations of CuH, AgH and AuH Desclaux:76c and other molecules Desclaux:02 . A large bond length contraction due to the relativistic effects was estimated. However, the accuracy of such calculations is limited in practice because the orbitals of the hydrogen atom are reexpanded on a heavy nucleus in the entire coordinate space. It is notable that the RECP and one-center expansion approaches were considered earlier as alternatives to each other Pitzer:79 ; Pyykko:79 .
The applicability of the discussed two-step algorithms for calculation of wavefunctions of molecules with heavy atoms is a consequence of the fact that the valence and core electrons may be considered as two subsystems, interaction between which is described mainly by some integrated properties of these subsystems. The methods for consequent calculation of the valence and core parts of electronic structure of molecules give us a way to combine the relative simplicity and accessibility both of molecular RECP calculations in gaussian basis set, and of relativistic finite-difference one-center calculations inside a sphere with the atomic core radius.
The first two-step calculations of the P,T-odd spin-rotational Hamiltonian parameters were performed for the PbF radical about 20 years ago Titov:85Dism ; Kozlov:87 , with a semiempirical accounting for the spin-orbit interaction. Before, only nonrelativistic SCF calculation of the TlF molecule using the relativistic scaling was carried out Hinds:80a ; Coveney:83 ; here the P,T-odd values were underestimated by almost a factor of three as compared to the later relativistic Dirac-Fock calculations. The latter were first performed only in 1997 by Laerdahlet al. Laerdahl:97 and by Parpia Parpia:97 . The next two-step calculation, for PbF and HgF molecules Dmitriev:92 , was carried out with the spin-orbit RECP part taken into account using the method suggested in Titov:92 .
Later we performed correlation GRECP/NOCR calculations of the core properties in YbF Titov:96b , BaF Kozlov:97 , again in YbF Mosyagin:98 and in TlF Petrov:02 . In 1998, first all-electron Dirac-Fock calculations of the YbF molecule were also performed by Quineyet al. Quiney:98 and by Parpia Parpia:98 . Recently we finished extensive two-step calculations of the P,T-odd properties and hyperfine structure of the excited states of the PbO molecule Isaev:04 ; Petrov:05a . One more two-step calculation of the electron EDM enhancement effect was performed very recently for the molecular ion HI<sup>+</sup> Isaev:04b .
We would like to emphasize here that the all-electron Dirac-Fock calculations on TlF and YbF are, in particular, important for checking the quality of the approximations made in the two-step method. The comparison of appropriate results, Dirac-Fock vs. RECP/SCF/NOCR, is, therefore, performed in papers Mosyagin:98 ; Petrov:02 and discussed in the present paper.
In this paper, the main features of the two-step method are presented and PNC calculations are discussed, both those without accounting for correlation effects (PbF and HgF) and those in which electron correlations are taken into account by a combined method of the second-order perturbation theory (PT2) and configuration interaction (CI), or “PT2/CI” Dzuba:96 (for BaF and YbF), by the relativistic coupled cluster (RCC) method Kaldor:97 ; Landau:01c (for TlF, PbO, and HI<sup>+</sup>), and by the spin-orbit direct-CI method Buenker:99 ; Alekseyev:04 ; Titov:01 (for PbO). In the ab initio calculations discussed here, the best accuracy of any current method has been attained for the hyperfine constants and P,T-odd parameters regarding the molecules containing heavy atoms.
## 5 Two-step method of calculation of core properties
The two-step method consists of a two-component molecular RECP calculation at the first step, followed by restoration of the proper four-component wave function in atomic cores at the second step. Though the method was developed originally for studying core properties in heavy-atom molecules, it can be efficiently applied to studying the properties described by the operators heavily concentrated in cores or on nuclei of light atoms in other computationally difficult cases, e.g., in many-atom molecules and solids. The details of these steps are described below.
### Generalized RECP
When core electrons of a heavy-atom molecule do not play an active role, the effective Hamiltonian with RECP can be presented in the form
$$𝐇^{\mathrm{Ef}}=\underset{i_v}{}[𝐡^{\mathrm{Schr}}(i_v)+𝐔^{\mathrm{Ef}}(i_v)]+\underset{i_v>j_v}{}\frac{1}{r_{i_vj_v}}.$$
(2)
This Hamiltonian is written only for a valence subspace of electrons which are treated explicitly and denoted by indices $`i_v`$ and $`j_v`$. In practice, this subspace is often extended by inclusion of some outermost core shells for better accuracy but we will consider them as the valence shells below if these outermost core and valence shells are not treated using different approximations. In Eq. (2), $`𝐡^{\mathrm{Schr}}`$ is the one-electron Schrödinger Hamiltonian
$$𝐡^{\mathrm{Schr}}=\frac{1}{2}\stackrel{}{}^2\frac{Z_{ic}}{r},$$
(3)
where $`Z_{ic}`$ is the charge of the nucleus decreased by the number of inner core electrons. $`𝐔^{\mathrm{Ef}}`$ in (2) is an RECP (relativistic pseudopotential) operator that is usually written in the radially-local (semi-local) Ermler:88 or separable (e.g., see Theurich:01 and references) approximations when the valence pseudospinors are smoothed in the heavy-atom cores. Contrary to the four-component wave function used in Dirac-Coulomb(-Breit) calculations, the pseudo-wave function in the RECP case can be both two- and one-component. The use of the effective Hamiltonian (2) instead of all-electron four-component Hamiltonians leads to the question about its accuracy. It was shown both theoretically and in many calculations (see Titov:99 ; Titov:02Dism and references) that the typical accuracy of the radially-local RECPs is within 1000–3000 cm<sup>-1</sup> for transition energies between low-lying states though otherwise is sometime stated (see Dolg:00aa ; Titov:00bmin ).
In our two-step calculations the generalized RECP operator Titov:99 ; Titov:00a is used that includes the operators of radially-local shape-consistent RECP, separable pseudopotential and Huzinaga-type model potential as its components. Additionally, the GRECP operator can include terms of other types, known as “self-consistent” and two-electron “term-splitting” corrections Titov:95 ; Titov:99 ; Titov:00a , which are important particularly for economical (but precise!) treatment of transition metals, lanthanides and actinides. With these terms, the accuracy provided by GRECPs can be even higher than the accuracy of the frozen core approximation (employing the same number of explicitly treated electrons) because they can account for relaxation of explicitly excluded (inner core) electrons Titov:99 ; Titov:02Dism . The theoretical background of the GRECP concept is developed in a series of papers Titov:99 ; Titov:00a ; Titov:02Dism ; Petrov:04b ; Mosyagin:05a ; Mosyagin:04a . In contrast to other RECP methods, GRECP employs the idea of separating the space around a heavy atom into three regions: inner core, outer core and valence, which are first treated employing different approximations for each. It allows one to attain practically any desired accuracy, while requiring moderate computational efforts since the overall accuracy is limited in practice by possibilities of correlation methods.
When innermost core shells must be treated explicitly, the four-component versions of the GRECP operator can be used, in principle, together with the all-electron relativistic Hamiltonians. The GRECP can describe here some quantum electrodynamics effects (self-energy, vacuum polarization etc.) thus avoiding their direct treatment. One more remark is that the two-component GRECP operator can be applied even to very light atoms when smoothing the large components of the four-component spinors only in the vicinity of the nucleus just to account for relativistic effects (the GRECP for hydrogen provides accuracy of treatment of very small relativistic contributions within 5%).
### Nonvariational One-Center Restoration
The electronic densities evaluated from the two-component pseudo-wave function very accurately reproduce the corresponding all-electron four-component densities in the valence and outer core regions not only for the state used in the GRECP generation but also for other states which differ by excitations of valence electrons. This is illustrated in figure 1 (see also tables 8 and 9 in Mosyagin:05a ), where the radial parts of the large components of the thallium bispinor and the corresponding pseudospinor are compared for the state averaged over the relativistic $`6s_{1/2}^26p_{1/2}^1`$ configuration, whereas the 21-electron GRECP is generated for the state averaged over the nonrelativistic $`6s^16p^16d^1`$ configuration. That is true also for the electronic densities obtained in the corresponding Dirac-Coulomb(-Breit) and GRECP calculations accounting for correlation.
In the inner core region, the pseudospinors are smoothed, so that the electronic density with the pseudo-wave function is not correct. When operators describing properties of interest are heavily concentrated near or on nuclei, their mean values are strongly affected by the wave function in the inner region. The four-component molecular spinors must, therefore, be restored in the heavy-atom cores.
All molecular spinors $`\varphi _p`$ can be restored as one-center expansions in the cores using the nonvariational one-center restoration (NOCR) scheme Titov:85Dism ; Kozlov:87 ; Dmitriev:92 ; Titov:96b ; Kozlov:97 ; Mosyagin:98 ; Petrov:02 ; Isaev:04 ; Petrov:05a that consists of the following steps:
* Generation of equivalent basis sets of one-center four-component spinors $`\left(\begin{array}{c}f_{nlj}(r)\chi _{ljm}\\ g_{nlj}(r)\chi _{2jl,jm}\end{array}\right)`$ and smoothed two-component pseudospinors $`\stackrel{~}{f}_{nlj}(r)\chi _{ljm}`$ in finite-difference all-electron Dirac-Fock(-Breit) and GRECP/SCF calculations of the same configurations of a considered atom and its ions. The nucleus is usually modeled by a uniform charge distribution within a sphere. The all-electron four-component hfdb HFDB ; Bratzev:77 ; Tupitsyn:02A and two-component grecp/hfj HFJ ; Tupitsyn:95 codes are employed to generate two equivalent numerical basis sets used at the restoration. These sets, describing mainly the atomic core region, are generated independently of the basis set used for the molecular GRECP calculations.
* The molecular pseudospinorbitals are then expanded in the basis set of one-center two-component atomic pseudospinors (for $`rR_{\mathrm{nocr}}`$, where $`R_{\mathrm{nocr}}`$ is the radius of restoration that should be sufficiently large for calculating core properties with required accuracy),
$$\stackrel{~}{\varphi }_p(𝐱)\underset{l=0}{\overset{L_{max}}{}}\underset{j=|l1/2|}{\overset{j=|l+1/2|}{}}\underset{n,m}{}c_{nljm}^p\stackrel{~}{f}_{nlj}(r)\chi _{ljm},$$
(4)
where $`𝐱`$ denotes spatial and spin variables. Note that for linear molecules only one value of $`m`$ survives in the sum for each $`\varphi _p`$.
* Finally, the atomic two-component pseudospinors in the molecular basis are replaced by equivalent four-component spinors and the expansion coefficients from Eq. (4) are preserved:
$$\varphi _p(𝐱)\underset{l=0}{\overset{L_{\mathrm{max}}}{}}\underset{j=|l1/2|}{\overset{j=|l+1/2|}{}}\underset{n,m}{}c_{nljm}^p\left(\begin{array}{c}f_{nlj}(r)\chi _{ljm}\\ g_{nlj}(r)\chi _{2jl,jm}\end{array}\right).$$
(5)
The molecular four-component spinors constructed this way are orthogonal to the inner core spinors of the atom, because the atomic basis functions used in Eq. (5) are generated with the inner core shells treated as frozen.
### Variational one-center restoration
In the variational technique of one-center restoration (VOCR) Titov:92A ; Titov:96 , the proper behavior of the four-component molecular spinors in the core regions of heavy atoms can be restored as an expansion in spherical harmonics inside the sphere with a restoration radius, $`R_{\mathrm{vocr}}`$, that should not be smaller than the matching radius, $`R_c`$, used at the RECP generation. The outer parts of spinors are treated as frozen after the RECP calculation of a considered molecule. This method enables one to combine the advantages of two well-developed approaches, molecular RECP calculation in a gaussian basis set and atomic-type one-center calculation in numerical basis functions, in the most optimal way. This technique is considered theoretically in Titov:96 and some results concerning the efficiency of the one-center reexpansion of orbitals on another atom can be found in Titov:02Dism .
The VOCR scheme can be used for constructing the core parts of the molecular spinors (orbitals) instead of using equivalent basis sets as in the NOCR technique. A disadvantage of the NOCR scheme is that molecular pseudoorbitals are usually computed in a spin-averaged GRECP/SCF molecular calculation (i.e. without accounting for the effective spin-orbit interaction) and the reexpansion of molecular pseudospinorbitals on one-center atomic pseudospinors can be restricted in accuracy, as was noticed in the GRECP/RCC/NOCR calculations Petrov:02 of the TlF molecule (see below). With the restored molecular bispinors, the two-electron integrals on them can be easily calculated. Thus, the four-component transfomation from the atomic basis that is a time-consuming stage of four-component calculations of heavy-atom molecules can be avoided. Besides, the VOCR technique developed in Titov:96b for molecular pseudospinors can be reformulated for the case of molecular pseudospinorbitals to reduce the complexity of the molecular GRECP calculation as is discussed in section 4.
However, the most interesting direction in the development of the two-step method is the possibility to “split” the correlation structure calculation of a molecule into two sequential correlation calculations: first, in the valence region, where the outer core and valence electrons are explicitly involved in the GRECP calculation; and then, in the core region, when the outer and inner core space regions are only treated at the restoration stage. Correlation effects in the valence and outer core regions (not only valence parts of molecular orbitals but also configuration coefficients) can be calculated, for example, by a combination of RCC and CI methods (see section 9) with very high accuracy. Then correlation effects in the inner and outer core regions (including the dipole-type relaxation of atomic inner core shells in a molecule) can be taken into account using the Dirac-Coulomb(-Breit) Hamiltonian and the one-center expansion. By increasing the one-center restoration radius $`R_{\mathrm{vocr}}`$ , one can take into account correlation effects in the intermediate region (outer core in our case) with the required accuracy. Roughly speaking, the computational efforts for the correlation structure calculations in the core and valence regions are added when using the two-step approach, whereas in the conventional one-step scheme, they have multiplicative nature.
### Two-step calculation of molecular properties
To evaluate one-electron core properties (hyperfine structure, P,T-odd effects etc.) employing the above restoraton schemes it is sufficient to obtain the one-particle density matrix, $`\{\stackrel{~}{D_{pq}}\}`$, after the molecular RECP calculation in the basis of pseudospinors $`\{\stackrel{~}{\varphi }_p\}`$. At the same time, the matrix elements $`\{W_{pq}\}`$ of a property operator $`𝑾(𝐱)`$ should be calculated in the basis of equivalent four-component spinors $`\{\varphi _p\}`$. The mean value for this operator can be then evaluated as
$$𝑾=\underset{pq}{}\stackrel{~}{D_{pq}}W_{pq}.$$
(6)
Only the explicitly treated valence shells are used for evaluating a core property when applying Eq. (6) since the atomic frozen core (closed) shells do not usually contibute to the properties of practical interest. However, correlations with the core electrons explicitly excluded from the RECP calculation can be also taken into account if the effective operator technique Lindgren:84 is applied to calculate the property matrix elements $`\{W_{pq}^{\mathrm{Ef}}\}`$ in the basis set of bispinors $`\{\varphi _p\}`$. Then, in expression (6) one should only replace $`\{W_{pq}\}`$ by $`\{W_{pq}^{\mathrm{Ef}}\}`$. Alternatively, the correlations with the inner core electrons can be, in principle, considered within the variational restoration scheme for electronic structure in the heavy atom cores. Strictly speaking, the use of the effective operators is correct when the molecular calculation is carried out with the “correlated” GRECP (see Mosyagin:04a ), in which the same correlations with the excluded core electrons are taken into account at the GRECP generation as they are in constructing $`\{W_{pq}^{\mathrm{Ef}}\}`$. Nevertheless, the use of the (G)RECP that does not account for the core correlations at the molecular calculation stage can be justified if these correlations do not influence dramatically the electronic structure in the valence region. The latter approximation was applied in the calculations of YbF and BaF molecules described in the following section.
When contributions to $`𝑾`$ are important both in the core and valence regions, the scheme for evaluating the mean value of $`𝑾(𝐱)`$ can be as follows:
* calculation of matrix elements with the molecular pseudospinorbitals (which are usually linear combinations of atomic gaussians) over the entire space region using conventional codes for molecular RECP calculations (although the operator $`𝑾`$ can be presented in different forms in calculations with the RECP and Dirac-Coulomb(-Breit) Hamiltonians),
$$\stackrel{~}{𝑾}=\underset{pq}{}\stackrel{~}{D_{pq}}\underset{r<\mathrm{}}{}\stackrel{~}{\varphi }_p(𝐱)𝑾(𝐱)\stackrel{~}{\varphi }_q(𝐱)𝑑𝐱;$$
(7)
* calculation of the inner part of the matrix element of the operator with the same molecular pseudospinorbitals using the one-center expansion for $`\{\stackrel{~}{\varphi }_p\}`$ ($`R_{\mathrm{ocr}}`$ stands for $`R_{\mathrm{nocr}}`$ or $`R_{\mathrm{vocr}}`$ below, $`R_{\mathrm{ocr}}R_c`$):
$$\stackrel{~}{𝑾}^<=\underset{pq}{}\stackrel{~}{D_{pq}}\underset{r<R_{\mathrm{ocr}}}{}\stackrel{~}{\varphi }_p(𝐱)𝑾(𝐱)\stackrel{~}{\varphi }_q(𝐱)𝑑𝐱;$$
(8)
* calculation of the inner part of the matrix element of the operator with the restored molecular four-component spinors using the one-center expansion for $`\{\varphi _p\}`$:
$$𝑾^<=\underset{pq}{}\stackrel{~}{D_{pq}}\underset{r<R_{\mathrm{ocr}}}{}\varphi _p^<(𝐱)𝐖(𝐱)\varphi _𝐪^<(𝐱)\mathrm{𝐝𝐱}.$$
(9)
The matrix element $`𝑾`$ of the $`𝑾(𝐱)`$ operator is evaluated as
$$𝑾=\stackrel{~}{𝑾}\stackrel{~}{𝑾}^<+𝑾^<.$$
(10)
Obviously, the one-center basis functions can be numerical (finite-difference) for better flexibility.
The mean values of the majority of operators for the valence properties can be calculated with good accuracy without accounting for the inner parts of the matrix elements (8) and (9). As discussed above, for calculating the mean values of the operators heavily concentrated on or near nuclei it is sufficient to account only for the inner parts (9) of the matrix elements of $`𝑾(𝐱)`$ on the restored functions $`\varphi _p^<(𝐱)`$, whereas the other contributions, (7) and (8), can be neglected.
Calculation of properties using the finite-field method Kunik:71 ; Monkhorst:77 and Eq. (6) within the coupled-cluster approach is described in section 8.
## 6 Calculations of PbF and HgF
The ground states of the diatomic radicals PbF and HgF are $`{}_{}{}^{2}\mathrm{\Pi }_{1/2}^{}`$ and $`{}_{}{}^{2}\mathrm{\Sigma }_{1/2}^{}`$, respectively. Here the superscript denotes spin multiplicity, $`\mathrm{\Pi }`$ and $`\mathrm{\Sigma }`$ are the projections of the electron orbital angular momentum on the molecular axis and the subscript is the projection of the total electron angular momentum on the molecular axis directed from the heavy atom to fluorine. It is convenient to describe the spin-rotational spectrum of the ground electronic state in terms of the effective spin-rotational Hamiltonian $`𝐇_{\mathrm{eff}}^{sr}`$, following Dmitriev:92 ; Kozlov:95 :
$`\begin{array}{cc}\hfill 𝐇_{\mathrm{eff}}^{sr}& =B\stackrel{}{J}^2+\mathrm{\Delta }\stackrel{}{S}^{^{}}\stackrel{}{J}+\stackrel{}{S}^{^{}}𝐀\stackrel{}{I}\hfill \\ & +\mu _0\stackrel{}{S}^{^{}}𝐆\stackrel{}{B}D\stackrel{}{\lambda }\stackrel{}{E}\hfill \\ & +W_\mathrm{A}k_\mathrm{A}\stackrel{}{\lambda }\times \stackrel{}{S}^{^{}}\stackrel{}{I}+(W_dd_e+W_\mathrm{S}k_\mathrm{S})\stackrel{}{S}^{^{}}\stackrel{}{\lambda }\hfill \end{array}`$ (14)
The first line in this expression describes the rotational structure with $`\omega `$\- or spin-doubling and the hyperfine interaction of the effective electron spin $`\stackrel{}{S}^{^{}}`$ with the nuclear spin $`\stackrel{}{I}`$. $`B`$ is the rotational constant, $`\stackrel{}{J}`$ is the electron-rotational angular momentum, $`\mathrm{\Delta }`$ is the $`\omega `$-doubling constant. The second line describes the interaction of the molecule with the external fields $`\stackrel{}{B}`$ and $`\stackrel{}{E}`$, ($`\stackrel{}{\lambda }`$ is the unit vector directed from the heavy nucleus to the light one). The last line corresponds to the P-odd electromagnetic interaction of the electrons with the anapole moment of the nucleus described by the constant $`k_\mathrm{A}`$ Khriplovich:97 , P,T-odd interaction of the electron EDM $`d_e`$ with the interamolecular field, and P,T-odd scalar interactions of the electrons with the heavy nucleus Dmitriev:92 .
The parameter $`\mathrm{\Delta }`$, tensors A and G, molecular dipole moment $`D`$ and the constants $`W_i`$ are expressed in terms of one-electron matrix elements; concrete expressions for the above parameters can be found in Dmitriev:92 , and for $`W_d`$ and $`A_{}`$ they are also given in the next sections. The results of the calculations are presented in Table 2.
In Dmitriev:92 the conclusion was made, that the accuracy in calculations of the parameters of the effective spin-rotational Hamiltonian is close to 20$`\%`$. However, further ab initio calculations showed the situation is more complicated.
As was understood in calculations of YbF Titov:96b , a fortuitous cancellation of effects of different types took place in the above calculations. Accounting for the polarization of the $`5s,5p`$-shells leads to an enhancement of the contributions of the valence shells to the $`A_{}`$, $`A_{}`$ and PNC constants on the level of 50% of magnitude. A contribution of comparable magnitude but of opposite sign takes place when the relaxation-correlation effects of the $`5d`$-shell are taken into account. This was confirmed in Mosyagin:XXa when accounting for electron correlation both in the valence and core regions of HgF as compared to the YbF case.
## 7 Calculations of YbF and BaF
The results of two-step calculations for the YbF molecule (1996,1998) Titov:96b ; Mosyagin:98 and for the BaF molecule (1997) Kozlov:97 are presented in Table 3 Table 3, where they are compared with other semiempirical and four-component results. For the isotropic hyperfine constant $`A=(A_{}+2A_{})/3`$, the accuracy of our calculation is about 3%, as determined by comparison to the experimental datum. The dipole constant $`A_\mathrm{d}=(A_{}A_{})/3`$ (which is much smaller in magnitude), though better than in all previous calculations known from the literature, is still underestimated by almost 23%. After a semiemprical correction for a perturbation of $`4f`$-shell in the core of Yb due to the bond making, this error is reduced to 8%. Our value for the effective electric field on the unpaired electron is $`W=W_d|\stackrel{}{S}^{^{}}\stackrel{}{\lambda }|`$ = 4.9 a.u.= $`2.5\times 10^{10}`$ V cm<sup>-1</sup> (see section 3 and Eq. (14)).
In Table 3 the values of the $`W_d`$ constant
$$W_dd_\mathrm{e}=2Wd_\mathrm{e}=2^2\mathrm{\Sigma }_{1/2}|\underset{i}{}H_d(i)|^2\mathrm{\Sigma }_{1/2},$$
(15)
where $`H_d`$ describes interaction of the electron EDM $`d_\mathrm{e}`$ with the internal molecular electric field $`𝐄^{\mathrm{mol}}`$:
$`H_d=2d_\mathrm{e}\left(\begin{array}{cc}0& 0\\ 0& \sigma \end{array}\right)𝐄^{\mathrm{mol}},`$ (18)
from various calculations are tabulated. These include the unrestricted Dirac-Fock calculation of Parpia (1998) Parpia:98 , four-component calculations of Quineyet al. (1998) Quiney:98 accounting for correlation, the most recent semiempirical calculation of Kozlov (1997) Kozlov:94 and our latest GRECP/RASSCF/EO calculation (EO stands for the effective operator technique in the framework of the second-order perturbation theory; see section 5 and paper Dzuba:96 for more details). All results are in very close agreement now. It should be noted that the valence electron contribution to $`W_d`$ in Parpia:98 differs by only 7.4% from the corresponding RECP-based calculation of Titovet al. (1996) Titov:96b .
A similar increase in the values for the hyperfine constants and parameters of the P,T-odd interactions when the correlations with the core shells (primarily, $`5s,5p`$) are taken into account is also observed for the BaF molecule Kozlov:97 , as one can see in Table 3. Of course, the corrections from the $`4f`$-electron excitations are not required for this molecule. The enhancement factor for the P,T-odd effects in BaF is three times smaller than in YbF mainly because of the smaller nuclear charge of Ba.
## 8 Calculation of <sup>205</sup>TlF molecule.
### Effective Hamiltonian.
Here we consider the P,T-odd interaction of the <sup>205</sup>Tl nucleus (which has one unpaired proton) with the electromagnetic field of the electrons in the <sup>205</sup>TlF molecule Petrov:02 . This molecule is one of the best objects for measurements of the proton EDM. The effective interaction with the EDM of the Tl nucleus in TlF is written in the form
$$H^{\mathrm{eff}}=(d^V+d^M)\stackrel{}{I}/I\stackrel{}{\lambda },$$
(19)
where $`\stackrel{}{I}`$ is the Tl nuclear spin operator; $`\stackrel{}{\lambda }`$ is the unit vector along $`z`$ (from Tl to F); $`d^V`$ and $`d^M`$ are the volume and magnetic constants Schiff:63 :
$$d^V=6SX=(d_pR+Q)X,$$
(20)
where $`S`$ is the nuclear Schiff moment; $`d_p`$ is the proton EDM; $`R`$ and $`Q`$ are the factors determined by the nuclear structure of <sup>205</sup>Tl (see Petrov:02 );
$$X=\frac{2\pi }{3}\left[\frac{}{z}\rho _\psi (\stackrel{}{r})\right]_{x,y,z=0};$$
(21)
$`\rho _\psi (\stackrel{}{r})`$ is the electronic density calculated from the electronic wavefunction $`\psi `$,
$$d^M=2\sqrt{2}(d_p+d_N)\left(\frac{\mu }{Z}+\frac{1}{2mc}\right)M,$$
(22)
where $`d_N`$ is the nuclear EDM arising due to P,T-odd forces between the nucleons; $`\mu `$, $`m`$ and $`Z`$ are the magnetic moment, mass and charge of the Tl nucleus; $`c`$ is the velocity of light;
$$M=\frac{1}{\sqrt{2}}\psi |\underset{i}{}\left(\frac{\stackrel{}{\alpha }_i\times \stackrel{}{𝐥}_i}{r_i^3}\right)_z|\psi ;$$
(23)
$`\stackrel{}{𝐥_𝐢}`$ is the orbital momentum of $`i`$-th electron; $`\stackrel{}{\alpha }_i`$ are its Dirac matrices. Accounting for $`H_{\mathrm{eff}}`$ leads to a difference in the hyperfine splitting of TlF in parallel and antiparallel electric and magnetic fields. The level shift $`h\nu =4(d^V+d^M)\stackrel{}{I}\stackrel{}{\lambda }`$/I is measured experimentally (for the latest results see Cho:91 ).
The parameters $`X`$ of Eq. (21) and $`M`$ of Eq. (23) are determined by the electronic structure of the molecule. They were calculated in 1997 for the ground $`0^+`$ state of TlF by Parpia Parpia:97 and by Laerdahl, Saue, Faegri Jr., and Quiney Laerdahl:97 using the Dirac–Fock method with gaussian basis sets of large sizes (see Table 4). Below we refer to paper Quiney:98b with the calculations presented in details and not to the brief communication Laerdahl:97 of the same authors. There was also a preliminary calculation of electronic structure in TlF performed by Wilsonet al. in Wilson:94 . In paper Petrov:02 the first calculation of <sup>205</sup>TlF that accounts for correlation effects was performed (see also Dzuba:02 where the limit on the Schiff moment of <sup>205</sup>Tl was recalculated).
### Results.
Calculations were carried out at two internuclear separations, the equilibrium $`R_e=2.0844`$ Å as in Ref. Parpia:97 , and 2.1 Å, for comparison with Ref. Quiney:98b . The relativistic coupled cluster (RCC) method Kaldor:99 ; Kaldor:04b with only single (RCC-S) or with single and double (RCC-SD) cluster amplitudes is used (for review of different coupled cluster approaches see also Paldus:99 ; Paldus:03 and references). The RCC-S calculations with the spin-dependent GRECP operator take into account effects of the spin-orbit interaction at the level of the one-configurational SCF-type method. The RCC-SD calculations include, in addition, the most important electron correlation effects.
The results obtained with the one-center expansion of the molecular spinors in the Tl core in either $`s;p`$, $`s;p;d`$ or $`s;p;d;f`$ partial waves are collected in Table 4. The first point to notice is the difference between spin-averaged SCF values and RCC-S values; the latter include spin-orbit interaction effects. These effects increase $`X`$ by 9% and decrease $`M`$ by 21%. The RCC-S function can be written as a single determinant, and results may therefore be compared with DF values, even though the RCC-S function is not variational. The GRECP/RCC-S values of $`M`$ indeed differ only by 1–3% from the corresponding DF values Parpia:97 ; Quiney:98b (see Table 4).
Much larger differences occur for the $`X`$ parameter. There are also large differences between the two DF calculations for $`X`$, which cannot be explained by the small change in internuclear separation. The value of $`X`$ may be expected to be less stable than $`M`$ Quiney:98b . The conclusion in Petrov:02 is that the RCC-S value for $`X`$, which is higher than that of Parpia:97 ; Quiney:98b , is more correct. The electron correlation effects are calculated by the RCC-SD method at the molecular equilibrium internuclear distance $`R_e`$. A major correlation contribution is observed, decreasing $`M`$ by 17% and $`X`$ by 22%.
The hyperfine structure constants of Tl $`6p_{1/2}^1`$ and Tl<sup>2+</sup> $`6s^1`$, which (like $`X`$ and $`M`$) depend on operators concentrated near the Tl nucleus, were also calculated. The errors in the DF values are 10–15% with respect to experiment and the RCC-SD results are within 1–4% of experiment. The improvement in $`X`$ and $`M`$ upon inclusion of correlation is expected to be similar.
## 9 Calculations of <sup>207</sup>PbO molecule.
As is noted in section 1, experiments on the excited $`a(1)`$ DeMille:00 or $`B(1)`$ Egorov:01 states of PbO having nonzero projection of total electronic momentum on the internuclear axis can be, in principle, sensitive enough to detect $`d_e`$ two or even four orders of magnitude lower than the current limit. Knowledge of the effective electric field, $`W`$, is required for extracting $`d_e`$ from the measurements (see section 3). In papers Isaev:04 ; Petrov:05a , $`W`$ for the $`a(1)`$ ($`{}_{}{}^{3}\mathrm{\Sigma }_{}^{+}`$ $`\sigma _1^2\sigma _2^2\sigma _3^2\pi _1^3\pi _2^1`$) and $`B(1)`$ ($`{}_{}{}^{3}\mathrm{\Pi }_{1}^{}`$ $`\sigma _1^2\sigma _2^2\sigma _3^1\pi _1^4\pi _2^1`$) states of the PbO molecule was calculated by the RCC-SD Kaldor:97 ; Landau:01c and configuration interaction (CI) Buenker:74 ; Buenker:99 ; Alekseyev:04 ; Titov:01 methods. To provide an accuracy check for the calculation of the electronic structure near the Pb nucleus the hyperfine constant, $`A_{}`$, was also calculated.
### Results.
CI calculations Petrov:05a were performed at two internuclear distances: $`R=3.8`$ a.u. (as in RCC calculations), and $`R=4.0`$ a.u. (which is close to the equilibrium distances of the $`a(1)`$ and $`B(1)`$ states). In the RCC calculations Isaev:04 the internuclear distance $`R=3.8`$ a.u. was used because of problems with convergence at larger distances. The calculated values with the one-center expansion of the molecular spinors in the Pb core in either $`s;p`$ or $`s;p;d`$ partial waves are collected in Table 5.
Let us consider the $`5s,5p,5d`$ orbitals of lead and $`1s`$ orbital of oxygen as the outercore and the $`\sigma _1`$, $`\sigma _2`$, $`\sigma _3`$, $`\pi _1`$, $`\pi _2`$ orbitals of PbO (consisting mainly of $`6s,6p`$ orbitals of Pb and $`2s,2p`$ orbitals of O) as valence. Although in the CI calculations we take into account only the correlation between valence electrons, the accuracy attained in the CI calculation of $`A_{}`$ is much better than in the RCC-SD calculation. The main problem with the RCC calculation was that the Fock-space RCC-SD version used there was not optimal in accounting for nondynamic correlations (see Isaev:00 for details of RCC-SD and CI calculations of the Pb atom). Nevertheless, the potential of the RCC approach for electronic structure calculations is very high, especially in the framework of the intermediate Hamiltonian formulation Landau:01c ; Kaldor:04b .
Next we estimate the contribution from correlations of valence electrons with outercore ones (which also account for correlations between outercore electrons) as the difference between the results of the corresponding 10- and 30-electron GRECP/RCC calculations (see also Isaev:00 where this correction is applied to the Pb atom). We designate such correlations in Table 5 as “outercore correlations”. When taking into account outercore contributions at the point $`R=4.0`$ a.u. we used the results of the RCC calculation at the point $`R=3.8`$ a.u. Since these contributions are relatively small and because the correlations with the outercore electrons are more stable than correlations in the valence region when the internuclear distance is changed, this approximation seems reasonable. Finally, we have linearly extrapolated the results of the calculations to the experimental equilibrium distances, $`R_e=4.06`$ a.u. for $`a(1)`$ Martin:88 and $`R_e=3.91`$ a.u. for $`B(1)`$ Huber:79 . The final results are: $`A_{}=3826`$ MHz, $`W=6.110^{24}\mathrm{Hz}/(e\mathrm{cm})`$ for $`a(1)`$ and $`A_{}=4887`$ MHz, $`W=8.010^{24}\mathrm{Hz}/(e\mathrm{cm})`$ for $`B(1)`$. The estimated error for the final $`W`$ value is 20% for the $`B(1)`$ state. For $`a(1)`$ the estimated error bounds put the actual $`W`$ value between 90% and 130% of our final value (for details see Petrov:05a ).
## 10 Conclusions
The P,T-parity nonconservation parameters and hyperfine constants have been calculated for the particular heavy-atom molecules which are of primary interest for modern experiments to search for PNC effects. It is found that a high level of accounting for electron correlations is necessary for reliable calculation of these properties with the required accuracy. The applied two-step (GRECP/NOCR) scheme of calculation of the properties described by the operators heavily concentrated in atomic cores and on nuclei has proved to be a very efficient way to take account of these correlations with moderate efforts. The results of the two-step calculations for hyperfine constants differ by less than 10% from the corresponding experimental data. A comparable level of accuracy is expected for the P,T-odd spin-rotational Hamiltonian parameters, which can not be obtained experimentally.
The precision of the discussed calculations is limited by the current possibilities of the correlation methods and codes and not by the GRECP and NOCR approximations, despite the fact that the GRECP/NOCR method allows one to reduce seriously the computational expenses of the correlation treatment as compared to conventional Dirac-Coulomb(-Breit) approaches when (1) using molecular spin-orbitals instead of spinors in (G)RECP calculation; (2) employing “correlated” GRECP versions Mosyagin:04a to account for correlations with the core electrons explicitly excluded from (G)RECP calculations; (3) combining gaussian basis functions at the molecular (G)RECP calculation with numerical functions at the one-center restoration; and (4) splitting the correlation treatment of a molecule into two sequential steps, first in valence and then in core regions.
In turn, the accuracy attained in the calculations presented above is sufficient for a reliable interpretation of the measurements in PNC experiments on these molecules.
### Acknowledgments.
We are very grateful to L.N.Labzowsky for initiating and supporting our activity in studying PNC properties in heavy-atom molecules over many years. We would like to thank our colleagues I.V.Abarenkov, A.B.Alekseyev, R.J.Buenker, E.Eliav, U.Kaldor, M.G.Kozlov, A.I.Panin, A.B.Tulub, and I.I.Tupitsyn for many stimulating discussions and fruitful collaboration on the relevant research.
The present work is supported by the U.S. CRDF grant RP2–2339–GA–02 and RFBR grant 03–03–32335. A.P. is grateful to Ministry of education of Russian Federation (grant PD 02–1.3–236) and to St.-Petersburg Committee of Science (grant PD 03-1.3-60). N.M. is also supported by the grants of Russian Science Support Foundation and the governor of Leningrad district. D.D. acknowledges additional support from NSF Grant No. PHY-0244927 and the David and Lucile Packard Foundation. A part of calculations of PbO was performed on computers of Boston University in the framework of the MARINER project. |
warning/0506/quant-ph0506192.html | ar5iv | text | # Quantum Scattering in Quasi-1D Cylindrical Confinement
## I Introduction
Matter at very small dimensions can be strongly affected by boundary and surface effects. Bulk 3D properties may then change substantially or even disappear. Reaching this regime poses a long term challenge motivated among others by the continuous technological drive towards ever smaller information processing devices itrs . Below a certain limit of small length scales, novel quantum electronic properties may appear and be useful in order to develop alternative devices itrs-erd .
Recent developments in atomic physics have addressed similar questions by using atom-optical devices (or atom-chips) made out of surface-patterned substrates weinstein1995a ; folman2002a ; reichel2002a ; hinds1999a ; fortagh2003a ; wang2005a or by means of laser manipulation of atoms for trapping grimm2000a and atomic lithography oberthaler2003a . Analogous to electronic matter waves, these devices and techniques have the potential to trap and manipulate atomic matter waves down to sub-micron scales. A vast range of phenomena with atoms can then be studied, several of them bearing close relationships to specific solid state systems.
Of particular interest for both bulk and atomic systems is the quasi-1D regime of quantum coherent *single mode* transport. In this ultimate limit, the degrees of freedom in the transversal trapping directions are effectively frozen in the quantum-limited ground state of the confining potential. Such one-dimensional system can show distinct physical properties e.g in quasi-1D quantum wires in two-dimensional electron gas (2DEG) systems itrs-erd ; datta1997a ; ferry1997a ; chu1989a ; bagwell1990a ; gurvitz1993a ; bardarson2004a ; indlekofer2005a and Tomonaga-Luttinger liquids tomonaga1950a ; luttinger1963a ; voit1994a . In elongated gases of ultracold atoms, the one dimensionality is revealed e.g. in magnetic guides weinstein1995a ; folman2002a ; reichel2002a ; hinds1999a ; fortagh2003a ; wang2005a , in phase fluctuations of quasi-condensates petrov2001a ; goerlitz2001a ; dettmer2001a , the Tonks-Girardeau gas of impenetrable bosons tonks1936 ; girardeau1960 ; lenard1966 ; tolra2004a ; paredes2004a ; kinoshita2004a or one dimensional optical lattices moritz2003a ; moritz2005a . In this broad context, we address two processes that are affected by the quasi-1D geometry. On one hand, the matter waves can be scattered off e.g. impurities or defects fixed along the guide. This *scattering by central fields* is an important variable in the transport properties and can be strongly affected by the confinement. This has been revealed by several studies of mesoscopic systems in the context of the Landauer approach to the quantized conductance of quantum point-contacts in the presence of impurities chu1989a ; bagwell1990a ; gurvitz1993a ; bardarson2004a . On the other hand, two-body *collision processes* should be dealt with when the matter wave consists of interacting particles. A full understanding of these low dimensional systems should therefore deal with both types of scattering. It should also derive *ad hoc* phenomenological parameters such as effective low dimensional coupling strengths from the real 3D character of the interactions. Another context in which such dimensional reduction plays a role is the transmission and reflection of electromagnetic bostroem1981a and elastic olsson1994a waves in transmission lines or resonators.
New light has recently been shed on this subject by studies of atomic collisions in 1D olshanii1998a ; moore2004a ; granger2004a ; bergeman2003 and 2D petrov2000b ; petrov2001b ; petrov2004a . Focusing on a parabolic cylindrical confinement, for which the center of mass can be eliminated, pioneering predictions for collision processes are possible, e.g. that of confinement induced *scattering resonances* (CIR) in 1D. In such a resonance olshanii1998a ; moore2004a ; granger2004a , a total reflection may take place between the colliding atoms. The $`s`$-wave (or zero-range) approximation to the scattering potential olshanii1998a ; moore2004a may be partially lifted, provided the coupling of orbital angular momenta due to the cylindrical geometry are not considered granger2004a . Another property induced by the confinement and closely related to CIR is the formation of weakly bound *quasi-molecules* with zero ($`s`$-wave) orbital angular momentum bergeman2003 , which have been recently observed with singlet fermionic atoms in tight laser traps moritz2005a . Coupling to the center of mass under non-parabolic confinement in the $`s`$-wave approximation is treated in peano2005a and can lead to the appearance of additional CIRs. In the independent and distinct context of impurity scattering studies, rectangular geometries predominate which are most suitable for the growth of semiconductor heterostructures. Similar properties of resonant scattering and weakly localized states are then predicted for non-interacting 2DEG systems. This is demonstrated in the $`s`$-wave approximation to the scattering potential in chu1989a . Finite range scatterers under general rectangular confinement are also considered, but in an already reduced two-dimensional space bagwell1990a ; gurvitz1993a ; bardarson2004a .
In this work, a detailed formalism is presented that extends the results for impurity scattering and collisions. It can treat general cylindrical confinement, parabolic or not, and incorporate all scattering phase-shifts beyond $`s`$-waves, as well as the full couplings of orbital angular momenta due to the broken spherical symmetry. A comprehensive assessment of the scattering process reveals the most important mechanisms and parameters at play. From this unified description, CIRs are seen to be a general low energy effect in quasi-1D geometries. It might be useful e.g. as an alternative gating mechanism in low power transistor-like devices that could incorporate quasi-1D structures itrs-erd ; datta1997a ; ferry1997a ; chu1989a ; bagwell1990a ; gurvitz1993a ; bardarson2004a ; indlekofer2005a ; effmass . The existence of weakly localized bound states and quasi-molecules can, in turn, be extended to include higher orbital angular momentum pairing, whose binding energies *and* wavefunctions can be systematically computed.
This paper is organized as follows. In section II the present problem of scattering in confined geometries is formulated for both scattering by central fields and collision processes. In section III, the $`s`$-wave approximation to the scattering potential is briefly reviewed. In section IV, after an overview of the free space scattering, our formalism in terms of Green’s functions and phase-shifts is developed. In section V, this formalism is applied to analyze two types of confinement induced effects, namely, CIR and weakly localized states. A discussion of the main approximations used is given in section VI, followed by the conclusions.
## II Scattering in Low Dimensionality
In a low energy scattering process under confinement, one expects the scattering potential to affect primarily the degrees of freedom not experiencing external confinement. The other degrees of freedom are then forced to occupy only one or very few confined states. However, the very existence of these confined states do play a role due to virtual transitions. The effects of the latter can become noticeable when properly examining the unconfined degrees of freedom. In this regard one is lead to calculate effective scattering amplitudes as given by Eq.(5b).
These amplitudes provide a complete description of the quasi-1D problem. For real energies, they describe the scattering properties of running waves. Weakly localized states are then found in the complex-energy plane.
### II.1 Hamiltonian under confinement
Consider first the scattering of a *single* particle of mass $`\mu `$ and coordinates $`𝒓=(x,y,z)`$ by a central field *fixed* at the origin. Let $`V(r)`$ be the spherically symmetric scattering potential and $`U(\rho )`$ the cylindrically symmetric confining potential, where $`𝝆=(x,y)`$ denotes the transverse coordinates. The Hamiltonian is then given by
$$H=\frac{\mathrm{}^2}{2\mu }^2+U(\rho )+V(r).$$
(1)
If $`R_V`$ and $`R_U`$ are the ranges of $`V(r)`$ and $`U(\rho )`$, respectively, then $`|V|V_0`$ for $`rR_V`$ and $`UU_0`$ for $`\rho R_U`$, where $`V_0`$ and $`U_0`$ are characteristic energy scales, such as the respective ground state energies (see fig. 1). Two particular examples are the parabolic confinement $`U(\rho )=\mu \omega _{}^2\rho ^2/2`$, for which $`R_U`$ is of the order of the oscillator length $`a_{}=(\mathrm{}/\mu \omega _{})^{1/2}`$ and the square-well (or hard-wall) such that $`U(\rho )+\mathrm{}`$ for $`\rho R_U`$ and zero otherwise. The first condition on the scattering and on the confining potentials assumed in this work is $`R_VR_U`$. In this sense, both the potentials $`V(r)`$ and $`U(\rho )`$ can be relatively *general* for this type of scattering process. This condition implies that there is a distinct region $`R_VrR_U`$ where spherical symmetry still prevails and one can define scattering phase-shifts. This is an (intermediate) asymptotic region, where the effect of the confinement $`U(\rho )`$ is felt only as a boundary condition to the scattering by $`V(r)`$ (see section IV.3). Fig. 1 depicts the geometry of the problem. Another important condition for the validity of the present approach to scattering in confined geometries is the condition $`k1/R_U`$ of low total scattering energy $`E\mathrm{}^2k^2/2\mu `$ (or Fermi energy) for the Hamiltonian $`H`$. Only few excited transverse states $`\phi _n`$ above the ground state (see Eq.(2)) can be effectively populated. The scattering process is thus assumed to occur at low temperatures or under sufficiently tight confinement.
Alternatively, consider the case of *collisions* between e.g. two cold atoms of masses $`m`$ and coordinates $`𝒓_i=(𝝆_i,z_i)`$, $`i=1,2`$. $`𝑹=(𝒓_1+𝒓_2)/2`$ and $`𝒓=𝒓_1𝒓_2`$ denote the center of mass (CM) and relative coordinates, respectively. The total Hamiltonian is
$$H_2=\underset{i=1}{\overset{2}{}}\left[\frac{\mathrm{}^2}{2m}_i^2+U_c(\rho _i)\right]+V(r),$$
where $`U_c`$ is the cylindrical confinement and $`V`$ the two-body interaction potential. For this scattering process $`U_c`$ should be *parabolic*, so that the CM can be separated exactly. The Hamiltonian for the relative coordinates is then given by Eq.(1), where $`\mu =m/2`$ becomes the (collisional) effective mass and $`U(\rho )=2U_c(\rho /2)`$. For non-parabolic confinement, additional scattering resonances that can occur even at weak coupling to the CM peano2005a is not included in Eq.(1). However, at ultracold temperatures for which only the lowest transversal states of the guide are occupied, the probability to find the atoms in the center of the guide is largest. Consequently, it might be a reasonable zeroth-order approximation to assume on average that the CM coordinates are also confined to the center of the guide. Therefore Eq.(1) is a starting point to the dynamics of the relative coordinates for a general interaction $`V(r)`$, in which the CM fluctuations around the mean value $`𝑹=0`$ are at the present stage neglected.
Because of the symmetry of $`U(\rho )+V(r)`$, the scattering solutions can be restricted to be axially symmetric. The corresponding axially symmetric and orthonormalized eigenstates of the confined variables $`\phi _n(\rho )`$, $`n=0,1,2,\mathrm{}`$, satisfy
$$\left[\frac{\mathrm{}^2}{2\mu }\left(\frac{^2}{x^2}+\frac{^2}{y^2}\right)+U(\rho )\right]\phi _n(\rho )=ϵ_n\phi _n(\rho )$$
(2)
where $`ϵ_n\mathrm{}^2q_n^2/2\mu >0`$ are the transversal energies, with $`q_01/R_U`$.
### II.2 1D scattering
In a purely 1D problem let the 1D interaction potential $`V_{1D}(z)`$ have the finite range $`R_{1D}`$ around the origin $`z=0`$, such that $`V_{1D}0`$ for $`|z|R_{1D}`$. The scattering problem for an incoming plane wave $`e^{ik_0z}`$ can be described by the amplitudes $`f^\pm `$ in the asymptotic behaviour of the solution $`\psi (z)`$ in the transmission and reflection regions
$$\psi (z)=\{\begin{array}{cc}e^{ik_0z}+f^+e^{ik_0z},\hfill & R_{1D}z,\hfill \\ e^{ik_0z}+f^{}e^{ik_0z},\hfill & zR_{1D}.\hfill \end{array}$$
Current conservation then yields
$$|1+f^+|^2=1|f^{}|^2.$$
(3)
Particular cases are resonant transmission ($`f^+0`$) and total reflection ($`f^+1`$) at finite $`k_0`$ and finite potentials.
### II.3 Effective 1D scattering
The above 1D picture corresponds to a fixed unperturbed transversal state, e.g., the ground state $`\phi _0(\rho )`$. The total wave function would then have the form $`\psi (z)\phi _0(\rho )`$ where scattering effects would occur only in $`\psi (z)`$. However, transitions to other states $`\phi _n(\rho )`$, $`n=1,2,\mathrm{}`$, even as virtual transitions, arise from the coupling between the $`\phi _n`$’s caused by the scattering potential $`V(r)`$.
Let $`\mathrm{\Psi }(𝒓)`$ be the full 3D axially symmetric scattering solution to Eq.(1),
$$\left[^2u(\rho )+k^2\right]\mathrm{\Psi }(𝒓)=v(r)\mathrm{\Psi }(𝒓),$$
(4a)
where $`E\mathrm{}^2k^2/2\mu >0`$ is the total scattering energy, $`u(\rho )2\mu U(\rho )/\mathrm{}^2`$ and $`v(r)2\mu V(r)/\mathrm{}^2`$. The above mentioned couplings and virtual transitions are best seen by expanding the solution in the cylindrical basis defined by the $`\phi _n`$’s,
$$\mathrm{\Psi }(𝒓)=\underset{n=0}{\overset{\mathrm{}}{}}\psi _n(z)\phi _n(\rho ),$$
(4b)
and solving for each $`\psi _n(z)`$. Substituting into Eq.(4a), multiplying by $`\phi _n^{}`$, integrating and using Eq.(2) yields
$$\left(\frac{d^2}{dz^2}+k^2q_n^2\right)\psi _n(z)=𝑑x𝑑y\phi _n^{}(\rho )v(r)\mathrm{\Psi }(𝒓).$$
For a given total energy $`E`$, one defines the longitudinal wave vectors $`k_n`$ by
$$k_n=\{\begin{array}{cc}\sqrt{k^2q_n^2},\hfill & \text{if }0nn_E,\hfill \\ i\sqrt{q_n^2k^2},\hfill & \text{if }n>n_E,\hfill \end{array}$$
where $`n_E`$ is the largest integer such that $`q_{n_E}<k<q_{1+n_E}`$ and denotes the highest “open channel”. In the single mode regime $`n_E=0`$. As usual, the general solution $`\psi _n(z)`$ can be expressed in terms of a homogeneous solution and a particular inhomogeneous solution, $`\psi _{f,n}(z)`$. This inhomogeneous solution can be written in terms of the 1D Green’s function
$$G_n(z)=\frac{e^{ik_n|z|}}{2ik_n},G_n^{\prime \prime }+k_n^2G_n=\delta (zz^{}),$$
describing an outward scattering for $`nn_E`$ and an exponentially decaying virtual state for the “closed channels” $`n>n_E`$. Then for all $`n`$
$$\psi _{f,n}(z)=d^3𝒓^{}G_n(zz^{})\phi _n^{}(\rho ^{})v(r^{})\mathrm{\Psi }(𝒓^{}).$$
The expansion Eq.(4b) of $`\mathrm{\Psi }(𝒓)`$ takes then the form
$$\mathrm{\Psi }(𝒓)=\underset{n=0}{\overset{n_E}{}}b_ne^{ik_nz}\phi _n(\rho )+\underset{n=0}{\overset{\mathrm{}}{}}\psi _{f,n}(z)\phi _n(\rho ),$$
(4c)
where the homogeneous part $`b_ne^{ik_nz}`$ of each $`\psi _n(z)`$, for some constant $`b_n`$, gives rise to the total incoming state and thus is limited to the open channels $`nn_E`$ only.
The asymptotic condition is obtained in the limit $`|z|\mathrm{}`$ by neglecting the exponentially decaying terms $`\psi _{f,n}(z)`$ in Eq.(4c) for $`n>n_E`$. Using then $`|zz^{}|=|z|z^{}`$ in $`G_n`$, for $`z\pm \mathrm{}`$, one has
$$\mathrm{\Psi }(𝒓)\underset{n=0}{\overset{n_E}{}}\left[b_ne^{ik_nz}+f_n^\pm e^{ik_n|z|}\right]\phi _n(\rho ),$$
(5a)
where the effective 1D scattering amplitudes $`f_n^\pm `$ are defined by (for $`nn_E`$)
$$f_n^\pm \frac{1}{2ik_n}d^3𝒓^{\mathbf{}}\left[e^{\pm ik_nz^{}}\phi _n(\rho ^{})\right]^{}v(r^{})\mathrm{\Psi }(𝒓^{}),$$
(5b)
for forward $`z>0`$ and backward $`z<0`$ scattering, respectively.
Note first that although $`|z||z^{}|`$, the phase $`k_nz^{}`$ in $`f_n^\pm `$ is not necessarily negligible since $`k_n|z^{}|2\pi `$ for sufficiently large momentum $`k_n`$ and depending on the range $`R_V`$ of the interaction $`v(r^{})`$. This phase is negligible only when the wavelength $`2\pi /k_n`$ of the incoming wave is not able to resolve the details of the potential $`v(r^{})`$ at low momenta. In this case, a zero-range or $`s`$-wave approximation to $`v(r^{})`$ is valid. However, this approximation should break down when the scattering occurs at short wavelengths or when higher angular momenta become necessary. Note that there is no clearly defined 1D range $`R_{1D}`$ as in the purely 1D case, since the asymptotics given by Eq.(5a) depends on the convergence behaviour of the series in Eq.(4c). If the convergence is slow, one expects an effective 1D range, say $`R_{1D}^{}`$, substantially larger than $`R_V`$ or even $`R_U`$. It is an effective range for the 1D collision and is of the order of the minimum separation from the scattering center, such that Eq.(5a) is the dominant contribution to the scattering wave function (see fig. 2). As regards collision processes, this effective range sets some limits on the average equilibrium distance between the colliding particles.
The present formalism to quantum scattering under confinement is approximate but non-perturbative, in the sense that there is no small parameter around which an expansion is performed. The determination of its range of validity is chosen here by imposing the conservation of probability on the final wavefunctions. The probability conservation condition in 3D and the identity $`_\mathrm{\Omega }d^3𝒓\text{div}𝒋\mathbf{(}𝒓\mathbf{)}=_{S_c+S_1+S_2}d^2𝑺𝒋\mathbf{(}𝒓\mathbf{)}`$ for a volume $`\mathrm{\Omega }`$ enclosed by a large cylinder of surface area $`S_c`$ and by two transversal discs of surface areas $`S_1`$ at $`z_1<0`$ and $`S_2`$ at $`z_2>0`$ yields
$`0`$ $`=`$ $`{\displaystyle _{S_c}}𝑑\varphi 𝑑z\rho \left[\widehat{𝝆}𝒋(𝒓)\right]`$
$`{\displaystyle _{S_1}}𝑑x𝑑y\left[\widehat{𝒌}𝒋(𝒓)\right]+{\displaystyle _{S_2}}𝑑x𝑑y\left[\widehat{𝒌}𝒋(𝒓)\right],`$
where $`\widehat{𝝆}`$ and $`\widehat{𝒌}`$ are cylindrical unit vectors. Applying $`\mathrm{\Psi }=_n\psi _n\phi _n`$, the first integral on the rhs vanishes for large radius $`\rho \mathrm{}`$ since each $`\phi _n(\rho )0`$. The other two integrals yield
$$0=\underset{n=0}{\overset{\mathrm{}}{}}\left[j_n(z_1)+j_n(z_2)\right],$$
where $`j_n(z)`$ are 1D currents calculated with $`\psi _n(z)`$. Setting $`z_{1,2}R_{1D}^{}`$ one can use the asymptotics for $`\psi _n`$ in Eq.(5a) so that the total conservation condition reads
$$0=\underset{n=0}{\overset{n_E}{}}\left(|b_n+f_n^+|^2+|f_n^{}|^2|b_n|^2\right)k_n.$$
(5c)
This equation extends the pure 1D result in Eq.(3). It also determines how the initial flux is distributed among the open channels. In other words, even when the initial state occupies only a single channel, i.e., $`b_n=\delta _{0,n}`$, new channels are occupied due to scattering, if the total energy is large enough $`n_E>0`$. This conservation condition can serve then as a gauge for the range of validity of our approximations to the amplitudes $`f_n^\pm `$.
## III Scattering in the $`s`$-wave approximation
In many circumstances, the low energy scattering in free space can be well described by approximating the interaction $`V(r)`$ by Fermi-Huang’s zero-range potential huang1987 ; fetter1971a
$$V_\delta (r)=\frac{2\pi \mathrm{}^2a}{\mu }\delta (𝒓)\frac{}{r}(r).$$
(6a)
This approximation singles out the $`l=0`$ angular momentum component of the collision and thus depends only on the $`s`$-wave scattering length $`a`$. The regularization $`(r)/r`$ assures that the singularity $`1/r`$ of the full scattering solution to this $`V_\delta (r)`$ is properly dealt with.
For collisions in a strong parabolic confinement and using Eq.(6a), the series in Eq.(4c) also has a singularity for short distances $`r0`$ very similar to $`1/r`$. By properly dealing with this singularity, with olshanii1998a or without any regularization moore2004a , the amplitudes $`f_n^\pm `$ in Eq.(5b) can be calculated. In the single mode regime ($`n_E=0`$) and for small longitudinal momenta $`k_0`$ one obtains
$$f_0^\pm =\frac{1}{1+i\left[\frac{a_{}^2}{2a}\left(1\frac{Ca}{a_{}}\right)\right]k_0},$$
(6b)
where $`C=1.4603\mathrm{}`$ and $`a_{}`$ is the harmonic oscillator length of the parabolic confinement olshanii1998a ; moore2004a . Here, a total reflection $`f_0^+=1`$ is predicted for large scattering lengths, on the order of the confinement length, such that $`a_{}=Ca`$, which is called the confinement induced resonance (CIR) olshanii1998a . This total reflection means that the colliding pair experiences a diverging effective 1D interaction along the longitudinal cylindrical axis (see section V.4.1).
It will be seen later that such a singularitiy should actually be the *same* singularity $`1/r`$. It arises from the 3D free space Green’s function and becomes dominant as $`r0`$ away from the confining boundaries (see sections IV.2 and VI). This latter identification is one of the elements that will allow us to incorporate all partial waves $`l1`$ besides the $`s`$-wave. In addition, the above conclusions concerning CIR under parabolic confinement will be seen to hold not only for general interactions $`V(r)`$, as first examined in granger2004a , but also for general *non-parabolic* confinement $`U(\rho )`$. In this sense the present formalism is of a more general character.
## IV Scattering Phase-Shifts
In the following, a general formalism is developed in order to calculate the scattering amplitudes for any given finite short-ranged and spherically symmetric scattering potential $`V(r)`$. This is done by expressing the amplitudes $`f_n^\pm `$ in terms of the scattering phase-shifts $`\delta _l`$ associated to this $`V(r)`$. Our formalism is also able to deal with a broad range of confinements $`U(\rho )`$ as well as scattering energies above the transversal ground state. Particularly relevant is the straightforward accounting of angular momenta couplings due to the confinement. We begin by providing the connection to scattering in free space.
### IV.1 Phase shifts in free space scattering
The standard formalism for free space scattering by a central potential is recast in a form suitable for later comparison. The main idea is to express all quantities relevant for scattering in terms of the phase-shifts. The rationale is that these phase-shifts can be considered as intrinsic to a given scattering potential $`V(r)`$, without much regard to details of external boundary conditions.
In the absence of any confinement, the scattering solution can be written as
$$\mathrm{\Phi }(𝒓)=\mathrm{\Phi }_i(𝒓)d^3𝒓^{}G(𝒓,𝒓^{})v(r^{})\mathrm{\Phi }(𝒓^{\mathbf{}}),$$
(7a)
where $`\mathrm{\Phi }_i(𝒓)`$ is the incoming state and $`G(𝒓,𝒓^{})`$ is the free space Green’s function. In terms of outward and inward scattered waves, the latter reads
$$G(𝒓,𝒓^{})=\gamma _+\frac{e^{ik|𝒓𝒓^{}|}}{4\pi |𝒓𝒓^{}|}+\gamma _{}\frac{e^{ik|𝒓𝒓^{}|}}{4\pi |𝒓𝒓^{}|},$$
(7b)
for some constants $`\gamma _\pm `$ obeying $`\gamma _++\gamma _{}=1`$. A single constant $`\gamma `$ can be introduced by setting
$$\gamma _\pm \frac{1\pm \gamma }{2}.$$
(7c)
In free space scattering, the inward component proportional to $`\gamma _{}`$ is usually absent (i.e., $`\gamma =1`$). However, one expects this *inflow* of particles if there is e.g. an exterior confinement that forces the scattered particles back towards the center. This interpretation will be validated in section IV.2 (see also Fig 2). $`\mathrm{\Phi }_i(𝒓)`$ and $`G(𝒓,𝒓^{})`$ are determined by the boundary conditions imposed on $`\mathrm{\Phi }(𝒓)`$, e.g. that of an axially symmetric solution.
The scattering variables of interest here are the amplitudes $`f_{\mathrm{out},\mathrm{in}}`$ of outward and inward waves, scattered along a direction given by $`\theta `$ and $`\varphi `$ at large distances $`rR_V`$ from the scattering center. Using the expansion $`|𝒓𝒓^{}|r𝒓𝒓^{}/r`$ in Eq.(7b), Eq.(7a) becomes
$$\mathrm{\Phi }(𝒓)\mathrm{\Phi }_i(𝒓)+f_{\mathrm{out}}\frac{e^{ikr}}{r}+f_{\mathrm{in}}\frac{e^{ikr}}{r},rR_V,$$
(8a)
with
$$f_{\mathrm{out},\mathrm{in}}=\frac{\gamma _\pm }{4\pi }d^3𝒓^{}\left[e^{\pm i𝒌𝒓^{}}\right]^{}v(r^{})\mathrm{\Phi }(𝒓^{}),$$
(8b)
where $`𝒌k𝒓/r`$ depends on the direction $`(\theta ,\varphi )`$ of $`𝒓`$. It can be clearly seen that these amplitudes have a shape similar to their equivalent effective 1D amplitudes $`f_n^\pm `$ in Eq.(5b). To express the amplitudes in terms of the phase-shifts an angular momentum decomposition of $`f_{\mathrm{out},\mathrm{in}}`$ is necessary. Following standard procedures morse1953 ; mott1965 , we obtain from Eq.(8b)
$$f_{\mathrm{out},\mathrm{in}}=\gamma _\pm \underset{l=0}{\overset{\mathrm{}}{}}(i)^l(2l+1)\tau _lP_l(\mathrm{cos}\theta ),$$
(9a)
with the definition
$$\tau _l\frac{1}{4\pi }d^3𝒓^{}\left[j_l(kr^{})P_l(\mathrm{cos}\theta ^{})\right]v(r^{})\mathrm{\Phi }(𝒓^{}),$$
(9b)
where $`j_l`$ and $`P_l`$ are the spherical Bessel functions and Legendre polynomials, respectively. The next step is to relate the elements $`\tau _l`$ to the phase-shifts. Decompose $`\mathrm{\Phi }_i(𝒓)`$ and $`G(𝒓,𝒓^{})`$ in spherical coordinates. Substituting then into Eq.(7a) for $`r>r^{}R_V`$ and carrying out the $`\varphi ^{}`$-integration one obtains for $`r>R_V`$
$`\mathrm{\Phi }(𝒓)`$ $`=`$ $`{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}(2l+1)(\beta _lik\gamma \tau _l)j_l(kr)P_l(\mathrm{cos}\theta )`$
$`+{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}(2l+1)(\beta _l^{}+k\tau _l)n_l(kr)P_l(\mathrm{cos}\theta ),`$
having used Eq.(7c) and Eq.(9b) and the constants $`\beta `$ and $`\beta ^{}`$ are related to $`\mathrm{\Phi }_i(𝒓)`$. On the other hand, the phase-shifts arise when directly decomposing this same solution $`\mathrm{\Phi }(𝒓)`$ in spherical coordinates and looking at the radial solutions outside the interaction region $`r>R_V`$,
$$\mathrm{\Phi }(𝒓)=\underset{l=0}{\overset{\mathrm{}}{}}c_l\left[\mathrm{cos}\delta _lj_l(kr)\mathrm{sin}\delta _ln_l(kr)\right]P_l(\mathrm{cos}\theta ).$$
Here the phase-shifts $`\delta _l`$ are calculated from the behaviour of the radial parts $`R_l(r)`$ inside the scattering region $`r<R_V`$ only, irrespective of the details of the outer boundary conditions mott1965 . The boundary conditions are matched by properly choosing the constants $`c_l`$, i.e. by superposing the spherical components. Comparing these last two forms of $`\mathrm{\Phi }(𝒓)`$ yields,
$`c_l`$ $`=`$ $`(2l+1){\displaystyle \frac{\beta _l+i\gamma \beta _l^{}}{\mathrm{cos}\delta _li\gamma \mathrm{sin}\delta _l}},`$ (10a)
$`\tau _l`$ $`=`$ $`{\displaystyle \frac{1}{k}}{\displaystyle \frac{\beta _l\mathrm{sin}\delta _l+\beta _l^{}\mathrm{cos}\delta _l}{\mathrm{cos}\delta _li\gamma \mathrm{sin}\delta _l}},l=0,1,\mathrm{}`$ (10b)
Considering the phase-shifts $`\delta _l`$ (or the scattering potential $`v(r)`$) as given, the second relation Eq.(10b) solves for the amplitudes $`f_{\mathrm{out},\mathrm{in}}`$ in Eq.(9a) (and also the integral equation Eq.(7a) at distances $`r>R_V`$). The first relation Eq.(10a) in turn determines how to superpose the spherical components in $`\mathrm{\Phi }=_lc_lR_lP_l`$ in order to match the boundary condition contained in Eq.(8a).
### IV.2 Green’s functions under confinement
Just as the free space 3D scattering amplitudes in Eqs.(8) can be written in terms of the phase-shifts via Eqs.(9) and Eq.(10b), a very similar procedure can be applied to calculate the amplitudes $`f_n^\pm `$ characterizing the 1D effective dynamics. In this section a key element of the formalism is worked out, namely, the Green’s function $`G_c`$ under confinement. It is shown that $`G_c`$ can be expressed in terms of the free space Green’s function $`G`$ as in Eq.(12). The final purpose, to be achieved in section IV.3, is to express it as a series in spherical coordinates, valid at least in the center of the guide where the spherical symmetry prevails over the cylindrical one due to the main assumption of a short range scattering potential $`v(r)`$ and a confining potential $`u(\rho )`$ that is sufficiently flat at the center, $`R_VR_U`$. Such a series allows then the straightforward introduction of the phase-shifts.
Consider the solution $`\mathrm{\Psi }`$ to Eq.(4a) given by the expansion in Eq.(4c). One can rewrite Eq.(4c) in the integral form (keeping again the $`\varphi ^{}`$-integration as in Eq.(9b) for convenience)
$$\mathrm{\Psi }(𝒓)=\mathrm{\Psi }_i(𝒓)d^3𝒓^{}G_c(𝒓,𝒓^{})v(r^{})\mathrm{\Psi }(𝒓^{}),$$
(11a)
where the incident state is
$$\mathrm{\Psi }_i(𝒓)=\underset{n=0}{\overset{n_E}{}}b_ne^{ik_nz}\phi _n(\rho ),$$
(11b)
and $`G_c`$, the axially symmetric Green’s function of the confining potential, has the form
$$G_c(𝒓,𝒓^{})=\underset{n=0}{\overset{\mathrm{}}{}}\phi _n(\rho )\phi _n^{}(\rho ^{})G_n(zz^{}).$$
(11c)
It is convenient to introduce a non-axially symmetric Green’s function $`G_u`$ under the confining potential $`u(\rho )`$. This $`G_u`$ then satisfies $`\left[^2u(\rho )+k^2\right]G_u(𝒓,𝒓^{})=\delta (𝒓𝒓^{})`$. If the $`\varphi `$-independent solution $`\mathrm{\Psi }(𝒓)`$ is written as an integral equation using this $`\varphi `$-dependent $`G_u`$, it follows from comparison with the integral equation Eq.(11a) that $`G_c`$ and $`G_u`$ are related by $`2\pi G_c(𝒓,𝒓^{})=𝑑\varphi ^{}G_u(𝒓,𝒓^{})`$. Compare now the differential equations of $`G_u`$ and that of the free space Green’s function $`G`$ of Eq.(7b) in the region of the scattering potential such that $`r,r^{}R_U`$. Since $`u(\rho )0`$ by assumption, they should differ at most by a homogeneous term, say $`\mathrm{\Delta }_u`$, so that $`G_u(𝒓,𝒓^{})G(𝒓,𝒓^{})+\mathrm{\Delta }_u(𝒓,𝒓^{})`$, for $`r,r^{}R_U`$, where $`\mathrm{\Delta }_u`$ satisfies the homogeneous equation $`(^2+k^2)\mathrm{\Delta }_u(𝒓,𝒓^{})=0`$ and we obtain
$$G_c(𝒓,𝒓^{})\frac{d\varphi ^{}}{2\pi }G(𝒓,𝒓^{})+\mathrm{\Delta }_c(𝒓,𝒓^{}),r,r^{}R_U,$$
(12)
where $`\mathrm{\Delta }_c𝑑\varphi ^{}\mathrm{\Delta }_u/2\pi `$. In order to obtain the function $`\mathrm{\Delta }_c`$ and $`\gamma `$ in Eq.(7c) in terms of the scattering parameters, we first write the closed channel ($`n>n_E`$) wavefunctions of $`G_c`$ in Eq.(11c) and $`G`$ in Eq.(7b) using cylindrical Bessel functions. The latter represents a good approximation to these transversal eigenstates for a given potential $`U(\rho )`$ close to the guide center $`\rho R_U`$ where $`U(\rho )0`$. By subsequently comparing $`G_c`$ and $`G`$ we can determine $`\mathrm{\Delta }_c`$.
First we express the transversal eigenstates $`\phi _n(\rho )`$ in terms of the normalized Bessel function according to
$$\phi _n(\rho )\frac{N_n}{\pi ^{1/2}R_U}J_0(q_n\rho ),\rho R_U,n>n_E,$$
(13a)
where $`N_n=|J_1(r_{n+1})|^1`$ and $`r_{n+1}`$ is the $`(n+1)`$-th root of $`J_0`$, related to the transverse momenta $`q_n`$ by $`q_nR_U=r_{n+1}`$. A good expression for the roots $`r_{n+1}`$ is given by the cosine approximation $`J_0(x)(2/\pi x)^{1/2}\mathrm{cos}(x\pi /4)`$. This leads to
$$q_n\left(n+\frac{3}{4}\right)\frac{\pi }{R_U},n>n_E.$$
(13b)
The corresponding energy dispersion relation $`ϵ_nn^2`$ is quadratic with respect to $`n`$ compared to the linear one $`ϵ_n2n`$ for, e.g., 2D harmonic oscillators (see e.g. olshanii1998a ). This should not cause any confusion since this excited spectrum $`n>n_E`$ is summed over in Eq.(11c) and leads to the confinement independent free space Green’s function $`G`$ in Eq.(12). It is the lower part of the spectrum $`nn_E`$, in turn, that relates to the term $`\mathrm{\Delta }_c`$.
Indeed, we substitute Eq.(13a) in the series for the closed channels in Eq.(11c) yielding for $`r,r^{}R_U`$
$`G_c(𝒓,𝒓^{}){\displaystyle \underset{n=0}{\overset{n_E}{}}}\phi _n(\rho )\phi _n^{}(\rho ^{}){\displaystyle \frac{e^{i\sqrt{k^2q_n^2}|zz^{}|}}{2(i)\sqrt{k^2q_n^2}}}`$
$`+`$ $`{\displaystyle \underset{n=1+n_E}{\overset{\mathrm{}}{}}}{\displaystyle \frac{q_n\delta q_n}{4\pi }}J_0(q_n\rho )J_0(q_n\rho ^{}){\displaystyle \frac{e^{\sqrt{q_n^2k^2}|zz^{}|}}{\sqrt{q_n^2k^2}}},`$
where $`\delta q_n=\pi /R_U`$ is the increment of $`q_n`$ and we used $`|J_1(r_{n+1})|^2=2/\pi q_nR_U`$. For $`r,r^{}R_U`$ a continuum approximation to the closed channel series can now be applied. Separating the real and imaginary parts of $`G_c`$ (supposing $`\phi _n(\rho )`$ real for $`nn_E`$ without much loss of generality) gives
$`G_c(𝒓,𝒓^{})`$ (14)
$``$ $`i{\displaystyle \underset{n=0}{\overset{n_E}{}}}\phi _n(\rho )\phi _n(\rho ^{}){\displaystyle \frac{\mathrm{cos}\left(\sqrt{k^2q_n^2}|zz^{}|\right)}{2\sqrt{k^2q_n^2}}}`$
$`{\displaystyle \underset{n=0}{\overset{n_E}{}}}\phi _n(\rho )\phi _n(\rho ^{}){\displaystyle \frac{\mathrm{sin}\left(\sqrt{k^2q_n^2}|zz^{}|\right)}{2\sqrt{k^2q_n^2}}}`$
$`+{\displaystyle _{q_{1+n_E}}^{\mathrm{}}}{\displaystyle \frac{qdq}{4\pi }}J_0(q\rho )J_0(q\rho ^{}){\displaystyle \frac{e^{\sqrt{q^2k^2}|zz^{}|}}{\sqrt{q^2k^2}}},`$
where $`q_{1+n_E}>k`$ designates the transversal momentum of the first “virtual” or closed channel $`n=1+n_E`$. As for $`G`$, one uses now the identity morse1953
$`{\displaystyle \frac{e^{ik|𝒓𝒓^{}|}}{4\pi |𝒓𝒓^{}|}}`$ $`=`$ $`{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}\left(2\delta _{0,m}\right)\mathrm{cos}[m(\varphi \varphi ^{})]`$
$`\times {\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{qdq}{4\pi }}J_m(q\rho )J_m(q\rho ^{}){\displaystyle \frac{e^{i\sqrt{k^2q^2}|zz^{}|}}{i\sqrt{k^2q^2}}},`$
and the correct branch $`0\mathrm{Arg}\sqrt{k^2q^2}<\pi `$ in order to obtain for the integrated free space Green’s function
$`{\displaystyle \frac{d\varphi ^{}}{2\pi }G(𝒓,𝒓^{})}`$ (15)
$`=`$ $`i\gamma {\displaystyle _0^k}{\displaystyle \frac{qdq}{4\pi }}J_0(q\rho )J_0(q\rho ^{}){\displaystyle \frac{\mathrm{cos}\left(\sqrt{k^2q^2}|zz^{}|\right)}{\sqrt{k^2q^2}}}`$
$`{\displaystyle _0^k}{\displaystyle \frac{qdq}{4\pi }}J_0(q\rho )J_0(q\rho ^{}){\displaystyle \frac{\mathrm{sin}\left(\sqrt{k^2q^2}|zz^{}|\right)}{\sqrt{k^2q^2}}}`$
$`+{\displaystyle _k^{\mathrm{}}}{\displaystyle \frac{qdq}{4\pi }}J_0(q\rho )J_0(q\rho ^{}){\displaystyle \frac{e^{\sqrt{q^2k^2}|zz^{}|}}{\sqrt{q^2k^2}}}.`$
To compare with the Green’s function $`G_c`$ in Eq.(14), we specialize Eq.(15) to $`r,r^{}R_U`$ and impose the condition for low total scattering energy.
The value of $`\gamma `$ follows from the comparison of the imaginary parts of Eqs.(14) and (15). Changing the integration variable from $`q`$ to $`\sqrt{k^2q^2}`$ in the first integration on the rhs of Eq.(15) and using the low energy condition in the imaginary parts of Eq.(14) and Eq.(15), $`\gamma `$ can be identified in the limit $`r,r^{}0`$
$$\gamma =2\pi \underset{n=0}{\overset{n_E}{}}\frac{|\phi _n(0)|^2}{k\sqrt{k^2q_n^2}}.$$
In other words, since $`r,r^{}R_U1/k`$, both imaginary parts are approximated by a single and the same constant $`\gamma k/4\pi `$, independent of $`𝒓`$ and $`𝒓^{\mathbf{}}`$. Although exact values of $`|\phi _n(0)|^2`$ can be employed for a given potential $`U(\rho )`$, it is more instructive to use general approximations such as $`|\phi _n(0)|^21/\pi R_U^2`$ obtained through normalization. An estimate that takes into account the spatial variation and the few nodes of $`\phi _n`$ is the square well approximation, namely, $`|\phi _n(0)|^2N_n^2/\pi R_U^2`$. Therefore, without loss of generality, one can thus write
$$\gamma =\underset{n=0}{\overset{n_E}{}}\frac{2N_n^2}{R_U^2k\sqrt{k^2q_n^2}}.$$
(16a)
It should be mentioned that one can improve on how these low energy poles $`q_n`$ of $`G_c`$ are treated. However, it suffices to use for now the above result in order to present the formalism.
By comparing the real parts of Eq.(14) and Eq.(15), it is seen that $`\mathrm{\Delta }_c`$ in Eq.(12) is given by
$`\mathrm{\Delta }_c(𝒓,𝒓^{})`$
$`=`$ $`{\displaystyle _k^{q_{1+n_E}}}{\displaystyle \frac{qdq}{4\pi }}J_0(q\rho )J_0(q\rho ^{}){\displaystyle \frac{e^{\sqrt{q^2k^2}|zz^{}|}}{\sqrt{q^2k^2}}}`$
$`+[{\displaystyle _0^k}{\displaystyle \frac{d\overline{p}}{4\pi }}J_0(\overline{q}\rho )J_0(\overline{q}\rho ^{})\mathrm{sin}(\overline{p}|zz^{}|)`$
$`{\displaystyle \underset{n=0}{\overset{n_E}{}}}\phi _n(\rho )\phi _n(\rho ^{}){\displaystyle \frac{\mathrm{sin}\left(k_n|zz^{}|\right)}{2k_n}}],`$
where $`k_n=\sqrt{k^2q_n^2}`$ and $`\overline{q}\sqrt{k^2\overline{p}^2}`$. The first term on the rhs stems from the offset between the lower limits of integration $`q_{1+n_E}`$ and $`k`$ in Eq.(14) and Eq.(15), respectively. It accounts for the discreteness of the low lying transversal states due to the confinement. The second term in the square brackets is of order $`(k^2/8\pi _{n=0}^{n_E}N_n^2/2\pi R_U^2)|zz^{}|`$ and thus is relatively smaller by a factor of $`|zz^{}|/R_U`$ compared to the first term, whose order of magnitude $`(q_{1+n_E}^2k^2)^{1/2}/4\pi 1/R_U`$ can be estimated by calculating the integral for $`𝒓,𝒓^{}0`$. Neglecting this second term yields
$$\mathrm{\Delta }_c(𝒓,𝒓^{})=_0^{p_c}\frac{dp}{4\pi }J_0(q\rho )J_0(q\rho ^{})e^{p|zz^{}|},$$
(16b)
where $`q(k^2+p^2)^{1/2}`$ and $`p_c(q_{1+n_E}^2k^2)^{1/2}`$. The relation Eq.(12) is then proved and will be the basis of an angular momenta decomposition of the effective amplitudes $`f_n^\pm `$ in section IV.3. In this regard, the weak dependence of $`\mathrm{\Delta }_c`$ on the coordinates $`𝒓`$ and $`𝒓^{}`$ is kept, since this dependence introduces couplings between orbital angular momenta. This will be seen as $`\mathrm{\Delta }_c`$ is decomposed in spherical coordinates in the following.
### IV.3 Angular momenta decomposition
We are now in the position to express the effective 1D scattering amplitudes $`f_n^\pm `$ in terms of the scattering phase-shifts $`\delta _l`$ of a given scattering potential. The first step is to decompose these amplitudes in Eq.(5b) in a similar way as was done in the free space case in Eqs.(9) by introducing “matrix” elements similar to $`\tau _l`$. Subsequently these quantities are related to the phase-shifts $`\delta _l`$ in a manner analogous to Eqs.(10).
Applying the condition $`R_VR_U`$, the transversal states $`\phi _n`$ in the integrand of the amplitudes Eq.(5b) should be well approximated by the Bessel functions Eq.(13a). This is because $`r^{}R_V`$ whereas $`|k_n|q_n1/R_U`$. It is enough then to find expansions of $`e^{\pm ik_nz}J_0(q_n\rho )`$ in spherical coordinates, with $`k^2=q_n^2+k_n^2`$. For this purpose one may invert the following identity morse1953
$`P_l(\mathrm{cos}\overline{\theta })j_l(kr)`$ $`=`$ $`{\displaystyle \frac{1}{2i^l}}{\displaystyle _0^\pi }𝑑w\mathrm{sin}we^{ikr\mathrm{cos}\overline{\theta }\mathrm{cos}w}`$
$`\times J_0(kr\mathrm{sin}\overline{\theta }\mathrm{sin}w)P_l(\mathrm{cos}w)`$
by using the orthonormality property of the Legendre polynomials. As a result, one obtains for $`rR_U`$
$$e^{ik_nz}\phi _n(\rho )=\underset{l=0}{\overset{\mathrm{}}{}}i^l(2l+1)\alpha _{ln}j_l(kr)P_l(\mathrm{cos}\theta ),$$
(17)
where $`\alpha _{ln}N_nP_l(k_n/k)/\pi ^{1/2}R_U`$. For more quantitative results, these coefficients $`\alpha _{ln}`$ can be evaluated numerically from the exact eigenfunctions $`\phi _n(\rho )`$. One can then express the amplitudes $`f_n^\pm `$ in Eq.(5b) as the following series in angular momenta for $`nn_E`$
$`f_n^\pm `$ $`=`$ $`f_{ng}\pm f_{nu}`$
$``$ $`{\displaystyle \underset{l\mathrm{even}}{}}{\displaystyle \frac{4\pi (2l+1)\alpha _{ln}}{2ik_n}}T_l\pm {\displaystyle \underset{l\mathrm{odd}}{}}{\displaystyle \frac{4\pi (2l+1)\alpha _{ln}}{2ik_n}}T_l`$
having used the parity property $`P_l(x)=()^lP_l(x)`$. The “matrix” element $`T_l`$ is defined by
$$T_l\frac{1}{i^l\mathrm{\hspace{0.17em}4}\pi }d^3𝒓^{}[j_l(kr^{})P_l(\mathrm{cos}\theta ^{})]v(r^{})\mathrm{\Psi }(𝒓^{}).$$
(18b)
The non-zero momenta $`l1`$ contributions to $`f_n^\pm `$ stem from the (small) dependence on $`z^{}`$ and $`\rho ^{}`$ of the integrand in the definition of $`f_n^\pm `$. Their neglect amounts to assuming a point like zero range interaction $`v(r)`$, for which only the $`s`$-wave remains. In the free space scattering case, the corresponding quantity equivalent to this $`T_l`$ is $`\tau _l`$ given in Eq.(9b). The key difference now is that each $`T_l`$ will depend on all other $`T_l`$’s, meaning that the confined scattering solution $`\mathrm{\Psi }(𝒓)`$ in the definition of $`T_l`$ couples the angular momenta, while the free space solution $`\mathrm{\Phi }(𝒓)`$ in the definition of $`\tau _l`$ does not.
The next step is to relate the $`T_l`$ to the phase-shifts $`\delta _l`$. From the results in Eqs.(16) for the approximation Eq.(12), the Green’s function $`G_c`$ can be decomposed into spherical coordinates in the region $`r^{},rR_U`$ by separately decomposing $`G`$ and $`\mathrm{\Delta }_c`$. The free space part $`G`$ can be decomposed following the procedures of section IV.1. Then, for $`r^{}<r`$,
$`{\displaystyle \frac{d\varphi ^{}}{2\pi }G(𝒓,𝒓^{})}`$ $`=`$ $`ik{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}j_l(kr^{})\left[\gamma j_l(kr)+in_l(kr)\right]`$ (19)
$`\times {\displaystyle \frac{2l+1}{4\pi }}P_l(\mathrm{cos}\theta )P_l(\mathrm{cos}\theta ^{}),`$
having used the expansion of $`e^{\pm ik|𝒓𝒓^{}|}/4\pi |𝒓𝒓^{}|`$ in spherical harmonics and $`\gamma _\pm =(1\pm \gamma )/2`$. To evaluate $`\mathrm{\Delta }_c`$ we use the analytic continuation $`\overline{\theta }=\pi /2i\theta ^{}`$ of the first identity of this section to obtain
$`e^{pz}J_0(q\rho )`$ $`=`$ $`{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}\left[i^l(2l+1)P_l(ip/k)\right]j_l(kr)P_l(\mathrm{cos}\theta ),`$
where $`p=k\mathrm{sinh}\theta ^{}`$ and $`q=k\mathrm{cosh}\theta ^{}`$ such that $`k^2=q^2p^2`$. The expansion of $`\mathrm{\Delta }_c`$ in Eq.(16b) for $`r,r^{}R_U`$ can then be written as
$`\mathrm{\Delta }_c(𝒓,𝒓^{})`$ $`=`$ $`{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}\left[{\displaystyle _0^{p_c}}{\displaystyle \frac{dp}{4\pi }}P_l(\sigma _{zz^{}}ip/k)e^{\sigma _{zz^{}}pz^{}}J_0(q\rho ^{})\right]`$ (20)
$`\times i^l(2l+1)j_l(kr)P_l(\mathrm{cos}\theta ),`$
where $`\sigma _{zz^{}}\mathrm{sign}(zz^{})`$. Inserting Eq.(20) and Eq.(19) back into Eq.(12), the scattering solution $`\mathrm{\Psi }(𝒓)`$ in Eq.(11a) becomes for $`R_VrR_U`$
$`\mathrm{\Psi }(𝒓)`$ $``$ $`{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}i^l(2l+1)\left[\alpha _l+\gamma _l(z)i\gamma kT_l\right]j_l(kr)P_l(\mathrm{cos}\theta )`$ (21)
$`+{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}i^l(2l+1)\left[kT_l\right]n_l(kr)P_l(\mathrm{cos}\theta ),`$
where $`\alpha _l_{n=0}^{n_E}b_n\alpha _{ln}`$ comes from the expansion of $`e^{ik_nz}\phi _n(\rho )`$ in the initial state $`\mathrm{\Psi }_i(𝒓)`$ in Eq.(11b) and $`\gamma _l(z)`$ is defined by
$`\gamma _l(z)`$ $``$ $`{\displaystyle _0^{p_c}}{\displaystyle \frac{dp}{4\pi }}{\displaystyle d^3𝒓^{}P_l(\sigma _{zz^{}}ip/k)}`$
$`\times \left[e^{\sigma _{zz^{}}pz^{}}J_0(q\rho ^{})\right]v(r^{})\mathrm{\Psi }(𝒓^{}).`$
The $`z`$-dependence of this $`\gamma _l(z)`$ arises from the sign of $`\sigma _{zz^{}}`$ in $`\mathrm{\Delta }_c`$. Because of this $`z`$-dependence, the above decomposition of $`\mathrm{\Psi }(𝒓)`$ is not fully spherically symmetric. It reflects the fact that an approximate radial solution on the left side of the guide at some $`z_1=|z_0|`$ needs not coincide necessarily with one on the right side at $`z_2=|z_0|`$ since the problem has no perfect spherical symmetry.
A better understanding of the role of this $`\gamma _l(z)`$ arises when one tries to write it in terms of the $`T_l`$’s that determine the amplitudes $`f_n^\pm `$ in Eq.(18) and the solution in Eq.(21). Expanding $`e^{\sigma _{zz^{}}pz^{}}J_0(q\rho ^{})`$ (see above) in the integrand in the definition of $`\gamma _l(z)`$ one obtains
$`\gamma _l(z)`$ $`=`$ $`{\displaystyle \underset{s=0}{\overset{\mathrm{}}{}}}(2s+1)\left[{\displaystyle _0^{p_c}}𝑑pP_l(ip/k)P_s(ip/k)\right]`$
$`\times {\displaystyle }d^3𝒓^{}{\displaystyle \frac{(\sigma _{zz^{}})^{l+s}}{i^s\mathrm{\hspace{0.17em}4}\pi }}[j_s(kr^{})P_s(\mathrm{cos}\theta ^{})]v(r^{})\mathrm{\Psi }(𝒓^{}),`$
assuming the summation and integration can be interchanged (see section V.3). If it happened that $`(\sigma _{zz^{}})^{l+s}=1`$, then $`\gamma _l(z)`$ would be a linear sum of the elements $`T_s`$ and the scattering solution in Eq.(21) would have spherical symmetry. Such a constant $`\gamma _l(z)`$ could be possible if one could neglect in the series for $`\gamma _l(z)`$ terms for which the parity of $`P_s`$ is different from that of $`P_l`$, such that for each $`l`$, only terms satisfying $`l+s=\mathrm{even}`$ would be retained. This condition can be related to the fact that the confining potential $`u(\rho )`$ only couples angular momenta with the same parity, in the sense that $`l|\widehat{u}|s𝑑\theta ^{}P_l(\mathrm{cos}\theta ^{})u(r\mathrm{sin}\theta ^{})P_s(\mathrm{cos}\theta ^{})=0`$ if $`l+s=\mathrm{odd}`$. Therefore, a consistent approximation to the scattering solution $`\mathrm{\Psi }(𝒓)`$ is found if one neglects odd parity combinations in the series defining $`\gamma _l(z)`$ and sets $`\gamma _l(z)\gamma _l`$, such that
$`\gamma _l(z)`$ $``$ $`\gamma _l`$ (22)
$``$ $`{\displaystyle \underset{s[l]}{}}(2s+1)\left[k{\displaystyle _0^{p_c/k}}𝑑xP_l(ix)P_s(ix)\right]T_s`$
$``$ $`{\displaystyle \underset{s[l]}{}}(2s+1)P_{ls}T_s,l=0,1,2,\mathrm{},`$
where $`s[l]`$ denotes the sum over even (odd) $`s`$ for even (odd) $`l`$.
Using Eq.(22), the scattering phase-shifts $`\delta _l`$ can now be introduced, since the expansion in Eq.(21) becomes spherically symmetric. Indeed, this expansion can be conveniently rewritten by the introduction of constants $`c_l^{}`$ and $`\delta _l`$ defined, for the moment, by ($`l=0,1,2,\mathrm{}`$)
$`c_l^{}\mathrm{cos}\delta _l`$ $``$ $`i^l(2l+1)\left[\alpha _l+\gamma _li\gamma kT_l\right],`$ (23a)
$`c_l^{}\mathrm{sin}\delta _l`$ $``$ $`i^l(2l+1)kT_l,`$ (23b)
such that our solution has the form
$$\mathrm{\Psi }(𝒓)\underset{l=0}{\overset{\mathrm{}}{}}c_l^{}\left[\mathrm{cos}\delta _lj_l(kr)\mathrm{sin}\delta _ln_l(kr)\right]P_l(\mathrm{cos}\theta ).$$
(24)
In order to identify the $`\delta _l`$, one notes, on one hand, that this wavefunction $`\mathrm{\Psi }(𝒓)`$ can be directly obtained as a spherically symmetric series $`_lc_lR_l(r)P_l(\mathrm{cos}\theta )`$ in the region $`rR_U`$, where $`u(\rho )`$ is negligible, if one solves the Schrödinger equation moving outward from the origin $`𝒓=0`$ (see Fig 2). The radial part $`R_l`$ can then be separately determined and for $`R_VrR_U`$ it should be given by Eq.(24), provided one chooses $`c_l=c_l^{}`$. On the other hand, the solution for the scattering potential $`v(r)`$ without confinement is given by a series with the same radial parts $`R_l`$ but different constants $`c_lc_l^{}`$ (see Section IV.1). The asymptotics of $`R_l`$ in this free space case is then also given by Eq.(24) so that $`\delta _l`$ are indeed the phase-shifts of the unconfined scattering problem. The difference between the confined and free space solutions, in the region $`R_VrR_U`$, is thus accounted for by the different constants $`c_l^{}`$ and $`c_l`$, respectively, related to distinct boundary conditions. Besides, $`\delta _l`$ depends solely on the solution of $`R_l`$ in the interior region $`r<R_V`$ mott1965 . The above defining relations for $`c_l^{}`$ and $`\delta _l`$ in Eqs.(23) are then *equations* that determine $`c_l^{}`$ and $`T_l`$ in terms of $`\delta _l`$, namely,
$$c_l^{}=\frac{(2l+1)(\alpha _l+\gamma _l)i^l}{\mathrm{cos}\delta _li\gamma \mathrm{sin}\delta _l},T_l=\frac{\alpha _l+\gamma _l}{i\gamma kk\mathrm{cot}\delta _l},$$
Finally, using Eq.(22) one obtains the matrix equation relating $`T_l`$ to the full ensemble of 3D free space scattering phase-shifts $`\delta _l`$ ($`l=0,1,2,\mathrm{}`$)
$$\left(i\gamma kk\mathrm{cot}\delta _l\right)T_l=\alpha _l+\underset{s[l]}{}(2s+1)P_{ls}T_s.$$
(25)
The coupling of angular momenta brought about by the confinement is a result of $`\mathrm{\Delta }_c`$, which accounts for the confining geometry and the discreteness of the low lying transversal states that should be resolved at low energies. The main equations that allow us the analysis of the effective 1D scattering in confined geometries are Eqs.(5), Eq.(18) and Eq.(25). The probability conservation condition Eq.(5c) serves to gauge the range of validity of the results. Although the above result is enough to present the formalism, improvements to Eq.(25) can be systematically made if necessary, e.g. by better dealing with the poles $`q_n`$ in Eq.(16a) and the constants $`\alpha _l`$.
## V Confinement induced effects
We now use the above formalism to analyze confinement induced phenomena that occur both in the scattering by a central field and in collision processes. Although related to each other, two important phenomena can be distinguished: resonant scattering (sections V.2 and V.4) and weakly localized states induced by the confinement (section V.5). In the present paper, our formalism is focused on the single mode regime. It is shown also that previous results of atom-atom collisions under parabolic confinement are recovered, regarding not only the scattering resonances (CIR) olshanii1998a ; moore2004a ; granger2004a but also the confinement induced weakly bound states bergeman2003 . CIR in atomic collisions involving excited transversal states at energies $`k^2>q_1^2`$ are treated in moore2004a for interactions in the $`s`$-wave approximation and in granger2004a for general $`V(r)`$. See also independent related works for 2DEG mesoscopic systems chu1989a ; bagwell1990a ; gurvitz1993a ; bardarson2004a .
### V.1 Single mode regime
In this regime, the total scattering energy $`E=\mathrm{}^2k^2/2\mu `$ allows for only the transversal ground state $`\phi _0(\rho )`$ to be effectively occupied $`k^2=q_0^2+k_0^2<q_1^2`$, i.e., to be an open channel. All excited states can only be virtually populated due to the scattering of the incoming state. Then $`n_E=0`$ and $`b_n=\delta _{n,0}`$. In terms of the even $`f_{0g}`$ and odd $`f_{0u}`$ angular momenta decomposition of $`f_0^\pm `$ given in Eq.(18), namely $`f_0^\pm =f_{0g}\pm f_{0u}`$, the current conservation can be written as
$$\left(\mathrm{Re}\{f_{0g}\}+|f_{0g}|^2\right)+\left(\mathrm{Re}\{f_{0u}\}+|f_{0u}|^2\right)=0.$$
(26)
A suitable parametrization of $`f_{0g,u}`$ is obtained by introducing 1D scattering parameters $`\delta _g`$ and $`\delta _u`$, such that
$`f_{0g}`$ $``$ $`{\displaystyle \frac{1}{1+i\mathrm{cot}\delta _g}},`$ (27a)
$`f_{0u}`$ $``$ $`{\displaystyle \frac{1}{1+i\mathrm{cot}\delta _u}}.`$ (27b)
In this form, Eq.(26) is verified if $`\delta _{g,u}`$ are arbitrary real numbers, since then $`\mathrm{Re}\{f_{0g,u}\}+|f_{0g,u}|^2=0`$ vanishes separately. There is by now no a priori reason for the parameters $`\delta _{g,u}`$ to be both real. However, in the case of general potentials $`v(r)`$ and $`u(\rho )`$ considered in this work, it will be seen in the next sections that Eqs.(27) with real $`\delta _{g,u}`$ provide a valuable means to establish the boson-fermion and fermion-boson mappings as well as the conservation condition Eq.(26).
### V.2 Boson-fermion and fermion-boson mappings
Consider now a collisional process between *identical* particles. It is convenient to split $`\mathrm{\Psi }(𝒓)`$ into symmetric $`\psi _g(z)\phi _0(\rho )`$ and antisymmetric $`\psi _u(z)\phi _0(\rho )`$ parts, with respect to $`𝒓𝒓`$,
$`\mathrm{\Psi }(𝒓)`$ $`=`$ $`\left[e^{ik_0z}+f_0^\pm e^{ik_0|z|}\right]\phi _0(\rho )`$
$`=`$ $`\left[\psi _g(z)+\psi _u(z)\right]\phi _0(\rho ),|z|>R_{1D}^{}.`$
$`\psi _g(z)`$ and $`\psi _u(z)`$ can be deduced by using the form Eq.(18) for the amplitude $`f_0^\pm `$ in terms of its even $`f_{0g}`$ and odd $`f_{0u}`$ sectors. For the symmetric part, we obtain,
$`\psi _g(z)`$ $`=`$ $`(1+f_{0g})\mathrm{cos}(k_0z)+if_{0g}\mathrm{sin}(k_0|z|)`$ (28a)
$`=`$ $`e^{i\delta _g}\mathrm{cos}(k_0|z|+\delta _g),|z|>R_{1D}^{},`$
where Eq.(27a) has been employed. The antisymmetric part is given by
$`\psi _u(z)`$ $`=`$ $`i(1+f_{0u})\mathrm{sin}(k_0z)\pm f_{0u}\mathrm{cos}(k_0z)`$ (28b)
$`=`$ $`ie^{i\delta _u}\mathrm{sin}(k_0z\pm \delta _u),|z|>R_{1D}^{},`$
where the plus (minus) sign in $`\psi _u(z)`$ refers to $`z>0`$ ($`z<0`$) and Eq.(27b) has been used.
If the colliding particles are e.g. spin polarized $`\mathrm{\Psi }(𝒓)`$ must be (anti-) symmetrized. For bosons, the correct effective 1D scattered wavefunction in the asymptotic region $`|z|>R_{1D}^{}`$ is $`\psi _g(z)`$. The resonance condition can be identified as $`f_{0g}=1`$ or $`\delta _g=|\pi /2|`$ and it follows from Eq.(28a) that the wavefunction becomes $`\psi _g(z)=i\mathrm{sin}(k_0|z|)`$, which is with respect to its modulus the wavefunction of a pair of non-interacting fermions, as can be seen by setting $`f_{0u}0`$ (or $`\delta _u0`$) in Eq.(28b). This is the well-known fermionic mapping of strongly interacting, inpenetrable bosons in 1D tonks1936 ; girardeau1960 ; lenard1966 ; paredes2004a . The bosons would not be allowed to be located at the same position $`z=0`$, supposing the asymptotic solution Eq.(28a) could be extended towards the origin $`|z|R_{1D}^{}`$.
For fermions, a reciprocal mapping exists at resonance
$$f_{0u}=1,\mathrm{or}\delta _u=|\pi /2|.$$
(29)
From Eq.(28b), the wavefunction becomes $`\psi _u(z)=\mathrm{cos}(k_0z)`$, which is the wavefunction of a pair of non-interacting bosons granger2004a . Clearly, from the exact, fully antisymmetric wavefunction of a pair of spin polarized fermions, the probability of finding both of them at the origin $`𝒓=0`$ is strictly zero. The bosonization implied by this confinement induced resonance means that the asymptotic behavior of the fermions far from the interaction region is that of free bosons.
These mappings between identical colliding particles are valid also under a longitudinal confinement along the $`z`$-axis. In order to neglect couplings to the center of mass, this longitudinal confinement can be chosen to be parabolic. The characteristic oscillator length $`R_{}`$ should then be flat enough, such that $`R_{}R_{1D}^{}`$. The asymptotics of the symmetric and antisymmetric eigenstates with (positive) energies $`E_{g,u}=\mathrm{}^2k_{g,u}^2/2\mu `$ are given by Eq.(28a) and Eq.(28b), respectively. The corresponding eigenvalues $`k_{g,u}^2=q_0^2+k_{0g,u}^2`$ can be estimated by solving the equations $`\psi _{g,u}(R_{})0`$ for $`k_{0g,u}`$, namely, $`\mathrm{cos}(k_{0g}R_{}+\delta _g)=0`$ and $`\mathrm{sin}(k_{0u}R_{}+\delta _u)=0`$. If the parameter $`\delta _g=\delta _g(k_{0g})`$ is made equal to $`\pm \pi /2`$ by varying, e.g., $`R_V`$, $`R_U`$, or $`R_{}`$, one sees that $`\psi _g(z)`$ is once more mapped onto an antisymmetric eigenstate without interaction $`v(r)0`$ (or $`\delta _u0`$), as already seen above for free longitudinal motion. Analogously, $`\psi _u(z)`$ is mapped onto a non-interacting symmetric eigenstate as $`\delta _u=\delta _u(k_{0u})\pm \pi /2`$. This latter type of bosonic mapping of fermionic eigenstates under longitudinal confinement was first verified numerically in granger2004a for the collision interaction $`V(r)=d/\mathrm{cosh}(r/b)^2`$.
Note that, if one is not restricted to the case of collisions, this result can be formally quite general regarding not only the interaction potential $`v(r)`$, but also some sufficiently flat transversal and longitudinal confinements.
### V.3 The effective amplitude $`f_0^\pm `$
The considerations in the preceding section assumed the general form Eqs.(27) for $`f_{0g,u}`$. In this section, $`f_0^\pm `$ is explicitly calculated by solving the matrix equation Eq.(25) for the elements $`T_l`$. From $`k1/R_U`$, it follows that $`kR_VR_V/R_U1`$, which is the condition of low scattering energy in 3D free space scattering and the phase-shifts
$$\mathrm{tan}\delta _l=\mathrm{tan}\delta _l(k)k^{2l+1}1/R_U^{2l+1}$$
are generally small mott1965 for large $`R_U`$ (or small $`R_V`$). One expects then that orbital angular momenta higher than the leading contributions, e.g. the $`s`$\- and $`p`$-waves, should not significantly change the main features arising from the latter. Indeed, the calculation of $`f_{0g}`$ and $`f_{0u}`$ can be done separately, since the even and odd angular momenta in Eq.(25) are uncoupled from each other. In the single mode regime, Eq.(16a) becomes $`\gamma =2/d_U^2kk_0`$ and $`d_UR_U/N_0`$, where $`k_0=(k^2q_0^2)^{1/2}`$ and $`f_0^\pm `$ in Eq.(18) depends only on the ratio $`t_lT_l/k_0`$. The matrix equation Eq.(25) in this case can be rewritten in terms of these $`t_l`$. Thus for all $`l`$,
$`\left\{2i\left[k\mathrm{cot}\delta _l+(2l+1)P_{ll}\right]d_U^2k_0\right\}t_l`$
$`=`$ $`\alpha _ld_U^2+k_0d_U^2{\displaystyle \underset{s[l]l}{}}(2s+1)P_{ls}t_s,`$
where the diagonal term $`s[l]=l`$ is excluded from the series on the rhs. Assume now in this equation for $`t_l`$ that this series converges. As $`l`$ increases, $`\alpha _l=\alpha _{l0}=P_l(k_0/k)/\pi ^{1/2}d_U`$ (see Eq.(21) and Eq.(17)) decrases, while $`P_{ls}`$ (see Eq.(22)) should decrease mildly due to the oscillation of $`P_l(ix)`$. Then the rhs of this equation for $`t_l`$ should decrease relatively smoothly with $`l`$. On the other hand, the phase $`\delta _l`$ decreases exponentially $`\delta _l1/R_U^{2l+1}`$ with $`l`$ at low energies, so that $`\mathrm{cot}\delta _l`$ increases exponentially. Then $`t_l`$ on the lhs should be exponentially small with increasing $`l`$. This is consistent with the assumption of convergence of the series on the rhs \[see also discussion on $`\gamma _l(z)`$ after Eq.(21)\]. Apart from this, one can use the smallness of the phase-shifts $`|\delta _2||\delta _0|`$ and $`|\delta _3||\delta _1|`$ in order to explicitly diagonalize the finite subspaces $`l=0,2`$ and $`l=1,3`$ separately. Then the couplings to $`l=2`$ and $`l=3`$ are seen to have a minor effect to the leading terms $`l=0`$ and $`l=1`$, respectively. Hence we put
$$t_l\frac{\frac{1}{2}d_U^2\alpha _l}{i\left[k\mathrm{cot}\delta _l+(2l+1)P_{ll}\right]\frac{1}{2}d_U^2k_0},l=0,1,\mathrm{}.$$
The even sector $`f_{0g}`$ is then given by the leading term, e.g, the $`s`$-wave, for which $`k\mathrm{cot}\delta _01/a`$, where $`a`$ is the $`s`$-wave 3D free space scattering length mott1965 . From Eq.(18) one obtains
$`f_{0g}`$ $``$ $`{\displaystyle \underset{l=0,2,\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{4\pi (2l+1)\alpha _{l0}}{2i}}t_l`$ (30a)
$``$ $`{\displaystyle \frac{1}{1+i\left[\frac{d_U^2}{2a}\left(1aP_{00}\right)\right]k_0}}.`$
The odd sector $`f_{0u}`$ is approximated by e.g. the $`p`$-wave. In this case, the relevant quantity suno2003a is the scattering volume $`V_p`$ related to $`\delta _1`$ by $`k\mathrm{cot}\delta _11/V_pk^2`$. Then
$`f_{0u}`$ $``$ $`{\displaystyle \underset{l=1,3,\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{4\pi (2l+1)\alpha _{l0}}{2i}}t_l`$ (30b)
$``$ $`{\displaystyle \frac{3[P_1(k_0/k)]^2}{1+i\left[\frac{d_U^2}{2}\left(\frac{1}{V_pk^2}3P_{11}\right)\right]k_0}}.`$
The range of validity of these results can be obtained by imposing e.g. the conservation of probability. As discussed in connection to Eq.(26) and Eqs.(27), the 1D scattering phase-shift $`\delta _g`$ is real for the result Eq.(30a). Then a sufficient condition for Eq.(26) to hold is
$$|k_0|k1/R_U$$
(31a)
since then $`3[P_1(k_0/k)]^21`$ so that $`f_{0u}0`$ and the parameter $`\delta _u0`$ can be considered also real. In this first region of low longitudinal momenta, the angular momenta couplings in the matrix equation for $`t_l`$ should become negligible due to the factor $`k_0`$ that multiplies the series $`d_U^2_{s[l]l}(2s+1)P_{ls}t_s`$ containing other angular momenta. The present uncoupled solution to $`t_l`$ should then be a good approximation. The second region is at relatively large momenta, when $`3[P_1(k_0/k)]^21`$, i.e., when
$$|k_0|q_0/\sqrt{2}1/R_U$$
(31b)
since then $`\delta _u`$ can also be taken as real. The decoupling between odd angular momenta rests then on the smallness of the phase-shifts $`\delta _l`$. These are, nevertheless, the most interesting regions for the momentum $`k_0`$ that are relevant to resonant scattering in symmetric and antisymmetric states, respectively. For more detailed applications these regions can be enlarged by improving on Eq.(25).
### V.4 Confinement induced resonances - CIR
As presented in sections V.1 and V.2 in general terms, the CIR is one remarkable effect arising due to the presence of a strong confinement. In the following, it is explicitly expressed in terms of the parameters determining a given scattering potential and a given confinement.
#### V.4.1 Even angular momenta
The symmetric wavefunction in the single mode regime is $`\psi _g(z)\phi _0(\rho )`$ and only the symmetric sector $`f_{0g}`$ contributes. It can be used to describe e.g. spin-polarized bosons or fermions in an antisymmetric spin state. In the expression Eq.(30a) for this $`f_{0g}`$, $`P_{00}`$ is given in Eq.(22), thus
$`P_{00}`$ $`=`$ $`k{\displaystyle _0^{p_c/k}}𝑑xP_0(ix)P_0(ix)`$ (32)
$`=`$ $`\sqrt{q_1^2k^2},`$
having used the definition of $`p_c`$ in Eq.(16b) and the single mode condition $`n_E=0`$.
At *small* scattering lengths $`|a|R_U`$, one can neglect this factor $`P_{00}`$ in Eq.(30a), since then $`|a|P_{00}1`$, and
$`f_{0g}`$ $``$ $`{\displaystyle \frac{1}{1+i\left[\frac{d_U^2}{2a}\right]k_0}},|a|R_U.`$
This is also true at “high” energies $`|k_0|`$ when $`kq_1`$ such that $`P_{00}0`$. One must, however, be careful about the validity of the present approximations in this region of large momenta $`|k_0|`$, as mentioned above in the context of Eq.(31b). The only possibility for resonance $`f_{0g}=1`$ is then at zero energy $`k_0=0`$.
At *large* scattering lengths $`R_V|a|R_U`$, this factor $`P_{00}`$ plays an important role in the scattering process at low momenta $`|k_0|1/R_U`$ for which $`P_{00}\sqrt{q_1^2q_0^2}`$ is maximum. Indeed, one obtains then, for $`|k_0|1/R_U`$,
$$f_{0g}\frac{1}{1+i\left[\frac{d_U^2}{2a}\left(1\frac{C^{}a}{d_U}\right)\right]k_0},$$
(33a)
where the constant $`C^{}`$ depends *only* on the first closed and the last open channels, namely
$`C^{}`$ $``$ $`d_U\sqrt{q_1^2q_0^2}`$ (33b)
$`=`$ $`{\displaystyle \frac{1}{(\mathrm{}^2/2\mu d_U^2)^{1/2}}}\sqrt{ϵ_1ϵ_0}.`$
Recalling that $`f_0^\pm f_{0g}`$, since $`f_{0u}0`$ for such small momenta $`|k_0|`$ (see Eqs.(30)), this result implies that the scattering can be described by means of an effective 1D interaction potential $`V_{1D}(z)`$ of zero range
$$V_{1D}(z)=g_{1D}\delta (z),g_{1D}=\frac{\mathrm{}^2}{\mu }\frac{2a}{d_U^2}\frac{1}{1C^{}a/d_U}.$$
In other words, Eq.(33a) is the solution to $`f_0^\pm `$ for the hypothetical pure 1D scattering problem under this potential $`V_{1D}(z)`$ and is valid also for distinguishable particles.
As discussed in section III for atom-atom collisions, this same 1D potential $`V_{1D}(z)`$ at low momenta $`|k_0|`$ was first shown in olshanii1998a to be a direct result of parabolic confinement and of the 3D zero range interaction potential $`V_\delta (r)`$ of Eq.(6a). The occurence of this collisional resonance at $`a_{}=Ca`$ for general interactions $`V(r)`$ was then shown in granger2004a .
For other types of scattering processes, e.g. in mesoscopic 2DEG systems chu1989a ; bagwell1990a ; gurvitz1993a ; bardarson2004a , Eq.(33a) shows therefore that at low longitudinal momenta $`k_0`$ the existence of a resonance of the type $`f_0^+f_{0g}=1`$ is characterized by an infinite effective 1D coupling strength
$$|g_{1D}|\mathrm{},d_UC^{}a,$$
(34)
and is indeed not restricted to atom-atom collisions under parabolic confinement or to two-dimensional quantum point-contacts with a single impurity gurvitz1993a . It is a general effect of single mode scattering at low velocities under quasi-1D confinement.
#### V.4.2 Odd angular momenta
For odd angular momenta, the procedure is similar to the previous analysis. The antisymmetric part $`\psi _u(u)`$ shows a scattering resonance when the associated amplitude $`f_{0u}=1`$. At this resonance, the system can be mapped to a non-interacting symmetric wavefunction as discussed in connection to Eq.(29). The term $`P_{11}`$ is given in Eq.(22), from which the actual condition of CIR in the amplitude $`f_{0u}`$ in Eq.(30b) can be worked out for given scattering and confining potentials. Note that due to the factor $`3[P_1(k_0/k)]^2`$ in the numerator of Eq.(30b), the odd component $`f_{0u}`$ is non-negligible only at relatively large longitudinal momenta $`|k_0|1/R_U`$, in contrast to the even amplitude $`f_{0g}`$ discussed above.
For atom-atom collisions under parabolic confinement, the existence of CIR beyond the $`s`$-wave was first predicted by neglecting couplings to orbital angular momenta and isolating a single partial wave $`l>0`$ and its phase-shift $`\delta _l`$ as the sole dominant contribution granger2004a . A direct consequence of contributions from non-zero partial waves $`l>0`$ is, among others, the possibility to define effective $`p`$-wave *zero-range* scattering potentials analogous to $`V_\delta (r)`$ and $`V_{1D}(z)`$, but which act only on antisymmetric wave functions of a colliding pair kanjilal2004a ; girardeau2004a .
### V.5 Weakly bound states
Besides CIR, a second effect of strong cylindrical confinement is the prediction of weakly bound and localized states. For atom-atom collisions this implies the formation of quasi-molecules. As first shown in bergeman2003 for zero-range atom-atom interactions under parabolic confinement, these bound states exist even when no bound state occurs in free space. They have recently been observed in optical traps loaded with <sup>40</sup>K atoms moritz2005a . Similar bound states localized around impurities are found in independent studies of mesoscopic systems chu1989a ; bagwell1990a ; gurvitz1993a ; bardarson2004a . In this section, a localized state for general cylindrical confinement is calculated from the even, more precisely zero, angular momentum sector.
#### V.5.1 $`s`$-wave binding energy and wavefunction
For the *free space* $`l=0`$ angular momentum, the bound state of the (attractive) scattering potential $`V(r)`$ has an energy $`E_{Bf}`$ given by
$$E_{Bf}\frac{\mathrm{}^2\kappa _{Bf}^2}{2\mu }\frac{\mathrm{}^2}{2\mu a^2},aR_V$$
whose relationship to the scattering length $`a`$ via $`\kappa _{Bf}1/a`$ holds only when $`aR_V>0`$ mott1965 . At negative scattering lengths $`a<0`$, one can speak of a “virtual” state close to be incorporated to the spectrum of $`V(r)`$. Under lateral confinement and along the radial direction of the cylindrical trap, the tail of its wavefunction $`e^{\kappa _{Bf}r}`$ is changed to vanish at a finite distance from the origin, i.e. at the edge $`r=\rho =R_U`$ of the guide when $`|a|R_U`$ (see Fig. 1). This lateral squeeze lifts $`E_{Bf}<0`$ by an amount $`ϵ_0`$. It can be sufficient for this state to pass the threshold without confinement $`E=0`$ as $`R_U`$ decreases further.
This confined localized state with energy $`E_B`$ satisfies Eq.(4a) with $`k^2`$ replaced by $`2\mu E_B/\mathrm{}^2`$. This replacement is equivalent to redefining $`k_0`$ to be the imaginary number $`k_{0B}\pm i(q_0^22\mu E_B/\mathrm{}^2)^{1/2}`$. Since the wave $`e^{ik_{0B}z}`$ should be absent from the solution $`\mathrm{\Psi }`$ and $`e^{ik_{0B}|z|}`$ should show a bound state like exponential decay, $`1/f_0^\pm (k_{0B})`$ must vanish at a *positive* imaginary value of $`k_{0B}`$, such that the interaction part $`e^{ik_{0B}|z|}`$ outweights $`e^{ik_{0B}z}`$. Let then $`k_{0B}ix_B/a`$, so that $`x_B>0`$ if $`a>0`$ and $`x_B<0`$ if $`a<0`$. In terms of this $`x_B`$, $`E_B`$ can be rewritten as
$$E_B\frac{\mathrm{}^2}{2\mu }\left(q_0^2\frac{x_B^2}{a^2}\right),ax_B>0.$$
(35)
From Eqs.(30), a *root* of $`1/f_0^\pm `$ arises either from the even sector $`1/f_{0g}`$ or from the odd sector $`1/f_{0u}`$. In order to recover the $`l=0`$ free space state $`E_{Bf}`$, this root should come from the even sector Eq.(30a). Using then the substitution $`k^22\mu E_B/\mathrm{}^2=q_0^2x_B^2/a^2`$ in the expression Eq.(32) for $`P_{00}`$, the equation for $`x_B`$ becomes
$$2s^2+\left[1\frac{a}{|a|}\sqrt{(sC^{})^2+x_B^2}\right]x_B=0,s\frac{a}{d_U},$$
(36a)
where $`|a|`$ arises as $`1/a^2`$ is taken out of the square root in $`P_{00}`$. Another enlarged form is the quartic equation
$$x_B^4+\left(C^2s^21\right)x_B^24s^2x_B4s^4=0.$$
(36b)
For each $`|s|`$, one should pick up the correct positive (negative) root $`x_B`$ if $`s>0`$ ($`s<0`$), as indicated in Eq.(35).
The weak confinement limit is $`|s|0`$. A finite root $`x_B+1`$ exists only for $`a>0`$ as can be seen from Eq.(36a). This means that $`E_BE_{Bf}`$, when $`R_U\mathrm{}`$ or when $`a0^+`$ as expected. For $`a<0`$, a trial Taylor expansion $`x_Bx_1s+x_2s^2`$ in Eq.(36b) yields $`x_1=0`$ and $`x_2=2`$. A localized state with energy
$$E_Bϵ_04\left(\frac{a}{d_U}\right)^4|E_{Bf}|,|a|d_U,$$
is thus expected also for negative scattering lengths. In this case, the confinement has been able to turn the“virtual” state $`E_{Bf}`$ of $`V(r)`$ into a real bound state localized around the scattering center.
The strong confinement limit is $`|s|1`$ (assuming the basic requirement $`R_VR_U`$ is not violated). Setting then $`x_Bx_{\mathrm{}}/s^1`$ in Eq.(36b) and reexpressing it in terms of $`s^1`$ gives for positive and negative $`a`$, respectively, $`x_{\mathrm{}}=\pm \{[(16+C^4)^{1/2}C^2]/2\}^{1/2}`$. For both positive and negative roots $`x_B`$, the same fraction of the transversal ground state energy $`ϵ_0`$ is reached as $`|s|\mathrm{}`$
$$E_B\left[1\left(\frac{x_{\mathrm{}}}{q_0d_U}\right)^2\right]ϵ_0,|a|d_U.$$
For parabolic guide $`d_U=a_{}`$ \[see Eq.(16a) with the exact value $`|\phi _n(0)|^2=1/\pi a_{}^2`$\]. Then $`C^{}=2`$ and $`q_0a_{}=\sqrt{2}`$, both obtained by directly using the exact energies $`ϵ_0\mathrm{}^2q_0^2/2\mu =\mathrm{}\omega _{}`$ and $`ϵ_1\mathrm{}^2q_1^2/2\mu =3\mathrm{}\omega _{}`$ in terms of $`\omega _{}`$. Hence $`E_B=(2\sqrt{2})ϵ_0=0.586ϵ_0`$ in good quantitative agreement with bergeman2003 ; moritz2005a . For square well confinement, it follows from Eq.(13b) and Eq.(33b) that $`C^{}=(20/3)^{1/2}=2.582`$ and $`q_0d_U=(3/2)^{1/2}=1.225`$, thus $`E_B=0.631ϵ_0`$. These energies are plotted in Fig. 3 as function of the confinement length scale $`d_U`$.
The spatial shape $`\mathrm{\Psi }_B(𝒓)`$ of this confinement induced weakly bound state in the region *far* from the scattering center is
$$\mathrm{\Psi }_B(𝒓)e^{x_B|z|/a}\phi _0(\rho ),|z|R_{1D}^{},$$
(37a)
setting aside overall constants and where $`x_B`$ is the proper root of Eq.(36). This asymptotic form is common to both positive and negative scattering lengths $`a`$, with $`ax_B>0`$. *Closer* to the origin, the wavefunction can be calculated by analytic continuation of the intermediate asymptotics in Eq.(21) and Eq.(22). There the function $`T_0=T_0(k_0)`$ diverges at $`k_0=ix_B/a`$, which is the condition for the pole of $`f_0^\pm `$. After this analytic continuation, the step towards Eq.(24) can no longer be taken (it was taken in order to calculate $`T_0=T_0(k_0)`$ for real $`k_0`$). In this series Eq(21), this divergence then singles out terms proportional to the constant $`T_0(ix_B/a)\mathrm{}`$. As a result, apart from overall constants, the wavefunction $`\mathrm{\Psi }_B(𝒓)`$ has the following form (for $`R_VrR_U`$)
$`\mathrm{\Psi }_B(𝒓)`$
$``$ $`\left[\sqrt{q_1^2\kappa _B^2}{\displaystyle \frac{2/d_U^2}{\sqrt{q_0^2\kappa _B^2}}}\right]j_0(\kappa _Br)+\kappa _Bn_0(\kappa _Br),`$
where $`\kappa _B^22\mu E_B/\mathrm{}^2`$ and $`E_B`$ is given in Eq.(35). For $`E_B<0`$, it can be seen that both imaginary values of $`\kappa _B`$ yield the same result. In particular, the free space bound state is recovered in the limit $`a/d_U0^+`$. Using $`E_BE_{Bf}\mathrm{}^2/2\mu a^2`$ in Eq.(37), one has
$$\mathrm{\Psi }_B(𝒓)\frac{e^{r/a}}{r},a/d_U0^+.$$
For negative scattering length $`a/d_U0^{}`$, the free space “virtual” bound state turns into a real one whose energy was calculated in the previous section, $`E_Bϵ_04(a/d_U)^4|E_{Bf}|`$. The wavefunction Eq.(37) becomes instead
$$\mathrm{\Psi }_B(𝒓)\frac{1}{|a|}j_0(q_0r)+q_0n_0(q_0r),a/d_U0^{}.$$
On the opposite limit of strong confinement $`|a|/d_U\mathrm{}`$, $`E_B`$ tends to a positive fraction $`\lambda `$ of $`ϵ_0`$, and the wavefunction becomes a superposition of both terms $`j_0(\sqrt{\lambda }q_0r)`$ and $`n_0(\sqrt{\lambda }q_0r)`$.
#### V.5.2 Higher angular momenta
Localized states from higher angular momenta contributions and their *couplings* among each other can in principle be systematically obtained by solving Eq.(25) and analyzing the poles of the full amplitude $`f_0^\pm =f_{0g}\pm f_{0u}`$. For instance, a $`d`$-wave pole should come from retaining the $`l=2`$ contribution. This and other poles from the odd sector $`f_{0u}`$ are treated in detail elsewhere.
## VI Discussion
Here the role of the continuum approximation in section IV.2 is discussed by comparing with previous models for scattering under confinement. In the interior of the guide $`r,r^{}R_U`$ where the confining potential is negligible, the Green’s function $`G_c(𝒓,𝒓^{})`$ under confinement approaches the 3D free space Green’s function $`G(𝒓,𝒓^{})`$, as shown in Eq.(12). Therefore, the singularity of $`G_c(𝒓,𝒓^{})`$ as $`r,r^{}0`$ is not only coincident with (see section III), but it is essentially *the* well-known singularity $`1/|𝒓𝒓^{}|`$ of the 3D scattering scenario without confinement. In addition, one expects that the details of the confinement, whether parabolic or not, should not be important for the physical understanding. In fact, the spectrum of *excited* transversal states sums up to yield $`G(𝒓,𝒓^{})`$ (integrated over the axial angle $`\varphi ^{}`$), as the continuum limit shows. The only condition is a short-ranged scattering potential $`R_VR_U`$ or a correspondingly “flat” confinement.
The remarkable effects induced by the confinement are in turn directly related to the *low lying* transversal states as also demonstrated in gurvitz1993a for a two-dimensional problem. The discreteness of these states is then captured in the factor $`\gamma `$ and in the correction term $`\mathrm{\Delta }_c(𝒓,𝒓^{})`$, which introduce renormalizations and couplings of orbital angular momenta not accounted for by $`G(𝒓,𝒓^{})`$ alone.
For a more quantitative analysis, consider the specific case of atom-atom collisions in the low energy $`s`$-wave approximation and under parabolic confinement olshanii1998a ; moore2004a ; bergeman2003 . As discussed in section III, the *discrete* summation over the transversal states can be dealt with in order to extract the singularity of the scattering solution and to determine the value of the constant $`C=1.4603\mathrm{}`$ in the effective scattering amplitude $`f_0^\pm `$. On the other hand, the corresponding value of $`C`$ in the *continuum* limit for this case is $`C\mathrm{lim}_s\mathrm{}(_0^s𝑑s^{}/\sqrt{s^{}}_{s^{}=1}^s1/\sqrt{s^{}})_0^1𝑑s^{}/\sqrt{s^{}}=2`$ (see Eq.(9) of olshanii1998a ), in agreement with $`C^{}=2`$ calculated in section V.5. This value $`C^{}=2`$ arises by using the exact harmonic oscillator energies $`ϵ_0=\mathrm{}\omega _{}`$ and $`ϵ_1=3\mathrm{}\omega _{}`$. It can be improved on by e.g. numerically computing $`\mathrm{\Delta }_c`$. As already pointed out in section IV.2, the exact excited spectrum $`ϵ_n2n`$ is substantially different from the “flat” potential approximation given in Eq.(13b). However, the discrete sum over this exact parabolic spectrum can be shown olshanii1998a ; moore2004a to have an equal singular behavior $`1/|z|`$ as in free space (see section III). Hence also within the discrete approach, this demonstrates the major role played by the lowest transversal states, whereas the sum over the excited states turns out to be qualitatively confinement-independent.
## VII Conclusions and Perspectives
An analytical treatment of quasi-1D quantum scattering by spherically symmetric and short but finite range potentials in general cylindrical confinement is developed. The full scattering wavefunction is calculated non-perturbatively without partially resuming perturbative series. All phase-shifts of the scattering potential can be readily incorporated. This formalism provides a unified physical picture of the process of confined quasi-1D scattering at low energies and can be applied to impurity scattering in mesoscopic 2DEG systems and two-body collisions. Following the reasoning related to Eq.(12) one expects to be able to treat non-cylindrical geometries and quasi-2D scattering in an analogous way. By computing then a few coefficients and functions such as $`\alpha _{ln}`$ and $`\mathrm{\Delta }_c`$ (or $`\mathrm{\Delta }_u`$) and only the lowest eigenstates of the confining potential, reasonable numerical results are also possible. Scattering resonances such as total reflexion, as well as weakly localized states can be determined. For the particular case of parabolic confinement, the formalism presented here can also be used to obtain e.g. confinement induced two-body weakly bound states. Of particular interest is the study of the poles of the scattering amplitude which correlate to orbital angular momenta beyond $`s`$-waves and how non-parabolic geometries can affect these unconventional pairings due to couplings to the center of mass coordinates.
###### Acknowledgements.
J.I.K. appreciates financial support from the Conselho Nacional de Desenvolvimento Científico e Tecnológico (CNPq) and the Alexander von Humboldt Foundation (AvH). J.S. thanks the DFG Schwerpunktprogramm: “Wechselwirkung in Ultrakalten Atom- und Molekülgasen” for financial support. |
warning/0506/cond-mat0506719.html | ar5iv | text | # Topological Coherent Modes in Trapped Bose Gas
## 1 Introduction
Trapped Bose gas, being dilute and cooled down to temperatures below the Bose-Einstein condensation temperature, forms a coherent system well described by the Gross-Pitaevskii equation \[1–5\]. Stationary solutions to the latter are called the topological coherent modes . The ground state, corresponding to the lowest energy level of the Gross-Pitaevskii stationary equation, represents the standard Bose-Einstein condensate. The states, associated with the higher energy levels, describe nonground-state condensates. The higher topological coherent modes can be generated by means of alternating fields, with the frequencies that are in resonance with the chosen transition frequencies .
The name topological coherent modes comes from the following. Different stationary solutions to the Gross-Pitaevskii equation, related to distinct energy levels, possess essentially dissimilar spatial shapes, with different number of zeros. The modes, displaying distinct spatial topology, can be named topological. These modes can be characterized as coherent due to the fact that the Gross-Pitaevskii equation can be interpreted as an exact equation for coherent states .
The topological coherent modes, described by the nonlinear Gross-Pitaevskii equation, should not be confused with elementary collective excitations, defined by the linear Bogolubov - de Gennes equations. The elementary excitations describe small oscillations around a given nonlinear topological mode and do not change the topology of the latter .
The general notion of the nonlinear coherent modes and the possibility of their generation by means of resonant fields was advanced in Ref. . A particular case of vortex creation was considered in Refs. . Various properties of topological coherent modes were also studied in theoretical papers \[10–27\]. A dipole topological mode was generated in experiment . The feasibility of resonant formation of multimode condensates, consisting of several topological coherent modes, was investigated .
Bose gas of trapped atoms with resonantly generated topological coherent modes exhibits a variety of unusual features. The aim of this report is to present a general picture of such a resonant system (Section 2) and to give a survey of its most interesting properties (Section 3).
## 2 Topological Coherent Modes
The topological coherent modes are defined as the stationary solutions to the eigenvalue problem
$$\widehat{H}\left[\phi _n\right]\phi _n\left(𝐫\right)=E_n\phi _n\left(𝐫\right),$$
(1)
which is represented by the stationary Gross-Pitaevskii equation with the nonlinear Hamiltonian
$$\widehat{H}\left[\phi \right]\frac{\mathrm{}^2}{2m}^2+U\left(𝐫\right)+NA_s\left|\phi \right|^2,$$
(2)
containing a trapping potential $`U\left(𝐫\right)`$ and the interaction intensity $`A_s4\pi \mathrm{}^2a_s/m`$, with $`m`$ being atomic mass; $`a_s`$ scattering length; and $`N`$ the total number of atoms.
In what follows, we shall use the notation for the scalar product $`(\phi _m,\phi _n)\phi _m^{}\left(𝐫\right)\phi _n\left(𝐫\right)𝑑𝐫`$. The eigenfunctions $`\phi _n\left(𝐫\right)`$, because of the nonlinearity of problem (1), are not necessarily orthogonal, so that $`(\phi _m,\phi _n)`$ is not compulsory the Kroneker delta $`\delta _{mn}`$. But the functions $`\phi _n\left(𝐫\right)`$ can always be normalized, with $`(\phi _n,\phi _n)\phi _n^2=1`$.
If the trapped Bose gas is initially in one of the coherent modes $`n`$, then to generate another mode requires applying an alternating field
$$V(𝐫,t)=V_1\left(𝐫\right)\mathrm{cos}\omega t+V_2\left(𝐫\right)\mathrm{sin}\omega t,$$
with a frequency $`\omega `$ being close to one of the transition frequencies
$$\omega _{mn}\frac{1}{\mathrm{}}\left(E_mE_n\right).$$
(3)
Say the modes with the energies $`E_1`$ and $`E_2`$ are connected, such that $`E_1<E_2`$, with the transition frequency being $`\omega _{21}`$. The related resonance condition reads as
$$\left|\frac{\mathrm{\Delta }\omega }{\omega }\right|1,\mathrm{\Delta }\omega \omega \omega _{21}.$$
(4)
With an applied time-dependent field, we have the temporal Gross-Pitaevskii equation
$$i\mathrm{}\frac{}{t}\phi (𝐫,t)=\left(\widehat{H}\left[\phi \right]+\widehat{V}\right)\phi (𝐫,t),$$
(5)
in which $`\widehat{V}=\widehat{V}(𝐫,t)`$. The alternating field can be represented as
$$V(𝐫,t)=\frac{1}{2}B\left(𝐫\right)e^{i\omega t}+\frac{1}{2}B^{}\left(𝐫\right)e^{i\omega t},$$
(6)
where $`B\left(𝐫\right)V_1\left(𝐫\right)iV_2\left(𝐫\right)`$. Experimentally, the alternating field can be realized as the modulation of the trapping potential. Another way could be by alternating the scattering length $`a_s\left(t\right)`$ by modulating an external magnetic field close to a Feshbach resonance \[31–33\].
One possibility of studying the resonant generation of coherent modes would be by a direct numerical solution of the temporal equation (5), which we have done in our works . However to get a deep physical insight into the problem, it is necessary to develop an analytical theory. Such a general theory, based on the averaging technique , was developed in Refs. \[6,12,15,20–24,29,30\]. In order to show that the basic equations of the analytical approach can be accurately derived and all consideration is well mathematically grounded, we describe, first of all, the main steps of this derivation.
We can look for the solution of Eq. (5) in the form of the mode expansion
$$\phi (𝐫,t)=\underset{n}{}c_n\left(t\right)\phi _n\left(𝐫\right)\mathrm{exp}\left(\frac{i}{\mathrm{}}E_nt\right),$$
(7)
with the coefficients $`c_n\left(t\right)`$ being slow functions of time, such that
$$\frac{\mathrm{}}{E_n}\left|\frac{dc_n}{dt}\right|1.$$
(8)
Then the functions $`c_n\left(t\right)`$ can be treated as temporal quasi-invariants with respect to the fast exponentials $`\mathrm{exp}\left(iE_nt/\mathrm{}\right)`$. Substituting expansion (7) into Eq. (5), we multiply the latter by the mentioned exponential and average over time according to the rule
$$\underset{T\mathrm{}}{lim}\frac{1}{T}_0^Tf(c_n,t)𝑑t,$$
where the quasi-invariants $`c_n`$ are kept fixed. Averaging in this way the equality $`\phi ^2=1`$, we have the normalization condition
$$\underset{n}{}\left|c_n\left(t\right)\right|^2=1.$$
(9)
For what follows, we need the notation for the transition amplitudes: one of them being due to atomic interactions,
$$\alpha _{mn}A_s\frac{N}{\mathrm{}}(\left|\phi _m\right|^2,\mathrm{\hspace{0.33em}2}\left|\phi _n\right|^2\left|\phi _m\right|^2),$$
(10)
and another related to the external field (6),
$$\beta _{mn}\frac{1}{\mathrm{}}(\phi _m,\widehat{B}\phi _n),$$
(11)
where $`\widehat{B}B\left(𝐫\right)`$. Note that from normalization (9) one has $`\left|c_n\right|^2=1_{m\left(n\right)}\left|c_m\right|^2`$. Using this, the nonlinear part, resulting from Eq. (5), can be represented as
$$A_s\frac{N}{\mathrm{}}\underset{m}{}\left(2\delta _{mn}\right)(\left|\phi _n\right|^2,\left|\phi _m\right|^2)=\underset{m\left(n\right)}{}\alpha _{nm}\left|c_m\right|^2c_n+\alpha _{nn}c_n.$$
Thus we come to the equation
$$i\frac{dc_n}{dt}=\underset{m\left(n\right)}{}\alpha _{nm}\left|c_m\right|^2c_n+\alpha _{nn}c_n+\frac{1}{2}\delta _{n1}\beta _{12}c_2e^{i\mathrm{\Delta }\omega t}+\frac{1}{2}\delta _{n2}\beta _{12}^{}c_1e^{i\mathrm{\Delta }\omega t}.$$
(12)
What is of physical interest is the behaviour of the fractional mode populations
$$w_n\left(t\right)\left|c_n\left(t\right)\right|^2.$$
(13)
These do not depend on the phase of $`c_n`$. Therefore, we may employ the gauge transformation
$$c_nc_n\mathrm{exp}\left(i\alpha _{nn}t\right).$$
(14)
Then Eq. (12) reduces to
$$i\frac{dc_n}{dt}=\underset{m\left(n\right)}{}\alpha _{nm}\left|c_m\right|^2c_n+\frac{1}{2}\delta _{n1}\beta _{12}c_2e^{i\mathrm{\Delta }t}+\frac{1}{2}\delta _{n2}\beta _{12}^{}c_1e^{i\mathrm{\Delta }t},$$
(15)
where
$$\mathrm{\Delta }\mathrm{\Delta }\omega +\alpha _{11}\alpha _{22}.$$
(16)
For neighboring modes, one has $`\alpha _{11}\alpha _{22}`$ and $`\mathrm{\Delta }\mathrm{\Delta }\omega `$. Otherwise, it is always possible to choose such a detuning $`\mathrm{\Delta }\omega `$ that $`\mathrm{\Delta }`$ be small, which is assumed in what follows, $`\left|\mathrm{\Delta }\right|\omega `$.
To preserve well-defined resonance, it is necessary, in analogy with the case of resonant atoms , that the transition amplitudes, involved in the process, be small. In the considered problem, such amplitudes are given by Eqs. (10) and (11). So, it is necessary that the transition amplitudes, due to atomic interactions, be small,
$$\left|\frac{\alpha _{12}}{\omega _{21}}\right|1,\left|\frac{\alpha _{21}}{\omega _{21}}\right|1,$$
(17)
as well as the amplitude determined by the modulating field,
$$\left|\frac{\beta _{12}}{\omega _{21}}\right|1.$$
(18)
From Eq. (15) it stems that if at the initial time $`c_n\left(0\right)=0`$ for $`n1,2`$, then always
$$c_n\left(t\right)=0,\left(n1,2\right)$$
(19)
for all $`t0`$. Hence Eq. (15) reduces to the system of two complex-valued equations
$$i\frac{dc_1}{dt}=\alpha _{12}\left|c_2\right|^2c_1+\frac{1}{2}\beta _{12}c_2e^{i\mathrm{\Delta }t},i\frac{dc_2}{dt}=\alpha _{21}\left|c_1\right|^2c_2+\frac{1}{2}\beta _{12}^{}c_1e^{i\mathrm{\Delta }t}.$$
(20)
Such equations, though resembling the case of two coupled electromagnetic modes , differ from that by the presence of the nonlinearity caused by atomic interactions.
The complex-valued system (20) is equivalent to four real-valued equations. However, because of the global gauge symmetry and due to the normalization condition $`\left|c_1\right|^2+\left|c_2\right|^2=1`$, the related dynamical system is, actually, two-dimensional. To show this, it is convenient to define $`c_j=\left|c_j\right|\mathrm{exp}\left(i\pi _jt\right)`$ and $`\beta _{12}\beta e^{i\gamma }`$, where $`\beta \left|\beta _{12}\right|`$. Let us also introduce the notation
$$\alpha \frac{1}{2}\left(\alpha _{12}+\alpha _{21}\right),\delta \mathrm{\Delta }+\frac{1}{2}\left(\alpha _{12}\alpha _{21}\right).$$
Then for the population difference
$$s\left|c_2\right|^2\left|c_1\right|^2$$
(21)
and effective phase difference
$$x\pi _2\pi _1+\gamma +\mathrm{\Delta },$$
(22)
we find the system of two equations
$$\frac{ds}{dt}=\beta \sqrt{1s^2}\mathrm{sin}x,\frac{dx}{dt}=\alpha s+\frac{\beta s}{\sqrt{1s^2}}\mathrm{cos}x+\delta .$$
(23)
The consideration can be generalized to the case of the multiple generation of topological coherent modes . This requires, instead of one modulating field (6), the action of several oscillating fields
$$V(𝐫,t)=\frac{1}{2}\underset{j}{}\left[B_j\left(𝐫\right)e^{i\omega _jt}+B_j^{}\left(𝐫\right)e^{i\omega _jt}\right],$$
(24)
whose frequencies are tuned to the resonance with the chosen transition frequencies. For instance, in the case of three coexisting modes, we obtain
$$i\frac{dc_1}{dt}=\left(a_{12}\left|c_2\right|^2+\alpha _{13}\left|c_3\right|^2\right)c_1+f_1,$$
$$i\frac{dc_2}{dt}=\left(a_{21}\left|c_1\right|^2+\alpha _{23}\left|c_3\right|^2\right)c_2+f_2,$$
$$i\frac{dc_3}{dt}=\left(a_{31}\left|c_1\right|^2+\alpha _{32}\left|c_2\right|^2\right)c_3+f_3,$$
(25)
where the forces $`f_j`$, related to the modulating fields, depend on the type of the mode generation, whether this is the cascade generation, $`V`$-type, or $`\mathrm{\Lambda }`$-type generation . The effective detuning $`\mathrm{\Delta }_{mn}\mathrm{\Delta }\omega _{mn}+\alpha _{nn}\alpha _{mm}`$, where $`\mathrm{\Delta }\omega _{mn}\omega _j\omega _{mn}`$, is again assumed to be small, $`\left|\mathrm{\Delta }_{mn}\right|\omega _{mn}`$.
In the same manner, one can derive the evolution equations for the amplitudes $`c_n\left(t\right)`$ for an arbitrary number of generated coherent modes, whose populations are given by $`\left|c_n\right|^2`$. The resulting equations, such as (20), (23), or (25), are nonlinear because of the binary atomic interactions. One could also include three-body interactions, which would yield the fifth-order nonlinearity with respect to $`\left|c_n\right|`$. Such three-body interactions can play an important role in describing the dissipation caused by three-body recombinations .
One could also take into account nonadiabatic corrections to the atomic evolution equations. However nonadiabatic description is of vital importance only for nonconfined motion of atoms , while for trapped atoms nonadiabatic corrections amount to at most a few percent .
## 3 Dynamic Resonant Effects
Trapped condensed Bose gas, with resonantly generated topological coherent modes, resembles a resonant atom, hence, such a resonant condensate should display the features typical of resonant finite-level atoms. But this resonant condensate is, in addition, a collective nonlinear system, because of which it can possess many other unusual properties, not existing in finite-level atoms.
(1) Interference Fringes. The total density $`\rho (𝐫,t)=\left|\phi (𝐫,t)\right|^2`$ of trapped atoms, with generated topological modes, is not simply the sum of the partial mode densities $`\rho _n(𝐫,t)=\left|c_n\left(t\right)\phi _n\left(𝐫\right)\right|^2`$, but the interference fringes arise, described by the interference density
$$\rho _{int}(𝐫,t)\rho (𝐫,t)\underset{n}{}\rho _n(𝐫,t),$$
with fast oscillation in time .
(2) Interference Current. Similarly, there exists a fastly oscillating interference current
$$𝐣_{int}(𝐫,t)𝐣(𝐫,t)\underset{n}{}𝐣_n(𝐫,t),$$
sometimes called the internal Josephson current .
(3) Mode Locking. Under this effect, the fractional mode populations $`w_n`$ are locked in the vicinity of their initial values, so that either
$$0w_n\left(t\right)\frac{1}{2}$$
for all $`t0`$, or
$$\frac{1}{2}w_n\left(t\right)1,$$
never crossing the line $`w_n=1/2`$, but being either below it or above it, depending on initial conditions .
(4) Dynamic Transition. Varying the system parameters, the dynamics of the mode populations can be qualitatively changed from the mode locked regime to the mode unlocked regime, when the mode populations fluctuate in the whole region
$$0w_n\left(t\right)1,$$
independently of their initial values .
(5) Critical Phenomena. On the manifold of the system parameters, there exists a critical surface in the vicinity of which the dynamics of the mode populations becomes unstable. Crossing the critical surface, when varying some parameters, changes the dynamics between the mode locked and mode unlocked regimes. Close to the surface, a tiny variation of some of the system parameters, say of the pumping amplitude or of the detuning, provokes drastic changes in the dynamics of the mode populations . The location of the critical surface also depends on the initial setup. Thus, for the case of two coexisting topological modes, we have the critical surface described by the relation
$$\frac{1}{2}\alpha s_0^2\beta \sqrt{1s_0^2}\mathrm{cos}x_0+\delta s_0=\beta \mathrm{sgn}\alpha ,$$
in which $`s_0s\left(0\right)`$ and $`x_0x\left(0\right)`$. For each given initial conditions, this is a surface in the three-dimensional space of the parameters $`\alpha ,\beta `$, and $`\delta `$. For $`s_0=1`$, the critical surface reads as
$$\beta \mathrm{sgn}\alpha \pm \delta =\frac{1}{2}\alpha .$$
Fixing $`s_0=1`$, we can reduce the above relation to two critical lines on the manifold of the parameters
$$b\frac{\beta }{\left|\alpha \right|},\epsilon \frac{\delta }{\left|\alpha \right|}.$$
These lines are
$$b+\epsilon =\frac{1}{2}\left(\alpha >0\right),$$
$$b\epsilon =\frac{1}{2}\left(\alpha <0\right).$$
A time-averaged system displays on the critical surface critical phenomena typical of statistical systems with the second-order phase transitions. For the averaged system it is possible to define an effective capacity and susceptibility, which diverge on the critical surface .
(6) Chaotic Motion. Fractional mode populations for a two-mode condensate are always periodic functions of time. But if the number of modes in the condensate is three or larger, then, depending on the system parameters, one has either quasiperiodic or chaotic motion. For instance, in the case of the three-mode condensate, for which $`\alpha _{ij}=\alpha `$, $`\beta _{ij}=\beta `$, and $`\mathrm{\Delta }_{ij}=0`$, chaotic motion appears when
$$\left|\frac{\beta }{\alpha }\right|0.639448,$$
that is, under a sufficiently strong pumping .
(7) Harmonic Generation. The generation of topological coherent modes may occur not solely under the direct resonance condition $`\omega =\omega _{21}`$, but also under the condition of harmonic generation
$$n\omega =\omega _{21}\left(n=1,2,\mathrm{}\right),$$
when just one modulating field, with a frequency $`\omega `$, is applied .
(8) Parametric Conversion. Another possibility of generating the topological modes is under the condition of parametric conversion
$$\underset{j}{}\left(\pm \omega _j\right)=\omega _{21},$$
when several alternating fields, with frequencies $`\omega _j`$, are involved .
(9) Atomic Squeezing. Squeezing is a quantum effect that does not exist for classical quantities. So, for treating it, one has to quantize the modes by considering the mode amplitudes $`c_n`$ as Bose operators. Then employing the pseudospin representation, one can define the squeezing factor
$$Q\frac{2\mathrm{\Delta }^2\left(S_z\right)}{\left|<S_\pm >\right|},$$
where $`\mathrm{\Delta }^2\left(S_z\right)`$ is the dispersion of the $`z`$-component of the total spin and $`S_\pm `$ are rising and lowering operators, respectively. The defined squeezing factor describes the relation between the dispersion, associated with the mode populations, and the dispersion of a relative phase. When $`Q<1`$, one says that the atomic squeezing occurs. Then the mode populations can be measured with a higher accuracy as compared to the measurement of current .
(10) Entanglement Production. In the condensate with topological coherent modes, entanglement can be produced. A general measure of entanglement production, valid for arbitrary systems, was introduced in Ref. . The measure of entanglement, generated by an operator $`\widehat{A}`$, is defined as
$$\epsilon \left(\widehat{A}\right)\mathrm{log}\frac{\widehat{A}_𝒟}{\widehat{A}^{}_𝒟},$$
where $`\left|\right|||_𝒟`$ implies the restricted norm over a set $`𝒟`$ of disentangled functions, and $`\widehat{A}^{}`$ is a nonentangling counterpart of $`\widehat{A}`$. Considering a $`p`$-particle reduced density operator $`\widehat{\rho }_p`$ for a multimode condensate, we find
$$\epsilon \left(\widehat{\rho }_p\right)=\left(1p\right)\mathrm{log}\underset{n}{sup}w_n,$$
where $`w_n=w_n\left(t\right)=\left|c_n\left(t\right)\right|^2`$. Hence the above measure $`\epsilon \left(\widehat{\rho }_p\right)\epsilon _p\left(t\right)`$ is a function of time, characterizing the entanglement evolution .
Summarizing, we have shown that by applying resonant modulating fields to a Bose-condensed gas of trapped atoms, one can create topological coherent modes. We have developed an analytical theory describing the condensate with such coherent modes and also made numerical calculations by directly solving the Gross-Pitaevskii equation . All results of the analytical theory are in very good agreement with numerical solutions. The condensate with topological coherent modes is a novel system possessing a rich variety of nontrivial properties which could be employed in many applications.
Acknowledgement
One of the authors (V.I.Y.) is grateful to the German Research Foundation for the Mercator Professorship and another author (E.P.Y.) appreciates financial support from the German Research Foundation under the DFG grant Be 142/72-1. |
warning/0506/hep-th0506178.html | ar5iv | text | # Underlying gauge symmetries of second-class constraints systems
## I Introduction
The specific feature of gauge theories is the occurrence of constraints which restrict the phase space of the system to a submanifold. A systematic study of the Hamiltonian formulation of gauge theories was made by Dirac DIRAC1 ; DIRAC2 who classified the constraints and developed the operator quantization schemes of the constraint Hamiltonian systems.
The Dirac’s theory of constraint systems was combined further with the Feynman path integral method. Faddeev FADD found an explicit measure on the constraint submanifolds, entering the path integrals. The Feynman diagram technique for non-abelian gauge theories was developed by Faddeev and Popov FAPO . The effective QCD Lagrangian involving ghost fields obeys the Becchi-Rouet-Stora (BRS) symmetry BRS . A completely general approach to quantization of gauge theories, in which the BRS transformation acquires an intrinsic meaning, is developed by Fradkin and his collaborators FRVI ; FRAD . An alternative symplectic scheme has been proposed by Faddeev and Jackiw FAJA ; JACK . An introduction to quantization of the gauge theories with the use of the path integral method can be found in Refs.HENN ; SLFA . The path integral representation for the evolution operator satisfies the unitarity condition and meets requirements of the relativistic covariance.
According to the Dirac’s classification, constraint equations like $`\mathrm{\Omega }_A=0`$, appearing in the gauge theories, are of the first class. The Poisson bracket of first-class constraint functions is a linear combination of constraint functions
$$\{\mathrm{\Omega }_A,\mathrm{\Omega }_B\}=𝒞_{AB}^D\mathrm{\Omega }_D$$
(I.1)
where $`𝒞_{AB}^D`$ are some structure functions. The quantization of gauge theories involves a set of gauge-fixing condition $`\chi _A=0`$. The functions $`\chi _A`$ must be admissible, i.e., the Poisson bracket of the gauge-fixing and constraint functions is non-degenerate:
$$det\{\chi _A,\mathrm{\Omega }_B\}0.$$
(I.2)
The gauge-fixing functions are independent
$$\{\chi _A,\chi _B\}=0.$$
(I.3)
Physical observables are gauge invariant and do not depend on the choice of $`\chi _A`$.
The wave functions satisfy the operator Dirac’s supplementary condition
$$\widehat{\mathrm{\Omega }}_B\mathrm{\Psi }=0.$$
(I.4)
Constraint functions $`𝒢_a`$ of second-class constraints systems obey the Poisson bracket relations which form a non-degenerate matrix
$$det\{𝒢_a,𝒢_b\}0.$$
(I.5)
The number of the constraints is even ($`a=1,\mathrm{},2m`$).
Among physical systems which belong to the second-class constraints systems are anomalous gauge theories ANOM ; AFAD ; FSHA1 ; FSHA2 ; JOOO , the $`O(4)`$ non-linear sigma model, which constitutes the lowest order of the chiral perturbation theory GALE , many-body systems involving collective and independent-particle degrees of freedom MBODY , and others.
The Dirac’s quantization method of such systems consists in constructing operators reproducing the Dirac bracket for canonical variables and taking constraints to be operator equations of the corresponding quantum theory.
Batalin and Fradkin BATA developed a quantization procedure for the second-class constraints systems, which converts constraints to the first class by introducing new canonical variables. The problem reduces thereby to the quantization of a first-class constraints system in an enlarged phase space. This method (see also BATY ) was found to be particularly useful for construction of the effective covariant Lagrangians in an extended configuration space BANE94 ; HONG04 .
The Hamilton-Jacobi scheme is also used for the quantization of constraint systems GUEL ; HAJA ; PIME ; BALE ; HONG04 .
Gauge systems are quantized by imposing gauge-fixing conditions. The initial system is reduced to a second-class constraints system. The evolution operator in the path representation can be written as FADD ; AFAD
$`Z`$ $`=`$ $`{\displaystyle \underset{t}{}\frac{d^nqd^np}{(2\pi \mathrm{})^n}(2\pi \mathrm{})^m\underset{a}{}\delta (𝒢_a)\sqrt{det\{𝒢_a,𝒢_b\}}}`$ (I.6)
$`\times \mathrm{exp}\left\{{\displaystyle \frac{i}{\mathrm{}}}{\displaystyle 𝑑t(p^i\dot{q}_i)}\right\}`$
where $``$ is the Hamiltonian. The path integral representation (I.6) allows not to solve the constraint equations explicitly and work in the unconstrained phase space.
The usefulness of the reduction of second-class constraints systems to equivalent gauged systems consists in getting the supplementary condition (I.4) which is not evident otherwise.
The second-class constraints systems to which gauge systems are reduced obey specific requirements:
(i) The constraint functions $`𝒢_a`$ split naturally into canonically conjugate pairs $`𝒢_a=(\chi _A,\mathrm{\Omega }_B)`$,
(ii) the wave functions satisfy Eq.(I.4),
(iii) the gauge-fixing functions $`\chi _A`$ are identically in involution Eq.(I.2), and
(iv) the constraint functions $`\mathrm{\Omega }_B`$ associated to the gauge invariance are first class Eq.(I.1).
As a consequence of Eq.(I.5), the matrix $`\{\chi _A,\mathrm{\Omega }_B\}`$ is non-degenerate. The constraint functions are not defined uniquely. In particular, a linear transformations of $`𝒢_a`$ by a non-degenerate matrix leads to constraint functions $`𝒢_a^{}`$ describing the same constraint submanifold.
A question arises if constraints of an arbitrary second-class constraints system can be redefined to fulfill (i) - (iv)?
This is the case for holonomic systems MFRF . Such systems can be treated as being obtained upon a gauge-fixing of the corresponding gauge invariant systems. Within the generalized Hamiltonian framework, constraints of holonomic systems are of the second class. In the Lagrangian formalism, the corresponding constraints do not depend on generalized velocities.
In this work, we review the underlying gauge symmetries of the holonomic systems and report new results connected to the quantization of more general second-class constraints systems.
The paper is organized as follows: In the next Section, we start from discussing a simple, but instructive example of a mechanical system under second-class holonomic constraints. The one-dimensional reduction of the $`O(n)`$ non-linear sigma model is discussed, which is equivalent to a mathematical pendulum on $`n1`$-dimensional sphere in an $`n`$-dimensional Euclidean space. Lagrangian $``$ of the system is transformed on the constraint submanifold of the configuration space into an equivalent Lagrangian $`_{}`$ to make explicit the appearance of an underlying gauge symmetry. The corresponding Hamiltonian, its constraints structure, and transformation properties are described in Sect. 3.
This example is analyzed further in Sect. 6 to construct a link with the Dirac’s quantization method.
In Sects. 4, 5, and 8, the quantization of systems under the second-class constraints is discussed.
The Poison bracket of the constraint functions $`𝒢_a`$ forms a symplectic structure in the space of constraint functions. The corresponding symplectic basis is suitable for splitting the constraints into canonically conjugate pairs ($`\chi _A`$,$`\mathrm{\Omega }_A`$). In Sect. 4, an algorithm is proposed for constructing the functions $`\chi _A`$ and $`\mathrm{\Omega }_A`$, describing the constraint submanifold of second-class constraints systems, for which the involution relations $`\{\chi _A,\chi _B\}=0`$, $`\{\mathrm{\Omega }_A,\mathrm{\Omega }_B\}=0`$, and $`\{\chi _A,\mathrm{\Omega }_B\}=\delta _{AB}`$ hold in the strong sense in an entire neighborhood of any fixed point of the constraint submanifold.
In Sect. 5, properties of the gradient projections into constraint submanifolds described by equations $`\chi _A=0`$, $`\mathrm{\Omega }_A=0`$, and onto their intersection are discussed. The gradient projection is useful for calculation of the Dirac bracket and for constructing symplectic basis for the constraint functions.
The Dirac’s quantization of the model of Sect. 2 is performed in Sect. 6 and the equivalence with the results of Sect. 3 is established. In Appendix A, the model of a mathematical pendulum on $`n1`$-dimensional sphere is embedded further into a space of higher dimension to study the constraints structure and the associated gauge invariance. The equivalence with the results of Sects. 3 and 6 is proved.
In quantum mechanics, systems are described by wave functions or, equivalently, by density matrices or the Wigner functions. The Wigner functions are the most closely related to the classical probability densities in the phase spaces. In Sect. 8, we analyze constraints to be imposed on the Wigner functions within the framework of the generalized Hamiltonian dynamics. Such constraints are not unique and depend decisively on the way the Wigner functions are extrapolated from the constraint submanifold into the unconstrained phase space. The main result of this Section consists in demonstrating the fact that constraints imposed on the Wigner functions can be taken in a symmetric fashion with regard to permutations $`\chi _A\mathrm{\Omega }_A`$, $`\chi _A\mathrm{\Omega }_A`$ and, more generally, canonical transformations preserving the Poisson brackets. Those constraints are not symmetric in the Dirac’s quantization scheme. We show that transition amplitudes are not affected by canonical transformations mixing the constraint functions.
In Sect. 9, the $`O(4)`$ non-linear sigma model is discussed as a field theory counterpart of the mathematical pendulum. We perform quantization in the straightforward way by solving the constraints from the outset and demonstrate the equivalence of the results with the method based on construction of the gauge-invariant model of the $`O(4)`$ non-linear sigma model. The Lagrange measure for the $`O(4)`$ non-linear sigma model, entering the path integral in the configuration space, is constructed and a parameterization is proposed for the pion fields, in terms of which the perturbation theory consistent with the mean field (MF) approximation can be developed. The covariance and the unitarity of the $`S`$-matrix are demonstrated.
The results are summarized in Sect.10.
## II Gauged Lagrangian for a spherical pendulum
We start from discussing a simple example of a mechanical system under holonomic constraints: A spherical pendulum on an $`n1`$-dimensional sphere $`S^{n1}`$ in $`n`$-dimensional Euclidean space. The trajectories of the system are determined as conditional extremals of the action functional $`𝒜=𝑑t`$ on the constraint submanifold $`\chi =0`$, with
$$\chi =\mathrm{ln}\varphi $$
(II.1)
where $`\varphi ^2=\varphi ^\alpha \varphi ^\alpha `$ and $`\alpha =1,\mathrm{},n`$. Lagrangian
$$=\frac{1}{2}\dot{\varphi }^\alpha \dot{\varphi }^\alpha $$
(II.2)
together with the constraint $`\chi =0`$ defines a pointlike massive particle on an $`n1`$-dimensional sphere $`S^{n1}`$ as a mechanical analogue of the $`O(n)`$ non-linear sigma model. The constraint $`\chi =0`$ implies
$`\varphi ^\alpha \dot{\varphi }^\alpha =0,`$ (II.3)
$`\dot{\varphi }^\alpha \dot{\varphi }^\alpha +\varphi ^\alpha \ddot{\varphi }^\alpha =0,`$ (II.4)
$`\mathrm{}`$
The dots stand for higher time derivatives of the constraint function $`\chi `$. Eq.(II.3) shows that the radial component of the velocity vanishes. The second equation shows that the radial component of the acceleration equals to the centrifugal force.
The equations of motion can be found using the D’Alembert-Lagrange variational principle for conditional extremals of the action functionals, or equivalently, the Euler-Lagrange variational principle for unconditional extremals with the constraints implemented through the Lagrange’s multipliers method.
Let us substitute $`_1=+\lambda \chi `$ and perform unconditional variations over the Lagrange’s multiplier $`\lambda `$ and the coordinates $`\varphi ^\alpha `$. One gets $`\chi =0`$ and $`\ddot{\varphi }^\alpha =\lambda \varphi ^\alpha /\varphi ^2`$. Multiplying further the last equation by $`\varphi ^\alpha `$ and substituting the result into Eq.(II.4), one gets $`\lambda =\dot{\varphi }^\alpha \dot{\varphi }^\alpha `$. Given that the $`\lambda `$ is fixed from Eq.(II.4), the higher time derivatives of $`\chi `$ vanish identically. The radial component of the acceleration is determined by the constraint, while the tangent component of the acceleration vanishes. The equations of motion can be presented in the form
$$\mathrm{\Delta }^{\alpha \beta }(\varphi )\frac{d^2}{dt^2}(\varphi ^\beta /\varphi )=0$$
(II.5)
where
$$\mathrm{\Delta }^{\alpha \beta }(\varphi )=\delta ^{\alpha \beta }\varphi ^\alpha \varphi ^\beta /\varphi ^2.$$
(II.6)
The tensor defined above obeys
$`\varphi ^\alpha \mathrm{\Delta }^{\alpha \beta }(\varphi )`$ $`=`$ $`0,`$ (II.7)
$`\mathrm{\Delta }^{\alpha \beta }(\varphi )\mathrm{\Delta }^{\beta \gamma }(\varphi )`$ $`=`$ $`\mathrm{\Delta }^{\alpha \gamma }(\varphi ).`$ (II.8)
It is invariant also with respect to dilatation
$$\mathrm{\Delta }^{\alpha \beta }(\varphi ^{})=\mathrm{\Delta }^{\alpha \beta }(\varphi ),$$
(II.9)
for
$$\varphi ^\alpha \varphi ^\alpha =\mathrm{exp}(\theta )\varphi ^\alpha $$
(II.10)
where $`\theta `$ is an arbitrary parameter.
Eqs.(II.5) tell that the particle moves without tangent acceleration. In general, the acceleration orthogonal to a constraint submanifold $`\chi =0`$ is fixed to keep the particle on it. In our case, the radial component of the acceleration is determined by the constraint. Eqs.(II.5) can also be derived using the D’Alembert-Lagrange variational principle.
On the submanifold $`\chi =0`$, the radial components of the velocities are equal to zero Eq.(II.3), so one can replace $`\dot{\varphi }^\alpha `$ by the tangent velocities $`\mathrm{\Delta }^{\alpha \beta }(\varphi )\dot{\varphi }^\beta `$ in $``$.
The conditional extremals of the action functional $`𝒜=𝑑t`$ do not change also if we divide $``$ by $`\varphi ^2`$ ($`=1`$). The conditional variational problem for Lagrangian
$$_{}=\frac{1}{2}\mathrm{\Delta }^{\alpha \beta }(\varphi )\dot{\varphi }^\alpha \dot{\varphi }^\beta /\varphi ^2$$
(II.11)
is thus completely equivalent to the conditional variational problem we started with.
The equations of motion (II.5) determine unconditional extremals of the action functional $`𝒜_{}=_{}𝑑t`$ on the configuration space $`=(\varphi ^\alpha )`$.
The extremals of the action functionals $`𝒜`$ and $`𝒜_{}`$ under the constraint $`\chi =0`$ coincide.
$`𝒜_{}`$ depends on the spherical coordinates $`\varphi ^\alpha /\varphi `$ which lie on an $`n1`$-dimensional sphere of a unit radius. Eqs.(II.5) for unconditional extremals of the $`𝒜_{}`$ in the coordinate space $`\varphi ^\alpha `$ coincide, as it should, with the D’Alembert-Lagrange equations for extremals of $`𝒜_{}`$ in the space formed by the spherical coordinates under the condition that they belong to an $`n1`$-dimensional sphere of a unit radius.
It is seen that Lagrangian (II.11) is invariant with respect to dilatation (II.10) where $`\theta `$ is an arbitrary function of time. In the context of the dynamics defined by (II.11) with no constraints imposed, $`\varphi `$ turns out to be an arbitrary function of time. It can always be selected to fulfill the constraint $`\chi =0`$ or some other admissible constraint.
The constraint $`\chi =0`$ can therefore be treated as a gauge-fixing condition, the function $`\varphi `$ as a gauge degree of freedom, the ratios $`\varphi ^\alpha /\varphi `$ as gauge invariant observables.
The equations of motion (II.5) are formulated in terms of the gauge invariant observables.
The above gauge symmetry is defined ”on-shell”, i.e., for $`\dot{\varphi }^\alpha `$ treated as a time derivative of $`\varphi ^\alpha `$. In the tangent bundle $`T=(\varphi ^\alpha ,\dot{\varphi }^\alpha )`$ where the coordinates and their derivatives are independent, $`_{}`$ is invariant with respect to a two-parameter set of transformations:
$`\varphi ^\alpha `$ $``$ $`\varphi ^\alpha =\mathrm{exp}(\theta )\varphi ^\alpha ,`$ (II.12)
$`\dot{\varphi }^\alpha `$ $``$ $`\dot{\varphi }^\alpha =\mathrm{exp}(\theta )\dot{\varphi }^\alpha `$ (II.13)
and
$`\varphi ^\alpha `$ $``$ $`\varphi ^\alpha =\varphi ^\alpha ,`$ (II.14)
$`\dot{\varphi }^\alpha `$ $``$ $`\dot{\varphi }^\alpha =\dot{\varphi }^\alpha +ϵ\varphi ^\alpha .`$ (II.15)
The last two equations describe the invariance under variation of the radial component of the particle velocities. If the $`\dot{\varphi }^\alpha `$ is treated on-shell, then $`ϵ\dot{\theta }`$.
A pointlike particle on an $`n1`$-dimensional sphere $`S^{n1}`$ has therefore underlying gauge symmetry connected with dilatation of the coordinates $`\varphi ^\alpha `$. Its physical origin is simple: We allow virtual displacements from the constraint submanifold and treat such displacements as unphysical gauge degrees of freedom. The constraint equation $`\chi =0`$ turns thereby into a gauge-fixing condition. The physical variables are specified by projections of the coordinates $`\varphi ^\alpha `$ onto an $`n1`$-dimensional sphere of a unit radius. Those projections are the spherical coordinates $`\varphi ^\alpha /\varphi `$.
Let us consider the system (II.11) within the generalized Hamiltonian framework.
## III The gauged spherical pendulum Hamiltonian
Lagrangian (II.11) is defined outside of the constraint submanifold $`\chi =0`$ and invariant with respect to the dilatation (II.10). The value $`\varphi `$ is an arbitrary function and becomes a gauge degree of freedom. The equations of motion (II.5) are derived using the constraint $`\chi =0`$. Being formulated, the equations of motion do not depend, however, on the constraint anymore and allow an extension to the unconstrained configuration space $``$. The invariance under the dilatation, Eq.(II.10), is a consequence of the two-parameter set of global symmetry transformations, Eqs.(II.12) - (II.15), of $`_{}`$ as a function defined on the tangent bundle $`T`$.
One can consider therefore $`𝒜_{}`$ without imposing any constraints and treat the equation $`\chi =0`$ as a gauge-fixing condition.
### III.1 Gauged spherical pendulum within generalized Hamiltonian framework
The canonical momenta corresponding to $`\dot{\varphi }^\alpha `$ are defined by
$$\pi ^\alpha =\frac{_{}}{\dot{\varphi }^\alpha }=\mathrm{\Delta }^{\alpha \beta }(\varphi )\dot{\varphi }^\beta /\varphi ^2.$$
(III.1)
They satisfy constraints
$$\pi ^\alpha \mathrm{\Delta }^{\alpha \beta }(\varphi )\pi ^\beta 0$$
(III.2)
which are equivalent to the one primary constraint:
$$\mathrm{\Omega }=\varphi \pi 0.$$
(III.3)
The primary Hamiltonian can be obtained with the use of the Legendre transform:
$$=\frac{1}{2}\varphi ^2\mathrm{\Delta }^{\alpha \beta }(\varphi )\pi ^\alpha \pi ^\beta .$$
(III.4)
For $`n=3`$, $``$ is proportional to the orbital momentum squared. The non-vanishing Poisson bracket relations for the canonical variables are defined by
$$\{\varphi ^\alpha ,\pi ^\beta \}=\delta ^{\alpha \beta }.$$
(III.5)
The constraint function $`\mathrm{\Omega }`$ is stable with respect to the time evolution:
$$\{\mathrm{\Omega },\}=0.$$
(III.6)
The relations
$`\{\varphi ^\alpha ,\mathrm{\Omega }\}`$ $`=`$ $`\varphi ^\alpha ,`$ (III.7)
$`\{\pi ^\alpha ,\mathrm{\Omega }\}`$ $`=`$ $`\pi ^\alpha `$ (III.8)
show that $`\mathrm{\Omega }`$ generates dilatation of $`\varphi ^\alpha `$ and $`\pi ^\alpha `$. The transformation law for the canonical coordinates Eq.(III.7) is in agreement with Eq.(II.12) in its infinitesimal form. The transformation law for the canonical momenta, which follows from Eqs.(II.12), (II.13), and (III.1),
$$\pi ^\alpha \pi ^\alpha =\mathrm{exp}(\theta )\pi ^\alpha ,$$
(III.9)
considered in its infinitesimal form, is in agreement with Eq.(III.8) either. The Hamiltonian $``$ is gauge invariant under the dilatation.
The roles of the gauge-fixing function $`\chi `$ and the gauge generator $`\mathrm{\Omega }`$ are similar. The function $`\chi `$ is identically in involution with the Hamiltonian:
$$\{\chi ,\}=0.$$
(III.10)
The Poisson bracket relations
$`\{\varphi ^\alpha ,\chi \}`$ $`=`$ $`0,`$ (III.11)
$`\{\pi ^\alpha ,\chi \}`$ $`=`$ $`\varphi ^\alpha `$ (III.12)
define the one-parameter set of transformations with respect to which $``$ is invariant. The function $`\chi `$ generates shifts of the longitudinal component of the canonical momenta. This symmetry is connected to the invariance of $`_{}`$ described by Eqs.(II.14) and (II.15) .
The gauge-fixing condition $`\chi =0`$ is admissible:
$$\{\chi ,\mathrm{\Omega }\}=1.$$
(III.13)
The equations of motion generated by the primary Hamiltonian look like
$`\dot{\varphi }^\alpha =\{\varphi ^\alpha ,\}`$ $`=`$ $`\varphi ^2\mathrm{\Delta }^{\alpha \beta }(\varphi )\pi ^\beta ,`$ (III.14)
$`\dot{\pi }^\alpha =\{\pi ^\alpha ,\}`$ $`=`$ $`\varphi ^\alpha \mathrm{\Delta }^{\beta \gamma }(\varphi )\pi ^\beta \pi ^\gamma .`$ (III.15)
### III.2 Quantization of spherical pendulum
The quantization of a mathematical pendulum on an $`n1`$-dimensional sphere is discussed in Ref.HONG04 where the standard Batalin-Fradkin-Tyutin (BFT) algorithm BATA ; BATY for second-class constraints systems is applied. The second-class constraints appear if one starts directly from $``$ and formulate the conditional variational problem for $`\chi =0`$. By constructing auxiliary fields, it is possible to pass over to an equivalent first-class constraint system.
In our approach, we start from $`𝒜_{}`$ which is gauge invariant explicitly, so the constraints appear to be of the first class from the start. One can therefore quantize the pendulum as a gauge-invariant system without introducing auxiliary fields.
From the point of view of the generalized Hamiltonian framework, the gauge-fixing conditions and the gauge generators play similar roles. The function $`\chi `$ does, however, not generate transformations in $`T`$, so it appears just as a candidate for gauge-fixing function. If we pass to the Lagrangian framework, we can verify that $`\chi =0`$ is the gauge-fixing condition indeed.
The system has the gauge invariance described by the generator $`\mathrm{\Omega }`$ and the admissible gauge-fixing condition $`\chi =0`$.
The standard procedure for gauge theories can therefore be applied for quantization of the spherical pendulum.
The system is quantized by the algebra mapping $`(\varphi ^\alpha ,\pi ^\alpha )(\widehat{\varphi }^\alpha ,\widehat{\pi }^\alpha )`$ and $`\{,\}i/\mathrm{}[,]`$. Consequently, to any symmetrized function in the phase space variables one may associate an operator function. The function is symmetrized in such a way that quantal image is a hermitian operator.
The quantum hermitian Hamiltonian has the form
$$\widehat{}=\frac{1}{2}\varphi \mathrm{\Delta }^{\alpha \beta }\widehat{\pi }^\beta \varphi \mathrm{\Delta }^{\alpha \gamma }\widehat{\pi }^\gamma .$$
(III.16)
The vector $`i\varphi \mathrm{\Delta }^{\alpha \beta }\widehat{\pi }^\beta `$ gives the angular part of the gradient operator. Although it is not conspicuous, $`\widehat{}`$ does not depend on the radial coordinate $`\varphi `$. The constraint operator can be defined as
$$\widehat{\mathrm{\Omega }}=(\varphi ^\alpha \widehat{\pi }^\alpha +\widehat{\pi }^\alpha \varphi ^\alpha )/2.$$
(III.17)
It acts only on the radial component of $`\varphi ^\alpha `$, so the relation
$$[\widehat{\mathrm{\Omega }},\widehat{}]=0$$
(III.18)
holds.
The physical subspace of the Hilbert space is singled out by imposing the Dirac’s supplementary condition
$$\widehat{\mathrm{\Omega }}\mathrm{\Psi }=0.$$
(III.19)
This condition implies
$$\mathrm{\Psi }=\varphi ^{n/2}\mathrm{\Psi }_1(\varphi ^\alpha /\varphi ).$$
(III.20)
The physical information is contained in $`\mathrm{\Psi }_1(\varphi ^\alpha /\varphi )`$.
The path integral for the evolution operator becomes
$$Z=\underset{t}{}\frac{d^n\varphi d^n\pi }{(2\pi \mathrm{})^{n1}}\delta (\chi )\delta (\mathrm{\Omega })\mathrm{exp}\left\{\frac{i}{\mathrm{}}𝑑t(\pi ^\alpha \dot{\varphi }^\alpha )\right\}.$$
(III.21)
Eqs.(III.19) and (III.21) solve the quantization problem for a mathematical pendulum on an $`n1`$-dimensional sphere $`S^{n1}`$.
It is desirable that the quantization procedure does not destroy the classical symmetries which results in having the supplementary condition (III.19) satisfied by the state $`\mathrm{\Psi }(t)`$ for any value of $`t`$. This feature can, in general, be violated either due to complex terms entering the Hamiltonian or by approximations adopted for treating the operator eigenvalues.
Given the initial state wave function $`\mathrm{\Psi }(0)`$ satisfying (III.19), the final wave function $`\mathrm{\Psi }(t)`$ can be found applying the evolution operator (III.21) on $`\mathrm{\Psi }(0)`$. Since $`\widehat{\mathrm{\Omega }}`$ commute with the Hamiltonian, the final state wave function obeys Eqs.(III.19).
In the standard canonical frame MASKAWA , the first canonical coordinates and momenta are identical with the constraint functions. If the Poisson bracket of the Hamiltonian with the constraint functions vanishes, the Hamiltonian is independent of the first canonical coordinates and momenta. It means that the quantized Hamiltonian is commutative with operators associated to the constraint functions.
By initiating the quantization procedure in the standard canonical frame, one gets the classical symmetries preserved and the Dirac’s supplementary condition fulfilled at any time.
The fact that the dilatation symmetry for the spherical pendulum is preserved on the quantum level is connected with the fact that the canonical frame we used is simply related to the standard canonical frame.
The integral over the canonical momenta Eq.(III.21) can be simplified to give
$$Z=\underset{t}{}\sqrt{\left(\chi /\varphi ^\alpha \right)^2}\delta (\chi )d^n\varphi \mathrm{exp}\left\{\frac{i}{\mathrm{}}𝑑t_{}\right\}.$$
(III.22)
Lagrange’s measure $`\sqrt{\left(\chi /\varphi ^\alpha \right)^2}\delta (\chi )d^n\varphi `$ coincides with the volume element of $`S^{n1}`$ sphere. It can be rewritten as an invariant volume of the configuration space, e.g., in terms of the angular variables $`(\phi _1\mathrm{}\phi _{n1})`$ with the help of the induced metric tensor. It is invariant under $`O(n)`$ rotations and remains the same for all functions $`\chi `$ vanishing at $`\varphi =1`$. For $`n=4`$, Lagrange’s measure matches Haar’s measure of the group $`SU(2)`$.
In gauge theories, the evolution operator is independent of the gauge-fixing conditions FADD . We can insert $`det\{\chi ,\mathrm{\Omega }\}=1`$ into the integrand of Eq.(III.22) to bring the measure into the form identical with gauge theories. In the path integral, $`\chi `$ can then be replaced with an arbitrary function. The condition (III.19) does not comprise the constraint $`\chi =0`$ also.
It is remarkable that physical observables depend on half of the number of the second-class constraints.
The main question to be raised here if the quantization method described above is general enough or specific only for the spherical pendulum? In Sect. 7 we show that such a method works for mechanical systems and field theories under holonomic constraints. The quantization of more general second-class constraints systems is discussed in Sect. 8.
## IV Local symplectic basis for second-class constraints functions
Two sets of the constraint functions are equivalent if they describe the same constraint submanifold. The Hamiltonian function admits transformations which do not change its value and its first derivatives on the constraint submanifold. This allows to make transformations of the constraint and Hamiltonian functions without changing the physical content of theory.
The second-class constraints satisfy $`det\{𝒢_a,𝒢_b\}0`$. The Poisson bracket defines therefore a non-degenerate symplectic linear structure in the vector space of the constraint functions $`𝒢_a`$. Indeed, any linear transformation $`𝒢_a^{}=𝒪_{ab}𝒢_b`$ with matrix $`𝒪_{ab}`$ depending on the canonical variables transforms accordingly the Poisson bracket: $`\{𝒢_a^{},𝒢_b^{}\}𝒪_{ac}𝒪_{bd}\{𝒢_c,𝒢_d\}`$. The Poisson bracket plays thereby a role of a skew-scalar product in the symplectic vector space of the constraint functions. Every symplectic space has a symplectic basis (see, e.g., ARNO ), so the constraint functions can be brought by linear transformations into the form
$$\{𝒢_a,𝒢_b\}_{ab}$$
(IV.1)
where
$$_{ab}=\begin{array}{cc}0\hfill & E_m\hfill \\ E_m\hfill & 0\hfill \end{array},$$
(IV.2)
with $`E_m`$ being the $`m\times m`$ identity matrix, $`_{ab}_{bc}=\delta _{ac}`$. Using representation $`𝒢_a=(\chi _A,\mathrm{\Omega }_A)`$, one has
$`\{\mathrm{\Omega }_A,\mathrm{\Omega }_B\}`$ $``$ $`0,`$ (IV.3)
$`\{\chi _A,\chi _B\}`$ $``$ $`0,`$ (IV.4)
$`\{\chi _A,\mathrm{\Omega }_B\}`$ $``$ $`\delta _{AB}.`$ (IV.5)
This basis is not unique. Indeed, there remains a group of symplectic transformations $`Sp(2m)`$, which keeps the Poison bracket of the constraint functions in the symplectic form.
At any given point of a neighborhood of the constraint submanifold one can find symplectic basis for second-class constraints functions in the weak form.
This result can be strengthened as shown below:
($`𝒜`$) If $`\xi `$ is close enough to the constrain submanifold, one has $`det\{\chi _A(\xi )`$,$`\mathrm{\Omega }_B(\xi )\}0`$. By continuity there exists a finite neighborhood $`\mathrm{\Delta }_\xi `$ of $`\xi `$, such that $`det\{\chi _A(\xi ^{})`$,$`\mathrm{\Omega }_B(\xi ^{})\}0`$ $`\xi ^{}\mathrm{\Delta }_\xi `$. Let us assume that the intersection of $`\mathrm{\Delta }_\xi `$ with the constraint submanifold is not empty, i.e., one can find $`\xi _1\mathrm{\Delta }_\xi `$ such that $`𝒢_a(\xi _1)=0`$. If it is not fulfilled, we start from another $`\xi `$ closer to the constraint submanifold.
Let us chose an equivalent set of the constraint functions $`\chi _A\chi _A^{}=\{\chi _A,\mathrm{\Omega }_B\}^1\chi _B`$ to ensure $`\{\chi _A^{},\mathrm{\Omega }_B\}\delta _{AB}`$ in the region $`\mathrm{\Delta }_\xi `$.
($``$) We replace further $`\mathrm{\Omega }_A\mathrm{\Omega }_A^{}=\mathrm{\Omega }_A\frac{1}{2}C_{AB}\chi _B^{}`$ to get equations
$$\{\mathrm{\Omega }_A^{},\mathrm{\Omega }_B^{}\}\{\mathrm{\Omega }_A,\mathrm{\Omega }_B\}C_{AB}+\frac{1}{4}C_{AD}C_{BF}\{\chi _D^{},\chi _F^{}\}0.$$
(IV.6)
These equations can be solved for matrix $`C_{AB}`$ in terms of a power series of the matrix $`\{\chi _D^{},\chi _F^{}\}`$:
$$C_{AB}=\underset{k=0}{\overset{\mathrm{}}{}}C_{AB}^{[k]}$$
(IV.7)
where $`C_{AB}^{[k]}=C_{BA}^{[k]}`$ and
$`C_{AB}^{[0]}`$ $`=`$ $`\{\mathrm{\Omega }_A,\mathrm{\Omega }_B\},`$
$`C_{AB}^{[k+1]}`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \underset{j=0}{\overset{k}{}}}C_{AD}^{[kj]}C_{BF}^{[j]}\{\chi _D^{},\chi _F^{}\},`$
with $`k=0,1,\mathrm{}`$ .
($`𝒞`$) The transform $`\chi _A^{}\chi _A^{\prime \prime }=\{\chi _A^{},\mathrm{\Omega }_B^{}\}^1\chi _B^{}`$ brings back the Poisson bracket $`\{\chi _A^{\prime \prime },\mathrm{\Omega }_B^{}\}`$ to the diagonal form $`\delta _{AB}.`$
($`𝒟`$) The last transform looks like $`\chi _A^{\prime \prime }\chi _A^{\prime \prime \prime }=\chi _A^{\prime \prime }\frac{1}{2}\{\chi _A^{\prime \prime },\chi _B^{\prime \prime }\}\mathrm{\Omega }_B^{}.`$ It provides $`\{\chi _A^{\prime \prime \prime },\chi _B^{\prime \prime \prime }\}0`$and keeps the relation $`\{\chi _A^{\prime \prime \prime },\mathrm{\Omega }_B^{}\}\delta _{AB}`$ unchanged.
Now, we remove primes from the notations. As a result of the steps ($`𝒜`$)-($`𝒟`$), we obtain weak equations
$$\{𝒢_a(\zeta ),𝒢_b(\zeta )\}_{ab}$$
(IV.8)
$`\zeta \mathrm{\Delta }_\xi `$. It is manifest that Eq.(IV.8) is valid in some neighborhood of any point $`\xi _1`$ of the constraint submanifold too.
The symplectic basis for second-class constraints functions in the weak form exists in an entire neighborhood of any given point of the constraint submanifold.
The existence of the local symplectic basis in the weak form is on the line with the Darboux’s theorem (see, e.g., RABJAR ) which states that around every point $`\xi `$ in a symplectic space there exists a coordinate system in $`\mathrm{\Delta }_\xi `$ such that $`\xi \mathrm{\Delta }_\xi `$ where the symplectic structure takes the standard canonical form. The symplectic space can be covered by such coordinate systems.
This is in sharp contrast to the situation in Riemannian geometry where the metric at any given point $`x`$ can always be made Minkowskian, but in any neighborhood of $`x`$ the variance of the Riemannian metric with the Minkowskian metric is, in general, $`\mathrm{\Delta }x^2`$. In other words, by passing to an inertial coordinate frame one can remove gravitation fields at any given point, but not in an entire neighborhood of that point. The Darboux’s theorem states, reversely, that the symplectic structure can be made to take the standard canonical form in an entire neighborhood $`\mathrm{\Delta }_\xi `$ of any point $`\xi \mathrm{\Delta }_\xi `$. In Riemannian spaces, locally means at some given point. In symplectic spaces, locally means at some given point and in an entire neighborhood of that point.
Locally, all symplectic spaces are indistinguishable. Any submanifold in a symplectic space, including any constraint submanifold, is a plane. The possibility of finding the standard canonical frame MASKAWA illustrates this circumstance.
In the view of this marked dissimilarity, the validity of Eqs.(IV.8) in a finite domain looks indispensable.
The global symplectic basis exists apparently for $`m=1`$ as, e.g., in the case of the spherical pendulum. The global existence of the basis (IV.8) has been proved for systems with one primary constraint MITRA1 , and also, assuming that $`det\{\chi _A,\mathrm{\Omega }_B\}0`$ holds globally VYTHE1 .
Admissible transformations on the second-class constraints functions allow to bring Eqs.(IV.8) into a strong form. The Hamiltonian can also be modified to convert the Poisson bracket relations with the second-class constraints functions into the strong form without changing the dynamics. The arguments given below follow closely our discussion of holonomic systems MFRF :
According to Eqs.(IV.8), at any given point of the constraint submanifold where $`det\{𝒢_a,𝒢_b\}0`$ one can select symplectic basis in which the constraint functions satisfy
$$\{𝒢_a,𝒢_b\}=_{ab}+𝒞_{ab}^c𝒢_c$$
(IV.9)
in an entire neighborhood of that point. The first-class Hamiltonian $``$ has relations
$$\{𝒢_a,\}=_a^b𝒢_b.$$
(IV.10)
The Jacobi identities for $`𝒢_a`$, $`𝒢_b`$, and $`𝒢_c`$ and for $`𝒢_a`$, $`𝒢_b`$, and $``$ imply
$`𝒞_{abc}+𝒞_{bca}+𝒞_{cab}`$ $``$ $`0,`$ (IV.11)
$`_{ab}_{ba}`$ $``$ $`0`$ (IV.12)
where $`𝒞_{abc}=_{cd}𝒞_{ab}^d`$ and $`_{ab}=_{ac}_a^c`$.
Let us define
$`𝒢_a^{}`$ $`=`$ $`𝒢_a+{\displaystyle \frac{1}{3}}𝒞_a^{bc}𝒢_b𝒢_c,`$ (IV.13)
$`^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}^{ab}𝒢_a𝒢_b.`$ (IV.14)
The constraint functions $`𝒢_a`$ and $`𝒢_a^{}`$ coincide on and in vicinity of the constraint submanifold up to the second order in $`𝒢_a`$ and generate accordingly identical phase flows on the submanifold $`𝒢_a=0`$. The Hamiltonian functions $``$ and $`^{}`$ coincide on and in vicinity of the constraint submanifold up to the second order in $`𝒢_a`$ and generate on the submanifold $`𝒢_a=0`$ identical Hamiltonian phase flows.
Using Eqs.(IV.11) and (IV.12) one gets
$`\{𝒢_a^{},𝒢_b^{}\}`$ $`=`$ $`_{ab}+𝒞_{ab}^{cd}𝒢_c^{}𝒢_d^{},`$ (IV.15)
$`\{𝒢_a^{},^{}\}`$ $`=`$ $`_a^{bc}𝒢_b^{}𝒢_c^{}`$ (IV.16)
where $`𝒞_{ab}^{cd}`$ and $`_a^{bc}`$ are new structure functions. The first-order terms in $`𝒢_a`$ do not appear in the right sides of these equations. The pair ($`𝒢_a^{}`$,$`^{}`$) describes the same Hamiltonian dynamics as ($`𝒢_a`$,$``$), being at the same time in a stronger involution around the constraint submanifold.
The above procedure can be repeated to remove the quadratic terms in $`𝒢_a`$, cubic terms in $`𝒢_a`$, etc. In general, assuming
$`\{𝒢_a^{[k]},𝒢_b^{[k]}\}`$ $`=`$ $`_{ab}+𝒞_{ab}^{c_1\mathrm{}c_k}𝒢_{c_1}^{[k]}\mathrm{}𝒢_{c_k}^{[k]},`$ (IV.17)
$`\{𝒢_a^{[k]},^{[k]}\}`$ $`=`$ $`_a^{b_1\mathrm{}b_k}𝒢_{b_1}^{[k]}\mathrm{}𝒢_{b_k}^{[k]},`$ (IV.18)
we get as a consequence of the Jacobi identities and of the symmetry of structure functions $`𝒞_{ab}^{c_1\mathrm{}c_k}`$ and $`_a^{b_1\mathrm{}b_k}`$ with respect to the upper indices:
$`𝒞_{abc_1\mathrm{}c_k}+𝒞_{bc_1a\mathrm{}c_k}+𝒞_{c_1ab\mathrm{}c_k}`$ $``$ $`0,`$ (IV.19)
$`_{ab_1\mathrm{}b_k}_{b_1a\mathrm{}b_k}`$ $``$ $`0.`$ (IV.20)
The next-order constraint functions and the Hamiltonian are given by
$`𝒢_a^{[k+1]}`$ $`=`$ $`𝒢_a^{[k]}+{\displaystyle \frac{1}{k+2}}𝒞_a^{c_0c_1\mathrm{}c_k}𝒢_{c_0}^{[k]}𝒢_{c_1}^{[k]}\mathrm{}𝒢_{c_k}^{[k]},`$ (IV.21)
$`^{[k+1]}`$ $`=`$ $`^{[k]}{\displaystyle \frac{1}{k+1}}^{b_0b_1\mathrm{}b_k}𝒢_{b_0}^{[k]}𝒢_{b_1}^{[k]}\mathrm{}𝒢_{b_k}^{[k]}.`$ (IV.22)
Using Eqs.(IV.17) and (IV.18), one can calculate the next-order structure functions $`𝒞_{ab}^{c_1\mathrm{}c_{k+1}}`$ and $`_a^{b_1\mathrm{}b_{k+1}}`$ and repeat the procedure. If structure functions vanish, we shift $`k`$ by one unit and check Eqs.(IV.17) and (IV.18) again. At each step, the Poisson bracket relations get closer to the normal form. In the limit $`k+\mathrm{}`$, we obtain
$`\{\stackrel{~}{𝒢}_a,\stackrel{~}{𝒢}_b\}`$ $`=`$ $`_{ab},`$ (IV.23)
$`\{\stackrel{~}{𝒢}_a,\stackrel{~}{}\}`$ $`=`$ $`0`$ (IV.24)
where $`\stackrel{~}{𝒢}_a=lim_k\mathrm{}𝒢_a^{[k]}`$ and $`\stackrel{~}{}=lim_k\mathrm{}^{[k]}`$. The matrix $`_{ab}`$ defines splitting of the constraints into two groups ($`\stackrel{~}{\chi }_A`$,$`\stackrel{~}{\mathrm{\Omega }}_A`$), such that
$`\{\stackrel{~}{\chi }_A,\stackrel{~}{\chi }_B\}`$ $`=`$ $`0,`$ (IV.25)
$`\{\stackrel{~}{\mathrm{\Omega }}_A,\stackrel{~}{\mathrm{\Omega }}_B\}`$ $`=`$ $`0,`$ (IV.26)
$`\{\stackrel{~}{\chi }_A,\stackrel{~}{\mathrm{\Omega }}_B\}`$ $`=`$ $`\delta _{AB}.`$ (IV.27)
The progress made consists in extending the validity of Eqs.(IV.1) from one point into its neighborhood Eqs.(IV.8), and further, in passing from the weak to the strong form of the Poisson bracket relations Eqs.(IV.23).
The symplectic basis for second-class constraints functions exists in the strong form in an entire neighborhood of any given point of the constraint submanifold. There exits a Hamiltonian function $`\stackrel{~}{}`$, describing the same dynamics on the constraint submanifold as the initial Hamiltonian function $``$, which is identically in involution with the constraint functions.
An independent geometric-based construction of $`\stackrel{~}{𝒢}_a`$ and $`\stackrel{~}{}`$ is given in Sect. 5.
A similar local basis exists for first-class constraints HENN ; SCHO49 . The arguments of Refs.HENN ; SCHO49 apply to second-class constraints under special restrictions which are discussed in Appendix B.
For systems of pointlike particles under holonomic constrains, Eqs.(IV.9) hold globally, so $`\stackrel{~}{𝒢}_a`$ and $`\stackrel{~}{}`$ exist globally also. Furthermore, $`\stackrel{~}{\mathrm{\Omega }}_A`$ is linear in the canonical momenta, $`\stackrel{~}{\chi }_B`$ does not depend on the canonical momenta, and $`\stackrel{~}{}`$ splits into a sum of a kinetic energy term quadratic in the canonical momenta and a potential energy term depending on the canonical coordinates MFRF .
## V Gradient projection
The concept of the gradient projection is useful for applications. It defines functions $`\xi _s(\xi )`$ which project an arbitrary point $`\xi `$ of the phase space onto a submanifold $`𝒢_a=0`$ ($`a=1,\mathrm{},2m`$) of the phase space along phase flows generated by the constraint functions $`𝒢_a`$.
### V.1 Full gradient projection
Let $`𝒢_a=0`$ be second-class constraints, $`det\{𝒢_a,𝒢_b\}0`$, $`\xi ^i`$ ($`i=1,\mathrm{},2n`$) are canonical variables. In vicinity of the submanifold $`𝒢_a=0`$, the projections can be constructed explicitly. Near the constraint submanifold, one can write $`\xi _s(\xi )=\xi +\{\xi ,𝒢_a\}\lambda _a`$. The small parameters $`\lambda _a`$ are determined by requiring $`𝒢_a(\xi _s(\xi ))=0`$ to the first order in $`𝒢_a`$. We get
$$\xi _s(\xi )=\xi \{\xi ,𝒢_a\}\{𝒢_a,𝒢_b\}^1𝒢_b.$$
(V.1)
It is seen that $`\{𝒢_a,\xi _s(\xi )\}=0`$, consequently $`\{𝒢_a,f(\xi _s(\xi ))\}=0`$ for any function $`f`$. This is natural, since $`𝒢_a`$ generate phase flows along which the projections $`\xi _s(\xi )`$ have been constructed.
The reciprocal statement is also true: If $`\{𝒢_a,f\}=0`$ for all $`𝒢_a`$, then $`f=f(\xi _s(\xi ))`$. Indeed, the coordinates on the constraint submanifold can be parameterized by $`\xi _s`$. The coordinates describing shifts from the constraint submanifold can be parameterized by $`𝒢_a`$. The functions $`f`$ can in general be written as $`f=f(\xi _s,𝒢_a)`$. If the Poisson brackets of $`f`$ with all $`𝒢_a`$ vanish, $`f`$ depends on $`\xi _s`$ only.
This can be summarized by
$$\{𝒢_a,f\}=0f=f(\xi _s(\xi )).$$
(V.2)
Beyond the lowest order in $`𝒢_a`$, the operation is unique provided the phase flows commute. This is always the case for the constraint functions taken to accomplish (IV.23) in a finite neighborhood of the constraint submanifold.
### V.2 Partial gradient projection
Let $`𝒢_a`$ split into $`\chi _A`$ and $`\mathrm{\Omega }_A`$. We wish to construct functions $`\xi _u(\xi )`$ which project an arbitrary point $`\xi `$ of the phase space onto the gauge-fixing surface $`\chi _A=0`$ with the use of the constraint functions $`\mathrm{\Omega }_A`$ associated to gauge transformations. In vicinity of the submanifold, one can write $`\xi _u(\xi )=\xi +\{\xi ,\mathrm{\Omega }_A\}\lambda _A`$. To the first order in $`\chi _A`$, the parameters $`\lambda _A`$ can be found from equation $`\chi _A(\xi _u(\xi ))=0`$:
$$\xi _u(\xi )=\xi \{\xi ,\mathrm{\Omega }_A\}\{\chi _A,\mathrm{\Omega }_B\}^1\chi _B.$$
(V.3)
The projection is made along the phase flows of $`\mathrm{\Omega }_A`$, so $`\{\mathrm{\Omega }_A,f(\xi _u(\xi ))\}=0`$ for any function $`f`$.
The reverse statement is also true. We write $`f=f(\xi _u,\chi _A)`$ and conclude from $`\{\mathrm{\Omega }_A,f\}=0`$ that the dependence on $`\chi _A`$ drops out.
It can be summarized as follows:
$$\{\mathrm{\Omega }_A,f\}=0f=f(\xi _u(\xi )).$$
(V.4)
The second partial gradient projection can be made onto the submanifold $`\mathrm{\Omega }_A=0`$ with the use of the constraint functions $`\chi _A`$. The result is similar to Eq.(V.3)
$$\xi _v(\xi )=\xi +\{\xi ,\chi _A\}\{\chi _A,\mathrm{\Omega }_B\}^1\mathrm{\Omega }_B.$$
(V.5)
The relation $`\{\chi _A,f(\xi _v(\xi ))\}=0`$ is valid for any function $`f`$. Furthermore,
$$\{\chi _A,f\}=0f=f(\xi _v(\xi )).$$
(V.6)
Combining the partial projections, e.g., $`\xi _s(\xi )=\xi _v(\xi _u(\xi ))`$, one gets the full gradient projection. To the first order in $`𝒢_a`$, the order in which the partial projections are applied does not matter, so $`\xi _v(\xi _u(\xi ))=\xi _u(\xi _v(\xi ))`$. The full gradient projection constructed in this way coincides with that given by Eq.(V.1).
### V.3 Example: Gauged spherical pendulum
The constraint function $`\mathrm{\Omega }`$ can be used to bring the vector $`\varphi ^\alpha `$ onto a sphere of a unit radius $`\varphi =1`$ ($`\chi =0`$). Such a transformation has a meaning of a gradient projection. The functions $`\varphi _u^\alpha `$ and $`\pi _u^\alpha `$ can be constructed in terms of the variables $`\varphi ^\alpha `$ and $`\pi _\alpha `$:
$`\varphi _u^\alpha `$ $`=`$ $`\mathrm{exp}(\theta )\varphi ^\alpha ,`$ (V.7)
$`\pi _u^\alpha `$ $`=`$ $`\mathrm{exp}(\theta )\pi ^\alpha .`$ (V.8)
The condition $`\chi (\varphi _u,\pi _u)=0`$ gives $`\mathrm{exp}(\theta )=1/\varphi `$.
The Poisson bracket relations for the projected variables (V.7) and (V.8) can be found to be
$`\{\varphi _u^\alpha ,\varphi _u^\beta \}`$ $`=`$ $`0,`$ (V.9)
$`\{\varphi _u^\alpha ,\pi _u^\beta \}`$ $`=`$ $`\mathrm{\Delta }_u^{\alpha \beta }(\varphi _u),`$ (V.10)
$`\{\pi _u^\alpha ,\pi _u^\beta \}`$ $`=`$ $`\varphi _u^\beta \pi _u^\alpha \varphi _u^\alpha \pi _u^\beta `$ (V.11)
where $`\mathrm{\Delta }_u^{\alpha \beta }(\varphi _u)=\delta ^{\alpha \beta }\varphi _u^\alpha \varphi _u^\beta `$.
The following properties are worthy of mention:
(i) The Poisson bracket of the projected canonical variables coincides with the Dirac bracket associated to the constraints $`\chi =0`$ and $`\mathrm{\Omega }=0`$ on the submanifold $`\chi =0`$.
The Dirac bracket for the canonical variables can be calculated using Eq.(V.4) to give
$`\{\varphi ^\alpha ,\varphi ^\beta \}_D`$ $`=`$ $`0,`$ (V.12)
$`\{\varphi ^\alpha ,\pi ^\beta \}_D`$ $`=`$ $`\mathrm{\Delta }^{\alpha \beta }(\varphi ),`$ (V.13)
$`\{\pi ^\alpha ,\pi ^\beta \}_D`$ $`=`$ $`(\varphi ^\beta \pi ^\alpha \varphi ^\alpha \pi ^\beta )/\varphi ^2.`$ (V.14)
The right sides of Eqs.(V.9) - (V.11) are reproduced at $`\chi =0`$.
(ii) The relations (V.9) - (V.11) define a Poisson algebra in the space of functions $`f(\varphi _u^\alpha ,\pi _u^\alpha )`$ depending on the projected canonical variables, so they can be used to generate consistently a Hamiltonian phase flow on the constraint submanifold $`\chi =0`$.
(iii) The Hamiltonian function
$$=\frac{1}{2}\mathrm{\Delta }_u^{\alpha \beta }(\varphi _u)\pi _u^\alpha \pi _u^\beta $$
(V.15)
coincides with Eq.(III.4): $``$$`=(\varphi _u^\alpha ,\pi _u^\alpha )`$$`=(\varphi ^\alpha ,\pi ^\alpha )`$. The Hamiltonian (III.4) is thus the function of the projected variables $`\xi _u(\xi )`$. It can be defined first on the submanifold $`\chi =0`$ and then extended to the unconstrained phase space using the gradient projection parallel to the phase flow associated to $`\mathrm{\Omega }`$. The relation $`\{\mathrm{\Omega },\}=0`$, Eq.(III.6), is the necessary and sufficient condition (see Eq.(V.4)) for $``$ to be a function of a fewer number of variables $`=(\varphi _u^\alpha ,\pi _u^\alpha )`$.
Let us consider the gradient projection onto the submanifold $`\mathrm{\Omega }=0`$ using the constraint function $`\chi `$. The constraint $`\chi =0`$ is responsible for shifts of the longitudinal component of the canonical momenta Eqs.(III.11) and (III.12). The functions $`\varphi _v^\alpha `$ and $`\pi _v^\alpha `$ have the form
$`\varphi _v^\alpha `$ $`=`$ $`\varphi ^\alpha ,`$ (V.16)
$`\pi _v^\alpha `$ $`=`$ $`\pi ^\alpha \varphi ^\alpha \varphi \pi /\varphi ^2.`$ (V.17)
Equation $`\mathrm{\Omega }(\varphi _v,\pi _v)=0`$ is fulfilled identically.
The Poisson bracket relations for the projected variables (V.16) and (V.17) can be found to be
$`\{\varphi _v^\alpha ,\varphi _v^\beta \}`$ $`=`$ $`0,`$ (V.18)
$`\{\varphi _v^\alpha ,\pi _v^\beta \}`$ $`=`$ $`\mathrm{\Delta }^{\alpha \beta }(\varphi _v),`$ (V.19)
$`\{\pi _v^\alpha ,\pi _v^\beta \}`$ $`=`$ $`(\varphi _v^\beta \pi _v^\alpha \varphi _v^\alpha \pi _v^\beta )/\varphi _v^2.`$ (V.20)
One can see again:
(i) The Poisson bracket of the projected variables coincides with the Dirac bracket Eqs.(V.12) - (V.14) associated to the constraints $`\chi =0`$ and $`\mathrm{\Omega }=0`$ on the submanifold $`\mathrm{\Omega }=0`$.
(ii) The Poisson bracket relations (V.18) - (V.20) are closed and define thereby a Poisson algebra in the space of functions depending on the projected canonical variables $`(\varphi _v^\alpha ,\pi _v^\beta )`$.
(iii) The Hamiltonian function
$$=\frac{1}{2}\varphi _v^2\mathrm{\Delta }^{\alpha \beta }(\varphi _v)\pi _v^\alpha \pi _v^\beta .$$
(V.21)
coincides with Eq.(III.4): $`=(\varphi _v^\alpha ,\pi _v^\alpha )=(\varphi ^\alpha ,\pi ^\alpha )`$. $``$ given by Eq.(III.4) is the function of the gradient variables $`\xi _v(\xi )`$. The relation $`\{\chi ,\}=0`$, Eq.(III.10), is the necessary and sufficient condition to present the Hamiltonian function as a function the projected variables: $`=(\varphi _v^\alpha ,\pi _v^\alpha )`$.
It is clear that the $``$ is defined finally on the intersection of the submanifolds $`\chi =0`$ and $`\mathrm{\Omega }=0`$, being thus a function of the $`\xi _s(\xi )`$. Eq.(III.4) represents its extension to the unconstrained phase space using the full gradient projection.
The statements (i) - (iii) are of the general validity for gradient projections. The statement (iii) has been proved as such above, the other two ones are proved below.
### V.4 Dirac bracket calculated by gradient projection
The phase flow associated to a function $`g=g(\xi )`$ defined on the submanifold $`𝒢_a=0`$ has an ambiguity since one can add to $`g=g(\xi )`$ a linear combination of the constraints $`\lambda _a𝒢_a`$ where $`\lambda _a`$ are undetermined parameters. The phase flow $`L_g[.]`$ applied to a function $`f=f(\xi )`$ suffers from this ambiguity also:
$$L_g[f]=\{f,g\}+\lambda _a\{f,𝒢_a\}.$$
(V.22)
The submanifold $`𝒢_a=0`$ should, however, be invariant, i.e., $`L_g[𝒢_a]=0`$. This equation allows to find $`\lambda _a=\lambda _a(\xi )`$. Substituting $`\lambda _a`$ into Eq.(V.22), one gets the Dirac bracket
$`L_g[f]`$ $`=`$ $`\{f,g\}\{f,𝒢_a\}\{𝒢_a,𝒢_b\}^1\{𝒢_b,g\}`$ (V.23)
$`=`$ $`\{f,g\}_D`$
where $`f`$, $`g`$, and $`𝒢_a`$ are functions of $`\xi `$. The Dirac bracket defines a phase flow generated by a function $`g=g(\xi )`$ within the constraint submanifold $`𝒢_a=0`$.
Using Eq.(V.1), one finds at $`𝒢_a=0`$
$`\{f(\xi ),g(\xi )\}_D`$ $`=`$ $`\{f(\xi ),g(\xi _s(\xi ))\}=\{f(\xi _s(\xi )),g(\xi )\}`$ (V.24)
$`=`$ $`\{f(\xi _s(\xi )),g(\xi _s(\xi ))\}.`$
This is the analogue of the statement (i) made in the previous subsection for the full gradient projection. The gradient projection can therefore be used to calculate the Dirac bracket.
There is an analogue of this statement for the partial gradient projections also. Let us suppose that the second-class constraints $`𝒢_a=0`$ split into the canonical pairs: $`\chi _A=0`$ and $`\mathrm{\Omega }_A=0`$ such that $`\{\chi _A,\chi _B\}=0`$, $`\{\mathrm{\Omega }_A,\mathrm{\Omega }_B\}=0`$, and $`det\{\chi _A,\mathrm{\Omega }_B\}0`$. The Dirac bracket becomes
$`\{f,g\}_D`$ $`=`$ $`\{f,g\}`$
$`+`$ $`\{f,\chi _A\}\{\chi _A,\mathrm{\Omega }_B\}^1\{\mathrm{\Omega }_B,g\}`$
$``$ $`\{f,\mathrm{\Omega }_A\}\{\chi _A,\mathrm{\Omega }_B\}^1\{\chi _B,g\}.`$
Let $`\xi _u(\xi )`$ and $`\xi _v(\xi )`$ be partial gradient projections such that $`\chi _A(\xi _u(\xi ))=0`$ and $`\mathrm{\Omega }_A(\xi _v(\xi ))=0`$ identically (cf. Eqs.(V.3) and (V.5)). If we replace arguments of the functions $`f`$ and $`g`$ with $`\xi _u(\xi )`$ or $`\xi _v(\xi )`$, the last two terms vanish due to (V.4) or (V.6), respectively. The Poisson bracket for the projected variables then coincides with the Dirac bracket for the canonical variables $`\xi `$ constrained to the submanifold $`\chi _A=0`$ or $`\mathrm{\Omega }_A=0`$:
$`\{f(\xi _u(\xi )),g(\xi _u(\xi ))\}_D`$ $`=`$ $`\{f(\xi _u(\xi )),g(\xi _u(\xi ))\},`$ (V.26)
$`\{f(\xi _v(\xi )),g(\xi _v(\xi ))\}_D`$ $`=`$ $`\{f(\xi _v(\xi )),g(\xi _v(\xi ))\}.`$ (V.27)
The arguments of the functions represent the partial gradient projections like in Eqs.(V.9) - (V.11) and (V.18) - (V.20).
Eqs.(V.26) and (V.27) are sufficient to calculate the Dirac bracket given that the partial gradient projections are constructed.
This completes the proof of the statement (i) from the previous subsection, extended to arbitrary Hamiltonian systems.
Turning to the point (ii), it is sufficient to notice that the Poisson bracket, e.g., $`\{\xi _s^i,\xi _s^j\}`$ determines a variation of the $`\xi _s^i`$ along the submanifold $`𝒢_a=0`$. This submanifold is parameterized with the $`\xi _s`$. The Poisson bracket is thus a function of the $`\xi _s`$ again. The involution relations for $`\{\xi _s^i,\xi _s^j\}`$ define therefore an algebra. Similar arguments apply for the partial gradient projections. Eqs.(V.9) - (V.11) and (V.18) - (V.20) represent therefore an illustration of this statement.
### V.5 Constraint functions $`\stackrel{~}{𝒢}_a`$ constructed by gradient projection
The statements of Sect. 4 can be proved using the gradient projection method.
The vector fields
$$I^{ij}\frac{𝒢_a}{\xi ^j}$$
(V.28)
determine phase flows associated to the constraint functions $`𝒢_a`$. These fields are non-singular, i.e., do not vanish in a neighborhood of the constraint submanifold $`=\{\xi :𝒢_a(\xi )=0a\}`$. The opposite would mean $`a`$ such that (V.28) vanish at some point $`\xi `$. It follows then that $`\{𝒢_a,𝒢_b\}=0`$ $`b`$. This is in contradiction with
$$det\{𝒢_a,𝒢_b\}0$$
(V.29)
which holds, by assumption, everywhere on $``$ and, by continuity, in a neighborhood of $``$. In Eq.(V.28), $`I^{ij}=I_{ij}`$ where
$$I_{ij}=\begin{array}{cc}0\hfill & E_n\hfill \\ E_n\hfill & 0\hfill \end{array},$$
(V.30)
with $`E_n`$ being the $`n\times n`$ identity matrix (cf. Eq.(IV.2)). In what follows, we denote phase flows (V.28) briefly as $`Id𝒢_a`$.
Let $`_a=\{\xi :𝒢_a(\xi )=0\}`$. The condition $`Id𝒢_a(\xi )0`$ $`\xi _a`$ is stronger than the one mentioned above. It looks evident, since any constraint function $`𝒢_a(\xi )`$ in a neighborhood of $`_a`$ can be redefined to assign the gradient a definite direction.
$`_a`$ is a subspace of the dimension $`2n1`$. The intersection of all $`_a`$ gives $``$. The tangent space to $`_a`$ at a point $`\xi `$, denoted as $`T_\xi _a`$, is skew-orthogonal to $`Id𝒢_a(\xi )`$. Indeed, if $`d𝒢_a=0`$ then $`d\xi T_\xi _a`$, and we obtain
$$I^{ij}\frac{𝒢_a}{\xi ^j}I_{ik}d\xi ^k=\frac{𝒢_a}{\xi ^i}d\xi ^i=0.$$
(V.31)
The space $`T_\xi _a`$ has the dimension $`2n1`$. Among the vectors of $`T_\xi _a`$ one can find $`Id𝒢_a(\xi )`$, since $`Id𝒢_a(\xi )`$ is skew-orthogonal to itself. $`T_\xi _a`$ is therefore a skew-orthogonal complement of $`Id𝒢_a(\xi )`$.
One can find $`b`$ such that $`\{𝒢_a(\xi ),𝒢_b(\xi )\}0`$. As discussed above,
$`Id𝒢_a(\xi )`$ $``$ $`T_\xi _a,`$ (V.32)
$`Id𝒢_b(\xi )`$ $``$ $`T_\xi _b,`$ (V.33)
and furthermore,
$`Id𝒢_a(\xi )`$ $``$ $`T_\xi _b,`$ (V.34)
$`Id𝒢_b(\xi )`$ $``$ $`T_\xi _a.`$ (V.35)
From the other side, any vector $`d\xi T_\xi _aT_\xi _b=T_\xi (_a_b)`$ is skew-orthogonal to $`Id𝒢_a(\xi )`$ and $`Id𝒢_b(\xi )`$. As a consequence of Eqs.(V.34) and (V.35), one gets
$`Id𝒢_a(\xi )T_\xi (_a`$ $``$ $`_b),`$ (V.36)
$`Id𝒢_b(\xi )T_\xi (_a`$ $``$ $`_b).`$ (V.37)
The subspace $`_a_b`$ and, respectively, $`T_\xi (_a_b)`$ have the dimension $`2n2`$. The vectors $`Id𝒢_a(\xi )`$ and $`Id𝒢_b(\xi )`$ are linearly independent and form a two-dimensional space $`𝒟_{ab}(\xi )`$ which is a skew-orthogonal complement of $`T_\xi (_a_b)`$ such that $`𝒟_{ab}(\xi )T_\xi (_a_b)=\mathrm{}`$.
Let us consider motion of a particle with a Hamiltonian function $`𝒢_b`$. In virtue of Eq.(V.35), the phase flow $`Id𝒢_b`$ does not lie in the submanifold $`_a`$ entirely and therefore crosses it. Let $`t_a(\zeta )`$ be time needed for particle located at $`\zeta _a`$ in a neighborhood of $`\xi _a`$ to cross $`_a`$ at some point $`\eta \xi `$. The equations of motion look like
$$\frac{d\zeta }{dt_a}=Id𝒢_b(\zeta ).$$
The derivative of $`t_a(\zeta )`$ along the phase flow $`Id𝒢_b(\zeta )`$ with respect to time is, by definition, equal to unity:
$$\{t_a(\zeta ),𝒢_b(\zeta )\}=1.$$
(V.38)
One may interpret, equivalently, $`𝒢_b(\zeta )`$ as a time needed to cross the submanifold $`_b`$ by a particle located at $`\zeta _b`$. The motion of such a particle is described by a Hamiltonian function $`t_a(\zeta )`$.
The function $`t_a(\zeta )`$ vanishes for $`\zeta =\eta _a`$. At any point of $`_a`$, $`d𝒢_a`$ and $`dt_a`$ vanish for $`d\eta T_\eta _a`$. There exists only one $`d\eta T_\eta _a`$ such that $`d𝒢_a0`$ and $`dt_a0`$. It means that $`d\eta `$ $`d𝒢_a`$ is proportional to $`dt_a`$ and $`Id𝒢_a`$ is in turn proportional to $`Idt_a`$.
The first canonical pair $`\stackrel{~}{𝒢}_a=t_a`$ and $`\stackrel{~}{𝒢}_b=𝒢_b`$ is thus constructed.
Let us consider the full gradient projection $`\xi _1(\xi )`$ onto the submanifold $`_a_b`$, using the constraint functions $`\stackrel{~}{𝒢}_a`$ and $`\stackrel{~}{𝒢}_b`$. One gets $`\stackrel{~}{𝒢}_a(\xi _1(\xi ))0`$, $`\stackrel{~}{𝒢}_b(\xi _1(\xi ))0`$, whereas the equations $`𝒢_c(\xi _1(\xi ))=0`$ for $`ca,b`$ are significant to determine the location of the constraint submanifold $``$, owing to shifts along the phase flows $`Id\stackrel{~}{𝒢}_a`$ and $`Id\stackrel{~}{𝒢}_b`$. A complete set of equations for $``$ can be taken to be
$`\stackrel{~}{𝒢}_a(\xi )`$ $`=`$ $`0,`$ (V.39)
$`\stackrel{~}{𝒢}_b(\xi )`$ $`=`$ $`0,`$ (V.40)
$`𝒢_c(\xi _1(\xi ))`$ $`=`$ $`0`$ (V.41)
for $`ca,b`$. In virtue of Eq.(V.2),
$`\{\stackrel{~}{𝒢}_a(\xi ),𝒢_c(\xi _1(\xi ))\}`$ $`=`$ $`0,`$ (V.42)
$`\{\stackrel{~}{𝒢}_b(\xi ),𝒢_c(\xi _1(\xi ))\}`$ $`=`$ $`0.`$ (V.43)
Eqs.(V.40)-(V.41) determine $``$ uniquely, so the determinant of the Poisson bracket relations between the $`2m`$ constraint functions is not zero. The functions $`\stackrel{~}{𝒢}_a`$ and $`\stackrel{~}{𝒢}_b`$ have the vanishing Poisson brackets with the rest ones, so
$$det\{𝒢_c(\xi _1(\xi )),𝒢_d(\xi _1(\xi ))\}0$$
(V.44)
where $`c,d`$ take $`2m2`$ values ($`c,da,b`$).
In the remaining set of the constraints, one can find $`c,d`$ such that $`\{𝒢_c(\xi _1(\xi )),𝒢_d(\xi _1(\xi ))\}0`$ and repeat the arguments we used earlier. The analogue of Eq.(V.38) looks like
$$\{t_c(\zeta ),𝒢_d(\xi _1(\zeta ))\}=1.$$
(V.45)
The Poisson brackets of the left side of this equation with $`\stackrel{~}{𝒢}_a(\zeta )`$ and $`\stackrel{~}{𝒢}_a(\zeta )`$ vanish. The Jacoby identity yields
$`\{\{t_c(\zeta ),\stackrel{~}{𝒢}_a(\zeta )\},𝒢_d(\xi _1(\zeta ))\}`$ $`=`$ $`0,`$ (V.46)
$`\{\{t_c(\zeta ),\stackrel{~}{𝒢}_b(\zeta )\},𝒢_d(\xi _1(\zeta ))\}`$ $`=`$ $`0.`$ (V.47)
The Poisson brackets of $`t_c(\zeta )`$ with $`\stackrel{~}{𝒢}_a(\zeta )`$ and $`\stackrel{~}{𝒢}_b(\zeta )`$ remain therefore constant along the phase flow $`Id𝒢_d(\xi _1(\zeta ))`$. At the submanifold $`_c^{}=\{\zeta :𝒢_c(\xi _1(\zeta ))=0\}`$, $`Idt_c(\zeta )`$ is proportional to $`Id𝒢_c(\xi _1(\zeta ))`$. Those brackets vanish at $`_c^{}`$, and furthermore, vanish for $`\zeta _c^{}`$. Eq.(V.2) suggests then $`t_c(\zeta )=t_c(\xi _1(\zeta ))`$.
The second canonical pair $`\stackrel{~}{𝒢}_c(\zeta )=t_c(\xi _1(\zeta ))`$ and $`\stackrel{~}{𝒢}_d(\zeta )=𝒢_d(\xi _1(\zeta ))`$ is thus constructed.
The proof can be completed by induction. We consider the full gradient projection $`\xi _2(\zeta )`$ onto the submanifold $`_a_b_c_d`$ along the commutative phase flows generated by $`\stackrel{~}{𝒢}_a(\zeta )`$, $`\stackrel{~}{𝒢}_b(\zeta )`$, $`\stackrel{~}{𝒢}_c(\zeta )`$, and $`\stackrel{~}{𝒢}_d(\zeta )`$. These constraint functions have the vanishing Poisson brackets with the $`2m4`$ remaining ones $`𝒢_e(\xi _2(\zeta ))`$ ($`ea,b,c,d`$). The latters constitute a complete non-degenerate set to determine the constraint submanifold $``$ uniquely, and so on.
At the end, one gets in a neighborhood of $`\xi `$ a symplectic basis (IV.23).
### V.6 Hamiltonian $`\stackrel{~}{}`$ constructed by gradient projection
The Hamiltonian $`\stackrel{~}{}`$ of Sect. 4 can also be constructed with the help of Eq.(V.2):
$$\stackrel{~}{}(\xi )=(\xi _s(\xi ))$$
(V.48)
where $`\xi _s(\xi )`$ are gradient projections defined by $`\stackrel{~}{𝒢}_a`$.
Eq.(V.48) and the algorithm described in Sect. 4 give, apparently, an equivalent Hamiltonian function, since the Hamiltonian flows on the constraint submanifold coincide. It can be demonstrated by the comparison of the Dirac brackets:
$$\{\xi ^i,\stackrel{~}{}(\xi )\}_D=\{\xi ^i,(\xi )\}_D$$
(V.49)
which holds due to Eq.(V.24).
Applications of the gradient projection method to constructing the quantum deformation of the Dirac bracket can be found in QDDB .
## VI Dirac quantization of spherical pendulum
In Sect. 2, we modified the initial Lagrangian (II.2) on the constraint submanifold $`\chi =\mathrm{ln}\varphi =0`$ to make more transparent the origin of the underlying dilatation gauge symmetry. Let us formulate the Hamiltonian dynamics of the spherical pendulum starting directly from Lagrangian $`+\lambda \mathrm{ln}\varphi `$. The straightforward follow-up to the Dirac’s scheme (see also ABDA01 ; MFRF ) leads to the Hamiltonian dynamics described in Sect. 3:
Using the Dirac’s scheme, we obtain the primary Hamiltonian $`_p=\frac{1}{2}\pi ^2\lambda \mathrm{ln}\varphi `$ and the primary constraint $`𝒢_0=\pi _\lambda 0`$, where $`\pi ^\alpha `$ are canonical momenta associated to the canonical coordinates $`\varphi ^\alpha `$ and $`\pi _\lambda `$ is the canonical momentum associated to the Lagrange multiplier $`\lambda `$. The canonical Hamiltonian becomes $`_c=_p+u\pi _\lambda `$ where $`u`$ is an undetermined function of time. The secondary constraints $`𝒢_{a+1}=\{𝒢_a,_c\}`$ can be found: $`𝒢_1=\mathrm{ln}\varphi `$, $`𝒢_2=\varphi \pi /\varphi ^2`$, $`𝒢_3=\pi ^2/\varphi ^22(\varphi \pi )^2/\varphi ^4+\lambda /\varphi ^2`$. The last constraint $`𝒢_3=0`$ allows to fix $`\lambda `$ in terms of $`\varphi ^\alpha `$ and $`\pi ^\alpha `$, no new constraints then appear.
The dimension of the phase space can be reduced by eliminating the canonically conjugate pair $`(\lambda ,\pi _\lambda )`$. We solve equations $`𝒢_0=0`$ with respect to $`\pi _\lambda `$ and $`𝒢_3=0`$ with respect to $`\lambda `$ and substitute solutions into $`_c`$. The result is
$$_c^{}=\frac{1}{2}\pi ^2+(\pi ^22(\varphi \pi )^2/\varphi ^2)\mathrm{ln}\varphi .$$
(VI.1)
There remain two constraint functions $`𝒢_1`$ and $`𝒢_2`$, such that $`\{𝒢_1,𝒢_2\}=1/\varphi ^2`$. These are the second-class constraints. The Hamiltonian $`_c^{}`$ is first-class: $`\{𝒢_1,_c^{}\}=𝒢_2𝒢_1\{𝒢_1,\lambda \}`$ and $`\{𝒢_2,_c^{}\}=𝒢_1\{𝒢_2,\lambda \}`$ where $`\lambda `$ is determined by $`𝒢_3=0`$.
In order to split the constraints into the gauge-fixing condition and the gauge generator, we should construct, according to the previous Section, $`\stackrel{~}{𝒢}_a`$ and $`\stackrel{~}{}`$ starting from $`𝒢_a`$ and $`_c^{}`$.
In Sect. 3, we already had the same constraint submanifold described by functions $`\chi =\mathrm{ln}\varphi `$ and $`\mathrm{\Omega }=\varphi \pi `$ satisfying $`\{\chi ,\mathrm{\Omega }\}=1`$, so one can set $`\stackrel{~}{𝒢}_1=\mathrm{ln}\varphi `$ and $`\stackrel{~}{𝒢}_2=\varphi \pi `$.
The Hamiltonian $`\stackrel{~}{}`$ can be constructed with the help of Eq.(V.48). Let us combine Eqs.(V.7), (V.8) and (V.16), (V.17) to get the full gradient projection:
$`\varphi _s^\alpha `$ $`=`$ $`\varphi ^\alpha /\varphi ,`$ (VI.2)
$`\pi _s^\alpha `$ $`=`$ $`\varphi \pi ^\alpha \varphi ^\alpha \varphi \pi /\varphi .`$ (VI.3)
Replacing $`\varphi ^\alpha `$ and $`\pi ^\alpha `$ in $`_c^{}`$ by variables $`\varphi _s^\alpha `$ and $`\pi _s^\alpha `$, respectively, we get
$$\stackrel{~}{}=\frac{1}{2}\pi _s^2$$
(VI.4)
which coincides with the Hamiltonian (III.4).
The difference $`\stackrel{~}{}_c^{}`$ is of the second order in the constraint functions. This is sufficient to have identical Hamiltonian flows on the constraint submanifold.
The equivalence is thus demonstrated, so further discussion can go in the parallel with Sect. 3. In the example considered, the constraint functions are of the second class, whereas in Sect. 3 they appear as the gauge-fixing condition, $`\chi =0`$, and the gauge generator, $`\mathrm{\Omega }`$. This suggests that the interpretation of second-class constraints of a holonomic system is a matter of convention.
## VII Second-class constraints as generators of gauge symmetries and gauge-fixing conditions
The type of constraints appearing within the Hamiltonian framework depends on the form of the corresponding Lagrangian. In case of the spherical pendulum, starting from Eq.(II.2) one arrives to the constraints $`\chi =0`$ and $`\mathrm{\Omega }=0`$ as to the second-class constraints of the Hamiltonian framework. Starting from Eq.(II.11), the constraint $`\chi =0`$ appears as a gauge-fixing condition, whereas the constraint $`\mathrm{\Omega }=0`$ as a gauge generator of the symmetry group. Lagrangians (II.2) and (II.11) are equivalent at the classical level, so they lead to the same classical dynamics.
It is possible, therefore, at least in the case of spherical pendulum, to interpret second-class constraints as a gauge-fixing condition and a gauge generator, and vice versa.
We wish to discuss whether this is a common situation. Any set of admissible gauge-fixing conditions and gauge generators can be treated as second-class constraints. This statement is widely used for quantization of gauge theories (see, e.g., SLFA ).
The reverse statement represents apparent interest also. The second-class constraints systems are analyzed from this point of view by Mitra and Rajaraman MITRA1 . In this Section, we discuss additional details about the underlying gauge symmetries of second-class constraints systems.
### VII.1 Converting a gauge system into a second-class constraints system
Let a primary Hamiltonian $``$ of a system be gauge invariant. The generators of gauge transformations $`\mathrm{\Omega }_A=0`$ are such that
$`\{\mathrm{\Omega }_A,\}`$ $`=`$ $`_A^B\mathrm{\Omega }_B,`$ (VII.1)
$`\{\mathrm{\Omega }_A,\mathrm{\Omega }_B\}`$ $`=`$ $`𝒞_{AB}^C\mathrm{\Omega }_C.`$ (VII.2)
The gauge-fixing conditions $`\chi _A=0`$ fulfill
$`\{\chi _A,\chi _B\}`$ $`=`$ $`0,`$ (VII.3)
$`det\{\chi _A,\mathrm{\Omega }_B\}`$ $``$ $`0.`$ (VII.4)
After imposing gauge-fixing conditions, the system becomes equivalent to a system described by a first-class Hamiltonian $`^{}`$ defined on the submanifold of second-class constraints $`𝒢_a=(\chi _A,\mathrm{\Omega }_A)`$. The Hamiltonian $`^{}`$ can be constructed from the canonical Hamiltonian of the original system, $`_c=+\lambda _B\mathrm{\Omega }_B`$, by requiring
$$\{\chi _A,_c\}=\{\chi _A,\}+\lambda _B\{\chi _A,\mathrm{\Omega }_B\}0.$$
(VII.5)
The gauge parameters $`\lambda _B`$ can be fixed from this equation in terms of the canonical variables. We thus get the first-class Hamiltonian $`^{}=_c`$ and second-class constraints $`𝒢_a`$ satisfying $`\{𝒢_a,^{}\}0`$ and $`det\{𝒢_a,𝒢_b\}0`$.
The gauge-fixing is equivalent to converting the original system into an equivalent one described by a first-class Hamiltonian and second-class constraints.
### VII.2 Converting a second-class constraints system into a gauge system
It was shown by Dirac DIRAC1 ; DIRAC2 that first-class constraints imply the presence of unphysical degrees of freedom the evolution of which is not fixed by the Hamilton’s equations. The dynamics is self-consistent provided the unfixed degrees of freedom do not belong to the set of physical observables. This is the case of gauge-invariant systems which have gauge degrees of freedom. According to the Dirac’s constraint dynamics, having first-class constraints is a necessary condition for the existence of a gauge symmetry. The reverse statement is less obvious:
Gauge transformations of the canonical coordinates have the form
$$\delta \varphi ^\alpha =\{\varphi ^\alpha ,\mathrm{\Omega }_A\}\theta _A$$
(VII.6)
where $`\theta _A`$ are infinitesimal functions of time. Since $`\mathrm{\Omega }_A`$ depend on the canonical coordinates and momenta, the variations $`\delta \varphi ^\alpha `$ depend on the coordinates and momenta also, unless $`\mathrm{\Omega }_A`$ are first degree polynomials of the canonical momenta. Subject to a Legendre transform, $`\delta \varphi ^\alpha `$ become, in general, functions of the coordinates and velocities. The gauge theories such as QED and QCD have the constraint functions $`\mathrm{\Omega }_A`$ as the first degree polynomials of the canonical momenta. In such cases $`\delta \varphi ^\alpha =\delta \varphi ^\alpha (\varphi ^\beta )`$ do not depend on velocities and define thereby transformations on the configuration space $`=(\varphi ^\alpha )`$ which induce, in turn, transformations in the tangent bundle $`T=(\varphi ^\alpha ,\dot{\varphi }^\alpha )`$.
The first-class constraints systems correspond to more general class of gauge-invariant systems with $`\delta \varphi ^\alpha =\delta \varphi ^\alpha (\varphi ^\beta ,\dot{\varphi }^\beta )`$ and, probably, new auxiliary variables. In what follows, we distinguish thereby between gauge transformations of the form $`\delta \varphi ^\alpha =\delta \varphi ^\alpha (\varphi ^\beta )`$ and generalized gauge transformations with velocity dependent parameters and/or new auxiliary variables.
Having a first-class constraint system is a necessary, but not a sufficient condition for an equivalent gauge-invariant system to exist in the original configuration space.
The possibility of constructing gauge-invariant systems in the unconstrained phase space, equivalent to second-class constraints systems upon a gauge-fixing, has been analyzed by Mitra and Rajaraman MITRA1 .
The applications discussed MITRA1 ; VYTHE1 ; VYTHE2 have constraint functions, associated to gauge transformations, as polynomials of the canonical momenta of the degree less or equal to unity. The gauge transformations are not velocity dependent, although involve new auxiliary variables. The gauged systems require generally an extended configuration space.
If global symplectic basis in the space of constraint functions can be found, equivalent generalized gauge systems for second-class constraints systems can be constructed in the unconstrained phase space.
By passing over to the standard canonical frame, one can always select $`\stackrel{~}{\mathrm{\Omega }}_A`$ as first canonical momenta. This is, however, not sufficient to have gauge invariance on the physical configuration space. Before doing a Legendre transform, one has to pass first to a canonical coordinate system where the coordinates $`\varphi ^\alpha `$ constitute a physical configuration space. This will mix up the canonical coordinates and momenta, preventing from having $`\stackrel{~}{\mathrm{\Omega }}_A`$ as the first degree polynomials.
It is worthwhile to notice that gauge invariant observables are not measurable if they involve auxiliary degrees of freedom. The sums of the vector potentials of the massive electrodynamics and the derivatives of the Stückelberg scalar are gauge invariant. They do not belong, however, to the set of physical observables. The equivalence with the ordinary gauge systems, where sets of gauge invariant quantities and physical observables coincide, is therefore not complete.
### VII.3 Gauge-invariant systems as holonomic systems
The quantization of gauge theories which appear under first-class constraints is studied in many details (see, e.g., Refs. HENN ; SLFA ). The second-class constraints systems are usually quantized by employing the BFT formalism HENN ; BATA ; BATY that converts second-class constraints into a first class by extending the original phase space to a higher dimension.
The BFT algorithm combined with the Fradkin-Vilkovisky quantization scheme FRVI ; FRAD if applied to a system in a $`2n`$-dimensional phase space ends up with an extended phase space of a dimension larger or equal to $`2n+12m`$, since the $`2m`$ constraints convert into first-class constraints, each of which requires an independent gauge-fixing condition. The remaining at least $`8m`$ degrees of freedom appear as ghost variables.
The problem of constructing gauge-invariant systems equivalent to second-class constraints systems in the original phase space is discussed in Ref. MITRA1 . The existence of the global symplectic basis for the constraint functions is specified as sufficient conditions for the existence of gauged partners associated to the second-class constraints systems. The method MITRA1 , being successful in removing auxiliary fields from the original phase space, generates auxiliary fields in the effective Lagrangians.
The systems of pointlike particles under holonomic constraints have the natural gauge counterparts both in the phase space and configuration space MFRF .
In holonomic systems, second-class constraints split into gauge-fixing conditions, $`\stackrel{~}{\chi }_A=0`$, and gauge generators, $`\stackrel{~}{\mathrm{\Omega }}_A`$. Such systems can be quantized further as gauge theories. The BFT method, if applied, would result to an extended phase space of the dimension $`2n+6m`$, containing $`2m`$ Lagrange multipliers and $`4m`$ ghost variables. The underlying gauge symmetry allows to reduce the number of auxiliary fields within the quantization scheme FRVI ; FRAD .
### VII.4 Frobenius’ condition for holonomic constraints
The holonomic constraints are defined by a constraint submanifold $`𝒩`$ in the configuration space $``$. It comes then out automatically that particle velocities belong to the tangent space $`T𝒩`$ of the constraint submanifold.
An $`nm`$-dimensional tangent subspace $`T𝒩T`$ of an $`n`$-dimensional configuration space $``$ can be defined in general through $`m`$ covector fields $`\omega _{A\alpha }`$ with $`A=1,\mathrm{},m`$ and $`\alpha =1,\mathrm{},n`$ by imposing constraints
$$\omega _{A\alpha }\dot{\varphi }^\alpha =0$$
(VII.7)
which are velocity dependent and therefore non-holonomic.
The Frobenius’ integrability conditions (see, e.g., ARNO ),
$$\frac{\omega _{A\alpha }}{\varphi ^\beta }(\dot{\varphi }_1^\alpha \dot{\varphi }_2^\beta \dot{\varphi }_2^\alpha \dot{\varphi }_1^\beta )=0,$$
(VII.8)
for tangent vectors $`\dot{\varphi }_1^\alpha `$ and $`\dot{\varphi }_2^\beta `$ satisfying (VII.7), if fulfilled, implies the existence of a submanifold $`𝒩`$ the tangent space of which coincides with the tangent subspace (VII.7). In such a case, the constraints (VII.7) can be replaced with holonomic constraints $`\chi _A=0`$, identifying the constraint submanifold $`𝒩`$, without modifying the dynamics.
Eqs.(VII.8) specify non-holonomic systems which have gauged counterparts in the original configuration space.
For an $`n=3`$ Euclidean space, Eqs.(VII.7) define, for $`m=1`$, a submanifold orthogonal to the vector $`\omega _\alpha `$. Eqs.(VII.8) tell us that the vector field $`\stackrel{}{\omega }`$ is a potential field provided that $`rot\stackrel{}{\omega }=0`$. One can find a potential function $`\chi `$ such that $`\stackrel{}{\omega }=\stackrel{}{}\chi `$. The sets of $`\chi =\mathrm{𝚌𝚘𝚗𝚜𝚝𝚊𝚗𝚝}`$ define, in our case, possible constraint submanifolds of an equivalent holonomic system.
The tangent subspace determined by Eq.(II.3) coincides with the tangent space of $`S^{n1}`$. The covector field $`\omega _\alpha =\varphi ^\alpha `$ satisfies the Frobenius’ condition.
## VIII Supplementary conditions for Wigner functions
The splitting of second-class constraints into constraints associated to gauge transformations and gauge-fixing conditions is not unique. Every pair ($`\chi _A`$,$`\mathrm{\Omega }_A`$) can be transformed, e.g., as $`\chi _A\mathrm{\Omega }_A`$, $`\mathrm{\Omega }_A\chi _A`$. The symplectic group $`Sp(2m)`$ of transformations mixes $`\chi _A`$ and $`\mathrm{\Omega }_A`$ further, keeping the Poisson brackets invariant. Canonical transformations, furthermore, do not modify the Poisson bracket of the constraints also. There exists therefore an affluent spectrum of representations and we have to demonstrate that the physical content of theory is invariant with respect to the choice of representation.
In the general case of second-class constraints, we do not have any criterion how to discriminate between first-class constraints associated to gauge transformations and gauge-fixing conditions. We show, conversely, that the way the second-class constraints are split is not important for quantization.
A useful hint towards that conclusion is delivered by the path integral representation for the evolution operator (I.6) which is invariant with respect to linear transformations of $`𝒢_a`$.
The problem is to demonstrate that the supplementary conditions $`\widehat{\mathrm{\Omega }}_A\mathrm{\Psi }=0`$ can be replaced by $`\widehat{\chi }_A\mathrm{\Psi }=0`$. In case of the spherical pendulum, this is almost evident, since $`\widehat{\mathrm{\Omega }}\mathrm{\Psi }=0`$ means $`\mathrm{\Psi }\mathrm{\Psi }(\varphi ^\alpha /\varphi )`$ whereas $`\widehat{\chi }\mathrm{\Psi }=0`$ means $`\mathrm{\Psi }\delta (\chi )\mathrm{\Psi }(\varphi ^\alpha )=\delta (\chi )\mathrm{\Psi }(\varphi ^\alpha /\varphi )`$. The essential dependence comes from the angular part of the wave function, while the radial part $`\delta (\chi )`$ is factorized, being commutative, Eq.(III.10), with the Hamiltonian and the $`S`$-matrix. It can be absorbed thereby by an overall normalization factor.
These arguments can apparently be extended to an arbitrary case. According to the Dirac’s prescription, the physical wave functions are annihilated by the constraint operators associated to gauge transformations
$$\widehat{\mathrm{\Omega }}_A\mathrm{\Psi }=0.$$
(VIII.1)
To simplify the matter, we pass to the standard canonical coordinate system and choose the constraint functions $`\chi _A`$ as first $`m`$ canonical coordinates $`q_A`$ and the constraint functions $`\mathrm{\Omega }_A`$ as first $`m`$ canonical momenta $`p_A`$. The remaining canonical variables are ($`q^{},p^{}`$) constitute the physical phase space $`\mathrm{\Gamma }_{}^{2(nm)}`$. Eq.(VIII.1) gives then, in the coordinate representation, a wave function of the form $`\mathrm{\Psi }=\mathrm{\Psi }(q^{})`$. Such a wave function has an infinite norm, since the integral $`d^nq|\mathrm{\Psi }(q^{})|^2`$ diverges. In the momentum representation, we would get $`\mathrm{\Psi }(p)=(2\pi \mathrm{})^m_A\delta (p_A)\mathrm{\Psi }(p^{})`$. $`\mathrm{\Psi }(p)`$ has an infinite norm $`(2\pi \mathrm{})^m\delta ^m(0)`$ provided $`\mathrm{\Psi }(p^{})`$ is normalized to unity. An infinite but numerical factor can be included into the norm of $`\mathrm{\Psi }(p)`$.
Let us check whether wave functions satisfying
$$\widehat{\chi }_A\mathrm{\Psi }=0.$$
(VIII.2)
have a physical sense. In the coordinate representation, one gets $`\mathrm{\Psi }=_A\delta (q_A)\mathrm{\Psi }(q^{})`$ and in the momentum space $`\mathrm{\Psi }=\mathrm{\Psi }(p^{})`$. These states have infinite norms. There is an apparent symmetry $`pq`$, $`qp`$ between wave functions satisfying (VIII.1) and (VIII.2).
Owing to the factor $`_A\delta (q_A)`$, conditions (VIII.1) and (VIII.2) single out the same set of functions, the nontrivial dependence of which comes from the physical variables $`q^{}`$ ($`p^{}`$) only.
The Wigner functions WIGNER provide complete information on quantum systems. The association rules of operators in the Hilbert space and functions in the phase space are discussed a long time WEYL ; WIGNER ; GROE ; MOYAL . The quantum mechanics can be reformulated using the Groenewold star-product GROE representing a deformation of the usual pointwise product of functions in the phase space. The Wigner functions and constraint functions are defined in the phase space, so it is natural to discuss supplementary conditions in terms of the Wigner functions. In this approach, the equivalence between Eqs.(VIII.1) and (VIII.2) becomes more transparent. In addition, more general group of supplementary conditions, equivalent to Eqs.(VIII.1) and (VIII.2), can be formulated in terms of the Wigner functions.
### VIII.1 Probability density localized on the constraint submanifold
Let us start from the standard canonical frame MASKAWA where the constraint functions $`\chi _A`$ are the first $`m`$ canonical coordinates $`q_A`$ and the constraint functions $`\mathrm{\Omega }_A`$ are the first $`m`$ canonical momenta $`p_A`$. The probability density in the phase space $`\mathrm{\Gamma }^{2n}`$ can be written as follows:
$$W(q,p)=(2\pi \mathrm{})^m\underset{A}{}\delta (q_A)\delta (p_A)W_{}(q^{},p^{}).$$
(VIII.3)
Identifying the $`W(q,p)`$ with the Wigner function in the unconstrained phase space and using the Wigner transform, one gets the density matrix
$`\rho (q,q^{})`$ $`=`$ $`{\displaystyle W(\frac{q+q^{}}{2},p)e^{\frac{i}{\mathrm{}}p(qq^{})}\frac{d^np}{(2\pi \mathrm{})^n}}`$ (VIII.4)
$`=`$ $`{\displaystyle \underset{A}{}}\delta ({\displaystyle \frac{q_A+q_A^{}}{2}})\rho _{}(q^{},q^{}).`$
It satisfies the operator equations
$`\widehat{q}_A\widehat{\rho }+\widehat{\rho }\widehat{q}_A=0,`$ (VIII.5)
$`\widehat{p}_A\widehat{\rho }+\widehat{\rho }\widehat{p}_A=0.`$ (VIII.6)
The first equation implies $`\rho \delta (q_A+q_A^{})`$. The second one implies that $`\rho `$ does not depend on $`q_Aq_A^{}`$.
The density matrix $`\rho _{}(q^{},q^{})`$ is normalized to unity, so
$$\rho (q,q)d^nq=1.$$
(VIII.7)
One can make a unitary transformation to pass to an arbitrary set of operators associated to the canonical variables. In terms of
$$\widehat{𝒢}_a=𝒰(\widehat{q}_A,\widehat{p}_A)𝒰^+,$$
and $`\widehat{\rho }`$ replaced with $`𝒰\widehat{\rho }𝒰^+`$, Eqs.(VIII.5) and (VIII.6) become
$$\widehat{𝒢}_a\widehat{\rho }+\widehat{\rho }\widehat{𝒢}_a=0.$$
(VIII.8)
Eqs.(VIII.8) are necessary and sufficient conditions to have the representation (VIII.4) in the standard canonical frame.
The supplementary conditions (VIII.8) cannot be formulated in terms of a wave function, since the density matrix in the unconstrained configuration space does not correspond to a pure state, even if system on the constraint submanifold is in a pure state.
In an arbitrary frame and in the classical limit, Eq.(VIII.3) looks like
$$W(\xi )=(2\pi \mathrm{})^m\underset{a}{}\delta (𝒢_a)\sqrt{det\{𝒢_a,𝒢_b\}}W_{}(\xi _s(\xi )).$$
(VIII.9)
Note that $`_a\delta (𝒢_a)`$ acts as a projection operator, so that $`_a\delta (𝒢_a)f(\xi )=_a\delta (𝒢_a)f(\xi _s(\xi ))`$ $`f(\xi )`$.
In the classical limit, the Wigner function satisfies
$$𝒢_a(\xi )W(\xi )=0.$$
(VIII.10)
The complete phase-space analogue of quantum Eqs.(VIII.8) can be formulated in terms of the symmetric part of the Groenewold star-product GROE to give
$$𝒢_a(\xi )W(\xi )=0.$$
(VIII.11)
The star-product has the following decomposition:
$$f(\xi )g(\xi )=f(\xi )g(\xi )+\frac{i\mathrm{}}{2}f(\xi )g(\xi )$$
(VIII.12)
where
$`f(\xi )g(\xi )`$ $`=`$ $`f(\xi )\mathrm{cos}({\displaystyle \frac{\mathrm{}}{2}}𝒫)g(\xi ),`$ (VIII.13)
$`f(\xi )g(\xi )`$ $`=`$ $`f(\xi ){\displaystyle \frac{2}{\mathrm{}}}\mathrm{sin}({\displaystyle \frac{\mathrm{}}{2}}𝒫)g(\xi )`$ (VIII.14)
and
$$𝒫=I^{ij}\stackrel{}{\frac{}{\xi ^i}}\stackrel{}{\frac{}{\xi ^j}}$$
is the so-called Poisson operator. In the limit $`\mathrm{}0`$,
$$\underset{\mathrm{}0}{lim}f(\xi )g(\xi )=\{f(\xi ),g(\xi )\}.$$
(VIII.15)
The Poisson bracket $`\{f(\xi ),g(\xi )\}`$ coincides generally with a function associated to the commutator $`i/\mathrm{}[\widehat{f},\widehat{g}]`$ to the lowest order in the Planck’s constant only. The skew-symmetric part of the Groenewold product provides a generalization of the Poisson bracket which is skew-symmetric with respect to two functions, real for real functions, coincides with the Poisson bracket to the lowest order in the Planck’s constant, satisfies the Jacoby identity, and keeps the association rule for commutators. This part of the Groenewold product is known also as the ”sine bracket” or Moyal bracket MOYAL .
In the limit $`\mathrm{}0`$, the quantum condition (VIII.11) recovers the classical one (VIII.10). These two conditions coincide in the standard canonical frame.
The normalization condition Eq.(VIII.7) holds in the quantum case. The Wigner function satisfies
$$W(\xi )\frac{d^{2n}\xi }{(2\pi \mathrm{})^n}=1.$$
(VIII.16)
The classical limit of the Wigner function Eq.(VIII.9) provides the conventional normalization on the constraint submanifold for $`W_{}(\xi _s(\xi ))`$.
The Groenewold star-product is non-local. The usual pointwise product of two functions like in Eq.(VIII.10) has no quantum counterpart. It can be treated as a classical limit of a quantum operator equation only. By contrary, Eq.(VIII.11) is the phase space analogue of the quantum operator equation (VIII.8).
Eq.(VIII.3) accomplishes a trivial extrapolation of the Wigner function from the constraint submanifold: The density is set equal to zero when $`\xi `$ does not belong to the constraint submanifold. The Wigner function is not a smooth function across the constraint submanifold, so it is hard to formulate using this approach an evolution equation similar to the Liouville equation in the unconstrained phase space.
### VIII.2 Probability density localized on and outside of the constraint submanifold
One can require that at any given point $`\xi `$ the density be the same as at $`\xi _s=\xi _s(\xi )`$. The gradient projections $`\xi _s`$ can be constructed, as discussed in Sect. 5, using phase flows generated by the constraint functions $`𝒢_a`$ to solve equations $`𝒢_a(\xi _s(\xi ))=0`$.
The analogue of Eq.(VIII.3) reads
$$W(\xi )=W_{}(\xi _s(\xi )).$$
(VIII.17)
Such a density has infinite norm, since there are directions in the phase space $`\mathrm{\Gamma }^{2n}`$, crossing the submanifold $`𝒢_a=0`$, along which the density remains constant. These directions are determined by the constraint functions.
Eq.(V.2) tells
$$\{𝒢_a(\xi ),W(\xi )\}=0.$$
(VIII.18)
In the standard canonical frame, where $`q_A=\chi _A`$ and $`p_A=\mathrm{\Omega }_A`$, the gradient projections can be easily constructed:
$$\xi _s(\xi )=(0,q^{},0,p^{})$$
(VIII.19)
where $`\xi =(q_A,q^{},p_A,p^{})`$. Eq.(VIII.17) then simplifies to
$$W(q,p)=W_{}(q^{},p^{}).$$
(VIII.20)
Note that
$`\xi _u(\xi )`$ $`=`$ $`(0,q^{},p_A,p^{}),`$ (VIII.21)
$`\xi _v(\xi )`$ $`=`$ $`(q_A,q^{},0,p^{}),`$ (VIII.22)
so that $`\xi _s(\xi )=\xi _u(\xi _v(\xi ))=\xi _v(\xi _u(\xi ))`$.
If we apply the Wigner transform, we get the density matrix
$`\rho (q,q^{})`$ $`=`$ $`{\displaystyle \rho (\frac{q+q^{}}{2},p)e^{\frac{i}{\mathrm{}}p(qq^{})}\frac{d^np}{(2\pi \mathrm{})^n}}`$ (VIII.23)
$`=`$ $`{\displaystyle \underset{A}{}}\delta (q_Aq_A^{})\rho _{}(q^{},q^{}).`$
It satisfies
$`\widehat{q}_A\widehat{\rho }\widehat{\rho }\widehat{q}_A`$ $`=`$ $`0,`$ (VIII.24)
$`\widehat{p}_A\widehat{\rho }\widehat{\rho }\widehat{p}_A`$ $`=`$ $`0`$ (VIII.25)
or, equivalently,
$$\widehat{𝒢}_a\widehat{\rho }\widehat{\rho }\widehat{𝒢}_a=0.$$
(VIII.26)
The phase space analogue of these operator equations looks like
$$𝒢_a(\xi )W(\xi )=0.$$
(VIII.27)
This condition is in agreement with its classical counterpart Eq.(VIII.18) by virtue of Eq.(VIII.15).
Eqs.(VIII.26) are distinct from Eqs.(VIII.8). This is a consequence of different extrapolation schemes of the density to the unconstrained phase space. Eqs.(VIII.26) are the necessary and sufficient conditions for Eq.(VIII.23).
The dependence of the density matrix on ($`q_A,q_A^{}`$) does not factorize, although $`\rho ^2=\rho `$ for $`\rho _{}^2=\rho _{}`$. The latter, even when fulfilled, is not sufficient to have a pure state. Eq.(VIII.23) describes, in particular, a noncoherent sum of pure states with different momenta $`p_A`$. The system is thereby identified as a mixed state. It cannot be described by a wave function. The constraint equations (VIII.26), respectively, cannot be formulated in terms of a wave function $`\mathrm{\Psi }(q)`$.
If one works with density matrices in the representation (VIII.3) or (VIII.20), it is not important how the second-class constraints were split. Eqs.(VIII.8) and (VIII.26) are symmetric explicitly with respect to the interchange $`\chi _A\mathrm{\Omega }_A`$, $`\mathrm{\Omega }_A\chi _A`$, linear transformations of the constraint functions, and furthermore, with respect to unitary transformations in the Hilbert space. In the classical limit, they are invariant with respect to canonical transformations.
The normalization condition of the Wigner function involves a projection operator
$$𝔓=\frac{d^{2m}\lambda }{(2\pi \mathrm{})^m}\underset{a=1}{\overset{2m}{}}\mathrm{exp}(\frac{i}{\mathrm{}}\widehat{𝒢}^a\lambda _a)$$
(VIII.28)
in terms of which the norm is calculated as follows QDDB
$$P(\xi )W(\xi )\frac{d^{2n}\xi }{(2\pi \mathrm{})^n}=1.$$
(VIII.29)
where $`P(\xi )`$ is the Weyl’s symbol of the operator (VIII.28). Note that $`𝔓`$ is commutative with $`\widehat{\rho }`$ in virtue of (VIII.26).
The Wigner function appears as a smooth function, so it is an appropriate object to describe an evolution of the system on line with the Liouville equation in the unconstrained phase space. A quantum extension of the the Liouville equation for $`W(\xi )`$ satisfying Eqs.(VIII.27) is given in Ref.QDDB .
### VIII.3 Probability density with mixed localization
Let us keep, e.g., canonical momenta on the constraint submanifold and use the gradient projection to extend the density as a function of the canonical coordinates away from the constraint submanifold. The Wigner function is given then by
$$W(q,p)=(2\pi \mathrm{})^m\underset{A}{}\delta (p_A)W_{}(q^{},p^{}).$$
(VIII.30)
Applying the Wigner transform, we get the density matrix
$`\rho (q,q^{})=\rho _{}(q^{},q^{})`$ (VIII.31)
which satisfies
$`\widehat{p}_A\widehat{\rho }`$ $`=`$ $`0,`$ (VIII.32)
$`\widehat{\rho }\widehat{p}_A`$ $`=`$ $`0`$ (VIII.33)
or, equivalently,
$`\widehat{\mathrm{\Omega }}_A\widehat{\rho }`$ $`=`$ $`0,`$ (VIII.34)
$`\widehat{\rho }\widehat{\mathrm{\Omega }}_A`$ $`=`$ $`0.`$ (VIII.35)
These equations provide the necessary and sufficient conditions for Eq.(VIII.31). No dependence on the ($`q_A`$,$`q_A^{}`$) appears. If $`\rho _{}^2=\rho _{}`$ in the coordinate space $`q^{}`$ on the constraint submanifold then $`\rho ^2=\rho `$ in the coordinate space $`q`$ (owing to an infinite normalization factor).
The system is allowed to appear in a pure state and can be described by a wave function $`\mathrm{\Psi }(q)`$ accordingly. Eqs.(VIII.34) and (VIII.35) can be reformulated in terms of wave functions to match the Dirac’s supplementary condition Eq.(VIII.1).
The mixed localization scheme breaks the symmetry $`\chi _A\mathrm{\Omega }_A`$, $`\mathrm{\Omega }_A\chi _A`$ from the outset. However, it allows to work with wave functions. The other methods we discussed lead, in the unconstrained configuration space, to mixed states.
The quantum analogue of Eqs.(VIII.34) and (VIII.35) in the phase space can be formulated in terms of the star-product:
$`\mathrm{\Omega }_A(\xi )W(\xi )`$ $`=`$ $`0,`$ (VIII.36)
$`W(\xi )\mathrm{\Omega }_A(\xi )`$ $`=`$ $`0.`$ (VIII.37)
As a consequence of Eqs.(VIII.34) and (VIII.35), we have $`\{\mathrm{\Omega }_A,W\}=0`$ in the classical limit. Eq.(V.4) tells then that $`W(\xi )=\rho (\xi _u(\xi ))`$. The second classical condition $`\mathrm{\Omega }_AW=0`$ results in the Wigner function $`W(\xi )=\delta (\mathrm{\Omega }_A(\xi ))\rho _{}(\xi _v(\xi ))`$. Combining these two equations, we obtain
$$W(\xi )=(2\pi \mathrm{})^m\delta (\mathrm{\Omega }_A(\xi ))W_{}(\xi _s(\xi )),$$
(VIII.38)
which is in agreement with Eq.(VIII.30). We used here $`\mathrm{\Omega }_A(\xi _u(\xi ))=\mathrm{\Omega }_A(\xi )`$ which is valid up to the second order in the constraint functions $`\chi _A`$.
It is possible to establish a relationship with the results of the previous subsection. Indeed, one can check that
$$\widehat{\rho }=𝔓_\mathrm{\Omega }\widehat{\rho }_s,$$
(VIII.39)
with $`\widehat{\rho }_s`$ being the density matrix from the previous subsection, satisfies Eqs.(VIII.34) and (VIII.35). The operator $`𝔓_\mathrm{\Omega }`$ is defined by
$$𝔓_\mathrm{\Omega }=\frac{d^m\lambda }{(2\pi \mathrm{})^m}\underset{A=1}{\overset{m}{}}\mathrm{exp}(\frac{i}{\mathrm{}}\widehat{\mathrm{\Omega }}_A\lambda _A)$$
(VIII.40)
In order to get the normalization condition for $`W(\xi )`$ one has to factorize $`\widehat{\rho }`$ according to (VIII.39), construct the Weyl’s symbol of $`\widehat{\rho }_s`$, and apply Eq.(VIII.29).
Let us consider the opposite localization:
$$W(q,p)=\underset{A}{}\delta (q_A)W_{}(q^{},p^{}).$$
(VIII.41)
The first $`m`$ canonical coordinates are kept on the constraint submanifold, whereas the first $`m`$ canonical momenta are projected.
The density matrix has the form
$`\rho (q,q^{})={\displaystyle \underset{A}{}}\delta (q_A)\delta (q_A^{})\rho _{}(q^{},q^{})`$ (VIII.42)
It satisfies $`\widehat{q}_A\widehat{\rho }=\widehat{\rho }\widehat{q}_A=0`$ or, equivalently,
$`\widehat{\chi }_A\widehat{\rho }`$ $`=`$ $`0,`$ (VIII.43)
$`\widehat{\rho }\widehat{\chi }_A`$ $`=`$ $`0.`$ (VIII.44)
These conditions are the necessary and sufficient conditions to have the density matrix of the form (VIII.42).
The system is allowed to appear in a pure state and can be described by a wave function. The constraints (VIII.43) and (VIII.44) can be reformulated in terms of the wave functions to give Eq.(VIII.2).
As a consequence of Eqs.(VIII.43) and (VIII.44), one gets in the classical limit
$$W(\xi )=\delta (\chi _A(\xi ))W_{}(\xi _s(\xi )),$$
(VIII.45)
in agreement with Eq.(VIII.41). We used here the relation $`\chi _A(\xi _v(\xi ))=\chi _A(\xi )`$ which is valid up to the second order in $`\mathrm{\Omega }_A`$.
The quantum analogue of Eqs.(VIII.43) and (VIII.44) is given by
$`\chi _A(\xi )W(\xi )`$ $`=`$ $`0,`$ (VIII.46)
$`W(\xi )\chi _A(\xi )`$ $`=`$ $`0.`$ (VIII.47)
There exists a relationship with the Wigner function of the previous subsection. One can check that
$$\widehat{\rho }=𝔓_\chi \widehat{\rho }_s$$
(VIII.48)
satisfies Eqs.(VIII.43) and (VIII.44). The operator $`𝔓_\chi `$ is defined by
$$𝔓_\chi =d^m\lambda \underset{A=1}{\overset{m}{}}\mathrm{exp}(\frac{i}{\mathrm{}}\widehat{\chi }_A\lambda _A)$$
(VIII.49)
Note that $`𝔓=𝔓_\mathrm{\Omega }𝔓_\chi `$ and $`\widehat{\mathrm{\Omega }}_A𝔓_\mathrm{\Omega }=\widehat{\chi }_A𝔓_\chi =0`$ in the operator sense.
The normalization of the Wigner function in arbitrary canonical coordinate system is quite involved:
Given a wave function in the unconstrained configuration space, one should construct the density matrix, factorize it according to Eqs.(VIII.39) or (VIII.48) and integrate the Wigner function associated to $`\rho _s`$ in order to extract the norm according to Eq.(VIII.29).
Respectively, the Hermitian product of wave functions is calculated using the off-diagonal Wigner function and integrating it over the unconstrained phase space according to the same prescription.
### VIII.4 Discussion
In the classical constraint systems, physical quantities depend on the probability densities localized on the constraint submanifold only. In the standard canonical frame, the Wigner functions of quantum systems are localized on the constraint submanifolds also. This is mandatory, since any unconstrained system can be treated as a constrained system with constraints imposed to remove the added unphysical degrees of freedom. In doing so, the unphysical degrees of freedom should not modify dynamics of the initial system. This is achieved by attributing physical sense to the Wigner functions on the constraint submanifold. How to extrapolate Wigner functions from the constraint submanifold into the unconstrained phase space is a matter of convention.
We discussed the most evident extrapolations. Among them are those which allow to describe systems in the original phase space as pure states (mixed localization). One of them can be constructed using the Dirac’s prescription (VIII.1). The dual condition (VIII.2) was found to be possible also. If we work with density matrices or the Wigner functions, one arrives at conditions (VIII.8) or (VIII.18), both are symmetric under the permutations $`\chi _A\mathrm{\Omega }_A`$, $`\mathrm{\Omega }_A\chi _A`$. The constraints imposed on the physical states do not depend on the splitting $`𝒢_a=(\chi _A,\mathrm{\Omega }_A)`$.
## IX The $`O(n)`$ non-linear sigma model
The $`O(n)`$ non-linear sigma model represents the field theory analogue of the spherical $`n1`$-dimensional pendulum. The $`n=4`$ case corresponds to the chiral non-linear sigma model due to the isomorphism of algebras $`su(2)_Lsu(2)_Rso(4)`$.
In Sect. 2, we started from the tangent bundle $`T=(\varphi ^\alpha ,\dot{\varphi }^\alpha )`$ of the dimension $`2n`$, defined over the configuration space $`=(\varphi ^\alpha )`$. Lagrangian (II.2) depends on $`n`$ velocities $`\dot{\varphi }^\alpha `$. On the constraint submanifold (II.1), it depends on $`n1`$ velocities $`\mathrm{\Delta }^{\alpha \beta }\dot{\varphi }^\beta `$ tangent to the constraint submanifold. The dynamics of $`\varphi `$ turns out to be independent on other variables. The constraint (II.1) reduces the effective number of the degrees of freedom to $`n1`$. Thus, a $`2n2`$ dimensional tangent bundle over the $`n1`$ configuration space can be constructed. It is described, e.g., by coordinates $`\vartheta ^i`$ analogous to the angular coordinates in three dimensional space. Lagrangian (II.11) is not degenerate with respect to $`\dot{\vartheta }^i`$. The coordinates $`\vartheta ^i`$ constitute the physical configuration space.
The path integral for the evolution operator in terms of the angular coordinates $`\vartheta ^i`$ can be treated as a reference point for comparison to more involved quantization methods.
### IX.1 Path integral for the $`O(n)`$ non-linear sigma model
Let us construct the path integral for the $`O(n)`$ non-linear sigma model in terms of the angular variables $`\vartheta ^i`$, i.e., by solving the constraint equations from the outset and compare it with the expression of Sect. 2 derived using the underlying gauge symmetry of the $`O(n)`$ non-linear sigma model.
The field theory extension of the spherical pendulum problem is rather straightforward. In what follows, the kinematic variables are functions of $`x_\mu =(t,𝐱)`$.
Let us construct a basis
$`e_i^\alpha `$ $`=`$ $`{\displaystyle \frac{}{\vartheta ^i}}\varphi ^\alpha ,`$ (IX.1)
$`e_i^\alpha e_j^\alpha `$ $`=`$ $`g_{ij},`$ (IX.2)
$`g^{ij}e_i^\alpha e_j^\beta `$ $`=`$ $`\delta ^{\alpha \beta }\varphi ^\alpha \varphi ^\beta /\varphi ^2,`$ (IX.3)
where $`g_{ij}=g_{ij}(\vartheta )`$ is an induced metric tensor on the submanifold $`\varphi (x)=1`$, $`detg_{ij}0`$. In terms of the coordinates $`\vartheta ^i`$, the field theory extension of Lagrangian (II.11) takes the form
$$_{}=\frac{1}{2}g_{ij}_\mu \vartheta ^i_\mu \vartheta ^j.$$
(IX.4)
The Legendre transformation of $`_{}`$ is well defined. In term of the canonical momenta
$$\varrho _i=\frac{_{}}{\dot{\vartheta }^i}=g_{ij}\dot{\vartheta }^j$$
(IX.5)
the Hamiltonian density can be found to be
$$_{}=\frac{1}{2}g^{ij}\varrho _i\varrho _j+\frac{1}{2}g_{ij}_a\vartheta ^i_a\vartheta ^j$$
(IX.6)
where $`a=1,2,3`$. The non-vanishing Poisson bracket for the canonical coordinates and momenta has the form
$$\{\vartheta ^i(t,𝐱),\varrho _j(t,𝐱^{})\}=\delta _j^i\delta (𝐱𝐱^{}).$$
(IX.7)
The Poisson bracket relations for coordinates $`\varphi ^\alpha `$ and momenta $`\pi ^\alpha `$ associated to the tangent velocities $`\pi _{}^\alpha =\varrho _ig^{ij}e_j^\alpha `$ agree with those discussed in Sect. 4.
The path integral in the phase space $`(\vartheta ^i,\varrho _i)`$ is given by
$$Z=\frac{d\vartheta ^id\varrho _i}{(2\pi \mathrm{})^{n1}}\mathrm{exp}\left\{\frac{i}{\mathrm{}}d^4x(\varrho ^i\dot{\vartheta }_i_{})\right\}.$$
(IX.8)
The Liouville measure $`d\vartheta ^id\varrho _i`$ is consistent with Eq.(IX.7). The integral over the canonical momenta in Eq.(IX.8) has a Gaussian form and can be calculated explicitly:
$$Z=\sqrt{detg_{ij}}d^{n1}\vartheta \mathrm{exp}\left\{\frac{i}{\mathrm{}}𝑑x^4_{}\right\}.$$
(IX.9)
The value $`\sqrt{detg_{ij}}d^{n1}\vartheta `$ gives volume of the configuration space, defined by the metric tensor $`g_{ij}`$. This measure is invariant under the $`O(n)`$ group.
The $`S`$-matrix (IX.9) can be written in an explicitly covariant form with respect to the $`O(n)`$ rotations and the Lorentz transformations. First, we rewrite Lagrangian density (IX.4) in terms of the coordinates $`\varphi ^\alpha `$
$$_{}=\frac{1}{2}\mathrm{\Delta }^{\alpha \beta }(\varphi )_\mu \varphi ^\alpha _\mu \varphi ^\beta /\varphi ^2$$
(IX.10)
and, second, rewrite the Lagrange measure
$$\sqrt{detg_{ij}}d^{n1}\vartheta =\sqrt{\left(\chi /\varphi ^\alpha \right)^2}\delta (\chi )d^n\varphi $$
(IX.11)
where $`\chi =\mathrm{ln}\varphi `$. The right side is the same for all functions vanishing at $`\varphi =1`$.
In this form we recover the result of Sect. 2.
### IX.2 Pion field parameterization in chiral sigma model
The $`n=4`$ case is especially interesting since it corresponds to the chiral non-linear sigma model. For $`n=4`$, the angular coordinates are defined by
$$\varphi ^\alpha =(\mathrm{cos}\psi ,\mathrm{sin}\psi \times (\mathrm{cos}\theta ,\mathrm{sin}\theta \times (\mathrm{cos}\phi ,\mathrm{sin}\phi )))$$
(IX.12)
where $`\vartheta ^1=\psi `$, $`\vartheta ^2=\theta `$, and $`\vartheta ^3=\phi `$.
The angular distance, $`\mathrm{\Theta }`$, between two vectors $`\varphi ^\alpha `$ and $`\varphi ^\alpha `$ is defined by scalar product $`\mathrm{cos}\mathrm{\Theta }=\varphi \varphi ^{}`$. The distance element becomes
$$d\mathrm{\Theta }^2=d\psi ^2+\mathrm{sin}^2\psi (d\theta ^2+\mathrm{sin}^2\theta d\phi ^2).$$
(IX.13)
The components of the metric $`g_{ij}`$ can be found using the expansion $`d\mathrm{\Theta }^2=g_{ij}d\vartheta ^id\vartheta ^j`$ or directly from Eqs.(IX.1) and (IX.2). The Lagrange measure of the path integral becomes
$`\sqrt{detg_{ij}}d^3\vartheta `$ $`=`$ $`\mathrm{sin}^2\psi \mathrm{sin}\theta d\psi d\theta d\phi `$ (IX.14)
$`=`$ $`{\displaystyle \frac{\mathrm{sin}^2\psi }{\psi ^2}}dV,`$
with $`dV=\psi ^2d\psi \mathrm{sin}\theta d\theta d\phi `$ being an element of the Euclidean volume. The corresponding Lagrangian can be found from Eq.(IX.4) to give
$$_{}=\frac{1}{2}(_\mu \psi )^2+\frac{\mathrm{sin}^2\psi }{2}((_\mu \theta )^2+\mathrm{sin}^2\theta (_\mu \phi )^2).$$
(IX.15)
The quantization of the chiral sigma model is made using an oscillator basis by expanding the $`_{}`$ around $`\vartheta ^i=0`$ breaking thereby the $`O(4)`$ symmetry down to its $`O(3)`$ subgroup. It can be successful provided that measure of the coordinate space is such that $`detg_{ij}=1`$. The path integrals convert then to the Gaussian integrals which can be calculated. The parameterization preserving the $`O(3)`$ symmetry and satisfying the above requirement is, apparently, unique. One should rescale the ”radius” $`\psi `$ according to
$$\mathrm{sin}^2\psi d\psi =\omega ^2d\omega .$$
(IX.16)
This elementary equation gives
$$\omega =\left(\frac{3}{2}(\psi \mathrm{sin}\psi \mathrm{cos}\psi )\right)^{1/3}.$$
(IX.17)
Lagrangian $`_{}`$ then becomes
$$_{}=\frac{\omega ^4}{2\mathrm{sin}^4\psi }(_\mu \omega )^2+\frac{\mathrm{sin}^2\psi }{2}((_\mu \theta )^2+\mathrm{sin}^2\theta (_\mu \phi )^2)$$
(IX.18)
where $`\psi `$ is a function of $`\omega `$ Eq.(IX.17). The mass term breaking the $`O(4)`$ symmetry looks like
$$_M=M^2(\mathrm{cos}\psi 1).$$
(IX.19)
The quadratic part of the Lagrangian used for the perturbation expansion can be selected as follows
$$_{}^{[2]}=\frac{1}{2}(_\mu \omega )^2+\frac{\omega ^2}{2}((_\mu \theta )^2+\mathrm{sin}^2\theta (_\mu \phi )^2)\frac{M^2}{2}\omega ^2.$$
In terms of the pion fields
$$\pi ^a=\omega \times (\mathrm{cos}\theta ,\mathrm{sin}\theta \times (\mathrm{cos}\phi ,\mathrm{sin}\phi )),$$
(IX.20)
it takes the standard form
$$_{}^{[2]}=\frac{1}{2}(_\mu \pi ^a)^2\frac{M^2}{2}(\pi ^a)^2,$$
(IX.21)
whereas the Lagrange measure is simply the Euclidean volume $`d^3\pi `$. The difference $`\delta _{}=_{}+_M_{}^{[2]}`$ can be considered as a perturbation.
The pion fields can be parameterized in various ways. The problem of ambiguities of the transition amplitudes, connected to the arbitrariness of that choice, was discussed first in Refs.CHIK ; KAME . It was shown that on-shell amplitudes do not depend on the choice of physical variables. This statement is known as the ”equivalence theorem”. The method proposed by Gasser and Leutwyller GALE associates the QCD Green functions to amplitudes of the effective chiral Lagrangian. Using this method, the QCD on- and off-shell amplitudes can be calculated in a way independent on the parameterization.
The $`S`$-matrix is invariant with respect to a symmetry group, if both the action functional and the Lagrange measure entering the path integral over canonical coordinates are invariant. The Liouville measure entering the path integral over canonical coordinates and momenta is always ”flat”, since we work in canonical basis. The quantization gives a non-trivial Lagrange measure, however (cf. Eq.(IX.9)). In case of the $`O(n)`$ non-linear sigma model, there is only one parameterization which makes the Lagrange measure flat, i.e., $`detg_{ij}=1`$. This requirement is useful for development of perturbation theory which uses an oscillator basis to convert path integrals into the Gaussian form.
The weight factor can always be exponentiated to generate an effective Lagrangian $`\delta _H`$, in which case $`\chi ^\alpha =\varphi ^\alpha `$ provides desired parameterization also. The exponential parameterization of the pion matrix
$$U(\varphi ^\alpha )=e^{i\tau ^\alpha \varphi ^\alpha }$$
(IX.22)
gives, in particular,
$$\delta _H=\frac{1}{a^4}\mathrm{ln}(\mathrm{sin}^2(\varphi )\varphi ^2)$$
(IX.23)
where $`a`$ is a lattice size. $`\delta _H`$ diverges in the continuum limit. The non-linear sigma model is not a renormalizable theory, so divergences cannot be absorbed into redefinition of $`F`$ and $`M_\pi `$. Using the mean filed (MF) approximation, it is usually possible to keep renormalizations finite. The exponentiation of a variable weight factor breaks, in general, selfconsistency of the MF approximation of the non-linear sigma model.
Divergences arising from $`\delta _H`$ could, however, be compensated by divergences coming from higher orders ChPT loops. From this point of view it looks naturally to attribute $`\delta _H`$ to higher orders ChPT loop expansion starting from one loop. The MF approximation for ChPT implies then the tree level approximation for the non-linear sigma model with $`\delta _H`$ neglected. Such an approximation, however, neglects the Haar measure from the start. It is therefore hard to expect that such an approximation describes correctly the high temperature regime where the chiral invariance is supposed to be restored.
The self-consistency of the MF approximation of the non-linear sigma model survives with the one parameterization only.
The invariance under the chiral transformations admits, within the MF approximation, only the parameterization given by Eq.(IX.17).
The parameterization based on the dilatation of $`\varphi ^\alpha `$ gives $`\delta _H=0`$, does not involve the higher orders ChPT loops, and allows to work in the continuum limit with finite quantities only.
The effective interaction terms in the Lagrangian which appear due to the presence of the Haar measure have been discussed earlier in QCD (see PARA and references therein).
The $`SU(2)`$ group has a finite group volume. The integration range of $`\omega `$ fields is therefore be restricted. According to the current paradigm, one can extend the integrals over $`\omega `$ from $`\mathrm{}`$ to $`+\mathrm{}`$ within a perturbation theory framework. The modification of the result is connected to the integration over large field fluctuations, so the variance has, apparently, a non-perturbative nature and does thereby not affect the perturbation series. As a matter of fact, this justifies the standard loop expansion in ChPT.
## X Summary
In this work, we discussed analogy between the second-class constraints systems and gauge theories with the equivalent structure of gauge generators and gauge-fixing conditions. Given the symplectic basis for the constraint functions exists globally, the second-class constraints systems can be interpreted as gauge invariant systems in the unconstrained phase space. Such systems can be quantized using the methods specific for gauge theories.
The second-class constraints $`𝒢_a`$ split in the symplectic basis into canonical pairs ($`\chi _A`$,$`\mathrm{\Omega }_B`$) satisfying $`\{\chi _A,\chi _B\}0`$, $`\{\mathrm{\Omega }_A,\mathrm{\Omega }_B\}0`$, and $`\{\chi _A,\mathrm{\Omega }_B\}\delta _{AB}`$. The constraint functions ($`\chi _A`$,$`\mathrm{\Omega }_B`$) can be transformed further, as discussed in Sects. 4 and 5, to fulfill the Poisson bracket relations in the strong sense in an entire neighborhood of any give point of the constraint submanifold. The Hamiltonian function can also be modified to be identically in involution with the constraints. The new constraints define the same constraint submanifold, whereas the new Hamiltonian and its first derivatives coincide with the original ones on the constraint submanifold. The constrained dynamics is thus not modified. The new constraint functions $`\chi _A`$ and $`\mathrm{\Omega }_A`$ are interpreted as gauge-fixing conditions and first-class constraints associated to gauge transformations.
We do not provide any criterion for what part of the constraints $`𝒢_a`$ describes the gauge-fixing conditions and what part describes the first-class constraints associated to gauge transformations. By contrary, we argue that transition amplitudes of the quantum theory do not depend on the interpretation of $`𝒢_a`$.
The Dirac’s supplementary conditions $`\widehat{\mathrm{\Omega }}_A\mathrm{\Psi }=0`$ depend on the way the constrains $`𝒢_a`$ were split. These conditions are equivalent to $`\widehat{\chi }_A\mathrm{\Psi }^{}=0`$, since the corresponding Wigner functions coincide on the constraint submanifold. The supplementary conditions for the Wigner functions, furthermore, can be made to be explicitly invariant with respect to possible transformations of the constraint functions. The ambiguity reflects the freedom in extrapolation of the Wigner functions from the constraint submanifold into the unconstrained phase space.
We showed, finally, that the proposed quantization scheme applies to an $`n1`$-dimensional spherical pendulum, which represents a mechanical version of the $`O(n)`$ non-linear sigma model. For this model, we demonstrated the existence of an underlying gauge symmetry which is the dilatation of the coordinates $`\varphi ^\alpha `$. The constraints appearing within the Hamiltonian framework are of the second class. If one starts, however, from an equivalent Lagrangian in which the underlying gauge symmetry is set up explicitly, the same constraints appear as a gauge-fixing condition and a constraint associated to the dilatation symmetry. It shows that the interpretation of second-class constraints is a matter of convention.
For holonomic systems, the quantization method based on construction of the gauged model does requires neither auxiliary canonical variables nor extended configuration space.
For second-class constraints systems, the underlying gauge symmetries induce in the configuration space transformations depending on velocities and involve auxiliary variables. The holonomic systems admit the natural gauged counterparts in the original configuration space.
The initial configuration space is, in general, too narrow to reflect the gauge invariance of a system described in terms of the generalized Hamiltonian dynamics. Gauge invariant quantities involving auxiliary variables, furthermore, do not belong to the set of physical observables, as distinct from the usual gauge theories. The equivalence of the first-class constraints systems with the ordinary gauge systems is therefore physically not complete and restricted to systems of point particles under holonomic constraints, non-holonomic constraints satisfying the Frobenius’ condition, and systems with one primary constraint only.
The path integral representation for the evolution operator of the $`O(n)`$ non-linear sigma model was constructed in Sect. 9 by solving the constraint equations. The equivalence with the quantization methods based on the reduction to the equivalent gauge systems was demonstrated.
After finishing this work we got to know about works NAKA84 ; NAKA89 ; NAKA93 ; NAKA01 where projection formalism is discussed.
###### Acknowledgements.
M.I.K. and A.A.R. wish to acknowledge kind hospitality at the University of Tübingen. This work is supported by DFG grant No. 436 RUS 113/721/0-2 and RFBR grant No. 06-02-04004.
## Appendix A Two-constraints form of gauged spherical pendulum
Let us consider another example. The system discussed in Sect. 3 is equivalent to a system described by a vector $`\varphi ^i=(\varphi ^0,\varphi ^\alpha )`$ on the ”light-cone” submanifold $`\varphi ^2=\varphi ^0\varphi ^0\varphi ^\alpha \varphi ^\alpha =0`$. The constraint $`\varphi ^\alpha \varphi ^\alpha =1`$ is equivalent to the constraint $`\varphi ^0=1`$, while the transformations (II.12) are equivalent to Lorentz boosts of the null vector $`\varphi ^i`$ along the vector $`\varphi ^\alpha `$,
$$\varphi ^i\varphi ^i=\mathrm{\Lambda }_j^i(\theta )\varphi ^j=\mathrm{exp}(\theta )\varphi ^i$$
(A.1)
where the ”boost velocity” $`v=\mathrm{tanh}(\theta )`$. The gauge invariance of the $`_{}`$ with respect to the dilatation reflects gauge invariance of the system with respect to the Lorentz boosts.
The ”light-cone” Lagrangian $`_2`$ can be constructed by considering the requirement of the conditional maximum of the action
$$\underset{\varphi ^\alpha ,\dot{\varphi }^\alpha }{\mathrm{max}}\{_{}𝑑t\}|_{\varphi ^\alpha \varphi ^\alpha =1}=\underset{\varphi ^i,\dot{\varphi }^i}{\mathrm{max}}\{_2𝑑t\}|_{\varphi ^2=0,\varphi ^0=1}.$$
(A.2)
If the $`\varphi ^\alpha \varphi ^\alpha `$ is treated as a gauge parameter, the problem simplifies. It is sufficient to require
$$\underset{\varphi ^\alpha ,\dot{\varphi }^\alpha }{\mathrm{max}}\{_{}𝑑t\}=\underset{\varphi ^i,\dot{\varphi }^i}{\mathrm{max}}\{_2𝑑t\}|_{\varphi ^2=0}.$$
(A.3)
The $`_2`$ can be chosen as a straightforward extension of $`_{}`$:
$$_2=\frac{1}{2}G_{ij}(\varphi )\dot{\varphi }^i\dot{\varphi }^j/((\eta \varphi )^2\varphi ^2)$$
(A.4)
where $`\eta =(1,0,\mathrm{},0)`$,
$$G_{ij}(\varphi )=g_{ij}\frac{(\eta \varphi )(\eta _i\varphi _j+\eta _j\varphi _i)\varphi ^2\eta _i\eta _j\varphi _i\varphi _j}{(\eta \varphi )^2\varphi ^2},$$
(A.5)
and $`g_{ij}=`$ diag$`(1,1,\mathrm{},1)`$, $`g^{ij}=g_{ij}`$.
The tensor $`G_{ij}`$ obeys
$`\eta ^iG_{ij}(\varphi )`$ $`=`$ $`\varphi ^iG_{ij}(\varphi )=0,`$ (A.6)
$`G_{ij}(\varphi )G_k^j(\varphi )`$ $`=`$ $`G_{ik}(\varphi ).`$ (A.7)
It is invariant with respect to the dilatation (A.1) and the shifts
$$\varphi ^i\varphi ^i=\varphi ^i+ϵ\eta ^i$$
(A.8)
where $`ϵ`$ is an arbitrary parameter.
Lagrangian (A.4) is well defined for $`\varphi ^20`$. It is invariant with respect to the dilatation (A.4) and the shifts (A.8).
The $`\varphi ^2`$ and the $`\eta \varphi `$ are thus gauge functions. They are not fixed by equations of motion and can be selected to satisfy admissible constraints. We consider therefore Lagrangian (A.4) without imposing any constraints. The initial and final conditions $`\varphi ^2=0`$ and $`\eta \varphi =1`$ are gauge-fixing conditions.
The canonical momenta corresponding to the $`\dot{\varphi }^i`$ are defined by
$$\pi _i=\frac{_c}{\dot{\varphi }^i}=G_{ij}(\varphi )\dot{\varphi }^j/((\eta \varphi )^2\varphi ^2).$$
(A.9)
They satisfy the primary constraints
$$\pi _iG_{ij}\pi ^j0$$
(A.10)
which are equivalent to two ones:
$`\mathrm{\Omega }_1`$ $`=`$ $`\varphi \pi 0,`$ (A.11)
$`\mathrm{\Omega }_2`$ $`=`$ $`\eta \pi 0.`$ (A.12)
The primary Hamiltonian can be obtained with the use of Legendre transformation:
$$=\frac{1}{2}((\eta \varphi )^2\varphi ^2)G^{ij}(\varphi )\pi _i\pi _j.$$
(A.13)
The Poisson bracket for canonical coordinates and momenta have the form
$$\{\varphi _i,\pi _j\}=g_{ij}.$$
(A.14)
The primary constraints are stable with respect to the time evolution:
$`\{\mathrm{\Omega }_1,\}`$ $`=`$ $`0,`$ (A.15)
$`\{\mathrm{\Omega }_2,\}`$ $`=`$ $`0.`$ (A.16)
The Hamiltonian $``$ is gauge invariant. The primary constraints are of the first class:
$$\{\mathrm{\Omega }_1,\mathrm{\Omega }_2\}=\mathrm{\Omega }_2.$$
(A.17)
The generators of gauge transformations constitute an algebra.
The relations
$`\{\varphi ^i,\mathrm{\Omega }_1\}`$ $`=`$ $`\varphi ^i,\{\varphi ^i,\mathrm{\Omega }_2\}=\eta ^i,`$ (A.18)
$`\{\pi _i,\mathrm{\Omega }_1\}`$ $`=`$ $`\pi _i,\{\pi _i,\mathrm{\Omega }_2\}=0`$ (A.19)
show that the $`\mathrm{\Omega }_1`$ generates the dilatation of the $`\varphi ^i`$ and the $`\pi _i`$, while the $`\mathrm{\Omega }_2`$ generates time-like shifts of $`\varphi ^i`$.
The gauge-fixing conditions
$`\chi _1`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{ln}((\eta \varphi )^2\varphi ^2),`$ (A.20)
$`\chi _2`$ $`=`$ $`\eta \varphi 1`$ (A.21)
generate the following transformations:
$`\{\varphi ^i,\chi _1\}=0,\{\varphi ^i,\chi _2\}=0,`$ (A.22)
$`\{\pi ^i,\chi _1\}={\displaystyle \frac{\varphi ^i\eta ^i\eta \varphi }{(\eta \varphi )^2\varphi ^2}},\{\pi ^i,\chi _2\}=\eta ^i.`$ (A.23)
They are identically in involution with the Hamiltonian:
$`\{\chi _1,\}`$ $`=`$ $`0,`$ (A.24)
$`\{\chi _2,\}`$ $`=`$ $`0.`$ (A.25)
The equations of motion generated by the primary Hamiltonian look like
$`\dot{\varphi }^i=\{\varphi ^i,\}=((\eta \varphi )^2\varphi ^2)G^{ij}(\varphi )\pi _j,`$ (A.26)
$`\dot{\pi }_i=\{\pi _i,\}(\varphi _i\eta _i(\eta \varphi ))G^{jk}(\varphi )\pi _j\pi _k.`$ (A.27)
The main inference is that starting from the different Lagrangian, an equivalent first-class constraints system was constructed. There is no principal distinction between the system described here and the one discussed in Sect. 3. In particular, one can solve the constraints $`\chi _2=0`$ and $`\mathrm{\Omega }_2=0`$ to remove the canonically conjugate pair ($`\varphi ^0`$,$`\pi _0`$) from the Hamiltonian. The system of Sect. 3 would be reproduced then explicitly.
## Appendix B Symplectic basis for first-class constraints
In Sect. 4, an equivalent system of second-class constraints satisfying involution relations Eqs.(IV.23) in the strong sense in an entire neighborhood of a given point of the constraint submanifold has been constructed. The existence of an equivalent Hamiltonian identically in involution with the new constraints has also been demonstrated Eq.(IV.24). The equivalence means that the constraint submanifolds and the phase space flows on the constraint submanifolds of the dynamical systems coincide.
Similar statements for first-class constraints systems are proved in Refs. SCHO49 ; HENN . The arguments of Refs.SCHO49 ; HENN cannot be extended to second-class constraints without additional assumptions. Let us discuss the restrictions.
To replace the constraint functions $`\mathrm{\Omega }_A`$ ($`A=1,\mathrm{},m`$) by equivalent constraint functions $`\stackrel{~}{\mathrm{\Omega }}_A`$ which are identically in involution $`\{\stackrel{~}{\mathrm{\Omega }}_A,\stackrel{~}{\mathrm{\Omega }}_B\}=0`$, one can solve equations $`\mathrm{\Omega }_A=0`$ with respect to first $`m`$ canonical momenta $`p_A=P_A`$ where $`P_A`$ are functions of the $`n`$ canonical coordinates and the remaining $`nm`$ canonical momenta. To have the constraints $`\stackrel{~}{\mathrm{\Omega }}_A=0`$ resolved, one has to require
$$det\frac{\stackrel{~}{\mathrm{\Omega }}_A}{p_B}0.$$
(B.1)
The phase space has a dimension $`2n`$, $`n>m`$. The functions $`\stackrel{~}{\mathrm{\Omega }}_A=p_AP_A`$ vanish on the submanifold $`\mathrm{\Omega }_A=0`$ only. The Poisson bracket $`\{\stackrel{~}{\mathrm{\Omega }}_A,\stackrel{~}{\mathrm{\Omega }}_B\}`$ vanishes weakly and does not depend on the first $`m`$ canonical momenta, therefore it vanishes identically.
The same prescription can be used to construct the functions $`\chi _A`$. The Poisson bracket $`\{\stackrel{~}{\chi }_A,\stackrel{~}{\chi }_B\}`$ vanishes identically. The new constraint functions satisfy $`det\{\stackrel{~}{\chi }_A,\stackrel{~}{\mathrm{\Omega }}_B\}0`$.
One can define an equivalent Hamiltonian $`\stackrel{~}{}`$. Let us substitute $`p_A=P_A`$ to the original Hamiltonian $`^{}`$ of the second-class constraints system. The resulting Hamiltonian $`\stackrel{~}{}`$ is first class with respect to the constraints $`𝒢_a=0`$, so $`\{\stackrel{~}{\mathrm{\Omega }}_A,\stackrel{~}{}\}0`$. The difference $`\stackrel{~}{}^{}`$ vanishes on the constraint submanifold $`\mathrm{\Omega }_A=0`$. The Hamiltonian $`\stackrel{~}{}`$ does not depend on the first $`m`$ canonical momenta, so the Poisson bracket does not depend on these canonical momenta either. The $`\stackrel{~}{}`$ is therefore identically in involution with the $`\stackrel{~}{\mathrm{\Omega }}_A`$. The similar procedure can be applied for the $`\stackrel{~}{\chi }_A`$.
We should get finally constraint functions $`\stackrel{~}{\chi }_A`$ and $`\stackrel{~}{\mathrm{\Omega }}_A`$ identically in involution with the $`\stackrel{~}{}`$.
The above arguments do not apply if some constraint functions do not depend on $`p_A`$ ($`q_A`$). The bracket $`\{\stackrel{~}{\mathrm{\Omega }}_A,\stackrel{~}{\mathrm{\Omega }}_B\}`$ can, e.g., be proportional to $`\chi _C`$ and $`\chi _C`$ can in turn be independent on $`p_A`$. In such a case, the weak equation $`\{\stackrel{~}{\mathrm{\Omega }}_A,\stackrel{~}{\mathrm{\Omega }}_B\}0`$ does not convert into the strong one, although the both sides do not depend on $`p_A`$. The similar restrictions appear in the construction of $`\stackrel{~}{\chi }_A`$ and $`\stackrel{~}{}`$. The systems under holonomic constraints have, in particular, constraints $`\chi _A=0`$ which do not depend on the canonical momenta. |
warning/0506/math-ph0506043.html | ar5iv | text | # Decomposition rules for conformal pairs associated to symmetric spaces and abelian subalgebras of ℤ₂-graded Lie algebras
## 1 Introduction
A pair $`(𝔰,𝔨)`$, where $`𝔰`$ is a finite-dimensional semisimple Lie algebra over $``$ and $`𝔨`$ is a reductive subalgebra of $`𝔰`$, such that the restriction of the Killing form of $`𝔰`$ to $`𝔨`$ is non-degenerate, is called a *conformal pair* if there exists an integrable highest weight module $`V`$ over the affine Kac–Moody algebra $`\widehat{𝔰}`$, faithful on each simple component of $`𝔰`$, such that the restriction to $`\widehat{𝔨}`$ of each weight space of the center of $`𝔨`$ in $`V`$ decomposes into a finite direct sum of irreducible $`\widehat{𝔨}`$-modules. In such a case $`𝔨`$ is called a *conformal subalgebra* of $`𝔰`$.
It is well-known that any integrable highest weight $`\widehat{𝔰}`$-module, when restricted to $`\widehat{𝔨}`$, decomposes into a direct sum of irreducible $`\widehat{𝔨}`$-modules , but almost always this decomposition is infinite.
The first cases of a finite decomposition, found in , are as follows. Let $`𝔤=𝔨𝔭`$ be the eigenspace decomposition of an inner involution of a simple Lie algebra $`𝔤`$ such that $`𝔨`$ is semisimple. This defines an embedding $`𝔨so(𝔭)`$. It was shown by Kac and Peterson in , by an explicit decomposition formula, that the restriction of the spinor representations of $`\widehat{so(𝔭)}`$ to $`\widehat{𝔨}`$ is a finite direct sum of irreducible $`\widehat{𝔨}`$-modules. Thus, the pair $`(so(𝔭),𝔨)`$ is conformal.
Due to their importance for string compactifications, a series of papers on conformal pairs appeared in the second half of the 1980s in physics literature. First of all, a connection to representation theory of the Virasoro algebra was established. Namely, it was found that the decomposition in question is finite if and only if the following numerical criterion holds: the central charges of the Sugawara construction of the Virasoro algebra for $`\widehat{𝔰}`$ and $`\widehat{𝔨}`$ are equal . This immediately has led to a conclusion: the decomposition in question has a chance to be finite only if the level of the $`\widehat{𝔰}`$-module $`V`$ is equal to $`1`$, and if it is finite for one of the $`\widehat{𝔰}`$-modules of level $`1`$, it is also finite for all others. Furthermore, Goddard, Nahm and Olive show that the observation of Kac and Peterson can be reversed. Namely all conformal pairs $`(so(𝔭),𝔨)`$ are obtained from an involution (not necessarily inner) of a semisimple Lie algebra $`𝔤`$, and all such pairs are conformal. However, they obtain this result using the above numerical criterion, and do not find actual decompositions.
All conformal subalgebras $`𝔨`$ for all simple Lie algebras $`𝔰`$ were classified in and by making use of the numerical criterion. Also, it was pointed out in that, using the conformal pairs $`(so_{2n},g\mathrm{}_n)`$ and $`(so_{4n},sp_{2n}\times s\mathrm{}_2)`$, one can reduce the study of conformal subalgebras in all classical Lie algebras to that in $`so_n`$.
Around the same time the general problem of restricting representations of affine Lie algebras to their subalgebras was treated, using the theory of modular forms. Namely it was observed in that the branching rules are described by certain modular functions, called branching functions, for which one can write down explicit transformation formulas. This idea was further developed in , where the above mentioned “modular constraints”, along with the “conformal constraints”, provided by the Virasoro algebra, allowed to compute easily branching functions (which are constants in the conformal pair case) in many interesting cases, and, in principle, in any given case. The technology, developed in was subsequently used in to find all the decompositions of all integrable highest weight modules of level $`1`$ over affine Lie algebras $`\widehat{𝔰}`$, restricted to affine subalgebras $`\widehat{𝔨}`$, where $`𝔨`$ is a conformal subalgebra of a simple exceptional Lie algebra $`𝔰`$. The branching rules of some other conformal embeddings were subsequently found in , and a few other papers, written in the 1990s.
The problem of finding a general conceptual formula for branching rules for level $`1`$ integrable highest weight modules over $`\widehat{so(𝔭)}`$, when restricted to $`\widehat{𝔨}`$, where the conformal embedding of $`𝔨`$ in $`so(𝔭)`$ is defined by the eigenspace decomposition of an involution of a semisimple Lie algebra $`𝔤=𝔨+𝔭`$, has remained an open problem.
In the present paper we completely solve this problem. The solution turned out to be intimately related to recent developments in the study of abelian subalgebras of simple Lie algebras, that began with a paper of Kostant and continued in , , , , , , , .
Let us explain our main observation on the example of a conformal embedding $`𝔨so(𝔨)`$, where $`𝔨`$ is a simple Lie algebra, via the adjoint representation of $`𝔨`$. In this case the restriction to $`\widehat{𝔨}`$ of the basic $`+`$ vector representations of $`\widehat{so(𝔨)}`$ decomposes into a direct sum of $`2^{\mathrm{rank}𝔨}`$ irreducible $`\widehat{𝔨}`$-modules (this decomposition was found already by Kac and Wakimoto ), and, remarkably, these $`2^{\mathrm{rank}𝔨}`$ $`\widehat{𝔨}`$-modules are in a canonical one-to-one correspondence with all abelian ideals of a Borel subalgebra of $`𝔨`$.
Our main result is that for all conformal pairs associated to symmetric spaces, the decomposition of level $`1`$ modules over $`\widehat{so(𝔭)}`$ is described in terms of a certain class of abelian subalgebras of semisimple $`_2`$-graded Lie algebras $`𝔤=𝔨𝔭`$, studied and classified recently by Cellini, Möseneder Frajria and Papi .
We hope that the connection of representation theory of affine Lie algebras to the theory of abelian subalgebras of simple Lie algebras will shed a new light on the latter theory as well. So far we obtained only partial results in this direction. Now we describe our results in a special case which might give the flavour of the general case. First remark that $`\widehat{so(𝔭)}`$ is an affine algebra of type $`B^{(1)}`$ or $`D^{(1)}`$ according to whether $`dim(𝔭)`$ is odd or even. Hence the level $`1`$ modules are the fundamental representations associated to the extremal nodes of the Dynkin diagram. They are the basic, vector and spin representations, and have been studied since a long time.
We consider in detail the case of the basic and vector representation of $`\widehat{𝔨}`$ and furthermore we assume that $`𝔨`$ is semisimple. Denote by $`L(\stackrel{~}{\mathrm{\Lambda }}_0)`$ the basic representation, by $`L(\stackrel{~}{\mathrm{\Lambda }}_1)`$ the vector representation and set $`L=L(\stackrel{~}{\mathrm{\Lambda }}_0)+L(\stackrel{~}{\mathrm{\Lambda }}_1)`$. The first step in our analysis consists in calculating the character of the $`\widehat{𝔨}`$-module $`L`$. This is done using the explicit description of the action given in . To be more precise we need to fix some notation. Let $`𝔥_0`$ be a Cartan subalgebra of $`𝔨`$; denote by $`\mathrm{\Delta }_𝔨`$ the set of $`𝔥_0`$-roots of $`𝔨`$ and by $`\mathrm{\Delta }(𝔭)`$ the set of $`𝔥_0`$-weights of $`𝔭`$. Fix a set of positive roots $`\mathrm{\Delta }_𝔨^+`$ and let $`𝔟_0`$ be the corresponding Borel subalgebra. Let $`\sigma `$ denote the involution which induces the decomposition $`𝔤=𝔨𝔭`$ and denote by $`\widehat{L}(𝔤,\sigma )`$ the extended loop algebra associated to the pair $`(𝔤,\sigma )`$ (see , Ch. 8) and by $`\widehat{W}`$ its Weyl group. The choice of $`\mathrm{\Delta }_𝔨^+`$ induces natural choices $`\widehat{\mathrm{\Delta }}^+,\widehat{\mathrm{\Delta }}_𝔨^+`$ for the positive roots of $`\widehat{L}(𝔤,\sigma ),\widehat{𝔨}`$ respectively. Here, as above, $`\widehat{𝔨}`$ denotes the untwisted affine algebra associated to $`𝔨`$; we also set $`\delta _𝔨`$ to be its fundamental imaginary root and $`\widehat{W}_𝔨`$ its Weyl group. Finally let $`\mathrm{\Delta }^+`$ denote a set of positive roots (w.r.t. the centralizer of $`𝔥_0`$ in $`𝔤`$) of $`𝔤`$ compatible with that of $`𝔨`$ (see 2.3). This allows us to define
$`\mathrm{\Delta }^+(𝔭)=\mathrm{\Delta }_{|𝔥_0}^+\mathrm{\Delta }(𝔭),`$
$`\widehat{\mathrm{\Delta }}^+(𝔭)=\{(m+{\displaystyle \frac{1}{2}})\delta _𝔨+\alpha \alpha \mathrm{\Delta }^+(𝔭),m\}.`$
Then it turns out that, up to an exponential factor, $`ch(L)`$ equals
$$\underset{\alpha \widehat{\mathrm{\Delta }}^+(𝔭)}{}\left(1+e^\alpha \right)^{mult(\alpha )}.$$
(see (3.4)). To extract from the previous formula information on the $`\widehat{𝔨}`$-module structure of $`L`$ we generalize an idea used in in the finite dimensional equal rank case. We introduce a natural map $`\psi _0:\widehat{𝔨}\widehat{L}(𝔤,\sigma )`$, whose transpose induces a bijection $`\psi _0^{}:\widehat{\mathrm{\Delta }}\widehat{\mathrm{\Delta }}_𝔨^+\widehat{\mathrm{\Delta }}^+(𝔭)`$. Using this map and the Weyl-Kac character formula we get the following decomposition into irreducible $`\widehat{𝔨}`$-modules
$$L(\stackrel{~}{\mathrm{\Lambda }}_ϵ)=\underset{\genfrac{}{}{0pt}{}{uW_{\sigma ,0}^{}}{\mathrm{}(u)ϵmod2}}{}L(\psi _0^{}(u\widehat{\rho })\widehat{\rho }_𝔨+\frac{1}{2}ϵ\delta _𝔨).$$
(1.1)
Here $`W_{\sigma ,0}^{}`$ is the set of minimal right coset representatives of $`(\psi _0^{})^1\widehat{W}_𝔨\psi _0^{}`$ in $`\widehat{W}`$, $`\widehat{\rho },\widehat{\rho }_𝔨`$ are the sum of fundamental weights in $`\widehat{\mathrm{\Delta }}^+,\widehat{\mathrm{\Delta }}_𝔨^+`$ and $`ϵ=0,1`$. A more accurate statement of formula (1.1) is given in Theorem 3.5. The combinatorial interpretation of formula (1.1) arises from the fact that, if $`C_1`$ denotes the fundamental alcove of $`\widehat{W}`$, then the set $`_{wW_{\sigma ,0}^{}}w\overline{C_1}`$ is the polytope studied in in connection with the problem of enumerating $`𝔟_0`$-stable abelian subalgebras of $`𝔭`$. This coincidence allows us to give a much more explicit rendering of formula (1.1). For instance, if $`\sigma `$ is an automorphism of type $`(0,\mathrm{},1,\mathrm{},0;1)`$ with $`1`$ in a position corresponding to a short root, then we have
$$L(\stackrel{~}{\mathrm{\Lambda }}_ϵ)=\underset{\genfrac{}{}{0pt}{}{A\mathrm{\Sigma }}{|A|ϵmod\mathrm{\hspace{0.17em}2}}}{}L\left(\mathrm{\Lambda }_{0,𝔨}+A\frac{1}{2}(|A|ϵ)\delta _𝔨\right),$$
where $`\mathrm{\Sigma }`$ is the set of $`𝔟_0`$-stable abelian subalgebras of $`𝔭`$, and for $`A\mathrm{\Sigma }`$, $`A,|A|`$ denote the sum (resp. the number) of the roots in $`A`$. The general case is completely described in Theorem 3.9.
The situation in the case of the spin representation ($`𝔨`$ semisimple) is much more complex. For instance, we need to use the factorization of the involution $`\sigma `$ as $`\mu \eta `$ with $`\eta `$ inner and $`\mu `$ diagram automorphism to write down the auxiliary map $`\psi _1:\widehat{𝔨}L^{}(𝔤,\sigma )`$ which plays the role of $`\psi _0`$ in the spin case and which is the crucial tool for manipulating the character. Moreover the target algebra $`L^{}(𝔤,\sigma )`$ is different according to whether we are considering the equal rank case, the $`A_{2n}^{(2)}`$ case or the remaining non equal rank cases. Surprisingly enough, we obtain a decomposition formula which is quite similar to (1.1):
$$ch(X_r)=2^{\frac{Nn}{2}}\underset{uW_{\sigma ,1}^{}}{}ch(L(a_0\psi _1^{}(u\widehat{\rho }^{})\widehat{\rho }_𝔨)).$$
(1.2)
We refer the reader to Proposition 4.7 for the undefined notation. As far as the combinatorial interpretation is concerned, we get again a description of formula (1.2) in terms of abelian subalgebras, but in the non equal rank the right class to consider is that of noncompact subalgebras (cf. Definition 4.3) which are stable under the Borel subalgebra $`𝔟_0𝔨𝔨_\mu `$ of the subalgebra $`𝔨𝔨_\mu `$ of $`\mu `$-fixed points in $`𝔨`$. The relevant results in this direction are Theorems 4.10, 4.12, 4.13.
A few words on the case in which $`𝔨`$ has a non-trivial center. In this case we describe the finite decomposition of an eigenspace of the center on the level $`1`$ modules. Also in this case there is a special subset of abelian stable subalgebras of $`𝔭`$ which plays an important role in the description of the decompositions. We describe in detail the finite decomposition of each eigenspace of the center (see Theorems 5.4, 5.5). The paper is organized as follows. In Section 2 we recall the necessary information on the structure and representation theory of affine Lie algebras, as well as the construction of all level $`1`$ modules over the affinization of orthogonal Lie algebras $`so_n`$. In Section 3 we find the decompositions of the basic $`+`$ vector representations of $`\widehat{so(𝔭)}`$, restricted to $`\widehat{𝔨}`$, $`𝔨`$ semisimple, and in Section 4 we solve the same problem for the spinor representations. In Section 5 we deal with the case when $`𝔨`$ has a non-trivial center. Finally in Section 6 we consider some concrete examples and discuss connections of the theory of abelian subalgebras to modular invariance.
## 2 Preliminaries
### 2.1 Lie algebra involutions and affine algebras
Let $`𝔤`$ be a semisimple Lie algebra over $``$, $`\sigma `$ an involutive automorphism of $`𝔤`$ and let $`𝔤=𝔨𝔭`$ be the corresponding eigenspace decomposition. Let $`𝔥_0`$ be a Cartan subalgebra of the reductive subalgebra $`𝔨`$ and let $`𝔷`$ be the centralizer of $`𝔥_0`$ in $`𝔤`$. Let $`(,)`$ denote a non-degenerate invariant symmetric bilinear form on $`𝔤`$. Then (see , Lemma 5.3 or Lemma 8.1), $`𝔷`$ is a Cartan subalgebra of $`𝔤`$ and $`(𝔨,𝔭)=0`$. In particular $`(,)_{|𝔨\times 𝔨}`$ is a nondegenerate invariant form on $`𝔨`$ and $`(,)`$ is nondegenerate when restricted to $`𝔥_0`$, so we can induce a form, still denoted by $`(,)`$ on $`𝔥_0^{}`$.
Consider the root system $`\mathrm{\Delta }_𝔨`$ of the pair $`(𝔨,𝔥_0)`$ and fix a subset of positive roots $`\mathrm{\Delta }_𝔨^+`$. Let $`𝔟_0`$ denote the corresponding Borel subalgebra of $`𝔨`$. We denote by $`\mathrm{\Delta }(𝔭)`$ the set of $`𝔥_0`$-weights of $`𝔭`$.
Let $`L(𝔤)`$ be the loop algebra of $`𝔤`$:
$$L(𝔤)=[t,t^1]𝔤.$$
Let $`L(𝔤,\sigma )`$ be the subalgebra
$$L(𝔤,\sigma )=\left(\underset{n,n\text{ even}}{}t^n𝔨\right)\left(\underset{n,n\text{ odd}}{}t^n𝔭\right)$$
and consider the extended loop algebra $`\widehat{L}(𝔤)=L(𝔤)K^{}d^{}`$ with bracket defined by
$`[t^n`$ $`X+\lambda K^{}+\mu d^{},t^mY+\lambda _1K^{}+\mu _1d^{}]=`$
$`=t^{n+m}[X,Y]+\delta _{n,m}n(X,Y)K^{}+\mu _1nt^nX+\mu mt^mY.`$ (2.1)
(The construction of the extended loop algebra is done in only for $`𝔤`$ simple, but everything extends to semisimple $`𝔤`$ in a straightforward way). Set $`\widehat{L}(𝔤,\sigma )=L(𝔤,\sigma )K^{}d^{}`$. Clearly $`\widehat{L}(𝔤,\sigma )`$ is a subalgebra of $`\widehat{L}(𝔤)`$.
Set $`\widehat{𝔥}=(1𝔥_0)K^{}d^{}`$ and let $`\widehat{\mathrm{\Delta }}`$ denote the set of nonzero $`\widehat{𝔥}`$-weights of $`\widehat{L}(𝔤,\sigma )`$. Define $`\delta ^{}\widehat{𝔥}^{}`$ by setting
$$\delta ^{}(1𝔥_0)=\delta ^{}(K^{})=0\delta ^{}(d^{})=1.$$
We identify $`𝔥_0^{}`$ and the subset $`\{\lambda \widehat{𝔥}^{}\lambda (d^{})=\lambda (K^{})=0\}`$ of $`\widehat{𝔥}^{}`$.
Notation. If $`\lambda \widehat{𝔥}^{}`$ we denote by $`\overline{\lambda }`$ its restriction to $`1𝔥_0`$.
The set of roots of $`\widehat{L}(𝔤,\sigma )`$ is
$`\widehat{\mathrm{\Delta }}=`$ $`\{k\delta ^{}+\alpha \alpha \mathrm{\Delta }_𝔨,k\text{ even}\}\{k\delta ^{}+\alpha \alpha \mathrm{\Delta }(𝔭),k\text{ odd}\}`$
$`\{k\delta ^{}k2,k0\}.`$
We set $`\widehat{\mathrm{\Delta }}^+=\mathrm{\Delta }_𝔨^+\{\alpha \widehat{\mathrm{\Delta }}\alpha (d^{})>0\}`$. Let $`\widehat{\mathrm{\Pi }}`$ be the corresponding set of simple roots. Denote by $`\widehat{W}`$ the Weyl group generated by $`\widehat{\mathrm{\Pi }}`$.
Assume now that $`\sigma `$ is indecomposable (i.e. $`𝔤`$ has no non-trivial $`\sigma `$-stable ideals). Then either $`𝔤`$ is simple, or $`𝔤`$ is a direct sum of two copies of a simple Lie algebra $`𝔨`$ and $`\sigma `$ permutes the summands. We call the latter the complex case.
The following two propositions give a summary of Exercises 8.1–8.4 from . The proof can be found in § 5 of , Ch.X.
###### Proposition 2.1.
1. $`|\widehat{\mathrm{\Pi }}|=n+1`$, where $`n`$ is the rank of $`𝔨`$.
2. If $`\widehat{\mathrm{\Pi }}=\{\alpha _0,\mathrm{},\alpha _n\}`$, then $`\overline{\alpha }_0,\mathrm{},\overline{\alpha }_n`$ span $`𝔥_0`$.
3. $`(\overline{\alpha }_i,\overline{\alpha }_i)>0`$ for all $`i`$ and
$$a_{ij}=2\frac{(\overline{\alpha }_i,\overline{\alpha }_j)}{(\overline{\alpha }_i,\overline{\alpha }_i)}_+,\text{ if }ij.$$
4. $`A=(a_{ij})`$ is a generalized Cartan matrix of an affine type.
We label the $`\alpha _i`$ so that the corresponding Dynkin diagram is one of those displayed at pp. 54–55 of .
If $`\alpha \widehat{\mathrm{\Delta }}`$, then $`\alpha `$ can be written uniquely as $`_{i=0}^nm_i(\alpha )\alpha _i`$ with $`m_i(\alpha )`$. Write $`\alpha _i=s_i\delta ^{}+\overline{\alpha }_i`$. By Proposition 2.1.3, $`\overline{\alpha }_i0`$. Let $`h_{\overline{\alpha }_i}`$ be the unique element of $`𝔥_0`$ such that $`\overline{\alpha }_i(h)=(h_{\overline{\alpha }_i},h)`$ for all $`h𝔥_0`$. Set $`h_i=\frac{2}{(\overline{\alpha }_i,\overline{\alpha }_i)}h_{\overline{\alpha }_i}`$ and fix $`t^{s_i}X_i\widehat{L}(𝔤,\sigma )_{\alpha _i}`$, $`t^{s_i}Y_i\widehat{L}(𝔤,\sigma )_{\alpha _i}`$ in such a way that $`(X_i,Y_i)=\frac{2}{(\overline{\alpha }_i,\overline{\alpha }_i)}`$. Then $`[X_i,Y_i]=h_i`$. It follows that
$$[t^{s_i}X_i,t^{s_i}Y_i]=\frac{2s_i}{(\overline{\alpha }_i,\overline{\alpha }_i)}K^{}+h_i.$$
Set $`\alpha _i^{}=\frac{2s_i}{(\overline{\alpha }_i,\overline{\alpha }_i)}K^{}+h_i`$ and $`\widehat{\mathrm{\Pi }}^{}=\{\alpha _0^{},\mathrm{},\alpha _n^{}\}`$. In the following proposition we use the notation of , Ch. 1.
###### Proposition 2.2.
The triple $`(\widehat{𝔥},\widehat{\mathrm{\Pi }},\widehat{\mathrm{\Pi }}^{})`$ is a realization of $`A`$ and the map
$$e_it^{s_i}X_if_it^{s_i}Y_i$$
extends to a Lie algebra isomorphism of the affine Kac-Moody algebra $`𝔤(A)`$ to the Lie algebra $`\widehat{L}(𝔤,\sigma ).`$
Let $`a_0,\mathrm{},a_n`$ (resp. $`a_0^{},\mathrm{},a_n^{}`$) be positive integers with $`G.C.D(a_0,\mathrm{},`$ $`a_n)=1`$ that are coefficients of a linear dependence between the rows (resp. columns) of the matrix $`A`$: $`\underset{i=0}{\overset{n}{}}a_i\overline{\alpha }_i=0`$ and $`\underset{i=0}{\overset{n}{}}a_i^{}h_i=0`$. Set $`\delta =_{i=0}^na_i\alpha _i`$ and notice that $`\delta =(_{i=0}^na_is_i)\delta ^{}`$. We also let $`K=_{i=0}^na_i^{}\alpha _i^{}`$ be the canonical central element.
###### Proposition 2.3.
Set $`k=\frac{2}{_{i=0}^na_is_i}`$. Then $`k=1`$ if $`\sigma `$ is inner and $`k=2`$ otherwise.
###### Proof.
Since $`2\delta ^{}=k\delta `$ is a root (cf. , Theorem 5.6 b)) we deduce that $`k`$ is an integer, hence $`k\{1,2\}`$. Since $`\delta ^{}=\frac{k}{2}\delta `$ is a root if and only if $`𝔷𝔥_0`$, we see that $`k=2`$ if and only if $`\sigma `$ is not of inner type. ∎
###### Remark 2.1.
If $`𝔤`$ is simple of type $`X_N`$, then $`\widehat{L}(𝔤,\sigma )`$ is an affine Kac-Moody algebra of type $`X_N^{(k)}`$. If $`𝔤=𝔨𝔨`$, where $`𝔨`$ is simple of type $`X_N`$, then $`\widehat{L}(𝔤,\sigma )`$ is of type $`X_N^{(1)}`$.
###### Remark 2.2.
If $`𝔤`$ is simple, using the terminology of , we have that $`\sigma `$ is an automorphism of type $`(s_0,\mathrm{},s_n;k)`$. In the complex case we have $`k=2,s_0=1`$.
If $`𝔤`$ is simple, we choose $`(,)=k^1(,)_n`$, where $`(,)_n`$ is the invariant form on $`𝔤`$ such that the square root lenght of a long root is $`2`$. We will call $`(,)_n`$ a normalized invariant form. If $`𝔤=𝔨\times 𝔨`$ we define $`(,)`$ by
$$((X,Y),(X^{},Y^{}))=\frac{1}{2}((X,X^{})_n+(Y,Y^{})_n),$$
where $`(,)_n`$ is the normalized invariant form on $`𝔨`$.
We define a standard invariant form $`(,)`$ on $`\widehat{L}(𝔤,\sigma )`$ by setting
$`(\alpha _i^{},h)={\displaystyle \frac{2}{(\overline{\alpha }_i,\overline{\alpha }_i)}}\alpha _i(h)\text{for }i=0,\mathrm{},n\text{ and }h\widehat{𝔥}`$ (2.2)
$`(d^{},d^{})=0.`$ (2.3)
We want to prove that the previous formulas define a normalized invariant form on $`\widehat{L}(𝔤,\sigma )`$, i.e. $`(\theta ,\theta )=2`$, where $`\theta =_{i=1}^na_i\alpha _i`$. Let $`\nu :\widehat{𝔥}\widehat{𝔥}^{}`$ be the induced isomorphism and let $`(,)`$ be the induced form on $`\widehat{𝔥}^{}`$. Since
$$\alpha _i(h)=\frac{(\overline{\alpha }_i,\overline{\alpha }_i)}{2}(\alpha _i^{},h)$$
we see that $`\nu ^1(\alpha _i)=\frac{(\overline{\alpha }_i,\overline{\alpha }_i)}{2}\alpha _i^{}`$. It follows that
$$(\alpha _i,\alpha _j)=(\overline{\alpha }_i,\overline{\alpha }_j),$$
hence, if $`\lambda ,\mu \text{Span }(\widehat{\mathrm{\Pi }})`$, then
$$(\lambda ,\mu )=(\overline{\lambda },\overline{\mu }).$$
(2.4)
We need to discuss the relationship between the roots of $`𝔤`$ and the roots of $`\widehat{L}(𝔤,\sigma )`$. Write $`𝔷=𝔥_0𝔞_{}`$, where $`𝔞_{}=𝔭𝔷`$. Let $`\mathrm{\Delta }`$ be the $`𝔷`$-root system of $`𝔤`$. There are three types of roots in $`\mathrm{\Delta }`$: those such that $`\alpha _{|𝔞_{}}=0`$ and whose root vector $`X_\alpha `$ is in $`𝔨`$, those such that $`\alpha _{|𝔞_{}}=0`$ and whose root vector $`X_\alpha `$ is in $`𝔭`$, and those such that $`\alpha _{|𝔞_{}}0`$. These are usually called respectively compact imaginary roots, noncompact imaginary (or singular imaginary) roots, and complex roots. To avoid confusion with standard Kac-Moody terminology we call them compact, noncompact, and complex. If $`\alpha `$ is a complex root, then the corresponding root vector decomposes as
$$X_\alpha =u_\alpha +v_\alpha $$
with $`u_\alpha 𝔨`$ and $`v_\alpha 𝔭`$. Then $`u_\alpha `$ is a root vector in $`𝔨`$ for the root $`\alpha _{|𝔥_0}`$ and $`v_\alpha `$ is a weight vector in $`𝔭`$ for the weight $`\alpha _{|𝔥_0}`$ in $`\mathrm{\Delta }(𝔭)`$. In particular $`\alpha `$ is a complex root if and only if $`\alpha _{|𝔥_0}\mathrm{\Delta }_𝔨\mathrm{\Delta }(𝔭)`$. It follows that $`\alpha \mathrm{\Delta }`$ is a compact root if and only if $`\alpha _{|𝔥_0}\widehat{\mathrm{\Delta }}`$ and $`\delta ^{}+\alpha _{|𝔥_0}\widehat{\mathrm{\Delta }}`$, $`\alpha \mathrm{\Delta }`$ is a noncompact root if and only if $`\alpha _{|𝔥_0}\widehat{\mathrm{\Delta }}`$ and $`\delta ^{}+\alpha _{|𝔥_0}\widehat{\mathrm{\Delta }}`$, and $`\alpha \mathrm{\Delta }`$ is a complex root if and only if $`\alpha _{|𝔥_0}\widehat{\mathrm{\Delta }}`$ and $`\delta ^{}+\alpha _{|𝔥_0}\widehat{\mathrm{\Delta }}`$. More precisely if $`k=1`$, then $`𝔥_0=𝔷`$ hence, if $`\alpha \widehat{\mathrm{\Delta }}`$, then $`\overline{\alpha }=\beta _{|𝔥_0}`$ with $`\beta `$ compact or noncompact. It follows that $`(\alpha ,\alpha )=(\overline{\alpha },\overline{\alpha })=(\beta ,\beta )_n`$. In particular, since $`\alpha _0`$ is a long root, $`(\alpha _0,\alpha _0)=2`$. If $`k=2`$ and $`𝔤`$ is simple, then $`\delta ^{}=\delta `$ hence, if $`\alpha \widehat{\mathrm{\Delta }}`$, then $`\overline{\alpha }=\beta _{|𝔥_0}`$ with $`\beta `$ compact or noncompact if and only if $`\alpha `$ is a long root. It follows that if $`\alpha `$ is a long root and $`\overline{\alpha }=\beta _{|𝔥_0}`$, then $`(\alpha ,\alpha )=(\overline{\alpha },\overline{\alpha })=k(\beta ,\beta )_n=4`$. In particular, since $`\alpha _0`$ is a short root, $`a_0(\alpha _0,\alpha _0)=2`$. In the complex case one checks directly that $`(\alpha _0,\alpha _0)=2`$ in this case too. We have proven
###### Lemma 2.4.
The form $`(,)`$ on $`\widehat{L}(𝔤,\sigma )`$ defined above is a normalized standard invariant form.
Since $`(d^{},K)=a_i^{}\frac{2}{(\alpha _i,\alpha _i)}\alpha _i(d^{})=a_i^{}\frac{2s_i}{(\alpha _i,\alpha _i)}`$, and $`K=_{i=0}^n\frac{2a_i^{}s_i}{(\alpha _i,\alpha _i)}K^{}`$ we see that $`(d^{},K^{})=1`$. Also remark that
$$(d^{},h)=0\text{ if }h𝔥_0$$
(2.5)
which is easily proved by observing that
$$(d^{},h_i)=(d^{},\alpha _i^{}2\frac{s_i}{(\alpha _i,\alpha _i)}K^{})=\frac{2}{(\alpha _i,\alpha _i)}\alpha _i(d^{})2\frac{s_i}{(\alpha _i,\alpha _i)}=0.$$
Let $`H`$ be the unique element of $`𝔥_0`$ such that $`\overline{\alpha _i}(H)=s_i`$ for $`i=1,\mathrm{},n`$. Then easy calculations show that $`d=\frac{a_0k}{2}(d^{}H\frac{1}{2}(H,H)K^{})`$ is a scaling element for $`\widehat{L}(𝔤,\sigma )`$ and $`(d,d)=0`$. It follows that, if we define $`\mathrm{\Lambda }_0\widehat{𝔥}^{}`$ by setting $`\mathrm{\Lambda }_0(\alpha _i^{})=\delta _{i0}`$ and $`\mathrm{\Lambda }_0(d)=0`$, then the form $`(,)`$ on $`\widehat{L}(𝔤,\sigma )`$ is given by the formulas of \[12, § 6.2\]. In particular we see that $`2\frac{s_ia_i^{}}{(\alpha _i,\alpha _i)}=a_is_i=\frac{2}{k}`$. Hence $`K=\frac{2}{k}K^{}`$.
### 2.2 The Lie algebra $`\widehat{𝔨}`$ and the character formula
In general $`𝔨`$ is a reductive Lie algebra, hence we can write $`𝔨=𝔨_0_{S=1}^M𝔨_S`$, where $`𝔨_0`$ is the center of $`𝔨`$ and $`𝔨_S`$ are the simple ideals of $`𝔨`$. If $`S>0`$, we denote by $`\mathrm{\Pi }_S`$ the set of simple roots of $`𝔨_S`$. Let also $`W_S`$ be the relative Weyl group: $`W_S=s_\alpha \alpha \mathrm{\Pi }_S`$, $`\mathrm{\Delta }_S`$ the relative root system, $`\mathrm{\Delta }_S=W_S\mathrm{\Pi }_S`$, and $`\theta _S`$ the highest root of $`\mathrm{\Delta }_S`$. We recall that the dimension of $`𝔨_0`$ is at most one.
We define the affine Lie algebra $`\widehat{𝔨}`$ as follows. Consider the standard loop algebra $`\stackrel{~}{𝔨}=L(𝔨)=_SL(𝔨_S)`$. On each simple ideal $`𝔨_S`$ let $`(,)_S`$ be the normalized invariant form. Set also $`(,)_0`$ to be the normalized invariant form of $`𝔤`$ restricted to $`𝔨_0`$. We then let $`𝔨_S^{}=L(𝔨_S)K_S`$ be the central extension of $`L(𝔨_S)`$ with bracket defined as usual as
$$[t^mX,t^nY]=t^{n+m}[X,Y]+\delta _{m,n}m(X,Y)_SK_S.$$
Set finally $`\widehat{𝔨}=(_S𝔨_S^{})d_𝔨`$, where $`d_𝔨`$ is the derivation $`t\frac{d}{dt}`$ on $`L(𝔨)`$ extended by setting $`[d_𝔨,K_S]=0`$. Set $`\widehat{𝔨}_S=𝔨_S^{}d_𝔨`$. We can extend the form $`(,)_S`$ on all of $`\widehat{𝔨}_S`$ by setting $`(K_S,𝔨_S)_S=(K_S,K_S)_S=(d_𝔨,𝔨_S)_S=(d_𝔨,d_𝔨)_S=0`$ and $`(d_𝔨,K_S)_S=1`$.
We denote by $`\widehat{W}_𝔨`$ the Weyl group of $`\widehat{𝔨}`$. It is a group of linear transformations on $`\widehat{𝔥}_𝔨^{}`$, where
$$\widehat{𝔥}_𝔨=1𝔥_0(_SK_S)d_𝔨.$$
If we define $`\delta _𝔨\widehat{𝔥}_𝔨^{}`$ setting
$$\delta _𝔨(d_𝔨)=1,\delta _𝔨(1𝔥_0)=\delta _𝔨(K_S)=0,$$
then the set of roots for $`\widehat{𝔨}`$ is
$$\widehat{\mathrm{\Delta }}_𝔨=\{n\delta _𝔨+\alpha \alpha \mathrm{\Delta }_𝔨,n\}\{n\delta _𝔨n,n0\},$$
where, as usual, we regard $`𝔥_0^{}`$ as a subset of $`\widehat{𝔥}_𝔨^{}`$ by extending $`\lambda 𝔥_0^{}`$ setting $`\lambda (d_𝔨)=\lambda (K_S)=0`$.
We set $`\mathrm{\Pi }_𝔨=_S\mathrm{\Pi }_S`$ and
$$\widehat{\mathrm{\Pi }}_𝔨=\mathrm{\Pi }_𝔨\{\delta _𝔨\theta _SS>0\}.$$
$`\widehat{\mathrm{\Pi }}`$ is a set of simple roots for $`\widehat{𝔨}`$ and we denote by $`\widehat{\mathrm{\Delta }}_𝔨^+`$ the corresponding subset of positive roots.
If $`\lambda \widehat{𝔥}_𝔨^{}`$ is a $`\widehat{\mathrm{\Delta }}_𝔨^+`$-dominant integral weight, we denote by $`L(\lambda )`$ the irreducible integrable $`\widehat{𝔨}`$ module of highest weight $`\lambda `$. We denote by $`\mathrm{\Lambda }_j^S`$ the fundamental weights of $`\widehat{𝔨}_S`$ and we set $`\widehat{\rho }_𝔨=\underset{j,S>0}{}\mathrm{\Lambda }_j^S`$. Recall the Weyl-Kac character formula for the character of $`L(\lambda ),\lambda \widehat{𝔥}_𝔨^{}`$:
$$ch(L(\lambda ))=\frac{_{w\widehat{W}_𝔨}ϵ(w)e^{w(\lambda +\widehat{\rho }_𝔨)\widehat{\rho }_𝔨}}{_{\alpha \widehat{\mathrm{\Delta }}_𝔨^+}(1e^\alpha )^{m_\alpha }}.$$
(2.6)
Here $`m_\alpha `$ denotes the multiplicity of the root $`\alpha `$.
### 2.3 Realization of level $`1`$ modules of $`\widehat{so(𝔭)}`$
If $`X𝔨`$ set $`ad_𝔭(X)=ad(X)_{|𝔭}`$. Since the action of $`𝔨`$ on $`𝔭`$ is orthogonal with respect to $`(,)`$, we have an inclusion $`𝔨so(𝔭)`$ defined by $`Xad_𝔭(X)`$. We let $`\widehat{so(𝔭)}`$ denote the affine Lie algebra $`L(so(𝔭))K_𝔭d_𝔭`$, where $`L(so(𝔭))K_𝔭`$ is the central extension of $`L(so(𝔭))`$ defined by setting
$$[t^mA,t^nB]=t^{n+m}[A,B]+\delta _{m,n}m<A,B>K_𝔭$$
and $`<A,B>=\frac{1}{2}tr(AB)`$. Note that $`\widehat{so(𝔭)}`$ is an affine algebra of type $`B^{(1)}`$ or $`D^{(1)}`$ according to whether $`dim(𝔭)`$ is odd or even.In the following we recall the realization of the level $`1`$ irreducible modules of $`\widehat{so(𝔭)}`$ described in .
Fix $`r`$, set $`r^{}=\frac{r}{2}`$ and consider the loop space $`\stackrel{~}{𝔭}=[t,t^1]𝔭`$. Define the bilinear form $`\mathrm{\Phi }_r`$ on $`\stackrel{~}{𝔭}`$ by setting
$$\mathrm{\Phi }_r(t^{m_1}X,t^{m_2}Y)=\delta _{r+m_1+m_2,1}(X,Y).$$
Let $`Cl_r(\stackrel{~}{𝔭})=Cl_r^+(\stackrel{~}{𝔭})Cl_r^{}(\stackrel{~}{𝔭})`$ be the corresponding Clifford algebra, decomposed into the sum of the even and odd part.
If $`m`$ set $`\stackrel{~}{𝔭}_m=_{im}(t^i𝔭)`$ and $`\stackrel{~}{𝔭}_m^{}=_{i<m}(t^i𝔭)`$. If $`r`$ is even set $`\stackrel{~}{U}_r=\stackrel{~}{𝔭}_r^{}`$. Then $`\stackrel{~}{U}_r`$ is a maximal isotropic subspace for $`\stackrel{~}{𝔭}`$ with respect to $`\mathrm{\Phi }_r`$. If $`r`$ is odd we choose a maximal isotropic subspace of $`\stackrel{~}{𝔭}`$ as follows. Recall (see \[26, § 9.3.1\]) that a set of positive roots $`\mathrm{\Delta }^+`$ for $`𝔤`$ is compatible with $`\mathrm{\Delta }_𝔨^+`$ if it is $`\sigma `$-stable and $`\mathrm{\Delta }^+\mathrm{\Delta }_𝔨\mathrm{\Delta }_𝔨^+`$. Let $`\mathrm{\Delta }^+`$ be such a positive system. Set $`\mathrm{\Delta }^+(𝔭)=\mathrm{\Delta }(𝔭)\mathrm{\Delta }^+|_{𝔥_0}`$ and
$$𝔭^\pm =\underset{\alpha \pm \mathrm{\Delta }^+(𝔭)}{}𝔭_\alpha .$$
Thus we can write
$$𝔭=𝔞_{}𝔭^+𝔭^{},$$
where $`𝔞_{}=𝔷𝔭`$. Choose a maximal isotropic subspace $`𝔞`$ of $`𝔞_{}`$. Set $`U=𝔞𝔭^+`$ and
$$\stackrel{~}{U}_r=\stackrel{~}{𝔭}_r^{}(t^{r^{}1}U).$$
Let $`\widehat{\sigma }_r`$ denote the left action of $`\widehat{so(𝔭)}`$ on the spin module (defined in ) $`s_r(\stackrel{~}{𝔭},\stackrel{~}{U}_r)=Cl_r(\stackrel{~}{𝔭})/Cl_r(\stackrel{~}{𝔭})\stackrel{~}{U}_r`$ .
If $`r`$ is even we let $`X_r`$ be the subalgebra of $`Cl_r(\stackrel{~}{𝔭})`$ generated by $`\stackrel{~}{𝔭}_r^{}^{}`$. If $`r`$ is odd, set $`L=dim𝔞_{}`$ and $`l=\frac{L}{2}=dim𝔞`$. Fix a basis $`\{v_i\}`$ of $`𝔞_{}`$ such that $`\{v_iil\}`$ is a basis of $`𝔞`$ and $`(v_i,v_{Lj+1})=\delta _{ij}`$. Then, if $`L`$ is even, we let $`X_r`$ be the subalgebra of $`Cl_r(\stackrel{~}{𝔭})`$ generated by
$$\stackrel{~}{𝔭}_{r^{}1}^{}(t^{r^{}1}(\text{Span }(v_ii>l)𝔭^{}))$$
while, if $`L`$ is odd, we let $`X_r=X_r^+`$ be the subalgebra of $`Cl_r(\stackrel{~}{𝔭})`$ generated by
$$\left(\stackrel{~}{𝔭}_{r^{}1}^{}(t^{r^{}1}(\text{Span }(v_ii>l+1)𝔭^{}))\right)(t^{r^{}1}v_{l+1}).$$
If $`rL`$ is even (resp. odd) then $`Cl_r(\stackrel{~}{𝔭})=Cl_r(\stackrel{~}{𝔭})\stackrel{~}{U}_rX_r`$ (resp. $`Cl_r^+(\stackrel{~}{𝔭})=Cl_r^{}(\stackrel{~}{𝔭})\stackrel{~}{U}_rX_r`$) therefore we can identify $`s_r(\stackrel{~}{𝔭},\stackrel{~}{U}_r)`$ and $`X_r`$. Set moreover $`X_r^\pm =X_rCl_r^\pm (\stackrel{~}{𝔭})`$. Set $`m=\frac{dim(𝔭)}{2}`$ and let $`\stackrel{~}{\mathrm{\Lambda }}_0,\mathrm{},\stackrel{~}{\mathrm{\Lambda }}_m`$ denote the fundamental weights of $`\widehat{so(𝔭)}`$ normalized by setting $`\stackrel{~}{\mathrm{\Lambda }}_i(d_𝔭)=0`$. Define an element $`\delta _𝔭`$ in the dual of the Cartan subalgebra of $`\widehat{so(𝔭)}`$ requiring that $`\delta _𝔭(d_𝔭)=1,\delta _𝔭(K_𝔭)=0,\delta _𝔭(x)=0`$ for any $`x`$ in the Cartan subalgebra of $`so(𝔭)`$.
###### Proposition 2.5.
1. The action $`\widehat{\sigma }_r`$ of $`\widehat{so(𝔭)}`$ on $`X_r`$ is described explicitly in Theorem 1 of .
2. If $`rL`$ is even we have $`X_r^+L(\stackrel{~}{\mathrm{\Lambda }}_0),X_r^{}L(\stackrel{~}{\mathrm{\Lambda }}_1\frac{1}{2}\delta _𝔭)`$ if $`r`$ is even and $`X_r^+L(\stackrel{~}{\mathrm{\Lambda }}_m),X_r^{}L(\stackrel{~}{\mathrm{\Lambda }}_{m1})`$ if $`r`$ is odd.
3. If both $`L`$ and $`r`$ are odd we have $`X_r^+L(\stackrel{~}{\mathrm{\Lambda }}_m)`$.
We shall conventionally refer to $`X_r`$ as the basic and vector representation if $`r`$ is even and as the spin representation if $`r`$ is odd.
Let $`\eta :\widehat{𝔨}\widehat{so(𝔭)}`$ be the Lie algebra homomorphism such that $`t^mXt^mad_𝔭(X)`$ and $`d_𝔨d_𝔭`$. Requiring that $`\eta `$ is a Lie algebra homomorphism fixes the value of $`\eta (K_S)`$: if $`h`$ is a nonzero element of $`𝔥_0𝔨_S`$, then
$$\eta ([th,t^1h])=\eta ((h,h)_SK_S)$$
while
$$[\eta (th),\eta (t^1h)]=[tad_𝔭(h),t^1ad_𝔭(h)]=<ad_𝔭(h),ad_𝔭(h)>K_𝔭$$
so
$$\eta (K_S)=\frac{<ad_𝔭(h),ad_𝔭(h)>}{(h,h)_S}K_𝔭.$$
Set
$$j_S=\frac{<ad_𝔭(h),ad_𝔭(h)>}{(h,h)_S}.$$
Let $`\kappa (,)`$ be the Killing form of $`𝔤`$ and $`\kappa _𝔨(,)`$ the Killing form of $`𝔨`$. We have
$$tr(ad_𝔭(h)ad_𝔭(h))=\kappa (h,h)\kappa _𝔨(h,h).$$
By Corollary 8.7 of we have that $`\frac{\kappa _𝔨(h,h)}{2(h,h)_S}=h_S^{}`$ where $`h_S^{}`$ is the dual Coxeter number of $`𝔨_S`$ if $`S>0`$ while $`h_0^{}=0`$. If $`𝔤`$ is simple we can apply Corollary 8.7 of obtaining $`\frac{\kappa (h,h)}{(h,h)}=2kh^{}`$ (here $`h^{}`$ denotes the dual Coxeter number of $`𝔤`$). It follows that
$$\frac{\kappa (h,h)}{(h,h)_S}=2kh^{}\frac{(h,h)}{(h,h)_S}.$$
(2.7)
In the complex case one checks directly that (2.7) still holds. The final outcome is that $`j_S`$ is independent of the choice of $`h`$ and that we can write
$$j_S=kh^{}\frac{(h,h)}{(h,h)_S}h_S^{}.$$
(2.8)
Notice that, if $`S>0`$, $`\frac{(h,h)}{(h,h)_S}=\frac{2}{(\theta _S,\theta _S)}`$. Setting $`n_S=\frac{2k}{(\theta _S,\theta _S)}=\frac{a_0k(\alpha _0,\alpha _0)}{(\theta _S,\theta _S)}`$ for $`S>0`$ and $`n_0=1`$, we see that, since $`a_0k(\alpha _0,\alpha _0)`$ is the length of a long root of $`\widehat{L}(𝔤,\sigma )`$, we can rewrite (2.8) as
$$j_S=n_Sh^{}h_S^{}$$
where $`n_S`$ is an integer (=1,2,3 or 4).
The map $`\eta `$ defines a representation $`\sigma _r`$ of $`\widehat{𝔨}`$ on $`X_r`$ by setting $`\sigma _r=\widehat{\sigma }_r\eta `$. Using Theorem 1 of , we can describe explicitly the action $`\sigma _r`$ of the Cartan subalgebra $`\widehat{𝔥}_𝔨`$ as follows. Fix weight vectors $`X_\alpha 𝔭_\alpha `$ such that $`(X_\alpha ,X_\alpha )=1`$ and set
$$\xi _{i,\alpha }=\{\begin{array}{cc}t^iX_\alpha \hfill & \text{if }rL\text{ is even}\hfill \\ \sqrt{2}(t^iX_\alpha )(t^{r^{}1}v_{l+1})\hfill & \text{if }rL\text{ is odd}\hfill \end{array}.$$
Set also
$$v_{i,j}=\{\begin{array}{cc}t^iv_j\hfill & \text{if }rL\text{ is even}\hfill \\ \sqrt{2}(t^iv_j)(t^{r^{}1}v_{l+1})\hfill & \text{if }rL\text{ is odd}\hfill \end{array}.$$
Set
$$J_{}=\{(i,\alpha )ir^{}1,\alpha \mathrm{\Delta }(𝔭)\}\{(i,j)ir^{}1,j=1,\mathrm{},L\}$$
if $`r`$ is even, and
$`J_{}=`$ $`\{(i,\alpha )i<r^{}1,\alpha \mathrm{\Delta }(𝔭)\}\{(r^{}1,\alpha )\alpha \mathrm{\Delta }^+(𝔭)\}`$
$`\{(i,j)i<r^{}1,j=1,\mathrm{},L\}\{(r^{}1,j)Lj+1l\}`$
if $`r`$ is odd. Putting any total order on $`J_{}`$, the vectors
$$v_{i_1,j_1}\mathrm{}v_{i_h,j_h}\xi _{m_1,\beta _1}\mathrm{}\xi _{m_k,\beta _k}$$
(2.9)
with $`(i_1,j_1)<\mathrm{}<(i_h,j_h)`$ and $`(m_1,\beta _1)<\mathrm{}<(m_k,\beta _k)`$ in $`J_{}`$ form a basis for $`X_r`$. If $`v`$ is a vector given by (2.9) and $`h𝔥_0`$ then
$$\sigma _r(h)v=(\underset{s=1}{\overset{k}{}}\beta _s(h)+\frac{1}{2}\delta _{r,2r^{}+1}\underset{\alpha \mathrm{\Delta }^+(𝔭)}{}\alpha (h))v$$
while
$$\sigma _r(d_𝔨)v=(\underset{s=1}{\overset{h}{}}(i_s+\frac{r+1}{2})+\underset{s=1}{\overset{k}{}}(m_s+\frac{r+1}{2}))v.$$
Since $`K_𝔭`$ acts as the identity on $`X_r`$ we find that
$$\sigma _r(K_S)v=j_Sv.$$
It follows that $`v`$ is a weight vector having weight
$$\underset{S}{}j_S\mathrm{\Lambda }_0^S+\underset{s=1}{\overset{h}{}}(i_s+\frac{r+1}{2})\delta _𝔨+\underset{s=1}{\overset{k}{}}((m_s+\frac{r+1}{2})\delta _𝔨+\beta _s)+\delta _{r,2r^{}+1}\rho _n,$$
(2.10)
where by definition $`\rho _n=\frac{1}{2}_{\alpha \mathrm{\Delta }^+(𝔭)}\alpha `$.
We shall use this formula in the next sections, treating separately the cases $`r`$ even, $`r`$ odd to obtain the decomposition of the $`\widehat{so(𝔭)}`$-modules $`X_r`$ with respect to the subalgebra $`\eta (\widehat{𝔨})`$.
## 3 Decomposition of the basic and vector representation (semisimple case)
Here we assume that $`r`$ is even and $`𝔨`$ is semisimple. Set $`c_S=\frac{(h,h)}{(h,h)_S}`$, where $`h`$ is any nonzero element of $`𝔥_0𝔨_S`$. As already observed $`c_S`$ does not depend on $`h`$. Define a linear map $`\psi _0:\widehat{𝔨}\widehat{L}(𝔤,\sigma )`$ by setting
$$\psi _0(t^nX)=t^{2n}X\psi _0(d_𝔨)=d^{}/2\psi _0(K_S)=2c_SK^{}=kc_SK.$$
(3.1)
Let $`\psi _0^{}:\widehat{𝔥}^{}(\widehat{𝔥}_𝔨)^{}`$ denote the transpose of $`\psi _0`$ (restricted to $`\widehat{𝔥}^{}`$). Clearly $`\delta _𝔨=2\psi _0^{}(\delta ^{})`$. Set
$$\widehat{\mathrm{\Delta }}^+(𝔭)=\{(m+1/2)\delta _𝔨+\alpha \alpha \mathrm{\Delta }(𝔭),m0\}.$$
Let us record the following facts.
###### Lemma 3.1.
1. The map $`\psi _0^{}`$ defines a bijection between $`\widehat{\mathrm{\Delta }}^+`$ and $`\widehat{\mathrm{\Delta }}_𝔨^+\widehat{\mathrm{\Delta }}^+(𝔭)`$.
2. $`\psi _0^{}(\widehat{𝔥})`$ is stable under the action of $`\widehat{W}_𝔨`$.
###### Proof.
For the first statement, remark that $`\psi _0^{}(m\delta ^{}+\alpha )=\frac{m}{2}\delta _𝔨+\alpha `$ for all $`\alpha \mathrm{\Delta }_𝔨\mathrm{\Delta }(𝔭)\{0\}`$. To prove the second statement, note that if $`\alpha =m\delta _𝔨+\beta `$ with $`\beta `$ a nonzero root in $`\mathrm{\Delta }_𝔨`$ and $`\lambda =\psi _0^{}(\mu )`$, then $`\alpha =\psi _0^{}(2m\delta ^{}+\beta )`$, hence
$$s_\alpha \lambda =\lambda \lambda (\alpha ^{})\alpha =\psi _0^{}(\mu \lambda (\alpha ^{})(2m\delta ^{}+\beta )).$$
Since $`\psi _0`$ is onto, $`\psi _0^{}`$ is bijective on its image so, for each $`w\widehat{W}_𝔨`$, we can define
$$\widehat{w}=(\psi _0^{})^1w\psi _0^{}.$$
Then we set $`\widehat{W}_{\sigma ,0}=\{\widehat{w}w\widehat{W}_𝔨\}`$.
###### Lemma 3.2.
Let $`\alpha \widehat{\mathrm{\Delta }}_𝔨`$ be a real root. If $`\alpha =\psi _0^{}(\gamma )`$ then $`\psi _0(\alpha ^{})=\gamma ^{}`$. In particular $`\widehat{s}_\alpha =s_\gamma `$.
###### Proof.
Assume that $`\alpha `$ is a real root of $`\widehat{𝔨}_S`$. We define a bilinear form $`\{,\}`$ on $`\widehat{𝔥}_𝔨\widehat{𝔨}_S`$ by setting
$$\{h,k\}=(\psi _0(h),\psi _0(k)).$$
Obviously, if $`h,k𝔥_0`$ then $`\{h,k\}=(h,k)=c_S(h,k)_S`$. We claim that
$$\{,\}=c_S(,)_S.$$
Indeed, if $`h𝔥_0𝔨_S`$, then $`\{K_S,h\}=kc_S(K,h)=0`$ and $`\{K_S,K_S\}=k^2c_S^2(K,K)=0`$. By (2.5) $`\{d_𝔨,h\}=\frac{1}{2}(d^{},h)=0`$, $`\{d_𝔨,d_𝔨\}=\frac{1}{4}(d^{},d^{})=0`$. Finally
$$\{d_𝔨,K_S\}=c_S(d^{},K^{})=c_S.$$
Let $`\nu _S:\widehat{𝔥}_𝔨\widehat{𝔨}_S(\widehat{𝔥}_𝔨\widehat{𝔨}_S)^{}`$ be the isomorphism induced by $`\{,\}`$. We claim that, if $`\alpha =\psi _0^{}(\gamma )\widehat{\mathrm{\Delta }}_S`$, then $`\psi _0(\nu _S^1(\alpha ))=\nu ^1(\gamma )`$. Indeed write $`\nu _S^1(\alpha )=h_S+aK_S`$ with $`h_S𝔨_S𝔥_0`$. If $`h\widehat{𝔥}`$, there exists $`h^{}\widehat{𝔥}_𝔨`$ such that $`h=\psi _0(h^{})`$. Write $`h^{}=_i(h_i^{}+b_iK_i)+cd_𝔨`$ with $`h_i^{}𝔨_i𝔥_0`$. Then
$`(\psi _0(\nu _S^1(\alpha )),h)`$ $`=(h_S+kc_SaK,{\displaystyle \underset{i}{}}(h_i^{}+kc_ib_iK)+\frac{1}{2}cd^{})`$
$`=(h_S+kc_SaK,h_S^{}+kc_Sb_SK+\frac{c}{2}d^{})`$
$`=(\psi _0(\nu _S^1(\alpha )),\psi _0(h_S^{}+b_SK_S+cd_𝔨))`$
$`=\{\nu _S^1(\alpha ),h_S^{}+b_SK_S+cd_𝔨\}`$
$`=\alpha (h^{})=\gamma (\psi _0(h^{}))=\gamma (h).`$
In particular
$$\psi _0(\alpha ^{})=\psi _0(\frac{2\nu _S^1(\alpha )}{\{\nu _S^1(\alpha ),\nu _S^1(\alpha )\}})=\frac{2\nu ^1(\gamma )}{(\nu ^1(\gamma ),\nu ^1(\gamma ))}=\gamma ^{}.$$
Let $`\mathrm{\Pi }_0=\{\alpha _i\widehat{\mathrm{\Pi }}s_i=0\}`$ and let $`\widehat{\mathrm{\Pi }}_{\sigma ,0}=\mathrm{\Pi }_0\{k\delta \theta _SS>0\}`$. Set also
$$\widehat{\mathrm{\Delta }}_{\sigma ,0}=\mathrm{\Delta }_𝔨+2\delta ^{}$$
(3.2)
An obvious consequence of Lemma 3.2 is the following fact.
###### Corollary 3.3.
$`\widehat{\mathrm{\Delta }}_{\sigma ,0}`$ is a root system and $`\widehat{W}_{\sigma ,0}`$ is the corresponding reflection group. In particular $`\widehat{W}_{\sigma ,0}`$ is the subgroup of $`\widehat{W}`$ generated by the reflections $`\{s_\alpha \alpha \widehat{\mathrm{\Pi }}_{\sigma ,0}\}`$.
###### Proof.
For the first statement observe that $`\psi _0^{}(\widehat{\mathrm{\Delta }}_{\sigma ,0})`$ is the set of real roots in $`\widehat{\mathrm{\Delta }}_𝔨`$. As for the second assertion, since $`k\delta =2\delta ^{}`$, it is clear that $`\psi _0^{}`$ is a bijection between $`\widehat{\mathrm{\Pi }}_{\sigma ,0}`$ and $`\widehat{\mathrm{\Pi }}_𝔨`$. ∎
Set $`N=rank(𝔤)`$ and for $`\alpha \widehat{\mathrm{\Delta }}^+(𝔭)`$ set $`m_\alpha =1`$ if $`\alpha =(m+\frac{1}{2})\delta _𝔨+\beta `$ with $`\beta \mathrm{\Delta }(𝔭)\{0\}`$, while we set $`m_\alpha =Nn`$ if $`\alpha =(m+\frac{1}{2})\delta _𝔨`$. Observe that, if $`\alpha =\psi _0^{}(\beta )`$, then $`m_\alpha `$ equals the multiplicity $`m_\beta `$ of $`\beta `$ as a root of $`\widehat{L}(𝔤,\sigma )`$ (see , Corollary 8.3).
Let $`\widehat{\rho }`$ be the element of $`\widehat{𝔥}^{}`$ such that $`\widehat{\rho }(d^{})=0`$ and $`\widehat{\rho }(\alpha _i^{})=1`$ for $`i=0,\mathrm{},n`$. Set
$$D_𝔤^{}=e^{_Sj_S\mathrm{\Lambda }_0^S+\widehat{\rho }_𝔨}\underset{\alpha \widehat{\mathrm{\Delta }}^+(𝔭)}{}(1e^\alpha )^{m_\alpha }\underset{\alpha \widehat{\mathrm{\Delta }}_𝔨^+}{}(1e^\alpha )^{m_\alpha }.$$
###### Lemma 3.4.
We have
$$D_𝔤^{}=e^{\psi _0^{}(\widehat{\rho })}\underset{\alpha \widehat{\mathrm{\Delta }}^+}{}(1e^{\psi _0^{}(\alpha )})^{m_\alpha }.$$
###### Proof.
By Lemma 3.1
$$D_𝔤^{}=e^{_Sj_S\mathrm{\Lambda }_0^S+\widehat{\rho }_𝔨}\underset{\alpha \widehat{\mathrm{\Delta }}^+}{}(1e^{\psi _0^{}(\alpha )})^{m_\alpha }.$$
It remains only to check that
$$\underset{S}{}j_S\mathrm{\Lambda }_0^S+\widehat{\rho }_𝔨=\psi _0^{}(\widehat{\rho }).$$
(3.3)
Since $`\psi _0(\alpha ^{})=\alpha ^{}`$ for $`\alpha \mathrm{\Pi }_𝔨`$ we see that $`\psi _0^{}(\widehat{\rho })(\alpha ^{})=1=(_Sj_S\mathrm{\Lambda }_0^S+\widehat{\rho }_𝔨)(\alpha ^{})`$ for $`\alpha \mathrm{\Pi }_𝔨`$. We defined $`\widehat{\rho }`$ so that $`\widehat{\rho }(d^{})=0`$ hence $`\psi _0^{}(\widehat{\rho })(d_𝔨)=0=(_Sj_S\mathrm{\Lambda }_0^S+\widehat{\rho }_𝔨)(d_𝔨)`$. It remains only to check that $`\psi _0^{}(\widehat{\rho })(K_S)=(_Sj_S\mathrm{\Lambda }_0^S+\widehat{\rho }_𝔨)(K_S)=j_S+h_S^{}`$, but, since $`\psi _0^{}(\widehat{\rho })(K_S)=kc_S\widehat{\rho }(K)`$, this follows immediately from (2.8) and the fact that $`\widehat{\rho }(K)=h^{}`$. ∎
By formula (2.10), the character of $`X_r`$ can be written as
$$ch(X_r)=e^{_Sj_S\mathrm{\Lambda }_0^S}\underset{\alpha \widehat{\mathrm{\Delta }}^+(𝔭)}{}(1+e^\alpha )^{m_\alpha }$$
(3.4)
hence
$$ch(X_r^+)ch(X_r^{})=e^{_Sj_S\mathrm{\Lambda }_0^S}\underset{\alpha \widehat{\mathrm{\Delta }}^+(𝔭)}{}(1e^\alpha )^{m_\alpha }$$
Applying Lemma 3.4 and setting
$$D_𝔨=e^{\widehat{\rho }_𝔨}\underset{\alpha \widehat{\mathrm{\Delta }}_𝔨^+}{}(1e^\alpha )^{m_\alpha },$$
we can write
$$ch(X_r^+)ch(X_r^{})=\frac{D_𝔤^{}}{D_𝔨}=\frac{_{\alpha \widehat{\mathrm{\Delta }}^+}(1e^{\psi _0^{}(\alpha )})^{m_\alpha }}{D_𝔨}$$
(3.5)
By the “denominator identity” (cf. , (10.4.4)), (3.5) can be rewritten as
$$\frac{_{w\widehat{W}}ϵ(w)e^{\psi _0^{}(w(\widehat{\rho }))}}{D_𝔨}$$
(3.6)
Let $`W_{\sigma ,0}^{}`$ be the set of minimal right coset representatives of $`\widehat{W}_{\sigma ,0}`$ in $`\widehat{W}`$. We can rewrite (3.6) as
$$\frac{_{uW_{\sigma ,0}^{}}ϵ(u)_{w\widehat{W}_𝔨}ϵ(w)e^{w\psi _0^{}(u(\widehat{\rho }))\widehat{\rho }_𝔨}}{_{\alpha \widehat{\mathrm{\Delta }}_𝔨^+}(1e^\alpha )^{m_\alpha }}$$
Using the Weyl-Kac Character formula in the formulation (2.6), the final outcome is that, if $`r`$ is even,
$$ch(X_r^+)ch(X_r^{})=\underset{uW_{\sigma ,0}^{}}{}ϵ(u)ch(L(\psi _0^{}(u\widehat{\rho })\widehat{\rho }_𝔨))$$
(3.7)
(cf. ). Using (3.3), we obtain the following result
###### Theorem 3.5.
If $`𝔨`$ is semisimple and $`r`$ is even then for $`ϵ=0`$ or $`1`$ one has:
$$ch(L(\stackrel{~}{\mathrm{\Lambda }}_ϵ))=\underset{\genfrac{}{}{0pt}{}{uW_{\sigma ,0}^{}}{\mathrm{}(u)ϵmod\mathrm{\hspace{0.17em}2}}}{}ch\left(L(\psi _0^{}(u\widehat{\rho }\widehat{\rho })+\underset{S}{}j_S\mathrm{\Lambda }_0^S+\frac{1}{2}ϵ\delta _𝔨)\right),$$
(3.8)
where $`W_{\sigma ,0}^{}`$ is the set of minimal right coset representatives of $`\widehat{W}_{\sigma ,0}`$ in $`\widehat{W}`$, $`\widehat{W}_{\sigma ,0}`$ being given by Corollary (3.3) and $`\psi _0`$ is defined by (3.1).
###### Proof.
If $`\lambda `$ is a weight of $`X_r^+`$ then $`\lambda (d_𝔨)`$, while, if $`\lambda `$ is a weight of $`X_r^{}`$ then $`\lambda (d_𝔨)\frac{1}{2}+`$ (cf. Proposition 2.5). It follows that $`X_r^+`$ and $`X_r^{}`$ do not have common components. ∎
### 3.1 Decomposition rules and combinatorics of roots
Let $`\mathrm{\Sigma }`$ denote the set of $`𝔟_0`$-stable abelian subspaces of $`𝔭`$. Each abelian subspace in $`\mathrm{\Sigma }`$ is a sum of $`𝔥_0`$-weight spaces. We identify $`𝔦\mathrm{\Sigma }`$ and the set $`A\mathrm{\Delta }(𝔭)`$ such that $`𝔦=_{\alpha A}𝔭_\alpha `$. In this section, we describe the connection between the subspaces in $`\mathrm{\Sigma }`$ and the decomposition of the basic and vector modules $`L(\stackrel{~}{\mathrm{\Lambda }}_ϵ)`$ stated in Theorem 3.5.
The set $`\mathrm{\Sigma }`$ has been studied in in the case of $`𝔤`$ simple. The results of that paper can be easily extended to the complex case. In the complex case, the subspaces in $`\mathrm{\Sigma }`$ turn out to correspond to more familiar objects. Indeed, we shall see that we can view $`\mathrm{\Sigma }`$ as the set of abelian ideals of $`𝔟_0`$. We deal at once with this special case.
We recall some general conventions and facts. Let $`𝔩`$ be a simple Lie algebra, $`\mathrm{\Phi }`$ its root system, $`𝔟_𝔩`$ a Borel subalgebra, $`\mathrm{\Phi }^+`$ and $`S`$ the corresponding set of positive roots and simple roots, respectively. A subset $`A`$ of $`\mathrm{\Phi }^+`$ is called abelian if $`\alpha +\beta \mathrm{\Phi }`$ for all $`\alpha ,\beta A`$. An abelian ideal of $`\mathrm{\Phi }^+`$ is an abelian set $`A`$ such that if $`\alpha A`$ and $`\gamma ,\alpha +\gamma \mathrm{\Phi }^+`$, then $`\alpha +\gamma A`$. If $`A`$ is an abelian ideal of $`\mathrm{\Phi }^+`$, then $`_{\alpha A}𝔩_\alpha `$ is an abelian ideal of $`𝔟_𝔩`$, and, conversely, each abelian ideal of $`𝔟_𝔩`$ is (uniquely) obtained in this way. Recently, there has been a great deal of work on these ideals by several authors (Kostant , Cellini-Papi , Panyushev , Suter ). There are various explicit descriptions of them and, in particular, we know that they are exactly $`2^{\mathrm{rank}(𝔩)}`$.
Now let $`\overline{𝔨}`$ be a simple Lie algebra, $`𝔤=\overline{𝔨}\overline{𝔨}`$, and $`\sigma :(x,y)(y,x)`$ be the switch automorphism of $`𝔤`$. Thus $`𝔨`$ is the diagonal copy of $`\overline{𝔨}`$ in $`𝔤`$, and $`𝔭=\{(x,x)x\overline{𝔨}\}`$. Then $`𝔭`$ is naturally isomorphic to $`𝔨`$ as a $`𝔨`$-module, so that what we are going to study is the decomposition of the basic and vector representations of $`\widehat{so(𝔨)}`$ with respect to $`\widehat{𝔨}`$, where $`𝔨`$ is any simple Lie algebra.
Clearly, $`\mathrm{\Delta }(𝔭)=\mathrm{\Delta }_𝔨\{0\}`$ and, through the natural isomorphism between $`𝔨`$ and $`𝔭`$, $`𝔭_\alpha `$ corresponds to $`𝔨_\alpha `$ for all $`\alpha \mathrm{\Delta }(𝔭)`$, where we intend $`𝔨_0=𝔥_0`$. In particular, if a subset $`A`$ of $`\mathrm{\Delta }(𝔭)`$ belongs to $`\mathrm{\Sigma }`$, then by definition $`_{\alpha A}𝔭_\alpha `$ is a $`𝔟_0`$ stable abelian subspace of $`𝔭`$ and therefore $`_{\alpha A}𝔨_\alpha `$ is a $`𝔟_0`$ stable abelian subspace of $`𝔨`$. It is easily seen that this implies $`A\mathrm{\Delta }_𝔨^+`$, and hence that $`A`$ is an abelian ideal of $`\mathrm{\Delta }_𝔨^+`$.
Thus, in this case, $`\mathrm{\Sigma }`$ is the set of abelian ideals of $`\mathrm{\Delta }_𝔨^+`$. The following theorem, which is an easy consequence of Theorem 3.5 and the results of , describes the decomposition of the basic and vector representations $`L(\stackrel{~}{\mathrm{\Lambda }}_0)`$ and $`L(\stackrel{~}{\mathrm{\Lambda }}_1)`$ of $`\widehat{so(𝔨)}`$ with respect to $`\widehat{𝔨}`$ in terms of $`\mathrm{\Sigma }`$ (cf. with , formula (4.2.13)). It is the nicest special case of Theorem 3.8 below.
Let us fix some notation. Set
$$\widehat{𝔥}_{}=\text{Span }_{}(\alpha _1^{},\mathrm{},\alpha _n^{})+K^{}+d^{}$$
and
$$\widehat{𝔥}_1^{}=\{x\widehat{𝔥}_{}^{}(x,\delta )=1\},\widehat{𝔥}_0^{}=\{x\widehat{𝔥}_{}^{}(x,\delta )=0\}.$$
Let $`\pi `$ be the canonical projection mod $`\delta `$ and set $`𝔥_1^{}=\pi \widehat{𝔥}_1^{}`$. We identify $`𝔥_0^{}`$ with $`\pi \widehat{𝔥}_0^{}`$.
For $`\alpha \widehat{\mathrm{\Delta }}^+`$ set
$$H_\alpha =\{x𝔥_1^{}(\alpha ,x)=0\}$$
and $`H_\alpha ^+=\{x𝔥_1^{}(\alpha ,x)0\}`$. Also, let $`C_1`$ be the fundamental alcove of $`\widehat{W}`$,
$$C_1=\{x𝔥_1^{}(\alpha ,x)0\alpha \widehat{\mathrm{\Pi }}\}.$$
It is well-known that there is a faithful action of $`\widehat{W}`$ on $`𝔥_1^{}`$ and that $`C_1`$ is a fundamental domain for this action.
For $`w\widehat{W}`$ we set
$$N(w)=\{\alpha \widehat{\mathrm{\Delta }}^+w^1(\alpha )\widehat{\mathrm{\Delta }}^{}\}.$$
Finally, for $`A\mathrm{\Sigma }`$, we denote by $`A`$ (resp. $`|A|`$) the sum (resp. the number) of elements in $`A`$.
###### Theorem 3.6.
Let $`ϵ=0\text{ or }1`$. Then one has the following decomposition of the basic and vector $`\widehat{so(𝔨)}`$-modules with respect to $`\widehat{𝔨}`$.
$$L(\stackrel{~}{\mathrm{\Lambda }}_ϵ)=\underset{\genfrac{}{}{0pt}{}{A\mathrm{\Sigma }}{|A|ϵmod\mathrm{\hspace{0.17em}2}}}{}L(h_𝔨^{}\mathrm{\Lambda }_0^𝔨+A\frac{1}{2}(|A|ϵ)\delta _𝔨)$$
(where $`h_𝔨^{}`$ and $`\mathrm{\Lambda }_0^𝔨`$ are respectively the dual Coxeter number and $`0`$th fundamental weight of $`\widehat{𝔨}`$).
Moreover, the highest weight vector $`v_A`$ of the submodule $`L(h_𝔨^{}\mathrm{\Lambda }_0^𝔨+A\frac{1}{2}(|A|ϵ)\delta _𝔨)`$ is, up to a constant factor, the following pure spinor (of the spin representation of $`Cl_0(\stackrel{~}{𝔨})`$):
$$v_A=\underset{\alpha A}{}(t^1e_\alpha ).$$
###### Proof.
Under our assumptions, the summation $`_Sj_S\mathrm{\Lambda }_0^S`$ in (3.5) has a single summand, which is $`j_𝔨\mathrm{\Lambda }_0^𝔨`$. Clearly, $`\widehat{L}(𝔤,\sigma )`$ is isomorphic to $`\widehat{𝔨}`$, hence $`h^{}=h_𝔨^{}`$. Since $`k=2`$, using (2.8), we obtain that $`j_𝔨\mathrm{\Lambda }_0^𝔨=h_𝔨^{}\mathrm{\Lambda }_0^𝔨`$.
Now we note that $`\widehat{\mathrm{\Delta }}=^{}\delta (\mathrm{\Delta }_𝔨+\delta )`$, and $`\widehat{\mathrm{\Delta }}_𝔨=\psi _0^{}(2^{}\delta (\mathrm{\Delta }_𝔨+2\delta ))`$. Hence, by Lemma 3.2, we obtain that $`\widehat{W}_{\sigma ,0}`$ is the subgroup of $`\widehat{W}`$ generated by the reflections with respect to roots in $`\mathrm{\Delta }_𝔨+2\delta `$. Then $`\widehat{W}_{\sigma ,0}`$ is isomorphic to $`\widehat{W}`$ itself and, moreover, it has $`2C_1`$ as an alcove. More precisely, $`2C_1`$ is the fundamental alcove of $`W_{\sigma ,0}`$ if we choose $`\mathrm{\Delta }_𝔨+2\delta \widehat{\mathrm{\Delta }}^+`$ as positive system for the real root system $`\mathrm{\Delta }_𝔨+2\delta `$. By general theory, it follows that
$$W_{\sigma ,0}^{}=\{w\widehat{W}wC_12C_1\}.$$
Let $`A`$ be an abelian ideal of $`\mathrm{\Delta }_𝔨^+`$ and consider the set $`A+\delta \widehat{\mathrm{\Delta }}^+`$. It is easy to prove that both $`A+\delta `$ and its complement in $`\widehat{\mathrm{\Delta }}^+`$ are closed under root addition, and hence that there exists a unique element $`w_A\widehat{W}`$ such that $`A+\delta =N(w_A)`$. Moreover, in it is proved that $`Aw_A`$ is a bijection between the set $`\mathrm{\Sigma }`$ of abelian ideals of $`\mathrm{\Delta }_𝔨^+`$ and the subset $`\{w\widehat{W}wC_12C_1\}`$. Now the claim follows directly from Theorem 3.5 and the following observations:
1. for $`w\widehat{W}`$, $`w(\widehat{\rho })\widehat{\rho }=N(w)`$ (see e.g. , Corollary 1.3.22);
2. for $`\alpha \widehat{\mathrm{\Delta }}^+`$, $`\psi _0^{}(\alpha +\delta )=\alpha +\frac{\delta _𝔨}{2}`$;
3. for $`A\mathrm{\Sigma }`$, $`ϵ(w_A)=(1)^{|A|}`$, and $`|A|=|A+\delta |=|N(w_A)|=\mathrm{}(w_A)`$.
The statement on highest weight vectors will be proved in Theorem 3.9.∎
We now turn to the combinatorial interpretation of the decomposition (3.8) for general $`𝔤=𝔨+𝔭`$ with $`𝔨`$ semisimple. We need to recall some results and notation from .
If $`𝔤`$ is simple, by the classification of Lie algebra involutions (see , Ch.8), we have that there exists an index $`p`$ such that $`s_p=1`$ and $`s_i=0`$ if $`ip`$.
Set
$$D_\sigma =\underset{w𝒲_{ab}^\sigma }{}wC_1,$$
(3.9)
where
$$𝒲_{ab}^\sigma =\{w\widehat{W}N(w)\{\alpha \widehat{\mathrm{\Delta }}m_p(\alpha )=1\}\}.$$
and, as above, $`N(w)=\{\alpha \widehat{\mathrm{\Delta }}^+w^1(\alpha )\widehat{\mathrm{\Delta }}^{}\}`$. If $`w𝒲_{ab}^\sigma `$ we shall say that $`w`$ is $`\sigma `$-minuscule.
Given $`w\widehat{W}`$, a root $`\beta \widehat{\mathrm{\Delta }}^+`$ belongs to $`N(w)`$ if and only if $`H_\beta `$ separates $`wC_1`$ and $`C_1`$. It follows that
$$D_\sigma =\underset{\genfrac{}{}{0pt}{}{\alpha \widehat{\mathrm{\Delta }}^+,}{m_p(\alpha )1}}{}H_\alpha ^+.$$
If $`\alpha \widehat{\mathrm{\Delta }}`$ we set $`\widehat{H}_\alpha =\{x\widehat{𝔥}_{}^{}/\delta (\alpha ,x)=0\}`$ and $`\widehat{H}_\alpha ^+=\{x\widehat{𝔥}_{}^{}/\delta (\alpha ,x)0\}`$. Set also
$$C_\sigma =\underset{\genfrac{}{}{0pt}{}{\alpha \widehat{\mathrm{\Delta }}^+,}{m_p(\alpha )1}}{}\widehat{H}_\alpha ^+.$$
Obviously
$$D_\sigma =C_\sigma 𝔥_1^{}.$$
Set $`\mathrm{\Phi }_\sigma =\widehat{\mathrm{\Pi }}_{\sigma ,0}\{k_p\delta +\alpha _p\},`$ where $`k_p=k,\mathrm{\hspace{0.17em}1}`$ according to whether $`\alpha _p`$ is long or short.
###### Proposition 3.7.
We have that
$$C_\sigma =\underset{\alpha \mathrm{\Phi }_\sigma }{}\widehat{H}_\alpha ^+.$$
Now let
$$P_\sigma =\underset{\alpha \widehat{\mathrm{\Pi }}_{\sigma ,0}}{}H_\alpha ^+.$$
(3.10)
It is a standard fact that the set of elements $`w\widehat{W}`$ such that $`wC_1`$ cover the polytope $`P_\sigma `$ is the set of minimal right coset representatives for the subgroup of $`\widehat{W}`$ generated by $`s_\alpha `$ with $`\alpha \widehat{\mathrm{\Pi }}_{\sigma ,0}`$, which, by Corollary 3.3, happens to be $`\widehat{W}_{\sigma ,0}`$. Since obviously $`D_\sigma P_\sigma `$, we have that $`𝒲_{ab}^\sigma W_{\sigma ,0}^{}`$. We elucidate the precise relation between $`𝒲_{ab}^\sigma `$ and $`W_{\sigma ,0}^{}`$ in the next proposition where, if $`\widehat{L}(𝔤,\sigma )`$ is simply laced, we regard all roots as long. Recall that we denote by $`𝔟_0`$ the Borel subalgebra of $`𝔨`$ associated to our initial choice of positive roots in $`\mathrm{\Delta }_𝔨`$. Recall that we are assuming that $`𝔨`$ is semisimple. If $`𝔤`$ is simple, let $`p`$ be the index such that $`s_p=1`$ and $`s_i=0`$ if $`ip`$. In the complex case (and only in this case) we have $`p=0`$ (see Remark 2.2). For $`\gamma Q^{}`$, we denote by $`t_\gamma `$ the translation by $`\gamma `$ (see \[12, (6.5.2)\]).
###### Proposition 3.8.
, 1). $`D_\sigma =P_\sigma `$ if and only if $`p=0`$ or $`\alpha _p`$ is short. If $`\alpha _p`$ is a long root, then $`P_\sigma D_\sigma `$ consists exactly of the alcove $`w_\sigma C_1`$, where $`w_\sigma =t_{k\alpha _p^{}}w_0^{}w_0`$, $`w_0`$ is the longest element of the parabolic subgroup of $`\widehat{W}`$ generated by $`\widehat{\mathrm{\Pi }}\{\alpha _p\}`$ and $`w_0^{}`$ is the longest element of the parabolic subgroup of $`\widehat{W}`$ generated by $`\widehat{\mathrm{\Pi }}\alpha _p^{}`$.
2). There is a bijection between $`𝔟_0`$-stable abelian subspaces in $`𝔭`$ and $`\sigma `$-minuscule elements, or, equivalently, alcoves paving $`D_\sigma `$. In this correspondence an element $`w𝒲_{ab}^\sigma `$ such that $`N(w)=\{\beta _1,\mathrm{},\beta _r\}`$ maps to $`\underset{i=1}{\overset{r}{}}𝔭_{\overline{\beta _i}}`$.
Set $`\mathrm{\Lambda }_{0,𝔨}=j_S\mathrm{\Lambda }_0^S`$ and let $`\mathrm{\Sigma }`$ denote the set of $`𝔟_0`$-stable abelian subspaces of $`𝔭`$. Identify $`𝔦\mathrm{\Sigma }`$ and the set $`A\mathrm{\Delta }(𝔭)`$ such that $`𝔦=_{\alpha A}𝔭_\alpha `$. We summarize the connection between abelian subspaces and the decomposition of $`X_r`$ in the following proposition.
###### Theorem 3.9.
(1) Assume that $`p=0`$ or $`\alpha _p`$ is a short root. Then
$$L(\stackrel{~}{\mathrm{\Lambda }}_ϵ)=\underset{\genfrac{}{}{0pt}{}{A\mathrm{\Sigma }}{|A|ϵmod\mathrm{\hspace{0.17em}2}}}{}L\left(\mathrm{\Lambda }_{0,𝔨}+A\frac{1}{2}(|A|ϵ)\delta _𝔨\right).$$
(2) Assume that $`\alpha _p`$ is a long root and $`p0`$. We have
$$L(\stackrel{~}{\mathrm{\Lambda }}_ϵ)=\underset{\genfrac{}{}{0pt}{}{A\mathrm{\Sigma }}{|A|ϵmod\mathrm{\hspace{0.17em}2}}}{}L(\mathrm{\Lambda }_{0,𝔨}+A\frac{1}{2}(|A|ϵ)\delta _𝔨))\nu L(\mathrm{\Lambda }_{0,𝔨}y+\frac{1}{2}ϵ\delta _𝔨),$$
where
$$y:=\psi _0^{}(N(w_\sigma ))=(\underset{\beta (\alpha _p+\mathrm{\Delta }_𝔨^+)\widehat{\mathrm{\Delta }}^+}{}\overline{\beta })+2\overline{\alpha }_p+\left(\frac{|(\alpha _p+\mathrm{\Delta }_𝔨^+)\widehat{\mathrm{\Delta }}^+|}{2}+2\right)\delta _𝔨$$
and $`\nu =\delta _{ϵ,\mathrm{}(w_\sigma )mod\mathrm{\hspace{0.17em}2}}`$.Moreover, in both cases, the highest weight vector of each component is, up to a constant factor, the pure spinor (of the spin representation of $`Cl_r(\stackrel{~}{𝔭})`$):
$$v_A=\underset{\alpha A}{}(t^{r^{}2}e_\alpha )$$
(3.11)
where $`𝔭_\alpha =e_\alpha `$. A highest weight vector in the component indexed by $`w_\sigma `$ is
$$v_\sigma =\underset{\beta (\alpha _p+\mathrm{\Delta }_𝔨^+)\widehat{\mathrm{\Delta }}^+}{}(t^{r^{}2}e_{\overline{\beta }})(t^{r^{}2}e_{\overline{\alpha }_p})(t^{r^{}3}e_{\overline{\alpha }_p}).$$
(3.12)
###### Proof.
It follows immediately from (3.3) that
$$\psi _0^{}(u\widehat{\rho })\widehat{\rho }_𝔨=\mathrm{\Lambda }_{0,𝔨}\psi _0^{}(N(u)).$$
(3.13)
hence we can rewrite formulas (3.8) as
$$ch(X_r)=ch(L(\stackrel{~}{\mathrm{\Lambda }}_0))+ch(L(\stackrel{~}{\mathrm{\Lambda }}_1))=\underset{wW_{\sigma ,0}^{}}{}ch(L(\mathrm{\Lambda }_{0,𝔨}\psi _0^{}(N(w))).$$
By Corollary 3.3 and Proposition 3.8, we can write
$$ch(X_r)=\underset{w𝒲_{ab}^\sigma }{}ch(L(\mathrm{\Lambda }_{0,𝔨}\psi _0^{}(N(w)))$$
if $`p=0`$ or $`\alpha _p`$ is short, while
$$ch(X_r)=\underset{w𝒲_{ab}^\sigma \{w_\sigma \}}{}ch(L(\mathrm{\Lambda }_{0,𝔨}\psi _0^{}(N(w)))$$
if $`p0`$ and $`\alpha _p`$ is long. If $`\alpha \widehat{\mathrm{\Delta }}`$ then $`\psi _0^{}(\alpha )=\frac{1}{2}m_p(\alpha )\delta _𝔨+\overline{\alpha }`$. If $`w=w_A`$ for some $`A\mathrm{\Sigma }`$ and $`\alpha N(w_A)`$ then $`m_p(\alpha )=1`$. Moreover, if $`w_A𝒲_{ab}^\sigma `$ encodes the subspace $`A`$, we have $`ϵ(w_A)=(1)^{\mathrm{}(w_A)}=(1)^{|A|}`$. This justifies the distribution of the summands in the basic and vector modules according to the parity of $`|A|`$.
The calculation of $`N(w_\sigma )`$ follows by a straightforward computation using standard properties of the sets $`N(w)`$ (see , 2.5). One gets
$$N(w_\sigma )=\left(\alpha _p+\mathrm{\Delta }_𝔨^+\right)\widehat{\mathrm{\Delta }}^+\{\alpha _p\}\{\alpha _p+k\delta \}.$$
(3.14)
Since $`k\delta =k_{i=0}^na_is_i\delta =2\delta ^{}`$ if we apply $`\psi _0^{}`$ to each element in the r.h.s. of (3.14) and take the sum we obtain the required expression for $`y`$.
We now check that $`v_A`$ is a highest weight vector. Set $`\lambda _A=\psi _0^{}(w_A(\widehat{\rho }))\widehat{\rho }_𝔨`$ be the corresponding highest weight. We will show that, if $`\alpha \widehat{\mathrm{\Pi }}_𝔨`$, then $`\lambda _A+\alpha `$ is not a weight of $`X_r`$. Indeed $`\lambda _A+\widehat{\rho }_𝔨=\psi _0^{}(w_A(\widehat{\rho }))`$ and $`\alpha =\psi _0^{}(\beta )`$ with $`\beta \widehat{\mathrm{\Pi }}_{\sigma ,0}`$ so we can write $`\lambda _A+\widehat{\rho }_𝔨+\alpha =\psi _0^{}(w_A(\widehat{\rho })+\beta )`$. We observe that
$$(w_A(\widehat{\rho })+\beta ,w_A(\widehat{\rho })+\beta )=(\widehat{\rho },\widehat{\rho })+(\beta ,\beta )+2(w_A(\widehat{\rho }),\beta ).$$
Since $`w_A(C_1)P_\sigma `$, we have that $`2(w_A(\widehat{\rho }),\beta )0`$, hence
$$(w_A(\widehat{\rho })+\beta ,w_A(\widehat{\rho })+\beta )>(\widehat{\rho },\widehat{\rho }).$$
If $`\lambda _A+\alpha `$ is a weight of $`X_r`$, then, according to (3.4), $`\lambda _A+\alpha +\widehat{\rho }_𝔨_Sj_S\mathrm{\Lambda }_0^S+\widehat{\rho }_𝔨\psi _0^{}(S)`$, where $`S`$ is the set of weights defined in Lemma 3.2.3 of , thus we can write that $`\psi _0^{}(w_A(\widehat{\rho })+\beta )\psi _0^{}(\widehat{\rho }S)`$. It follows that $`w_A(\widehat{\rho })+\beta \widehat{\rho }S`$. Applying Lemma 3.2.4 of (with $`\mu =w_A(\widehat{\rho })+\beta \widehat{\rho }`$), we find a contradiction. Obviously the same argument applies also to $`w_\sigma `$. ∎
## 4 Decomposition of the spin representation (semisimple case)
We now consider the case when $`r`$ is odd and $`𝔨`$ is semisimple. We distinguish two cases: $`𝔤`$ not simple (the complex case) and $`𝔤`$ simple.
### 4.1 Complex case
We consider here the case when $`𝔤=𝔨\times 𝔨,\sigma (X,Y)=(Y,X)`$, $`𝔨`$ is simple and embeds in $`𝔤`$ diagonally. We have that $`\mathrm{\Delta }(𝔭)=\mathrm{\Delta }_𝔨\{0\}`$ and we can choose $`\mathrm{\Delta }^+(𝔭)=\mathrm{\Delta }_𝔨^+`$. In this case $`𝔨`$ is simple, so, by (2.8), the sum $`_Sj_S\mathrm{\Lambda }_{0,S}`$ reduces to one summand, which equals $`\widehat{\rho }_𝔨`$. By (2.10) the character of $`X_r`$ is
$`ch(X_r)=`$
$`=e^{\widehat{\rho }_𝔨}2^{\frac{n}{2}}\left({\displaystyle \underset{\alpha \widehat{\mathrm{\Delta }}_𝔨^+}{}}(1+e^\alpha )^{m_\alpha }\right)=e^{\widehat{\rho }_𝔨}2^{\frac{n}{2}}{\displaystyle \frac{\left(_{\alpha \widehat{\mathrm{\Delta }}_𝔨^+}(1e^{2\alpha })^{m_\alpha }\right)}{\left(_{\alpha \widehat{\mathrm{\Delta }}_𝔨^+}(1e^\alpha )^{m_\alpha }\right)}}`$
$`=e^{\widehat{\rho }_𝔨}2^{\frac{n}{2}}{\displaystyle \frac{_{w\widehat{W}}ϵ(w)e^{2w\widehat{\rho }_𝔨2\widehat{\rho }_𝔨}}{\left(_{\alpha \widehat{\mathrm{\Delta }}_𝔨^+}(1e^\alpha )^{m_\alpha }\right)}}=2^{\frac{n}{2}}{\displaystyle \frac{_{w\widehat{W}}ϵ(w)e^{w(\widehat{\rho }_𝔨+\widehat{\rho }_𝔨)\widehat{\rho }_𝔨}}{\left(_{\alpha \widehat{\mathrm{\Delta }}_𝔨^+}(1e^\alpha )^{m_\alpha }\right)}}`$
$`=2^{\frac{n}{2}}L(\widehat{\rho }_𝔨).`$
Thus, we obtain
###### Proposition 4.1.
(, 4.2.2). In the complex case the spin representation of $`𝔨\times 𝔨`$ restricts to $`2^{\frac{rk(𝔨)}{2}}`$ copies of the $`\widehat{𝔨}`$-module $`L(\widehat{\rho }_𝔨)`$.
### 4.2 $`𝔤`$ simple case
We assume now that $`r`$ is odd and $`𝔨`$ is a semisimple symmetric subalgebra of a simple algebra $`𝔤`$.
#### Structure theory.
By the classification of Lie algebra involutions (see , Ch.8), we have that there exists $`p\{0,\mathrm{},n\}`$ such that $`ka_p=2`$ and $`s_p=1`$ while $`s_i=0`$ for $`ip`$. Set $`\varpi _p`$ to be the unique element of $`𝔥_0`$ such that $`\overline{\alpha }_i(\varpi _p)=\delta _{ip}`$ for $`i=1,\mathrm{},n`$. Set
$$\mu =\sigma \mathrm{exp}(\pi iad(\varpi _p)).$$
Let $`𝔨_\mu `$ denote the set of $`\mu `$-fixed points in $`𝔤`$.
It is easy to show that $`𝔥_0`$ is a Cartan subalgebra of $`𝔨_\mu `$: if $`𝔥^{}`$ is a Cartan subalgebra of $`𝔨_\mu `$ containing $`𝔥_0`$ then $`[𝔥^{},\varpi _p]=0`$, so $`𝔥^{}𝔨`$. This implies $`𝔥^{}=𝔥_0`$. If $`m`$ is a positive integer such that $`\sigma ^m=\mu ^m=id`$, then, by Proposition 8.5 of (with notation therein) the map $`t^{\frac{m}{2}\varpi _p}`$ is an isomorphism $`L(𝔤,\mu ,m)L(𝔤,\sigma ,m)`$. In particular the linear map $`t_p:\overline{\alpha }+i\delta ^{}\overline{\alpha }\frac{m}{2}(i\overline{\alpha }(\varpi _p))\delta ^{}`$ defines a bijection between $`\widehat{\mathrm{\Delta }}`$ and the set of $`\widehat{𝔥}`$-roots of $`\widehat{L}(𝔤,\mu ,m)`$. It follows that $`t_p(\widehat{\mathrm{\Pi }})`$ is a set of simple roots for $`t_p(\widehat{\mathrm{\Delta }}^+)`$.
If $`z`$, set $`L(𝔤,\mu ,m)_z=\{xL(𝔤,\mu ,m)d^{}x=zx\}`$. Since $`t_p(\alpha _i)=\overline{\alpha }_i`$ if $`i>0`$, $`t_p(\alpha _0)=\frac{m}{ka_0}\delta ^{}+\overline{\alpha }_0`$ and $`𝔨_\mu =L(𝔤,\mu ,m)_0`$, we see that the set of $`𝔥_0`$-roots of $`𝔨_\mu `$ is $`\mathrm{\Delta }_f:=\{\overline{\alpha }\alpha \widehat{\mathrm{\Delta }},m_0(\alpha )=0\}`$.
Clearly $`\mathrm{\Pi }_f=\{\overline{\alpha }_1,\mathrm{},\overline{\alpha }_n\}`$ is a set of simple roots for $`\mathrm{\Delta }_f`$ and the corresponding set of positive roots is $`\mathrm{\Delta }_f^+=\{\overline{\alpha }\alpha \widehat{\mathrm{\Delta }}^+,m_0(\alpha )=0\}`$.
#### Explicit description of $`\mathrm{\Delta }_𝔨`$ and $`\mathrm{\Delta }(𝔭)`$.
Set
$$\mathrm{\Delta }_f^0=\{\alpha \mathrm{\Delta }_f,\alpha (\varpi _p)0mod2\},\mathrm{\Delta }_f^1=\{\alpha \mathrm{\Delta }_f\alpha (\varpi _p)1mod2\}$$
and let $`\mathrm{\Delta }_{f,s}`$ and $`\mathrm{\Delta }_{f,l}`$ be, respectively, the set of short and long roots in $`\mathrm{\Delta }_f`$. We let $`\mathrm{\Delta }_x^ϵ=\mathrm{\Delta }_x\mathrm{\Delta }_f^ϵ`$ ($`x=f,l`$ or $`f,s`$; $`ϵ=0,1`$).
Recall from Section 2 our classification of $`𝔷`$-roots of $`𝔤`$ into complex, compact, and noncompact roots. Set
$`\mathrm{\Delta }_{cx}=\{\alpha \mathrm{\Delta }(𝔭)\alpha =\beta _{|𝔥_0},\text{ }\beta \text{ complex}\},`$
$`\mathrm{\Delta }_{ci}=\{\alpha \mathrm{\Delta }_𝔨\alpha =\beta _{|𝔥_0},\text{ }\beta \text{ compact}\},`$
$`\mathrm{\Delta }_{ni}=\{\alpha \mathrm{\Delta }(𝔭)\alpha =\beta _{|𝔥_0},\text{ }\beta \text{ noncompact}\};`$
$`\widehat{\mathrm{\Delta }}_{cx}=\{i\delta _𝔨+\alpha i,\alpha \mathrm{\Delta }_{cx}\},`$
$`\widehat{\mathrm{\Delta }}_{ci}=\{i\delta _𝔨+\alpha i,\alpha \mathrm{\Delta }_{ci}\},`$
$`\widehat{\mathrm{\Delta }}_{ni}=\{i\delta _𝔨+\alpha i,\alpha \mathrm{\Delta }_{ni}\}.`$
If $`k=1`$ then $`𝔷=𝔥_0`$ and $`\sigma `$ is of inner type. It follows that $`\sigma =\mathrm{exp}(\pi ih)`$ for some $`h𝔥_0`$. Since $`\sigma (X_j)=e^{\pi i\overline{\alpha }_j(h)}X_j=e^{\pi is_j}X_j`$ for $`j=1\mathrm{},n`$, we see that $`\sigma =\mathrm{exp}(\pi iad(\varpi _p))`$ and $`\mu =id`$. Hence, in this case,
$$\mathrm{\Delta }_{cx}=\mathrm{},\mathrm{\Delta }(𝔭)=\mathrm{\Delta }_f^1=\mathrm{\Delta }_{ni},\mathrm{\Delta }_𝔨=\mathrm{\Delta }_f^0=\mathrm{\Delta }_{ci}.$$
(4.1)
Suppose now that $`k=2`$, so that $`\delta ^{}=\delta `$. Recall from 2.1 that $`\alpha \mathrm{\Delta }`$ is a noncompact root if and only if $`\delta +\alpha _{|𝔥_0}`$ is a long root of $`\widehat{\mathrm{\Delta }}`$, $`\alpha `$ is compact if and only if $`\alpha _{|𝔥_0}`$ is a long root of $`\widehat{\mathrm{\Delta }}`$, and $`\alpha `$ is complex if and only if $`\alpha _{|𝔥_0}\widehat{\mathrm{\Delta }}`$ and it is not a long root.
Assume that $`k=2`$ and $`\widehat{L}(𝔤,\sigma )`$ is not of type $`A_{2n}^{(2)}`$. The following relations hold.
$$\mathrm{\Delta }_{|𝔥_0}=(\mathrm{\Delta }(𝔭)\{0\})\mathrm{\Delta }_𝔨=\overline{\widehat{\mathrm{\Delta }}}\{0\}=\mathrm{\Delta }_f$$
The first equality is clear, the second depends on the fact that $`\widehat{\mathrm{\Delta }}`$ is the set of roots of $`\widehat{L}(𝔤,\sigma )`$, whereas the third follows from the explicit description of $`\widehat{\mathrm{\Delta }}`$ given in Proposition 6.3 a) of . From the above discussion it follows that
$$\mathrm{\Delta }_{cx}=\mathrm{\Delta }_{f,s},\mathrm{\Delta }_{ci}=\mathrm{\Delta }_{f,l}^0,\mathrm{\Delta }_{ni}=\mathrm{\Delta }_{f,l}^1.$$
(4.2)
Moreover
$$\mathrm{\Delta }(𝔭)=\mathrm{\Delta }_{ni}\mathrm{\Delta }_{cx}=\mathrm{\Delta }_{f,l}^1\mathrm{\Delta }_{f,s},\mathrm{\Delta }_𝔨=\mathrm{\Delta }_{ci}\mathrm{\Delta }_{cx}=\mathrm{\Delta }_{f,l}^0\mathrm{\Delta }_{f,s}.$$
If $`\widehat{L}(𝔤,\sigma )`$ is of type $`A_{2n}^{(2)}`$, then $`\mathrm{\Delta }_𝔨`$ is the subsystem of $`\widehat{\mathrm{\Delta }}`$ generated by $`\{\alpha _0,\mathrm{},\alpha _{n1}\}`$. It follows that $`\mathrm{\Delta }_𝔨`$ does not contain long roots of $`\widehat{\mathrm{\Delta }}`$, hence $`\mathrm{\Delta }_𝔨=\mathrm{\Delta }_{cx}`$. Since $`\overline{\widehat{\mathrm{\Delta }}}=\mathrm{\Delta }_f\frac{1}{2}\mathrm{\Delta }_{f,l}\{0\}`$ (see again , Prop. 6.3 b)), arguing as above we have
$$\mathrm{\Delta }_{ni}=\mathrm{\Delta }_{f,l}\mathrm{\Delta }_{cx}=\mathrm{\Delta }_𝔨=\frac{1}{2}\mathrm{\Delta }_{f,l}\mathrm{\Delta }_{f,s}.$$
(4.3)
As we have seen in Section 2.3, the explicit realization of the spin module depends on the choice of a set of positive roots $`\mathrm{\Delta }`$ for $`𝔤`$ that is compatible with $`\mathrm{\Delta }_𝔨^+`$. We make a particular choice that we now explain.
Let $`u`$ be the longest element in the Weyl group of $`𝔨`$, $`u^{}`$ the longest element in the parabolic subgroup corresponding to $`\mathrm{\Pi }_𝔨\{\alpha _0\}`$, and $`w_0=uu^{}`$. Clearly $`w_0`$ stabilizes both $`\mathrm{\Delta }_𝔨`$ and $`\mathrm{\Delta }(𝔭)`$, hence $`\mathrm{\Delta }_{|𝔥_0}w_0(\frac{1}{2}\mathrm{\Delta }_f\mathrm{\Delta }_f)`$. It is easy to see that $`\mathrm{\Delta }_𝔨^+w_0(\frac{1}{2}\mathrm{\Delta }_f^+\mathrm{\Delta }_f^+)`$. It follows that
$$\mathrm{\Delta }^+=\{\alpha \mathrm{\Delta }\alpha _{|𝔥_0}w_0(\frac{1}{2}\mathrm{\Delta }_f^+\mathrm{\Delta }_f^+)\}$$
is a positive set of roots for $`\mathrm{\Delta }`$ compatible with $`\mathrm{\Delta }_𝔨^+`$. Recall that we set $`\mathrm{\Delta }^+(𝔭)=\mathrm{\Delta }_{|𝔥_0}^+\mathrm{\Delta }(𝔭)`$. We let
$`\mathrm{\Delta }_{cx}^+=\mathrm{\Delta }_{cx}\mathrm{\Delta }_𝔨^+,\mathrm{\Delta }_{ci}^+=\mathrm{\Delta }_{ci}\mathrm{\Delta }_𝔨^+,\mathrm{\Delta }_{ni}^+=\mathrm{\Delta }_{ni}\mathrm{\Delta }^+(𝔭)`$
$`\widehat{\mathrm{\Delta }}_a^+=\mathrm{\Delta }_a^+\{j\delta _𝔨+\alpha j>0,\alpha \mathrm{\Delta }_a\}(a=cx,ci,ni).`$
#### The algebras $`L^{}(𝔤,\sigma )`$.
Recall that $`(,)_n`$ denotes the normalized invariant form on $`𝔤`$. Since there is $`\overline{\alpha }_i\mathrm{\Delta }_f`$ such that $`\alpha _i`$ is long in $`\widehat{\mathrm{\Delta }}`$, it follows that $`(,)_n_{|𝔨_\mu }`$ is the normalized invariant form on $`𝔨_\mu `$. If $`\mathrm{\Delta }_f`$ is a root system of type $`Y_n`$, we can realize the affine Lie algebra of type $`Y_n^{(1)}`$ as the subalgebra $`\widehat{𝔨}_\mu =L(𝔨_\mu )K^{}d^{}`$ of $`\widehat{L}(𝔤)`$. We set $`(,)=(,)_n`$ in (2.1), so that $`K^{}`$ is the canonical central element of $`\widehat{𝔨}_\mu `$. We denote by $`\widehat{\mathrm{\Delta }}_\mu `$ the set of roots of $`\widehat{𝔨}_\mu `$ with respect to $`\widehat{𝔥}`$ and by $`\widehat{W}_{𝔨_\mu }`$ its Weyl group. If $`\overline{\theta }_f`$ is the highest root of $`\mathrm{\Delta }_f`$ with respect to $`\mathrm{\Pi }_f`$, then $`\widehat{\mathrm{\Pi }}_\mu =\{\overline{\theta }_f+\delta ^{},\overline{\alpha }_1,\mathrm{},\overline{\alpha }_n\}`$ is a set of simple roots of $`\widehat{𝔨}_\mu `$ with respect to $`\widehat{𝔥}`$. With this choice of the simple roots, the set of positive roots is
$$\widehat{\mathrm{\Delta }}_\mu ^+=\mathrm{\Delta }_f^+((\mathrm{\Delta }_f\{0\})+^+\delta ^{}).$$
Let $`\mathrm{\Lambda }_\mu `$ be the linear functional on $`\widehat{𝔥}`$ which maps $`K^{}`$ to $`1`$ and $`𝔥_0+d^{}`$ to $`0`$.
Let $`(,)^\mu `$ be the normalized invariant form of $`\widehat{𝔨}_\mu `$ such that $`(\mathrm{\Lambda }_\mu ,\mathrm{\Lambda }_\mu )^\mu =0`$ and let $`\nu :\widehat{𝔥}\widehat{𝔥}^{}`$ be the isomorphism induced by $`(,)^\mu `$. For any subset $`R`$ of real roots in $`\widehat{\mathrm{\Delta }}_\mu `$ we set $`R^{}=\{\nu (\alpha ^{})\alpha R\}`$. Then it is clear that, if $`A`$ is the generalized Cartan matrix of $`\widehat{𝔨}_\mu `$, then $`(\widehat{𝔥},\widehat{\mathrm{\Pi }}_\mu ^{},\nu ^1(\widehat{\mathrm{\Pi }}_\mu ))`$ is a realization of the Cartan matrix $`{}_{}{}^{t}A`$. Let $`\widehat{𝔨}_\mu ^{}=𝔤({}_{}{}^{t}A)`$ be the twisted affine algebra corresponding to the given realization of $`{}_{}{}^{t}A`$. By general theory of root systems, the set of real roots of $`\widehat{𝔨}_\mu ^{}`$ is $`\widehat{\mathrm{\Delta }}_{\mu ,re}^{}`$, where $`\widehat{\mathrm{\Delta }}_{\mu ,re}`$ is the set of real roots of $`\widehat{𝔨}_\mu `$. Since $`(,)^\mu `$ is a normalized form on $`\widehat{𝔨}_\mu `$, we have that the set of imaginary roots for $`\widehat{𝔨}_\mu ^{}`$ is $`^{}\delta ^{}`$. It follows that the set of roots of $`\widehat{𝔨}_\mu ^{}`$ is $`\widehat{\mathrm{\Delta }}_\mu ^{}:=\widehat{\mathrm{\Delta }}_{\mu ,re}^{}^{}\delta ^{}`$. Observe that if $`\widehat{L}(𝔤,\sigma )`$ is of type $`X_N^{(2)}=A_{2l1}^{(2)},D_{l+1}^{(2)},E_6^{(2)}`$, then $`\widehat{𝔨}_\mu ^{}`$ is of type $`X_N^{}^{(2)}=D_{l+1}^{(2)},A_{2l1}^{(2)},E_6^{(2)}`$ respectively (see , 13.9). Moreover, since $`\nu (\alpha ^{})=\frac{2}{(\alpha ,\alpha )^\mu }\alpha `$, the Weyl group of $`\widehat{𝔨}_\mu ^{}`$ is $`\widehat{W}_{𝔨_\mu }`$.
###### Remark 4.1.
If $`rk(𝔤)=N`$, then the number of short roots in $`\mathrm{\Pi }_f^{}`$ is $`2nN`$, therefore, as a root of $`\widehat{𝔨}_\mu ^{}`$, $`j\delta ^{}`$ has multiplicity $`2nN`$ if $`j`$ is odd and $`n`$ if $`j`$ is even.
We define the Lie algebra $`L^{}(𝔤,\sigma )`$ as follows
$$L^{}(𝔤,\sigma )=\{\begin{array}{cc}\widehat{𝔨}_\mu \hfill & \text{if }k=1\text{,}\hfill \\ \widehat{𝔨}_\mu ^{}\hfill & \text{if }k=2\text{ and }a_0=1\text{,}\hfill \\ \widehat{L}(𝔤,\sigma )^{}\hfill & \text{if }k=a_0=2\text{.}\hfill \end{array}$$
In the last case $`\widehat{L}(𝔤,\sigma )`$ is of type $`A_{2n}^{(2)}`$ and $`\widehat{L}(𝔤,\sigma )^{}`$ is realized with a construction analogous to that performed for $`\widehat{𝔨}_\mu ^{}`$, using the normalized invariant form of $`\widehat{L}(𝔤,\sigma )`$. In particular, the set of roots of $`L^{}(𝔤,\sigma )`$ is
$$(\frac{1}{2}\mathrm{\Delta }_{f,l}+\frac{1}{2}(21)\delta ^{})(\mathrm{\Delta }_{f,s}+\delta ^{})(\mathrm{\Delta }_{f,l}+(2)\delta ^{})^{}\delta ^{}.$$
(4.4)
We will denote by $`\widehat{\mathrm{\Delta }}^{}`$ the set of roots of $`L^{}(𝔤,\sigma )`$ in all cases. We choose $`(\widehat{\mathrm{\Delta }}^{})^+=(\frac{1}{2}\widehat{\mathrm{\Delta }}_\mu ^+\widehat{\mathrm{\Delta }}_\mu ^+)\widehat{\mathrm{\Delta }}^{}`$ as a set of positive roots and notice that the corresponding set $`\widehat{\mathrm{\Pi }}^{}`$ of simple roots is $`\widehat{\mathrm{\Pi }}_\mu `$, $`\widehat{\mathrm{\Pi }}_\mu ^{}`$, and $`\{\frac{1}{2}(\delta ^{}\overline{\theta }_f),\overline{\alpha }_1,\mathrm{},\overline{\alpha }_n\}`$ if $`a_0k=1`$, $`2`$, and $`4`$ respectively. Let $`\widehat{\rho }^{}`$ denote the sum of the fundamental weights of $`L^{}(𝔤,\sigma )`$. Observe that the Weyl group of $`L^{}(𝔤,\sigma )`$ is $`\widehat{W}_{𝔨_\mu }`$.
#### The map $`\psi _1:\widehat{𝔥}_𝔨\widehat{𝔥}`$.
We already observed that $`(,)_n{}_{|𝔨_\mu }{}^{}=(,)_{|𝔨_\mu }^\mu `$. Also recall that we let $`c_S=\frac{(h,h)}{(h,h)_S}`$, where $`h`$ is any nonzero element of $`𝔥_0𝔨_S`$. It follows from the discussion preceding Lemma 2.4 that
$$(h,h)^\mu =kc_S(h,h)_S.$$
Consider the linear map $`\varphi :\widehat{𝔥}_𝔨\widehat{𝔥}`$ defined by
$$\varphi _{|𝔥_0}=id_{𝔥_0},\varphi (d_𝔨)=d^{},\varphi (K_S)=kc_SK^{}.$$
We define
$$\psi _1=\varphi w_0^1$$
(4.5)
so that
$$\psi _1^{}=w_0\varphi ^{}.$$
It is clear that $`\psi _1`$ is surjective, hence $`\psi _1`$ is injective. We denote by $`\psi _1^{}^1`$ the inverse of $`\psi _1^{}:\widehat{𝔥}^{}\psi _1^{}(\widehat{𝔥}^{})`$.
It is immediate from the definition of $`\psi _1`$ that
$`\psi _1^{}(\mathrm{\Lambda }_\mu )={\displaystyle \underset{S}{}}kc_S\mathrm{\Lambda }_0^S,`$ (4.6)
$`\psi _1^{}(\delta ^{})=\delta _𝔨,`$ (4.7)
$`\psi _1^{}(\lambda )=w_0(\lambda )\text{ for }\lambda 𝔥_0^{}.`$ (4.8)
Note that, by (4.7), (4.8) and relation $`w_0(\mathrm{\Delta }_𝔨)\mathrm{\Delta }_𝔨`$ we have that $`\widehat{\mathrm{\Delta }}_𝔨\psi _1^{}(\widehat{𝔥}^{})`$, hence $`\psi _1^{}(\widehat{𝔥}^{})`$ is $`\widehat{W}_𝔨`$-stable.
###### Lemma 4.2.
For $`\alpha \widehat{\mathrm{\Delta }}_𝔨`$, let $`\beta `$ be the unique element of $`\widehat{\mathrm{\Delta }}_\mu `$ such that $`\psi _1^{}(\beta )`$ is a multiple of $`\alpha `$. Let $`s_\alpha :\widehat{𝔥}_𝔨^{}\widehat{𝔥}_𝔨^{}`$ be the reflection with respect to $`\alpha `$ and $`s_\beta ^{}:\widehat{𝔥}^{}\widehat{𝔥}^{}`$ the reflection with respect to $`\beta `$. Then
$$\psi _{1}^{}{}_{}{}^{1}s_\alpha \psi _1^{}=s_\beta ^{}.$$
###### Proof.
The proof is the same as for Lemma 3.2.∎
###### Remark 4.2.
We set $`\widehat{W}_{\sigma ,1}=(\psi _1^{})^1\widehat{W}_𝔨\psi _1^{}`$. Lemma 4.2 says that $`\widehat{W}_{\sigma ,1}`$ is a subgroup of $`\widehat{W}_{𝔨_\mu }`$.
If $`\widehat{L}(𝔤,\sigma )`$ is not of type $`A_{2n}^{(2)}`$ then, by Lemma 4.2, $`\widehat{W}_{\sigma ,1}`$ is generated by the reflection $`s_\alpha `$ with $`\alpha `$ a real root in $`\psi _1^{}{}_{}{}^{1}(\widehat{\mathrm{\Delta }}_𝔨)`$. By (4.7), and (4.1)–(4.2), we have that the set of real roots in $`\psi _1^{}{}_{}{}^{1}(\widehat{\mathrm{\Delta }}_𝔨)`$ is
$$\widehat{\mathrm{\Delta }}_{\sigma ,1}:=(\mathrm{\Delta }_{cx}\mathrm{\Delta }_{ci})+\delta ^{}.$$
(4.9)
If $`\widehat{L}(𝔤,\sigma )`$ is of type $`A_{2n}^{(2)}`$, then, by Lemma 4.2 and (4.3), we have that $`\widehat{W}_{\sigma ,1}`$ is the subgroup of $`\widehat{W}_{𝔨_\mu }`$ generated by all reflections with respect to roots in
$$\widehat{\mathrm{\Delta }}_{\sigma ,1}:=(\mathrm{\Delta }_{f,s}+\delta ^{})(\mathrm{\Delta }_{f,l}+2\delta ^{}).$$
(4.10)
#### The character.
We set
$$\widehat{\mathrm{\Delta }}_{re}^+(𝔭)=\mathrm{\Delta }^+(𝔭)\{\alpha +j\delta _𝔨\alpha \mathrm{\Delta }(𝔭)\{0\},j^+\}.$$
From (2.10) we obtain directly
$$ch(X_r)=2^{\frac{Nn}{2}}e^{_Sj_S\mathrm{\Lambda }_0^S+\rho _n}\underset{j>0}{}(1+e^{j\delta _𝔨})^{Nn}\underset{\alpha \widehat{\mathrm{\Delta }}_{re}^+(𝔭)}{}(1+e^\alpha ).$$
(4.11)
Recall that
$$D_𝔨=e^{\widehat{\rho }_𝔨}\underset{i>0}{}(1e^{i\delta _𝔨})^n\underset{\alpha (\widehat{\mathrm{\Delta }}_𝔨^+)_{re}}{}(1e^\alpha ),$$
and set
$$\rho ^{}=\underset{S}{}j_S\mathrm{\Lambda }_0^S+\rho _n+\widehat{\rho }_𝔨.$$
Then dividing and multiplying (4.11) by $`D_𝔨`$ yields
$$ch(X_r)=2^{\frac{Nn}{2}}D_𝔤^+/D_𝔨,$$
(4.12)
where
$`D_𝔤^+`$ $`=e^\rho ^{}{\displaystyle \underset{i>0}{}}(1+e^{i\delta _𝔨})^{Nn}{\displaystyle \underset{i>0}{}}(1e^{i\delta _𝔨})^n`$
$`\times {\displaystyle \underset{\alpha \widehat{\mathrm{\Delta }}_{ni}^+}{}}(1+e^\alpha ){\displaystyle \underset{\alpha \widehat{\mathrm{\Delta }}_{cx}^+}{}}(1e^{2\alpha }){\displaystyle \underset{\alpha \widehat{\mathrm{\Delta }}_{ci}^+}{}}(1e^\alpha ).`$
If $`\widehat{L}(𝔤,\sigma )`$ is not of type $`A_{2n}^{(2)}`$ set
$`D_𝔤^{}`$ $`=e^\rho ^{}{\displaystyle \underset{i>0}{}}(1+e^{i\delta _𝔨})^{Nn}{\displaystyle \underset{i>0}{}}(1e^{i\delta _𝔨})^n`$
$`\times {\displaystyle \underset{\alpha \widehat{\mathrm{\Delta }}_{ni}^+\widehat{\mathrm{\Delta }}_{ci}^+}{}}(1e^\alpha ){\displaystyle \underset{\alpha \widehat{\mathrm{\Delta }}_{cx}^+}{}}(1e^{2\alpha }).`$
Observe that $`D_𝔤^{}`$ differs from $`D_𝔤^+`$ just in the product over $`\widehat{\mathrm{\Delta }}_{ni}^+`$.
If $`\widehat{L}(𝔤,\sigma )`$ is of type $`A_{2n}^{(2)}`$ then set
$`\widehat{\mathrm{\Delta }}_{ni}^{even}=\mathrm{\Delta }_{ni}^+(\mathrm{\Delta }_{ni}+2^+\delta _𝔨)`$
$`\widehat{\mathrm{\Delta }}_{ni}^{odd}=\widehat{\mathrm{\Delta }}_{ni}^+\widehat{\mathrm{\Delta }}_{ni}^{even}`$
$`\widehat{\mathrm{\Delta }}_{f,cx}^+=(\mathrm{\Delta }_{cx}^+\mathrm{\Delta }_f)(\mathrm{\Delta }_{cx}\mathrm{\Delta }_f+^+\delta _𝔨).`$
Recalling that in this case $`N=2n`$ and that, by (4.3), $`\mathrm{\Delta }_{cx}=\frac{1}{2}\mathrm{\Delta }_{ni}(\mathrm{\Delta }_{cx}\mathrm{\Delta }_f)`$, we can rewrite $`D_𝔤^+`$ as
$`D_𝔤^+`$ $`=e^\rho ^{}{\displaystyle \underset{i>0}{}}(1+e^{i\delta _𝔨})^n{\displaystyle \underset{i>0}{}}(1e^{i\delta _𝔨})^n`$
$`\times {\displaystyle \underset{\alpha \widehat{\mathrm{\Delta }}_{ni}^{odd}}{}}(1+e^\alpha ){\displaystyle \underset{\alpha \widehat{\mathrm{\Delta }}_{ni}^{even}}{}}(1+e^\alpha ){\displaystyle \underset{\alpha \widehat{\mathrm{\Delta }}_{ni}^{even}}{}}(1e^\alpha ){\displaystyle \underset{\alpha \widehat{\mathrm{\Delta }}_{f,cx}^+}{}}(1e^{2\alpha })`$
$`=e^\rho ^{}{\displaystyle \underset{i>0}{}}(1e^{2i\delta _𝔨})^n{\displaystyle \underset{\alpha \widehat{\mathrm{\Delta }}_{ni}^{odd}}{}}(1+e^\alpha ){\displaystyle \underset{\alpha \widehat{\mathrm{\Delta }}_{ni}^{even}}{}}(1e^{2\alpha }){\displaystyle \underset{\alpha \widehat{\mathrm{\Delta }}_{f,cx}^+}{}}(1e^{2\alpha }).`$
In this case we set
$$D_𝔤^{}=e^\rho ^{}\underset{i>0}{}(1e^{2i\delta _𝔨})^n\underset{\alpha \widehat{\mathrm{\Delta }}_{ni}^{odd}}{}(1e^\alpha )\underset{\alpha \widehat{\mathrm{\Delta }}_{ni}^{even}}{}(1e^{2\alpha })\underset{\alpha \widehat{\mathrm{\Delta }}_{f,cx}^+}{}(1e^{2\alpha }),$$
that differs from $`D_𝔤^+`$ just in the product over $`\widehat{\mathrm{\Delta }}_{ni}^{odd}`$.
First we show how to compute $`D_𝔤^{}`$ and then we shall compute $`D_𝔤^+`$ from $`D_𝔤^{}`$.
###### Lemma 4.3.
$$\rho ^{}=\psi _1^{}(a_0\widehat{\rho }^{})$$
###### Proof.
We start from the equal rank case. In this case $`\widehat{\rho }^{}=h^{}\mathrm{\Lambda }_\mu +\rho `$, where $`\rho `$ is half the sum of the roots in $`\mathrm{\Delta }_f^+`$. It follows that $`\psi _1^{}(\widehat{\rho }^{})=h^{}\psi _1^{}(\mathrm{\Lambda }_\mu )+\psi _1^{}(\rho )`$. Since
$$\psi _1^{}(\rho )=w_0(\rho )=\frac{1}{2}\underset{\alpha \mathrm{\Delta }_𝔨^+}{}\alpha +\frac{1}{2}\underset{\alpha \mathrm{\Delta }^+(𝔭)}{}\alpha =\rho _𝔨+\rho _n$$
we can write that $`\psi _1^{}(\widehat{\rho }^{})=h^{}\psi _1^{}(\mathrm{\Lambda }_\mu )+\rho _n+\rho _𝔨`$. By (4.6), $`\psi _1^{}(\mathrm{\Lambda }_\mu )=c_S\mathrm{\Lambda }_0^S`$ Hence, by (2.8), we conclude that $`\psi _1^{}(\widehat{\rho }^{})=j_S\mathrm{\Lambda }_0^S+\rho _n+\widehat{\rho }_𝔨`$ as desired.
If $`k=2`$ and $`\widehat{L}(𝔤,\sigma )`$ is not of type $`A_{2n}^{(2)}`$, denoting by $`(h^{})^{}`$ the dual Coxeter number of $`L^{}(𝔤,\sigma )`$, we have $`\widehat{\rho }^{}=(h^{})^{}\mathrm{\Lambda }_\mu +\rho ^{}`$, where $`\rho ^{}`$ is half the sum of the roots in $`(\mathrm{\Delta }_f^+)^{}`$ and in turn $`\psi _1^{}(\widehat{\rho }^{})=(h^{})^{}\psi _1^{}(\mathrm{\Lambda }_\mu )+\psi _1^{}(\rho ^{})`$. Since
$$\psi _1^{}(\rho ^{})=w_0(\rho ^{})=\frac{1}{2}\underset{\alpha \mathrm{\Delta }_{ni}^+}{}\alpha +\frac{1}{2}\underset{\alpha \mathrm{\Delta }_{ci}^+}{}\alpha +\underset{\alpha \mathrm{\Delta }_{cx}^+}{}\alpha =\rho _𝔨+\rho _n$$
we need only to check that $`\psi _1^{}((h^{})^{}\mathrm{\Lambda }_\mu )=_S(j_S+h_S^{})\mathrm{\Lambda }_0^S`$. But $`\psi _1^{}((h^{})^{}\mathrm{\Lambda }_\mu )=_S(h^{})^{}2c_S\mathrm{\Lambda }_0^S`$ which equals, by (2.8), $`(h^{})^{}_S\frac{j_s+h_S^{}}{h^{}}\mathrm{\Lambda }_0^S`$. The claim follows because $`h^{}=(h^{})^{}`$.
Finally, if $`\widehat{L}(𝔤,\sigma )`$ is of type $`A_{2n}^{(2)}`$, we have $`2\widehat{\rho }^{}=h^{}\mathrm{\Lambda }_\mu +2\rho `$. Now, from (4.3), we obtain that
$`\psi _1^{}(2\rho )`$ $`={\displaystyle \underset{\alpha \mathrm{\Delta }_{f,s}^+}{}}\alpha +{\displaystyle \underset{\alpha \mathrm{\Delta }_{f,l}^+}{}}\alpha `$
$`={\displaystyle \frac{1}{2}}({\displaystyle \underset{\alpha \mathrm{\Delta }_{f,l}^+}{}}\alpha +{\displaystyle \underset{\alpha \frac{1}{2}\mathrm{\Delta }_{f,l}^+}{}}\alpha +{\displaystyle \underset{\alpha \mathrm{\Delta }_{f,s}^+}{}}\alpha )`$
$`+{\displaystyle \frac{1}{2}}({\displaystyle \underset{\alpha \mathrm{\Delta }_{f,s}^+}{}}\alpha +{\displaystyle \underset{\alpha \frac{1}{2}\mathrm{\Delta }_{f,l}^+}{}}\alpha )`$
$`=\rho _n+\rho _𝔨.`$
Finally, by (2.8), $`\psi _1^{}(h^{}\mathrm{\Lambda }_\mu )=h^{}_S2c_S\mathrm{\Lambda }_0^S=h^{}_S\frac{j_S+h_S^{}}{h^{}}\mathrm{\Lambda }_0^S`$ and we conclude as above. ∎
###### Proposition 4.4.
$$D_𝔤^{}=e^{\psi _1^{}(a_0\widehat{\rho }^{})}\underset{\alpha (\widehat{\mathrm{\Delta }}^{})^+}{}(1e^{\psi _1^{}(a_0\alpha )})^{m_\alpha },$$
where $`m_\alpha `$ is the multiplicity of $`\alpha `$ as a root of $`L^{}(𝔤,\sigma )`$.
###### Proof.
If $`\widehat{L}(𝔤,\sigma )`$ is not of type $`A_{2n}^{(2)}`$, then formulas (4.6)–(4.8) imply that $`\psi _1^{}`$ is a bijection between the set $`(\widehat{\mathrm{\Delta }}_{re}^{})^+`$ of positive real roots in $`\widehat{\mathrm{\Delta }}^{}`$ and $`\widehat{\mathrm{\Delta }}_{ni}^+\widehat{\mathrm{\Delta }}_{ci}^+2\widehat{\mathrm{\Delta }}_{cx}^+`$. Hence
$$D_𝔤^{}=e^\rho ^{}\underset{i>0}{}(1+e^{i\delta _𝔨})^{Nn}\underset{i>0}{}(1e^{i\delta _𝔨})^n\underset{\alpha (\widehat{\mathrm{\Delta }}_{re}^{})^+}{}(1e^{\psi _1^{}(\alpha )}).$$
Next we observe that
$`{\displaystyle \underset{i>0}{}}(1+e^{i\delta _𝔨})^{Nn}{\displaystyle \underset{i>0}{}}(1e^{i\delta _𝔨})^n=`$ $`{\displaystyle \underset{i>0}{}}(1e^{2i\delta _𝔨})^{Nn}{\displaystyle \underset{i>0}{}}(1e^{i\delta _𝔨})^{2nN}=`$
$`{\displaystyle \underset{i>0}{}}(1e^{2i\delta _𝔨})^n{\displaystyle \underset{i>0}{}}(1e^{(2i1)\delta _𝔨})^{2nN}.`$ (4.13)
hence, using Remark 4.1,
$$D_𝔤^{}=e^\rho ^{}\underset{\alpha (\widehat{\mathrm{\Delta }}^{})^+}{}(1e^{\psi _1^{}(\alpha )})^{m_\alpha }.$$
Applying Lemma 4.3, we obtain the result in this case.
If $`\widehat{L}(𝔤,\sigma )`$ is of type $`A_{2n}^{(2)}`$, then, by (4.4), $`\psi _1^{}`$ defines a bijection between $`2(\widehat{\mathrm{\Delta }}_{re}^{})^+`$ and $`\widehat{\mathrm{\Delta }}_{ni}^{odd}2\widehat{\mathrm{\Delta }}_{ni}^{even}2\widehat{\mathrm{\Delta }}_{f,cx}^+`$ hence
$$D_𝔤^{}=e^\rho ^{}\underset{i>0}{}(1e^{2i\delta _𝔨})^n\underset{\alpha (\widehat{\mathrm{\Delta }}_{re}^{})^+}{}(1e^{2\psi _1^{}(\alpha )}).$$
By Remark 4.1, $`m(j\delta ^{})=n`$ for all $`j`$, hence
$$D_𝔤^{}=e^\rho ^{}\underset{\alpha (\widehat{\mathrm{\Delta }}^{})^+}{}(1e^{2\psi _1^{}(\alpha )})^{m_\alpha }.$$
Lemma 4.3 implies the result in this case too. ∎
Applying Weyl-Kac denominator formula we readily obtain
###### Corollary 4.5.
$$D_𝔤^{}=e^\rho ^{}\underset{w\widehat{W}_{𝔨_\mu }}{}ϵ(w)e^{a_0(\psi _1^{}(w(\widehat{\rho }^{})\widehat{\rho }^{}))}.$$
#### Decomposition of $`X_r`$.
We now show how to compute $`D_𝔤^+`$ from $`D_𝔤^{}`$. This will allow us to compute the decomposition of $`X_r`$.
For $`\gamma _1,\mathrm{},\gamma _ta_0\psi _1^{}((\widehat{\mathrm{\Delta }}^{})^+)`$, we set
$$ϵ_p(\gamma _1,\mathrm{},\gamma _t)=\{\begin{array}{cc}(1)^{|\{\gamma _1,\mathrm{},\gamma _t\}\widehat{\mathrm{\Delta }}_{ni}^{odd}|}\hfill & \text{ if }\widehat{L}(𝔤,\sigma )\text{ is of type }A_{2n}^{(2)}\hfill \\ (1)^{|\{\gamma _1,\mathrm{},\gamma _t\}\widehat{\mathrm{\Delta }}_{ni}^+|}\hfill & \text{ otherwise}.\hfill \end{array}$$
(4.14)
Set
$$h_\sigma =\{\begin{array}{cc}d_𝔨\hfill & \text{ if }\widehat{L}(𝔤,\sigma )\text{ is of type }A_{2n}^{(2)}\hfill \\ \varpi _p\hfill & \text{otherwise}.\hfill \end{array}$$
By the explicit description of $`a_0\psi _1^{}((\widehat{\mathrm{\Delta }}^{})^+)`$ given in the proof of Proposition 4.4 it is clear from (4.1)–(4.3) that $`(1)^{(\gamma _1+\mathrm{}+\gamma _t)(h_\sigma )}=ϵ_p(\gamma _1,\mathrm{},\gamma _t)`$. In particular, if we define a function $`ϵ_p`$ on the $``$-lattice $`L`$ generated by $`a_0\psi _1^{}((\widehat{\mathrm{\Delta }}^{})^+)`$ by setting
$$ϵ_p(\lambda )=(1)^{\lambda (h_\sigma )}.$$
then, if $`\lambda =\gamma _1+\mathrm{}+\gamma _t`$ with $`\gamma _ia_0\psi _1^{}((\widehat{\mathrm{\Delta }}^{})^+)`$, we have
$$ϵ_p(\lambda )=ϵ_p(\gamma _1,\mathrm{},\gamma _t).$$
(4.15)
###### Lemma 4.6.
We have
$$D_𝔤^{}=e^\rho ^{}\underset{\lambda L}{}a_\lambda e^\lambda ,$$
with $`a_\lambda `$. Moreover
$$D_𝔤^+=e^\rho ^{}\underset{\lambda L}{}ϵ_p(\lambda )a_\lambda e^\lambda .$$
###### Proof.
By Corollary 4.5
$$D_𝔤^{}=e^\rho ^{}\underset{w\widehat{W}_{𝔨_\mu }}{}ϵ(w)e^{a_0(\psi _1^{}(<N^{}(w)>))},$$
where $`N^{}(w)=\{\alpha (\widehat{\mathrm{\Delta }}^{})^+w^1(\alpha )<0\}`$, hence the first assertion follows. The second statement follows directly from (4.15) and the definition of $`D_𝔤^+`$ and $`D_𝔤^{}`$. ∎
We set $`\widehat{\mathrm{\Delta }}_{\sigma ,1}^{}=\widehat{\mathrm{\Delta }}_{\sigma ,1}`$ (see (4.9)) if $`k=1`$ or $`\widehat{L}(𝔤,\sigma )`$ is of type $`A_{2n}^{(2)}`$, while we set $`\widehat{\mathrm{\Delta }}_{\sigma ,1}^{}=(\widehat{\mathrm{\Delta }}_{\sigma ,1})^{}`$ in the other cases. We notice that $`\widehat{\mathrm{\Delta }}_{\sigma ,1}^{}`$ is a root system contained in $`\widehat{\mathrm{\Delta }}^{}`$ and its associated reflection group is $`\widehat{W}_{\sigma ,1}`$. By general theory of reflection groups (see ) the set
$$W_{\sigma ,1}^{}=\{u\widehat{W}_{𝔨_\mu }N^{}(u)\widehat{\mathrm{\Delta }}^{}\widehat{\mathrm{\Delta }}_{\sigma ,1}^{}\}$$
is a set of minimal coset representatives of $`\widehat{W}_{\sigma ,1}\backslash \widehat{W}_{𝔨_\mu }`$.
For $`w\widehat{W}_{𝔨_\mu }`$ set $`N^{}(w)=N^{}(w)\widehat{\mathrm{\Delta }}_{\sigma ,1}^{}`$. Set also $`\mathrm{}(w)=|N^{}(w)|`$ and, if $`v\widehat{W}_{\sigma ,1}`$, $`\mathrm{}^{}(v)=|N^{}(v)|`$. Now assume $`v\widehat{W}_{\sigma ,1}`$ and $`uW_{\sigma ,1}^{}`$. Since $`v(\widehat{\mathrm{\Delta }}_{\sigma ,1}^{})=\widehat{\mathrm{\Delta }}_{\sigma ,1}^{}`$, we have that $`vN^{}(u)\widehat{\mathrm{\Delta }}^{}\widehat{\mathrm{\Delta }}_{\sigma ,1}^{}`$. It is a standard fact that $`N^{}(v)\widehat{\mathrm{\Delta }}_{\sigma ,1}^{}`$. In particular $`N^{}(vu)=N^{}(v)\dot{}(vN^{}(u)(\mathrm{\Delta }_\mu ^{})^+)`$ (disjoint union), whence, $`N^{}(vu)\widehat{\mathrm{\Delta }}_{\sigma ,1}^{}=N^{}(v)\widehat{\mathrm{\Delta }}_{\sigma ,1}^{}=N^{}(v)`$.
If $`ϵ`$ and $`ϵ^{}`$ denote the sign functions in $`\widehat{W}_{𝔨_\mu }`$ and $`\widehat{W}_{\sigma ,1}`$, respectively, then $`ϵ(w)=(1)^{\mathrm{}(w)}`$ and $`ϵ^{}(v)=(1)^{\mathrm{}^{}(v)}`$. Notice that the set of real roots in $`\widehat{\mathrm{\Delta }}^{}\widehat{\mathrm{\Delta }}_{\sigma ,1}^{}`$ maps under $`a_0\psi _1^{}`$ bijectively onto $`\widehat{\mathrm{\Delta }}_{ni}^{odd}`$ if $`\widehat{L}(𝔤,\sigma )`$ is of type $`A_{2n}^{(2)}`$, and onto $`\widehat{\mathrm{\Delta }}_{ni}`$ in all the other cases, therefore
$$ϵ(vu)=ϵ^{}(v)ϵ_p(a_0\psi _1^{}N^{}(vu)).$$
(4.16)
It follows from Corollary 4.5 and (4.16) that
$$D_𝔤^{}=e^\rho ^{}\underset{v\widehat{W}_{\sigma ,1}}{}\underset{uW_{\sigma ,1}^{}}{}ϵ^{}(v)ϵ_p(a_0\psi _1^{}N^{}(vu))e^{a_0\psi _1^{}(N^{}(vu))},$$
and, by Lemma 4.6,
$$D_𝔤^+=e^\rho ^{}\underset{v\widehat{W}_{\sigma ,1}}{}\underset{uW_{\sigma ,1}^{}}{}ϵ^{}(v)e^{a_0\psi _1^{}N^{}(vu)}.$$
Clearly, if $`v\widehat{W}_𝔨`$, then $`ϵ^{}(\psi _1^{}{}_{}{}^{1}v\psi _1^{})=det(\psi _1^{}{}_{}{}^{1}v\psi _1^{})=det(v)`$. Therefore from the above equation we obtain
$$D_𝔤^+=\underset{v\widehat{W}_{\sigma ,1}}{}\underset{uW_{\sigma ,1}^{}}{}ϵ^{}(v)e^{a_0\psi _1^{}(vu\widehat{\rho }^{})}=\underset{uW_{\sigma ,1}^{}}{}\underset{v\widehat{W}_𝔨}{}det(v)e^{v(a_0\psi _1^{}(u\widehat{\rho }^{}))},$$
and, since $`ch(X_r)=2^{\frac{Nn}{2}}\frac{D_𝔤^+}{D_𝔨}`$, from (2.6) we deduce the following result.
###### Proposition 4.7.
If $`𝔨`$ is semisimple and $`r`$ is odd, then
$$ch(X_r)=2^{\frac{Nn}{2}}\underset{uW_{\sigma ,1}^{}}{}ch(L(a_0\psi _1^{}(u\widehat{\rho }^{})\widehat{\rho }_𝔨)),$$
(4.17)
where $`\psi _1`$ is defined by (4.5).
### 4.3 Combinatorial interpretation of decompositions of spin modules.
We will use the following general facts. Let $`\widehat{𝔤}_1,\widehat{𝔤}_2`$ be two affine Kac-Moody algebras, and for $`i=1,2`$, let $`\widehat{𝔥}_i`$ be a Cartan subalgebra of $`\widehat{𝔤}_i`$, $`\widehat{\mathrm{\Delta }}_i`$ be the corresponding root system, and $`\widehat{W}_i`$ be the Weyl group. Endow $`\widehat{𝔤}_i`$, with a fixed arbitrary invariant form and $`\widehat{𝔥}_i^{}`$ with the form induced by this choice. We say that a linear isomorphism $`f:\widehat{𝔥}_1^{}\widehat{𝔥}_2^{}`$ is an extended isomorphism of root systems if $`f`$ is an isometry and $`f\widehat{\mathrm{\Delta }}_1=\widehat{\mathrm{\Delta }}_2`$. For $`\alpha \widehat{\mathrm{\Delta }}_i`$ let $`s_\alpha `$ be the reflection on $`\widehat{𝔥}_i`$ with respect to $`\alpha `$. The following result is clear.
###### Lemma 4.8.
If $`f`$ is an extended isomorphism of root systems, then, for all $`\alpha \widehat{\mathrm{\Delta }}_1`$, $`fs_\alpha f^1=s_{f\alpha }`$. In particular, if $`A\widehat{\mathrm{\Delta }}_1`$, and $`W_A`$ is the subgroup of $`\widehat{W}_1`$ generated by the reflections with respect to elements of $`A`$, then $`fW_Af^1`$ is the subgroup of $`\widehat{W}_2`$ generated by the reflections with respect to elements in $`fA`$.
###### Lemma 4.9.
Let $`f:𝔥_1^{}\delta _1𝔥_2^{}\delta _2`$ be a linear map such that $`f\widehat{\mathrm{\Delta }}_1=\widehat{\mathrm{\Delta }}_2`$ and, for all $`\alpha ,\beta \widehat{\mathrm{\Delta }}_1`$, $`(\alpha ,\beta )_1=(f\alpha ,f\beta )_2`$. Then there exists a unique extension of $`f`$ to an extended isomorphism of root systems.
###### Proof.
Since $`\delta _2`$ is the orthogonal subspace of $`𝔥_2^{}\delta _2`$ in $`\widehat{𝔥}_2^{}`$, and since $`\widehat{\mathrm{\Delta }}_2`$ spans $`𝔥_2^{}\delta _2`$, the conditions $`(f\mathrm{\Lambda }_0^1,f\alpha )_2=(\mathrm{\Lambda }_0^1,\alpha )_1`$ for all $`\alpha \widehat{\mathrm{\Delta }}_1`$ determine $`f\mathrm{\Lambda }_0^1`$ modulo $`\delta _2`$. By a direct computation we see that the further condition $`(f\mathrm{\Lambda }_0^1,f\mathrm{\Lambda }_0^1)_2=0`$ determines the component in $`\delta _2`$ of $`f\mathrm{\Lambda }_0^1`$. ∎
###### Definition 4.3.
Let us say that a $`𝔥_0`$-stable subspace $`S`$ of $`𝔭`$ is noncompact if all weights of $`𝔥_0`$ on $`S`$ are in $`\mathrm{\Delta }_{ni}`$.
We will describe the decomposition of $`X_r`$ in terms of certain noncompact subspaces of $`𝔭`$. For the sake of a better exposition we discuss various cases separately: we consider the equal rank case, the case when $`\widehat{L}(𝔤,\sigma )`$ is of type $`A_{2n}^{(2)}`$ and the remaining non equal rank cases.
#### Equal rank case.
In the equal rank case all $`𝔥_0`$-stable subspaces of $`𝔭`$ are noncompact, for $`\mathrm{\Delta }(𝔭)`$ is equal to $`\mathrm{\Delta }_{ni}`$, henceforth the final outcome will be very similar to decomposition of the basic and vector representations.
Recall that in this case $`\mu =Id`$, so $`L^{}(𝔤,\sigma )=\widehat{L}(𝔤)`$ and $`\widehat{\mathrm{\Delta }}^{}=\widehat{\mathrm{\Delta }}_\mu `$. The isomorphism $`t^{\varpi _p}:L(𝔤,\mu ,2)L(𝔤,\sigma )`$ induces a linear isomorphism $`g:𝔥_0^{}+\delta ^{}𝔥_0^{}+\delta ^{}`$ such that $`g(\widehat{\mathrm{\Delta }}_\mu )=\widehat{\mathrm{\Delta }}`$. Explicitly
$$g:\lambda +j\delta ^{}\lambda +(2j+\lambda (\varpi _p))\delta ^{}.$$
(4.18)
By (2.4) it is clear that $`g`$ preserves scalar products of roots.
By (4.1), $`\mathrm{\Delta }_𝔨=\mathrm{\Delta }_f^0`$ so $`g(\widehat{\mathrm{\Delta }}_{\sigma ,1})=\mathrm{\Delta }_f^0+2\delta ^{}`$. Comparing this with (3.2) we see that
$$g(\widehat{\mathrm{\Delta }}_{\sigma ,1})=\widehat{\mathrm{\Delta }}_{\sigma ,0}.$$
(4.19)
By Lemma 4.8, $`\widehat{W}=g\widehat{W}_{𝔨_\mu }g^1`$ and, by (4.19), $`g\widehat{W}_{\sigma ,1}g^1=\widehat{W}_{\sigma ,0}`$. Recall that in this case we have that $`p>0`$ and $`a_p=2`$, hence $`g(\delta ^{}\overline{\theta })=\overline{\theta }=\overline{\alpha }_0=\alpha _0`$ and $`g(\overline{\alpha }_i)=\alpha _i`$ for $`i=1,\mathrm{},n`$. It follows that $`g(\widehat{\mathrm{\Delta }}_\mu ^+)=\widehat{\mathrm{\Delta }}^+`$ hence $`N(gug^1)=g(N^{}(u))`$, for all $`u\widehat{W}_{𝔨_\mu }`$. It follows that
$$W_{\sigma ,1}^{}=g^1W_{\sigma ,0}^{}g.$$
and $`g(\widehat{\rho }^{})=\widehat{\rho }`$.
Recalling that $`N=n`$ in this case, we can rewrite the decomposition of $`X_r`$ given in (4.17) as
$`X_r`$ $`={\displaystyle \underset{uW_{\sigma ,1}^{}}{}}L({\displaystyle \underset{S}{}}j_S\mathrm{\Lambda }_0^S+\rho _n+\psi _1^{}(u(\widehat{\rho }^{})\widehat{\rho }^{}))`$
$`={\displaystyle \underset{uW_{\sigma ,1}^{}}{}}L({\displaystyle \underset{S}{}}j_S\mathrm{\Lambda }_0^S+\rho _n\psi _1^{}(N^{}(u)))`$
$`={\displaystyle \underset{uW_{\sigma ,0}^{}}{}}L({\displaystyle \underset{S}{}}j_S\mathrm{\Lambda }_0^S+\rho _n\psi _1^{}(g^1N(u)))`$
Applying the discussion of 3.1 we deduce the analog of Theorem 3.9 for the spin representation in the equal rank case:
###### Theorem 4.10.
Set $`m=\frac{dim(𝔭)}{2}`$.
(1) Assume that $`\alpha _p`$ is a short root. Then
$$L(\stackrel{~}{\mathrm{\Lambda }}_{mϵ})=\underset{\genfrac{}{}{0pt}{}{A\mathrm{\Sigma }}{|A|ϵmod\mathrm{\hspace{0.17em}2}}}{}L\left(\mathrm{\Lambda }_{0,𝔨}+\rho _n+w_0A\right).$$
(2) Assume that $`\alpha _p`$ is a long root. We have
$`L(\stackrel{~}{\mathrm{\Lambda }}_{mϵ})`$ $`={\displaystyle \underset{\genfrac{}{}{0pt}{}{A\mathrm{\Sigma }}{|A|ϵmod\mathrm{\hspace{0.17em}2}}}{}}L(\mathrm{\Lambda }_{0,𝔨}+\rho _n+w_0Ak_A\delta _𝔨))`$
$`{\displaystyle \nu L\left(\mathrm{\Lambda }_{0,𝔨}+\rho _nw_0(y)(k_y+1)\delta _𝔨\right)},`$
where $`y`$ and $`\nu `$ are as in Theorem 3.9 (2) and $`k_A=|w_0(A)\mathrm{\Delta }^+(𝔭)|`$, $`k_y=|(\alpha _p+\mathrm{\Delta }_𝔨^+)(\delta ^{}\mathrm{\Delta }_f^+)|`$. Moreover, in both cases, the highest weight vector of each component indexed by $`A\mathrm{\Sigma }`$ is, up to a constant factor, the pure spinor (of the spin representation of $`Cl_r(\stackrel{~}{𝔭})`$):
$$v_A=\underset{\alpha w_0(A)\mathrm{\Delta }^+(𝔭)}{}(t^{r^{}2}e_\alpha )\underset{\alpha w_0(A)(\mathrm{\Delta }^+(𝔭))}{}t^{r^{}1}e_\alpha ,$$
(4.20)
where $`𝔭_\alpha =e_\alpha `$. An highest weight vector for the component indexed by $`w_\sigma `$ is
$`({\displaystyle }`$ $`{}_{\beta (\alpha _p+\mathrm{\Delta }_𝔨^+)(\delta ^{}\mathrm{\Delta }_f^+)}{}^{}t_{}^{r^{}2}e_{\overline{\beta }})`$ (4.21)
$`({\displaystyle \underset{\beta (\alpha _p+\mathrm{\Delta }_𝔨^+)(\delta ^{}+\mathrm{\Delta }_f^+)}{}}t^{r^{}1}e_{\overline{\beta }})(t^{r^{}1}e_{\overline{\alpha }_p})(t^{r^{}2}e_{\overline{\alpha }_p}).`$
###### Proof.
First of all observe that a weight vector $`v`$ is in $`X_r^+`$ if and only if its weight is equal to $`_Sj_S\mathrm{\Lambda }_0^S+\rho _n+\lambda `$, where $`\lambda `$ is a sum of an even number of elements of $`\widehat{\mathrm{\Delta }}_{ni}`$, hence $`L(_Sj_S\mathrm{\Lambda }_0^S\psi _1^{}(g^1N(u)))`$ occurs in $`X_r^+`$ if and only if $`\mathrm{}(u)`$ is even.
The rest of the result now follows as in Theorem 3.9. Only the coefficient of $`\delta _𝔨`$ needs checking. If $`A\mathrm{\Sigma }`$ and $`\alpha N(w_A)`$, then $`m_p(\alpha )=1`$ hence $`\alpha =\delta ^{}\pm \overline{\alpha }`$ with $`\overline{\alpha }\mathrm{\Delta }_f^1\mathrm{\Delta }_f^+`$. If $`\alpha =\overline{\alpha }+\delta ^{}`$ then $`g^1(\alpha )=\overline{\alpha }`$, while, if $`\alpha =\overline{\alpha }+\delta ^{}`$ then $`g^1(\alpha )=\overline{\alpha }+\delta ^{}`$. Write $`N(w_A)=\{\overline{\gamma }_1+\delta ^{},\mathrm{}\overline{\gamma }_s+\delta ^{},\overline{\beta }_1+\delta ^{},\mathrm{},\overline{\beta }_r+\delta ^{}\}`$ with $`\overline{\beta }_i,\overline{\gamma }_i\mathrm{\Delta }_f^1`$. Hence $`\psi _1^{}(g^1(N(w_A)))=w_0(\overline{\beta _i}\overline{\gamma _i})+s\delta _𝔨`$. Since $`A=\overline{N(w_A)}=\{\overline{\gamma _1},\mathrm{}\overline{\gamma _s},\overline{\beta _1},\mathrm{},\overline{\beta _r}\}`$ we have that $`s=|w_0(A)\mathrm{\Delta }^+(𝔭)|=k_A`$. The coefficient $`k_y`$ is computed similarly.
It remains to check that $`k_A=0`$ for all $`A\mathrm{\Sigma }`$ if and only if $`\alpha _p`$ is short. Let $`W_f`$ denote the Weyl group of $`\mathrm{\Delta }_f`$. Clearly $`\{wW_fw(\mathrm{\Delta }_f^+)\psi _1^{}{}_{}{}^{1}(\mathrm{\Delta }_𝔨^+)\}W_{\sigma ,1}^{}`$. By \[4, Theorem 5.12\], $`|W_{\sigma ,1}^{}|=|\{wW_fw(\mathrm{\Delta }_f^+)\psi _1^{}{}_{}{}^{1}(\mathrm{\Delta }_𝔨^+)\}|`$ if and only if $`\alpha _p`$ is short. The result follows. ∎
#### The non equal rank case with $`𝐚_\mathrm{𝟎}=\mathrm{𝟏}`$.
It is clear that $`𝔨_\mu `$ is $`\sigma `$-stable, hence we can consider the subalgebra of $`\widehat{𝔨}_\mu `$
$$\widehat{L}(𝔨_\mu ,\sigma _{|𝔨_\mu })=L(𝔨_\mu ,\sigma _{|𝔨_\mu })K^{}d^{}.$$
Clearly,
$$𝔨^{}=𝔨𝔨_\mu ,𝔭^{}=𝔭𝔨_\mu $$
are, respectively, the $`1`$ and $`1`$ eigenspaces of $`\sigma _{|𝔨_\mu }`$ on $`𝔨_\mu `$. Let us denote by $`\mathrm{\Delta }_𝔨^{}`$ the $`𝔥_0`$-roots of $`𝔨^{}`$, by $`\mathrm{\Delta }(𝔭^{})`$ the set of weights of $`𝔥_0`$ on $`𝔭^{}`$.
Since $`\sigma _{|𝔨_\mu }=\mathrm{exp}(\pi iad(\varpi _p))`$, it is clear that
$$\mathrm{\Delta }_𝔨^{}=\mathrm{\Delta }_f^0,\mathrm{\Delta }(𝔭^{})\{0\}=\mathrm{\Delta }_f^1.$$
(4.22)
For $`w\widehat{W}_{𝔨_\mu }`$, let $`N_\mu (w)=\{\alpha \widehat{\mathrm{\Delta }}_\mu ^+w^1(\alpha )<0\}`$. Observe that $`W_{\sigma ,1}^{}=\{u\widehat{W}_{𝔨_\mu }N_\mu (u)\widehat{\mathrm{\Delta }}_\mu \widehat{\mathrm{\Delta }}_{\sigma ,1}\}`$. This is because both $`W_{\sigma ,1}^{}`$ and $`\{u\widehat{W}_{𝔨_\mu }N_\mu (u)\widehat{\mathrm{\Delta }}_\mu \widehat{\mathrm{\Delta }}_{\sigma ,1}\}`$ are the set of minimal length coset representatives. We notice that the set of real roots in $`\widehat{\mathrm{\Delta }}_\mu \widehat{\mathrm{\Delta }}_{\sigma ,1}`$ equals the set of real roots in $`\widehat{\mathrm{\Delta }}^{}\widehat{\mathrm{\Delta }}_{\sigma ,1}^{}`$. The above observation implies that $`N_\mu (u)=N^{}(u)`$ for $`uW_{\sigma ,1}^{}`$. In particular
$$(u\widehat{\rho }^{}\widehat{\rho }^{})=u\widehat{\rho }_\mu \widehat{\rho }_\mu ,$$
where $`\widehat{\rho }_\mu `$ denotes the sum of the fundamental weights of $`\widehat{𝔨}_\mu `$.
This time the isomorphism
$$t^{\varpi _p}:L(𝔨_\mu ,id,2)\widehat{L}(𝔨_\mu ,\sigma _{|𝔨_\mu })$$
induces a linear isomorphism $`g:𝔥_0\delta ^{}𝔥_0\delta ^{}`$, still given by (4.18), such that $`g(\widehat{\mathrm{\Delta }}_\mu )`$ is the set $`\widehat{\mathrm{\Delta }}_{\mu ,\sigma }`$ of roots of $`\widehat{L}(𝔨_\mu ,\sigma _{|𝔨_\mu })`$. By Lemma (4.9), we can uniquely extend $`g`$ to an extended isomorphism of $`\widehat{\mathrm{\Delta }}_\mu `$ with $`\widehat{\mathrm{\Delta }}_{\mu ,\sigma }`$, which we still denote by $`g`$. We choose $`g\widehat{\mathrm{\Pi }}_\mu `$ as a set of simple roots for $`\widehat{\mathrm{\Delta }}_{\mu ,\sigma }`$, and denote by $`\widehat{\mathrm{\Delta }}_{\mu ,\sigma }^+`$ the corresponding positive system of roots. Then it is clear that $`g`$ maps $`\widehat{\mathrm{\Delta }}_\mu ^+`$ onto $`\widehat{\mathrm{\Delta }}_{\mu ,\sigma }^+`$. We denote by $`\widehat{W}_{𝔨_\mu ,\sigma }`$ the Weyl group of $`\widehat{L}(𝔨_\mu ,\sigma _{|𝔨_\mu })`$ and, for $`w\widehat{W}_{𝔨_\mu ,\sigma }`$, we denote by $`N_\sigma `$ its negative set with respect to $`\widehat{\mathrm{\Delta }}_{\mu ,\sigma }^+`$. By Lemma 4.8, $`\widehat{W}_{𝔨_\mu ,\sigma }=g\widehat{W}_{𝔨_\mu }g^1`$. Moreover, it is clear that $`N_\sigma (gug^1)=gN_\mu (u)`$, for all $`u\widehat{W}_{𝔨_\mu }`$.
Since $`W_{\sigma ,1}^{}=\{u\widehat{W}_{𝔨_\mu }N_\mu (u)\mathrm{\Delta }_{ni}+\delta ^{}\}`$, by (4.2), we have that $`gW_{\sigma ,1}^{}g^1=\{v\widehat{W}_{𝔨_\mu ,\sigma }N_\sigma (v)\mathrm{\Delta }_{f,l}^1+(1+2)\delta ^{}\}`$. Since the set of real roots in $`\widehat{\mathrm{\Delta }}_{\mu ,\sigma }\backslash (\mathrm{\Delta }(𝔨^{})+2\delta ^{})`$ is $`\mathrm{\Delta }_f^1+(1+2)\delta ^{}`$ we have in particular that $`gW_{\sigma ,1}^{}g^1`$ is precisely the set of all elements $`v`$ in $`W_{\sigma _{|𝔨_\mu },0}^{}`$ such that $`N_\sigma (v)\mathrm{\Delta }_{ni}+(1+2)\delta ^{}`$. We actually have a stronger result.
###### Lemma 4.11.
If $`v\widehat{W}_{𝔨_\mu ,\sigma }`$ is such that $`N_\sigma (v)\mathrm{\Delta }_{ni}+(1+2)\delta ^{}`$, then $`v`$ is $`\sigma _{|𝔨_\mu }`$-minuscule. In particular
$$gW_{\sigma ,1}^{}g^1=\{v𝒲_{ab}^{\sigma _{|𝔨_\mu }}\overline{N_\sigma (v)}\mathrm{\Delta }_{ni}\}.$$
###### Proof.
We recall (see ) that, if $`\sigma _{|𝔨_\mu }`$ is of type $`(s_0,\mathrm{},s_n;k)`$, then $`v`$ is $`\sigma _{|𝔨_\mu }`$-minuscule if $`ht_{\sigma _{|𝔨_\mu }}(\alpha )=1`$ for all $`\alpha N_\sigma (v),`$ where $`ht_{\sigma _{|𝔨_\mu }}(\alpha )=s_im_i(\alpha )`$.
We use the well known fact that in a finite root system a long root is the sum of two short roots. Suppose now that $`N_\sigma (v)\mathrm{\Delta }_{ni}+(1+2)\delta ^{}`$ and that $`(2m+1)\delta ^{}+\alpha `$ is in $`N_\sigma (v)`$. By (4.2) $`\alpha \mathrm{\Delta }_f^1`$, thus we can write $`\alpha =\beta +\gamma `$ with $`\beta \mathrm{\Delta }_{f,s}^1`$ and $`\gamma \mathrm{\Delta }_{f,s}^0`$. It follows that $`(2m+1)\delta ^{}+\alpha =(2m\delta ^{}+\gamma )+(\delta ^{}+\beta )`$, hence, by the biconvexity property of $`N_\sigma (v)`$, we find that $`\delta ^{}+\beta N_\sigma (v)`$ unless $`m=0`$ and $`\gamma \widehat{\mathrm{\Delta }}_{\mu ,\sigma }^+`$. If we write $`\alpha =_{i=1}^nm_i\overline{\alpha }_i`$, then $`m_p=\pm 1`$. Since $`\delta ^{}+\alpha =_{i=1}^nm_i\alpha _i`$ if $`m_p=1`$ and $`\delta ^{}+\alpha =2\delta ^{}+_{i=1}^nm_i\alpha _i`$ if $`m_p=1`$, we see that, in any case, $`ht_{\sigma _{|𝔨_\mu }}(\delta ^{}+\alpha )=1`$. ∎
We identify $`\mathrm{\Delta }_𝔨^{}`$ with the roots in $`\alpha \widehat{\mathrm{\Delta }}_{\mu ,\sigma }`$ such that $`\alpha (d^{})=0`$ and choose $`\mathrm{\Delta }_𝔨^{}^+=\widehat{\mathrm{\Delta }}_{\mu ,\sigma }^+\mathrm{\Delta }_𝔨^{}`$ as a set of positive roots for $`𝔨^{}`$. We denote by $`𝔟^{}`$ the corresponding Borel subalgebra of $`𝔨^{}`$.
###### Remark 4.4.
By the definition of $`g`$ we see that the set of simple roots for $`𝔨^{}`$ is given by
$$\mathrm{\Pi }_𝔨^{}=\{\begin{array}{cc}\{\overline{\alpha }_ii0,p\}\hfill & \text{if }\overline{\theta }_f(\varpi _p)<2\hfill \\ \{\overline{\theta }_f\}\{\overline{\alpha }_ii0,p\}\hfill & \text{if }\overline{\theta }_f(\varpi _p)=2.\hfill \end{array}$$
It follows that $`𝔟^{}=𝔟_0𝔨^{}`$.
Combining Lemma 4.11 and Remark 4.4 with the results of exposed in § 3.1, we find the analogue of Theorem 3.9 for this case. Let $`\mathrm{\Sigma }_{ni}^{}`$ be the the set of $`𝔟^{}`$-stable abelian noncompact subspaces of $`𝔭^{}`$. Recall that in section 2.3 we set $`L=Nn`$ and $`l=\frac{Nn}{2}`$.
###### Theorem 4.12.
Set $`m=\frac{dim𝔭}{2}`$.
(1) Assume that $`m`$ is even. Then
$$L(\stackrel{~}{\mathrm{\Lambda }}_{m1})=L(\stackrel{~}{\mathrm{\Lambda }}_m)=2^{l1}\underset{A\mathrm{\Sigma }_{ni}^{}}{}L\left(\mathrm{\Lambda }_{0,𝔨}+\rho _n+w_0Ak_A\delta _𝔨\right).$$
(2) Assume that $`m`$ is odd. We have
$$L(\stackrel{~}{\mathrm{\Lambda }}_m)=2^l\underset{A\mathrm{\Sigma }_{ni}^{}}{}L(\mathrm{\Lambda }_{0,𝔨}+\rho _n+w_0Ak_A\delta _𝔨)).$$
In both cases $`k_A=|w_0(A)\mathrm{\Delta }^+(𝔭)|`$. Moreover the highest weight vectors of each component indexed by $`A\mathrm{\Sigma }_{ni}^{}`$ are, up to a constant factor, the pure spinor (of the spin representation of $`Cl_r(\stackrel{~}{𝔭})`$)
$$\underset{s=0}{\overset{l}{}}\underset{l+1j_1<\mathrm{}<j_sL}{}v_{r^{}1,j_k}\underset{\alpha w_0(A)\mathrm{\Delta }^+(𝔭)}{}(t^{r^{}2}e_\alpha )\underset{\alpha w_0(A)(\mathrm{\Delta }^+(𝔭))}{}(t^{r^{}1}e_\alpha )$$
(4.23)
if $`l`$ is even, while, if $`l`$ is odd, they are
$`{\displaystyle \underset{s=0}{\overset{l}{}}}`$ $`{\displaystyle \underset{l+1<j_1<\mathrm{}<j_sL}{}}v_{r^{}1,j_k}({\displaystyle \underset{\alpha w_0(A)\mathrm{\Delta }^+(𝔭)}{}}\left((t^{r^{}2}e_\alpha )(v_{r^{}1,l+1})\right)`$
$`{\displaystyle \underset{\alpha w_0(A)(\mathrm{\Delta }^+(𝔭))}{}}\left((t^{r^{}1}e_\alpha )(t^{r^{}1}v_{r^{}1,l+1})\right).`$
###### Proof.
By a direct computation, we see that $`g^1(\alpha +\delta ^{})=\alpha +\delta ^{}`$ if $`\alpha \mathrm{\Delta }_f^+`$, while $`g^1(\alpha +\delta ^{})=\alpha `$ if $`\alpha \mathrm{\Delta }_f^+`$. We can therefore apply the proof of Theorem 4.10. We need only to check the decomposition of $`X_r^+=L(\stackrel{~}{\mathrm{\Lambda }}_m)`$ and $`X_r^{}=L(\stackrel{~}{\mathrm{\Lambda }}_{m1})`$ when $`m`$ is even, but this follows readily from the description of the highest vectors. ∎
#### The $`A_{2n}^{(2)}`$-case.
Recall that in this case $`L^{}(𝔤,\sigma )=\widehat{L}(𝔤,\sigma )^{}`$, and that we chose $`\widehat{\mathrm{\Pi }}^{}=\{\frac{1}{2}(\delta ^{}\overline{\theta }_f),\overline{\alpha }_1,\mathrm{},\overline{\alpha }_n\}`$ as root basis for $`L^{}(𝔤,\sigma )`$. From the explicit description of $`\widehat{\mathrm{\Delta }}_{\sigma ,1}`$, we obtain that $`W_{\sigma ,1}^{}`$ is the set of all elements $`v\widehat{W}_{𝔨_\mu }`$ such that $`N^{}(v)`$ is included in the set of short roots of $`\widehat{\mathrm{\Delta }}^{}`$. If we choose $`(\widehat{\mathrm{\Pi }}^{})^{}=\{\delta ^{}\overline{\theta }_f,\overline{\alpha }_1,\mathrm{},\overline{\alpha }_{n1},\frac{1}{2}\overline{\alpha }_n\}`$ as root basis for $`\widehat{L}(𝔤,\sigma )`$, we obtain $`(\widehat{\mathrm{\Delta }}^{}{}_{}{}^{+})^{}`$ as positive system for $`\widehat{L}(𝔤,\sigma )`$. Observe that $`\widehat{W}=\widehat{W}_{𝔨_\mu }`$. It is clear that if we regard $`v\widehat{W}_{𝔨_\mu }`$ as an element of $`\widehat{W}`$ and denote by $`N^{}(v)`$ the negative set of $`v`$ with respect to this choice of the positive roots, we obtain that $`N^{}(v)=(N^{}(v))^{}`$. In particular, for $`vW_{\sigma ,1}^{}`$, $`N^{}(v)=2N^{}(v)`$. Therefore, as a subset on $`\widehat{W}`$, $`W_{\sigma ,1}^{}`$ is the set of all elements in $`v`$ such that $`N^{}(v)`$ is included in the set of long roots of $`\widehat{\mathrm{\Delta }}`$. Now we observe that $`\mathrm{\Delta }_𝔨=\frac{1}{2}\mathrm{\Delta }_{f,l}\mathrm{\Delta }_{f,s}=(\mathrm{\Delta }_f)^{}`$, so $`\{\overline{\alpha }_1,\mathrm{},\overline{\alpha }_{n1},\frac{1}{2}\overline{\alpha }_n\}`$ is the set of simple roots corresponding to $`\mathrm{\Delta }_𝔨(\frac{1}{2}\mathrm{\Delta }_f^+\mathrm{\Delta }_f^+)`$. It follows that $`w_0(\{\overline{\alpha }_1,\mathrm{},\overline{\alpha }_{n1},\frac{1}{2}\overline{\alpha }_n\})=\{\alpha _0,\mathrm{},\alpha _{n1}\}`$. Since $`\overline{\theta }_f=2\overline{\alpha }_1+\mathrm{}+2\overline{\alpha }_{n1}+\overline{\alpha }_n`$ we see that $`w_0(\delta ^{}\overline{\theta }_f)=\delta ^{}+\overline{\alpha }_n=\alpha _n`$, hence $`w_0((\widehat{\mathrm{\Pi }}^{})^{})=\widehat{\mathrm{\Pi }}`$ and $`w_0((\widehat{\mathrm{\Delta }}^{}{}_{}{}^{+})^{})=\widehat{\mathrm{\Delta }}^+`$. This says that $`w_0W_{\sigma ,1}^{}w_0^1`$ is the set of elements of $`\widehat{W}`$ such that $`N(v)`$ is included in the set of long roots of $`\widehat{\mathrm{\Delta }}`$. Since the set of long roots of $`\widehat{\mathrm{\Delta }}`$ is $`\mathrm{\Delta }_{ni}+(1+2)\delta ^{}`$ we have that $`w_0^1W_{\sigma ,1}^{}w_0W_{\sigma ,0}^{}`$. Lemma 4.11 applies, so we can conclude that
$$w_0^1W_{\sigma ,1}^{}w_0=\{v𝒲_{ab}^\sigma \overline{N(v)}\mathrm{\Delta }_{ni}\}.$$
Arguing as in the previous twisted cases, we finally obtain the analogous of Theorem 3.9 for this case. Set $`\mathrm{\Sigma }_{ni}`$ to be the set of noncompact $`𝔟_0`$-stable abelian subspaces of $`\mathrm{\Delta }(𝔭)`$.
###### Theorem 4.13.
Set $`m=\frac{dim𝔭}{2}`$.
(1) Assume that $`m`$ is even. Then
$$L(\stackrel{~}{\mathrm{\Lambda }}_{m1})=L(\stackrel{~}{\mathrm{\Lambda }}_m)=2^{\frac{n}{2}1}\underset{A\mathrm{\Sigma }_{ni}}{}L\left(\mathrm{\Lambda }_{0,𝔨}+\rho _n+A|A|\delta _𝔨\right).$$
(2) Assume that $`m`$ is odd. We have
$$L(\stackrel{~}{\mathrm{\Lambda }}_m)=2^{\frac{n}{2}}\underset{A\mathrm{\Sigma }_{ni}}{}L(\mathrm{\Lambda }_{0,𝔨}+\rho _n+A|A|\delta _𝔨)).$$
Moreover the highest weight vectors of each component indexed by $`A\mathrm{\Sigma }_{ni}`$ are, up to a constant factor, the pure spinor (of the spin representation of $`Cl_r(\stackrel{~}{𝔭})`$)
$$(\underset{s=0}{\overset{\frac{n}{2}}{}}\underset{\frac{n}{2}+1j_1<\mathrm{}<j_sn}{}v_{r^{}1,j_k})\underset{\alpha A}{}(t^{r^{}2}e_\alpha )$$
if $`n`$ is even, while, if $`n`$ is odd, are
$`{\displaystyle \underset{s=0}{\overset{\frac{n}{2}}{}}}`$ $`{\displaystyle \underset{\frac{n+1}{2}<j_1<\mathrm{}<j_sn}{}}v_{r^{}1,j_k}{\displaystyle \underset{\alpha A}{}}(t^{r^{}2}e_\alpha )(t^{r^{}1}v_{r^{}1,l+1}),`$
where $`l=\frac{n}{2}`$.
###### Proof.
We know that
$$ch(X_r)=\underset{uW_{\sigma ,1}^{}}{}ch(L(\mathrm{\Lambda }_{0,𝔨}+\rho _na_0\psi _1^{}(N^{}(u))).$$
By the above discussion $`a_0\psi _1^{}(N^{}(u))=\psi _1^{}(N(w_0^1uw_0))`$ so we can write
$$ch(X_r)=\underset{A\mathrm{\Sigma }_{ni}}{}ch(L(\mathrm{\Lambda }_{0,𝔨}+\rho _n\psi _1^{}(N(w_A))).$$
The coefficient of $`\delta _𝔨`$ is computed as in 3.1. The rest of the proof follows as in the previous cases. ∎
## 5 The Hermitian symmetric case
In this section we discuss the decomposition of a conformal pair $`(so(𝔭),𝔨)`$ when $`𝔤=𝔨𝔭`$ is an infinitesimal Hermitian symmetric space. In this case there exists a node $`i0`$ such that $`a_i=1`$, $`s_0=s_i=1`$, and $`s_j=0`$ for $`j0,i`$. It turns out that $`𝔨`$ is an equal rank subalgebra of $`𝔤`$ and it is not semisimple. We can write $`𝔨=_{S>0}𝔨_S𝔨_0`$, where $`𝔨_0=\varpi _i`$ and $`\varpi _i`$ is the unique element of $`𝔥_0`$ such that $`\overline{\alpha }_j(\varpi _i)=\delta _{ij}`$ for $`j>0`$. Recall that in this case
$$\widehat{𝔨}=\widehat{[𝔨,𝔨]}\widehat{𝔨}_0,$$
where $`\widehat{𝔨}_0=[t,t^1]𝔨_0K_0`$ with bracket defined by
$$[t^nH+aK_0,t^mH+bK_0]=\delta _{n,m}(H,H)_nK_0.$$
As before $`(,)_n`$ is the normalized invariant form of $`𝔤`$. Let $`\overline{r}=0`$ if $`r`$ is even, $`\overline{r}=1`$ if $`r`$ is odd, and $`\psi _{\overline{r}}^{}`$, $`\widehat{W}_{\sigma ,\overline{r}}`$ be defined as in Section 3 or 4, according to the parity of $`r`$ (note that in this case $`a_0=k=1`$). Then let $`\psi _{\overline{r}}^{}:\widehat{𝔥}^{}(\widehat{𝔥}_𝔨)^{}`$ denote the transpose of the map $`\psi _{\overline{r}}`$ restricted to $`\widehat{𝔥}_𝔨`$.
The same computation performed in the equal rank case for $`𝔨`$ semisimple would give
$$ch(X_r^\pm )=\underset{\genfrac{}{}{0pt}{}{uW_{\sigma ,\overline{r}}^{}}{\mathrm{}(u)ϵmod\mathrm{\hspace{0.17em}2}}}{}ch(L(\psi _{\overline{r}}^{}(u\widehat{\rho }_{\overline{r}})\widehat{\rho }_𝔨))),$$
(5.1)
where $`ϵ=0,\mathrm{\hspace{0.17em}1}`$ according to whether we are considering the $`+`$ or $``$ case. Moreover $`W_{\sigma ,\overline{r}}^{}`$ is the set of minimal right coset representatives for $`\widehat{W}/\widehat{W}_{\sigma ,0}`$ if $`\overline{r}=0`$ and for $`\widehat{W}_{𝔨_\mu }/\widehat{W}_{\sigma ,1}`$ if $`\overline{r}=1`$, and $`\widehat{\rho }_0=\widehat{\rho },\widehat{\rho }_1=\widehat{\rho }^{}`$.
We first deal with the basic and vector case and then we transfer our results to the spin case via map $`g`$ defined in $`(\text{4.18})`$. So we assume $`\overline{r}=0`$. The starting point to provide a more explicit form of (5.1) is a remarkable subset of stable subpaces which has been introduced in , Section 6. Recall from (3.9) the definition of the polytope encoding $`𝔟_0`$-stable abelian subspaces of $`𝔭`$ and set
$$D_\sigma ^{}=D_\sigma \{x𝔥_1^{}(x,\alpha _i)<0\}.$$
$`D_\sigma ^{}`$ corresponds exactly to the set of stable abelian subspaces of $`𝔭`$ which include $`𝔤_{\alpha _i}`$. Let $`\omega _i^{}`$ be the unique element in $`Span_{}(\alpha _1,\mathrm{},\alpha _n)`$ such that $`(\alpha _j,\omega _i^{})=\delta _{ij}`$. From the proof of Lemma 6.1 of we deduce the following fact.
###### Lemma 5.1.
Consider the group of translations $`T_{\omega _i^{}}=\{t_{j\omega _i^{}}j\}`$. Then $`\overline{D}_\sigma ^{}`$ is a fundamental domain for the action of $`T_{\omega _i^{}}`$ on $`_{wW_{\sigma ,0}^{}}w\overline{C}_1`$.
Therefore there exists a “special” subset of stable subspaces of $`𝔭`$ such that the translates of the corresponding alcoves cover the domain $`W_{\sigma ,0}^{}\overline{C_1}`$. At this point this fact gives little information on the weights appearing in the decomposition (5.1), since $`T_{\omega _i^{}}`$ is not included in $`\widehat{W}`$. This requires some more work, which we perform in a general setting.
Let $`\widehat{𝔩}=𝔤(A)`$, where $`A`$ is a generalized Cartan matrix of affine tipe $`X_m^{(1)}`$, $`\widehat{𝔥}_𝔩`$ its Cartan subalgebra, $`\widehat{\mathrm{\Pi }}_𝔩=\{\beta _0,\beta _1,\mathrm{}\beta _m\}`$ and $`\widehat{\mathrm{\Pi }}_𝔩^{}=\{\beta _0^{},\beta _1^{},\mathrm{}\beta _m^{}\}`$ the sets of simple roots and coroots. Moreover, let $`\widehat{W}_𝔩`$ be the Weyl group of $`\widehat{𝔩}`$, $`\mathrm{\Lambda }_0^𝔩,\mathrm{\Lambda }_1^𝔩,\mathrm{},\mathrm{\Lambda }_m^𝔩`$ be the fundamental weights, and $`\widehat{\rho }_𝔩=\mathrm{\Lambda }_0^𝔩+\mathrm{}+\mathrm{\Lambda }_m^𝔩`$.
As usual we assume that $`\mathrm{\Pi }_𝔩=\{\beta _1,\mathrm{}\beta _m\}`$ has Dynkin diagram of finite type $`X_m`$, and we denote by $`𝔩`$ the corresponding finite dimensional simple Lie subalgebra of $`\widehat{𝔩}`$. Also, we denote by $`W_𝔩`$ the Weyl group of $`𝔩`$, by $`\mathrm{\Delta }_𝔩^+`$ its set of positive roots, by $`\theta _𝔩`$ its highest root, and we set $`\delta _𝔩=\beta _0\theta _𝔩`$.
Identify $`\widehat{𝔥}_𝔩`$ and $`\widehat{𝔥}_𝔩^{}`$ via the normalized invariant form. Let $`\omega _1^{},\mathrm{},\omega _m^{}`$ be the fundamental coweights of $`𝔩`$ and, for $`i\{1,\mathrm{},m\}`$, let $`w_iW_𝔩`$ be such that $`N(w_i)=\{\alpha \mathrm{\Delta }_𝔩^+(\alpha ,\omega _i^{})0\}`$. It is well-known that $`w_i`$ exists (and it is unique). We denote by $`\stackrel{~}{W}`$ the extended affine Weyl group of $`𝔩`$, i.e. $`\stackrel{~}{W}=T_{P_𝔩^{}}W_𝔩`$, where $`P_𝔩^{}`$ is the coweight lattice. We regard $`\stackrel{~}{W}`$ as a group of transformations on $`\widehat{𝔥}_𝔩^{}`$. Moreover, we set
$$Z=\{t_{\omega _i^{}}w_ii\{1,\mathrm{},m\},(\theta _𝔩,\omega _i^{})=1\}\{1\}.$$
It is well known that $`Z`$ is exactly the subgroup of all elements in $`\stackrel{~}{W}`$ that map the fundamental alcove of $`𝔩`$ to itself (see ). We may identify the fundamental alcove of $`𝔩`$ with $`C_{\widehat{𝔩}}𝔥_1^{}`$, where $`C_{\widehat{𝔩}}`$ is the fundamental chamber of $`\widehat{𝔩}`$, and $`𝔥_1^{}=(\mathrm{\Lambda }_0^𝔩+\text{Span}_{}\{\beta _0,\mathrm{},\beta _m\})/\delta _𝔩`$ (see , Section 6.6, or , Section 1). Since the restriction to $`𝔥_1^{}`$ is a faithful representation of $`\stackrel{~}{W}`$, we obtain
$$Z=\{v\stackrel{~}{W}v\widehat{\mathrm{\Pi }}_𝔩=\widehat{\mathrm{\Pi }}_𝔩\}.$$
###### Lemma 5.2.
For all $`vZ`$,
$$v\widehat{\rho }_𝔩=\widehat{\rho }_𝔩.$$
###### Proof.
We fix $`vZ\{1\}`$ and set $`v^1(\beta _i)=\beta _{j_i}`$ for $`i\{0,1\mathrm{},m\}`$. We denote by $`(,)`$ the form induced on $`\widehat{𝔥}_𝔩^{}`$ by the normalized invariant form of $`\widehat{𝔩}`$ and we recall that $`(,)`$ is invariant under $`\stackrel{~}{W}`$. Then, for $`i=0,\mathrm{},m`$,
$$v\widehat{\rho }_𝔩(\beta _i^{})=\frac{2(v\widehat{\rho }_𝔩,\beta _i)}{(\beta _i,\beta _i)}=\frac{2(\widehat{\rho }_𝔩,v^1\beta _i)}{(v^1\beta _i,v^1\beta _i)}=\widehat{\rho }_𝔩(\beta _{j_i}^{})=1.$$
It follows that $`v\widehat{\rho }_𝔩\widehat{\rho }_𝔩mod\delta _𝔩`$.
It remains to prove that $`(v\widehat{\rho }_𝔩,\mathrm{\Lambda }_0^𝔩)=0`$. We assume that $`v=t_{\omega _i^{}}w_i`$. Since $`W_𝔩`$ fixes $`\mathrm{\Lambda }_0^𝔩`$, by formula (6.5.2) of we have
$$v\mathrm{\Lambda }_0^𝔩=t_{\omega _i^{}}\mathrm{\Lambda }_0^𝔩=\mathrm{\Lambda }_0^𝔩+\omega _i^{}\frac{1}{2}|\omega _i^{}|^2\delta _𝔩.$$
(5.2)
Since $`\widehat{\rho }_𝔩=\rho _𝔩+h_𝔩^{}\mathrm{\Lambda }_0^𝔩`$, where $`\rho _𝔩`$ is the sum of fundamental weights of $`𝔩`$ and $`h_𝔩^{}`$ is its dual Coxeter number, we obtain
$$v\widehat{\rho }_𝔩=v\rho _𝔩+h_𝔩^{}(\mathrm{\Lambda }_0^𝔩+\omega _i^{}\frac{1}{2}|\omega _i^{}|^2\delta _𝔩).$$
(5.3)
But $`v\widehat{\rho }_𝔩(\rho _𝔩+h_𝔩^{}\mathrm{\Lambda }_0^𝔩)\delta _𝔩`$, hence
$$v\rho _𝔩=\rho _𝔩h_𝔩^{}\omega _i^{}+x\delta _𝔩,$$
(5.4)
for some $`x`$. It follows that
$$w_i\rho _𝔩=t_{\omega _i^{}}(\rho _𝔩h_𝔩^{}\omega _i^{}+x\delta _𝔩)=\rho _𝔩h_𝔩^{}\omega _i^{}+x\delta _𝔩(\rho _𝔩h_𝔩^{}\omega _i^{}+x\delta _𝔩,\omega _i^{})\delta _𝔩,$$
and since the component of $`w_i\rho _𝔩`$ in $`\delta _𝔩`$ is zero, we obtain
$$x=(\widehat{\rho }_𝔩,\omega _i^{})+h_𝔩^{}|\omega _i^{}|^2.$$
(5.5)
Combining equations (5.3), (5.4), and (5.5), we have
$$(v\widehat{\rho }_𝔩,\mathrm{\Lambda }_0^𝔩)=(\widehat{\rho }_𝔩,\omega _i^{})+\frac{1}{2}h^{}|\omega _i^{}|^2.$$
(5.6)
Now, since $`Z`$ is a group, there exists $`i^{}\{1\mathrm{},m\}`$ such that $`v^1=t_{\omega _i^{}^{}}w_i^{}`$. Therefore, as in (5.2), we obtain
$$v^1\mathrm{\Lambda }_0^𝔩=t_{\omega _i^{}^{}}\mathrm{\Lambda }_0^𝔩=\mathrm{\Lambda }_0^𝔩+\omega _i^{}^{}\frac{1}{2}|\omega _i^{}^{}|^2\delta _𝔩.$$
Hence
$$(v\widehat{\rho }_𝔩,\mathrm{\Lambda }_0^𝔩)=(\widehat{\rho }_𝔩,v^1\mathrm{\Lambda }_0^𝔩)=(\widehat{\rho }_𝔩,\omega _i^{}^{})\frac{1}{2}h^{}|\omega _i^{}^{}|^2.$$
(5.7)
Since $`T_{P_𝔩^{}}`$ is normal in $`\stackrel{~}{W}_𝔩`$, $`v^1=w_i^1t_{\omega _i^{}}=t_{w_i^1\omega _i^{}}w_i^1`$ and since $`T_{P_𝔩^{}}W`$ is a semidirect product,
$$w_i^1\omega _i^{}=\omega _i^{}^{}$$
(5.8)
and $`w_i^1=w_i^{}`$. This implies, in particular, that $`|\omega _i^{}|^2=|\omega _i^{}^{}|^2`$, and therefore from (5.6) and (5.7) we obtain that
$$(\widehat{\rho }_𝔩,\omega _i^{})+\frac{1}{2}h^{}|\omega _i^{}|^2=(\widehat{\rho }_𝔩,\omega _i^{}^{})\frac{1}{2}h^{}|\omega _i^{}|^2.$$
At this point, in order to conclude, it suffices to prove that
$$(\widehat{\rho }_𝔩,\omega _i^{})=(\widehat{\rho }_𝔩,\omega _i^{}^{}).$$
(5.9)
By equation (5.8), we have that $`(\widehat{\rho }_𝔩,\omega _i^{}^{})=(w_i\widehat{\rho }_𝔩,\omega _i^{})`$, hence
$$(\widehat{\rho }_𝔩,\omega _i^{})(\widehat{\rho }_𝔩,\omega _i^{}^{})=(2\widehat{\rho }_𝔩+w_i\widehat{\rho }_𝔩\widehat{\rho }_𝔩,\omega _i^{})=(\mathrm{\Delta }_𝔩^+N(w_i),\omega _i^{}).$$
By the definition of $`w_i`$, the last term of the above equalities is zero. This proves (5.9) and hence the lemma. ∎
Denote by $`\mathrm{\Sigma }^{}`$ the set of abelian $`𝔟_0`$-stable subspaces of $`𝔭`$ whose corresponding alcoves lie in $`D_\sigma ^{}`$. The previous lemma is the key to read the weight of a factor appearing in (5.1) in terms of the weight of a subspace in $`\mathrm{\Sigma }^{}`$.
###### Proposition 5.3.
If $`A=wC_1,wW_{\sigma ,0}^{}`$, then there exists a unique $`k`$ and a unique $`I\mathrm{\Sigma }^{}`$ such that
$`\psi _0^{}(w(\widehat{\rho }))\widehat{\rho }_𝔨=`$
$`\mathrm{\Lambda }_{0,𝔨}+I+kh^{}\varpi _i+({\displaystyle \frac{1}{2}}dim(I)+k(|I^+||I^{}|){\displaystyle \frac{k^2}{4}}dim(𝔭))\delta _𝔨,`$ (5.10)
where $`I^\pm =I\pm \mathrm{\Delta }^+(𝔭)`$.
###### Proof.
By Lemma 5.1 we have $`A=t_{k\omega _i^{}}(A^{})`$ for a unique $`k`$ and a unique alcove $`A^{}D_\sigma ^{}`$. Suppose that $`A^{}=w^{}C_1,w^{}\widehat{W}`$. Then $`wC_1=t_{k\omega _i^{}}w^{}C_1`$, and hence there exists a unique $`zZ`$ such that $`w=t_{k\omega _i^{}}wz`$. By Lemma 5.2 and formula (6.5.2) of we thus obtain
$`\psi _0^{}(w(\widehat{\rho }))\widehat{\rho }_𝔨`$ $`=\psi _0^{}(t_{k\omega _i^{}}w^{}z(\widehat{\rho }))\widehat{\rho }_𝔨=\psi _0^{}(t_{k\omega _i^{}}w^{}(\widehat{\rho }))\widehat{\rho }_𝔨`$
$`=\psi _0^{}(t_{k\omega _i^{}}(w^{}(\widehat{\rho })\widehat{\rho }))+\psi _0^{}(t_{k\omega _i^{}}(\widehat{\rho }))\widehat{\rho }_𝔨`$
$`=\psi _0^{}(t_{k\omega _i^{}}(I)dim(I)\delta ^{})+\psi _0^{}(t_{k\omega _i^{}}(\widehat{\rho }))\widehat{\rho }_𝔨,`$ (5.11)
where $`I`$ is the ideal in $`\mathrm{\Sigma }^{}`$ corresponding to $`w^{}`$. Note that $`\psi _0^{}(\delta ^{})=\frac{1}{2}\delta _𝔨`$ and that $`\psi _0^{}(\omega _i^{})=\nu (\varpi _i)+\frac{|\varpi _i|^2}{2}\delta _𝔨`$. Also remark that
$$(I,\omega _i^{})=|I^+||I^{}|,(\widehat{\rho },\omega _i^{})=\frac{dim(𝔭)}{4}.$$
(5.12)
Combining (5), (5.12) and formula (5.6) we get (5.3). ∎
Denote by $`c_{I,k}`$ the coefficient of $`\delta _𝔨`$ in formula (5.3). For $`q`$, denote by $`L(\stackrel{~}{\mathrm{\Lambda }}_ϵ)_q`$ the eigenspace of eigenvalue $`q`$ for the action of $`𝔨_0`$ on $`L(\stackrel{~}{\mathrm{\Lambda }}_ϵ)`$.
###### Remark 5.1.
$`L(\stackrel{~}{\mathrm{\Lambda }}_ϵ)_q`$ is non zero if and only if $`qϵmod2`$ . In fact, by (2.10), the weights of $`L(\stackrel{~}{\mathrm{\Lambda }}_ϵ)`$ are of the form $`\mathrm{\Lambda }_{0.𝔨}_{j=1}^s\gamma _j,\gamma _j\widehat{\mathrm{\Delta }}^+(𝔭),sϵmod2`$. Since $`\mathrm{\Delta }^+(𝔭)=\mathrm{\Delta }_f^1`$ in this case, we have $`(\mathrm{\Lambda }_{0.𝔨}_{j=1}^s\gamma _j)(\varpi _i)=s`$.
From the previous Proposition it follows that
###### Theorem 5.4.
$$L(\stackrel{~}{\mathrm{\Lambda }}_ϵ)_q=\underset{\genfrac{}{}{0pt}{}{I\mathrm{\Sigma }^{}}{|I^+||I^{}|qmod\frac{dim(𝔭)}{2}}}{}L(\mathrm{\Lambda }_{0,𝔨}+I+k_Ih^{}\nu (\varpi _i)+(c_{I,k_I}+\frac{ϵ}{2})\delta _𝔨),$$
where $`k_I=\frac{2(q|I^+|+|I^{}|)}{dim(𝔭)}`$.
###### Proof.
Consider the sum
$$\underset{wW_{\sigma ,0}^{}}{}ch(L(\psi _0^{}(w(\widehat{\rho }))\widehat{\rho }_𝔨)).$$
This sum makes sense because, given a weight $`\mu `$, there is only a finite number of elements $`wW_{\sigma ,0}^{}`$ such that $`\mu `$ is a weight of $`L(\psi _0^{}(w(\widehat{\rho }))\widehat{\rho }_𝔨)`$. Indeed, if $`\mu `$ occurs in $`L(\psi _0^{}(w(\widehat{\rho }))\widehat{\rho }_𝔨)`$, then $`\mu =\psi _0^{}(w(\widehat{\rho }))\widehat{\rho }_𝔨_{\alpha \widehat{\mathrm{\Delta }}_𝔨^+}n_\alpha \alpha `$, hence $`\mu (\varpi _i)=(\psi _0^{}(w(\widehat{\rho }))\widehat{\rho }_𝔨)(\varpi _i)`$. It follows from Proposition 5.3 that
$$\mu (\varpi _i)=I(\varpi _i)+kh^{}\nu (\varpi _i)(\varpi _i)$$
and there is only a finite number of $`I\mathrm{\Sigma }^{}`$ and $`k`$ that satisfy this equation.
We can therefore write
$`{\displaystyle \underset{wW_{\sigma ,0}^{}}{}}ch(L(\psi _0^{}(w(\widehat{\rho }))\widehat{\rho }_𝔨))`$ $`={\displaystyle \underset{wW_{\sigma ,0}^{}}{}}{\displaystyle \frac{_{u\widehat{W}_𝔨}ϵ(u)e^{u\psi _0^{}(w\widehat{\rho })}}{D_𝔨}}`$
$`={\displaystyle \frac{_{wW_{\sigma ,0}^{}}_{u\widehat{W}_𝔨}ϵ(u)e^{u\psi _0^{}(w\widehat{\rho })}}{D_𝔨}}`$
$`={\displaystyle \frac{D_𝔤^+}{D_𝔨}}=ch(X_r).`$
Thus we can write
$$L(\stackrel{~}{\mathrm{\Lambda }}_ϵ)=\underset{k}{}\underset{\genfrac{}{}{0pt}{}{I\mathrm{\Sigma }^{}}{|I|ϵmod\mathrm{\hspace{0.17em}2}}}{}L(\mathrm{\Lambda }_{0,𝔨}+I+kh^{}\nu (\varpi _i)+(c_{I,k}+\frac{1}{2}ϵ)\delta _𝔨).$$
(5.13)
Observe now that
$`(\mathrm{\Lambda }_{0,𝔨}+I+kh^{}\nu (\varpi _i)+c_{I,k}\delta _𝔨)(\varpi _i)`$ $`=|I^+||I^{}|+kh^{}|\varpi _i|^2`$
$`=|I^+||I^{}|+k{\displaystyle \frac{dim(𝔭)}{2}}.`$
The result follows by collecting in (5.13) the terms with eigenvalue $`q`$. ∎
Arguing as in the semisimple equal rank case we obtain, for the spin representations, the following result.
###### Theorem 5.5.
Set $`m=\frac{dim(𝔭)}{2}`$. The eigenvalues of $`\varpi _i`$ on $`L(\stackrel{~}{\mathrm{\Lambda }}_{mϵ})`$ are of the form $`\frac{dim(𝔭)}{4}+q,q,qϵmod\mathrm{\hspace{0.17em}2}`$. The corresponding eigenspaces decompose as
$$L(\stackrel{~}{\mathrm{\Lambda }}_{mϵ})_{\frac{dim(𝔭)}{4}+q}=\underset{\genfrac{}{}{0pt}{}{I\mathrm{\Sigma }^{}}{|I^+||I^{}|qmod\frac{dim(𝔭)}{2}}}{}L(\mathrm{\Lambda }_{0,𝔨}+I+\rho _n+k_Ih^{}\nu (\varpi _i)+c_{I,k_I}^{}\delta _𝔨),$$
where $`k_I=\frac{2(q|I^+|+|I^{}|)}{dim(𝔭)}`$ and
$$c_{I,k_I}^{}=(k1)|I^+|k|I^{}|+(k^2k)\frac{dim(𝔭)}{4}.$$
## 6 Examples and applications
### 6.1 Combinatorial interpretation of decompositions in type $`C`$
We want to give a combinatorial interpretation of Theorem 3.5 and Theorem 4.10 for the pair $`𝔤=sp(V_1V_2)sp(V_1)sp(V_2)=𝔨`$, where $`V_1,V_2`$ are complex vector spaces of dimension $`2m,2n`$ respectively. It turns out that in this specific case (and indeed only in this) the decomposition formulas afford bijections between level $`m`$ representations of $`\widehat{sp(2n)}`$ and level $`n`$ representations of $`\widehat{sp(2m)}`$. This result, in the case of the spin representation, appears as Proposition 2 in . In our general setting we are considering the case of a Lie algebra $`𝔤`$ of type $`C_{n+m}`$ endowed with an involution $`\sigma `$ of type $`(0,\mathrm{}0,1,0\mathrm{}.0;1)`$, where $`1`$ appears in position $`m`$.
Let $`P_{n,m}`$ denote the set of $`(m+1)`$-weak compositions of $`n`$, i.e. ordered $`(m+1)`$-tuples $`(k_0,\mathrm{},k_m)`$ of non negative integers such that $`_{i=0}^mk_i=n`$. Let also $`S_{h,k}`$ denote the set of $`h`$ elements subsets of $`\{1,\mathrm{},k\}`$. The map $`(k_0,\mathrm{},k_m)\{k_0+1,k_0+k_1+2,\mathrm{},k_0+\mathrm{}+k_{m1}+m\}`$ is a bijection $`\zeta _{n,m}:P_{n,m}S_{m,m+n}`$. If $`c:S_{m,m+n}S_{n,m+n}`$ is the map which associates to an $`m`$-element subset of $`\{1,\mathrm{},m+n\}`$ its complement, the map $`\zeta _{m,n}^1c\zeta _{n,m}:P_{n,m}P_{m,n}`$ is a bijection, which we denote by $`(k_0,\mathrm{},k_m)(k_0^{},\mathrm{},k_n^{})`$. Set also $`k_i^{\prime \prime }=k_{ni}^{},\mathrm{\hspace{0.17em}0}in`$.
Let $`\dot{\mathrm{\Lambda }}_0,\mathrm{},\dot{\mathrm{\Lambda }}_m,\ddot{\mathrm{\Lambda }}_0,\mathrm{},\ddot{\mathrm{\Lambda }}_n`$ be the fundamental weights of the simple ideals of $`\widehat{𝔨}`$, assuming that both components have the Dynkin diagram displayed as in , §4, Table $`AffI`$.
###### Proposition 6.1.
Let $`𝔤,𝔨`$ be as above. The following decomposition formulas for the level $`1`$ modules of $`\widehat{so(𝔭)}`$ into irreducible $`[\widehat{𝔨},\widehat{𝔨}]`$-modules hold ($`ϵ=0,1`$):
$`L(\stackrel{~}{\mathrm{\Lambda }}_{lϵ})`$ $`=`$ $`{\displaystyle \underset{\genfrac{}{}{0pt}{}{(k_0,\mathrm{},k_m)P_{n,m}}{_{i=0}^mik_iϵmod\mathrm{\hspace{0.17em}2}}}{}}L(k_0\dot{\mathrm{\Lambda }}_0+\mathrm{}+k_m\dot{\mathrm{\Lambda }}_m)L(k_0^{}\ddot{\mathrm{\Lambda }}_0+\mathrm{}+k_n^{}\ddot{\mathrm{\Lambda }}_n),`$
$`L(\stackrel{~}{\mathrm{\Lambda }}_ϵ)`$ $`=`$ $`{\displaystyle \underset{\genfrac{}{}{0pt}{}{(k_0,\mathrm{},k_m)P_{n,m}}{_{i=0}^mik_iϵmod\mathrm{\hspace{0.17em}2}}}{}}L(k_0\dot{\mathrm{\Lambda }}_0+\mathrm{}+k_m\dot{\mathrm{\Lambda }}_m)L(k_0^{\prime \prime }\ddot{\mathrm{\Lambda }}_0+\mathrm{}+k_n^{\prime \prime }\ddot{\mathrm{\Lambda }}_n).`$
The key remark to deduce 6.1 from 3.5 and 4.10 is the following combinatorial interpretation of the sets $`N(w),wW_{\sigma ,\overline{r}}^{}`$. Consider the following rectangle $`R_{n,m}`$ filled with roots (of $`\widehat{L}(𝔤,\sigma )`$) as displayed in the following figure for $`m=2,n=3`$:
$`\alpha _2+\alpha _3+\alpha _4`$ $`\alpha _1+\alpha _2+\alpha _3+\alpha _4`$ $`\alpha _1+\alpha _2+\alpha _3`$ $`\alpha _1+\alpha _2`$ $`\alpha _2+\alpha _3`$ $`\alpha _2`$
Then the sets $`N(w),wW_{\sigma ,\overline{r}}^{}`$ can be described as the sets roots lying in the boxes under any lattice path from the South-West corner of the rectangle to the North-East corner. This is readily checked observing that these sets are biconvex (hence are of the form $`N(w)`$, for some $`w\widehat{W}`$), that they are either void or intersect $`\mathrm{\Pi }`$ exactly in $`\alpha _m`$ (hence are of the form $`N(w)`$, for some $`wW_{\sigma ,\overline{r}}^{}`$), and finally that they are as many as the above lattice paths, hence $`\left(\genfrac{}{}{0pt}{}{n+m}{n}\right)=|W_{\sigma ,\overline{r}}^{}|`$ in number (see , Table 5.1). Now the proposition follows by direct computation taking into account that $`\mathrm{\Lambda }_{0,𝔨}=n\dot{\mathrm{\Lambda }}_0+m\ddot{\mathrm{\Lambda }}_0,\rho _n+\mathrm{\Lambda }_{0,𝔨}=n\dot{\mathrm{\Lambda }}_m+m\ddot{\mathrm{\Lambda }}_0`$, $`\psi _{\overline{r}}^{}(\alpha _m)=\dot{\mathrm{\Lambda }}_1+\ddot{\mathrm{\Lambda }}_1`$ ($`r`$ even) and $`w_0=s_ms_{m1}s_m\mathrm{}s_1s_2\mathrm{}s_m`$, $`\psi _{\overline{r}}^{}(\alpha _m)=w_0(\dot{\mathrm{\Lambda }}_1)+\ddot{\mathrm{\Lambda }}_1`$ ($`r`$ odd) . More explicitely, it is not difficult to see that if $`p_w`$ is the lattice path associated to $`wW_{\sigma ,\overline{r}}^{}`$ and $`p_w(a_1,\mathrm{},a_m)(b_1,\mathrm{},b_n)`$, where $`0a_1a_2\mathrm{}n`$ (resp. $`mb_1b_2\mathrm{}0`$) are the lengths of the rows and (resp. columns) of the subdiagram of $`R_{n,m}`$ whose bottom border is $`p_w`$, counted from bottom to top (resp. from left to right), then
$$\mathrm{\Lambda }_{0,𝔨}\psi _{\overline{r}}^{}(N(w))=\underset{i=1}{\overset{n}{}}\dot{\mathrm{\Lambda }}_{mb_i}+\underset{i=1}{\overset{m}{}}\ddot{\mathrm{\Lambda }}_{na_i}$$
for $`r`$ even and
$$\rho _n+\mathrm{\Lambda }_{0,𝔨}\psi _{\overline{r}}^{}(N(w))=\underset{i=1}{\overset{n}{}}\dot{\mathrm{\Lambda }}_{b_i}+\underset{i=1}{\overset{m}{}}\ddot{\mathrm{\Lambda }}_{na_i}.$$
for $`r`$ odd.
### 6.2 A special case
Suppose that $`\sigma `$ is an automorphism of type $`(0,\mathrm{},1,\mathrm{},0;1)`$ with $`1`$ in a position corresponding to a long simple root (say $`\alpha _p`$). We show below how to calculate the $`\widehat{𝔨}`$-decomposition of the basic and vector representations in terms of a special class of representatives and how to get information on asymptotic dimension. We set for shortness $`\widehat{W}=\widehat{W}_{\sigma ,0},W^{}=W_{\sigma ,0}^{}`$. Let $`W_f`$ the Weyl group generated by $`s_1,\mathrm{},s_n`$. Let $`a(\mathrm{\Lambda })`$ denote the asymptotic dimension of a module $`L(\mathrm{\Lambda })`$ (see \[16, (2.1.5)\] for the definition).
###### Proposition 6.2.
1. The map $`ww_\sigma w`$ is an involution $`i`$ on $`W^{}`$.
Moreover we have that $`i(W^{}W_f)=W^{}(W^{}W_f).`$
2. Denote by $`\mathrm{\Lambda }_w=_{i=0}^nb_i\mathrm{\Lambda }_i`$ the weight of the factor indexed by $`wW^{}`$ in formula (3.8). If $`wW^{}W_f`$, then $`\mathrm{\Lambda }_{i(w)}=_{i=0}^nb_i\mathrm{\Lambda }_{\pi (i)}`$, where $`\pi `$ is a suitable permutation of $`\{1,\mathrm{},n\}`$. In particular $`a(\mathrm{\Lambda }_w)=a(\mathrm{\Lambda }_{i(w)})`$.
###### Proof.
Consider the set $`P_\sigma `$ defined in (3.10). By \[4, Lemma 5.9\], , we have $`w_\sigma (P_\sigma )P_\sigma `$, hence left multiplication by $`w_\sigma `$ gives a map $`i:W^{}W^{}`$. By \[4, Lemma 5.11\], we deduce that $`0`$ does not belong to $`w_\sigma \overline{C}_1`$; in particular $`w_\sigma W^{}W_f`$. This easily implies that $`i(W^{}W_f)W^{}(W^{}W_f)`$. It is clear that $`i`$ is injective. Proposition 5.8 and Theorem 5.12 of give $`|W^{}|=2|W^{}W_f|`$, hence $`i(W^{}W_f)=W^{}(W^{}W_f)`$. Finally $`i`$ is an involution since $`w_\sigma `$ is an involution. Indeed $`w_\sigma `$ is defined in as the product of certain elements of the extended Weyl groups of the irreducible components of the extended Dynkin diagram of $`𝔤`$ minus the $`p`$th node; in the action of these elements is completely worked out. This explicit description proves both that $`w_\sigma `$ is an involution and that it acts on each simple component $`\widehat{𝔨}_S`$ by permuting the fundamental weights. The assertion on asymptotic dimension follows from the fact that this quantity is invariant under the action of certain elements in the extended affine Weyl group. More precisely the invariance follows from \[16, (2.2.15-16)\] taking into account that $`w_\sigma `$ is a product of elements in $`W_0^+`$ (in the notation of ). ∎
### 6.3 More examples
The following examples should make clear how to use our decomposition formulas in explicit cases. To avoid cumbersome notation we describe the decomposition as $`[\widehat{𝔨},\widehat{𝔨}]`$ modules. In other words we consider the weights of the $`\widehat{𝔨}`$-modules appearing in the decompositions modulo $`\delta _𝔨`$.
1. We describe the decomposition of $`X_1`$ when $`𝔤`$ is of type $`G_2`$ and $`\sigma `$ of type $`(0,1,0;1)`$. In this case $`\widehat{𝔨}`$ is of type $`A_1^{(1)}\times A_1^{(1)}`$. $`W_{\sigma ,1}`$ is generated inside $`\widehat{W}`$ by $`s_0,s_2,s_1s_2s_1s_2s_1,s_1s_2s_1s_0s_1s_2s_1s_0s_1s_2s_1`$ and
$$W_{\sigma ,1}^{}=\{id,s_1,s_1s_0,s_1s_2,s_1s_2s_0,s_1s_2s_0s_1\}.$$
According to formula 4.10, the highest weights of the irreducible components are of the form $`2\dot{\mathrm{\Lambda }}_0+10\ddot{\mathrm{\Lambda }}_0+\rho _n\psi _1^{}(N(u))`$, where $`u`$ ranges over $`W_{\sigma ,1}^{}`$. Here and in the following $`\dot{\mathrm{\Lambda }}_i`$ denotes the $`i`$-th fundamental weight for the first copy of $`A_1^{(1)}`$ whereas $`\ddot{\mathrm{\Lambda }}_i`$ denotes the $`i`$-th fundamental weight for the other copy. Since $`\rho _n=w_0(2\alpha _1+3\alpha _2)`$ and $`\overline{\alpha }_1=\frac{1}{2}(\overline{\alpha }_0+3\overline{\alpha }_2)`$ we have
$`X_1=`$ $`L(2\dot{\mathrm{\Lambda }}_1)L(10\ddot{\mathrm{\Lambda }}_0)`$
$`L(\dot{\mathrm{\Lambda }}_0+\dot{\mathrm{\Lambda }}_1)L(7\ddot{\mathrm{\Lambda }}_0+3\ddot{\mathrm{\Lambda }}_1)`$
$`L(2\dot{\mathrm{\Lambda }}_1)L(4\ddot{\mathrm{\Lambda }}_0+6\ddot{\mathrm{\Lambda }}_1)`$
$`L(2\dot{\mathrm{\Lambda }}_0)L(6\ddot{\mathrm{\Lambda }}_0+4\ddot{\mathrm{\Lambda }}_1)`$
$`L(\dot{\mathrm{\Lambda }}_0+\dot{\mathrm{\Lambda }}_1)L(3\ddot{\mathrm{\Lambda }}_0+7\ddot{\mathrm{\Lambda }}_1)`$
$`L(2\dot{\mathrm{\Lambda }}_0)L(10\ddot{\mathrm{\Lambda }}_1).`$
2. We describe the decomposition of $`X_0`$ when $`𝔤`$ is of type $`D_4`$ and $`\sigma `$ of type $`(0,1,0,0;2)`$. In this case $`\widehat{𝔨}`$ is of type $`A_1^{(1)}\times C_2^{(1)}`$. $`W_{\sigma ,0}`$ is generated inside $`\widehat{W}`$ by $`s_0,s_2,s_3,s_1s_0s_1s_2s_1s_0s_1,`$ $`s_1s_2s_3s_2s_1s_0s_1s_2s_3s_2s_1`$ and we have
$`W_{\sigma ,0}^{}=\{Id,s_1,s_1s_0,s_1s_2,s_1s_0s_1,s_1s_0s_2,s_1s_2s_3,s_1s_0s_2s_3,s_1s_2s_3s_2,`$
$`s_1s_0s_2s_3s_2,s_1s_2s_3s_2s_1,s_1s_0s_2s_3s_2s_1\}.`$
According to formula 4.12, the highest weights of the irreducible components are of the form $`10\dot{\mathrm{\Lambda }}_0+3\ddot{\mathrm{\Lambda }}_0\psi _0^{}(N(u))`$, where $`u`$ ranges over $`W_{\sigma ,0}^{}`$. Taking into account that $`\overline{\alpha }_1=(\overline{\alpha }_0+\overline{\alpha }_2+\overline{\alpha }_3)`$,we get
$`X_0=`$ $`L(10\dot{\mathrm{\Lambda }}_0)L(3\ddot{\mathrm{\Lambda }}_0)`$
$`L(8\dot{\mathrm{\Lambda }}_0+2\dot{\mathrm{\Lambda }}_1)L(2\ddot{\mathrm{\Lambda }}_0+\ddot{\mathrm{\Lambda }}_2)`$
$`L(6\dot{\mathrm{\Lambda }}_0+4\dot{\mathrm{\Lambda }}_1)L(\ddot{\mathrm{\Lambda }}_0+2\ddot{\mathrm{\Lambda }}_1)`$
$`L(4\dot{\mathrm{\Lambda }}_0+6\dot{\mathrm{\Lambda }}_1)L(\ddot{\mathrm{\Lambda }}_0+2\ddot{\mathrm{\Lambda }}_1)`$
$`L(2\dot{\mathrm{\Lambda }}_0+8\dot{\mathrm{\Lambda }}_1)L(2\ddot{\mathrm{\Lambda }}_0+\ddot{\mathrm{\Lambda }}_2)`$
$`L(10\dot{\mathrm{\Lambda }}_1)L(3\ddot{\mathrm{\Lambda }}_0)`$
$`L(10\dot{\mathrm{\Lambda }}_0)L(3\ddot{\mathrm{\Lambda }}_2)`$
$`L(8\dot{\mathrm{\Lambda }}_0+2\dot{\mathrm{\Lambda }}_1)L(\ddot{\mathrm{\Lambda }}_0+2\ddot{\mathrm{\Lambda }}_2)`$
$`L(6\dot{\mathrm{\Lambda }}_0+4\dot{\mathrm{\Lambda }}_1)L(2\ddot{\mathrm{\Lambda }}_1+\ddot{\mathrm{\Lambda }}_2)`$
$`L(4\dot{\mathrm{\Lambda }}_0+6\dot{\mathrm{\Lambda }}_1)L(2\ddot{\mathrm{\Lambda }}_1+\ddot{\mathrm{\Lambda }}_2)`$
$`L(2\dot{\mathrm{\Lambda }}_0+8\dot{\mathrm{\Lambda }}_1)L(\ddot{\mathrm{\Lambda }}_0+2\ddot{\mathrm{\Lambda }}_2)`$
$`L(10\dot{\mathrm{\Lambda }}_1)L(3\ddot{\mathrm{\Lambda }}_2).`$
3. It is easy to see from our formulas that if $`𝔤`$ is of type $`D_{l+1}`$ and $`\sigma `$ is of type$`(1,0,\mathrm{},0;2)`$ then both the spin and the basic and vector representations restrict to the spin and basic and vector representations for $`B_l^{(1)}`$.
4. Finally we consider the decomposition of the spin representation $`X_1`$ for $`𝔤`$ of type $`D_4`$ and $`\sigma `$ of type $`(0,1,0,0;2)`$. As in example 2, $`\widehat{𝔨}`$ is of type $`A_1^{(1)}\times C_2^{(1)}`$. $`\widehat{W}_{𝔨_\mu }`$ is an affine Weyl group of type $`B_3`$. Recall that we chose $`\widehat{\mathrm{\Pi }}_\mu =\{\theta _f+\delta ^{},\overline{\alpha }_1,\overline{\alpha }_2,\overline{\alpha }_3\}`$ as a set of positive roots for $`\widehat{\mathrm{\Delta }}_\mu `$. Set $`\beta _0=\theta _f+\delta ^{},\beta _i=\overline{\alpha }_i,i=1,2,3,s_i=s_{\beta _i},i=0,1,2,3.`$ Then
$$\psi _{1}^{}{}_{}{}^{1}(\widehat{\mathrm{\Pi }}_𝔨)=\{\beta _2,\beta _3,\beta _0+\beta _2+\beta _3,\beta _0+\beta _1+\beta _2,\beta _1+\beta _2+\beta _3\},$$
hence $`\widehat{W}_{\sigma ,1}`$ is generated by $`s_2,s_3,s_0s_2s_3s_2s_0,s_0s_1s_2s_1s_0,s_1s_2s_3s_2s_1.`$ A set of minimal right coset representatives is
$$W_{\sigma ,1}^{}=\{Id,s_0,s_1,s_1s_0,s_1s_2,s_0s_2\}.$$
Taking into account that $`\overline{\alpha }_1=(\overline{\alpha }_0+\overline{\alpha }_2+\overline{\alpha }_3)`$, and that $`\rho _n=5\dot{\mathrm{\Lambda }}_0+\ddot{\mathrm{\Lambda }}_1`$, we get
$`X_1=`$ $`L(5\dot{\mathrm{\Lambda }}_0+5\dot{\mathrm{\Lambda }}_1)L(2\ddot{\mathrm{\Lambda }}_0+\ddot{\mathrm{\Lambda }}_1)`$
$`L(3\dot{\mathrm{\Lambda }}_0+7\dot{\mathrm{\Lambda }}_1)L(\ddot{\mathrm{\Lambda }}_0+\ddot{\mathrm{\Lambda }}_1+\ddot{\mathrm{\Lambda }}_2)`$
$`L(7\dot{\mathrm{\Lambda }}_0+3\dot{\mathrm{\Lambda }}_1)L(\ddot{\mathrm{\Lambda }}_0+\ddot{\mathrm{\Lambda }}_1+\ddot{\mathrm{\Lambda }}_2)`$
$`L(5\dot{\mathrm{\Lambda }}_0+5\dot{\mathrm{\Lambda }}_1)L(\ddot{\mathrm{\Lambda }}_1+2\ddot{\mathrm{\Lambda }}_2)`$
$`L(\dot{\mathrm{\Lambda }}_0+9\dot{\mathrm{\Lambda }}_1)L(3\ddot{\mathrm{\Lambda }}_1)`$
$`L(9\dot{\mathrm{\Lambda }}_0+\dot{\mathrm{\Lambda }}_1)L(3\ddot{\mathrm{\Lambda }}_1).`$
### 6.4 Connections with modular invariance.
We now try to use the formulas developed in the previous sections to obtain information on the action of $`SL(2,)`$ on modified characters described in . Here we consider the very special case when $`\sigma `$ comes from an automorphism of the diagram of $`𝔤`$. This implies that $`𝔤`$ is either simple of type $`A,D,E`$ or of complex type. Furthermore $`𝔨`$ is simple. We shall also assume that $`𝔤`$ is not of type $`A_{2n}`$. These are precisely the cases in which $`W_{\sigma ,1}^{}=\{1\}`$.
Let $`h_𝔨^{}`$ denote the dual Coxeter number of $`𝔨`$ and set $`j=h^{}h_𝔨^{}`$. We denote by $`\dot{\mathrm{\Lambda }}_i`$ the $`i`$-th fundamental weight of $`\widehat{𝔨}`$ and by $`P_+^j`$ the set of dominant weights for $`\widehat{𝔨}`$ of level $`j`$. Recall that $`N=rk𝔤`$ while $`n=rk𝔨`$. By (3.7) we have that
$$ch(L(\stackrel{~}{\mathrm{\Lambda }}_0))ch(L(\stackrel{~}{\mathrm{\Lambda }}_1))=\underset{wW_{\sigma ,0}^{}}{}ϵ(w)ch(L(\psi _0^{}(w\widehat{\rho })\widehat{\rho }_𝔨)).$$
(6.1)
Formula (4.17) becomes in our case
$$ch(L(\stackrel{~}{\mathrm{\Lambda }}_m))=2^{\frac{Nn}{2}}ch(L(j\dot{\mathrm{\Lambda }}_0+\rho _n))$$
(6.2)
if $`Nn`$ is odd, and
$$ch(L(\stackrel{~}{\mathrm{\Lambda }}_{m1}))+ch(L(\stackrel{~}{\mathrm{\Lambda }}_m))=2^{\frac{Nn}{2}}ch(L(j\dot{\mathrm{\Lambda }}_0+\rho _n)),$$
(6.3)
if $`Nn`$ is even. Here $`m=\frac{dim(𝔭)}{2}`$.
Denote by $`\chi _\mathrm{\Lambda }`$ is the modified character of $`L(\mathrm{\Lambda })`$ (see \[16, (1.5.11)\]), and set $`Y=\{h\widehat{𝔥}^{}Re\delta _𝔨(h)>0\}`$. Moreover we write $`\mathrm{\Lambda }_w`$ for $`\psi _0^{}(w\widehat{\rho })\widehat{\rho }_𝔨`$. Since the pair $`(so(𝔭),𝔨)`$ is conformal, relation (6.1) translates into
$$(\chi _{\stackrel{~}{\mathrm{\Lambda }}_0}\chi _{\stackrel{~}{\mathrm{\Lambda }}_1})_{|Y}=\underset{wW_{\sigma ,0}^{}}{}ϵ(w)\chi _{\mathrm{\Lambda }_w},$$
(6.4)
whereas (6.2) gives
$$(\chi _{\stackrel{~}{\mathrm{\Lambda }}_m})_{|Y}=2^{\frac{Nn}{2}}\chi _{j\dot{\mathrm{\Lambda }}_0+\rho _n}$$
(6.5)
($`Nn`$ odd), and (6.3) gives
$$(\chi _{\stackrel{~}{\mathrm{\Lambda }}_{m1}}+\chi _{\stackrel{~}{\mathrm{\Lambda }}_m})_{|Y}=2^{\frac{Nn}{2}}\chi _{j\dot{\mathrm{\Lambda }}_0+\rho _n}$$
(6.6)
($`Nn`$ even). Recall from \[16, Remark 4.2.2\] that if $`Nn`$ is odd,
$$\chi _{\stackrel{~}{\mathrm{\Lambda }}_m}(\frac{1}{\tau })=\frac{1}{\sqrt{2}}(\chi _{\stackrel{~}{\mathrm{\Lambda }}_0}\chi _{\stackrel{~}{\mathrm{\Lambda }}_1})(\tau )$$
(6.7)
and, if $`Nn`$ is even,
$$(\chi _{\stackrel{~}{\mathrm{\Lambda }}_{m1}}+\chi _{\stackrel{~}{\mathrm{\Lambda }}_m})(\frac{1}{\tau })=(\chi _{\stackrel{~}{\mathrm{\Lambda }}_0}\chi _{\stackrel{~}{\mathrm{\Lambda }}_1})(\tau ).$$
(6.8)
By modular invariance of modified characters,
$$\chi _{j\dot{\mathrm{\Lambda }}_0+\rho _n}(\frac{1}{\tau })=\underset{\mathrm{\Lambda }P_+^j}{}a(\mathrm{\Lambda },j\dot{\mathrm{\Lambda }}_0+\rho _n)\chi _\mathrm{\Lambda }.$$
(6.9)
(here $`a(,)`$ is the function $`P_+^j\times P_+^j`$ defined in \[16, (2.1.7)\]). Assume $`Nn`$ even and use (6.4),(6.8),(6.6),(6.9) obtaining
$`{\displaystyle \underset{wW_{\sigma ,0}^{}}{}}ϵ(w)\chi _{\mathrm{\Lambda }_w}(\tau )`$ $`=(\chi _{\stackrel{~}{\mathrm{\Lambda }}_0}\chi _{\stackrel{~}{\mathrm{\Lambda }}_1})(\tau )=(\chi _{\stackrel{~}{\mathrm{\Lambda }}_{m1}}+\chi _{\stackrel{~}{\mathrm{\Lambda }}_m})({\displaystyle \frac{1}{\tau }})`$
$`=2^{\frac{Nn}{2}}\chi _{j\dot{\mathrm{\Lambda }}_0+\rho _n}({\displaystyle \frac{1}{\tau }})=2^{\frac{Nn}{2}}{\displaystyle \underset{\mathrm{\Lambda }P_+^j}{}}a(\mathrm{\Lambda },j\dot{\mathrm{\Lambda }}_0+\rho _n)\chi _\mathrm{\Lambda }(\tau ).`$
The case $`Nn`$ odd is analogous. We can deduce the following
###### Proposition 6.3.
We have $`a(\mathrm{\Lambda },j\dot{\mathrm{\Lambda }}_0+\rho _n)=0`$ unless there exists $`wW_{\sigma ,0}^{}`$ such that $`\mathrm{\Lambda }+\widehat{\rho }_𝔨=\psi _0^{}(w\widehat{\rho })`$. In such a case $`a(\mathrm{\Lambda },j\dot{\mathrm{\Lambda }}_0+\rho _n)=2^{\frac{Nn}{2}}(1)^{\mathrm{}(w)}`$.
###### Remark 6.1.
In the complex case this result was obtained in the same way in , (4.2.14).
###### Remark 6.2.
Recall that, if $`\mathrm{\Sigma }`$ is the set of $`𝔟_0`$-stable abelian subspaces of $`𝔭`$, then, according to Theorem 3.9, the set $`\mathrm{\Sigma }`$ parametrizes the irreducible components of $`X_0`$. By \[16, (2.2.3)\], we know that
$$\underset{\mathrm{\Lambda }P_+^j}{}|a(\mathrm{\Lambda },j\dot{\mathrm{\Lambda }}_0+\rho _n)|^2=1.$$
We can therefore deduce that $`|\mathrm{\Sigma }|=2^{Nn}`$ in these cases. This fact was first proved in by a different method. In the complex case we have yet another proof of Peterson’s $`2^{\mathrm{rank}}`$ abelian ideals Theorem (see again ).
###### Remark 6.3.
If $`𝔤`$ is of type $`D_N`$, then $`𝔨`$ is of type $`B_{N1}`$ and one only obtains again that $`a(\dot{\mathrm{\Lambda }}_{N1},\dot{\mathrm{\Lambda }}_0)=a(\dot{\mathrm{\Lambda }}_{N1},\dot{\mathrm{\Lambda }}_1)=\frac{1}{\sqrt{2}}`$ and $`a(\dot{\mathrm{\Lambda }}_{N1},\dot{\mathrm{\Lambda }}_{N1})=0`$.
P.C.: Dipartimento di Scienze, Università di Chieti-Pescara, Viale Pindaro 42, 65127 Pescara, ITALY;
cellini@sci.unich.it
P.MF.: Politecnico di Milano, Polo regionale di Como, Via Valleggio 11, 22100 Como, ITALY;
frajria@mate.polimi.it
V.K.: Department of Mathematics, Rm 2-165, MIT, 77 Mass. Ave, Cambridge, MA 02139;
kac@math.mit.edu
P.P.: Dipartimento di Matematica, Università di Roma “La Sapienza”, P.le A. Moro 2, 00185, Roma , ITALY;
papi@mat.uniroma1.it |
warning/0506/astro-ph0506284.html | ar5iv | text | # COSMOGRAPHY, DECELERATING PAST, AND COSMOLOGICAL MODELS: LEARNING THE BAYESIAN WAY
## 1 INTRODUCTION
It is usual in cosmology that when new data come in, we need to readjust our parameter values or shall even attempt to modify our theoretical model itself. But each time while doing so, we are guided by our knowledge gathered so far about the universe. This kind of gradual learning is characteristic of sciences like cosmology, where one cannot conduct laboratory experiments. A very interesting turn of events which unfolded in cosmology in the past ten years exemplifies this learning process. Till the early 1990s, a Friedmann cosmological model with either radiation or dust as energy component was conceived to be the standard model in cosmology and it was generally believed that this standard model described the observed universe (at least from the time of nucleosynthesis) to a very good accuracy. Most of the widely discussed cosmological problems, which ought to have surfaced in the classical epoch, were known to disappear when we incorporate into the standard model the theory of inflation, which envisages a flat universe. The value of the Hubble parameter $`H_0100h`$ km s<sup>-1</sup> Mpc<sup>-1</sup> quoted during that period was $`0.4<h<1`$ (Kolb & Turner 1990). There was no serious age problem and no real need for a cosmological constant $`\mathrm{\Lambda }`$ in the late universe. The beginning of the present phase of $`\mathrm{\Lambda }`$-term or dark energy (Padmanabhan 2003; Peebles & Ratra 2003) in cosmology is the measurement of the Hubble parameter reported in Freedman et al. (1994); Pierce et al. (1994) and summarized in Jaffe (1996), that gave a more specific and high range $`h>0.7`$. This, along with flatness imposed by the theory of inflation created a short spell of age problem, which required a $`\mathrm{\Lambda }`$-term for its solution. But a nonzero $`\mathrm{\Lambda }`$ was always behind the curtains and hence its appearance was not conceived to be a major deviation in the standard model scenario. But then came the Type Ia supernovae (SNe Ia) data (Perlmutter et al. 1999; Riess et al. 1998), which required the presence of a $`\mathrm{\Lambda }`$-term large enough to cause an accelerated expansion. The most recent release of SNe data (Tonry et al. 2003; Knop et al. 2003; Riess et al. 2004) prompts many cosmologists even to speculate that some extremely unphysical energy densities, such as those with $`w<1`$ in the equation of state $`p=w\rho `$, are required to explain the data. Though this change was gradual, what happened here is a U-turn from the standpoint of the early nineties, since the dynamics, energy content, equation of state, etc. of the universe have now become totally speculative. But this is typical of every learning process and cannot be termed unscientific.
It has long been recognized that the application of Bayes’s theorem in physical problems represents learning. \[For a recent review, see D’Agostini (2003).\] This theorem tells us how to adjust our plausibility assumptions regarding a hypothesis when our state of knowledge changes through the acquisition of new data. Classical mathematicians such as Bernoulli, Bayes, Laplace and Gauss have found Bayes theorem useful in problems such as those in astronomy, thanks to its ability to learn. Later, because of difficulties with assigning prior probabilities (which were mistakenly considered to be purely subjective expressions of a person’s opinions about hypotheses), Bayes’s probability theory has gone out of favor in physical sciences and was replaced by the more apparently objective ‘frequentist’ approach. But the realization that the frequentist definition of probability is as subjective as the Bayesian has called forth a re-examination of this controversy. If some simple desideratum such as ‘equivalent states of knowledge should be represented by equivalent probability assignments’ (which is termed as Jaynes’s consistency) is followed, the Bayesian approach will help to quantify the collective wisdom of scientists and hence can be made less subjective. When compared to the improper application of frequentist probability theory, the Bayesian approach is most powerful in problems such as those in cosmology, where the process of learning and also the quantification and readjustment of plausibility assessments by the scientific community are very important.
The apparent magnitude-redshift ($`mz`$) data of SNe Ia are the only qualitative signature of an accelerated expansion and hence it is very important in understanding the dynamics of the universe. This accelerated expansion is confirmed (John 2004), provided there are no evolutionary effects for SNe (Drell et al. 2000). For analyzing the data, we need an expression for the scale factor $`a(t)`$. If the attempt is to obtain $`a(t)`$ as the solution of Einstein equation, one must know the energy densities present in the universe and also their equations of state. But here the data indicate that known forms of energies are unable to account for it. Under such circumstances, the use of the traditional Friedmann solutions of the scale factor $`a(t)`$ obtained by assuming the presence of known energies in cosmological problems, and particularly in the analysis of SNe data, can be very misleading. Hence these SNe data are to be analyzed cosmographically, without making any specific assumptions on the energy densities in the universe. Such an analysis of SNe data in John (2004) assumes only homogeneity and isotropy for the universe; i.e., the universe is assumed to have a Robertson-Walker (RW) metric. The scale factor of the universe can most naturally be expanded into a Taylor series in $`t`$ about the present epoch and we attempt to find its coefficients from observation, as it was done to evaluate the Hubble parameter and deceleration parameter in the original cosmographic approach (Weinberg 1972). Terms up to fifth-order were kept in the above series and it was assumed that they make a good approximation. The results of this and other analyses of SNe data do not give reason to believe that there were only standard model energy densities or even any other energies with equation of state of the form $`p=w\rho `$ in the present universe. Instead, the likelihoods for the various expansion rates obtained in our calculation are very broadly peaked and these indicate that there are a variety of choices for the energy densities in the present epoch. The Taylor series expansion approach is extended to third-order and fourth-order in Sahni et al. (2003) and Visser (2004), respectively. Also Daly & Djorgovski (2003); Wang & Mukherjee (2003); Daly & Djorgovski (2004) have attempted similar model-independent analyses of SNe data.
When ensembles and repeated experiments are not possible, a natural and useful procedure in cosmology is to compare how best different models can account for the data, using the Bayesian method. In Bayesian model comparison, one computes odds ratios between different models. Jaffe (1996) and Hobson et al. (2002) have used Bayesian theory to test the relative merits of different cosmological models. In one such application of this technique, John & Narlikar (2002) have compared the standard and inflationary models having nonzero $`\mathrm{\Lambda }`$ with a new simple model having the scale factor $`a(t)t`$. This comparison was made using the SNe data set in Perlmutter et al. (1999) and by assuming flat priors. Flat priors indicate that we are having no prior information regarding parameters of the models except that they lie in some fiducial range. In the present case, this is equivalent to stating that one depends only on the present SNe data for making the comparison. In John & Narlikar (2002), it was found that the then available apparent magnitude-redshift data alone were not very much discriminatory between these different models. However, to be true to the spirit of Bayesian theory, our plausibility assignments should be updated with the acquisition of new data. In this paper, we attempt to do this, by using the recent release of SNe data.
Another important question we try to answer in this paper is whether the data really endorse that the universe was decelerating in the past. For this, a model with some decelerating phase in the past is compared with another model having no such phase, by using the Bayesian approach. In both cases, we assume the present universe to be accelerating. It is claimed that this method of comparison is more robust than other investigations seeking evidence for a decelerating past since both the Bayesian and cosmographic approaches mentioned above are used here. As mentioned earlier, the scale factor is expanded into a fifth-order polynomial in time (fifth-order is required for sufficient accuracy) and then the combinations of various coefficients in this expansion were separated into those which correspond only to acceleration during the entire period and those which had at least some decelerating phase in this period. Considering these as rival models, the Bayes factor is calculated. As an example of the Bayesian model comparison technique, we compare the general relativistic and Newtonian explanations of deflection of light by sun, to show how powerful is the available observed data on deflection of light in discriminating these explanations. We argue that, in contrast to this case, the SNe data are not capable of discriminating the “always accelerating” and “decelerating in the past” cosmological models. The results have serious implications for modern theoretical cosmology for it is almost entirely built on the firm belief that the universe was decelerating in the past.
The paper is organized as follows. In section 2, the general formalism of Bayesian model comparison employed here is given. The question of whether there was a decelerating past for the universe is discussed in section 3. Comparison of some Friedmann cosmological models is discussed in section 4 and section 5 summarizes our conclusions.
## 2 BAYESIAN MODEL COMPARISON
The Bayes’s theorem helps to evaluate the posterior (i.e., after analyzing the data) probability $`p(H_i|D,I)`$ for a hypothesis $`H_i`$ given the data $`D`$ and the truth of some background information $`I`$, as
$$p(H_i|D,I)=\frac{p(H_i|I)p(D|H_i,I)}{p(D|I)}.$$
(1)
$`p(H_i|I)`$ is called the prior (i.e., before analyzing the data) probability and is the probability for $`H_i`$, given the truth of $`I`$ alone. $`p(D|H_i,I)`$ is the probability for obtaining the data $`D`$ if the hypothesis $`H_i`$ and $`I`$ were true and is called the likelihood for the hypothesis. The factor in the denominator serves the purpose of normalization.
In Bayesian model comparison, one finds the odds ratios; i.e., the ratios between the posterior probabilities for different models. If we have to compare rival models $`M_i`$ and $`M_j`$, take the truths of these as the hypotheses $`H_i`$ and $`H_j`$, respectively and write, using Bayes’s theorem,
$$\frac{p(M_i|D,I)}{p(M_j|D,I)}=\frac{p(M_i|I)p(D|M_i,I)}{p(M_j|I)p(D|M_j,I)}O_{ij}.$$
(2)
To some extent, evaluating the prior probability is subjective, since it depends only on the prior information $`I`$ and this may vary from person to person. But Bayesians view this theory as an attempt to quantify the collective wisdom of researchers working in the field and hence finding fiducial prior probabilities would be highly rewarding. However, if the information $`I`$ does not prefer one model over the other, the prior probabilities get cancel out and the odds ratio is simply
$$O_{ij}=\frac{p(D|M_i,I)}{p(D|M_j,I)}B_{ij}.$$
(3)
As mentioned above, the probability $`p(D|M_i,I)`$ for the data $`D`$, given that the model $`M_i`$ and $`I`$ are true, is called the likelihood for the model $`M_i`$ and is denoted as $`(M_i)`$. For parameterized models, with parameters $`\alpha ,\beta ,..`$, this quantity can be evaluated as
$$p(D|M_i,I)(M_i)=𝑑\alpha 𝑑\beta \mathrm{}p(\alpha ,\beta ,\mathrm{}|M_i)_i(\alpha ,\beta ,\mathrm{}),$$
(4)
where $`p(\alpha ,\beta ,\mathrm{}|M_i)`$ is the prior probability for the set of parameter values $`\alpha ,\beta ,..`$ and $`_i(\alpha ,\beta ,\mathrm{})`$ is the likelihood for the combination. The latter quantity is often taken to be
$$_i(\alpha ,\beta ,\mathrm{})=\mathrm{exp}[\chi _i^2(\alpha ,\beta ,..)/2],$$
(5)
where
$$\chi ^2=\mathrm{\Sigma }_k\left(\frac{\widehat{A}_kA_k(\alpha ,\beta ,..)}{\sigma _k}\right)^2$$
(6)
is the $`\chi ^2`$-statistic. Here $`\widehat{A}_k`$ are the measured values of the observable $`A`$, $`A_k(\alpha ,\beta ,..)`$ are its expected values (from theory) and $`\sigma _k`$ are the uncertainties in the measurement of the observable.
In a certain volume $`V`$ of the parameter space, one can assign flat prior probability for all the parameters by taking $`p(\alpha ,\beta ,\mathrm{}|M_i)=`$ constant throughout this volume. When normalized, this prior is simply $`1/V`$. In those special cases where there are no adjustable parameters in the model, we have $`\delta `$-function prior and from equations (4) and (5),
$$(M_i)=\mathrm{exp}(\chi _i^2/2).$$
(7)
$`B_{ij}`$ in equation (3) is referred to as the Bayes factor. The interpretation of this quantity is as follows (Drell et al. 2000): If $`1<B_{ij}<3`$, there is evidence against model $`M_j`$, but it is not worth more than a bare mention. If $`3<B_{ij}<20`$, this evidence is positive. If $`20<B_{ij}<150`$, it is strong and if $`B_{ij}>150`$, the evidence is very strong.
To appreciate the significance of this interpretation of the Bayes factor, let us compare the general relativistic and Newtonian explanations of the deflection of a light ray that just grazes the sun’s surface. The theoretical prediction made by general relativity ($`M_1`$) in this case is $`\theta _0=1.75^{\prime \prime }`$ whereas in the purely Newtonian case ($`M_2`$), it is $`\theta _0=0.875^{\prime \prime }`$. If we use the data $`\widehat{\theta }_0=1.98\pm 0.16^{\prime \prime }`$ obtained from the classic 1919 eclipse expedition to the island of Sobral (Dyson et al. 1920a, b), the Bayes factor calculated using equations (3)-(7) above is $`B_{12}10^{10}`$. It may be noted that the $`\chi ^2`$-values and hence the Bayes factor crucially depend on the error bars. Though the error bars in the above case are now felt to be grossly underestimated, this published data took the world by storm in favor of GR in 1919 and the reason can be understood from the huge value of $`B_{12}`$ we obtained above. If all the 9 observational data points provided with error bars in Table 8.1 of Weinberg (1972) were used, one gets $`B_{12}10^{82}`$, a spectacularly large value for the Bayes factor to settle the issue in favor of general relativity. These kinds of results are common in tests of quantum theory too. These examples show that the above requirement $`B_{ij}>150`$ for an evidence to be very strong is really modest.
## 3 ACCELERATING VS. DECELERATING PAST
We compare models of an always accelerating universe from time $`t_0+T_p`$ to the present ($`t_0`$ is the present time and $`T_p`$ is negative) with those which were decelerating for at least some time in this period. Computations were performed for various values of $`T_p`$ ranging from $`1\times 10^{17}`$ s to $`5\times 10^{17}`$ s, the absolute value of the latter ($`15`$ Gyr) being usually quoted as the upper limit for the age of the universe.
In John (2004), it was assumed that the scale factor of the universe can be approximated by a fifth-order polynomial in time. With the present value of the scale factor as $`a_0`$, the present deceleration parameter as $`q_0`$, and the present values of other parameters related to higher order derivatives as $`r_0`$, $`s_0`$ and $`u_0`$, this expression for $`a(t)`$ is
$`a(t_0+T)`$ $`=`$ $`a_0\left[1+H_0T{\displaystyle \frac{q_0H_0^2}{2!}}T^2+{\displaystyle \frac{r_0H_0^3}{3!}}T^3{\displaystyle \frac{s_0H_0^4}{4!}}T^4+{\displaystyle \frac{u_0H_0^5}{5!}}T^5\right]`$ (8)
$``$ $`a_0\left[1+a_{(1)}T+a_{(2)}T^2+a_{(3)}T^3+a_{(4)}T^4+a_{(5)}T^5\right].`$
In addition to the above parameters, we have $`k=0,\pm 1`$, which is the curvature constant appearing in the RW metric and $`M`$, the absolute luminosity of SNe. The parameters $`a_0`$ and $`M`$ have reasonable flat priors for $`a_0>3000`$ Mpc and $`19.6<M<19.1`$ magnitudes, respectively. The upper bound for $`a_0`$ is chosen as 8000 Mpc, large enough to incorporate spatially flat models. Consequently, $`k=0`$ need not be included in the calculations. In John (2004), marginal likelihoods were evaluated for other parameters and found that the contributing ranges are $`0.6<h<0.8`$, $`2<q_0<1`$, $`15<r_0<15`$, $`65<s_0<65`$, and $`150<u_0<150`$. However, some combination of these parameter values were found to give no solution to the equation $`1+z=a(t_0)/a(t_0+T)`$ even for a time as past as $`T=10\times 10^{17}`$ and those values were excluded.
In the Bayesian model comparison to find evidence for a decelerating past for the universe using new data, one should accept the above ranges as the prior information obtained from the previous analysis. However, a modification is suggested to the effect that we shall begin by accepting the present universe to be accelerating; i.e., $`q_0<0`$. Considering the various other cosmological observations and the general perception among cosmologists, it is only reasonable to include this into the information $`I`$ we have, before analyzing the present data. The two hypotheses we want to compare may now be explicitly stated: (1) The universe is always accelerating from time $`t_0+T_p`$ to the present epoch (model $`M_1`$) and (2) There is at least one decelerating epoch for the universe during this period (model $`M_2`$). The factors $`p(\alpha ,\beta ,..|M_i)`$ in equation (4), which are the prior probabilities for the parameters (given the truth of the respective models), are taken in this case to be the flat probabilities $`1/V_1`$ and $`1/V_2`$ for models $`M_1`$ and $`M_2`$, respectively, where $`V_1`$ and $`V_2`$ are the volumes in the parameter space corresponding to each of them. For any particular combination of parameter values, a sure test for the occurrence of deceleration during $`T_p<T<0`$ is to plot
$$\ddot{a}(t_0+T)=a_0[2a_{(2)}+6a_{(3)}T+12a_{(4)}T^2+20a_{(5)}T^3]$$
(9)
for this interval and to see whether it becomes negative at any time during the period. The Bayes factor is then
$$B_{12}=\frac{(1/V_1)\mathrm{\Sigma }_k_{V_1}\mathrm{exp}(\chi _1^2/2)𝑑a_0𝑑h𝑑q_0𝑑r_0𝑑s_0𝑑u_0𝑑M}{(1/V_2)\mathrm{\Sigma }_k_{V_2}\mathrm{exp}(\chi _2^2/2)𝑑a_0𝑑h𝑑q_0𝑑r_0𝑑s_0𝑑u_0𝑑M}.$$
(10)
Here, $`\chi ^2`$ for a model is calculated using equation (6), with the replacement of $`A`$ by $`m`$, the apparent magnitude of the SN and the parameters $`\alpha ,\beta ,..`$ are $`a_0`$, $`h`$, $`q_0`$, $`r_0`$, $`s_0`$, $`u_0`$, $`M`$, and $`k`$. The expression for $`m`$ to be used is
$$m=5\mathrm{log}\frac{D}{1\text{Mpc}}+25+M.$$
(11)
$`D/1`$Mpc refers to the luminosity distance $`D=r_1a_0(1+z)`$, expressed in megaparsecs. The comoving coordinate $`r_1`$ can be found from
$$_{T_1}^0\frac{cdT}{a(t_0+T)}=_0^{r_1}\frac{dr}{1kr^2}S_k^1(r_1),$$
(12)
where $`t_1=t_0+T_1`$ is the time at which an SN at $`r_1`$ emits the light and $`S_k^1(r_1)`$ is equal to $`\mathrm{sin}^1(r_1)`$ for $`k=+1`$, and $`\mathrm{sinh}^1(r_1)`$ for $`k=1`$.
An important part of the calculation is the solution of the following equation, used to find $`T_1`$ in terms of $`z`$, for each combination of parameter values. This is done in a direct and purely numerical way \[and differently from the way it was done in John (2004)\]:
$$1+z=\frac{a(t_0)}{a(t_0+T_1)}=\frac{1}{1+a_{(1)}T_1+a_{(2)}T_1^2+a_{(3)}T_1^3+a_{(4)}T_1^4+a_{(5)}T_1^5}.$$
(13)
Another improvement is that this numerical solution for $`T_1`$ is found with the more reliable regula falsi method, rather than the Newton-Raphson method employed in the above paper. As a consequence of these modifications, $`r_1`$ has to be obtained by numerical integration in equation (12). Though it now requires more computation time, these changes have made the analysis free of the truncations that weakened the approximations for certain parameters in the previous case. No point in the parameter space is now left out for the reason of breaking down of the approximations. As mentioned earlier, the only points left out are those which do not have a solution for equation (13) for all $`z`$ in the data set, even for the past $`10^{18}`$s. They are not included in the volumes $`V_1`$ and $`V_2`$ either. The values of the Bayes factor obtained for various values of $`T_p`$ are tabulated in Table 1. The values obtained while using the 54 “All SCP” SNe (data $`D_1`$) in Knop et al. (2003) (as reproduced in John (2004)) and the 157 “gold” data points (data $`D_2`$) in Riess et al. (2004), respectively are given in two separate columns. It may also be noted that while using $`D_2`$, one does not have to include the parameter $`M`$ in the calculations since the data give $`mM`$, the distance modulus in place of $`m`$.
The envelopes of the marginal likelihoods for parameters obtained in John (2004) using $`D_1`$ were found to be mostly unchanged. However, it can be seen that the afore-said modifications have made significant improvement in the quality of the curves and for the purpose of verification, the new likelihoods for $`h`$, $`q_0`$, $`r_0`$, $`s_0`$, and $`u_0`$ are given here in Figures 1-5. The fluctuations have now disappeared and we have very smooth curves. These likelihoods were computed for ranges wider than the ones in the previous case, but the prior ranges for these parameters remained the same (information $`I`$) in the computations of Bayes factors and other marginal likelihoods.
The results show that except in the case of using data $`D_2`$ for $`T_p=2\times 10^{17}`$s, there is hardly any evidence for a decelerating phase in the past 15 Gyrs. All the other Bayes factors in the table are slightly greater than unity and this shows that if at all there is evidence, it is in favor of an always accelerating universe during this period. However, in sharp contrast to the model comparison exercise using data of light deflection by sun (discussed in sec. 2), here the evidences are too weak and hence it is safer to conclude that the data are unable to discriminate these two hypotheses.
## 4 COMPARISON OF SOME FRIEDMANN MODELS
In this section, we use the method of model comparison to some Friedmann cosmological models. The popular models compared are (1) the Friedmann-Lamaitre-Robertson-Walker (FLRW) model (Model $`M_1`$) having a Robertson-Walker (RW) metric with $`k=\pm 1,0`$, which contains ordinary matter and a “constant” cosmological constant and (2) the same FLRW model but with inflation, which implies a $`k=0`$ RW metric (Model $`M_2`$). A comparison between these two models itself is an interesting problem, even though the SNe data do not contain direct imprints of inflation, except through its prediction of flat space sections or the equivalent condition $`\mathrm{\Omega }_m+\mathrm{\Omega }_\mathrm{\Lambda }=1`$.
Another model (Model $`M_3`$) we compare with the above is a Friedmann model having the simplest evolution with time; i.e., the one whose scale factor $`a(t)t`$ throughout the history of the universe. This kind of a coasting evolution can be arrived at in a number of ways. The quantity $`\rho +3p`$, where we denote the total energy density as $`\rho `$ and the total pressure of the universe as $`p`$, is sometimes called the gravitational charge. In all Friedmann models, if the gravitational charge vanishes, the evolution is coasting. Another approach is to generalize a dimensional argument by Chen and Wu to deduce that the energy densities in the universe, which are not separately conserved, vary as $`a^2(t)`$ (John & Babu Joseph 2000). In this case too, a coasting evolution follows. In the above paper, it was noted that when we consider the universe as containing more than one component, eg., matter/radiation and $`\mathrm{\Lambda }`$, the coasting model has a unique feature in which it predicts a ratio between the density parameters; $`\mathrm{\Omega }_m/\mathrm{\Omega }_\mathrm{\Lambda }=2/(1+3w)`$ ($`w=0`$ or 1/3). Hence it has no coincidence problem. It was shown that this model is devoid of other problems like the horizon, flatness, monopole, size, age of the universe and the generation of density perturbations on scales well above the present Hubble radius in the classical epoch. An added advantage is that here one can consider the generation of the observed density perturbations as a late-time classical behavior too. The solution of the coincidence, age and density perturbation problems deserve special mention since these problems are not solvable in an inflationary scenario. It was also shown that the evolution of temperature in the model is nearly the same as that in the standard big bang model and this will enable nucleosynthesis to proceed in an identical manner, provided the total density parameter $`\mathrm{\Omega }2`$. In John & Narlikar (2002), the afore-said condition was stated to be a problem with model $`M_3`$.
The above coasting evolution can also be obtained by a quite different route (John & Babu Joseph 1996, 1997). Here we consider the analytic continuation of a real analytic manifold (the spacetime) into the complex, producing a complex spacetime. There are a number of instances of the use of complex numbers or complex analytic functions in general theory of relativity (Flaherty 1980; Newman 1988). This model is arrived at by extending the idea of a possible ‘signature change’ in the universe (Ellis et al. 1992; Mars et al. 2001), a widely discussed speculation which involves some basic issues in the general relativity. This extension leaves us in an unphysical universe, but it was noticed that a proper interpretation of the theory will enable us to obtain a cosmological model, with the essential features as summarized above. The evolution of $`a(t)`$ is obtained as $`a^2(t)=a_{min}^2+c^2t^2`$ and a quantum cosmological calculation shows that $`a_{min}l_p`$, the Planck length. Thus after the Planck epoch, $`at`$ and the evolution is the same as in the above case. In addition to the above advantages, this approach solves the singularity problem too.
However, the analysis in this paper is intended to provide and improve an example of Bayesian model comparison (John & Narlikar 2002) and the coasting model is suitable for this purpose due to its simplicity. This model is not a realistic one since it requires some more fundamental derivation such as that in the scalar field dark energy models. As in the previous work, we assume that the model is given and then proceed with the evaluation.
One observes that when the Bayes factor is near unity, the prior odds $`p(M_i|I)/p(M_j|I)`$ in equation (2) becomes very important. While comparing models $`M_1`$, $`M_2`$, and $`M_3`$, it was observed in the above paper that whereas models $`M_1`$ and $`M_2`$ are plagued by the large number of cosmological problems, $`M_3`$ suffers from the heuristic nature of its derivation and the problem with nucleosynthesis, as mentioned earlier. But one can argue that the inclusion of a small $`\mathrm{\Lambda }`$ in the present universe in models $`M_1`$ and $`M_2`$ itself is done heuristically. However, a balance was set between the prior probabilities and these models were compared with prior odds equal to unity. The results thus obtained in the previous work, while using the $`mz`$ data (Perlmutter et al. 1999) for comparison were $`B_{13}3.1`$ and $`B_{23}5`$ and the conclusion made was that there is some but not strong evidence against model $`M_3`$.
The data used in the present analysis are the same as those used in the calculations in the previous section. The parameters in the models are $`H_0`$, $`M`$, $`\mathrm{\Omega }_m`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$. But in $`M_2`$ and $`M_3`$, their numbers are effectively reduced by one each, due to the conditions $`\mathrm{\Omega }_m+\mathrm{\Omega }_\mathrm{\Lambda }=1`$ and $`\mathrm{\Omega }_m/\mathrm{\Omega }_\mathrm{\Lambda }=2/(1+3w)`$ ($`w=0`$ or 1/3), respectively, in these models. While using data $`D_1`$, one can consider $`H_0`$ and $`M`$ as a single parameter by suitably combining them, but while using data $`D_2`$, we do not need to include $`M`$. However, in the present analysis using $`D_1`$, we consider $`H_0`$ and $`M`$ as independent parameters, as we did in section 3 and they are assumed to have flat priors in the ranges $`0.6<h<0.8`$ and $`19.6<M<19.1`$, respectively. This is a modification of the procedure in John & Narlikar (2002), but this alone does not affect the results in any way. Since the data consist of several common SNe and many of them are refined values of those used in the above paper, we decide to make their analysis independently of the one in that work. However, the information $`I`$ we have regarding the ranges of $`\mathrm{\Omega }_m`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$ are modified on the basis of the previous results. The new ranges chosen are $`0<\mathrm{\Omega }_m<1`$ and $`0<\mathrm{\Omega }_\mathrm{\Lambda }<1`$. One can see that small variations in these ranges do not affect our conclusions drastically.
The luminosity distance $`D`$ in these cases can generally be written as $`D=(c/H_0)d(z;\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda })`$. Thus the apparent magnitude is $`m=5\mathrm{log}(c/H_0)+5\mathrm{log}d(z;\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda })+25+M`$. The expression for $`d`$ to be used for $`M_1`$ and $`M_2`$ is
$$d=5\mathrm{log}\{(1+z)|\mathrm{\Omega }_k|^{1/2}S_k[|\mathrm{\Omega }_k|^{1/2}I(z)]\},$$
where $`\mathrm{\Omega }_k=1\mathrm{\Omega }_m\mathrm{\Omega }_\mathrm{\Lambda }`$ and $`S_k(x)=\mathrm{sin}x`$ for $`\mathrm{\Omega }_m+\mathrm{\Omega }_\mathrm{\Lambda }>1`$, $`S_k(x)=\mathrm{sinh}x`$ for $`\mathrm{\Omega }_m+\mathrm{\Omega }_\mathrm{\Lambda }<1`$ and $`S_k(x)=x`$ for $`\mathrm{\Omega }_m+\mathrm{\Omega }_\mathrm{\Lambda }=1`$. Also
$$I(z)=_0^z[(1+z^{})^2(1+\mathrm{\Omega }_mz^{})z^{}(2+z^{})(\mathrm{\Omega }_\mathrm{\Lambda })]^{1/2}𝑑z^{}.$$
For Model 3, the function $`d`$ can be written as
$$d=5\mathrm{log}\{m(1+z)S_k[\frac{1}{m}\mathrm{ln}(1+z)]\},$$
where $`m=\sqrt{2k/(3\mathrm{\Omega }_m2)}`$ for the nonrelativistic era and $`S_k(x)=\mathrm{sin}x`$ for $`\mathrm{\Omega }_m>2/3`$, $`S_k(x)=\mathrm{sinh}x`$ for $`\mathrm{\Omega }_m<2/3`$ and $`S_k(x)=x`$ for $`\mathrm{\Omega }_m=2/3`$.
The Bayes factors are evaluated as in the previous section and the results obtained in the present analysis are tabulated in Table 2. These show that the FLRW and inflationary models are at a more advantageous position than in John & Narlikar (2002), as per the interpretation of the Bayes factor.
## 5 DISCUSSION
Our attempt in this paper is to apply the Bayesian model comparison method to some vexing problems in cosmology in the light of latest $`mz`$ data of SNe Ia. These data are the results of potential landmark observations in cosmology, after the discovery of the Hubble’s law. Since they point to an accelerating universe today, it is natural to ask for how long the universe remained in this phase. The Bayesian analysis shows that except in one case, there is no evidence from SNe data to conclude that a changeover from deceleration to acceleration occurred anywhere in the past $`5\times 10^{17}`$ s.
But a rider we add to the above conclusion is even more important. The odds ratios between ‘always accelerating’ and ‘decelerating’ models in the given period, with impartial prior odds, are obtained to be close to unity, and that too in favor of the former model in most cases. However, as per the interpretation of this ratio, it is highly objectionable to state that one or the other of the models analyzed is ruled out, when it is on the basis of obtaining this ratio close to unity. To highlight this point, the model comparison exercise was performed to one of the classic tests of general relativity, viz., the deflection of light by sun, as it just grazes the sun’s surface. The huge value of the odds ratio (more strictly, the Bayes factor, since the prior odds was taken to be unity) obtained tells us how the Newtonian explanation of this phenomenon is ruled out by this observation. One can see that epoch-making discoveries are always accompanied by such large values for the Bayes factor. We wish anybody who claims that some particular model is ‘ruled out’ by some data to take note of this fact. Bayesian theory also teaches the important lesson that the price for overlooking this fact would be unavoidable and frequent U-turns. The present analysis rules out neither the accelerating nor the decelerating models; instead, we safely conclude that the data cannot discriminate these models.
The present analysis makes use of a significant improvement in the cosmographic approach adopted in John (2004). In order to overcome the problem of truncations in the series for $`1/a(t)`$ used in the equation $`1+z=a_0/a(t)`$ in the previous work, the solution of this equation is performed here in a purely numerical way. Though this consumes more computation time, it is advantageous that the accuracy is not compromised. As a consequence of the above modification, we had to perform a numerical integration in the expression for $`r_1`$. Though these amount to application of brute force in the analysis, the procedure has become much transparent. When compared to our own and other model-independent analyses of SNe data carried out previously (mentioned in the introduction), the present one is simpler and uses the least amount of clumsy formulas and this makes the treatment more intelligible. The procedure also makes use of the marginal likelihood method in the analysis. Since we are basically interested in the coefficients in the series expansion and the odds ratio, it is better to keep the treatment simple, to ensure accuracy at every stage. Considerable enhancement in such accuracy can be noticed from the plots of new marginal likelihoods for $`h`$, $`q_0`$, $`r_0`$ $`s_0`$,and $`u_0`$ (shown in Figs. 1-5, respectively, which are drawn using the data $`D_1`$), when compared with the corresponding curves in John (2004).
It may be recalled that the likelihood for the model $`(M_i)`$ computed using equation (4) is in fact the probability for the data, given the model and $`I`$ are true and is not exactly the ‘probability for the model’. Similarly, in the Bayesian scheme, the marginal likelihood for a parameter is not to be considered as the probability distribution for the parameter. However, we have computed the mean and $`\sigma `$ values from the above marginal likelihoods. Their new values are $`h=0.68\pm 0.06`$, $`q_0=0.90\pm 0.65`$, $`r_0=2.7\pm 6.7`$, $`s_0=36.5\pm 52.9`$, and $`u_0=142.7\pm 320`$. The large $`\sigma `$-values show that indeed there are a variety of choices for the values of the parameters, which may be considered as ‘good fit’ in the conventional way.
As stated earlier, the $`\chi ^2`$-values and hence the Bayes factor crucially depend on the error bars. The fact that the data cannot clearly discriminate different models in cosmology implies that the error bars in the data are quite large. Hopefully, future observations will have sufficiently small errors so that the Bayes factor between different models may become large. But some important points one should check in such cases is whether the errors are truly Gaussianly distributed and also whether we know accurately the standard deviation of these errors. Any deviation from these conditions will affect the validity of the model comparison. In cases where we are not sure about these two assumptions, one can resort to some other version of probability theory; for instance, the median statistics advocated by Gott, III et al. (2001) for analyzing cosmological datasets. These examples demonstrate how important are detailed considerations of the underlying assumptions in making judgments on what the data tells in cosmology.
It should be admitted that the Taylor expansion of the scale factor has more parameters to be constrained than realistic models and thereby it weakens the power of the data. But the cosmographic approach, which is basically a kinematic one, is not an alternative to realistic cosmological models. Instead, it helps to consolidate the evidences for such models. For example, evaluating the expansion parameters like $`h`$, $`q_0`$, $`r_0`$, $`s_0`$, etc. will help to constrain these models, as in the traditional use of the values of $`h`$ and $`q_0`$. It is also of use in answering questions like how long into the past the accelerating phase prevailed, as we have attempted to do in this paper. When we have more terms in the Taylor expansion, our extrapolation into the past becomes more accurate. Truncating at the fifth-order is due to practical considerations. Adding one more term would increase the computation time at least by an order of magnitude. We have seen that for extrapolating to the past corresponding to $`z1.75`$, up to fifth-order term in the expansion is required.
In this paper, we have also performed the comparison between some Friedmann models such as FLRW, inflationary and coasting models. This is viewed as an extension of a previous work using new and refined data. This analysis uses the previous one in modifying the flat priors. We use flat priors for the parameters $`\mathrm{\Omega }_m`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$ in the intervals $`0<\mathrm{\Omega }_m<1`$ and $`0<\mathrm{\Omega }_\mathrm{\Lambda }<1`$. The Bayes factors show that when data $`D_1`$ is used, the FLRW and inflationary models are at par with each other, and there is positive evidence against the coasting models. But in the case of using $`D_2`$, when compared to the the first model, there is some evidence against the inflationary one and a strong evidence against the coasting one. Also when compared to inflationary models, there is a positive evidence against the coasting model. But we again recall that these evidences are not very strong to rule out any of them. It is also to be noted that there is no contradiction between the two analyses in sections 3 and 4. The results in the latter section simply states that when compared to models $`M_1`$ and $`M_2`$, which are largely accelerating, the non-accelerating model $`M_3`$ is disfavored. Thus the evaluation of Bayes factor helps to quantify our knowledge, even when we are not aware of the full story.
These lessons are important in keeping cosmology a science. So far in the history of physics, though there were scientific revolutions, we always find that the new theories contain the old ones as limiting cases. This is because experimental evidences do not contradict the older theories in the realms in which they are applicable. For example, in a laboratory experiment where general relativistic effects are negligible, the Bayes factor between Newtonian and Einsteinian theories will be close to unity. Our results show that cosmological observations have not yet passed this stage and it is our opinion that the claims in cosmology, including the old and new ones we have compared in this paper, are not convincing enough since they are not supported by satisfactorily large Bayes factors. In physics we also see that in the realms where new theories are indispensable, they make unambiguous predictions, which are then verified experimentally. It can be seen that in such cases, new theories are often supported by huge Bayes factors, as in the example of light deflection discussed above in this paper. To achieve such goals in cosmology, we should tread with caution, begin with the cosmographic approach since it is the most fundamental one, and Bayes theory will help us to look for hard evidences. We have made such an attempt in this paper and the main conclusion, which is consistent with other approaches, is that the available data cannot clearly discriminate the cosmological models analyzed.
The author wishes to thank the UGC for a Research Grant and IUCAA, Pune, where part of this work was done, for hospitality. |
warning/0506/cond-mat0506300.html | ar5iv | text | # Wavepacket Dynamics, Quantum Reversibility and Random Matrix Theory
## 1 Introduction
In recent years there has been an increasing interest in understanding the theory of driven quantized chaotic systems . Driven systems are described by a Hamiltonian $`(Q,P,x(t))`$, where $`x(t)`$ is a time dependent parameter and $`(Q,P)`$ are some generalized actions. Due to the time dependence of $`x(t)`$, the energy of the system is not a constant of motion. Rather the system makes ”transitions” between energy levels, and therefore absorbs energy. This irreversible loss of energy is known as dissipation. To have a clear understanding of quantum dissipation we need a theory for the time evolution of the energy distribution.
Unfortunately, our understanding on quantum dynamics of chaotic systems is still quite limited. The majority of the existing quantum chaos literature concentrates on understanding the properties of eigenfunctions and eigenvalues. One of the main outcomes of these studies is the conjecture that Random Matrix Theory (RMT) modeling, initiated half a century ago by Wigner , can capture the universal aspects of quantum chaotic systems . Due to its large success RMT has become a major theoretical tool in quantum chaos studies , and it has found applications in both nuclear and mesoscopic physics (for a recent review see ). However, its applicability to quantum dynamics was left unexplored .
This paper extends our previous reports on quantum dynamics, both in detail and depth. Specifically, we analyze two dynamical schemes: The first is the so-called wavepacket dynamics associated with a rectangular pulse of strength $`+ϵ`$ which is turned on for a specified duration; The second involves an additional pulse followed by the first one which has a strength $`ϵ`$ and is of equal duration. We define this latter scheme as driving reversal scenario. We illuminate the direct relevance of our study with the studies of quantum irreversibility of energy spreading and consequently with quantum dissipation. We investigate the conditions under which maximum compensation is succeeded and define the notion of compensation (echo) time. To this end we rely both on numerical calculations performed for a chaotic system and on analytical considerations based on Linear Response Theory (LRT). The latter constitutes the leading theoretical framework for the analysis of driven systems and our study aims to clarify the limitations of LRT due to chaos. Our results are always compared with the outcomes of RMT modeling. We find that the RMT approach fails in general, to give the correct picture of wave-evolution. RMT can be trusted only to the extend that it gives trivial results that are implied by perturbation theory. Non-perturbative effects are sensitive to the underlying classical dynamics, and therefore the $`\mathrm{}0`$ behavior for effective RMT models is strikingly different from the correct semiclassical limit.
The structure of this paper is as follows: In the next section we discuss the notion of irreversibility which is related to driving reversal schemes and distinguish it from micro-reversibility which is associated with time reversal experiments. In Section 3 we discuss the driving schemes that we are using and we introduce the various observables that we will study in the rest of the paper. In Section 4, the model systems are introduced and an analysis of the statistical properties of the eigenvalues and the Hamiltonian matrix is presented. The Random Matrix Theory modeling is presented in Subsection 4.4. In Section 5 we introduce the concept of parametric regimes and exhibit its applicability in the analysis of parametric evolution of eigenstates . Section 6 extends the notion of regimes in dynamics and presents the results of Linear Response Theory for the variance and the survival probability. The Linear Response Theory (LRT) for the variance is analyzed in details in the following Subsection 6.1. In this subsection we also introduce the notion of restricted quantum-classical correspondence (QCC) and show that, as far as the second moment of the evolving wavepacket is concerned, both classical and quantum mechanical LRT coincides. In 6.5 we present in detail the results of LRT for the survival probability for the two driving schemes that we analyze. The following Sections 7 and 8 contain the results of our numerical analysis together with a critical comparison with the theoretical predictions obtained via LRT. Specifically in Sect.7, we present an analysis of wavepacket dynamics and expose the weakness of RMT strategy to describe wavepacket dynamics. In Sect.8 we study the evolution in the second half of the driving period and analyze the Quantum Irreversibility in energy spreading, where strong non-perturbative features are found for RMT models . Section 9 summarizes our findings.
## 2 Reversibility
The dynamics of either a classical or a quantum mechanical system is generated by a Hamiltonian $`(Q,P;x(t))`$ where $`x=(X_1,X_2,X_3,\mathrm{})`$ is a set of parameters that can be controlled from the ”outside”. In principle $`x`$ stands for the infinite set of parameters that describe the electric and magnetic fields acting on the system. But in practice the experimentalist can control only few parameters. A prototype example is a gas of particles inside a container with a piston. Then $`X_1`$ may be the position of the piston, $`X_2`$ may be some imposed electric field, and $`X_3`$ may be some imposed magnetic field. Another example is electrons in a quantum dot where some of the parameters $`X`$ represent gate voltages.
What do we mean by reversibility? Let us assume that the system evolves for some time. The evolution is described by
$$U[x]=\text{Exp}\left(\frac{i}{\mathrm{}}_0^t(x(t^{}))𝑑t^{}\right),$$
(1)
where Exp stands for time ordered exponentiation. In the case of the archetype example of a container with gas particles, we assume that there is a piston (position $`X`$) that is translated outwards ($`X^A(t)`$ increasing). Then we would ”undo” the evolution, by displacing the piston ”inwards” ($`X^B(t)`$ decreasing). In such a case the complete evolution is described by $`U[x]=U[x^B]U[x^A]`$. If we get $`U=1`$ (up to a phase factor), then it means that it is possible to bring the system back to its original state. In this case we say that the process $`U[x]`$ is reversible.
In the strict adiabatic limit the above described process is indeed reversible. What about the non-adiabatic case? In order to have a well posed question we would like to distinguish below between ”time reversal” and ”driving reversal” schemes.
### 2.1 Time reversal scheme
Obviously we are allowed to invent very complicated schemes in order to ”undo” the evolution. The ultimate scheme (in the case of the above example) involves reversal of the velocities. Assume that this operation is represented by $`U_T`$, then the reverse evolution is described by
$$U_{\text{reverse}}=U_TU[x^B]U_T,$$
(2)
where in $`x_B(t)`$ we have the time reversed piston displacement ($`X(t)`$) together with the sign of the magnetic field (if it exists) should be inverted. The question is whether $`U_T`$ can be realized. If we postulate that any unitary or anti-unitary transformation can be realized, then it follows trivially that any unitary evolution is ”micro-reversible”. But when we talk about reversibility (rather than micro-reversibility) we allow control over a restricted set of parameters (fields). Then the question is whether we can find a driving scheme, named $`x^T`$, such that
$$U_T=U[x^T]\text{???}$$
(3)
With such restriction it is clear that in general the evolution is not reversible.
Recently it has been demonstrated in an actual experiment that the evolution of spin system (cluster with many interacting spins) can be reversed. Namely, the complete evolution was described by $`U[x]=U[x^T]U[x^A]U[x^T]U[x^A]`$, where $`U[x^A]`$ is generated by some Hamiltonian $`_A=_0+\epsilon 𝒲`$. The term $`_0`$ represents the interaction between the spins, while the term $`𝒲`$ represents some extra interactions. The unitary operation $`U[x^T]`$ is realized using NMR techniques, and its effect is to invert the signs of all the couplings. Namely $`U[x^T]_0U[x^T]=_0`$. Hence the reversed evolution is described by
$$U_{\text{reverse}}=\mathrm{exp}\left(\frac{i}{\mathrm{}}t(_0+\epsilon 𝒲)\right),$$
(4)
which is the so-called Loschmidt Echo scenario. In principle we would like to have $`\epsilon =0`$ so as to get $`U=1`$, but in practice we have some un-controlled residual fields that influence the system, and therefore $`\epsilon 0`$. There is a huge amount of literature that discusses what happens in such scenario .
### 2.2 Driving reversal scheme
The above described experiment is in fact exceptional. In most cases it is possible to invert the sign of only one part of the Hamiltonian, which is associated with the driving field. Namely, if for instance $`U[x^A]`$ is generated by $`_A=H_0+\epsilon 𝒲`$, then we can realize
$$U_{\text{reverse}}=\mathrm{exp}\left(\frac{i}{\mathrm{}}t(_0\epsilon 𝒲)\right),$$
(5)
whereas Eq.(4) cannot be realized in general. We call such a typical scenario “driving reversal” in order to distinguish it from “time reversal” (Loschmidt Echo) scenario.
The study of “driving reversal” is quite different from the study of “Loschmidt Echo”. A simple minded point of view is that the two problems are formally equivalent because we simply permute the roles of $`_0`$ and $`𝒲`$. In fact there is no symmetry here. The main part of the Hamiltonian has in general an unbounded spectrum with well defined density of states, while the perturbation $`𝒲`$ is assumed to be bounded. This difference completely changes the “physics” of dynamics.
To conclude the above discussion we would like to emphasize that micro-reversibility is related to “time reversal” experiment which in general cannot be realized, while the issue of reversibility is related to “driving reversal”, which in principle can be realized. Our distinction reflects the simple observation that not any unitary or anti-unitary operation can be realized.
## 3 Object of the Study
In this paper we consider the issue of irreversibility for quantized chaotic systems. We assume for simplicity one parameter driving. We further assume that the variation of $`x(t)`$ is small in the corresponding classical system so that the analysis can be carried out with a linearized Hamiltonian. Namely,
$$(Q,P;x(t))_0+\delta x(t)𝒲,$$
(6)
where $`_0(Q,P;x(0))`$ and $`\delta x=x(t)x(0)`$. For latter purposes it is convenient to write the perturbation as
$$\delta x(t)=\epsilon \times f(t),$$
(7)
where $`\epsilon `$ controls the ”strength of the perturbation”, while $`f(t)`$ is the scaled time dependence. Note that if $`f(t)`$ is a step function, then $`\epsilon `$ is the ”size” of the perturbation, while if $`f(t)t`$ then $`\epsilon `$ is the ”rate” of the driving. In the representation of $`_0`$ we can write
$$=𝑬+\delta x(t)𝑩,$$
(8)
where by convention the diagonal terms of $`𝑩`$ are absorbed into the diagonal matrix $`𝑬`$. From general considerations that we explain later it follows that $`𝑩`$ is a banded matrix that looks random. This motivates the study of an Effective Banded Random Matrix (EBRM) model, as well as its simplified version which is the standard Wigner Banded Random Matrix (WBRM) model. (See detailed definitions in the following).
In order to study the irreversibility for a given driving scenario, we have to introduce measures that quantify the departure from the initial state. We define a set of such measures in the following subsections.
### 3.1 The evolving distribution $`P_t(n|n_0)`$
Given the Hamiltonian $`(Q,P;x)`$, an initial preparation at state $`|n_0`$, and a driving scenario $`x(t)`$, it is most natural to analyze the evolution of the probability distribution
$$P_t(n|n_0)=|n|U(t)|n_0|^2.$$
(9)
We always assume that $`x(t)=x(0)`$.
By convention we order the states by their energy. Hence we can regard $`𝒫_t(n|n_0)`$ as a function of $`r=nn_0`$, and average over the initial preparation, so as to get a smooth distribution $`𝒫_t(r)`$.
The survival probability is defined as
$$𝒫(t)=|n_0|U(t)|n_0|^2=P_t(n_0|n_0),$$
(10)
and the energy spreading is defined as
$$\delta E(t)=\sqrt{\underset{n}{}P_t(n|n_0)(E_nE_{n_0})^2}.$$
(11)
These are the major measures for the characterization of the distribution. In later sections we would like to analyze their time evolution.
The physics of $`\delta E(t)`$ is very different from the physics of $`𝒫(t)`$ because the former is very sensitive to the tails of the distribution. Yet, the actual ”width” of the distribution is not captured by any of these measures. A proper measure for the width can be defined as follows:
$$\delta E_{\text{core}}(t)=[n_{75\%}n_{25\%}]\mathrm{\Delta },$$
(12)
where $`\mathrm{\Delta }`$ is the mean level spacing and $`n_q`$ is determined through the equation $`_nP_t(n|n_0)=q`$. Namely it is the width of the main body of the distribution. Still another characteristic of the distribution is the participation ratio $`\delta n_{\text{IPR}}(t)`$. It gives the number of levels that are occupied at time $`t`$ by the distribution. The ratio $`\delta n_{\text{IPR}}/(n_{75\%}n_{25\%})`$ can be used as a measure for sparsity. We assume in this paper strongly chaotic systems, so sparsity is not an issue and $`\delta n_{\text{IPR}}\delta E_{\text{core}}/\mathrm{\Delta }`$.
### 3.2 The compensation time $`t_r`$
In this paper we consider two types of driving schemes. Both driving schemes are presented schematically in Figure 1.
The first type of scheme is the wavepacket dynamics scenario for which the driving is turned-on at time $`t=0`$ and turned-off at a later time $`t=T`$.
The second type of scenario that we investigate is what we call driving reversal. In this scenario the initial rectangular pulse is followed by a compensating pulse of equal duration. The total period of the cycle is $`T`$.
In Figure 9 we show representative results for the time evolution of $`\delta E(t)`$ in a wavepacket scenario, while in Figure 12 we show what happens in case of a driving reversal scenario. Corresponding plots for $`𝒫(t)`$ are presented in Figure 13. We shall define the models and we shall discuss the details of these figures later on. At this stage we would like to motivate by inspection of these figures the definition of ”compensation time”.
We define the compensation time $`t_r`$, as the time after the driving reversal, when maximum compensation (maximum return) is observed. If it is determined by the maximum of the survival probability kernel $`𝒫(t)`$, then we denote it as $`t_r^P`$. If it is determined by the minimum of the energy spreading $`\delta E(t)`$ then we denote it as $`t_r^E`$. It should be remembered that the theory of $`𝒫(t)`$ and $`\delta E(t)`$ is not the same, hence the distinction in the notation. The time of maximum compensation is in general not $`t_r=T`$ but rather
$$T/2<t_r<T.$$
(13)
We emphasize this point because the notion of ”echo”, as used in the literature, seems to reflect a false assertion .
For the convenience of the reader we concentrate in the following table on the major notations in this paper:
| Notation | explanation | reference |
| --- | --- | --- |
| $`(Q,P;x(t))`$ | classical linearized Hamiltonian | Eq.(6) |
| $`(t)`$ | generalized force | Eq.(17) |
| $`C(\tau )`$ | correlation function | Eq.(18) |
| $`\tau _{\text{cl}}`$ | correlation time | |
| $`\stackrel{~}{C}(\omega )`$ | fluctuation spectrum | Eq.(19) |
| $`=𝑬+\delta x𝑩`$ | The Hamiltonian matrix | Eq.(8) |
| 2DW | the physical model system | Eq.(14) |
| EBRM | the corresponding RMT model | |
| WBRM | the Wigner RMT model | |
| $`\mathrm{\Delta }`$ | mean level spacing | Eq.(23) |
| $`\mathrm{\Delta }_b`$ | energy bandwidth | Eq.(24) |
| $`\sigma `$ | RMS of near diagonal couplings | |
| $`P_{\text{spacings}}(s)`$ | energy spacing distribution | Eq.(16) |
| $`P_{\text{couplings}}(q)`$ | distribution of couplings | Eq.(25) |
| $`E_n(x)`$ | eigen-energies of the Hamiltonian | |
| $`E_nE_mr\mathrm{\Delta }`$ | estimated energy difference for $`r=nm`$ | |
| $`P(n|m)`$ | overlaps of eigenstates given a constant perturbation $`\epsilon `$ | Eq.(27) |
| $`P(r)`$ | smoothed version of $`P(n|m)`$ | |
| $`\mathrm{\Gamma }(\delta x)`$ | the number of levels that are mixed non-perturbatively | |
| $`\delta E_{\text{cl}}\delta x`$ | the classical width of the LDoS | Eq.(29) |
| $`\delta x=\epsilon f(t)`$ | driving scheme | Eq.(7) |
| $`T`$ | The period of the driving cycle (if applicable) | |
| $`P_t(n|m)`$ | the transition probability | Eq.(9) |
| $`P_t(r)`$ | smoothed version of $`P_t(n|n_0)`$ | |
| $`𝒫(t)`$ | the survival probability $`P_t(n_0|n_0)`$ | Eq.(10) |
| $`p(t)=1𝒫(t)`$ | total transition probability | Eq.(47) |
| $`\delta E(t)`$ | energy spreading | Eq.(11) |
| $`\delta E_{\text{core}}(t)`$ | the ”core” width of the distribution | Eq.(12) |
| $`t_r^P`$ | compensation time for the survival probability | |
| $`t_r^E`$ | compensation time for the energy spreading | |
| $`t_{\text{prt}},t_{\text{sdn}},t_{\text{erg}}`$ | various time scales in the dynamics | Eq.(59,62,70) |
| $`\epsilon _c,\epsilon _{\text{prt}}`$ | borders between regimes | Eq.(31,33) |
| $`P_{\text{FOPT}},P_{\text{prt}},P_{\text{sc}}`$ | various approximations to $`P()`$ | Eq.(30,32,34,35) |
## 4 Modeling
We are interested in quantized chaotic systems that have few degrees of freedom. The dynamical system used in our studies is the Pullen-Edmonds model . It consists of two harmonic oscillators that are nonlinearly coupled. The corresponding Hamiltonian is
$$(Q,P;x)=\frac{1}{2}\left(P_1^2+P_2^2+Q_1^2+Q_2^2\right)+xQ_1^2Q_2^2.$$
(14)
The mass and the frequency of the harmonic oscillators are set to one. Without loss of generality we set $`x(0)=x_0=1`$. Later we shall consider classically small deformations ($`\delta x1`$) of the potential. One can regard this model (14) as a description of a particle moving in a two dimensional well (2DW). The energy $`E`$ is the only dimensionless parameter of the classical motion. For high energies $`E>5`$ the motion of the Pullen-Edmonds model is ergodic. Specifically it was found that the measure of the chaotic component on the Poincaré section deviates from unity by no more than $`10^3`$.
In Figure 2 we display the equipotential contours of the model Hamiltonian (14) with $`x_0=1`$. We observe that the equipotential surfaces are circles but as the energy is increased they become more and more deformed leading to chaotic motion. Our analysis is focused on an energy window around $`E3`$ where the motion is mainly chaotic. This is illustrated in the right panel of Figure 2 where we report the Poincaré section (of the phase space) of a selected trajectory, obtained from $`_0`$ at $`E=3`$. The ergodicity of the motion is illustrated by the Poincaré section, filling the plane except from some tiny quasi- integrable islands.
The perturbation is described by $`𝒲=Q_1^2Q_2^2`$. In the classical analysis there is only one significant regime for the strength of the perturbation. Namely, the perturbation is considered to be classically small if
$$\delta x\epsilon _{\text{cl}},$$
(15)
where $`\epsilon _{\text{cl}}=1`$. This is the regime where (classical) linear analysis applies. Namely, within this regime the deformation of the energy surface $`_0=E`$ can be described as a linear process (see Eq. (29)).
### 4.1 Energy levels
Let us now quantize the system. For obvious reasons we are considering a de-symmetrized 1/8 well with Dirichlet boundary conditions on the lines $`Q_1=0`$, $`Q_2=0`$ and $`Q_1=Q_2`$. The matrix representation of $`_0`$ in the basis of the un-coupled system is very simple. The eigenstates of the Hamiltonian $`_0`$ are then obtained numerically.
As mentioned above, we consider the experiments to take place in an energy window $`2.8<E<3.1`$ which is classically small and where the motion is predominantly chaotic. Nevertheless, quantum mechanically, this energy window is large, i.e., many levels are found therein. The local mean level spacing $`\mathrm{\Delta }(E)`$ at this energy range is given approximately by $`\mathrm{\Delta }4.3\mathrm{}^2`$. The smallest $`\mathrm{}`$ that we can handle is $`\mathrm{}=0.012`$ resulting in a matrix size of about $`4000\times 4000`$. Unless stated otherwise, all the numerical data presented below correspond to a quantization with $`\mathrm{}=0.012`$.
As it was previously mentioned in the introduction, the main focus of quantum chaos studies has so far been on issues of spectral statistics . In this context it turns out that the sub -$`\mathrm{}`$ statistical features of the energy spectrum are ”universal”, and obey the predictions of RMT. In particular we expect that the level spacing distribution $`P(s)`$ of the ”unfolded” (with respect to $`\mathrm{\Delta }`$) level spacings $`s_n=(E_{n+1}E_n)/\mathrm{\Delta }`$ will follow with high accuracy the so-called *Wigner* *surmise*. For systems with time reversal symmetry it takes the form
$$P_{\text{spacings}}(s)=\frac{\pi }{2}s\text{e}^{\frac{\pi }{4}s^2},$$
(16)
indicating that there is a linear repulsion between nearby levels. Non-universal (i.e. system specific) features are reflected only in the large scale properties of the spectrum and constitute the fingerprints of the underlying classical chaotic dynamics.
The de-symmetrized 2DW model shows time reversal symmetry, and therefore we expect the distribution to follow Eq.(16). The analysis is carried out only for the levels contained in the chosen energy window around $`E=3`$. Instead of plotting $`P(s)`$ we show the integrated distribution $`I(s)=_0^sP(s^{})𝑑s^{}`$, which is independent of the bin size of the histogram. In Figure 3 we present our numerical data for $`I(s)`$ while the inset shows the deviations from the theoretical prediction (16). The agreement with the theory is fairly good and the level repulsion is clearly observed. The observed deviations have to be related on the one hand to the tiny quasi-integrable islands that exist at $`E=3`$ as well as to rather limited level statistics.
### 4.2 The band-profile
In this subsection we explain that the band-structure of $`𝑩`$ is related to the fluctuations of the classical motion. This is the major step towards RMT modeling.
Consider a given ergodic trajectory $`(Q(t),P(t))`$ on the energy surface $`(Q(0),P(0);x_0)=E`$. An example is shown in Fig. 2b. We can associate with it a stochastic-like variable
$$(t)=\frac{}{x}(Q(t),P(t),x(t)),$$
(17)
which for our linearized Hamiltonian is simply the perturbation term $`=𝒲=Q_1^2Q_2^2`$. It can be interpreted as the generalized force that acts on the boundary of the 2D well. It may have a non-zero average (“conservative” part) but below we are interested only in its fluctuations.
In order to characterize the fluctuations of $`(t)`$ we introduce the autocorrelation function $`C(\tau )`$
$$C(\tau )=(t)(t+\tau )^2.$$
(18)
The angular brackets denote an averaging which is either micro-canonical over some initial conditions $`(Q(0),P(0))`$ or temporal (due to the assumed ergodicity). The power spectrum for the 2D well model is shown in Fig.4 (see solid line).
For generic chaotic systems (described by smooth Hamiltonians), the fluctuations are characterized by a short correlation time $`\tau _{\text{cl}}`$, after which the correlations are negligible. In generic circumstances $`\tau _{\text{cl}}`$ is essentially the ergodic time. For our model system $`\tau _{\text{cl}}1`$.
The power spectrum of the fluctuations $`\stackrel{~}{C}(\omega )`$ is defined by a Fourier transform:
$$\stackrel{~}{C}(\omega )=_{\mathrm{}}^{\mathrm{}}C(\tau )\mathrm{exp}(i\omega \tau )𝑑\tau .$$
(19)
This power spectrum is characterized by a cut-off frequency $`\omega _{\text{cl}}`$ which is inverse proportional to the classical correlation time
$$\omega _{\text{cl}}=\frac{2\pi }{\tau _{\text{cl}}}.$$
(20)
Indeed in the case of our model system we get $`\omega _{\text{cl}}7`$ which is in agreement with Fig.4.
The implication of having a short but non-vanishing classical correlation time $`\tau _{\text{cl}}`$ is having large but finite bandwidth in the perturbation matrix $`𝑩`$. This follows from the identity
$$\stackrel{~}{C}(\omega )=\underset{n}{}|B_{nm}|^22\pi \delta \left(\omega \frac{E_nE_m}{\mathrm{}}\right),$$
(21)
which implies
$$|𝑩_{nm}|^2=\frac{\mathrm{\Delta }}{2\pi \mathrm{}}\stackrel{~}{C}\left(\omega =\frac{E_nE_m}{\mathrm{}}\right).$$
(22)
Hence the matrix elements of the perturbation matrix $`𝑩`$ are extremely small outside of a band of width $`b=\mathrm{}\omega _{\text{cl}}/\mathrm{\Delta }`$.
In the inset of Figure 5 we show a snapshot of the perturbation matrix $`|𝑩_{nm}|^2`$. It clearly shows a band-structure. At the same figure we also display the band-profiles for different values of $`\mathrm{}`$. A good agreement with the classical power spectrum $`\stackrel{~}{C}(\omega )`$ is evident.
It is important to realize that upon quantization we end up with two distinct energy scales. One is obviously the mean level spacing (see previous subsection)
$$\mathrm{\Delta }\mathrm{}^d,$$
(23)
where the dimensionality is $`d=2`$ in case of our model system. The other energy scale is the bandwidth
$$\mathrm{\Delta }_b=\frac{2\pi \mathrm{}}{\tau _{\text{cl}}}=b\mathrm{\Delta }.$$
(24)
This energy scale is also known in the corresponding literature as the ”non-universal” energy scale , or (in case of diffusive motion) as the Thouless energy <sup>1</sup><sup>1</sup>1The dimensionless parameter $`b`$ scales like $`b\mathrm{}^{(d1)}`$ and in the frame of mesoscopic systems is recognized as the dimensionless Thouless conductance.. One has to notice that deep in the semiclassical limit $`\mathrm{}0`$ these two energy scales differ enormously from one another (provided $`d2`$). We shall see in the following sections that this scale separation has dramatic consequences in the theory of driven quantum systems.
### 4.3 Distribution of couplings
We investigate further the statistical properties of the matrix elements $`𝑩_{nm}`$ of the perturbation matrix, by studying their distribution. RMT assumes that upon appropriate unfolding they must be distributed in a Gaussian manner. The ’unfolding’ aims to remove system specific properties and reveal the underlying universality. It is done by normalizing the matrix elements with the local standard deviation $`\sigma =\sqrt{|𝐁_{nm}|^2}`$ related through Eq. (22) with the classical power spectrum $`\stackrel{~}{C}(\omega )`$.
The existing literature is not conclusive about the distribution of the normalized matrix elements $`q=𝑩_{nm}/\sigma `$. Specifically, Berry and more recently Prosen , claimed that $`𝒫(q)`$ should be Gaussian. On the other hand, Austin and Wilkinson have found that the Gaussian is approached only in the limit of high quantum numbers while for small numbers, i.e., low energies, a different distribution applies, namely
$$P_{\text{couplings}}(q)=\frac{\mathrm{\Gamma }(\frac{N}{2})}{\sqrt{\pi N}\mathrm{\Gamma }(\frac{N1}{2})}\left(1\frac{q^2}{N}\right)^{(N3)/2}.$$
(25)
This is the distribution of the elements of an $`N`$-dimensional vector, distributed randomly over the surface of an $`N`$-dimensional sphere of radius $`\sqrt{N}`$. For $`N\mathrm{}`$ this distribution approaches a Gaussian.
The distribution $`𝒫(q)`$ for our model is reported in Figure 6. The solid line corresponds to a Gaussian of unit variance while the dashed-dotted line is obtained by fitting Eq. (25) to the numerical data using $`N`$ as a fitting parameter. We observe that the Gaussian resembles better our numerical data although deviations, especially for matrix elements close to zero, can be clearly seen. We attribute these deviations to the existence of the tiny stability islands in the phase space. Trajectories started in those islands cannot reach the chaotic sea and vice versa. Quantum mechanically the consequence of this would be vanishing matrix elements $`𝑩_{nm}`$ which represent the classically forbidden transitions.
### 4.4 RMT modeling
It was the idea of Wigner more than forty years ago, to study a simplified model, where the Hamiltonian is given by Eq. (8), and where $`𝑩`$ is a Banded Random Matrix (BRM) . The diagonal matrix $`𝑬`$ has elements which are the ordered energies $`\{E_n\}`$, with mean level spacing $`\mathrm{\Delta }`$. The perturbation matrix $`𝑩`$ has a *rectangular* band-profile of band-size $`b`$. Within the band $`0<|nm|b`$ the elements are independent random variables given by a Gaussian distribution with zero mean and a variance $`\sigma ^2=|𝑩_{nm}|^2`$. Outside the band they vanish. We refer to this model as the *Wigner BRM* model (WBRM).
Given the band-profile, we can use Eq.(22) in reverse direction to calculate the correlation function $`C(\tau )`$. For the WBRM model we get
$$C(\tau )=2\sigma ^2b\mathrm{sinc}\left(\tau /\tau _{\text{cl}}\right),$$
(26)
where $`\tau _{\text{cl}}=\mathrm{}/\mathrm{\Delta }_b`$. Thus, there are three parameters $`(\mathrm{\Delta },b,\sigma )`$ that define the WBRM model.
The WBRM model can be regarded as a simplified local description of a true Hamiltonian matrix. This approach is attractive both analytically and numerically. Analytical calculations are greatly simplified by the assumption that the off-diagonal terms can be treated as independent random numbers. Also from a numerical point of view it is quite a tough task to calculate the true matrix elements of the $`𝑩`$ matrix. It requires a preliminary step where the chaotic $`_0`$ is diagonalized. Due to memory limitations one ends up with quite small matrices. For the Pullen-Edmonds model we were able to handle matrices of final size $`N=4000`$ maximum. This should be contrasted with the WBRM simulations, where using self -expanding algorithm we were able to handle system sizes up to $`N=100000`$ along with significantly reduced CPU time.
We would like to stress again that the underlying assumption of WBRM, namely that the off-diagonal elements are uncorrelated random numbers, has to be treated with extreme care.
The WBRM model involves an additional simplification. Namely, one assumes that $`𝑩`$ has a *rectangular* band-profile. A simple inspection of the band-profile of our model Eq. (14) shows that this is not the case (see Fig. 5). We eliminate this simplification by introducing a RMT model that is even closer to the dynamical one. To this end, we randomize the signs of the off-diagonal elements of the perturbation matrix $`𝑩`$ keeping its band-structure intact. This procedure leads to a random model that exhibits only universal properties while it lack any semiclassical limit. We will refer to it as the *effective* banded random matrix model (EBRM).
## 5 The Parametric Evolution of the Eigenfunctions
As we change the parameter $`\delta x`$ in the Hamiltonian Eq. (8), the instantaneous eigenstates $`\{|n(x)\}`$ evolve and undergo structural changes. In order to understand the actual dynamics, it is important to understand these structural changes. This leads to the introduction of
$$P(n|m)=|n(x)|m(x_0)|^2,$$
(27)
which is easier to analyze than $`P_t(n|n_0)`$. Up to some trivial scaling and shifting $`P(n|m)`$ is essentially the local density of states (LDoS):
$$P(E|m)=\underset{n}{}|n(x)|m(x_0)|^2\delta (EE_n).$$
(28)
The averaged distribution $`P(r)`$ is defined in complete analogy with the definition of $`P_t(r)`$. Namely, we use the notation $`r=nm`$, and average over several $`m`$ states with roughly the same energy $`E_mE`$.
Generically $`P(r)`$ undergoes the following structural changes as a function of growing $`\delta x`$. We first summarize the generic picture, which involves the parametric scales $`\epsilon _c`$ and $`\epsilon _{\text{prt}}`$. and the approximations $`P_{\text{FOPT}}`$, $`P_{\text{prt}}`$, and $`P_{\text{sc}}`$. Then we discuss how to determine these scales, and what these approximations are.
* The first order perturbative theory (FOPT) regime is defined as the range $`\delta x<\epsilon _c`$ where we can use FOPT to get an approximation that we denote as $`P()P_{\text{FOPT}}`$.
* The (extended) perturbative regime is defined as the range $`\epsilon _c<\delta x<\epsilon _{\text{prt}}`$ where we can use perturbation theory (to an infinite order) to get a meaningful approximation that we denote as $`P()P_{\text{prt}}`$. Obviously $`P_{\text{prt}}`$ reduces to $`P_{\text{FOPT}}`$ in the FOPT regime.
* The non-perturbative regime is defined as the range $`\delta x>\epsilon _{\text{prt}}`$ where perturbation theory becomes non-applicable. In this regime we have to use either RMT or semiclassics in order to get an approximation that we denote as $`P()P_{\text{sc}}`$.
Irrespective of these structural changes, it can be proved that the variance of $`P(r)`$ is strictly linear and given by the expression
$$\delta E(\delta x)=\sqrt{C(0)}\delta x\delta E_{\text{cl}}.$$
(29)
The only assumption that underlines this statement is $`\delta x\epsilon _{\text{cl}}`$. It reflects the linear departure of the energy surfaces.
### 5.1 Approximations for $`P(n|m)`$
The simplest regime is obviously the FOPT regime where, for $`P(n|m)`$, we can use the standard textbook approximation $`P_{\text{FOPT}}(n|m)1`$ for $`n=m`$, while
$$P_{\text{FOPT}}(n|m)=\frac{\delta x^2|𝑩_{nm}|^2}{(E_nE_m)^2},$$
(30)
for $`nm`$. If outside of the band we have $`𝑩_{nm}=0`$, as in the WBRM model, then $`P_{\text{FOPT}}(r)=0`$ for $`|r|>b`$. To find the higher order tails (outside of the band) we have to go to higher orders in perturbation theory. Obviously this approximation makes sense only as long as $`\delta x<\epsilon _c`$ where
$$\epsilon _c=\mathrm{\Delta }/\sigma \mathrm{}^{(1+d)/2},$$
(31)
and $`d`$ is the degrees of freedom of our system ($`d=2`$ for the 2D well model).
If $`\delta x>\epsilon _c`$ but not too large then we still have tail regions which are described by FOPT. This is a non-trivial observation which can be justified by using perturbation theory to infinite order. Then we can argue that a reasonable approximation is
$$P_{\text{prt}}(n|m)=\frac{\delta x^2|𝑩_{nm}|^2}{(E_nE_m)^2+\mathrm{\Gamma }^2},$$
(32)
where $`\mathrm{\Gamma }`$ is evaluated by imposing normalization of $`P_{\mathrm{prt}}(n|m)`$. In the case of WBRM model $`\mathrm{\Gamma }=(\sigma \delta x/\mathrm{\Delta })^2\times \mathrm{\Delta }`$. The appearance of $`\mathrm{\Gamma }`$ in the above expression cannot be obtained from any finite-order perturbation theory: Formally it requires summation to infinite order. Outside of the bandwidth the tails decay faster than exponentially. Note that $`P_{\text{prt}}(n|m)`$ is a Lorentzian in the case of a flat bandwidth (WBRM model), while in the general case it can be described as a ”core-tail” structure.
Obviously the above approximation makes sense only as long as $`\mathrm{\Gamma }(\delta x)<\mathrm{\Delta }_b`$. This expression assumes that the bandwidth $`\mathrm{\Delta }_b`$ is sharply defined, as in the WBRM model. By elimination this leads to the determination of $`\epsilon _{\text{prt}}`$, which in case of the WBRM model is simply
$$\epsilon _{\text{prt}}=\sqrt{b}\epsilon _c\frac{\mathrm{}}{\tau _{\text{cl}}\sqrt{C(0)}}.$$
(33)
In more general cases the bandwidth is not sharply defined. Then we have to define the perturbative regime using a practical numerical procedure. The natural definition that we adopt is as follows. We calculate the spreading $`\delta E(\delta x)`$, which is a linear function. Then we calculate $`\delta E_{\text{prt}}(\delta x)`$, using Eq.(32)). This quantity always saturates for large $`\delta x`$ because of having finite bandwidth. We compare it to the exact $`\delta E(\delta x)`$, and define $`\epsilon _{\text{prt}}`$ for instance as the $`80\%`$ departure point.
What happens if perturbation theory completely fails? In the WBRM model the LDoS becomes semicircle:
$$P_{\text{sc}}(n|m)=\frac{1}{2\pi \mathrm{\Delta }}\sqrt{4\left(\frac{E_nE_m}{\mathrm{\Delta }}\right)^2},$$
(34)
while in systems that have a semiclassical limit we expect to get
$$P_{\text{sc}}(n|m)=\frac{dQdP}{(2\pi \mathrm{})^d}\rho _n(Q,P)\rho _m(Q,P),$$
(35)
where $`\rho _m(Q,P)`$ and $`\rho _n(Q,P)`$ are the Wigner functions that correspond to the eigenstates $`|m(x_0)`$ and $`|n(x)`$ respectively.
### 5.2 The $`P(n|m)`$ in practice
There are some findings that go beyond the above generic picture and, for completeness, we mention them. The first one is the ”localization regime” which is found in the case of the WBRM model for $`\epsilon >\epsilon _{\text{loc}}`$. where
$$\epsilon _{\text{loc}}=b^{3/2}\epsilon _c.$$
(36)
In this regime it is important to distinguish between the non-averaged $`P(n|m)`$ and the averaged $`P(r)`$ because the eigenfunctions are non-ergodic but rather localized. This localization is not reflected in the LDoS which is still a semicircle. A typical eigenstate is exponentially localized within an energy range $`\delta E_\xi =\xi \mathrm{\Delta }`$ much smaller than $`\delta E_{\text{cl}}`$. The localization length is $`\xi b^2`$. In actual physical applications it is not clear whether there is such a type of localization. The above scenario for the WBRM model is summarized in Fig. 7 where we plot $`P(n|m)`$ in the various regimes. The localized regime is not an issue in the present work and therefore we will no further be concerned with it.
The other deviation from the generic scenario, is the appearance of a non-universal ”twilight regime” which can be found for some quantized systems . In this regime a co-existence of a perturbative and a semiclassical structure can be observed. For the Pullen-Edmonds model (14) there is no such distinct regime.
For the Hamiltonian model described by Eq. (14) the borders between the regimes can be estimated . Namely $`\epsilon _c3.8\mathrm{}^{3/2}`$ and $`\epsilon _{\text{prt}}5.3\mathrm{}`$. In Fig. 8 we report the parametric evolution of the eigenstates for the Hamiltonian model of Eqs. (14) and we compare the outcomes with the results of the EBRM model . Despite the overall quantitative agreement, some differences can be detected:
* In the FOPT regime (see Fig. 8a), the RMT strategy fails in the far tails regime $`\mathrm{\Delta }\times |r|>\mathrm{\Delta }_b`$ where system specific interference phenomena become important.
* In the extended perturbative regime (see Fig. 8b) the line-shape of the averaged wavefunction $`P(n|m)`$ is different from Lorentzian. Still the general features of $`P_{\text{prt}}`$ (core-tail structure) can be detected. In a sense, Wigner’s Lorentzian (32) is a special case of core-tail structure. Finally, as in the standard perturbative regime one observes that the far-tails are dominated by either destructive interference (left tail), or by constructive interference (right tail).
* Deep in the non-perturbative regime ( $`\epsilon >\epsilon _{\text{prt}}`$ ) the overlaps $`P(n|m)`$ are well approximated by the semiclassical expression. The exact shape is determined by simple classical considerations . This is in contrast to the WBRM model which does not have a classical limit.
## 6 Linear Response Theory
The definition of regimes for driven systems is more complicated than the corresponding definition in case of LDoS theory. It is clear that for short times we always can use time-dependent FOPT. The question is, of course, what happens next. There we have to distinguish between two types of scenarios. One type of scenario is wavepacket dynamics for which the dynamics is a transient from a preparation state to some new ergodic state. The second type of scenario is persistent driving, either linear driving ($`\dot{x}=\epsilon `$) or periodic driving ($`x(t)=\epsilon \mathrm{sin}(\mathrm{\Omega }t)`$). In the latter case the strength of the perturbation depends on the rate of the driving, not just on the amplitude. The relevant question is whether the long time dynamics can be deduced from the short time analysis. To say that the dynamics is of perturbative nature means that the short time dynamics can be deduced from FOPT, while the long time dynamics can be deduced on the basis of a Markovian (stochastic) assumption. The best known example is the derivation of the exponential Wigner law for the decay of metastable state. The Fermi-Golden-Rule (FGR) is used to determine the initial rate for the escaping process, and then the long-time result is extrapolated by assuming that the decay proceeds in a stochastic-like manner. Similar reasoning is used in deriving the Pauli master equation which is used to describe the stochastic-like transitions between the energy levels in atomic systems.
A related question to the issue of regimes is the validity of Linear Response Theory (LRT). In order to avoid ambiguities we adopt here a practical definition. Whenever the result of the calculation depends only on the two point correlation function $`C(\tau )`$, or equivalently only on the band-profile of the perturbation (which is described by $`\stackrel{~}{C}(\omega )`$), then we refer to it as ”LRT”. This implies that higher order correlations are not expressed. There is a (wrong) tendency to associate LRT with FOPT. In fact the validity of LRT is not simply related to FOPT. We shall clarify this issue in the next section.
For both $`\delta E(t)`$ and $`𝒫(t)`$ we have ”LRT formulas” which we discuss in the next sections. Writing the driving pulse as $`\delta x(t)=\epsilon f(t)`$ for the spreading we get:
$$\delta E^2(t)=\epsilon ^2\times _{\mathrm{}}^{\mathrm{}}\frac{d\omega }{2\pi }\stackrel{~}{C}(\omega )\stackrel{~}{F}_t(\omega ),$$
(37)
while for the survival probability we have
$$𝒫(t)=\mathrm{exp}\left(\epsilon ^2\times _{\mathrm{}}^{\mathrm{}}\frac{d\omega }{2\pi }\stackrel{~}{C}(\omega )\frac{\stackrel{~}{F}_t(\omega )}{(\mathrm{}\omega )^2}\right).$$
(38)
Two spectral functions are involved: One is the power spectrum $`\stackrel{~}{C}(\omega )`$ of the fluctuations defined in Eq. (19), and the other $`\stackrel{~}{F}_t(\omega )`$ is the spectral content of the driving pulse which is defined as
$$\stackrel{~}{F}_t(\omega )=\left|_0^t𝑑t^{}\dot{f}(t^{})\mathrm{e}^{i\omega t^{}}\right|^2.$$
(39)
Here we summarize the main observations regarding the nature of wavepacket dynamics in the various regimes:
* FOPT regime: In this regime $`𝒫(t)1`$ for all time, indicating that all probability is all the time concentrated on the initial level. An alternative way to identify this regime is from $`\delta E_{\text{core}}(t)`$ which is trivially equal to $`\mathrm{\Delta }`$.
* Extended perturbative regime: The appearance of a core-tail structure which is characterized by separation of scales $`\mathrm{\Delta }\delta E_{\text{core}}(t)\delta E(t)\mathrm{\Delta }_b`$. The core is of non-perturbative nature, but the variance $`\delta E^2(t)`$ is still dominated by the tails. The latter are described by perturbation theory.
* Non-perturbative regime: The existence of this regime is associated with having the finite energy scale $`\mathrm{\Delta }_b`$. It is characterized by $`\mathrm{\Delta }_b\delta E_{\text{core}}(t)\delta E(t)`$. As implied by the terminology, perturbation theory (to any order) is not a valid tool for the analysis of the energy spreading. Note that in this regime, the spreading profile is characterized by a single energy scale ($`\delta E\delta E_{\text{core}}`$).
### 6.1 The energy spreading $`\delta E(t)`$
Of special importance for understanding quantum dissipation is the theory for the variance $`\delta E^2(t)`$ of the energy spreading. Having $`\delta E(t)ϵ`$ means linear response. If $`\delta E(t)/ϵ`$ depends on $`ϵ`$, we call it “non-linear response”. In this paragraph we explain that linear response theory (LRT) is based on the “LRT formula” Eq.(37) for the spreading. This formula has a simple classical derivation (see Subsection 6.2 below).
From now on it goes without saying that we assume the classical conditions for the validity of Eq.(37) are satisfied (no $`\mathrm{}`$ involved in such conditions). The question is what happens to the validity of LRT once we “quantize” the system. In previous publications, we were able to argue the following:
The LRT formula can be trusted in the perturbative regime, with the exclusion of the adiabatic regime. In the sudden limit the LRT formula can also be trusted in the non-perturbative regime. In general the LRT formula cannot be trusted in the non-perturbative regime. The LRT formula can be trusted deep in the non-perturbative regime, provided the system has a classical limit.
For a system that does not have a classical limit (Wigner model) we were able to demonstrate that LRT fails in the non-perturbative regime. Namely, for the WBRM model the response $`\delta E(t)/ϵ`$ becomes $`ϵ`$ dependent for large $`ϵ`$, meaning that the response is non-linear. Hence the statement in item (C) above has been established. We had argued that the observed non-linear response is the result of a quantal non-perturbative effect. Do we have a similar type of non-linear response in the case of quantized chaotic systems? The statement in item (D) above seems to suggest that the observation of such non-linearity is not likely. Still, it was argued in that this does not exclude the possibility of observing a “weak” non-linearity.
The immediate (naive) tendency is to regard LRT as the outcome of quantum mechanical first order perturbation theory (FOPT). In fact the regimes of validity of FOPT and of LRT do not coincide. On the one hand we have the adiabatic regime where FOPT is valid as a leading order description, but not for response calculation. On the other hand, the validity of Eq.(37) goes well beyond FOPT. This leads to the (correct) identification of what we call the “perturbative regime”. The border of this regime is determined by the energy scale $`\mathrm{\Delta }_b`$, while $`\mathrm{\Delta }`$ is not involved. Outside of the perturbative regime we cannot trust the LRT formula. However, as we further explain below, the fact that Eq.(37) is not valid in the non-perturbative regime, does not imply that it fails there.
We stress again that one should distinguish between “non-perturbative response” and “non-linear response”. These are not synonyms. As we explain in the next paragraph, the adiabatic regime is “perturbative” but “non-linear”, while the semiclassical limit is “non-perturbative” but “linear”.
In the adiabatic regime, FOPT implies zero probability to make a transitions to other levels. Therefore, to the extent that we can trust the adiabatic approximation, all probability remains concentrated on the initial level. Thus, in the adiabatic regime, Eq.(37) is not a valid formula: It is essential to use higher orders of perturbation theory, and possibly non-perturbative corrections (Landau-Zener ), in order to calculate the response. Still, FOPT provides a meaningful leading order description of the dynamics (i.e. having no transitions), and therefore we do not regard the adiabatic non-linear regime as “non-perturbative”.
In the non-perturbative regime the evolution of $`P_t(n|m)`$ cannot be extracted from perturbation theory: not in leading order, neither in any order. Still it does not necessarily imply a non-linear response. On the contrary: The semiclassical limit is contained in the deep non-perturbative regime . There, the LRT formula Eq.(37) is in fact valid. But its validity is not a consequence of perturbation theory, but rather the consequence of quantal-classical correspondence (QCC).
In the next subsection we will present a classical derivation of the general LRT expression (37). In Subsection 6.3 we derive it using first order perturbation theory (FOPT). In Subsection 6.5 we derive the corresponding FOPT expression for the survival probability.
### 6.2 Classical LRT derivation for $`\delta E(t)`$
The classical evolution of $`E(t)=(Q(t),P(t))`$ can be derived from Hamilton equations. Namely,
$$\frac{\text{d}E(t)}{\text{d}t}=[,]_{\text{PB}}+\frac{}{t}=\epsilon \dot{f}(t)(t),$$
(40)
where $`[]_{\mathrm{PB}}`$ indicates the Poisson Brackets. Integration of Eq. (40) leads to
$$E(t)E(0)=\epsilon _0^t(t^{})\dot{f}(t^{})𝑑t^{}.$$
(41)
Taking a micro-canonical average over initial conditions we obtain the following expression for the variance
$$\delta E^2(t)=\epsilon ^2_0^tC(t^{}t^{\prime \prime })\dot{f}(t^{})\dot{f}(t^{\prime \prime })𝑑t^{}𝑑t^{\prime \prime },$$
(42)
which can be re-written in the form of (37).
One extreme special case of Eq.(37) is the sudden limit for which $`f(t)`$ is a step function. Such evolution is equivalent to the LDoS studies of Section 5. In this case $`F_t(\omega )=1`$, and accordingly
$$\delta E_{\text{cl}}=\epsilon \times \sqrt{C(0)}\text{[“sudden” case]}.$$
(43)
Another extreme special case is the response for persistent (either linear or periodic) driving of a system with an extremely short correlation time. In such case $`F_t(\omega )`$ becomes a narrow function with a weight that grows linearly in time. For linear driving ($`f(t)=t`$) we get $`F_t(\omega )=t\times 2\pi \delta (\omega )`$. This implies diffusive behavior:
$$\delta E(t)=\sqrt{2D_Et}\text{[“Kubo” case]},$$
(44)
where $`D_Eϵ^2`$ is the diffusion coefficient. The expression for $`D_E`$ as an integral over the correlation function is known in the corresponding literature either as Kubo formula, or as Einstein relation, and is the corner stone of the Fluctuation-Dissipation relation.
### 6.3 Quantum LRT derivation for $`\delta E(t)`$
The quantum mechanical derivation looks like an exercise in first order perturbation theory. In fact a proper derivation that extends and clarifies the regime where the result is applicable requires infinite order. If we want to keep a complete analogy with the classical derivation we should work in the adiabatic basis . (For a brief derivation see Appendix D of ).
In the following presentation we work in a ”fixed basis” and assume $`f(t)=f(0)=0`$. We use the standard textbook FOPT expression for the transition probability from an initial state $`m`$ to any other state $`n`$. This is followed by integration by parts. Namely,
$`P_t(n|m)`$ $`=`$ $`{\displaystyle \frac{\epsilon ^2}{\mathrm{}^2}}|𝑩_{nm}|^2\left|{\displaystyle \underset{0}{\overset{t}{}}}\text{d}t^{}f(t^{})\text{e}^{i(E_nE_m)t^{}/\mathrm{}}\right|^2`$ (45)
$`=`$ $`{\displaystyle \frac{\epsilon ^2}{\mathrm{}^2}}|𝑩_{nm}|^2{\displaystyle \frac{\stackrel{~}{F}_t(\omega _{nm})}{(\omega _{nm})^2}},`$
where $`\omega _{nm}=(E_nE_m)/\mathrm{}`$. Now we calculate the variance and use Eq. (22) so as to get
$`\delta E^2(t)`$ $`=`$ $`{\displaystyle \underset{n}{}}P_t(n|m)(E_nE_m)^2`$ (46)
$`=`$ $`\epsilon ^2{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{\text{d}\omega }{2\pi }}\stackrel{~}{C}(\omega )\stackrel{~}{F}_t(\omega ).`$
### 6.4 Restricted QCC
The FOPT result for $`\delta E(t)`$ is *exactly* the same as the classical expression Eq. (37). It is important to realize that there is no $`\mathrm{}`$-dependence in the above formula. This correspondence does not hold for the higher $`k`$moments of the energy distribution. If we use the above FOPT procedure we get that the latter scale as $`\mathrm{}^{k2}`$.
We call the quantum-classical correspondence for the second moment ”restricted QCC”. It is a very robust correspondence . This should be contrasted with ”detailed QCC” that applies only in the semiclassical regime where $`P_t(n|m)`$ can be approximated by a classical result (and not by a perturbative result).
### 6.5 Quantum LRT derivation for $`𝒫(t)`$
With the validity of FOPT assumed we can also calculate the time-decay of the survival probability $`𝒫(t)`$. From Eq. (45) we get:
$$p(t)\underset{n(n_0)}{}P_t(n|m)=\epsilon ^2_{\mathrm{}}^{\mathrm{}}\frac{d\omega }{2\pi }\stackrel{~}{C}(\omega )\frac{\stackrel{~}{F}_t(\omega )}{(\mathrm{}\omega )^2}.$$
(47)
Assuming that $`𝒫(t)=1p(t)`$ can be extrapolated in a ”stochastic” fashion we get Eq. (38). Another way to write the final formula is as follows:
$$𝒫(t)=\mathrm{exp}\left[\frac{1}{\mathrm{}^2}_0^t_0^tC(t^{}t^{\prime \prime })\delta x(t^{})\delta x(t^{\prime \prime })𝑑t^{}𝑑t^{\prime \prime }\right].$$
(48)
For constant perturbation (wavepacket dynamics) and assuming long times we obtain the Wigner decay,
$$𝒫(t)=\mathrm{exp}\left[\left(\frac{ϵ}{\mathrm{}}\right)^2\stackrel{~}{C}(\omega =0)\times t\right],$$
(49)
which can be regarded as a special case of Fermi-Golden-Rule.
### 6.6 Note on $`𝒫(t)`$ for a time reversal scenario
The ”LRT formula” for $`𝒫(t)`$ in the case of ”driving reversal scenario” is
$$𝒫_{\text{DR}}(t)=\mathrm{exp}\left[\left(\frac{\epsilon }{\mathrm{}}\right)^2_0^T_0^TC(t^{}t^{\prime \prime })f(t^{})f(t^{\prime \prime })𝑑t^{}𝑑t^{\prime \prime }\right],$$
(50)
where we assumed the simplest scenario with $`f(t)=1`$ for $`0<t<(T/2)`$ and $`f(t)=1`$ for $`(T/2)<t<T`$. It is interesting to make a comparison with the analogous result in case of ”time reversal scenario”.
The well known Feynman-Vernon influence functional has the following approximation:
$`F[x_A,x_B]`$ $`=`$ $`\mathrm{\Psi }|U[x_B]^1U[x_A]|\mathrm{\Psi }`$
$`=`$ $`\mathrm{exp}\left[{\displaystyle \frac{1}{2\mathrm{}^2}}{\displaystyle _0^t}{\displaystyle _0^t}C(t^{}t^{\prime \prime })(x_B(t^{})x_A(t^{\prime \prime }))^2𝑑t^{}𝑑t^{\prime \prime }\right].`$
This expression is in fact exact in the case of harmonic bath, and assuming thermal averaging over the initial state. Otherwise it should be regarded as an extrapolated version of leading order perturbation theory (as obtained in the interaction picture). What people call nowadays ”fidelity” or “Loschmidt echo” is in fact a special case of the above expression which is defined by setting $`t=T/2`$ and $`x_A=\epsilon /2`$ while $`x_B=\epsilon /2`$. Thus
$`𝒫_{\text{TR}}(t)`$ $`=`$ $`|F[x_A,x_B]|^2`$
$`=`$ $`\mathrm{exp}\left[\left({\displaystyle \frac{\epsilon }{\mathrm{}}}\right)^2{\displaystyle _0^{T/2}}{\displaystyle _0^{T/2}}C(t^{}t^{\prime \prime })𝑑t^{}𝑑t^{\prime \prime }\right].`$
Assuming a very short correlation time one obtains
$$𝒫_{\text{TR}}(T)=\mathrm{exp}\left[\frac{1}{2}\left(\frac{ϵ}{\mathrm{}}\right)^2\stackrel{~}{C}(\omega =0)\times T\right],$$
(53)
which again can be regarded as a special variation of the Fermi-Golden-Rule (but note the pre-factor $`1/2`$).
### 6.7 The survival probability and the LDoS
For constant perturbation it is useful to remember that $`𝒫(t)`$ LDoS as follows:
$`𝒫(t)`$ $``$ $`\left|n(x_0)|e^{i(x)t/\mathrm{}}|n(x_0)\right|^2`$ (54)
$`=`$ $`|{\displaystyle \underset{m}{}}e^{iE_m(x)t/\mathrm{}}|m(x)|n(x_0)|^2|^2`$
$`=`$ $`|{\displaystyle _{\mathrm{}}^{\mathrm{}}}P(E|m)\text{e}^{iEt/\mathrm{}}dE|^2.`$
This implies that a Wigner decay is associated with a Lorentzian approximation for the LDoS. In the non-perturbative regime the LDoS is not a Lorentzian, and therefore one should not expect an exponential. In the semiclassical regime the LDoS shows system specific features and therefore the decay of $`𝒫(t)`$ becomes non-universal.
## 7 Wavepacket Dynamics for Constant Perturbation
The first evolution scheme that we are investigating here is the so-called *wavepacket dynamics*. The classical picture is quite clear : The initial preparation is assumed to be a micro-canonical distribution that is supported by the energy surface $`_0(Q,P)=E(0)`$. Taking $``$ to be a generator for the classical dynamics, the phase-space distribution spreads away from the initial surface for $`t>0`$. ‘Points’ of the evolving distribution move upon the energy surfaces of $`(Q,P)`$. Thus, the energy $`E(t)=_0(Q(t),P(t))`$ of the evolving distributions spreads with time. Using the LRT formula Eq.(39) for rectangular pulse $`f(t^{})=1`$ for $`0<t^{}<t`$ we get
$$\stackrel{~}{F}_t(\omega )=\left|1\text{e}^{i\omega t}\right|^2=(\omega t)^2\mathrm{sinc}\left(\frac{\omega t}{2}\right),$$
(55)
and hence
$$\delta E_{\text{cl}}(t)=\epsilon \times \sqrt{2(C(0)C(t))}.$$
(56)
For short times $`t\tau _{\text{cl}}`$ we can expand the correlation function as $`C(t)C(0)\frac{1}{2}C^{\prime \prime }(0)t^2`$, leading to a ballistic evolution. Then, for $`t\tau _{\text{cl}}`$, due to ergodicity, a ‘steady-state distribution’ appears, where the evolving ‘points’ occupy an ‘energy shell’ in phase-space. The thickness of this energy shell equals $`\delta E_{\text{cl}}`$. Thus, we have a crossover from ballistic energy spreading to saturation:
$$\delta E(t)\{\begin{array}{ccc}\sqrt{2}(\delta E_{\text{cl}}/\tau _{\text{cl}})t\hfill & \text{for}& \hfill t<\tau _{\text{cl}}\\ \sqrt{2}\delta E_{\text{cl}}\hfill & \text{for}& \hfill t>\tau _{\text{cl}}\end{array}.$$
(57)
Figure 9 shows the classical energy spreading (heavy dashed line) for the 2DW model. In agreement with Eq. (57) we see that $`\delta E_{\text{cl}}(t)`$ is first ballistic and then saturates. The classical dynamics is fully characterized by the two classical parameters $`\tau _{\text{cl}}`$ and $`\delta E_{\text{cl}}`$.
### 7.1 The quantum dynamics
Let us now look at the quantized 2DW model. The quantum mechanical data are reported in Fig. 9 (left panel) where different curves correspond to various perturbation strengths $`\epsilon `$. As in the classical case (heavy dashed-line) we observe an initial ballistic-like spreading followed by saturation. This could lead to the wrong impression that the classical and the quantum spreading are of the same nature. However, this is definitely not the case.
In order to detect the different nature of quantum ballistic-like spreading, one has to inquire measures that are sensitive to the structure of the profile, such as the core-width $`\delta E_{\text{core}}(t)`$. In Fig. 10 we present our numerical data for the 2DW model. If the spreading were of a classical type, it would imply that the spreading profile is characterized by a single energy scale. In such a case we would expect that $`\delta E_{\text{core}}(t)\delta E(t)`$. Indeed this is the case for $`\epsilon >\epsilon _{\text{prt}}`$ with the exclusion of very short times: The larger $`\epsilon `$ is the shorter the quantal transient becomes. In the perturbative regimes, in contrast to the semiclassical regime, we have a separation of energy scales $`\delta E_{\text{core}}(t)\delta E(t)`$. In the perturbative regimes $`\delta E(t)`$ is determined by the tails, and it is not sensitive to the size of the ‘core’ region.
Using the LRT formula for $`𝒫(t)`$ we get, for short times ($`t\tau _{\text{cl}}`$) during the ballistic-like stage
$$𝒫(t)=\mathrm{exp}\left(C(\tau =0)\times \left(\frac{ϵt}{\mathrm{}}\right)^2\right),$$
(58)
while for long times ($`t\tau _{\text{cl}}`$) we have the FGR decay of Eq.(49). Can we trust these expressions? Obviously FOPT can be trusted as long as $`𝒫(t)1`$. This can be converted into an inequality $`t<t_{\text{prt}}`$ where
$$t_{\text{prt}}=\left(\frac{\epsilon _{\text{prt}}}{\epsilon }\right)^{\nu =1,2}\tau _{\text{cl}}.$$
(59)
The power $`\nu =1`$ applies to the non-perturbative regime where the breakdown of $`𝒫(t)`$ happens to be before $`\tau _{\text{cl}}`$. The power $`\nu =2`$ applies to the perturbative regime where the breakdown of $`𝒫(t)`$ happens after $`\tau _{\mathrm{cl}}`$ at $`t_{\text{prt}}=\mathrm{}/\mathrm{\Gamma }`$, i.e. after the ballistic-like stage.
The long term behavior of $`𝒫(t)`$ in the non-perturbative regime is not the Wigner decay. It can be obtained by Fourier transform of the LDoS. In the non-perturbative regime the LDoS is characterized by the single energy scale $`\delta E_{\text{cl}}\delta x`$. Hence the decay in this regime is characterized by a semiclassical time scale $`2\pi \mathrm{}/\delta E_{\text{cl}}`$.
### 7.2 The EBRM dynamics
Next we investigate the applicability of the RMT approach to describe wavepacket dynamics and specifically the energy spreading $`\delta E(t)`$. At first glance, we might be tempted to speculate that RMT should be able, at least as far as $`\delta E(t)`$ is concerned, to describe the actual quantum picture. After all, we have seen in Subsection 6.1 that the quantum mechanical LRT formula (46) for the energy spreading involves as its only input the classical power spectrum $`\stackrel{~}{C}(\omega )`$. Thus we would expect that an effective RMT model with the same band-profile would lead to the *same* $`\delta E(t)`$.
However, things are not so trivial. In Figure 9 we show the numerical results for the EBRM model <sup>2</sup><sup>2</sup>2The same qualitative results were found also for the prototype WBRM model, see .. In the standard and in the extended perturbative regimes we observe a good agreement with Eq.(46). This is not surprising as the theoretical prediction was derived via FOPT, where correlations between off-diagonal elements are not important. In this sense the equivalence of the 2DW model and the EBRM model is trivial in these regimes. But as soon as we enter the non-perturbative regime, the spreading $`\delta E(t)`$ shows a qualitatively different behavior from the one predicted by LRT: After an initial ballistic spreading, we observe a premature crossover to a diffusive behavior
$$\delta E(t)=\sqrt{2D_Et}.$$
(60)
The origin of the diffusive behavior can be understood in the following way. Up to time $`t_{\text{prt}}`$ the spreading $`\delta E(t)`$ is described accurately by the FOPT result (46). At $`tt_{\text{prt}}`$ the evolving distribution becomes as wide as the bandwidth, and we have $`\delta E_{\text{core}}\delta E\mathrm{\Delta }_b`$ rather than $`\delta E_{\text{core}}\delta E\mathrm{\Delta }_b`$. We recall that in the non-perturbative regime FOPT is subjected to a breakdown before reaching saturation. The following simple heuristic picture turns out to be correct. Namely, once the mechanism for ballistic-like spreading disappears, a stochastic-like behavior takes its place. The stochastic energy spreading is similar to a random-walk process where the step size is of the order $`\mathrm{\Delta }_b`$, with transient time $`t_{\text{prt}}`$. Therefore we have a diffusive behavior $`\delta E(t)^2=2D_\text{E}t`$ with
$$D_\text{E}=C\mathrm{\Delta }_b^2/t_{\text{prt}}=C\mathrm{\Delta }^2b^{5/2}\epsilon \sigma /\mathrm{}\mathrm{}$$
(61)
where $`C`$ is some numerical pre-factor. This diffusion is not of classical nature, since in the $`\mathrm{}0`$ limit we get $`D_E0`$. The diffusion can go on until the energy spreading profile ergodically covers the whole energy shell and saturates to a classical-like steady state distribution. The time $`t_{\text{erg}}`$ for which we get ergodization is characterized by the condition $`(D_\text{E}t)^{1/2}<\delta E_{\text{cl}}`$, leading to
$$t_{\text{erg}}=b^{3/2}\mathrm{}\epsilon \sigma /\mathrm{\Delta }^2\mathrm{\hspace{0.17em}\hspace{0.17em}1}/\mathrm{}.$$
(62)
For completeness we note that for $`\epsilon >\epsilon _{loc}`$ there is no ergodization but rather dynamical (”Anderson” type) localization. Hence, in the latter case, $`t_{\text{erg}}`$ is replaced by the break-time $`t_{\text{brk}}`$. The various regimes and time scales are illustrated by the diagram presented in Fig. 11.
## 8 Driving Reversal Scenario
A thorough understanding of the one-period driving reversal scenario is both important within itself, and for constituting a bridge towards a theory dealing with the response to periodic driving . In the following subsection we present our results for the prototype WBRM model, while in Subsection 8.2 we consider the 2DW model and compare it to the corresponding EBRM model. The EBRM is better for the purpose of making comparisons with the 2DW, while the WBRM is better for the sake of quantitative analysis (the ”physics” of the EBRM and the WBRM models is, of course, the same).
The quantities that monopolize our interest are the energy spreading $`\delta E(t)`$ and the survival probability $`𝒫(t)`$. In Figs. 12 and 13 we present representative plots. From a large collection of such data that collectively span a very wide range of parameters, we extract results for $`\delta E(T)`$, for $`𝒫(T)`$, and for the corresponding compensation times. These are presented in Figs. 12,13,14,15,16 and 17.
### 8.1 Driving Reversal Scenario: RMT Case
#### 8.1.1 LRT for the energy spreading
Assuming that the driving reversal happens at $`t=T/2`$, the spectral content $`\stackrel{~}{F}_t(\omega )`$ for $`T/2<t<T`$ is
$$\stackrel{~}{F}_t(\omega )=\left|12\text{e}^{i\omega T/2}+\text{e}^{i\omega t}\right|^2.$$
(63)
Inserting Eq. (63) into Eq. (37) we get
$$\delta E(t)=\epsilon \times \sqrt{6C(0)+2C(t)4C(\frac{T}{2})4C(t\frac{T}{2})}.$$
(64)
For the WBRM model we can substitute in Eq. (64) the exact expression Eq. (26) for the correlation function, and get
$$\delta E(t)=2\epsilon \sigma \times \sqrt{3b+b\mathrm{sinc}(\frac{t}{t_{\mathrm{cl}}})2b\mathrm{sinc}(\frac{T}{2\tau _{\mathrm{cl}}})2b\mathrm{sinc}(\frac{t\frac{T}{2}}{\tau _{\mathrm{cl}}})}.$$
(65)
We can also find the compensation time $`t_r^E`$ by minimizing Eq. (64) with respect to $`t`$. For the WBRM model we have
$$\frac{2\mathrm{cos}\left[\frac{T/2t}{\tau _{\text{cl}}}\right]}{\tau _{\mathrm{cl}}(T/2t)}+\frac{2\mathrm{sin}\left[\frac{T/2t}{\tau _{\text{cl}}}\right]}{(T/2t)^2}+\frac{\mathrm{cos}\left[\frac{t}{\tau _{\text{cl}}}\right]}{t\tau _{\mathrm{cl}}}=\frac{1}{t^2}\mathrm{sin}\left[\frac{t}{\tau _{\text{cl}}}\right],$$
(66)
which can be solved numerically to get $`t_r^E`$.
The spreading width at the end of the period is
$$\delta E(T)=\epsilon \times \sqrt{6C(0)+2C(T)8C(\frac{T}{2}))}.$$
(67)
It is important to realize that the dimensional parameters in this LRT analysis are determined by the time scale $`\tau _{\text{cl}}`$ and by the energy scale $`\delta E_{\text{cl}}`$. This means that we have a scaling relation (using units such that $`\sigma =\mathrm{\Delta }=\mathrm{}=1`$)
$$\frac{\delta E(T)}{\sqrt{b}\epsilon }=h_{\text{LRT}}^E\left(bT\right).$$
(68)
Deviation from this scaling relation implies a non-perturbative effect that goes beyond LRT.
The LRT scaling is verified nicely by our numerical data (see upper panels of Fig. 14 ). The values of perturbation strength for which the LRT results are applicable correspond to $`\epsilon <\epsilon _{\text{prt}}`$. In the same Figure we also plot the whole analytical expression (67) for the spreading $`\delta E(T)`$. Similarly in Fig. 14 (lower panels) we present our results for the compensation time $`t_r^E`$. All the data fall one on top of the other once we rescale them. It is important to realize that the LRT scaling relation implies that the compensation time $`t_r^E`$ is *independent* of the perturbation strength $`\epsilon `$. It is determined only by the *classical* correlation time $`\tau _{\text{cl}}`$. In the same figure we also present the resulting analytical result (heavy-dashed line) which had been obtained via Eq.(66). An excellent agreement with our data is evident.
#### 8.1.2 Energy spreading in the non-perturbative regime
We turn now to discuss the dynamics in the non-perturbative regime, which is our main interest. In the absence of driving reversal (see Subsection 7.2) we obtain diffusion ($`\delta E(t)\sqrt{t}`$) for $`t>t_{\text{prt}}`$, where
$$t_{\text{prt}}=\mathrm{}/(\sqrt{b}\sigma \epsilon ).$$
(69)
If $`(T/2)<t_{\text{prt}}`$, this non-perturbative diffusion does not have a chance to develop, and therefore we can still trust Eq. (64). So the interesting case is $`(T/2)>t_{\text{prt}}`$, which means large enough $`\epsilon `$. In the following analysis we distinguish between two stages in the non-perturbative diffusion process. The first stage ($`t_{\text{prt}}<t<t_{\text{sdn}}`$) is reversible, while the second stage ($`t>t_{\text{sdn}}`$) is irreversible. For much longer time scales we have recurrences or localization, which are not the issue of this paper. The new time scale ($`t_{\text{sdn}}`$) did not appear in our ”wavepacket dynamics” study, because it can be detected only by time driving reversal experiment.
The determination of the time scale $`t_{\text{sdn}}`$ is as follows. The diffusion coefficient is $`D_\text{E}=\mathrm{\Delta }^2b^{\text{5/2}}\sigma \epsilon /\mathrm{}`$ up to a numerical pre-factor. The diffusion law is $`\delta E^2(t)=D_\text{E}t`$. The diffusion process is reversible as long as $`𝐄`$ does not affect the relative phases of the participating energy levels. This means that the condition for reversibility is $`(\delta E(t)\times t)/\mathrm{}1`$. The latter inequality can be written as $`t<t_{\text{sdn}}`$, where
$$t_{\text{sdn}}=\left(\frac{\mathrm{}^2}{D_\text{E}}\right)^{1/3}=\left(\frac{\mathrm{}^3}{\mathrm{\Delta }^2b^{5/2}\sigma \epsilon }\right)^{1/3}.$$
(70)
It is extremely important to realize that without reversing the driving, the presence or the absence of $`𝐄`$ in the Hamiltonian cannot be detected. It is only by driving reversal that we can easily determine (as in the upper panels of Fig.12) whether the diffusion process is reversible or irreversible.
The dimensional parameters in this analysis are naturally the time scale $`t_{\text{sdn}}`$ and the resolved energy scale $`\mathrm{}/T`$. Therefore we expect to have instead of the LRT scaling, a different ”non-perturbative” scaling relation. Namely, $`\delta E(T)/(\mathrm{}/T)`$ should be related by a scaling function to $`T/t_{\text{sdn}}`$. Equivalently (using units such that $`\sigma =\mathrm{\Delta }=\mathrm{}=1`$) it can be written as
$$\frac{\delta E(T)}{b^{5/6}\epsilon ^{1/3}}=h_{\text{nprt}}^E\left(b^{5/6}\epsilon ^{1/3}T\right).$$
(71)
Obviously the non-perturbative scaling with respect to $`ϵ^{1/3}`$ goes beyond any implications of perturbation theory. It is well verified by our numerical data (see upper right panel of Fig. 14). The values of perturbation strength for which this scaling applies correspond to $`\epsilon >\epsilon _{\text{prt}}`$. The existence of the $`t_{\text{sdn}}`$ scaling can also be verified in the lower right panel of Fig. 14, where we show that $`t_r/T`$ is by a scaling function related to $`b^{5/6}\epsilon ^{1/3}T`$.
#### 8.1.3 Decay of $`𝒫(t)`$ in the FOPT regime
We can substitute Eq. (63) for the spectral content $`\stackrel{~}{F}_t(\omega )`$ of the driving into the LRT formula Eq. (47), and come out with the following expression for the survival probability at the end of the period $`t=T`$
$$𝒫(T)\mathrm{exp}(\epsilon ^2T^4b^3).$$
(72)
This is a super-Gaussian decay, which is quite different from the standard Gaussian decay Eq. (58) or any other results on reversibility that appear in literature . We have verified that this expression is valid in the FOPT regime. See Figure 15a.
For the WBRM model, we get the following expression for $`p(t)`$ after substituting the spectral content $`\stackrel{~}{F}_t(\omega )`$ given from Eq. (63)
$$p(t)=\frac{(ϵ\sigma )^2}{\mathrm{\Delta }\mathrm{}}\times _{\omega _{\mathrm{cl}}}^{\omega _{\mathrm{cl}}}𝑑\omega \frac{64[\mathrm{cos}(\frac{\omega T}{2})+\mathrm{cos}(\omega (\frac{T}{2}t))]+2cos(\omega t)}{\omega ^2}.$$
(73)
The corresponding compensation time $`t_r^P`$ can be found after minimizing the above expression (corresponding to the maximization of $`𝒫(t)=1p(t)`$) with respect to time $`t`$. This results in the following equation
$$\text{ si}\left(\omega _{cl}t\right)=2\text{ si}\left(\omega _{cl}\left(t\frac{T}{2}\right)\right),$$
(74)
which has to be solved numerically in order to evaluate $`t_r^P`$. Above $`\text{si}(x)=_0^x\frac{\mathrm{sin}x}{x}`$. Our numerical data are reported in Fig. 16 together with the theoretical prediction (74).
#### 8.1.4 Decay of $`𝒫(t)`$ in the Wigner regime
We now turn to discuss $`𝒫(t)`$ in the ”Wigner regime”. By this we mean $`\epsilon _c<\epsilon <\epsilon _{\text{prt}}`$. This distinction does not appear in the $`\delta E(t)`$ analysis. The time evolution of $`\delta E(t)`$ is dominated by the tails of the distribution and does not affect the ”core” region. Therefore $`\delta E(t)`$ also agreed with LRT outside of the FOPT regime in the whole (extended) perturbative regime. But this is not the case with $`𝒫(t)`$, which is mainly influenced by the ”core” dynamics. As a result in the ”Wigner regime” we get different behavior compared with the FOPT regime.
We look at the survival probability $`𝒫(T)`$ at the end of the driving period. In the Wigner regime, instead of the LRT-implied super-Gaussian decay, we find a Wigner-like decay:
$$𝒫(T)\text{e}^{\mathrm{\Gamma }(\epsilon )T},$$
(75)
where $`\mathrm{\Gamma }\epsilon ^2/\mathrm{\Delta }`$. In Figure 15b we present our numerical results for various perturbation strengths in this regime. A nice overlap is observed once we rescale the time axis as $`\epsilon ^2\times T`$. We would like to emphasize once more that both in the standard and in the extended perturbative regimes the scaling law involves the perturbation strength $`\epsilon `$. This should be contrasted with the LRT scaling of $`\delta E(t)`$.
What about the compensation time $`t_r^P`$? A reasonable assumption is that it will exhibit a different scaling in the FOPT regime and in the Wigner regime (as is the case of $`𝒫(t)`$). Namely, in the FOPT regime we would expect ”LRT scaling” with $`\tau _{\text{cl}}`$, while in the Wigner regime we would expect ”Wigner scaling” with $`t_{\text{prt}}=\mathrm{}/\mathrm{\Gamma }`$. The latter is of non-perturbative nature and reflects the ”core” dynamics. To our surprise we find that this is not the case. Our numerical data presented in Fig. 16 show beyond any doubt that the ”LRT scaling” applies within the whole (extended) perturbative regime, as in the case of $`t_r^E`$, thus not invoking the perturbation strength $`\epsilon `$. We see that the FOPT expression (74) for $`t_r^P`$ shown as a heavy-dashed line describes the numerical findings.
We conclude that the compensation time $`t_r`$ is mainly related to the dynamics of the tails, and hence can be deduced from the LRT analysis.
#### 8.1.5 Decay of $`𝒫(t)`$ in the non-perturbative regime
Let us now turn to the non-perturbative regime (see Fig. 15c). As in the case of the spreading kernel $`\delta E(T)`$, the decay of $`𝒫(T)`$ is no longer captured by perturbation theory. Instead, we observe the same non-universal scaling with respect to $`\epsilon ^{1/3}\times T`$ as in the case of $`\delta E(T)`$.
$$𝒫(T)=h_{\text{nprt}}^P\left(b^{5/6}\epsilon ^{1/3}T\right).$$
(76)
The reason is that in the non-perturbative regime the two energy scales $`\mathrm{\Gamma }`$ and $`\mathrm{\Delta }_b`$, which were responsible for the difference between $`𝒫(T)`$ and $`\delta E(T)`$, lose their meaning. As a consequence, the spreading process involves only one time scale and the behavior of both $`𝒫(T)`$ and $`\delta E(T)`$ becomes similar, leading to the same scaling behavior.
### 8.2 Driving Reversal Scenario: 2DW Case
In the representative simulations of the 2DW model in Fig. 12 (upper left panel) we see that the spreading $`\delta E(t)`$ for $`T=0.48`$ and various perturbation strengths $`\epsilon `$ follows the LRT predictions very well. Fig. 12 (lower left panel) shows that the agreement with the LRT is observed for any value of the period $`T`$. This stands in clear contrast to the EBRM model shown in Fig. 12 (right panels).
The agreement with LRT in the non-perturbative regime, as in the case of wavepacket dynamics, reflects detailed QCC. We recall that ”to get into the non-perturbative regime” and ”to make $`\mathrm{}`$ small” means the same. All our simulations are done in a regime where LRT can be trusted at the classical ($`=`$non-perturbative) limit. It is only for RMT models that we observe a breakdown of LRT in the non-perturbative regime.
What about $`𝒫(T)`$? This quantity has no classical analogue. Therefore QCC considerations are not applicable. Also LDoS considerations cannot help here. The one-to-one correspondence between the LDoS and the survival probability applies to the simple wavepacket dynamics scenario (constant perturbation).
It is practically impossible to make a quantitative analysis of $`𝒫(T)`$ in the case of a real model because the band-profile is very structured and there are severe numerical limitations. Rather, what we can easily do is to compare the 2DW with the corresponding EBRM. Any difference between the two constitutes an indication for a non-perturbative effect. Representative simulations are presented in Figure 13.
In Figure 17 we show the dependence of the compensation time $`t_r^P`$ on $`T`$ for the EBRM model. We see very nice scaling behavior that indicates that our numerics (as far as $`𝒫(T)`$ is concerned!) is limited to the perturbative regime. We emphasize again that the physics of $`𝒫(t)`$ is very different from the physics of $`\delta E(t)`$. Therefore, this finding by itself should not be regarded as very surprising. A sharp crossover to a non-perturbative behavior can be expected for a “sharp” band-profile only (which is the WBRM and not the EBRM– see Fig.16).
Now we switch from the EBRM model to the 2DW model. Do we see any deviation from LRT scaling? The answer from Fig. 17 is clearly yes, as reflected by the $`ϵ`$ dependence of the curve. The effect is small, but ”it is there”. It indicates that the “body” of the probability distribution, in the case of the 2DW dynamics, does not evolve the same way as in the EBRM case. Indeed we know that the main part of the distribution evolves faster (in a ballistic fashion rather than diffusively), and therefore we observe lower values of $`t_r^P`$.
Assuming that the decay of $`𝒫(T)`$ is given by the exponential law, we extract the corresponding decay rates $`\gamma `$. It should be clear that the fitting is done merely in order to extract a numeric measure for the behavior of the decay. We would not like to suggest that the decay looks strictly exponential. The results are reported in Figure 18. We find that for $`\epsilon <\epsilon _{\text{prt}}`$, the decay rate $`\gamma (\epsilon )\epsilon ^2`$, as expected by Wigner’s theory, while for $`\epsilon >\epsilon _{\text{prt}}`$ we find that $`\gamma \epsilon `$. This linear dependence on $`\epsilon `$ is essentially the same as in the corresponding wavepacket dynamics scenario. There it is clearly associated with the width $`\delta E_{\text{cl}}\epsilon `$ of the LDoS.
As far as $`\gamma `$ is concerned the behavior of 2DW and the EBRM models are the same, and there is an indication of the crossover from the perturbative to the non-perturbative regime, as implied (in a non-rigorous fashion) by the LDoS theory. It is $`t_r^P`$ rather than $`\gamma `$ that exhibits sensitivity to the nature of the dynamics. This is because $`t_r^P`$ is sensitive to the evolution of the main part of the distribution. We already had made this observation on the basis of the analysis of the WBRM model (see previous Subsection 8.1). Here we see another consequence of this observation.
## 9 Conclusions
There is a hierarchy of challenges in the study of quantum dynamics. The simple way to explain this hierarchy is as follows: Let us assume that there are two Hamiltonians, $`_1`$ and $`_2`$, that differ slightly from each other. Let us then quantify the difference by a parameter $`\epsilon `$. Let us distinguish between a FOPT regime, Wigner regime, and non-perturbative (semicircle or semiclassical) regime according to the line shape of the LDoS. Do we have enough information to say something about the dynamics?
In the conventional wavepacket dynamics, one Hamiltonian is used for preparation and for measurement, while the other for propagation. It is well known that the Fourier transform of the LDoS gives the survival amplitude and hence $`𝒫(t)`$. But what about other features of the dynamics. What about the energy spreading $`\delta E(t)`$ for example? It turns out that the answer requires more than just knowing the LDoS. In particular we observe that in the non-perturbative regime physical models differ from the corresponding RMT model. In the former case we have ballistic spreading while in the latter we have diffusion.
Is there any new ingredient in the study of driving reversal dynamics? Maybe it is just a variation on conventional wavepacket dynamics? The answer turns out to be interesting. There is a new ingredient in the analysis. This becomes very clear in the RMT analysis where we find a new time scale that distinguishes between a stage of ”reversible diffusion” and a stage of ”irreversible diffusion”. This time scale ($`t_{\text{sdn}}`$) can only be probed in a driving reversal experiment. It is absent in the study of conventional wavepacket dynamics.
Things become more interesting, and even surprising, once we get into details. Let us summarize our main findings. We start with the conventional wavepacket dynamics, and then turn to the driving reversal scenario.
The main observations regarding wavepacket dynamics are summarized by the diagrams in Fig. 11. We always have an initial ballistic-like stage which is implied by FOPT. During this stage the first order (in-band) tails of the energy distribution grow like $`t^2`$. We call this behavior ”ballistic-like” because the second moment $`\delta E(t)`$ grows like $`t^2`$. It is not a genuine ballistic behavior because the $`r`$th moment does not grow like $`t^r`$ but rather all the moments of this FOPT distribution grow like $`t^2`$.
The bandwidth $`\mathrm{\Delta }_b`$ is resolved at the time $`\tau _{\text{cl}}`$. In the perturbative regime this happens before the breakdown of perturbation theory, while in the non-perturbative regime the breakdown $`t_{\text{prt}}`$ happens before $`\tau _{\text{cl}}`$. As a result, in the non-perturbative regime we can get a non-trivial spreading behavior which turns out to be ”ballistic” or ”diffusive”, depending on whether the system has a classical limit or is being RMT modeled.
Once we consider a driving reversal scenario, it turns out to be important to mark the time $`t_{\text{sdn}}`$ when the energy distribution is resolved. The question is ill-defined in the perturbative regime because there the energy distribution is characterized by two energy scales (the ”bandwidth” and the much smaller ”core width”). But the question is well-defined in the non-perturbative regime where the distribution is characterized by one energy scale. It is not difficult to realize that for ballistic behavior $`t_{\text{sdn}}\tau _{\text{cl}}`$ which is also the classical ergodic time. But for diffusion we get separation of time scales $`t_{\text{prt}}t_{\text{sdn}}\tau _{\text{cl}}`$. Thus we conclude that the diffusion has two stages: One is reversible while the other is irreversible.
But the second moment does not fully characterize the dynamics. In the other extreme we have the survival probability. Whereas $`\delta E(t)`$ is dominated by the tails, $`𝒫(t)`$ is dominated by the ”core” of the distribution. Therefore it becomes essential to distinguish between the FOPT regime where the ”core” is just one level, and the rest of the perturbative regime (the ”Wigner” regime) where the core is large (but still smaller compared with the bandwidth).
The main findings regarding the driving reversal scenario are summarized by the following table:
| Regime | perturbation strength | $`𝒫\left(T\right)`$ behavior | $`t_r`$ behavior | $`\delta E\left(T\right)`$ behavior |
| --- | --- | --- | --- | --- |
| 1st order perturbative | $`\epsilon <\epsilon _c`$ | LRT | LRT | LRT |
| | | (super-Gaussian) | | (ballistic-like) |
| extended perturbative | $`\epsilon _c<\epsilon <\epsilon _{\text{prt}}`$ | Wigner | LRT(!) | LRT |
| (”Wigner”) | | (Exponential) | | (ballistic-like) |
| non-perturbative | $`\epsilon >\epsilon _{\text{prt}}`$ | non-perturbative | non-perturbative | non-perturbative |
| | | (non-universal) | (non-universal) | (diffusive/ballistic) |
| for the WBRM we have diffusion while for the 2DW model we have ballistic behaviour as implied by classical LRT | | | | |
| --- | --- | --- | --- | --- |
As expected we find that $`𝒫(T)`$ obeys FOPT behavior in the FOPT regime, which turns out to be super-Gaussian decay. In the Wigner regime $`\delta E(T)`$ still obeys LRT because the tails obey FOPT, while the non-perturbative core barely affects the second moment. But in contrast to that $`𝒫(T)`$ is sensitive to the core, and therefore we find Wigner (exponential) decay rather than FOPT (super-Gaussian) behavior. However, when we look more carefully at the whole $`𝒫(t)`$ curve, we find that this is not the whole story. We can characterize $`𝒫(t)`$ by the compensation time $`t_r`$. It turns out that $`t_r`$ is sensitive to the nature of the dynamics. Consequently it obeys ”LRT scaling” rather than ”Wigner scaling”. This has further consequences that are related to quantal-classical correspondence. Just by looking at $`𝒫(T)`$ we cannot tell whether we look at the ”real simulation” or on its RMT modeling. But looking on $`t_r`$ we can find a difference. It turns out that in the physical model $`t_r`$ exhibits $`\epsilon `$ dependence, while in the case of RMT modeling $`t_r`$ is independent of $`\epsilon `$ and exhibits ”Wigner scaling”.
Finally we come to the non-perturbative regime. Here we have, in a sense a simpler situation. We have only one energy scale, and hence only one time scale, and therefore $`\delta E(T)`$ and $`𝒫(t)`$ essentially obeys the same scaling. Indeed we have verified that the non-perturbative scaling with $`t_{\text{sdn}}`$ in WBRM simulations is valid for both the second moment and the survival probability.
Finally we would like to emphasize that the notion of ”non-perturbative” behavior should not be confused with ”non-linear” response. In case of quantized models, linear response of the energy spreading $`\delta E(T)`$ is in fact a consequence of non-perturbative behavior. This should be contrasted with the WBRM model, where QCC does not apply, and indeed deviations from the linear response appear once we enter the non-perturbative regime.
The study of irreversibility in a simple driving reversal scenario is an important step towards the understanding of irreversibility and dissipation in general. The analysis of dissipation reduces to the study of energy spreading for time dependent Hamiltonians $`(Q,P;x(t))`$. In generic circumstances the rate of energy absorption is determined by a diffusion-dissipation relation: The long time process of dissipation is determined by the short time diffusion process. The latter is related to the fluctuations $`\stackrel{~}{C}(\omega )`$ via what we call ”LRT formula”. Thus the understanding of short time dynamics is the crucial step in establishing the validity of the fluctuation-dissipation relation.
## Acknowledgments
This research was supported by a grant from the GIF, the German-Israeli Foundation for Scientific Research and Development, and by the Israel Science Foundation (grant No.11/02). |
warning/0506/quant-ph0506169.html | ar5iv | text | # Entanglement and criticality in translational invariant harmonic lattice systems with finite-range interactions
## Abstract
We discuss the relation between entanglement and criticality in translationally invariant harmonic lattice systems with non-randon, finite-range interactions. We show that the criticality of the system as well as validity or break-down of the entanglement area law are solely determined by the analytic properties of the spectral function of the oscillator system, which can easily be computed. In particular for finite-range couplings we find a one-to-one correspondence between an area-law scaling of the bi-partite entanglement and a finite correlation length. This relation is strict in the one-dimensional case and there is strog evidence for the multi-dimensional case. We also discuss generalizations to couplings with infinite range. Finally, to illustrate our results, a specific 1D example with nearest and next-nearest neighbor coupling is analyzed.
LABEL:FirstPage1
Due to the development of powerful tools to quantify entanglement there is a growing interest in the relation between entanglement and criticality in quantum many-body systems. For a variety of spin models it was shown that in the absence of criticality, there is a strict relation between the von-Neumann entropy of a compact sub-set of spins in the ground state and the surface area of the ensemble. E.g. it was shown in GVidal-PRL-2003 ; Korepin ; Calabrese ; Mezzadri that the entanglement in non-critical one-dimensional spin chains approaches a constant value, while it grows logarithmically in the critical case, where the correlation length diverges. Employing field theoretical methods it was argued that in $`d`$ dimensions the entropy grows as a polynomial of power $`d1`$ under non-critical conditions, thus establishing an area theorem. A similar relation was suggested for harmonic lattice models in Bombelli-PRD-1986 and Srednicki-PRL-1993 . Very recently, employing methods of quantum information for Gaussian states, Plenio et al. Plenio gave a derivation of the area theorem for harmonic lattice models with nearest-neighbor couplings. All these findings suggest a general correspondence between entanglement and criticality for non-random potentials. Yet recently special cases have been found for spin chains with Ising-type interactions Duer-PRL-2005 and for harmonic lattices systems Eisert-preprint where the correlation length diverges but the entanglement obeys an area law. Thus the relation between entanglement scaling and criticality remains an open question. It should also be noted that in disorderd systems, i.e. systems with random couplings the relation between entanglement area law and criticality is broken.
In the present paper we show that for harmonic lattice systems with translational invariant, non-random, and finite-range couplings both entanglement scaling and criticality are determined by the analytic properties of the so-called spectral function. For finite-range interactions we find that the properties of the spectral function lead to a one-to-one correspondence between entanglement and criticality. To illustrate our results we discuss a specific one-dimensional example with nearest and next-nearest couplings. Despite the finite range of the coupling this model undergoes a transition from area-law behavior to unbounded logarithmic growth of entanglement.
Let us first consider a one-dimensional system, i.e. a chain of $`N`$ harmonic oscillators described by canonical variables $`(q_i,p_i)`$, $`i=1,2,\mathrm{}N`$. The oscillators are coupled by a translational invariant quadratic Hamiltonian
$$H=\frac{1}{2}\underset{i=1}{\overset{N}{}}p_i^2+\frac{1}{2}\underset{i,j=1}{\overset{N}{}}V_{ij}q_iq_j$$
(1)
where $`V`$ is a real, non-random, symmetric matrix with positive eigenvalues. For a translational invariant system $`V`$ is a Toeplitz matrix, i.e. its elements depend only on the difference of the indexes $`V_{ij}V_k=V_k`$. For a finite system translational invariance implies furthermore periodic boundary conditions $`V_k=V_{Nk}`$. We assume in the following that the interactions are of finite range, i.e that $`V_k0`$ for $`kR`$, where $`R`$ is a finite number independent on $`N`$. As we will show at the end of the paper some generalizations to infinite range couplings are possible. Being positive definite, $`V`$ has a unique positive square root $`V^{1/2}`$ and its inverse $`V^{1/2}`$ which completely determine the ground-state in position (or momentum) representation, $`\mathrm{\Psi }_0(Q)(detV^{1/2})^{1/4}\mathrm{exp}\{\frac{1}{4}Q|V^{1/2}|Q\}`$, Bombelli-PRD-1986 ; Srednicki-PRL-1993 .
The most important characteristics of the oscillator system is the spectrum of $`V`$. Since $`V`$ is a ciculant matrix its eigenvalues can be expressed in terms of complex roots of unity $`z_j=\mathrm{exp}\{i2\pi j/N\}=\mathrm{exp}\{i\theta _j\}`$, $`(j=1,\mathrm{},N)`$:
$`\lambda _j`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{R1}{}}}V_k\left(z_j\right)^k={\displaystyle \frac{1}{2}}V_0+{\displaystyle \frac{1}{2}}{\displaystyle \underset{k=(R1)}{\overset{R1}{}}}V_k\left(\mathrm{e}^{i\frac{2\pi }{N}j}\right)^k.`$ (2)
Eq.(2) together with the positivity of $`V`$ permits the representation $`\lambda _j=h^2(z_j)=|h^2(z_j)|`$ where $`h(z_j)`$ is a polynomial of order $`(R1)/2`$ in $`z_j`$ (assuming for simplicity that $`R`$ is an odd number). Thus $`|h(z)|_{l=1}^{(R1)/2}\left|z\stackrel{~}{z}_l\right|.`$ where $`\stackrel{~}{z}_l\mathrm{exp}\{i\alpha _l\}`$ are the zeroth of $`h(z)`$ which are either real or complex conjugate pairs with $`|\stackrel{~}{z}_l|1`$ Polya . Let $`QR`$ be the number of real zeroth $`\stackrel{~}{z}_r`$ with multiplicity $`m_r\{0,1,\mathrm{}\}`$. Then $`_{r=1}^Q\left|z_j\stackrel{~}{z}_r\right|^{m_r}=_{r=1}^Q\left(22\mathrm{cos}(\theta _j\alpha _r)\right)^{m_r/2},`$ and
$$\lambda (z_j)=\lambda _j=\lambda _0(z_j)\underset{r=1}{\overset{Q}{}}\left(22\mathrm{cos}(\theta _j\alpha _r)\right)^{m_r}.$$
(3)
$`\lambda (z)`$ is the so-called spectral function. $`\lambda _0(z)`$ is called the regular part of $`\lambda `$. It is a polynomial of the complex variable $`z`$ which has zeroth outside the unit circle. As a consequence its inverse $`\lambda _0^1(z)`$ is analytic on and inside the unit circle. $`_{r=1}^Q\left(22\mathrm{cos}(\theta \alpha _r)\right)^{m_r}`$ is called the singular part. If $`\lambda `$ is singular, i.e. if in the thermodynamic limit $`V`$ has eigenvalues arbitrarily close to zero, the total Hamiltonian, eq.(1), has a vanishing energy gap between the ground and first excited state.
To evaluate sums of eigenvalues in the limit $`N\mathrm{}`$, one can interpret eq.(2) as a Fourier series of $`\lambda (\theta )`$. Thus $`V_k=\frac{1}{2\pi }_0^{2\pi }d\theta \lambda (\theta )\mathrm{e}^{i\theta k}`$. This integral representation is also valid for a finite number $`N`$ of oscillators up to an error $`𝒪(1/N)`$ as long as $`k(N+1)/2`$. Due to the periodic boundary conditions $`V^{\pm 1/2}`$ are also Toeplitz matrices, and their elements $`V_{ij}^{\pm 1/2}=V_k^{\pm 1/2}`$ can be expressed in terms of $`\lambda ^{\pm 1/2}`$ for $`k(N+1)/2`$
$$\left(V^{\pm 1/2}\right)_k=\frac{1}{2\pi }_0^{2\pi }d\theta \lambda ^{\pm 1/2}(\theta )\mathrm{e}^{i\theta k}.$$
(4)
Since for the spatial correlation of an oscillator system holds $`q_iq_{i+l}V_l^{1/2}`$ Reznik , the analytic properties of $`\lambda ^{1/2}`$ determine the spatial correlation length $`\xi `$:
$`\xi ^1`$ $``$ $`\underset{l\mathrm{}}{lim}{\displaystyle \frac{1}{l}}\mathrm{ln}\left|q_iq_{i+l}\right|=\underset{l\mathrm{}}{lim}{\displaystyle \frac{1}{l}}\mathrm{ln}\left|V_l^{1/2}\right|`$ (5)
$`=`$ $`\underset{l\mathrm{}}{lim}{\displaystyle \frac{1}{l}}\mathrm{ln}\left|{\displaystyle \frac{1}{2\pi }}{\displaystyle _0^{2\pi }}d\theta \lambda ^{1/2}(\theta )\mathrm{e}^{i\theta l}\right|.`$
If some derivative of $`\lambda ^{1/2}(\theta )`$, say the $`m`$th one, does not exists, partial integrations shows that the integral has a contribution proportional to $`l^m`$. In this case the correlation length $`\xi `$ is infinite, defining a critical system. If on the other hand $`\lambda ^{1/2}(\theta )`$ is smooth the integral decays faster than any polynomial in $`l^1`$. In this case the correlation length is finite, corresponding to a non-critical system. From the form of $`\lambda (\theta )`$ given in eq.(3) it is clear that a regular spectral function implies a finite correlation length, i.e. a noncritical behavior and a singular one an infinite correlation length, i.e. a critical behavior.
In the following we will show that the analytic properties of $`\lambda `$ also determine the entanglement scaling of the oscillator system. The bi-partite entanglement of a compact block of $`N_1`$ oscillators (inner partition $``$) with the rest (outer partition $`𝒪`$) is determined by the $`N_1`$-dimensional sub-matrices $`A`$ and $`D`$ Bombelli-PRD-1986 ; Srednicki-PRL-1993 ; Reznik ; Plenio
$$V^{1/2}=\left[\begin{array}{cc}A& B\\ B^T& C\end{array}\right],V^{1/2}=\left[\begin{array}{cc}D& E\\ E^T& F\end{array}\right],$$
(6)
$`C`$ and $`F`$ are here $`\left(NN_1\right)\times \left(NN_1\right)`$ matrices. The entropy is given by the eigenvalues $`\mu _i1`$ of the matrix product $`AD`$ Plenio :
$`S={\displaystyle \underset{i=1}{\overset{N_1}{}}}f\left(\sqrt{\mu _i}\right),`$ (7)
where $`f(x)=\frac{x+1}{2}\mathrm{ln}\frac{x+1}{2}\frac{x1}{2}\mathrm{ln}\frac{x1}{2}`$. Despite the simplicity of its form, (7) cannot be evaluated in general. This is in contrast to spin systems where $`AD`$ is itself a Toeplitz matrix Korepin ; Mezzadri .
An upper bound to $`S`$ can be found from the logarithmic negativity $`\mathrm{ln}\rho ^\mathrm{\Gamma }`$, where $`\rho ^\mathrm{\Gamma }`$ is the partial transpose of the total ground state $`\rho `$ and $`||||`$ denotes the trace norm. As shown in Plenio ; Eisert-preprint the logarithmic negativity is bounded by the square root of the maximum eigenvalue of $`V`$ and a sum of absolute values of matrix elements of $`V_{ij}^{1/2}`$ between all sites $`i`$ and $`j𝒪`$.
$$S4\lambda _{\mathrm{max}}^{1/2}\underset{i}{}\underset{j𝒪}{}\left|V_{ij}^{1/2}\right|.$$
(8)
A lower bound to the entropy can be found making use of $`\frac{x+1}{2}\mathrm{ln}\frac{x+1}{2}\frac{x1}{2}\mathrm{ln}\frac{x1}{2}>\mathrm{ln}x`$ . This yields
$$S>\frac{1}{2}\underset{i=1}{\overset{N_1}{}}\mathrm{ln}\mu _i=\frac{1}{2}\mathrm{ln}\left(detAD\right).$$
(9)
This estimate has a simple and very intuitive meaning. To see this we first note that the matrix $`D`$ can be expressed in the form $`D=(ABCB^{})^1`$. Thus
$`S`$ $`>`$ $`{\displaystyle \frac{1}{2}}\mathrm{ln}det\left(\mathrm{𝟏}BC^1B^TA^1\right)`$
$`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{ln}det\left(\begin{array}{cc}A& B\\ B^T& C\end{array}\right)\left(\begin{array}{cc}A^1& 0\\ 0& C^1\end{array}\right)`$
$`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{ln}{\displaystyle \frac{detAdetC}{detV^{1/2}}}={\displaystyle \frac{1}{2}}\mathrm{ln}{\displaystyle \frac{detFdetD}{detV^{1/2}}}.`$
where the last equation was obtained by expressing $`A`$ in terms of $`D,E`$, and $`F`$. The last line of (Entanglement and criticality in translational invariant harmonic lattice systems with finite-range interactions) is just Shannon’s classical mutual information $`I(Q_1:Q_2)`$ or $`I(P_1:P_2)`$ respectively, where $`Q_1=(q_1,q_2,\mathrm{},q_{N_1})`$ and $`Q_2=(q_{N_1+1},\mathrm{},q_N)`$ are the position vectors of the two subsystems and $`P_{1,2}`$ the respective momentum vectors. $`I(Q_1:Q_2)`$ is defined as
$$I(Q_1:Q_2)=\mathrm{d}^NQp(Q_1,Q_2)\mathrm{ln}\frac{p(Q_1,Q_2)}{p_1(Q_1)p_2(Q_2)}$$
(15)
where $`p(Q_1,Q_2)=|\mathrm{\Psi }_0|^2`$ is the total and $`p_{1,2}(Q_{1,2})`$ the reduced probability density in position space. Straight forward calculation shows
$$I(Q_1:Q_2)=\frac{1}{2}\mathrm{ln}\frac{detAdetC}{detV^{1/2}}S.$$
(16)
In order to evaluate Shannon’s mutual information in the form given in eq.(9) we want to make use of the asymptotic properties of Toeplitz matrices. For this we note that since $`V^{\pm 1/2}`$ are Toeplitz matrices, so are $`A`$ and $`D`$. Their elements $`A_k`$, and $`D_k`$ can be obtained from $`\lambda ^{\pm 1/2}`$ by (4) if $`N_1(N+1)/2`$.
If $`\lambda (\theta )`$ is regular, we can apply the strong Szegö theorem Szegoe , which states:
$`det(D)\mathrm{exp}\left\{c_0N_1+{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}k|c_k|^2\right\},`$ (17)
for $`N_1\mathrm{}`$. Here the $`c_k`$ are Fourier-coefficients of $`\mathrm{ln}\lambda ^{1/2}(\theta )`$, i.e. $`c_k=\frac{1}{2\pi }_0^{2\pi }d\theta \mathrm{ln}\lambda ^{1/2}(\theta )\mathrm{e}^{i\theta k}`$. Noting that the corresponding coefficients for $`A`$ have opposite sign, we find the lower bound
$$S\frac{1}{2}\mathrm{ln}(det(A))+\frac{1}{2}\mathrm{ln}(det(D))=\underset{k=0}{\overset{\mathrm{}}{}}k|c_k|^2.$$
(18)
To find an upper bound to $`S`$ we make use of eq.(8). For a finite-range interaction there is always a maximum eigenvalue $`\lambda _{\mathrm{max}}^{1/2}`$. Furthermore since $`\lambda _0^{1/2}(\theta )`$ is smooth, eq.(4) implies an exponential bound to the matrix elements of $`V^{1/2}`$. I.e. $`|V_{ij}^{1/2}|K\mathrm{exp}\{\alpha |ij|\}`$, for $`|ij|(N+1)/2`$, where $`K,\alpha >0`$. With this we find
$`{\displaystyle \underset{i}{}}{\displaystyle \underset{j𝒪}{}}\left|V_{ij}^{1/2}\right|`$ $`=`$ $`2N_1{\displaystyle \underset{k=N_1+1}{\overset{(N+1)/2}{}}}|V_k^{1/2}|+2{\displaystyle \underset{k=1}{\overset{N_1}{}}}k|V_k^{1/2}|`$ (19)
$`<`$ $`{\displaystyle \frac{2K\mathrm{e}^\alpha }{(1\mathrm{e}^\alpha )^2}}`$
for $`N,N_1\mathrm{}`$. Thus $`S`$ has also a finite upper bound in 1D. One recognizes that for one dimensional harmonic chains with a regular spectral function $`\lambda (\theta )`$ the entropy has a lower and an upper bound independent on the number of oscillators, which implies an area theorem. Furthermore, as shown above, the spatial correlation length is finite, i.e. the system is non-critical.
Let us now consider a singular function $`\lambda `$. In this case we can calculate the asymptotic behavior of the Toeplitz determinants using Widom’s theorem Widom . This theorem states that for $`N_1\mathrm{}`$ and for $`m_r>1`$:
$$detD\mathrm{exp}\left\{c_0N_1\right\}N_1^{_rm_r^2/4}.$$
(20)
Widom’s theorem cannot be applied to $`A`$, since $`\lambda ^{1/2}(\theta )_r(22\mathrm{cos}(\theta \alpha _r))^{m_r/2}`$ involves negative exponents. We thus employ the alternative expression (Entanglement and criticality in translational invariant harmonic lattice systems with finite-range interactions) containing the matrices $`D`$ and $`F`$. Since the elements of $`D`$ and $`F`$ can only be obtained by the Fouriertransform (4) if their dimension is at most $`(N+1)/2`$, there is only one particular decomposition which we can consider, namely $`N_1=(N1)/2`$ and $`N_2=(N+1)/2`$. For the same reason it is not possible to apply Widom’s theorem to $`V^{1/2}`$ as a whole. $`detV^{1/2}`$ can however easily be calculated directly from the discrete eigenvalues (3). After a lengthy but straight forward calculation we eventually obtain the following expression for the mutual information with $`N_1=(N1)/2`$ and $`N_2=(N+1)/2`$
$`I`$ $`=`$ $`\left({\displaystyle \underset{r=1}{\overset{Q}{}}}{\displaystyle \frac{m_r^2}{4}}\right)\mathrm{ln}N+\mathrm{const}..`$ (21)
Thus a singular spectral function $`\lambda ^{1/2}(\theta )`$ in the case of half/half partitioning leads to a lower bound to the entropy that grows logarithmically with the number of oscillators stating a break-down of the area law of entanglement. As shown above a singular spectral function also implies a diverging spatial correlation length, defining a critical system.
The above discussion can be extended to $`d`$ dimensions. In this case one would consider the entropy $`S`$ of a hypercube of oscillators with dimensions $`N_1\times N_2\times \mathrm{}\times N_d`$. Since we are interested in the thermodynamic limit we can again assume $`N_i(N+1)/2`$. In this case the matrices $`A`$ and $`D`$ are Toeplitz matrices with respect to each spatial direction and their elements $`A_{k_1,k_2,\mathrm{},k_d}`$ can be obtained from the square root of the $`d`$-dimensional function $`\lambda (\theta _1,\mathrm{},\theta _d)=_{k_1=0}^{N_11}\mathrm{}_{k_d=0}^{N_d1}V_{k_1,\mathrm{},k_d}\mathrm{exp}\left\{i_{j=1}^d\theta _jk_j\right\}`$. If $`\lambda ^{1/2}`$ is regular, the $`d`$-dimensional Szegö theorem holds d-dimSzegoe , which asserts that the Toeplitz determinant of dimension $`n_1\times n_2\times \mathrm{}\times n_d`$ has the asymptotic form
$`detD\mathrm{exp}\left\{c_0n_1\mathrm{}n_d+{\displaystyle \underset{j=1}{\overset{d}{}}}{\displaystyle \frac{n_1\mathrm{}n_d}{n_j}}|C_j|\right\},`$ (22)
where $`c_0=\frac{1}{(2\pi )^d}_0^{2\pi }d\theta _1\mathrm{}_0^{2\pi }d\theta _d\mathrm{ln}\left(\lambda ^{1/2}(\theta _1,\mathrm{},\theta _d)\right)`$, and the $`C_i`$ are some constants, whose explicit form is of no interest here. We see that under the above conditions for the $`d`$-dimensional characteristic function $`\lambda ^{1/2}(\theta _1,\mathrm{},\theta _d)`$ the entropy has the lower bound
$$S>\underset{j=1}{\overset{d}{}}\frac{n_1n_2\mathrm{}n_d}{n_j}C_jn^{d1}$$
(23)
which is again proportional to the surface area. We note that the lower bound (23) to the entropy given by the multi-dimensional Szegö theorem is more general than the estimates given in Plenio and Eisert-preprint , which are restricted to nearest neigbor interactions. From the exponential bound to the matrix elements of $`V^{1/2}`$ one can also find an upper bound to the entropy using eq.(8)
$`{\displaystyle \underset{i}{}}{\displaystyle \underset{j𝒪}{}}\left|V_{ij}^{1/2}\right|K{\displaystyle \underset{j=1}{\overset{d}{}}}{\displaystyle \frac{n_1n_2\mathrm{}n_d}{n_j}}n^{d1}.`$ (24)
Eqs.(23) and (24) establish an area law for arbitrary dimensions in the case of a regular spectral function.
In order to obtain a lower bound to the entropy for a singular spectral function in more than one dimension and to show a corresponding break-down of the entanglement area law, one would need a multi-dimensional generalization of Widoms theorem Widom . Although no such generalization is known to us, there is strong evidence for a break down of the area law in higher dimensions. First of all for an interaction matrix that is separabel in the $`d`$ dimensions, i.e. whose elements can be written as products $`V_{i_1,j_1}V_{i_2,j_2}\mathrm{}V_{i_d,j_d}`$, the 1D discussion can straight-forwardly extended to $`d`$ dimensions. Secondly Widom has given a generalization of his matrix theorem to operator functions $`f(A)`$ on $`R^d`$ Widom-2 . The proof given in Widom-2 makes however use of strong conditions on $`f`$ that are not fulfilled for the case we are interested here.
To illustrate validity and break-down of the area theorem let us consider the Hamiltonian $`H=\frac{1}{2}_{i=1}^Np_i^2+\frac{1}{2}_{i=1}^N\left(2\eta q_i+q_{i+1}+q_{i1}\right)^2`$ with periodic boundary conditions. The square root of the spectral function reads in this case $`\lambda ^{1/2}(\theta )=\left|2\eta 2\mathrm{cos}\theta \right|`$. For $`\eta >1`$, $`\lambda ^{1/2}`$ is regular and the correlation length is finite. For $`\eta <1`$, $`\lambda ^{1/2}(\theta )`$ can be written as $`\lambda ^{1/2}(\theta )=\left(22\mathrm{cos}\left(\theta +\theta _0\right)\right)^{1/2}\left(22\mathrm{cos}\left(\theta \theta _0\right)\right)^{1/2}`$, with $`\eta =\mathrm{cos}\theta _0`$, and thus is singular. In this case the correlation length is infinite. We have numerically calculated the entropy for this system for different values of $`\eta `$. The results are shown in fig.1. One recognizes an unlimited logarithmic growth of $`S`$ for $`\eta <1`$ and a saturation for $`\eta >1`$.
In the present paper we discussed the relation between entanglement and criticality in translational invariant harmonic lattice systems with finite-range couplings. We have shown that upper and lower bounds to the entropy of entanglement as well as the correlation length are solely determined by the analytic properties of the spectral function. If the spectral function is regular, the entanglement obeys an area law and the system is non-critical. If the spectral function has a singular part, the area law breaks down and the system is critical. Thus for harmonic lattice systems with translational invariant, non-random and finite-range couplings there is a one-to-one correspondence between entanglement and criticality. We note that some of our results apply also to more general couplings. For the estimates of the entropy it is sufficient that the number of roots of $`h(z)`$ on the unit circle is finite. This is always fulfilled for banded coupling matrices $`V`$ but also holds under more general conditions. For couplings of infinite range, the regular part of the spectral function $`\lambda _0`$ is no longer a polynomial. Thus $`\lambda _0^{1/2}`$ may not be smooth anymore and could have a singularity in a derivative of some order. In such a case the spectral function could be regular, allowing for an entanglement area theorem, and at the same time the correlation length would be infinite, i.e. the system would be critical as in the example of Ref.Eisert-preprint .
The authors would like to thank J. Eisert and M. Cramer for many stimulating discussions. This work was supported by the DFG through the SPP “Quantum Information” as well as the European network QUACS. |
warning/0506/cond-mat0506754.html | ar5iv | text | # High efficiency deterministic Josephson Vortex Ratchet
## Abstract
We investigate experimentally a Josephson vortex ratchet — a fluxon in an asymmetric periodic potential driven by a deterministic force with zero time average. The highly asymmetric periodic potential is created in an underdamped annular long Josephson junction by means of a current injector providing efficiency of the device up to 91%. We measured the ratchet effect for driving forces with different spectral content. For monochromatic high-frequency drive the rectified voltage becomes quantized. At high driving frequencies we also observe chaos, sub-harmonic dynamics and voltage reversal due to the inertial mass of a fluxon.
Long Josephson junction, sine-Gordon, Josephson vortex ratchet, fluxon ratchet, soliton ratchet
The ratchet effect, i.e., the net unidirectional motion of a particle in a spatially asymmetric periodic potential in the presence of deterministic or stochastic forces with zero time average, received a lot of attention during the 20-th century. The second law of thermodynamics does not allow to extract useful work out of equilibrium thermal fluctuations, as was didactically demonstrated by FeynmanFeynman et al. (1963). Thus, the only way to produce useful work is to supply non-native fluctuations (usually colored noise), which is the basic principle of operation for any ratchet.
Particularly during the last decade ratchets were receiving a lot of attentionJülicher et al. (1997); Reimann (2002); Hänggi et al. (2005). Several new implementations, in particular based on the motion of the Josephson phase in SQUIDsSterck et al. (2002) or vortices in long Josephson junctions (LJJ)Goldobin et al. (2001); Carapella (2001); Carapella and Costabile (2001); Carapella et al. (2002) or Josephson junction arrays (JJA)Falo et al. (1999); Trías et al. (2000); Lee (2003), were suggested and tested. The investigation of quantum ratchetsLinke et al. (1999); Grifoni et al. (2002); Majer et al. (2003), i.e., a quantum particle moving/tunneling quantum mechanically in an asymmetric potential, is a fascinating new field not very well developed up to now especially experimentally. Advantages of Josephson junction based ratchets are: (I) directed motion results in an average dc voltage which is easily detected experimentally; (II) Josephson junctions are very fast devices which can operate (capture and rectify noise) in a broad frequency range from dc to $`100\mathrm{GHz}`$, thus capturing a lot of spectral energy; (III) by varying junction design and bath temperature both overdamped and underdamped regimes are accessible; and (IV) one can operate Josephson ratchets in the quantum regimeMajer et al. (2003).
In this letter we investigate experimentally the *deterministic underdamped Josephson vortex ratchet* (JVR), in which a Josephson vortex (fluxon) moves along a LJJ. We implemented a novel, effective way to construct a strongly asymmetric potential by means of a current injector and systematically study a quasi-statically driven ratchet with different spectral content of the driver. For non-adiabatic drive we observe quantized rectification, *voltage reversal*, sub-harmonic, and chaotic dynamics.
Our system can be described by the following perturbed sine-Gordon equationGoldobin et al. (2001)
$$\varphi _{xx}\varphi _{tt}\mathrm{sin}\varphi =\alpha \varphi _t\gamma (x)\xi (t),$$
(1)
where $`\varphi `$ is the Josephson phase, the curvilinear coordinate $`x`$ along the LJJ and the time $`t`$ are normalized to the Josephson penetration depth $`\lambda _J`$ and inverse plasma frequency $`\omega _p^1`$, accordingly, $`\alpha `$ is the dimensionless damping parameter and $`\gamma (x)=j_{\mathrm{inj}}(x)/j_c`$ and $`\xi (t)=j(t)/j_c`$ are the bias current densities normalized to the critical current density of the LJJ. $`\gamma (x)`$ has zero spatial average and is used to create an asymmetric potential (see below). $`\xi (t)`$ is a spatially homogenous deterministic (or stochastic) drive with zero time average. The ultimate aim of ratchet operation is to rectify $`\xi (t)`$ to produce non-zero voltage $`\varphi _t0`$. In the absence of the r.h.s., the solitonic solution of Eq. (1) is a Josephson vortex (sine-Gordon kink) $`\varphi (x)=4\mathrm{arctan}\mathrm{exp}\left[(xx_0(t))/\sqrt{1u^2}\right]`$, situated at $`x_0(t)`$ and moving with velocity $`u=dx_0(t)/dt`$. The r.h.s. of Eq. (1) is usually considered as a perturbation.McLaughlin and Scott (1978) It does not change drastically the vortex shape, but defines its dynamics, e.g., equilibrium velocityMcLaughlin and Scott (1978). Such an approximation essentially treats the vortex as a rigid object, and its dynamics can be reduced to the dynamics of a relativistic underdamped point-like particleCarapella (2001) (cf. non-relativistic case Borromeo et al. (2002)). In this terms the ratchet should rectify $`\xi (t)`$ to produce a nonzero average velocity $`u0`$.
To build a JVR, a fluxon should move in an asymmetric periodic potential $`U(x_0)`$. Opposite to solitonSalerno and Quintero (2002) or stringMarchesoni (1996) ratchets where the *asymmetric* potential $`V(\varphi )`$ is a primitive function of the current-phase relation (CPR), in our case, the CPR is sinusoidal, i.e., $`V(\varphi )=1\mathrm{cos}(\varphi )`$ is a *symmetric* function of $`\varphi `$. Instead, we construct an asymmetric periodic potential $`U(x_0)`$, which is a function of the fluxon *coordinate*. In contrast to other JVRs Falo et al. (1999); Trías et al. (2000); Lee (2003); Carapella (2001); Carapella and Costabile (2001); Carapella et al. (2002), we construct $`U(x_0)`$ using a single current injector. As we proposed earlier (see the last paragraph before Sec. III in Ref. Goldobin et al., 2001) the current injection with profile $`\gamma (x)`$ is equivalent to an applied nonuniform magnetic field $`h(x)`$ such that $`h_x(x)=\gamma (x)`$. Then the potential $`U(x_0)2\pi wh(x_0)`$, where $`w`$ is the LJJ’s widthGoldobin et al. (2001). The periodicity of the potential is provided by using an annular LJJ (ALJJ) and $`\gamma (x)`$ with zero spatial average.
If we apply $`\gamma (x)`$ using a current injector of width $`\mathrm{\Delta }w`$ situated at $`x=x_{\mathrm{inj1}}`$ and extract the current *from the same electrode* along the rest of the LJJ, i.e. $`\gamma (x)=\gamma _1`$ for $`|xx_{\mathrm{inj1}}|<\mathrm{\Delta }w/2`$ and $`\gamma (x)=\gamma _2`$ for all other $`x`$, then $`U(x_0)`$ looks like an asymmetric saw-tooth potential with the steep slope $`\gamma _1`$ and the gentle slope $`\gamma _2`$. Note that, zero spatial average of $`\gamma (x)`$ requires $`j_c\gamma _1\mathrm{\Delta }ww=j_c\gamma _2(L\mathrm{\Delta }w)w=I_{\mathrm{inj1}}`$, where $`L`$ is the LJJ circumference. By changing $`I_{\mathrm{inj1}}`$ one modulates the amplitude of the potential. This allows operation as a flashing ratchet too. Below we focus on the rocking ratchet, i.e., when the potential $`I_{\mathrm{inj1}}`$ is (almost) constant while $`\xi (t)0`$.
Similar systems, but with magnetic field induced potential $`U(x_0)`$, were already studied for the case of quasi-static deterministic and stochastic driveCarapella and Costabile (2001); Carapella (2001) and for deterministic high frequency driveCarapella et al. (2002). Ratchetlike systems based on an asymmetry of the driver (rather than the potential) were also investigatedMarchesoni (1986); Flach et al. (2002); Ustinov et al. (2004).
Experiments have been performed with Nb-AlO<sub>x</sub>-Nb ALJJs that have geometry shown in Fig. 1. A pair of current injectors (inj2) attached to the top electrode and separated by distance $`\mathrm{\Delta }x`$ can be used to insert a fluxon in the ALJJ.Ustinov (2002); Malomed and Ustinov (2004) A single injector (inj1) in the bottom layer is used to create the asymmetric potential as described above. We investigated several samples with different normalized circumferences $`l=L/\lambda _J`$ ($`L=2\pi R`$; $`R`$ is a mean radius, see Fig. 1). The device parameters are summarized in Tab. 1. For all samples $`\lambda _Jw=5\mu \mathrm{m}`$ and $`\lambda _J\mathrm{\Delta }w=\mathrm{\Delta }x=5\mu \mathrm{m}`$ ($`10\mu \mathrm{m}`$ for G3), i.e., we can treat our ALJJ like a one dimensional LJJ and inj2 like an almost ideal discontinuityGaber et al. ; Malomed and Ustinov (2004). The quantity $`\nu _0=\overline{c}_0/L`$ is the maximum revolution frequency of a fluxon ($`\overline{c}_0`$ is the Swihart velocity) and $`V_1=\mathrm{\Phi }_0\nu _0=\mathrm{\Phi }_0\omega _p/l`$ is the corresponding voltage, i.e., the asymptotic voltage of the first fluxon step.
Measurements were performed in a shielded cryostat at $`T=4.2\mathrm{K}`$ unless stated otherwise. Before operating the ratchet, each ALJJ was characterized and injectors were calibrated. All junctions listed in Tab. 1 showed good $`I`$-$`V`$ characteristics (IVC) and nice symmetric $`I_c(H)`$ dependences (not shown). The dependence $`I_c(I_{\mathrm{inj2}})`$ looks like a Fraunhofer pattern (not shown) in accord with theory Malomed and Ustinov (2004). The first minimum is reached at $`I_{\mathrm{inj2}}=\pm 3.9\mathrm{mA}`$ and corresponds to the phase twisted by $`\pm 2\pi `$ in a tiny region between the injectors and to a free (anti)fluxon inserted into the ALJJ outside inj2.
To calibrate inj1, we measure $`I_c(I_{\mathrm{inj1}})`$. The current $`I_{\mathrm{inj1}}`$ is applied between inj1 and the bottom electrode of the ALJJ. In Fig. 2 gray symbols show $`I_c(I_{\mathrm{inj1}})`$ with no fluxons trapped in the junction. Starting from the main maximum, the $`I_c(I_{\mathrm{inj1}})`$ curve has different slopes for increasing or decreasing $`I_{\mathrm{inj1}}`$. The ratio of these slopes corresponds to the ratio of slopes of the $`h(x)`$ saw-toothGoldobin et al. (2001). Then we applied $`I_{\mathrm{inj2}}=\pm 3.9\mathrm{mA}`$ to inject a fluxon, and measured $`I_c(I_{\mathrm{inj1}})`$ again, as shown by black symbols in Fig. 2. We see that $`I_c(0)`$ dropped down almost to zero (the residual pinning is present due to the finite inj2 sizes $`\mathrm{\Delta }x`$ and $`\mathrm{\Delta }w`$Malomed and Ustinov (2004)). Applying a finite $`I_{\mathrm{inj1}}`$ creates a potential which pins the fluxon stronger and one needs to apply a larger bias current $`I`$ to let the fluxon move around the ALJJ and generate a voltage. For small $`I_{\mathrm{inj1}}`$ the depinning current $`I_c`$ grows almost linearly with $`I_{\mathrm{inj1}}`$, but it is *asymmetric* for positive and negative direction of the bias current (driving force). This corresponds to different slopes of the potential when the fluxon tries to move to the left or to the right out of the well. The ratio of slopes of the $`I_c(I_{\mathrm{inj1}})`$ dependence for one trapped vortex reflects the asymmetry of the potential $`U(x_0)`$Goldobin et al. (2001) and for sample C3 is about $`6.6`$. Figure 2 also defines our “working area”, i.e., the reasonable range of $`|I_{\mathrm{inj1}}|0.6\mathrm{mA}`$ for ratchet operation, shown by the dashed lines.
If we apply $`I_{\mathrm{inj2}}=3.9\mathrm{mA}`$ (insert one fluxon) and $`I_{\mathrm{inj1}}=0.2\mathrm{mA}`$ (create a potential of intermediate amplitude) a fluxon step appears on the IVC (not shown), corresponding to the rotation of a fluxon around the ALJJ with $`u\overline{c}_0`$. The depinning $`I_c`$ and return current $`I_r`$ of the fluxon step depends on the polarity of the applied bias current as well as on $`I_{\mathrm{inj1}}`$. To demonstrate rectification of the rocking ratchet we apply the periodic bias current (deterministic driving force) $`I(t)=\xi (t)j_cLw=I_{ac}\mathrm{sin}(2\pi \nu t)`$ with frequency $`\nu =100\mathrm{Hz}\nu _0`$ (quasi-static regime) and measure $`V_{dc}(I_{ac})=V(I_{ac})`$ by averaging the voltage over $`10\mathrm{ms}`$ ($`1000`$ data points sampled at $`100\mathrm{kHz}`$) — one period of the ac drive. For ALJJ C3 $`V_{dc}(I_{ac})=V_1u(I_{ac})`$ is shown in Fig. 3 for different values of $`I_{\mathrm{inj1}}`$ (open circles). All $`V_{dc}(I_{ac})`$ curves have similar features. For small $`I_{ac}`$ the driving force acting on a fluxon is not sufficient to push the fluxon out of the well in either direction so that $`V_{dc}=0`$. At higher amplitude the bias is able to push the fluxon in one direction, which results in $`V_{dc}0`$. The value of $`I_{ac}`$ at which rectification first appears grows linearly with $`I_{\mathrm{inj1}}`$ as it should be according to Fig. 2. Further, $`V_{dc}`$ grows with $`I_{ac}`$, but this dependence due to relativistic saturation of fluxon velocity is weaker than the linear dependence predicted for a non-relativistic particleBorromeo et al. (2002). At an even larger amplitude $`I_{ac}`$ the junction switches into the resistive state, generating a high positive or negative dc voltage. The latter regime is not discussed here as it has nothing to do with the operation of the JVR. We note that we very rarely observed a decrease of $`V_{dc}`$ at higher $`I_{ac}`$ — a typical behavior for many ratchet systems. In our case, the asymmetry is so high, that a negative fluxon step does not appear in most cases.
We have also investigated the influence of the shape (spectral content) of the driving force $`I(t)`$ on the performance of our ratchet. As an example, in Fig. 3 we also show rectification curves $`V_{dc}(I_{ac})`$ for *time symmetric* saw-tooth pulses and for a *time symmetric* rectangular drive at $`I_{\mathrm{inj1}}=600\mu \mathrm{A}`$. One can see that rectangular pulses result in higher performance (rectification) because the ALJJ spends more time at the fluxon step than in the case of a sinusoidal drive. Similarly, a saw-tooth drive results in lower rectification. Usually it is believed that a drive with compact spectrum (e.g., monochromatic) is more efficient in terms of rectification than a drive with broad spectrum (e.g., white noise, which most frequently provides no rectification). Our results show that this is not always true: a rectangular drive with a rather broad discrete spectrum is more efficient than a monochromatic one. On the other hand, a saw-tooth drive also with broad spectrum, is less efficient than a sinusoidal one.
A typical value of rectified voltage $`V_{dc}20\mu \mathrm{V}V_1/2`$) and is higher than reported earlierCarapella and Costabile (2001). In principle, $`V_11/L`$, so one can get higher $`V_{dc}`$ for a shorter LJJ. However, when the LJJ becomes too short, the asymmetry of the potential vanishes due to the convolution with the fluxon shape, see Eq. (17) of Ref. Goldobin et al., 2001.
To drive the ratchet at higher frequencies ($`1\mathrm{GHz}`$) we placed an emitting antenna connected to an rf generator close to the ALJJ, so that the bias leads act as a pickup antenna. Note that this induces no signal in inj2 and a rather small signal in inj1, so that our ratchet has $`90\mathrm{}95\%`$ of a “rocking” and $`5\mathrm{}10\%`$ of a “flashing” potential. The dc voltage was averaged over $`10\mathrm{ms}`$ ($`10^7`$ periods) and recorded vs. applied power $`P`$ as shown in Fig. 4 for sample C3. Rather than growing smoothly, the dc voltage is quantized, $`V_n=n\mathrm{\Phi }_0\nu `$ (Shapiro-like steps). Each step corresponds to an integer number $`n`$ of turns of a fluxon around the ALJJ per period of ac drive. The curves in Fig. 4 are quite noisy because their voltage is comparable to the residual noise of our measurement setup. Note that the $`P`$-axis shows the applied power at the output of the generator, so the curves corresponding to different frequencies may appear shifted along the $`P`$-axis due to the frequency dependent coupling of microwaves to the junction. At higher applied power we sometimes (e.g., for sample C3 at $`\nu =0.5\mathrm{GHz}`$) see a stepwise decrease of the rectified dc voltage, see Fig. 4. This feature can be reproduced in simulations and is related to the bifurcation to the period two dynamics. In general, multiple discrete voltage steps like in Fig. 4 can be observed for frequencies up to $`3\mathrm{GHz}`$.
When the time $`t_0`$ required for one revolution of a fluxon around the ALJJ (in the best case $`t_0=\overline{c}_0/L`$) becomes comparable with the period of the driving force $`1/\nu `$, i.e., the voltage $`V=\mathrm{\Phi }_0\nu `$ is approaching $`V_1`$, the fluxon has time only for one revolution, so only one step can be observed on the $`V_{dc}(P)`$ curves.
For frequencies $`\nu >4\mathrm{GHz}`$ we used a $`(0.7\times 2.5\times 1.5)\mathrm{cm}^3`$ copper box with the first mode at $`6\mathrm{GHz}`$ to achieve better coupling and avoid multiple low frequency resonances. To make the system less chaotic we increase the damping by measuring at slightly higher temperature $`T6\mathrm{K}`$. In Fig. 5 we show $`V_{dc}(P)`$ for ALJJs C4, G3, and C2 (see Tab. 1). The curves show a single quantized voltage step corresponding to one fluxon revolution. For sample C4 we achieved operation of the ratchet up to very high frequency of $`29\mathrm{GHz}`$ ($`V_{dc}/V_1=u=0.88`$). For sample G3, which is longer, we reach $`12\mathrm{GHz}`$ ($`u=0.91`$). Our normalized average velocity $`u`$ is considerably larger than 0.22Ustinov et al. (2004) or $`0.33`$Carapella et al. (2002), reported recently for similar ratchet(-like) systems. The curve for sample C2 shows a negative dc voltage step, as expected for $`I_{\mathrm{inj1}}=470\mu \mathrm{A}>0`$, see Fig. 2. Then $`V_{dc}`$ averages to zero over a large range of $`P`$, but for $`7.5\mathrm{dBm}<P<12\mathrm{dBm}`$ there appears a quantized dc voltage of opposite polarity. This *voltage reversal* (or particle’s *current reversal*) corresponds to the motion of the fluxon in an unnatural direction and was predicted for point-like underdamped deterministically driven particleBarbi and Salerno (2000); Jung et al. (1996); Mateos (2000). Note the two ingredients essential for current reversal in our case: (a) high driving frequencies and (b) low damping. For low frequencies we observe no current reversal, see Fig. 3. In the overdamped regime a monochromatically driven ratchet shows no current reversal eitherBartussek et al. (1994). At high power our JVRs show unquantized $`V_{dc}`$ regions, corresponding to chaotic dynamics. Such regimes can be reproduced in simulations. The details will be presented elsewhere.
In conclusion, we investigated experimentally *a relativistic underdamped Josephson vortex ratchet* where a strongly asymmetric potential is created using a current injector. We observed quantized rectification of a deterministic signal at high frequencies up to $`29\mathrm{GHz}`$, average velocity $`u`$ up to $`0.91`$, voltage reversal, as well as sub-harmonic and chaotic regimes. Deterministic quasi-static signals with broad spectra may be rectified better or worse than monochromatic ones.
We thank P. Hänggi for discussions. This work is supported by the Deutsche Forschungsgemeinschaft. |
warning/0506/astro-ph0506003.html | ar5iv | text | # H2 Pure Rotational Lines in the Orion Bar
## 1 Introduction
A substantial fraction of the dense interstellar medium resides in clouds where far-ultraviolet photons emitted by hot stars dominate both the energetics and the chemistry of the primarily neutral gas (Hollenbach & Tielens 1999). In the dense ISM, the neutral photodissociation region (PDR) material includes both extended molecular clouds with only modest column densities (N$`{}_{\mathrm{H}_2}{}^{}<`$10$`{}_{}{}^{22}\mathrm{cm}_{}^{2}`$) (Plume et al. 1999; Jansen et al. 1995) and surface layers of clumps within higher column density molecular cores (Stutzki et al 1988). Within photodissociation regions, the material makes the transition from hot, ionized gas to cold, molecular gas as attenuation of the far-ultraviolet field increases farther from the cloud surface. As one moves from the ionization front deeper into the molecular cloud, hydrogen goes from atomic to molecular form and carbon goes from C<sup>+</sup> to C<sup>o</sup> and then to CO (Tielens & Hollenbach 1985b; Black & van Dishoeck 1987). The rich chemistry of the photodissociated gas differs significantly from classical dark-cloud chemistry driven by cosmic-ray ionization (Sternberg & Dalgarno 1995).
The thermal balance in the photodissociation region is intimately connected with the radiative transfer for UV photons that drive the chemistry and energetics and with the chemical state of the material in different layers of the structure (Draine & Bertoldi 1999). The most important heating mechanisms include ejection of photoelectrons from dust grains (Bakes & Tielens 1994; Weingartner & Draine 2001, and refs therein) and the collisional deexcitation of H<sub>2</sub> molecules initially excited by UV photons (Sternberg & Dalgarno 1989). Deeper into the regions, gas-grain collisions may also heat the gas. Fine-structure lines of neutral atoms or of singly ionized species with low ionization potential provide the cooling in the outer layers of the PDR while CO rotational lines cool the predominantly molecular inner regions. Quadrupole rotational lines of H<sub>2</sub> can contribute to the cooling at intermediate depths.
The cooling of PDRs produces a broad variety of line and dust feature emission, each arising in a particular layer of the photodissociated structure and each with its own dependence on the density of the region and on the strength of the incident radiation field. From the outer, predominantly atomic layers, one sees emission in the far-IR fine-structure lines of \[CII\] (158 $`\mu `$m) and \[OI\] (63 and 145 $`\mu `$m), as well as lines of \[FeII\] and \[SiII\] at somewhat shorter wavelengths. The 3.3 $`\mu `$m feature attributed to PAHs also appears to arise in this zone. Deeper into the cloud, carbon becomes neutral and the 370$`\mu `$m and 609 $`\mu `$m \[CI\] fine-structure lines become important emitters. Farther into the neutral zone, emission from a few molecules, notably CN and HCO<sup>+</sup>, is significantly enhanced over dark-cloud values. Because of the enhanced temperatures, high-J lines of CO can also trace the distribution of PDR material.
The rotational and ro-vibrational transitions of molecular hydrogen are particularly useful tracers of the properties of PDRs. In the outer portions of the PDR, H<sub>2</sub> formed on dust grains is excited by UV photons. A radiative cascade through the ro-vibrational levels of the ground electronic state follows this fluorescent excitation and produces a distinctive pattern of emergent line strengths (Black & van Dishoeck 1987; Hasegawa et al. 1987). At densities greater than $`5\times `$10<sup>4</sup> cm<sup>-3</sup>, collisional de-excitation modifies the line ratios as it heats the gas (Sternberg & Dalgarno 1989; Luhman et al. 1997). Models show that self shielding allows hydrogen to make a transition from predominantly atomic to predominantly molecular form at a depth into the PDR where the temperature can be high enough to produce significant populations up to $`J=68`$ in the H<sub>2</sub> ground rotational state (Draine & Bertoldi 1996). This warm molecular gas then produces strong emission in the mid-IR rotational lines of H<sub>2</sub> and may even contribute significantly to emission in transitions arising from low J levels of the first excited vibrational state. The low critical density of the ground state rotational transitions makes ratios of mid-IR lines good probes of the temperature in the layers where they arise.
Unfortunately, the line emission from PDRs rarely arranges itself into the neat stratified pattern one might expect from a homogeneous, plane-parallel region. Furthermore, the extent of emission tracing the surface layers of clouds, particularly CII and H<sub>2</sub> emission, is often larger than one would naively expect based on likely column densities and dust opacities. The usual explanation invoked to resolve this discrepancy is that the PDR material actually resides on the surface layers of clumps and that UV radiation reaches these surfaces by propagating through a much more tenuous interclump medium (Stutzki et al 1988; Howe et al. 1991; Meixner & Tielens 1993; Spaans 1996).
The first observations of PDRs in the H<sub>2</sub> ground-state rotational lines produced a surprising result: the bulk of the H<sub>2</sub> emission was coming from very warm gas. Observations of the 0-0 S(1) and S(2) emission from the Orion Bar implied temperatures of 400-1000 K (Parmar, Lacy, & Achtermann 1991). Lower spectral and spatial resolution ISO observations of PDRs in S140, NGC 7023, and NGC 2023 in a larger number of H<sub>2</sub> rotational transitions are consistent with temperatures in a similar range (Timmermann et al. 1996; Fuente et al. 1999; Bertoldi et al. 2000). These observations have served as a testbed for models of the thermal balance and chemistry of PDRs (Draine & Bertoldi 1999; Bertoldi et al. 2000). The large number of transitions detected in the ISO observations help constrain the models. The poor spatial and spectral resolution of these observations, however, leave a number of questions open. The width and velocity of the H<sub>2</sub> lines cannot be measured accurately enough to compare them to CO, CS, or HCN emission from farther into the PDRs. The large beams inevitably lead to an averaging of flux from different layers within the PDR structure. This averaging could destroy important information about the physics and structure of the PDR.
The PDR known as the Orion Bar lies $``$2′ southeast of the Trapezium stars in Orion at the interface of the HII region formed by the Trapezium and dense gas associated with the Orion Molecular Cloud. The nearly edge-on geometry of the Orion Bar (Marconi et al. 1998) lends itself well to the study of PDRs. Because of its favorable geometry and relatively close distance ($``$450 pc, Hoogerwerf, de Bruijne, & de Zeeuw 2000), it is possible to observe stratification in the Orion Bar (Figure 1). The ionization front defined by the sharp edge of the radio continuum emission (Felli et al. 1993) lies immediately to the northwest of the maximum emission in the 3.3 $`\mu `$m PAH feature (Bregman et al. 1994). Approximately 15″ southeast of the ionization front, there is a strong maximum in the distribution of H<sub>2</sub> 1-0 S(1) line flux (van der Werf et al. 1996). The peak of the high column density molecular ridge, as seen in CO, lies an additional 10″ to the southeast (see Fig.1, Tielens et al. 1993). In this work, we present new maps of the 0-0 S(1) and 0-0 S(2) transitions of H<sub>2</sub> and selected observations of the 0-0 S(4) line toward the Orion Bar. These observations have both high spatial (2″) and spectral ($``$4 km s<sup>-1</sup>) resolution. We compare these results to observations of the 1-0 S(1) line at similar spatial resolution (van der Werf et al. 1996). The high spatial resolution probes a critical scale in Orion (A<sub>V</sub>=1 at n<sub>H</sub>=3$`\times `$10<sup>4</sup> cm<sup>-3</sup> corresponds to 10″ at the distance of Orion) and the spectral resolution of our observations is sufficient to resolve lines with the same width as the CO lines arising in the bar.
By obtaining high spatial resolution H<sub>2</sub> pure-rotational line maps, in a region with nearly edge-on geometry, we hope to resolve the thermal and chemical structure of PDRs. We will look not only at the temperature distribution in the Orion Bar, but also the column densities of H<sub>2</sub> and compare the distribution of H<sub>2</sub> rotational emission to other tracers of gas in PDRs. We will use the data to test PDR models of the warm region at the H/H<sub>2</sub> interface.
## 2 Observations and Data Reduction
We mapped the H<sub>2</sub> v = 0-0 S(1), S(2), and S(4) lines at 17.03 $`\mu `$m, 12.28 $`\mu `$m, and 8.03 $`\mu `$m (Table 1) toward the Orion Bar in 2002 December. We made the observations using the Texas Echelon Cross Echelle Spectrograph (TEXES, Lacy et al. 2002) on the 3m NASA Infrared Telescope Facility (IRTF). The spatial scale on the 256<sup>2</sup> Si:As array was 0.35″ pixel<sup>-1</sup>. The spectral resolution of TEXES was determined from emission lines (C<sub>2</sub>H<sub>2</sub> and C<sub>2</sub>H<sub>6</sub>) toward Titan, assumed to be unresolved, observed during the same observing run. Table 2 lists the slit width, length, resolving power, and total integration time for each line. We oriented the slit parallel to the Orion Bar ionization front, at a position angle of 45 or northeast-southwest.
With TEXES, the standard method for producing maps of spectral lines is to step the slit across the object and a portion of adjacent sky, taking a spectrum at each position, with the telescope secondary mirror held fixed (i.e., without chopping). For the S(1) and S(2) lines, we mapped the Orion Bar by stepping the telescope from northwest to southeast in 1/2 slit width steps (0.7″ for S(2) and 1″ for S(1)) to create 40″ long scans. The 0,0 position for the maps is at R.A. = 5<sup>h</sup> 35<sup>m</sup> 19.7<sup>s</sup>, Dec. = -5 25′ 28.3″ (J2000.0) and the mapped region runs from 13″ northwest to 27″ southeast of this position (Figure 2). Spectra taken at the end of each scan, where no line emission is present, are then used as sky frames in the data reduction (Lacy et al. 2002). In Orion, we used the last 7″ of our scans, the positions farthest from the ionization front, for sky subtraction. In order to cover the same area with the shorter S(2) slit as we did for the S(1) line, we mapped the Bar at 12.28 $`\mu `$m by making two scans offset by $`\pm `$2″ northeast-southwest with respect to the center of our 0-0 S(1) map. Since autoguiding is unavailable while using TEXES in scan mode, we made note of the telescope drift during the scan as the guide star passed through the boresight at the (0,0) position, and recentered the telescope on the boresight before each set of four scans.
We determined our absolute positional uncertainty by scanning the Becklin-Neugebauer object (BN) multiple times during our observations. We found that the maximum position drift for the step scans was 1.3″ over the duration of our observations. The relative positional uncertainty of our observations was determined from the offset of the guide star observed during each scan. The RMS drift was 0.7″ with a maximum of 2.1″. The main effect of these drifts would have been to smear out our beam. However, by noting these drifts as they occurred, we were able to correct for them during the summation that produced the final maps, thereby making the smearing negligible compared to our final 2″ to 2.5″ resolution. We therefore estimate the relative positional uncertainty in our 0-0 S(1) and S(2) maps to be less than 0.5″.
For the S(4) line, we took spectra at 6 different depths into the bar spaced 1.4″ apart. The slit positions for these observations are superposed on the map in Figure 3. We performed sky subtraction by nodding the telescope every 8-16 seconds between the source position and a reference position 80″ south.
We reduced the raw images of cross-dispersed spectra using the standard TEXES pipeline reduction program (Lacy et al. 2002). This program removes artifacts, flat-fields the frames using images of ambient and sky loads, extracts the spectra and removes telluric absorption. The telluric lines (principally from H<sub>2</sub>O, CO<sub>2</sub> and CH<sub>4</sub>) are also used to provide a wavelength calibration accurate to $``$1 km s<sup>-1</sup>.
Unlike other visible and infrared spectrometers, TEXES can provide an absolute radiometric calibration of line intensities. The method makes use of observations of ambient and sky loads to compute the system throughput and gain, as well as the absorptivity of the atmosphere (Lacy et al. 2002). For maps of point sources or for line sources with extents significantly larger than the slit width, this radiometric calibration does not need to be corrected for slit losses.
We used maps of $`\beta `$Gem and $`\alpha `$Tau to confirm the absolute calibration scale. Scans across $`\beta `$Gem and $`\alpha `$Tau (at 12.3 $`\mu `$m) yielded radiometric flux densities of 8.0 and 44.9 $`\times `$ 10<sup>-22</sup> ergs cm<sup>-2</sup> s<sup>-1</sup> Hz<sup>-1</sup>, respectively. These values agree to within 5% with flux densities derived from the N Band magnitudes of $`\beta `$Gem and $`\alpha `$Tau (Tokunaga 1984) extrapolated to 12.3 $`\mu `$m (8.4 and 43.4 $`\times `$ 10<sup>-22</sup> ergs cm<sup>-2</sup> s<sup>-1</sup> Hz<sup>-1</sup>).
Once properly calibrated spectra were available for each position along the 40″ scans (or, in the case of the S(4) line, each nod pair), we subtracted a linear baseline from each spectrum. The individual spatial scans were then shifted in position to correct for telescope drift and spatially coincident spectra (40 scans in S(2) and 80 scans in S(1)) were averaged together. To smooth the data, we resampled the S(1) data onto the same grid as our S(2) observations and smoothed both datasets with a 1.5″ Gaussian. The final spatial resolution along the slit for all 3 ground state H<sub>2</sub> lines is $``$2.0″, while the resolution in the direction of our scan is 2.5″ for the 0-0 S(1) data and 2.0″ for the 0-0 S(2) data.
We determined the statistical uncertainty of the integrated intensities in the summed spectra at each position by calculating the rms noise integrated over a stretch of baseline adjacent to the line emission and comparable in width to the region containing most of the line flux. Table 2 lists the average uncertainty in the integrated intensities for the different lines. The highest signal to noise ratio for any individual 2″ resolution element is 58 for the 0-0 S(1) spectra and 86 for the 0-0 S(2) spectra.
## 3 Results
### 3.1 Maps of H<sub>2</sub>
The first two panels of Figure 2 show the distribution of v=0-0 S(1) and S(2) intensity. The final panel shows the distribution of intensity in the 2.12 $`\mu `$m v = 1-0 S(1) line (van der Werf et al. 1996) resampled onto the same grid at the same spatial resolution as our observations of the ground-state rotational lines. For all three, up is northwest. The dominant feature in the maps is the bright horizontal (northeast-southwest) ridge centered at y$``$-2″. The overall thickness of the bright ridge is $``$8″ (0.017 pc). The brightness of the side of the ridge facing the ionization front and $`\theta ^1`$C Ori rises very steeply in the ground-state rotational lines. The cuts shown in Figure 4 indicate that the 0-0 S(1) and S(2) intensity increases from $`<`$10% to 50% of the peak value in 2″–3″, or barely more than a single resolution element. In two of the three cuts, the rise toward the peak of the ridge is much more gradual in the 1-0 S(1) line with emission present at 20-50% of the peak value $``$5″ in front of the peak.
There is remarkable agreement in the intensity distribution, not only between the 0-0 S(1) and 0-0 S(2) lines but also between the distributions in these lines from the vibrational ground state and the distribution of v=1-0 S(1) emission (Figure 2). The three H<sub>2</sub> line maps agree to within the uncertainties about the location and width of the bright ridge, as well as the presence and extent of lower-level extended emission behind the ridge. Further, there is substructure within the bright ridge that is present in all three maps. There are numerous other small-scale features in the map, with sizes ranging from unresolved up to $``$10″. The smaller scale structure in the intensity distribution is also evident in the cuts shown in Figure 4. Along these cuts, displaced by only 3″ to the northeast or southwest, there are substantial differences at adjacent points.
Typical signal-to-noise for the integrated intensities along the bright ridge is better than 25 and the peak intensities are 9.6 and 8.2 $`\times `$ 10<sup>-4</sup> ergs s<sup>-1</sup> cm<sup>-2</sup> sr<sup>-1</sup> in 0-0 S(1) and S(2) respectively. Along the ridge, the line intensities do not drop below 60% of the peak The smaller scale structures behind the ridge have intensities that range from 10 to 50% of the peak. The weakest detected emission in both lines is $``$0.6 $`\times `$ 10<sup>-4</sup> ergs s<sup>-1</sup> cm<sup>-2</sup> sr<sup>-1</sup>, centered around (+3,+10). The signal to noise in the individual 0-0 S(4) spectra is $``$20 at the brightest positions.
### 3.2 Spectral Properties
In order to study the details in the spectral line shapes and to compare line shapes from one part of the H<sub>2</sub> emission region to another, we need to average over larger areas to improve the signal to noise. We have chosen a set of areas that encompass different morphological features in the maps of integrated line strength. We have not placed any of these areas along the leading edge of the bright ridge where the steep intensity gradients combined with even our modest pointing uncertainties make comparisons uncertain. Figure 3 shows the outlines of the areas and assigns labels to them. These labels appear then in Figure 5 next to the 0-0 S(1) and 0-0 S(2) (and, for areas A and B, 0-0 S(4)) spectra created by integrating over the designated areas. With the spectral resolution available with TEXES, we are able to resolve the lines. Measured FWHM line widths are 5-8 km s<sup>-1</sup> (Table 3). Deconvolving the instrument profile and assuming that the intrinsic line shape is Gaussian, we derive physical FWHM line widths of 4-6 km s<sup>-1</sup> (Table 3), in agreement with linewidths expected for optically thin, thermal gas at $``$1000 K. The instrumental profile, however, may be closer to a Lorenzian, which could imply line widths as low as 2-4 km s<sup>-1</sup>. Line widths for the 0-0 S(1) and S(2) transitions agree to within 1-2 km s<sup>-1</sup>, with neither line being systematically wider than the other, while physical linewidths for the 0-0 S(4) spectra are wider than the 0-0 S(1) and S(2) linewidths for areas A and B by 1.2 km s<sup>-1</sup> and 1.5 km s<sup>-1</sup>. In the molecular gas deeper in the cloud, Hogerheijde, Jansen, & van Dishoeck (1995) and Batrla & Wilson (2003) find linewidths of 1.5 to 5 km s<sup>-1</sup> (in NH<sub>3</sub>, CS, and isotopes of CO), although deconvolved widths of H<sup>13</sup>CN lines toward individual clumps are $`<`$1 km s<sup>-1</sup> (Lis & Schilke 2003). Closer to the ionization front, Wyrowski et al. (1997) measured linewidths of 2 to 2.5 km s<sup>-1</sup> for C91$`\alpha `$ emission. The radial velocity of the H<sub>2</sub> emission (V<sub>LSR</sub>) is 10-11 km s<sup>-1</sup>, in agreement with published radial velocities for both molecular lines (NH<sub>3</sub>, CS, and isotopes of CO) and carbon recombination lines (Hogerheijde, Jansen, & van Dishoeck 1995; van der Werf et al. 1996; Batrla & Wilson 2003). Peak velocities of the H<sub>2</sub> lines vary by no more than 2 km s<sup>-1</sup>.
### 3.3 Line Ratios and Temperatures
Density estimates for the H<sub>2</sub> emission zone in the Orion Bar are 5-25 $`\times `$ 10<sup>4</sup> cm<sup>-3</sup> derived from C91$`\alpha `$ observations (Wyrowski et al. 1997). At densities in this range, the upper states of the 0-0 S(1), S(2), and S(4) transitions should all be thermally populated (Mandy & Martin 1993; Bertoldi & Draine 1996). Table 4 lists line intensities and excitation temperatures derived from the 0-0 S(1)/0-0 S(2) line intensity ratios (assuming that foreground extinction is negligible and that the lines are optically thin) for the six areas shown in Figure 3 and the excitation temperature derived from the 0-0 S(4)/S(2) line intensity ratios for areas A and B. We use the ratio of the 0-0 S(1)/S(2) intensities, together with the measured 0-0 S(1) line intensity to derive the column density of warm molecular hydrogen toward regions A through F. We assume that the states are populated according to a thermal distribution, which for T $``$ 300 K implies an ortho-to-para ratio of 3. Even if the averaged regions contain cloud material with temperature gradients, these derived, single temperature LTE values provide a way of comparing aggregate properties in the different regions.
Derived excitation temperatures range from roughly 400 to 600 K. Along the bright ridge, the temperature varies between 400 and 500 K. The two areas in Figure 3 illustrate this with the bright part of the ridge (position A) having T\[S(2)/S(1)\]= 460 K while the fainter region along the ridge (position B) has a slightly lower temperature (430 K). The derived column densities at the two ridge positions are very similar ($`9\times 10^{20}`$ cm<sup>-3</sup>, Table 4). T\[S(4)/S(2)\], however, is cooler in position A (503 K) than in position B (572 K). There is an apparent trend (with the exception of position F) toward higher temperatures at greater depth into the molecular cloud (farther from $`\theta ^1`$C).
Our mean temperature for the gas emitting the H<sub>2</sub> rotational lines is similar to that derived by Habart et al. (2004) from ISO observations and by Parmar, Lacy, & Achtermann (1991). We find only a small overall range in 0-0 S(2)/S(1) temperatures and do not see the gradient with distance from $`\theta ^1`$C that was present at modest significance in the Parmar, Lacy, & Achtermann (1991) results. Parmar, Lacy, & Achtermann (1991) present spectra integrated along a 2″$`\times `$ 10″ slit oriented at the same position angle as our map but took data along a different cut through the ionization front. Their spectra sample the H<sub>2</sub> distribution every 5″ perpendicular to the front. Our maps indicate that Parmar et al.’s high temperature point, which lies closest to the ionization front, is very sensitive to relative pointing errors because of the sharp intensity gradients. If the front position is indeed right at the leading edge of the H<sub>2</sub> distribution, the second point in the earlier map may lie behind the brightest part of the H<sub>2</sub> ridge.
Comparing our results to the 1-0 S(1) intensities derived from the maps of van der Werf et al. (1996), we find values of I(1-0 S(1))/I(0-0 S(1)) of 0.42 and 0.44 at positions A and B, respectively. Shifting the 1-0 S(1) line map by 1.5″ produces less than 10% changes in the line intensities at positions A and B. At position A, the v=2-1 S(1)/1-0 S(1) intensity ratio is $`0.14`$ (van der Werf et al. 1996). At the positions behind the ridge, I(1-0 S(1))/I(0-0 S(1)) varies from $`<`$0.07 to 0.65. Note, however, that the intensity calibration of the v=2-1 S(1) results is less reliable than the 1-0 S(1) intensities toward positions A and B. The total area mapped in v=2-1 S(1) by van der Werf et al. (1996) was only 40″ by 40″ . Unlike the mid-IR spectroscopic mapping results, the on and off line wavelengths were observed separately in the near-IR so that zero point offsets could occur that would have a stronger effect on the line calibration at points where the lines are weak. If, as appears to be the case, the intensities were set to zero at the edge of the field mapped in v=2-1 S(1), then any flux on scales greater than 40″ would be absent. This problem is less acute for the v=1-0 S(1) map both because the intensity distribution peaks more sharply and because the van der Werf et al. (1996) map of this line covers a considerably larger area. Usuda et al. (1996) also used Fabry-Perot imaging to determine the 2-1 S(1)/1-0 S(1) ratio. Although their spatial resolution of 8″ may fail to resolve some of the structure, their larger map area results in a more reliable flux ratio. The Usuda et al. (1996) v=2-1 S(1) map implies that the edges of the van der Werf et al. (1996) field lie at about the 50% level of the line intensity distribution. Given the observing and data analysis techniques, the two v=2-1 S(1) results are consistent with one another, but the Usuda et al. (1996) value for the 2-1 S(1)/1-0 S(1) ratio is more reliable. We therefore use their value, $`0.25`$, for the 2-1 S(1)/1-0 S(1) intensity ratio. We discuss the observed H<sub>2</sub> line intensity ratios in the context of a realistic model of the temperature structure in §4.
The formal uncertainties in the LTE temperatures, based on the signal to noise of our ground vibrational state H<sub>2</sub> spectra are small. The maximum uncertainty in temperature was 40 K for the temperature derived from the 0-0 S(1)/0-0 S(2) line ratio in region D (Figure 3). Regions with brighter lines (areas A, B, & E) had much lower uncertainties (4-10 K). Registration and positional errors had a modest effect on the derived average temperatures. We shifted the 0-0 S(1) map by +0.7″ and -0.7″ (greater than our quoted position uncertainty of 0.5″). The maximum change in temperature due to these shifts was 80 K for temperatures derived from the 0-0 S(1)/0-0 S(2) line ratio for the summed spectra in region C. Typical changes in derived temperature from the 0-0 S(1)/0-0 S(2) line ratio due to a comparable position shift were 25 K. At the leading edge of the bright ridge, however, temperatures could change from $``$500 to $``$900 K with relative position shifts of $``$3″.
We do not correct for extinction in our derivation of temperature, as the effects of reddening are expected to be quite low. In Orion OMC-1, Rosenthal, Bertoldi & Drapatz (2000) find values of $`\mathrm{A}_\lambda /\mathrm{A}_\mathrm{K}`$ of 0.525, 0.527, and 0.441 for the 0-0 S(1), S(2), and S(4) lines respectively. Using these values for A<sub>λ</sub>, and assuming A<sub>K</sub>/A<sub>V</sub>=0.112 (Rieke & Lebofsky 1985), we find that an error in T of 50 K corresponds to A$`{}_{\mathrm{V}}{}^{}`$20, for temperatures derived from the 0-0 S(1)/S(2) and 0-0 S(4)/S(2) line intensity ratios.
A single temperature cannot describe our observed line intensities (to within the uncertainties) for the pure rotational lines. The deviation from a single temperature fit becomes more severe as we include vibrationally excited lines. Predictions of the 1-0 S(1) line intensity based on the temperatures and column densities derived from the rotational lines for areas A & B underestimate the observed 1-0 S(1) line intensity by factors of 4 to 75. The enhanced 1-0 S(1) line intensities could be caused by a temperature gradient in the region of H<sub>2</sub> emission or by fluorescent excitation of the vibrational lines. Clearly, more complex analysis of line ratios is necessary. We discuss the inputs and results of our PDR model analysis in the next section.
## 4 Discussion
### 4.1 Inputs to a PDR Model for the Orion Bar
The Orion bar is a particularly useful test site for PDR models because of its closeness (at 450 pc 2″=1.4$`\times 10^{16}`$ cm, Hoogerwerf, de Bruijne, & de Zeeuw 2000), its high gas density and strong incident UV field (leading to bright line and continuum emission), and its distinctive geometry. These favorable properties have inspired a significant number of observational studies over the past two decades which can provide valuable inputs into any model of the region. The key region-specific parameters in modeling photodissociation regions are the strength of the incident UV field and the density of the neutral gas as a function of distance from the front. In comparing the models to observations, it is also important to understand the source geometry, in particular the tilt of the source with respect to the line of sight.
The usual practice in modeling photodissociation regions has been to assume that the dust-related parameters in the models: dust opacity, photoelectric heating rate, and H<sub>2</sub> formation rate, are fixed and to search through a family of models in an attempt to find the best match when the incident UV field and the density (often taken to be uniform) of the PDR are varied (Burton, Hollenbach, & Tielens 1990). Because the UV field incident on the Orion Bar and the pressure at the ionized boundary of the PDR are well constrained by observations, we can reverse the usual procedure, take the UV field and pressure as givens, and use Orion as a testbed for the dust-related parameters going into models of high density PDRs with high incident UV fields. In a companion paper (Draine et al. 2005, Paper II), we discuss this model study in detail. In the current work, we discuss the derivation of the initial conditions from observations and compare the results of our observational study to the best-fit PDR model derived in Paper II.
The dominant source of UV photons for the Orion Bar PDR is the O6 star $`\theta ^1`$C Ori. At the position of our H<sub>2</sub> map, $`\theta ^1`$C Ori lies 120″ from the ionization front where the direction of the incident radiation is only 20 degrees from the normal to the front projected onto the plane of the sky. It is convenient to measure the intensity of the radiation incident at the ionization front by $`\chi `$, the ratio of specific energy density at 1000 Å to the value $`u_\lambda `$= 4$`\times 10^{17}`$ ergs cm<sup>-3</sup> Å<sup>-1</sup> estimated by Habing et al. (1968) for the mean interstellar radiation field. Based both on simple geometric dilution and on the strength of the far-IR radiation emitted by the warm dust in the Orion bar after being heated by the UV and visible radiation from $`\theta ^1`$C Ori, the far-UV flux incident on the PDR is $``$3$`\times `$10<sup>4</sup> times the mean interstellar radiation field (Herrmann et al. 1997). Marconi et al. (1998) used observations of the OI 1.317 $`\mu `$m line to infer the incident UV intensity at 1040 $`\AA `$. If the PDR is inclined with cos$`\theta `$= 0.1, where $`\theta `$ is the angle between the line of sight and normal to the PDR, the Marconi results imply $`\chi `$= 2.9$`\times `$10<sup>4</sup>. For all of the models explored in Paper II, we take $`\chi `$= 3$`\times `$10<sup>4</sup> at the ionization front.
$`\theta ^1`$ Ori C, A, and E combined have a 1–10 keV luminosity $`L_X(110\mathrm{keV})=2.4\times 10^{32}\mathrm{erg}\mathrm{s}^1`$ (Schulz et al. 2003). At a distance of $`8\times 10^{17}`$ cm distance, 1–10 keV X-rays from the Trapezium stars will contribute an ionization rate $`10^{16}`$ s<sup>-1</sup> at the Bar. Lower energy X-rays from the Trapezium, and X-rays from young stars that are less luminous but are closer to the Bar will contribute additional ionization. In addition, the cosmic ray ionization rate may be enhanced in this region by nonthermal particle acceleration in stellar wind shocks. We adopt a nonthermal ionization rate $`\zeta _{CR}1\times 10^{15}`$ s<sup>-1</sup> for gas in the Orion Bar PDR, but we stress that our results do not depend sensitively on this rate. We note that McCall et al. (2003) inferred an ionization rate $`1.2\times 10^{15}`$ s<sup>-1</sup> in the molecular gas toward $`\zeta `$ Per.
Based on the emission measure derived from radio continuum observations (Felli et al. 1993), and $`n_e`$ determined from \[S II\]$`I(6716)/I(6731)`$ (Pogge et al. 1992), the thermal pressure at the ionization front $`nT6\times 10^7`$ cm<sup>-3</sup> K. The gas has been accelerated away from the PDR, however, so the pressure in the PDR should be somewhat larger. For the PDR models of the Orion Bar, we have taken the total pressure to be uniform through the PDR at $`P/k`$ = 8$`\times `$10<sup>7</sup> cm<sup>-3</sup> K. In most previously published PDR models dealing with H<sub>2</sub> excitation (Black & van Dishoeck 1987; Sternberg & Dalgarno 1989; Burton, Hollenbach, & Tielens 1990), the assumption has been that turbulent or magnetic pressure dominate the gas pressure throughout the region and that an assumption of constant density is therefore reasonable. The Draine et al. (2005) models explicitly calculate density as a function of depth into the PDR assuming a constant pressure that includes a non-thermal contribution<sup>1</sup><sup>1</sup>1The non-thermal pressure is taken to be $`p_{\mathrm{nt}}=\rho v_{\mathrm{nt}}^2`$, with $`v_{\mathrm{nt}}=1\mathrm{km}\mathrm{s}^1`$. that is fairly small in the outer parts of the region. In the outer parts of the PDR, constant densities derived from older PDR models and various observations are generally consistent with the pressure value used in our models, albeit with a large spread. The densities derived by Wyrowski et al. (1997) from the C91$`\alpha `$ results (5–25 $`\times `$ 10<sup>4</sup> cm<sup>-3</sup>) are consistent with the assumed pressure if T$``$1000 K in the carbon line formation region. Non-LTE excitation modeling of millimeter and submillimeter molecular line ratios is not dependent on PDR model results and yields densities ranging from a few 10<sup>5</sup> cm<sup>-3</sup> to a few 10<sup>6</sup> cm<sup>-3</sup> (Burton, Hollenbach, & Tielens 1990; Tauber et al. 1995; Hogerheijde, Jansen, & van Dishoeck 1995; Young Owl et al. 2000) in the gas farther behind the ionization front where molecular line ratios and brightness temperatures imply $`T120`$ K. The higher densities are in good agreement with a virial analysis of the brightest HCN clumps in the molecular ridge (Lis & Schilke 2003).
The pressure assumed for our models is also consistent with densities derived by a second line of argument, based on geometry and the chemical stratification shown in older constant density models of PDRs. There is a clear stratification of emission zones, manifested by a shift in the observed location of the peak emission, as one goes farther from $`\theta ^1`$C into the PDR, albeit with some overlap of what should be, from a theoretical point of view, distinct regions within the PDR (Figure 1). Along a line perpendicular to the bright ridge in Figure 2, the ionization front lies at $`y17`$″ (Felli et al. 1993). Within the neutral gas, there are successive emission zones for the FeII 1.64 $`\mu `$m line (Marconi et al. 1998, $`y16`$″), the 3.3 $`\mu `$m PAH feature (Bregman et al. 1994, $`y12`$″), the H<sub>2</sub> rotational and ro-vibrational transitions (this paper, $`y2`$″), submillimeter continuum emission (Lis et al. 1998, $`y+5`$″), and various millimeter and submillimeter lines of HCO<sup>+</sup>, CO, and HCN (Tauber et al. 1994; Hogerheijde, Jansen, & van Dishoeck 1995; Young Owl et al. 2000; Lis & Schilke 2003, $`y+8`$″). The peak of the H<sub>2</sub> rotational line emission in Figure 2 is at $`y2`$″, or $``$15″ ($``$9$`\times 10^{16}`$cm) from the ionization front. By comparing the linear displacements of these peaks to the column density peaks in plane-parallel PDR models with appropriately chosen incident UV fields, Tielens et al. (1993); Tauber et al. (1994) and Simon et al. (1997) estimate the density for a homogeneous medium to be 5$`\times `$10<sup>4</sup> to 3$`\times 10^5`$ cm<sup>-3</sup>, consistent with the range of densities derived from physical measurements and with an appropriate density range for our constant pressure models.
Many authors have suggested that propagation of far-UV radiation through a clumpy medium offers an explanation for the range in derived densities, particularly in the molecular part of PDR’s. In such a picture, the attenuation scale length in the low density medium or the distance to reach a clump area filling factor of unity set the size scale for the PDR (Stutzki et al 1988). There is, however, no unambiguous evidence of a clumpy structure in the region of the Orion Bar where H<sub>2</sub> emits. The fairly straight leading edge of the bright H<sub>2</sub> ridge and the absence of strong variations in either column density or in the ratio of 1-0 S(1) line to 0-0 S(1) line intensity along the ridge argue that the H<sub>2</sub> emission arises either from a uniform medium, from a low density PDR component, or from an ensemble of clumps with a line of sight filling factor significantly larger than one. Clumps with densities of $``$6 $`\times `$ 10<sup>6</sup> cm<sup>-3</sup> have been observed deep into the molecular gas behind the PDR via HCN, HCO and their isotopomers (Lis & Schilke 2003; Young Owl et al. 2000). However, for the neutral gas closest to the ionization front, Marconi et al. (1998) use the relative strengths of near-IR Fe II lines to exclude the presence of clumps with densities $`>`$10<sup>6</sup> cm<sup>-3</sup>. In our models (Paper II), we therefore assume a uniform medium.
Apart from the temperature within the H<sub>2</sub> bright ridge of $``$450 K, there are several other temperature measurements that can serve as constraints on the thermal balance through the PDR. If the carbon recombination lines arising from the part of the PDR closest to the ionization front are purely thermally broadened, the temperatures in that layer are 1000-1600 K (Wyrowski et al. 1997). CO 6-5 brightness temperatures $``$10″ in front of the molecular peak ($``$5″ behind the H<sub>2</sub> ridge) imply a kinetic temperature of 120-180K (Lis et al. 1998). Farther into the cloud, Batrla & Wilson (2003) use NH<sub>3</sub> line ratios to derive a temperature of 120 K for what they argue are the surfaces of the dense clumps in the molecular ridge.
In comparing the PDR model intensities to observed intensities, we must tilt the models correctly with respect to the line of sight and account for radiative transfer through the inclined PDR slabs. There are two lines of evidence that argue that the Orion bar is highly inclined from the plane of the sky. Molecular line observers, who see low-level emission both in front of and behind the bar, conclude that the roughly factor of 10 enhancement in column density seen for optically thin lines from the bar implies an inclination of only a few degrees from the line of sight (Tauber et al. 1994; Hogerheijde, Jansen, & van Dishoeck 1995). The steepness of the dropoff in radio continuum at the ionization front and the steep rise in the rotational H<sub>2</sub> line emission on the leading edge of the ridge form the second argument for the almost edge-on orientation of the PDR (Felli et al. 1993, this paper). Walmsley et al. (2000) have estimated that the Orion bar is approximately plane-parallel and that we view it from a direction with 1/cos$`\theta `$ = 10 where $`\theta `$ is the angle between the line-of-sight and the normal to the PDR. In calculating models of the emergent line intensities, we adopt this as a plausible estimate for the enhancement of the surface brightness of optically thin lines relative to the face-on surface brightness. Accordingly, in all of the models in Paper II, Draine et al. attempt to reproduce the line intensities observed at positions A and B with a plane-parallel PDR viewed from an angle such that 1/cos$`\theta `$ = 10. The model line intensities include attenuation by dust within the PDR.<sup>2</sup><sup>2</sup>2The dust attenuation cross section is taken to be $`\sigma _\lambda =(A_\lambda /A_{1000\mathrm{\AA }})\sigma _{1000}`$, where $`A_\lambda `$ is the extinction at wavelength $`\lambda `$ for an $`R_V=5.5`$ extinction law. For 1/cos$`\theta `$ = 10, internal extinction significantly attenuates the emission at the shorter wavelengths. For example, in the PDR model discussed below the 1-0 S(1) line is attenuated by a factor of 0.26, the equivalent of A<sub>K</sub> = 1.4 . The lower surface brightness seen at positions C-F may be due to viewing with a different inclination angle (Hogerheijde, Jansen, & van Dishoeck 1995), or possibly this emission arises from a region physically separate from the location of the edge-on Bar.
### 4.2 Results for the Best PDR Model
We present here a comparison of our new observational results for the Orion Bar and a theoretical model for this high-excitation PDR. This model may be relevant not only for a single dense PDR illuminated by O and B stars in galactic star forming regions but also for studies of physical conditions over large areas in the inner regions of starburst galaxies. The 0-0 S(2)/0-0 S(1) line ratio for the nucleus of NGC 253 (Devost et al. 2004) is the same as observed in the Orion Bar and the intensity averaged over an 800 pc $`\times `$ 700 pc region of NGC 253 is fully 50% of that observed at peak A in the Orion Bar. If the emission in NGC 253 originates from PDRs, these must be both very intense and have a high surface filling factor.
In paper II, Draine et al. present a grid of models for the Orion Bar PDR near position A. Table 5 gives the values of the pressure ($`P`$), radiation intensity ($`\chi `$), rate of ionization of H by cosmic rays or X-rays ($`\zeta _{CR}`$), and abundances of the coolants C, O, Si, and Fe used in the models. Gas-phase abundances for C, O, Si, and Fe are taken from Jenkins (2004) for gas with “depletion factor” $`F=1`$, corresponding to approximately the level of depletion seen in the diffuse molecular cloud toward $`\zeta `$ Oph. The vibrational line emission from the models is sensitive to the rate coefficients for vibrational deexcitation of H<sub>2</sub>, particularly by collisions with H atoms. Usuda et al. (1996) found that in the Orion Bar, the 2-1 S(1)/1-0 S(1) intensity ratios were anticorrelated with the 1-0 S(1) line intensities, and were usually lower than the 2-1 S(1)/1-0 S(1) intensity ratio ($``$0.6) expected for pure UV-fluorescence. Our values adopted for these rates are discussed by Draine & Bertoldi (2005, in preparation). Table 6 compares different estimates for the $`T=1000`$ K rate coefficients, $`k_{\mathrm{vdexc}.}(v,J)`$, for vibration deexcitation by H atom collisions of the $`(v,J)=(1,3)`$ and $`(2,3)`$ levels of H<sub>2</sub> (the levels responsible for 1–0S(1) and 2–1S(1) line emission). Our adopted rates are an order of magnitude smaller than the vibrational deexcitation rates adopted by Sternberg & Dalgarno (1989), but exceed the rate coefficients calculated by Mandy & Martin (1993), by factors of 8 and 2, respectively. Our rates are a factor of 150 times larger than the rates recommended by Le Bourlot et al. (1999).
The model grid explores variations in the dust ultraviolet attenuation cross section, the H<sub>2</sub> formation rate, and the photoelectric heating rate. Table 7 compares a model from this grid (Model 1) to the observations at position A. The model is within $``$10% of the 0-0 S(1), 0-0 S(2), 0-0 S(4), and 1-0 S(1) intensities for the average of positions A and B. The model 0-0 S(1) intensity is $``$2% below 0-0 S(1) at position A, and $``$9% above the value at B. For 0-0 S(2), the model is $``$18% below A and $``$1% below B. For 0-0 S(4), the model is $``$7% above A, and $``$14% below B. For 1-0 S(1), the model is $``$6% above A, and $``$12% above B.
The 2-1 S(1)/1-0 S(1) intensity ratio for the model is 0.23. There is some uncertainty concerning the observed line ratio. As discussed in §3.3, we use the results of Usuda et al. (1996), 2-1 S(1)/1-0 S(1) = 0.25 at position A, which we take as the best observational determination.
Figure 6 shows the temperature profile for Model 1. The ionization front is defined to be the point where $`n(\mathrm{H}^+)=n(\mathrm{H}^0)`$; at this point, the gas temperature is $``$9000 K, but the temperature drops rapidly with distance from the ionization front, as heating due to photoionization of H declines and the fractional ionization drops. Model 1 successfully reproduces the observed $`9\times 10^{16}`$ cm separation between the ionization front and the peak of the H<sub>2</sub> line emission (Figure 6). In fact, the figure shows that the model even has an extended (2-3″) tail on the ionization front side of the 1-0 S(1) peak, consistent with the observed cuts shown in Figure 4. Most of the 1-0 S(1) emission in this model arises from collisional excitation of (1,3), the $`v=1,J=3`$ state. For example, at $`RR_{\mathrm{IF}}=6.0\times 10^{16}`$ cm, $`n(1,3)/n(0,3)0.023`$, which is essentially the thermal ratio ($`e^{5936\mathrm{K}/T}`$) at the local temperature $`T=1600`$K. The density is not high enough to fully thermalize the vibrational levels – ultraviolet pumping contributes in part to the population of (1,3), and accounts for most of the population of (2,3). The rise in $`n(1,3)`$ to a local maximum at $`8.5\times 10^{16}`$ cm reflects competition between increasing $`n(\mathrm{H}_2)`$ and declining $`T`$. The second maximum at $`9.2\times 10^{16}`$ cm is due to UV-pumping: the decline in $`T`$ and the drop in $`n(\mathrm{H})/n_\mathrm{H}`$ lead to a drop in the rates for collisional deexcitation of the vibrationally-excited states, so that the $`v=1`$ levels are no longer collisionally deexcited, but going deeper into the cloud, the UV pumping rates drop and therefore so does $`n(1,3)`$. Note this maximum of $`n(1,3)`$ coincides with the maximum in $`n(2,3)`$.
The rotational levels of $`v=0`$ are thermalized and arise in the zone where the temperature gradient is quite steep. The lower $`J`$ levels \[e.g., (0,2)\] peak farther from the ionization front than the higher $`J`$ levels \[e.g., (0,6)\].
### 4.3 Parameters for the Best PDR Model
The best-fit model from Paper II (Model 1) uses a rate coefficient for formation of H<sub>2</sub> on dust grains <sup>3</sup><sup>3</sup>3With the usual definition: $`(dn(\mathrm{H}_2)/dt)_{\mathrm{form}}=R_{\mathrm{H}_2}n_\mathrm{H}n(\mathrm{H})`$. with a value $`R_{\mathrm{H}_2}=3.8\times 10^{17}`$cm$`^3`$s<sup>-1</sup> at $`T=1000`$ K – similar to the value $`3\times 10^{17}\mathrm{cm}^3\mathrm{s}^1`$ found by Habart et al. (2004) for the Orion Bar PDR. However, Model 1 implies significant deviations from the standard values adopted for other dust-related parameters in the Orion Bar PDR.
In order to achieve the agreement in the separation of the ionization front and H<sub>2</sub> peak, Model 1 adopts a dust attenuation cross section at 1000Å $`\sigma _{1000}=0.48\times 10^{21}`$cm<sup>2</sup> – if a significantly higher value of $`\sigma _{1000}`$ is used, the increased FUV attenuation brings the H<sub>2</sub> peak too close to the ionization front. The adopted $`\sigma _{1000}`$ is significantly smaller than the $`\lambda =1000`$Å extinction cross section $`2.3\times 10^{21}`$cm<sup>2</sup> inferred from the Fitzpatrick (1999) parametrization of the interstellar reddening law for sightlines with $`R_VA_V/E(BV)5.5`$, if we take $`N_\mathrm{H}/E(BV)5.8\times 10^{21}\mathrm{cm}^2`$ from Bohlin, Savage & Drake (1978). At 1000Å the dust albedo is estimated to be $`0.4`$ (Draine 2003; Gordon 2004), implying an attenuation cross section $`0.6\times 2.3\times 10^{21}\mathrm{cm}^2=1.4\times 10^{21}\mathrm{cm}^2`$ – 3 times larger than the value adopted for Model 1. The dust now in the Orion PDR might have undergone extensive coagulation during the long time it spent in cold, dense molecular gas prior to the arrival of the photodissociation front. Such coagulation would lower the far-ultraviolet scattering and absorption per H nucleon.
As described above, the PDR model corrects for internal absorption in the PDR assuming a $`R_V=5.5`$ reddening law, and therefore the K band attenuation coefficient has been scaled down by the same factor of $`3`$ as the UV extinction. This would be appropriate if the reduced extinction were due to an overall deficiency of dust grains, but it would not be correct if the low UV extinction were due to dust coagulation, as coagulation of small grains would not decrease the K band extinction unless the coagulation resulted in grains larger than $`1\mu `$m. As noted above, even the reduced extinction assumed in the model has attenuated the 1-0 S(1) line intensity by a factor 0.26 because we assume that we are observing the PDR from a direction with $`1/\mathrm{cos}\theta =10`$ – if the K band attenuation coefficient were significantly larger than the (reduced) value in the model, it would be very difficult to reproduce the observed H<sub>2</sub> line intensities. The degree to which the observed line intensities have been affected by extinction in the Orion Bar is an important question; additional observational studies of the reddening using H<sub>2</sub> emission lines would be of great value.
In order to lower the 2-1 S(1)/1-0 S(1) line ratio from the pure fluorescence value $``$0.6 to the observed value $``$0.25, the atomic zone of the PDR must have a gas temperature $`T1000`$K – this is required so that (1) the rate coefficients for collisional deexcitation are large enough to suppress 2-1 S(1) emission by collisionally deexciting H<sub>2</sub> in the $`v=2,J=3`$ state fast enough to compete with spontaneous decay, and (2) to collisionally excite 1-0 S(1) emission. Although the models explicitly include heating from collisional deexcitation of vibrationally-excited H<sub>2</sub>, the dominant heating process is photoelectric heating from dust. With the cooling processes that are present, the only way to produce the required high temperature is for the heating rate to be substantially larger than the photoelectric heating rate predicted by existing models of photoelectric emission from dust (Bakes & Tielens 1994; Weingartner & Draine 2001). Draine et al. (2005) provide this additional heating by means of an ad-hoc increase in the photoelectric heating rate in the atomic portion of the PDR by a factor $`3`$ relative to the estimate of Weingartner & Draine (2001, hereafter WD01) for $`R_V=5.5`$ dust.
Above it has been argued that the separation of the peak of the H<sub>2</sub> emission from the PDR requires a reduction in the FUV absorption by the dust; since photoelectric heating cannot occur without absorption of UV photons, one might have expected a corresponding reduction in the dust photoelectric heating rate, whereas Model I posits an increased heating rate. The increased photoelectric heating rate may be regarded as a proxy for some other heating process that may be present, or perhaps it is indicative of overestimation of the fine structure cooling (dominated by \[OI\]63$`\mu `$m and \[CII\]158$`\mu `$m emission). In any event, it indicates that there is a substantial error in our account of the heating and cooling in the atomic zone of the PDR.
However, although high temperatures, and therefore an enhanced photoelectric heating rate or its effective equivalient, are required to suppress 2-1 S(1) emission in the region where UV pumping of H<sub>2</sub> is taking place, this enhanced photoelectric heating rate cannot be present in the regions that are predominantly molecular – otherwise there would be too much emission in 0-0 S(4), 0-0 S(2), and 0-0 S(1). Draine et al. (2005) therefore adopt an ad-hoc photoelectric heating rate that is reduced to $`0.4`$ of the WD01 heating rate where $`2n(\mathrm{H}_2)/n_\mathrm{H}=0.5`$, and $``$0.1 of the WD01 heating rate where $`2n(\mathrm{H}_2)/n_\mathrm{H}=0.9`$. Such variation in the grain photoelectric heating properties could perhaps come about if the grains from the cold dense molecular cloud enter the PDR in some state (perhaps coated or clumped) yielding a low photoelectric heating rate. As these grains enter the PDR and are exposed to both the $`\lambda <1100`$Å radiation radiation that dissociates H<sub>2</sub> and atomic H, perhaps the grain properties are altered (e.g., dispersal of clumps, or photolysis of coatings) so as to increase the photoelectric yields.
There are no velocity-resolved observations of 1-0 S(1) toward the Orion Bar. In the PDR associated with the reflection nebula NGC 7023, this line has a width of 3.4 km s<sup>-1</sup> (Lemaire et al. 1999), consistent with our measurements of the ground state lines in the Orion Bar. The narrow linewidths of our pure rotational H<sub>2</sub> lines indicate that the gas, if shocked, must be shocked at a very low velocity. Observed linewidths in regions with even moderate (v<sub>s</sub> $``$ 20 km s<sup>-1</sup>) shock velocites are greater than $``$ 30 km s<sup>-1</sup> (Parmar, Lacy, & Achtermann 1994; Tedds, Brand, & Burton 1997). Shock models of Draine, Roberge, & Dalgarno (1983) and Kaufman & Neufeld (1996) predict H<sub>2</sub> v=1-0 S(1) intensities greater than $`10^4`$ ergs s<sup>-1</sup> cm<sup>-2</sup> sr<sup>-1</sup> in face on PDRs from shocks with velocities greater than 20 km s<sup>-1</sup>. Tielens et al. (1993) calculate the heating input by a shock in the Orion Bar and find that heating of the gas by the FUV field (65 erg cm<sup>-2</sup> s<sup>-1</sup>) exceeds shock heating unless the shock velocity is greater than 10 km s<sup>-1</sup>. Given the morphological similarities of the 1-0 S(1) and 0-0 S(1) emission, it is also unlikely that shock excitation could appreciably contribute to the observed 1-0 S(1) intensity.
In addition to explaining the PDR structure and H<sub>2</sub> line intensities, future modeling must address the high temperatures in the CO/HCO<sup>+</sup>/NH<sub>3</sub> zone of the Orion Bar. The models in Draine et al. (2005, in prep) do not calculate the thermal balance realistically that far into the PDR. Figure 6 shows, however, that the model temperature has already dropped to 50 K, far below the temperatures of 100-120 K derived from observations of the molecular zone, even before that zone is reached. Clearly, other heating mechanisms must be in play within that zone as well. We emphasize that our models of the Orion Bar are optimized for the region in the PDR where H<sub>2</sub> emits, and do not apply to regions deeper into the molecular cloud.
## 5 Summary
We obtained high resolution (R = 75,000 to 100,000) spectral maps of H<sub>2</sub> v = 0-0 S(1) and S(2) covering a 12″ by 40″ region in the Orion Bar PDR. Linewidths for the spectra in our maps are 4-6 km s<sup>-1</sup> with V<sub>LSR</sub> ranging from 10.2 to 11.5 km s<sup>-1</sup>. Comparison of our maps with v = 1-0 S(1) observations (van der Werf et al. 1996) reveals exceptional similiarity in the line intensity distributions.
To model our line intensities (detailed in Draine et al. 2005, in prep), we use estimates of the FUV field, pressure and inclination angle from the literature (§4.1), and allow flexibility in dust-related parameters (dust opacity, photoelectric heating rate and H<sub>2</sub> formation rate). The best-fit model matches the distance between the H<sub>2</sub> line emission and the ionization front, the observed intensities of the H<sub>2</sub> $`v=10`$ S(1) and $`v=21`$ S(1) lines, and the intensities of the ground state rotational lines.
In order to reproduce the observed separation between the ionization front and the H<sub>2</sub> emission peak, the model requires a reduction in the FUV attenuation cross section, by a factor of $``$3 relative to a priori estimates. This model also requires an enhanced heating rate in the atomic region of the PDR, corresponding to a factor $`3`$ increase in the photoelectric heating rate (or a corresponding reduction in the radiative cooling) in order to maintain $`T1500`$K in the atomic zone. Though the uniqueness of our solution has not been tested, it is apparent that the standard dust-related parameters used in PDR models do not allow for a reasonable match to our observations.
The spatial resolution of our observations (0.002 pc) is roughly the thickness of the H<sub>2</sub> emission region, according to our best-fit model. Thus, we were not able to spatially resolve the temperature structure of the PDR, and our observed width and steepness of the bright ridge is due primarily to inclination effects. Line intensities and ratios for extended emission behind the main H<sub>2</sub> ridge can be explained if 1/cos($`\theta `$) decreases as one moves further into the cloud. This geometry would agree with Hogerheijde, Jansen, & van Dishoeck (1995).
Acknowledgements: We would like to thank Tommy Greathouse for support of the observing, Paul van der Werf for providing his results in electronic form, and David Hollenbach along with our anonymous referee for helpful comments. This work was supported by NSF grant AST-0205518. BTD was supported in part by NSF grant AST-9988126 and in part by grants from the W.M. Keck Foundation and the Monell Foundation. MJR was supported by NSF grant AST-0307497 and NASA grant NNG04GG92G. TEXES was built and operated by grants from the National Science Foundation and the Texas Advanced Research Program. |
warning/0506/quant-ph0506243.html | ar5iv | text | # Contents
## Chapter 1 Introduction
Despite the unsurpassed predictive success of quantum theory, there is, since its inception almost 80 years ago, a persistent problem with its conventional interpretation, namely the measurement problem.
The problem arises as follows. Quantum theory was developed in order to explain the behavior of ‘microscopic’ systems. With each microscopic system, quantum theory associates a wavefunction $`\psi `$. According to the conventional interpretation of quantum theory,<sup>1</sup><sup>1</sup>1With the ‘conventional interpretation’ we mean the Dirac-von Neumann approach which can be found in most standard textbooks. this wavefunction provides the most complete specification of the microscopic system. Further, the dynamics of the wavefunction $`\psi `$ is governed by two different laws. First, there is the dynamical evolution according to the Schrödinger equation, which is deterministic. Given the initial wavefunction one can uniquely determine the wavefunction at a later time. There is also another type of evolution of the wavefunction, which is the collapse of the wavefunction. The collapse rule is introduced in quantum theory in order to explain the definite outcome that is obtained when a measurement is performed. In this respect, the collapse of the wavefunction is said to occur when a measurement is performed by a ‘macroscopic observer’ (human or not) on the ‘microscopic system’ described by this wavefunction. The result is a replacement of the wavefunction $`\psi `$ by another wavefunction which from that time on provides the (complete) description of the microscopic system. Contrary to the dynamical law given by the Schrödinger equation, the collapse law is not deterministic.
Considered separately both laws of dynamical evolution are unambiguously defined. On the other hand it is unclear what exactly is meant by a ‘microscopic’ and a ‘macroscopic’ systems, or what exactly is meant by an ‘observer’ and a ‘system’. Hence it is unclear when the wavefunction evolves according to which of the two dynamical laws. The ambiguity becomes most striking in the following example. When a macroscopic observer performs a measurement on a microscopic system the collapse law should apply. But if the macroscopic system is regarded as a collection of microscopic systems, then the wavefunction of the total system, which consists of the observer and the system under observation, should evolve in time according to the Schrödinger equation and the collapse law should not be invoked. It is obvious, that these two ways of describing the measurement process are mutually incompatible if the wavefunction is to be regarded as the most complete specification of the system.
In practical situations the difference between the ‘macroscopic observer’ and the ‘microscopic system’ is of course sufficiently large so that one can often say with certainty whether or not the collapse has occurred. Nevertheless the ambiguous distinction between the ‘macroscopic observer’ and the ‘microscopic system’ presents an obvious logical flaw which is intolerable if one wants to regard quantum theory as a fundamental theory describing Nature. Because the ambiguous distinction is needed for the collapse law, and because the collapse law is invoked to describe the measurement process, the problem is generally referred to as the measurement problem.
A possible resolution for the measurement problem resides in the view that the complete specification of a microscopic system is not only provided by the wavefunction, but also by some extra variables.<sup>2</sup><sup>2</sup>2Of course this is not the only way in which the measurement problem can be solved. One could for example also adopt an approach in which the wavefunction is dismissed altogether in the description of a quantum system, or an approach where the wavefunction still gives the complete description of a quantum system, but where the Schrödinger equation is modified, as in spontaneous collapse models (for a review see ). However, such theories will not be dealt with in the thesis. These extra variables should have an objective existence, irrespective of the fact whether or not a measurement is performed. They should also determine the outcome in experiments, so that the collapse law becomes superfluous. There is then no distinction needed between microscopic and macroscopic systems; both are described by these extra variables together with the wavefunction. A theory in which the system is described by such additional variables is called a realistic theory. If this realistic theory accounts for the same empirical predictions as quantum theory, it is also called an interpretation of quantum theory. The extra variables are usually termed hidden variables. However, because the reason for introducing these hidden variables is usually to give a definite account for the outcome in experiments, the term ‘hidden variables’ is a kind of misnomer. For this reason Bell preferred to term these extra variables as beables . This is a term which we shall use frequently further on.
An example of such a theory was presented by Louis de Broglie in 1927 at the Solvay Congress in Brussels (cf. , and references therein). De Broglie called his theory the pilot-wave theory. In fact, de Broglie regarded his pilot-wave theory only as a truncated version of his theory of the double solution which he had been working on since 1923, the time he proposed the idea of associating wave properties to massive particles, the key idea which led to quantum theory. However, because of the unfavorable reception of his pilot-wave approach at the Solvay Congress and because of objections raised by Pauli, de Broglie abandoned his ideas. It was only after David Bohm reinvented the ideas of pilot-wave theory in 1952 (although from a different perspective) that de Broglie returned to his original ideas and that he was able to answer Pauli’s criticism.
In de Broglie’s pilot-wave theory the description of a quantum system by means of the wavefunction is extended by considering point particles which follow definite trajectories. The velocity field of these particles is fully determined by the wavefunction. Given the initial positions of the particles, their trajectories are fully determined by this velocity field. In this sense the particles are ‘piloted’ by the wavefunction, hence the name pilot-wave theory. With an ensemble of quantum systems (all described by the same wavefunction) there corresponds a distribution of the actual positions of the particles. With a particular assumption on the initial distribution (i.e. the particles should initially be distributed according to the quantum distribution), pilot-wave theory reproduces the quantum probabilities for the ensemble. Hence, with this assumption pilot-wave theory can be considered as an interpretation of quantum theory. The hidden variables or the beables in pilot-wave theory are the particle positions.
In order to make the distinction between the notion of a particle in the standard interpretation of quantum theory and the notion of a particle in the theory of de Broglie and Bohm, we will term the latter the particle beable. Note that Bell himself used the term ‘beables’ to refer to the particle positions instead of to the particles themselves , however in the literature the notion of beable is often extended to cover both interpretations.
Not only does the pilot-wave theory provide us with a logically unambiguous theory because it is devoid of the measurement problem, it also provides a clear picture of what the theory is about . The standard interpretation is not so clear about what physical entities are associated with the mathematics. The standard interpretation is certainly not about particles, because the most complete description is given by the wavefunction. Pilot-wave theory is unambiguous in this respect. In pilot-wave theory, matter is built of point-particles, the particle beables, moving in three dimensional ‘physical’ space and these particles are causally influenced by the wavefunction which is grounded in configuration space.
To give credit to both of its inventors we will term the theory of de Broglie and Bohm, the de Broglie-Bohm pilot-wave theory or for short the pilot-wave theory . In the literature other names for the theory can be found. Although these different names may also carry nuances in the interpretation of the theory, the basic mathematical structure is the same. For example Bohm and Hiley and Holland refer to the theory as the ontological or causal interpretation. The group around Dürr, Goldstein and Zanghì prefers to call the theory Bohmian mechanics.
The principles of the de Broglie-Bohm pilot-wave formalism are most easily sketched in the case of non-relativistic quantum theory. We do this in the next section. In the following section we then discuss in detail how the pilot-wave interpretation solves the measurement problem. We end the introductory chapter with an outline of the thesis.
### 1.1 The pilot-wave interpretation
The standard quantum mechanical description of a system of $`N`$ spinless particles is given by means of a wavefunction $`\psi (𝐱_1,\mathrm{},𝐱_N,t)`$ in configuration space $`\mathrm{}^{3N}`$, which satisfies the non-relativistic Schrödinger equation
$$i\mathrm{}\frac{\psi (𝐱_1,\mathrm{},𝐱_N,t)}{t}=\left(\underset{k=1}{\overset{N}{}}\frac{\mathrm{}^2_k^2}{2m_k}+V\right)\psi (𝐱_1,\mathrm{},𝐱_N,t),$$
(1.1)
with $`m_k`$ the mass of the $`k^{\text{th}}`$ particle and $`V(𝐱_1,\mathrm{},𝐱_N)`$ a potential.
In the standard quantum interpretation the wavefunction is used to calculate detection probability distributions for observables. In particular, the probability density to make a joint detection of the $`N`$ particles at the configuration $`(𝐱_1,\mathrm{},𝐱_N)`$ at a particular time $`t`$ is given by $`|\psi (𝐱_1,\mathrm{},𝐱_N,t)|^2`$. The continuity equation for this distribution is given by
$$\frac{|\psi |^2}{t}+\underset{k=1}{\overset{N}{}}\mathbf{}_k𝐣_k=0,$$
(1.2)
with
$$𝐣_k=𝐯_k|\psi |^2$$
(1.3)
the 3-vector probability current and
$$𝐯_k=\frac{\mathrm{}}{2im_k|\psi |^2}\left(\psi ^{}\mathbf{}_k\psi \psi \mathbf{}_k\psi ^{}\right)=\frac{\mathrm{}}{m_k}\text{Im}\frac{\mathbf{}_k\psi }{\psi }.$$
(1.4)
The continuity equation expresses the conservation of the detection probability distribution $`P`$.
In the pilot-wave interpretation , the $`N`$-particle wavefunction $`\psi `$ is not regarded as providing the complete description of a quantum system. One also assumes the existence of $`N`$ point particles (the particle beables) which have definite positions at all times in physical space $`\mathrm{}^3`$. If we represent the position of the $`k^{\text{th}}`$ particle beable with the 3-vector $`𝐗_k`$, then the trajectories $`𝐗_k(t)`$ are solutions to the differential equations
$`{\displaystyle \frac{d𝐗_k}{dt}}`$ $`=`$ $`𝐯_k(𝐗_1,\mathrm{},𝐗_N,t)`$ (1.5)
$`=`$ $`{\displaystyle \frac{\mathrm{}}{m_k}}\text{Im}{\displaystyle \frac{\mathbf{}_k\psi (𝐱_1,\mathrm{},𝐱_N,t)}{\psi (𝐱_1,\mathrm{},𝐱_N,t)}}|_{𝐱_j=𝐗_j}.`$
In this way the wavefunction acts as a guiding wave, the pilot-wave, which governs the motion of the particle beables; there is no back-reaction of the particles onto the wavefunction. The equations (1.5) are called the guidance equations.<sup>3</sup><sup>3</sup>3In fact Bohm presented the pilot wave interpretation as a second order formalism . This involved a Newtonian-like force law for the particle beable, including an extra potential, the quantum potential. The reason for Bohm’s preference for this second order formalism rests in his observation that the Schrödinger equation could be written, by separation of real and imaginary parts, as a Hamilton-Jacobi-like equation together with the continuity equation. In this thesis we adopt the view held by de Broglie , in which the guidance law (1.5) is regarded as the fundamental dynamical equation for the particle beables. Amongst the main advocates of this view are Bell , Dürr et al. and Valentini . Nevertheless, because of the close connection to the classical Hamilton-Jacobi formulation, the second order formulation may be a valuable aid in the study of the emergence of classical mechanics out of quantum theory .
According to the pilot-wave interpretation, it is the position of the particle beable that is revealed when a position measurement is performed. In the following section, where we deal with the pilot-wave description of a measurement process, we consider this more carefully.
If we now consider an ensemble of $`N`$-particle systems, all described by the same wavefunction $`\psi `$, then this ensemble determines a probability distribution $`\rho (𝐗_1,\mathrm{},𝐗_N,t)`$ of the actual position vectors of the $`N`$ particle beables. Because the motion of the particle beables is governed by the guidance equations (1.5), their distribution $`\rho `$ satisfies the same continuity equation as the quantum mechanical probability density $`|\psi |^2`$. Therefore, if the densities $`\rho `$ and $`|\psi |^2`$ are equal at a certain time $`t_0`$, i.e.
$$\rho (𝐱_1,\mathrm{},𝐱_N,t_0)=|\psi (𝐱_1,\mathrm{},𝐱_N,t_0)|^2,$$
(1.6)
then the equality will hold for all times $`t`$, i.e.
$$\rho (𝐱_1,\mathrm{},𝐱_N,t)=|\psi (𝐱_1,\mathrm{},𝐱_N,t)|^2.$$
(1.7)
In the pilot-wave interpretation one assumes that initially, before a measurement is performed, the distribution of the particle beables $`\rho `$ (over the ensemble) is given by the quantum mechanical distribution $`|\psi |^2`$ (this assumption is also called the quantum equilibrium hypothesis and the distribution $`\rho =|\psi |^2`$ is called the equilibrium distribution ). The densities will then remain equal during the experiment, and pilot-wave theory and standard quantum mechanics will predict the same detection probabilities for the particle positions.
Because most quantum measurements boil down to position measurements, pilot-wave theory and standard quantum mechanics will in general yield the same detection probabilities. The situation is different, if one considers for example measurements involving time related quantities, such as time of arrival, tunneling times etc. Pilot-wave theory makes unambiguous predictions for such measurements, but in conventional quantum theory there is no consensus about what these quantities should be (see e.g. and references therein).
The quantum equilibrium hypothesis is introduced here to match the empirical distributions predicted by quantum theory. However, there exist some possible justifications for the quantum equilibrium hypothesis. One possible justification was presented by Dürr et al. . But it would take to far to repeat their analysis here. Another possible justification was presented by Valentini . By a sub-quantum $`H`$-theorem Valentini was able to show that in reasonable circumstances a non-equilibrium distribution $`\rho |\psi |^2`$ for the particle beables (all guided by the same wavefunction) may approach the equilibrium distribution on a certain coarse grained level.<sup>4</sup><sup>4</sup>4Recently this was illustrated by numerical simulations . This suggest that quantum theory can be regarded as an equilibrium theory. Although this yields an interesting research program, we will not consider the possibility of non-equilibrium in this thesis. Our main goal is rather to study in how far a pilot-wave interpretation is possible to cover other domains in quantum theory.
We also want to note that there were recent claims by Ghose , and which were later adopted by Golshani and Akhavan , that standard quantum mechanics and pilot-wave theory would predict incompatible results for some specific experiments. However, we argued elsewhere that these claims are flawed . We indicated that a non-equilibrium density for the particle beables is implicitly assumed from the outset. The density of the particle beables then remains in non-equilibrium during the experiment and hence it is obvious that one arrives at incompatible predictions for the two theories. Despite our comment (and that of others ) the experiment proposed by Ghose was recently performed by Brida et al. . The result of the experiment was that standard quantum theory was confirmed (as expected). A correct analysis of the experiment in terms of pilot-wave theory would have led to the same predictions as standard quantum theory.
Pilot-wave theory has many features which are not present in quantum theory. The most striking property of pilot-wave theory is that it is nonlocal. In fact, as shown by Bell, any realistic theory which leads to the same statistical predictions as standard quantum theory must be nonlocal . Yet, quantum theory remains local in the sense that one cannot use quantum theory to send faster than light signals.<sup>5</sup><sup>5</sup>5We can state this more correctly as follows. As shown by Ghirardi et al. , the standard quantum theory of measurement cannot be used for superluminal transmission of signals. Hence, if the velocity of probability flow does not exceed the speed of light, which is for example always the case for the relativistic theory spin-1/2 of Dirac (the velocities along the flowlines of the particle probability density are bounded by the speed of light), then quantum theory does not allow faster than light signals. In non-relativistic quantum theory there is in fact no restriction on the velocity of probability flow. But of course the domain of applicability of non-relativistic quantum theory is limited to ‘low’ speeds. In the pilot-wave interpretation, the nonlocality manifests itself by the fact that the position of one-particle beable may depend on the positions of other particle beables (by the guidance law (1.5)). This dependence is instantaneous no matter how far the other particle beables may be located. It is also important to note that this nonlocality is not a consequence of dealing with a non-relativistic description of quantum phenomena. The nonlocality of pilot-wave theory is inescapable, even at the level of relativistic quantum theory. A very illustrative example of the nonlocality present in pilot-wave theory was given by Rice .
In conclusion, we arrive at pilot-wave theory only by a minor shift in interpretation (although with far reaching consequences). Instead of interpreting
$$|\psi (𝐱_1,\mathrm{},𝐱_N,t)|^2d^3x_1\mathrm{}d^3x_N$$
(1.8)
as the probability of finding the particles in a volume element $`d^3x_1\mathrm{}d^3x_N`$ around the configuration $`(𝐱_1,\mathrm{},𝐱_N)`$, at a certain time $`t`$, as in the standard interpretation of quantum mechanics, we interpret it in pilot-wave theory as the probability of the particles being in a volume element $`d^3x_1\mathrm{}d^3x_N`$ around the configuration $`(𝐱_1,\mathrm{},𝐱_N)`$ at the time $`t`$. Inspired by the analogy with the continuity equation in hydrodynamics, we derive the velocity field (1.4) for the particle beables from the quantum continuity equation (1.2) for the density $`|\psi |^2`$. This is the scheme that we will adopt in the rest of the thesis. When constructing a pilot-wave theory for quantum theory we shall try to identify a continuity equation that can be seen as a conservation equation for a density, be it a density of particles or fields,<sup>6</sup><sup>6</sup>6We will argue that a field ontology seems preferred over a particle ontology in a pilot-wave interpretation for quantum field theory. Instead of introducing particle beables we will then introduce field beables. and then from this continuity equation we shall try to construct a guidance equation for the particles or fields.
### 1.2 The measurement process
In this section we describe the standard quantum mechanical measurement process in terms of the pilot-wave interpretation, along the lines of the presentations that can be found in . In the pilot-wave interpretation, the measurement process is treated just as any other quantum process; there is no privileged role for the observer or measurement apparatus.
Suppose a system which is described by the wavefunction $`\psi ^{(s)}`$. The system may consist of $`N`$ particles, so that the wavefunction lives in $`3N`$-dimensional configuration space. There also correspond $`N`$ particle beables with the system, the positions of which we denote by the $`3N`$-dimensional vector $`x^{(s)}`$. Suppose similarly an apparatus with wavefunction $`\psi ^{(a)}`$ and a collection of particle beables at the configuration $`x^{(a)}`$. The apparatus is introduced to measure some property of the system. In the standard interpretation of quantum theory this property is represented by an operator $`\widehat{A}`$ and the possible outcomes of the measurement correspond to the eigenvalues of this operator.
Initially the system under observation and the measurement apparatus have not interacted yet, so that they may be described by the product wavefunction $`\psi ^{(s)}\psi ^{(a)}`$. As time evolves the system under observation gets coupled to the apparatus and the total wavefunction $`\psi ^{(s)}\psi ^{(a)}`$ evolves to the entangled wavefunction $`_i\psi _i^{(s)}\psi _i^{(a)}`$, where the states $`\psi _i^{(s)}`$ are eigenstates of the operator $`\widehat{A}`$. This evolution is determined by the Schrödinger equation, which should contain a particular interaction Hamiltonian depending on the observable that is being measured (e.g. this interaction Hamiltonian could be a von Neumann type of interaction Hamiltonian). The states $`\psi _i^{(a)}`$ are assumed to be non-overlapping in configuration space, i.e. $`\psi _i^{(a)}\psi _j^{(a)}0`$ for $`ij`$.<sup>7</sup><sup>7</sup>7Note that the condition that the states are non-overlapping is stronger than the condition that they are orthogonal. If states are non-overlapping they are orthogonal, but not vice versa. For example different plane waves are orthogonal but are overlapping. In fact it is sufficient that the overlap of the states is minimal. This property is generally satisfied in an ordinary measurement. We will present the reason for this below.
Now if the apparatus would be known to be in one particular state, say $`\psi _k^{(a)}`$, then the system under observation would be in the state $`\psi _k^{(s)}`$. In standard quantum theory, the collapse rule is introduced to reduce the state of the total state $`_i\psi _i^{(s)}\psi _i^{(a)}`$ to the state $`N\psi _k^{(s)}\psi _k^{(a)}`$ (with $`N`$ some normalization factor). The result of the measurement is then the eigenvalue of $`\widehat{A}`$ corresponding to the eigenstate $`\psi _k^{(s)}`$.
In fact, the collapse law does not have to be invoked at this stage yet. A second apparatus may also be introduced, which measures the first apparatus, and so on. This chain of apparatuses getting correlated may then for example be ended by a final observer for which the collapse law may be invoked. As explained in the introduction, the point where the collapse law should apply is not well defined and presents the core of the measurement problem.
By contrast the pilot-wave description of the measurement process is unambiguous. Because the different terms $`\psi _i^{(s)}\psi _i^{(a)}`$ are non-overlapping, they can be seen as defining ‘channels’ in configuration space, the channels being the non-overlapping supports of the different terms $`\psi _i^{(s)}\psi _i^{(a)}`$. The configuration $`(x^{(s)},x^{(a)})`$, which has the positions of the particle beables as components, enters one of the channels during the interaction. Suppose the configuration has entered the channel corresponding to $`\psi _k^{(s)}\psi _k^{(a)}`$. The other terms $`\psi _i^{(s)}\psi _i^{(a)}`$ with $`ik`$ are called the empty waves. If the different terms $`\psi _i^{(s)}\psi _i^{(a)}`$ do not overlap again at a later time (this is in principle accomplished by coupling the system with a large number of particles, which leads to decoherence, so that the probability of re-overlap becomes minimal), then as far as the particle beables are concerned, the empty wavepackets have no further influence on the motion of the beables $`(x^{(s)},x^{(a)})`$ and hence these wavepackets may then be dismissed in the future description of the particle beables. This corresponds to the collapse of the wavefunction in the standard interpretation of quantum mechanics.
Because we further assume the beables to be distributed according to the quantum distribution, one can easily verify that the probability for the particle beables to enter the channel corresponding to $`\psi _k^{(s)}\psi _k^{(a)}`$ is given by
$$\psi _k^{(s)}\psi _k^{(a)}|\underset{i}{}\psi _i^{(s)}\psi _i^{(a)}=𝑑\mathrm{\Omega }\left|\psi _k^{(s)}\psi _k^{(a)}\right|^2,$$
(1.9)
where the integral on the right hand side ranges over the whole configuration space ($`d\mathrm{\Omega }`$ is the measure on the configuration space). Hence we recover the quantum probabilities in the pilot-wave description of the measurement process.
So we can describe the measurement process in the context of pilot-wave theory. Essential in our treatment is that at some stage in the measurement chain, the wavefunctions of the apparatus are non-overlapping in configuration space, because this allows us to dismiss the empty wavepackets. This situation is obtained in an ordinary measurement. For example the different states could correspond to macroscopic needles pointing in different directions and one can easily convince oneself that these states are non-overlapping in configuration space. The reason why the total wavefunction of system and apparatus actually evolves to such a superposition is a purely quantum mechanical one.
### 1.3 Summary and organization of the thesis
We have seen how we can give a pilot-wave description of a spinless, non-relativistic quantum system. In the main part of the thesis we will study how this pilot-wave interpretation can be extended to cover other domains in quantum theory. We will successively consider the domain of non-relativistic quantum theory, relativistic quantum theory and quantum field theory. More advanced domains, such as quantum gravity and string theory, are not considered.
The extension of the pilot-wave formulation a spinless, non-relativistic quantum system to include spin does not present any difficulties. We deal with this extension in Chapter 2.
The construction of a pilot-wave formulation for relativistic quantum theory is more problematic. In Chapter 3, we will consider relativistic wave equations and we will consider the question to which extent it is possible to construct a pilot-wave model for these relativistic wave equations. It will turn out that a pilot-wave interpretation with point particle as beables, analogous to the pilot-wave interpretation for non-relativistic quantum mechanics, is in general impossible. The reason is that, already at the standard quantum mechanical level, a particle interpretation in analogy with the one for non-relativistic quantum mechanics is in general not possible.
It seems that only for the Dirac theory for spin-1/2, and under restricted circumstances, such a quantum mechanical particle interpretation may be provided. For example for sufficiently low energies a one-particle interpretation is possible. This is because there exists a positive density (proportional to the charge density) which is the time component of a future-causal four-vector and which can hence be interpreted as a probability density. For higher energies and if only electromagnetic interaction is considered, one can maintain the particle interpretation, albeit a many-particle one, by the introduction of a Dirac sea. For other types of interaction, such as weak interaction, it is unknown how to extend the notion of a Dirac sea and hence it is unknown how to continue the particle approach. In the domain where the quantum mechanical particle interpretation is applicable for the Dirac theory, a pilot-wave interpretation can be devised. This was already clear to Bohm, who originated the pilot-wave interpretation for the Dirac theory.
After reviewing Bohm’s pilot-wave interpretation for the Dirac theory, we consider the alleged pilot-wave models for the Duffin-Kemmer-Petiau (DKP) theory and the Harish-Chandra (HC) theory, which were initiated by Ghose et al. The DKP wave equation is a first-order relativistic equation for massive spin-0 and spin-1, but is nevertheless completely equivalent to the second order Klein-Gordon equation in the spin-0 representation and to the Proca equations in the spin-1 representation. The HC equation is the massless counterpart of the DKP equation, which is equivalent to the massless Klein-Gordon equation in the spin-0 representation, and to Maxwell’s equations for the electromagnetic field in the spin-1 representation. As is well known there is no quantum mechanical particle interpretation for these wave equations because of the lack of a conserved, future-causal current (contrary to the Dirac theory, the charge current is not always future-causal for spin-0 and spin-1 bosons). In fact there is not even a quantum mechanical interpretation, because the lack of a positive definite inner product blocks the setup of a Hilbert space (again this is related to the fact that the charge currents for both spin-0 and spin-1 are not always future-causal). Ghose et al. tried to give a quantum mechanical particle interpretation by constructing a conserved, future-causal current from the energy–momentum tensor. With this quantum mechanical particle interpretation they could also as associate a pilot-wave model. The resulting equations look very similar to the ones for the Dirac theory. However, despite this similarity, we show that the suggested quantum mechanical particle interpretation, and hence also the associated pilot-wave model, suffer from some problems, which make this approach in general untenable.
Although the pilot-wave model suggested by Ghose et al. can not be treated as valid model describing physical reality, we think that the model still has value as an illustrative model. For this reason, we further consider the extension of the model to many particles.
Of course, as is well known, relativistic wave equations are not suitable to describe high energy quantum systems. The theory describing high energy quantum systems is quantum field theory. Quantum field theory is not about particles in physical 3-space, as was non-relativistic quantum theory; strong localization of particles in physical 3-space leads to problems with causality (i.e. superluminal spread of localized states ). Instead quantum field theory can be seen as describing fields in physical 3-space.<sup>8</sup><sup>8</sup>8Only in momentum space, when using Fock space, one can recover the notion of particles. This is clear in the functional Schrödinger picture, where the quantum states are described by wavefunctionals, which are defined on a configuration space of fields. Instead of searching for a pilot-wave interpretation for relativistic quantum theory in terms of particle beables, we will therefore consider the possibility of a pilot-wave interpretation in terms of field beables, a view strongly supported by Valentini.
It will appear that the construction of a pilot-wave theory in terms of field beables presents no difficulty in the bosonic case. This will be illustrated in Chapter 4 with a discussion on the construction of a pilot-wave theory for the massive spin-0 field, the massive spin-1 field, the electromagnetic field and then also for the massive spin-0 field coupled to the electromagnetic field (i.e. scalar quantum electrodynamics). In particular we will discuss in detail the two existing models for the electromagnetic field, namely the one by Bohm and Kaloyerou and the one by Valentini. The main difference between the two models is that Bohm and Kaloyerou only introduce beables for gauge independent variables, whereas Valentini also introduces beables for gauge variables. We will show that the guidance equations for the beables corresponding to the gauge variables are rather meaningless, because they only express the fact that these beables are stationary. In addition, inclusion of these beables for gauge dependent variables also makes that the densities of field beables are non-normalizable. In order to avoid this problem, we think it is preferable to adopt the approach by Bohm and Kaloyerou.
A pilot-wave interpretation in terms of field beables for fermionic field theory seems less straightforward. In Chapter 5 we reconsider the idea of Valentini to construct a pilot-wave interpretation in terms of field beables, with the field beables being elements of the Grassmann algebra. This approach looks very promising at first sight, because there exists a functional Schrödinger picture for fermionic fields in terms of Grassmann variables. However, closer inspection reveals that it is not possible to associate a pilot-wave interpretation with it.
Hence for fermionic field theory a different approach should be taken. There do exist different approaches. There is for example the pilot-wave approach by Holland and the Bell-type model by Dürr, Goldstein, Tumulka and Zanghì. The Bell-type model differs from pilot-wave models in the fact that it involves an element of stochasticity. Although these models look very interesting, we do not consider them in detail in the thesis.
It is important to note that when the field approach is taken as fundamental, which we do in this thesis, then this approach is incompatible with the particle approach which was so successful for non-relativistic quantum theory. The field beables that are introduced in field theory do not reside into the particle beables in the non-relativistic limit. Hence, also in the non-relativistic case the actual beables should be regarded to be fields.
Nevertheless, the pilot-wave interpretation in terms of particle beables may serve well for illustrative purposes. An example of this is given in Chapter 6. There we show that the pilot-wave interpretation in terms of particle beables may serve as a theoretical underpinning for otherwise rather ad hoc trajectories that are used for describing some experiments concerning optical imaging.
## Chapter 2 Particle beables for non-relativistic quantum mechanics
### 2.1 Non-relativistic quantum mechanics
In the preceding chapter, the pilot-wave interpretation for a system consisting of non-relativistic spinless particles was introduced. In this section we consider the extension of this pilot-wave interpretation to include spin. We only consider the one-particle case. The extension to many particles is straightforward. We consider charged particles which move under the influence of an external electromagnetic field.
In the standard interpretation of quantum mechanics, a non-relativistic particle with spin $`s`$ is described by means of a $`(2s+1)`$-component wavefunction $`\psi `$, who’s index transforms according to the $`(2s+1)`$-dimensional representation of the rotation group. The wavefunction $`\psi `$ satisfies the wave equation (cf. \[43, p. 471\] and )<sup>1</sup><sup>1</sup>1In this chapter, and in subsequent chapters, we adopt the summation convention of Einstein.
$$i\mathrm{}\frac{\psi _\alpha (𝐱,t)}{t}=\frac{\mathrm{}^2D^2}{2m}\psi _\alpha (𝐱,t)\frac{eg}{2mc}𝐒_{\alpha \beta }^{(s)}𝐁\psi _\beta (𝐱,t)+(eV_0+V)\psi _\alpha (𝐱,t),$$
(2.1)
where $`e`$ is the charge of the particle. $`V_i`$ and $`V_0`$ are the electromagnetic potentials, with corresponding magnetic field $`𝐁=\mathbf{}\times 𝐕`$ and $`V`$ denotes an additional scalar potential. $`D_i=_i\frac{ie}{\mathrm{}c}V_i`$ is the covariant derivative. The three $`(2s+1)`$-dimensional matrices $`S_i^{(s)}`$ ($`i=1,2,3`$) are the generators of the rotation group in the $`(2s+1)`$-dimensional representation. They satisfy the commutation relations
$$[S_i^{(s)},S_j^{(s)}]=i\mathrm{}\epsilon _{ijk}S_k^{(s)}.$$
(2.2)
The constant $`g`$ is the gyromagnetic factor. Hurley derived the wave equation under the assumptions of Galilean covariance and ‘minimality’ , and he found the gyromagnetic factor
$$g=\{\begin{array}{cc}0\hfill & \text{for spin }0\hfill \\ \frac{1}{s}\hfill & \text{for spin }s\hfill \end{array}.$$
(2.3)
This implies that the correct gyromagnetic factor for an elementary particle can be found even without considering relativistic wave equations (although the correct prediction of the gyromagnetic factor for the electron is generally regarded as a success of the Dirac equation). If the particle has an internal structure, then the gyromagnetic factor may of course differ from $`1/s`$.
The conservation equation for the particle detection probability density $`\psi ^{}\psi `$ reads
$$\frac{\psi ^{}\psi }{t}+\mathbf{}𝐣=0$$
(2.4)
with $`𝐣`$ the 3-vector probability current which can be written as the sum
$$𝐣=𝐣_c+𝐣_s,$$
(2.5)
of a current which formally resembles the conventional Schrödinger current
$$𝐣_c=\frac{\mathrm{}}{2mi}\left(\psi ^{}𝑫\psi \left(𝑫\psi \right)^{}\psi \right)=\frac{\mathrm{}}{2mi}\left(\psi ^{}\mathbf{}\psi \left(\mathbf{}\psi ^{}\right)\psi \right)\frac{e}{mc}𝐕\psi ^{}\psi $$
(2.6)
and a ‘spin current’
$$𝐣_s=\frac{g}{2m}\mathbf{}\times \left(\psi ^{}𝐒^{(s)}\psi \right).$$
(2.7)
With the introduction of the spin 3-vector
$$𝐬=\frac{\psi ^{}𝐒^{(s)}\psi }{\psi ^{}\psi }$$
(2.8)
and the magnetic moment 3-vector
$$𝐦=\frac{ge}{2mc}\psi ^{}\psi 𝐬=\frac{ge}{2mc}\psi ^{}𝐒^{(s)}\psi ,$$
(2.9)
the spin term in the current can also be written as
$$𝐣_s=\frac{g}{2m}\mathbf{}\times \left(\psi ^{}\psi 𝐬\right)=\frac{c}{e}\mathbf{}\times 𝐦.$$
(2.10)
Within a factor $`e`$, corresponding to a change from the particle probability current to the charge current, the spin current $`𝐣_s`$ formally resembles a magnetization current for a classical polarized medium.
We can construct a pilot-wave interpretation by introducing a structureless particle (the particle beable), who’s motion is governed by the wavefunction according to the guidance equation<sup>2</sup><sup>2</sup>2Note that in the previous chapter we used a different notation for the position vector of the particle beable and the argument of the wavefunction. Because this can in fact not lead to possible confusion, we use from now on the same notation for both.
$$\frac{d𝐱}{dt}=\frac{𝐣}{\psi ^{}\psi }.$$
(2.11)
For an ensemble of spin-$`s`$ particles, all described by the same wavefunction $`\psi `$, the equilibrium distribution for the particle beables is given by $`\psi ^{}\psi `$.
### 2.2 Examples
We now consider some examples:
Spin-0: In the spin-0 case, the generators of the rotation group $`S_i^{(0)}`$ are zero and the wave equation (2.1) is simply the Schrödinger wave equation for a spinless particle
$$i\mathrm{}\frac{\psi (𝐱,t)}{t}=\frac{\mathrm{}^2D^2}{2m}\psi (𝐱,t)+(eV_0+V)\psi (𝐱,t).$$
(2.12)
The corresponding current is
$$𝐣=\frac{\mathrm{}}{2mi}\left(\psi ^{}𝑫\psi \left(𝑫\psi \right)^{}\psi \right).$$
(2.13)
The pilot-wave interpretation is the one given originally by de Broglie and Bohm (cf. Section 1.1). In particular the guidance equation reads
$$\frac{d𝐱}{dt}=\frac{\mathrm{}}{2mi|\psi |^2}\left(\psi ^{}𝑫\psi \left(𝑫\psi \right)^{}\psi \right).$$
(2.14)
Spin-1/2: In the spin-$`1/2`$ case, the generators of the rotation group are proportional to the $`2\times 2`$ Pauli matrices $`\sigma _i`$ ($`i=1,2,3`$), i.e. $`S_i^{(1/2)}=\mathrm{}\sigma _i/2`$. The Schrödinger equation (2.1) for the two component wavefunction $`\psi `$ is then the Pauli equation
$$i\mathrm{}\frac{\psi _\alpha (𝐱,t)}{t}=\frac{\mathrm{}^2D^2}{2m}\psi _\alpha (𝐱,t)\frac{eg\mathrm{}}{4mc}𝝈_{\alpha \beta }𝐁\psi _\beta (𝐱,t)+(eV_0+V)\psi _\alpha (𝐱,t).$$
(2.15)
In the corresponding guidance law for the particle beable, there is a spin contribution arising from the nonzero spin term in the current
$$𝐣_s=\frac{g\mathrm{}}{4m}\mathbf{}\times \left(\psi ^{}𝝈\psi \right),$$
(2.16)
so that the total current reads
$$𝐣=\frac{\mathrm{}}{2mi}\left(\psi ^{}𝑫\psi \left(𝑫\psi \right)^{}\psi \right)+\frac{g\mathrm{}}{4m}\mathbf{}\times \left(\psi ^{}𝝈\psi \right).$$
(2.17)
The corresponding guidance equation reads
$$\frac{d𝐱}{dt}=\frac{\mathrm{}}{2mi\psi ^{}\psi }\left(\psi ^{}𝑫\psi \left(𝑫\psi \right)^{}\psi \right)+\frac{g\mathrm{}}{4m\psi ^{}\psi }\mathbf{}\times \left(\psi ^{}𝝈\psi \right).$$
(2.18)
The trajectories for a non-relativistic spin-1/2 particle (with the additional spin term) were recently studied for specific systems. Holland and Philippidis studied the particle paths for spin-1/2 eigenstate for the two slit experiment . Colijn and Vrscay studied the spin-1/2 particle paths for hydrogen eigenstates and transitions between them .
Spin-1: In the spin-$`1`$ case, the generators of the rotation group are the $`3\times 3`$ matrices $`(S_j^{(1)})_{ik}=i\mathrm{}ϵ_{ijk}`$. The Schrödinger equation (2.1) for the three component wavefunction $`\psi `$ is then the non-relativistic spin-1 equation
$$i\mathrm{}\frac{\psi _i(𝐱,t)}{t}=\frac{\mathrm{}^2D^2}{2m}\psi _i(𝐱,t)\frac{ieg\mathrm{}}{2mc}\epsilon _{ijk}B_j\psi _k(𝐱,t)+(eV_0+V)\psi _i(𝐱,t),$$
(2.19)
with $`i=1,2,3`$. In the corresponding guidance law for the particle beable, there is a spin contribution arising from the nonzero spin term in the current
$$(𝐣_s)_i=\frac{g}{2m}\left(\mathbf{}\times (\psi ^{}𝐒^{(1)}\psi )\right)_i=\frac{g\mathrm{}}{2m}\text{Im}\left(\psi _i^{}_j\psi _j\psi _j^{}_j\psi _i\right),$$
(2.20)
so that the total current reads
$$j_i=\frac{\mathrm{}}{2mi}\left(\psi ^{}D_i\psi \left(D_i\psi \right)^{}\psi \right)+\frac{g\mathrm{}}{2m}\text{Im}\left(\psi _i^{}_j\psi _j\psi _j^{}_j\psi _i\right).$$
(2.21)
The corresponding guidance equation reads
$$\frac{dx_i}{dt}=\frac{\mathrm{}}{2mi\psi ^{}\psi }\left(\psi ^{}D_i\psi \left(D_i\psi \right)^{}\psi \right)+\frac{g\mathrm{}}{2m\psi ^{}\psi }\text{Im}\left(\psi _i^{}_j\psi _j\psi _j^{}_j\psi _i\right).$$
(2.22)
### 2.3 Spin eigenstates
The Schrödinger equation for a spinless particle (2.12) can be derived from the wave equation (2.1) by considering a spin eigenstate. However the particle current for a spin eigenstate will in general not reduce to the current for a spinless particle (2.13); as pointed out before in the case of spin-1/2 , there will be an additional spin contribution to the current.
Let us consider this in more detail. If we consider a particle with arbitrary spin $`s`$ for which the wavefunction is a spin eigenstate, then the wavefunction can be written as the product of a space dependent part and a spin dependent part, i.e.
$$\psi _\alpha (𝐱,t)=\psi ^{}(𝐱,t)\chi _\alpha ,\chi _\alpha ^{}\chi _\alpha =1,$$
(2.23)
with $`\psi ^{}(𝐱,t)`$ a scalar wavefunction. For vanishing electromagnetic potentials, $`V_i=V_0=0`$, it follows from the wave equation (2.1) that $`\psi ^{}(𝐱,t)`$ satisfies the Schrödinger equation for a spinless particle. The corresponding particle current for the spin eigenstate reads
$$𝐣=\frac{\mathrm{}}{2mi}\left(\psi ^{}\mathbf{}\psi ^{}\left(\mathbf{}\psi ^{}\right)\psi ^{}\right)+\frac{g}{2m}\mathbf{}\times (|\psi ^{}|^2𝐬)$$
(2.24)
and the corresponding guidance equation is
$$\frac{d𝐱}{dt}=\frac{\mathrm{}}{2mi|\psi ^{}|^2}\left(\psi ^{}\mathbf{}\psi ^{}\left(\mathbf{}\psi ^{}\right)\psi ^{}\right)+\frac{g}{2m|\psi ^{}|^2}\mathbf{}\times (|\psi ^{}|^2𝐬).$$
(2.25)
The spin vector $`𝐬`$ is now a constant vector, given by
$$𝐬=\chi ^{}𝐒^{(s)}\chi .$$
(2.26)
So, even though the spin dependent part $`\chi `$ can be factored out of the wave equation, leading to the Schrödinger equation (2.12) for $`\psi ^{}`$, it still appears non-trivially in the current and hence in the guidance equation. As a result, the non-relativistic description of a particle in a spin eigenstate with nonzero spinvector $`𝐬`$, should include the spin term in the current. This spin term potentially plays for example a role in time of arrival measurements (see below).
### 2.4 Note on the uniqueness of the pilot-wave interpretation
#### 2.4.1 On the uniqueness of the particle current and corresponding guidance equation
Note that the current $`𝐣`$ is not uniquely determined by the continuity equation (2.4). It is determined only up to a divergenceless vector. For example, one can construct a new current $`\overline{𝐣}`$ by adding the divergenceless current $`𝐣_a`$ to the current $`𝐣`$. The newly defined current $`\overline{𝐣}=𝐣+𝐣_a`$ then also satisfies the continuity equation, with the same probability density $`\psi ^{}\psi `$. Hence, we are left with an apparent ambiguity in the definition of the particle probability current. This ambiguity is unobservable when the quantum probabilities are derived solely from the density $`\psi ^{}\psi `$ (such as spin measurements with a Stern-Gerlach setup), because this density is unaltered by the addition of $`𝐣_a`$ to the current. Nevertheless, the additional current may in principle lead to observable effects in measurements involving time (see below). In any case, the guidance equation in the pilot-wave interpretation is derived from the quantum mechanical particle current and the possibility of an additional contribution to the particle current leads to an ambiguity at the level of the pilot-wave interpretation. For example the spin current $`𝐣_s`$ is divergenceless and hence could in principle be relinquished in the definition of the current, and hence also in the definition of the guidance equation.
Nevertheless, there exist some arguments in favor of (2.5) as the definition for the particle current.<sup>3</sup><sup>3</sup>3Appeal to Noether’s theorem in order to derive the correct current is of no use in this case. Because if the charge current is derived as the conserved current corresponding to global phase invariance, then it is only determined up to a total divergence and hence Noether’s theorem is not decisive in whether or not we should in include the spin term in the current. A first argument rests on the observation that it is the charge current $`ej^\mu =e(c\psi ^{}\psi ,𝐣)`$ that couples to the electromagnetic field in the Maxwell equations $`_\mu F^{\mu \nu }=ej^\mu /c`$, with $`F^{\mu \nu }`$ is the electromagnetic field tensor.<sup>4</sup><sup>4</sup>4The Maxwell equations $`_\mu F^{\mu \nu }=j^\mu /c`$ can be derived from the fully coupled Lagrangian density
$$=\frac{i\mathrm{}}{2}\left(\psi ^{}\frac{\psi }{t}\frac{\psi ^{}}{t}\psi \right)\frac{\mathrm{}^2}{2m}(𝐃\psi )^{}𝐃\psi +\frac{eg}{2mc}\psi ^{}𝐒^{(s)}𝐁\psi (eV_0+V)\psi ^{}\psi \frac{1}{4}F_{\mu \nu }F^{\mu \nu }.$$
(2.27) Hence, if the particle current is chosen to be proportional to the charge current, then it should be given by (2.5).
In the case of spin-1/2, Holland considered still other arguments to establish the uniqueness of the particle current . Holland first showed that the particle current in the relativistic spin-1/2 Dirac theory is unique (under reasonable assumptions). With the demand that the non-relativistic spin-1/2 particle current should be obtained by taking the non-relativistic limit of the Dirac current, this non-relativistic particle current is also unique. The resulting non-relativistic spin-1/2 current is the one that is presented in Section 2.2. This choice for the non-relativistic current also implies that the pilot-wave interpretation that can be provided for the Dirac equation (see Section 3.2) reduces to the pilot-wave interpretation for non-relativistic quantum theory as presented here in the non-relativistic limit.
In we used similar arguments as Holland’s in order to obtain the uniqueness of the particle model for relativistic spin-0 and spin-1 proposed by Ghose et al. . By taking the non-relativistic limit, the particle currents in the particle model of Ghose et al. reduce to the non-relativistic currents given in Section 2.2 (with the gyromagnetic factor as given in (2.3)). However the model proposed by Ghose et al. contains some features which make it in general untenable to maintain the model as a valid description of physical reality (we will discuss this in detail in Section 3.3, albeit without the discussion on the uniqueness). Therefore, this uniqueness in the case of spin-0 and spin-1 should definitely not be regarded as conclusive. Nevertheless, with Holland’s result for spin-1/2, we could require on the basis of uniformity that the spin term should be included for all values of spin.
In fact, Deotto and Ghirardi were among the first to consider different currents compatible with the continuity equation from which different guidance laws for the particle beables could be derived. However, they came to the conclusion that the requirement of Galilean covariance of the current was insufficient to impose its uniqueness.
Although the additional contribution in the guidance equation may not be detectable in quantum probabilities derived from $`\psi ^{}\psi `$, it plays a role in the case of time of arrival measurements . In the pilot-wave interpretation, the distribution of arrival times of the particles at $`𝐱`$ for free particles is given by $`|𝐣(𝐱,t)|`$ and the corresponding mean arrival time at $`𝐱`$ is
$$\overline{t}=\frac{_0^+\mathrm{}|𝐣(𝐱,t)|t𝑑t}{_0^+\mathrm{}|𝐣(𝐱,t)|𝑑t}.$$
(2.28)
Because these quantities depend on the current $`𝐣`$ and not on the probability density $`\psi ^{}\psi `$, the spin contribution in the current may in principle lead to an observable effect. Recently it was argued that for free spin eigenstates, spin contributions with a gyromagnetic factor $`g=0`$ or $`g=1/2`$ would in principle be experimentally distinguishable for time of arrival distributions . Hence time of arrival measurements might contribute to the determination of the correct particle current and hence to the correct guidance equation in the pilot-wave interpretation. Of course the question remains whether the difference for mean arrival times, which are calculated with different guidance equations, is experimentally observable.
The definition of the time of arrival depends in fact not only on the particular choice for the guidance equation. It also depends strongly on the particular pilot-wave model we adhere to. In the next section we give examples of such alternative models. For some of these models the question whether of not we should include a spin term in the guidance equation becomes irrelevant, simply because they involve a different ontology. It may also very well be that for some of these models, the definition of ‘time of arrival’ depends on how exactly we model time measurements. In such a case the definition of time of arrival may not be so unambiguous anymore.
#### 2.4.2 Alternative pilot-wave models
In the pilot-wave interpretation we presented here, the particle beable is a structureless point particle; it does not carry spin degrees of freedom. Spin is solely a property of the wavefunction. The pilot-wave interpretation still solves the measurement problem, because measurements can in general be reduced to position measurements, the exceptions being measurements involving time.
There also exist pilot-wave models in which spin degrees of freedom are also attached to the particle beables, with the particles beables being point particles or rigid bodies, see e.g. Bohm et al. and Holland . Although these models are very interesting and may be very illustrative in some cases, they tend to complicate things. Especially if one only wants the pilot-wave model to reproduce the quantum probabilities, the assumption of an additional structure of the particle beables is unnecessary.
The view of particle beables as structureless point particles is also the one advocated by Bell , Dürr et al. and later also by Bohm and Hiley . The only difference with the model presented here is that Bell and Dürr et al. do not include the spin part, which arises from the spin current $`𝐣_s`$, in the guidance equation. Bohm and Hiley consider the spin term when they discuss non-relativistic spin-1/2 particles.
There exist still other, completely different, ontologies. For example, in Chapters 4 and 5 we will see that a pilot-wave interpretation with fields as beables seems more natural in quantum field theory. However this field ontology does not reduce to a particle ontology in the non-relativistic limit. This means that if the field ontology is taken as fundamental then even at the level of non-relativistic quantum systems the beables are fields. With a field beable approach, the issue of identifying a unique guidance equation from the continuity equation then of course acquires a new character. We do not know of arguments leading to a preferred choice for the guidance equation in this case. Therefore we will not return to the question of uniqueness when dealing with field theory.
## Chapter 3 Particle beables for relativistic quantum mechanics
### 3.1 Introduction
In the preceding chapter we have seen how a pilot-wave interpretation could be constructed for a non-relativistic particle with arbitrary spin. The success of this pilot-wave formulation can be regarded as a consequence of the fact that the non-relativistic wave equations already admitted a quantum mechanical particle interpretation at the level of the standard interpretation. The key feature that allowed for this particle interpretation was the existence of a positive definite density which is conserved. This density could then be successfully identified with a particle density.
The situation is totally different in relativistic quantum theory. When we try to formulate a particle interpretation for relativistic wave equations along the lines of the particle interpretation in non-relativistic quantum theory, we encounter some serious difficulties. The difficulties have nothing to do with the construction of relativistic wave equations itself. Wave equations that are both Lorentz covariant and causal can be found for any value of spin (see e.g. the Bhabha wave equations below).<sup>1</sup><sup>1</sup>1As shown by Velo and Zwanziger covariant wave equations may suffer from noncausal propagations , but the Bhabha equations do not suffer from this problem . However, the problem is rooted in the fact that the wave equations in general lack an associated current which has all the mathematical properties required for a particle current.
A notable exception is the spin-1/2 theory of Dirac. In the Dirac theory there is a conserved current (proportional to the charge current) which has a positive time component in every Lorentz frame. For energies below the threshold of pair creation, this time component can then be interpreted as the particle probability density. If only electromagnetic interaction is considered, then one can extend this particle interpretation to higher energies by using Dirac’s original suggestion of the Dirac sea and by passing to the many-particle description. Because one has a quantum mechanical particle interpretation in this case, one can also devise a pilot-wave interpretation, in the same way as the pilot-wave interpretation for non-relativistic quantum mechanics. The first to present this pilot-wave interpretation was Bohm. We recall this pilot-wave interpretation in Section 3.2.
This success for the Dirac equation is not repeated for other types of wave equations. For example if we consider the Bhabha wave equations, which are Dirac-type equations for massive particles with arbitrary spin,<sup>2</sup><sup>2</sup>2For an elaborate review see and references therein. then one can show that the charge density is only positive in the case of spin-1/2, where the Bhabha equation is the Dirac equation. Hence, only in the spin-1/2 representation one can construct a particle interpretation for the Bhabha wave equation starting from the charge current. Because the Bhabha equations are derived from fairly basic principles such as Lorentz covariance and restriction to first-order space and time derivatives,<sup>3</sup><sup>3</sup>3In fact the Bhabha equations can be seen as the relativistic counterparts of the non-relativistic wave equations (2.1), which can be derived from general principles such as Galilean covariance and ‘minimality’ . we do not think that the issue can be easily resolved by resorting to other relativistic wave equations.
If we consider the energy density for the class of Bhabha wave equations, then one can show that the energy density is only positive in the case of spin-0 and spin-1.<sup>4</sup><sup>4</sup>4The proof that the charge density is not positive definite for integer spin representations and that the energy density is not positive definite for half-integer spin representations can be found in . That the charge density is not positive definite for spin higher than 1/2, is discussed in . Akhiezer and Berestetskii mention without proof that for spins higher than one, neither the charge density nor energy density is positive definite \[73, p. 240\]. In the spin-0 case and the spin-1 case the Bhabha wave equation reduces to the first-order Duffin-Kemmer-Petiau (DKP) wave equation , which is completely equivalent with the familiar second order Klein-Gordon equation and Proca equations.
Hence, it could be tempting to try to construct a particle interpretation from the energy–momentum tensor for spin-0 and spin-1. This is indeed what has been done by Ghose et al. . Together with a quantum mechanical particle interpretation, Ghose et al. then also devised a pilot-wave model for the DKP equation. Both the quantum mechanical particle interpretation and the pilot-wave model display formal similarities with the equations for the Dirac theory.
Ghose et al. also extended the particle interpretation to massless spin-0 and spin-1 particles. In this case the particles are described by the first-order Harish-Chandra theory , which can be obtained from the DKP theory only by minimal modifications, and which is equivalent with the massless Klein-Gordon theory and Maxwell’s theory.
We will discuss the particle interpretation of Ghose et al. in detail in Section 3.3. In particular, we will indicate some features which make it hard to maintain this particle interpretation as a valid description for relativistic spin-0 and spin-1 particles. We argue that the quantum mechanical particle interpretation for the DKP equation can at best be regarded valid for sufficiently low energies, because in the non-relativistic limit it reduces to the particle interpretation for non-relativistic spin-0 and spin-1 particles. This reasoning does of course not apply to the massless case, for which the wavefunctions propagate at the speed of light.
As a result, also the corresponding pilot-wave interpretation for the DKP equation can only be considered as a valid approximation in the non-relativistic limit. In general the pilot-wave interpretation, or we better call it a trajectory model from now on, should rather be regarded as describing the tracks of energy flow and all the results reported for this trajectory model should be reinterpreted as such. In this way the tracks of energy flow merely provide a visualization of processes, they do not entail some new interpretation for quantum theory. But because such visualizations are interesting on their own, we extend the trajectory model to many particles.
Of course, as is well known, the problem of assigning particle probability densities to relativistic wave equations and the associated difficulties with the interpretation of the negative energy states (localization of a particle within its Compton wavelength implies the appearance of negative energy states), led to the conception of field theory. In field theory the notion of fields rather than the notion of particles is fundamental. Hence, as in the standard interpretation, a pilot-wave interpretation in terms of field beables might be better suited for dealing with high energy phenomena. We will see in the following chapter that such an interpretation in terms of fields is perfectly possible. Although the pilot-wave interpretation for fermionic fields still has to be developed further, the simplicity of the pilot-wave interpretation for bosons in terms of field beables is in striking contrast with the difficulties that are encountered when trying to develop a pilot-wave interpretation in terms of particle beables. Because of the striking simplicity of the field description and because a particle interpretation has even been abandoned in favor of a field interpretation already at the level of standard quantum theory,<sup>5</sup><sup>5</sup>5Of course one still has the notion of particles in Fock space. But in Fock space states are composed of states with definite momenta and if one tries to construct states which are localized in physical 3-space one encounters violations of causality . we consider the pilot-wave approach in terms of fields as fundamental. Nevertheless, a particle interpretation may still serve well as an illustrative model for the description of low energy phenomena.
### 3.2 Spin-1/2 relativistic quantum mechanics
#### 3.2.1 Massive spin-1/2: The Dirac formalism
The Dirac equation reads<sup>6</sup><sup>6</sup>6In this chapter we work in units in which $`\mathrm{}=c=1`$. The Lorentzian indices, which are denoted by $`\mu ,\nu ,\mathrm{}`$, are raised and lowered by the metric $`g_{\mu \nu }=\text{diag}(1,1,1,1)`$. The index $`0`$ denotes the time index and the indices $`i,j,\mathrm{}`$ denote the spatial index.
$$(i\gamma ^\mu _\mu m)\psi =0,$$
(3.1)
with $`m`$ the mass of the particle and $`\gamma ^\mu `$ the Dirac matrices which satisfy the commutation relations $`\gamma ^\mu \gamma ^\nu +\gamma ^\nu \gamma ^\mu =2g^{\mu \nu }`$. Equivalently one can write the Dirac equation in the Schrödinger form
$$i_0\psi =(i\alpha ^i_i+m\beta )\psi ,$$
(3.2)
with $`\alpha ^i=\gamma ^0\gamma ^i`$ and $`\beta =\gamma ^0`$.
The current $`j^\mu =\overline{\psi }\gamma ^\mu \psi `$ (which differs from the charge current by a factor $`e`$) is a conserved and future-causal Lorentz 4-vector . Hence, this current has a positive time component in every Lorentz frame. This allowed Dirac to interpret the density $`j^0(𝐱,t)=\psi ^{}(𝐱,t)\psi (𝐱,t)`$ as the probability density for a spin-1/2 particle to be detected at the position $`𝐱`$ at the time $`t`$. The integral curves of the 4-vector $`j^\mu `$ can then be seen as the flowlines of the ‘particle detection probability’. Because the current $`j^\mu `$ is future-causal it is guaranteed that these probability flowlines are time-like. Hence, the charge current can be a given a particle interpretation, which meets all basic relativistic requirements.
The pilot-wave interpretation now proceeds in the same way as in the non-relativistic case . A structureless point particle (the particle beable) is introduced for which the 4-vector velocity field is given by $`u^\mu =j^\mu /\sqrt{j^\nu j_\nu }`$. The possible trajectories $`x^\mu (\tau )`$ are the integral curves of the velocity field, i.e. they are solutions to the guidance equation
$$\frac{dx^\mu }{d\tau }=u^\mu .$$
(3.3)
The trajectory of a particle beable is uniquely determined by the specification of an initial configuration $`x^\mu (\tau _0)`$. Equivalently one can write the trajectories as curves $`𝐱(t)`$ which are found by solving the guidance equation
$$\frac{dx^i}{dt}=\frac{u^i}{u^0}=\frac{j^i}{j^0}.$$
(3.4)
The trajectory of a particle beable is then uniquely determined by the specification of an initial position $`𝐱(t_0)`$. Because the vector $`u^\mu `$ is future-causal, the motions are future-causal. In an ensemble the probability density of the particle beables is given by $`j^0=\psi ^{}\psi `$.
In the non-relativistic limit the pilot-wave model for the Dirac equation reduces to the pilot-wave interpretation for non-relativistic quantum mechanics as presented in Chapter 1 .
If an interaction with an electromagnetic field $`V_\mu `$ is introduced, through the minimal coupling prescription $`_\mu D_\mu =_\mu +ieV_\mu `$, then the pilot-wave interpretation as described above is still applicable . This is because the charge current is still conserved (an electromagnetic field does not posses charge and hence cannot exchange charge with a charged particle) and because the charge current is still future-causal (the charge current contains no derivatives and hence this current retains its form after minimal coupling).
It now seems that the trajectory interpretation presents no problem at all. As noted by Holland , the particle trajectories are always well defined, regardless of whether or not the state contains negative energy contributions. However, this success is only a deceptive appearance. Also in the pilot-wave interpretation negative energy states require a meaningful interpretation. This is because the problem of a possible radiation catastrophe manifests itself already at the level of the wavefunction. Hence, in order to prevent a positive energy wavefunction to lose energy by radiative transitions to lower and lower energies, one has to give a meaningful interpretation to the negative energy states. One could for example adopt the original idea by Dirac and assume a Dirac sea where every negative energy state is occupied, so that due to the Pauli principle, a transition to a negative energy state is impossible. This then requires a many-particle approach and in the corresponding pilot-wave model we should then consider a many-particle wavefunction describing an infinite, although fixed, number of particles, i.e. all the particles in the Dirac sea (the negative energy particles) plus the number of positive energy particles that one wants to describe. This is also the way Bohm and Hiley viewed the pilot-wave interpretation for the Dirac theory . When dealing with specific problems, the total wavefunction will factorize and it will be sufficient to consider a finite number of particles (the many-particle case is reviewed in the following section).<sup>7</sup><sup>7</sup>7Recently the pilot-wave interpretation which explicitly incorporates every particle in the Dirac sea was re-derived by Colin as the continuum limit of the stochastic Bell model . However, although the idea of a Dirac sea may serve well to describe spin-1/2 particles with electromagnetic interaction, it remains the question of how this model could be adopted to give an account for other types of interaction such as weak interaction (recall the standard example of beta decay \[83, p. 144\]). Of course, particle creation and annihilation finds a natural home in quantum field theory (there is no need for a Dirac sea), and hence a pilot-wave interpretation for quantum field theory is desired.
We turn to quantum field theory in the next two chapters. For the moment we continue to discuss the non-quantized Dirac theory a little more. In the rest of the section we recall the many-particle Dirac formalism. In the context of this formalism we can make some statements about the possibility of formulating a Lorentz invariant pilot-wave model. Similar statements will also apply in the context of quantum field theory.
#### 3.2.2 The many-particle Dirac formalism
The one-particle Dirac formalism is extended to many-particles by the introduction of an $`N`$-particle wavefunction $`\psi _{r_1\mathrm{}r_N}(𝐱_1,\mathrm{},𝐱_N,t)`$ with $`N`$ spin indices. The wavefunction is assumed to be anti-symmetric in order to satisfy the Pauli principle. We also introduce operators $`\gamma _{(r)}^\mu ,\alpha _{(r)}^i`$ and $`\beta _{(r)}`$ which operate only on the $`r^{\text{th}}`$ spin index, belonging to the $`r^{\text{th}}`$ particle, e.g. $`\gamma _{(r)}^\mu =\mathbb{𝟙}\mathrm{}\gamma ^\mu \mathrm{}\mathbb{𝟙}`$ with $`\gamma ^\mu `$ at the $`r^{\text{th}}`$ place in the product . The Schrödinger form of the Dirac equation for the $`N`$-particle system then reads
$$i_0\psi =\underset{r=1}{\overset{N}{}}\left(i\alpha _{(r)}^i_i^{(r)}+\beta _{(r)}m_r\right)\psi $$
(3.5)
where $`_i^{(r)}=/(𝐱_r)^i`$ and $`m_r`$ is the mass of the $`r^{\text{th}}`$ particle. The tensor current is defined as
$$j^{\mu _1\mathrm{}\mu _N}=\psi ^{}\gamma _{(1)}^0\gamma _{(1)}^{\mu _1}\mathrm{}\gamma _{(N)}^0\gamma _{(N)}^{\mu _N}\psi $$
(3.6)
and satisfies the conservation equation
$$_0j^{0_1\mathrm{}0_N}+\underset{r=1}{\overset{N}{}}_{i_r}^{(r)}j^{0_1\mathrm{}i_r\mathrm{}0_N}=0.$$
(3.7)
with $`j^{0_1\mathrm{}0_N}=\psi ^{}\psi `$ a positive quantity and $`j^{0_1\mathrm{}i_r\mathrm{}0_N}=\psi ^{}\alpha _{(r)}^i\psi `$.
In the pilot-wave interpretation the velocity field for the $`r^{\text{th}}`$ particle beable is given by
$$u_r^{\mu _r}=\frac{j^{0_1\mathrm{}\mu _r\mathrm{}0_N}}{\sqrt{j^{0_1\mathrm{}\nu _r\mathrm{}0_N}j_{0_1\mathrm{}\nu _r\mathrm{}0_N}}},$$
(3.8)
so that the corresponding guidance equation reads
$$\frac{dx_r^i}{dt}=\frac{\psi ^{}\alpha _{(r)}^i\psi }{\psi ^{}\psi }.$$
(3.9)
Because
$$j^{0_1\mathrm{}\mu _r\mathrm{}0_N}j_{0_1\mathrm{}\mu _r\mathrm{}0_N}0$$
(3.10)
the motion of the particles is future-causal. The distribution of the particle beables in an ensemble is assumed to be the equilibrium distribution, i.e. $`j^{0_1\mathrm{}0_N}=\psi ^{}\psi `$.
#### 3.2.3 Note on Lorentz covariance
In the one-particle case, the pilot-wave interpretation is covariant. First, the trajectories have a covariant meaning because the defining velocity field $`u^\mu `$ transforms as a Lorentz 4-vector under Lorentz transformations. Second, also the probability interpretation is covariant. To see this, consider an arbitrary spacelike hypersurface $`\sigma (x)`$, with $`n_\mu (x)`$ the future-oriented unit normal. The positive Lorentz scalar $`j^\mu (x)n_\mu (x)`$ is then interpreted as the probability for a particle to cross the spacelike hypersurface $`\sigma (x)`$ at $`x`$ in any frame . In the case of an equal-time hyperplane the crossing probability is given by $`j^0`$.
In the multi-particle case, the pilot-wave interpretation is not Lorentz covariant. The velocity fields (3.8) do not transform as 4-vectors. In addition, as shown by Berndl et al. , equilibrium can in general not hold simultaneously in all Lorentz frames. By this is meant the following. If we assume the density of crossings through an equal-time hyperplane in a particular frame to be given by $`\psi ^{}\psi `$, then the density of crossings of crossings will be given by $`\psi ^{}\psi `$ for any other equal-time hyperplane in this frame. But the density of crossings of crossings through an equal-time hyperplane in another Lorentz frame will in general not equal $`\psi ^{}\psi ^{}`$, with $`\psi ^{}`$ the wavefunction in the other frame. Berndl et al. showed this feature must hold, not only for the pilot-wave model presented above, but for any pilot-wave model for the many-particle Dirac theory. An exception occurs when the wavefunction is a product wavefunction of one-particle wavefunctions.
Bohm and Hiley accepted the idea of a preferred Lorentz frame to which the pilot-wave interpretation should be formulated . The same opinion is hold by Valentini. Valentini argues that the natural symmetry of pilot-wave theory is Aristotelian invariance (because of the first-order character of the guidance law) and that hence the search for a Lorentz covariant pilot-wave model is misguided . According to this view, the Lorentz covariance of the one-particle Dirac theory and the Galilean covariance of the pilot-wave interpretation of non-relativistic quantum mechanics are,<sup>8</sup><sup>8</sup>8The Galilean invariance of pilot-wave theory for non-relativistic quantum theory is discussed in \[12, pp. 122-124\]. despite appearances, not the fundamental symmetries. With Aristotelian invariance as the fundamental symmetry, there should be a preferred class of Aristotelian inertial frames<sup>9</sup><sup>9</sup>9This is a class of reference frames which are connected by Aristotelian transformations. Aristotelian transformations being time-independent transformations of 3-space, such as translations or rotations. relative to which the pilot-wave interpretation should be formulated. In this class of Aristotelian inertial frames, the interactions between the different particles may be instantaneous (due to the nonlocality of pilot-wave theory at the subquantum level), but no causality paradoxes arise, because the notions of cause and effect only make sense with respect to this frame.
On the other hand, Berndl et al. express the desire for a Lorentz invariant pilot-wave model . In they provide an example of how some notion of covariance could be introduced. In their model they introduce a particular space-time foliation as an additional space-time structure. This particular space-time foliation would then be determined by some covariant law, possibly depending on the Dirac wavefunction. Then, instead of writing the pilot-wave interpretation with respect to a particular frame as in the view of Bohm, Hiley and Valentini, it should be written with respect to this particular space-time foliation. If the equilibrium density is assumed on one of the leaves of the foliation, it is guaranteed that the model reproduces the predictions of standard quantum theory.
There exist still other pilot-wave models which are Lorentz covariant . However these models lack a probability interpretation, which makes it difficult to relate these models to standard quantum theory. The model by Squires could be ruled out from the start simply because it is a local model and hence in contradiction with the empirically verified violations of the Bell inequalities.
We will further not consider the possibility of constructing a Lorentz covariant pilot-wave model. The merit of the model by Berndl et al. is that it provides a counter example to claims that a covariant pilot-wave interpretation is impossible. But as Berndl et al. indicate themselves , any theory can be made Lorentz covariant by the incorporation of additional structure. If the pilot-wave interpretation is then devised to reproduce the quantum statistics, it remains a mere guessing what the additional structure should look like. In the pilot-wave models for field theory that we consider similar remarks apply. The pilot-wave models are not Lorentz covariant at the subquantum level, but at the quantum level they will yield the same statistics as the conventional interpretation and hence at the quantum level the models are Lorentz covariant.
### 3.3 Spin-0 and spin-1 relativistic quantum mechanics
In this section we consider the particle interpretation for relativistic spin-0 and spin-1 bosons proposed by Ghose et al. Because the Duffin-Kemmer-Petiau (DKP) formalism and the Harish-Chandra (HC) formalism are perhaps less known, we review the essential elements in the following two sections. A more elaborate discussion of the DKP formalism and the HC formalism can respectively be found in and , or in the review by Ghose .
#### 3.3.1 Massive spin-0 and spin-1: The Duffin-Kemmer-Petiau formalism
The DKP equation for a particle with mass $`m`$ reads
$$(i\beta ^\mu _\mu m)\psi =0,$$
(3.11)
with adjoint equation
$$i^\mu \overline{\psi }\beta _\mu +m\overline{\psi }=0,$$
(3.12)
where $`\overline{\psi }=\psi ^{}\eta _0`$ with $`\eta _0=2\beta _0^21`$, and the DKP matrices $`\beta ^\mu `$ satisfy the commutation relations
$$\beta ^\mu \beta ^\nu \beta ^\lambda +\beta ^\lambda \beta ^\nu \beta ^\mu =\beta ^\mu g^{\nu \lambda }+\beta ^\lambda g^{\nu \mu }.$$
(3.13)
There are three inequivalent irreducible representations of the $`\beta ^\mu `$, one is $`10\times 10`$ and describes spin-1 bosons, another one is $`5\times 5`$ which describes spin-0 bosons and the third one is the trivial $`1\times 1`$ representation.
The DKP equations can be derived from the Lagrangian density
$$_{DKP}=\frac{i}{2}(\overline{\psi }\beta _\mu ^\mu \psi ^\mu \overline{\psi }\beta _\mu \psi )m\overline{\psi }\psi .$$
(3.14)
The corresponding conserved symmetrized energy–momentum tensor reads
$$\mathrm{\Theta }_{DKP}^{\mu \nu }=m\overline{\psi }(\beta ^\mu \beta ^\nu +\beta ^\nu \beta ^\mu g^{\mu \nu })\psi .$$
(3.15)
The conserved charge current is
$$s_{DKP}^\mu =e\overline{\psi }\beta ^\mu \psi $$
(3.16)
with $`e`$ the charge of the particle.
The DKP equation (3.11) can be written in the following equivalent Hamiltonian form
$`i_0\psi `$ $`=`$ $`(i\stackrel{~}{\beta }^i_i+m\beta _0)\psi ,`$ (3.17)
$`i\beta ^i\beta _0^2_i\psi `$ $`=`$ $`m(1\beta _0^2)\psi `$ (3.18)
with $`\stackrel{~}{\beta }^i=\beta ^0\beta ^i\beta ^i\beta ^0`$. The first equation is a Schrödinger-like equation and the second equation has to be regarded as an additional constraint on the wavefunction $`\psi `$. Only when the two equations (3.17) and (3.18) are taken together, they are equivalent with the covariant form (3.11).
That the constraint equation is compatible with the equations of motion can be seen by writing (3.17) and (3.18) in operator form. From the Schrödinger-like equation (3.17) we find the free Hamiltonian operator
$$\widehat{H}=\stackrel{~}{𝜷}\widehat{𝐩}+m\beta _0,$$
(3.19)
with $`\widehat{𝐩}=i\mathbf{}`$ the momentum operator. Physical states have to satisfy the additional constraint (3.18), which can be written as
$$\widehat{C}\psi =0,$$
(3.20)
with $`\widehat{C}`$ an idempotent operator, i.e. $`\widehat{C}^2=\widehat{C}`$, given by
$$\widehat{C}=\frac{1}{m}𝜷\widehat{𝐩}\beta _0^2+1\beta _0^2=1\frac{1}{m}\widehat{H}\beta _0.$$
(3.21)
Because $`\widehat{H}\beta _0\widehat{H}=m\widehat{H}`$, we have that $`\widehat{C}\widehat{H}=0`$. If a wavefunction initially, say at time $`t=0`$, satisfies the constraint (3.18), i.e. $`\widehat{C}\psi (𝐱,0)`$, then the wavefunction will also satisfy the constraint at a later time, because in the Heisenberg picture we have $`\widehat{C}\psi (𝐱,t)=\widehat{C}\mathrm{exp}\left(i\widehat{H}t\right)\psi (𝐱,0)=0`$. Hence the constraint (3.18) is compatible with the equation of motion (3.17).
With the help of (3.20) one can further show that
$$\widehat{H}^2\psi =(\widehat{𝐩}^2+m^2)\psi .$$
(3.22)
Hence every component of the DKP wavefunction satisfies the massive Klein-Gordon equation. Contrary to the Dirac case, the equality is not valid on the operator level, i.e. $`\widehat{H}^2(\widehat{𝐩}^2+m^2)`$. Rather one has the property
$$\widehat{H}^3=\widehat{H}(\widehat{𝐩}^2+m^2).$$
(3.23)
The equivalence of the DKP wave equation with the Klein-Gordon equation in the spin-0 representation and the Proca equations in the spin-1 representation can be shown in a representation independent way . However, it is instructive to see how this equivalence arises in the explicit matrix representations which are given in Appendix A.
In the spin-0 representation, the constraint equation (3.18) implies a DKP wavefunction $`\psi `$ with the following five components: $`\omega `$, $`_1\varphi /\sqrt{m}`$, $`_2\varphi /\sqrt{m}`$, $`_3\varphi /\sqrt{m}`$, $`\sqrt{m}\varphi `$ with $`\omega `$ and $`\varphi `$ two scalar wavefunctions. Equation (3.17) further implies $`\omega =_0\varphi /\sqrt{m}`$. In this way the number of independent components of the wavefunction $`\psi `$ is reduced to one and we can write
$$\psi =\frac{1}{\sqrt{m}}\left(\begin{array}{c}_\mu \varphi \\ m\varphi \end{array}\right).$$
(3.24)
The Schrödinger equation (3.17) then reduces to the massive Klein-Gordon (KG) equation for $`\varphi `$
$$\mathrm{}\varphi +m^2\varphi =0.$$
(3.25)
By substitution of the wavefunction (3.24) into the DKP Lagrangian (3.14), the DKP energy–momentum tensor (3.15) and the DKP charge current (3.16), we obtain respectively the KG Lagrangian for the KG wavefunction $`\varphi `$
$$_{KG}=^\alpha \varphi _\alpha \varphi ^{}m^2\varphi ^{}\varphi ,$$
(3.26)
the KG energy–momentum tensor
$$\mathrm{\Theta }_{KG}^{\mu \nu }=^\mu \varphi ^\nu \varphi ^{}+^\mu \varphi ^{}^\nu \varphi g^{\mu \nu }_{KG}$$
(3.27)
and the KG charge current
$$s_{KG}^\mu =ie\left(\varphi ^{}^\mu \varphi \varphi ^\mu \varphi ^{}\right).$$
(3.28)
In the spin-1 representation we can take the ten components of the DKP wavefunction as:
$$\psi =(𝐄,𝐁,m𝐀,mA_0)^T/\sqrt{m}.$$
(3.29)
Equation (3.18) then implies the following relations
$$\mathbf{}𝐄=m^2A_0,𝐁=\mathbf{}\times 𝐀.$$
(3.30)
The equation (3.17) leads to the relations
$$_0𝐄=\mathbf{}\times 𝐁+m^2𝐀,_0𝐁=\mathbf{}\times 𝐄,𝐄=\mathbf{}A_0_0𝐀.$$
(3.31)
The equations in (3.30) and (3.31) are recognized as the Proca equations
$$_\mu G^{\mu \nu }=m^2A^\nu $$
(3.32)
for $`A^\mu =(A_0,𝐀)`$, with $`G^{\mu \nu }=^\mu A^\nu ^\nu A^\mu `$.
If the Kemmer wavefunction (3.29) is substituted into the DKP Lagrangian (3.14), the DKP energy–momentum tensor (3.15) and the DKP charge current (3.16), we obtain with the help of the relations (3.30) and (3.31) respectively the Proca Lagrangian
$$_P=\frac{1}{2}G_{\mu \nu }^{}G^{\mu \nu }+m^2A_\mu ^{}A^\mu ,$$
(3.33)
the symmetrized Proca energy–momentum tensor
$$\mathrm{\Theta }_P^{\mu \nu }=G^{\mu \alpha }G_\alpha ^\nu G^{\mu \alpha }G_\alpha ^\nu +m^2(A^\mu A^\nu +A^\mu A^\nu )g^{\mu \nu }_P$$
(3.34)
and the Proca charge current
$$s_P^\mu =ie\left(G^{\mu \nu }A_\nu G^{\mu \nu }A_\nu ^{}\right).$$
(3.35)
#### 3.3.2 Massless spin-0 and spin-1: The Harish-Chandra formalism
We cannot just take the limit $`m0`$ in the DKP theory to describe massless bosons. Nevertheless, one can describe massless spin-0 and spin-1 bosons in a first-order formalism in close analogy with the DKP formalism. This formalism was developed by Harish-Chandra (HC). In this section we will only review the essential elements of this theory. This theory can be cast in a representation independent form, but as in the case of the DKP theory, we will often fall back to the explicit representation given in Appendix A.
In order to describe massless spin-0 and spin-1 bosons, Harish-Chandra proposed a modification of the DKP equations by replacing the mass $`m`$ by $`m\gamma `$,
$`(i\beta ^\mu _\mu m\gamma )\psi `$ $`=`$ $`0,`$ (3.36)
$`i^\mu \overline{\psi }\beta _\mu +m\overline{\psi }\gamma `$ $`=`$ $`0,`$ (3.37)
with $`\gamma `$ a matrix that satisfies
$`\gamma ^2`$ $`=`$ $`\gamma ,`$ (3.38)
$`\gamma \beta ^\mu +\beta ^\mu \gamma `$ $`=`$ $`\beta ^\mu .`$ (3.39)
From (3.36), (3.38) and (3.39) one can derive the second order wave equation
$$\mathrm{}(\gamma \psi )=0.$$
(3.40)
Hence the wavefunction $`(\gamma \psi )`$ describes a massless boson.
In both the 10-dimensional and 5-dimensional representation we have two inequivalent choices for $`\gamma `$. For each representation we only presented one particular choice in Appendix A. With our choice for $`\gamma `$ in the 10-dimensional representation, massless spin-1 particles are described. If the wavefunction $`\psi `$ is real then $`\gamma \psi `$ describes an uncharged, massless spin-1 particle.<sup>10</sup><sup>10</sup>10In a general representation we cannot describe uncharged particles simply by assuming $`\psi `$ is real. A more general condition is needed then, which can be found in . Harish-Chandra called this condition the reality condition. In this case the HC theory is equivalent with the Maxwell’s theory for the electromagnetic field. With our choice for $`\gamma `$ in the 5-dimensional representation, massless spin-0 particles are described. As in the spin-1 case, an uncharged particle is described by a real wavefunction. The other choice for $`\gamma `$ in the 10-dimensional representation also describes massless spin-0 particles. The other choice for $`\gamma `$ in the 5-dimensional representation is unphysical.
The wave equations are invariant under transformations $`\psi \psi +(1\gamma )\stackrel{~}{\psi }`$ with $`\stackrel{~}{\psi }`$ satisfying
$$i\beta ^\mu _\mu \stackrel{~}{\psi }=0.$$
(3.41)
This invariance corresponds to the gauge invariance in the spin-1 case. The wave equations can be derived from the following gauge invariant Lagrangian density
$$_{HC}=\frac{i}{2}(\overline{\psi }\beta _\mu ^\mu \psi ^\mu \overline{\psi }\beta _\mu \psi )m\overline{\psi }\gamma \psi .$$
(3.42)
The symmetrized energy–momentum tensor is given by
$`\mathrm{\Theta }_{HC}^{\mu \nu }`$ $`=`$ $`m\overline{\psi }(\beta ^\mu \beta ^\nu +\beta ^\nu \beta ^\mu g^{\mu \nu })\gamma \psi `$ (3.43)
$`=`$ $`m\psi ^{}\gamma \eta _0(\beta ^\mu \beta ^\nu +\beta ^\nu \beta ^\mu g^{\mu \nu })\gamma \psi .`$
Note that the factor $`m`$ does not represent the mass of the particle (the mass is zero), but a constant that can be removed by a suitable normalization of $`\gamma \psi `$. The conserved charge current is also formally the same as in the massive case
$$s_{HC}^\mu =e\overline{\psi }\beta ^\mu \psi .$$
(3.44)
For uncharged bosons $`\psi `$ the charge current is zero. Not only the charge $`e`$ is then zero, but because $`\psi `$ satisfies the reality condition also the current $`\overline{\psi }\beta ^\mu \psi `$ is identically zero.
The HC equation (3.36) can be written in the following equivalent Hamiltonian form
$`i_0(\gamma \psi )`$ $`=`$ $`i\stackrel{~}{\beta }^i_i(\gamma \psi ),`$ (3.45)
$`i\beta ^i\beta _0^2_i\psi `$ $`=`$ $`m(1\beta _0^2)\gamma \psi .`$ (3.46)
Similarly as in the massive case, the first equation is a Schrödinger-like equation and the second equation has to be regarded as an additional constraint on the wavefunction $`\psi `$. Although it follows from (3.46) that
$$i\beta ^i\beta _0^2_i(\gamma \psi )=0,$$
(3.47)
only the set (3.45), (3.46) and not the set (3.45), (3.47) is equivalent with the wave equation (3.36). Similarly as in the massive case, one can show that the constraint (3.46) is compatible with the Schrödinger equation (3.45) by using (3.47). The Hamiltonian operator $`\widehat{H}`$ and the constraint operator $`\widehat{C}`$ can be obtained from the corresponding operators in the massive case, cf. (3.19) and (3.20), simply by putting $`m=0`$.
Just as in the massive case the equivalence of the HC theory with the massless Klein-Gordon theory in the spin-0 representation and with the Maxwell theory in the uncharged spin-1 representation is easily shown with the explicit representation for the matrices $`\gamma `$ in Appendix A. The action of the matrices $`\gamma `$ on the wavefunctions (3.24) and (3.29), results in a projection on the mass independent components of the wavefunctions, i.e. the massless states $`\gamma \psi `$ are obtained from the massive states $`\psi `$ (given by (3.24) and (3.29) just by putting the mass $`m`$ equal to zero. This simplicity is the reason why we only presented the particular representations for $`\gamma `$ given in Appendix A.
### 3.4 Trajectory models for spin-0 and spin-1 bosons
In the context of the massive Klein-Gordon theory, de Broglie initially entertained the idea of constructing a pilot-wave model for spin-0 particles by giving a particle interpretation to the charge current . Later, Vigier gave a particle interpretation to the charge current in the DKP formalism, thereby extending the pilot-wave model of de Broglie to account also for spin-1 particles .
However, as is well known, a quantum mechanical particle interpretation for the charge current for spin-0 or spin-1 bosons is in general untenable because the current is not always future-causal. Even a restriction of the positive energy part of the Hilbert space presents no solution. There exist examples of superpositions of positive energy eigenstates for which the charge current is spacelike, or for which the charge density may become negative in certain space-time regions . In the case of an uncharged boson the situation is even worse because then the charge current is zero.
In order to circumvent the problems in associating a particle interpretation to the charge current, Ghose et al. proposed to start instead from the energy–momentum tensor, from which future-causal 4-vectors can be constructed. They did this as follows. Let the tensor $`\mathrm{\Theta }^{\mu \nu }`$ represent the DKP energy–momentum tensor $`\mathrm{\Theta }_{DKP}^{\mu \nu }`$ in the massive case and the CH energy–momentum tensor $`\mathrm{\Theta }_{HC}^{\mu \nu }`$ in the massless case and let $`n^\mu `$ be a constant future-causal 4-vector (below possible examples for the vector $`n^\mu `$ are presented). By contraction of the energy–momentum tensor $`\mathrm{\Theta }^{\mu \nu }`$ and the 4-vector $`n^\mu `$ we obtain a 4-vector
$$j^\mu =\mathrm{\Theta }^{\mu \nu }n_\nu $$
(3.48)
which is future-causal and conserved.<sup>11</sup><sup>11</sup>11The fact that $`j^\mu `$ is future-causal can be derived as follows. We use the notation $`\stackrel{~}{\psi }`$ to represent either the wavefunction $`\psi `$ for massive bosons or $`\gamma \psi `$ for massless bosons. Because $`\mathrm{\Theta }^{00}=m\stackrel{~}{\psi }^{}\stackrel{~}{\psi }0`$ and $`\mathrm{\Theta }^{0\mu }\mathrm{\Theta }_{0\mu }0`$ (which can be verified in the explicit representations used in Appendix A), the vector $`\mathrm{\Theta }^{0\mu }=\delta _\nu ^0\mathrm{\Theta }^{\nu \mu }`$ is future-causal. Because the product of two future-causal vectors is positive, $`j^0=\mathrm{\Theta }^{0\mu }n_\mu 0`$ for $`n^\mu `$ future-causal. The fact that $`j^\mu j_\mu 0`$ can be seen if we perform a Lorentz transformation such that $`\mathrm{\Lambda }_\nu ^\mu n^\nu =\delta _0^\mu `$ because then $`j^\mu j_\mu =\mathrm{\Theta }^{0\mu }\mathrm{\Theta }_{0\mu }^{}`$, with $`\mathrm{\Theta }^{\mu _1\mu _2}=\mathrm{\Lambda }_{\nu _1}^{\mu _1}\mathrm{\Lambda }_{\nu _2}^{\mu _2}\mathrm{\Theta }^{\nu _1\nu _2}`$ and $`\mathrm{\Theta }^{0\mu }\mathrm{\Theta }_{0\mu }^{}`$ is positive as it has the same form as the positive quantity $`\mathrm{\Theta }^{0\mu }\mathrm{\Theta }_{0\mu }`$ (just replace $`\stackrel{~}{\psi }(𝐱,t)`$ in the energy–momentum tensor by $`\stackrel{~}{\psi }^{}(𝐱^{},t^{})`$, where the accents refer to quantities in the new Lorentz frame).
Hence, the current $`j^\mu `$ satisfies all the properties required for a particle current. Ghose et al. interpret the current $`j^\mu `$ as a particle current, and associate a pilot-wave interpretation to it along the lines of the pilot-wave interpretation of the Dirac theory. The velocity field of the particle beable is given by $`u^\mu =j^\mu /\sqrt{j^\nu j_\nu }`$, so that the possible trajectories $`𝐱(t)`$ are solutions to the guidance equation
$$\frac{dx^i}{dt}=\frac{u^i}{u_0}=\frac{j^i}{j^0}.$$
(3.49)
For an ensemble of particles all described by the same DKP or HC wavefunction, the probability density of the particle beables is then given by $`j^0`$.
In the case that $`n^\mu =\delta _0^\mu `$, the probability density is given by the energy density $`j^0=\mathrm{\Theta }^{00}=m\psi ^{}\psi `$ and the guidance equation reads
$$\frac{dx^i}{dt}=\frac{\mathrm{\Theta }^{i0}}{\mathrm{\Theta }^{00}}=\{\begin{array}{cc}\frac{\psi ^{}\stackrel{~}{\beta }^i\psi }{\psi ^{}\psi }\hfill & \text{in the massive case}\hfill \\ \frac{\psi ^{}\gamma \stackrel{~}{\beta }^i\gamma \psi }{\psi ^{}\gamma \psi }\hfill & \text{in the massless case}\hfill \end{array}$$
where $`\stackrel{~}{\beta }^i=\beta ^0\beta ^i\beta ^i\beta ^0`$. Hence, in the case that $`n^\mu =\delta _0^\mu `$ the pilot-wave equations look similar to the pilot-wave equations in the Dirac theory.
Let us now turn to the definition of $`n^\mu `$. Ghose et al. proposed the vector $`n^\mu `$ to be the vector normal to a particular spacelike foliation . The vector $`n^\mu `$ is then generally space-time dependent, but the program of associating trajectories to the current $`j^\mu `$ can still be carried out, provided some minor modifications are taken into account, such as the introduction of a covariant derivative on the foliation. However, in the following we will restrict our attention to foliations in terms of hyperplanes so that the vector $`n^\mu `$ is a constant vector. In this case one can alternatively regard the constant vector $`n^\mu `$ as the 4-velocity $`a^\mu `$ of some observer.
In fact it was Holland who initiated the construction of a conserved current from the energy–momentum tensor by contracting it with the four-velocity of an observer, but in the context of Maxwell’s theory and the KG theory . However, Holland was reluctant to regard the conserved current as a particle probability current and regarded the guidance equation in (3.49) as the defining formula for the tracks of energy flow. Holland gave a series of arguments supporting this view. Most of the arguments were against the notion of a photon as a localized object and do not apply to massive bosons for which the particle aspect is well accepted. An other argument by Holland is that, if for example $`a^\mu =\delta _0^\mu `$, then the boson density $`j^0`$ (for photons or massive spin-0 particles) would be given by the energy density, which it is manifestly not in the analogous formula for the quantized boson field.
We agree with Holland that the current $`j^\mu `$ should not be interpreted as a particle current. But we think that the main reasons for not doing so are the following.
Even if we adopt the view that the energy density can be seen as the particle density (we can take the energy density as an observable quantity, whether or not it is interpreted as a particle density), it is only in the frame at rest relative to the observer with velocity $`a^\mu `$ that we have that $`a^\mu =\delta _0^\mu `$ and that the probability density is the energy density. In general the quantity $`j^0=\mathrm{\Theta }^{0\nu }a_\nu `$ does not correspond to a known observable quantity. Even for another choice of the vector $`n^\mu `$ it is unclear to which known observable quantity $`j^0=\mathrm{\Theta }^{0\nu }n_\nu `$ could correspond.
Second, the current depends on the arbitrary choice of the observer (or for that matter on the arbitrary choice of the foliation), hence although the current $`j^\mu =\mathrm{\Theta }^{\mu \nu }a_\nu `$ is written in a covariant form it is not covariant in content. For example if we consider two observers, $`O`$ and $`O^{}`$, which describe the same system relative to the frame at which they are at rest, then observer $`O`$ associate the following velocity field to the system
$$u^\mu (x)=\frac{\mathrm{\Theta }^{\mu 0}(x)}{\sqrt{\mathrm{\Theta }^{\nu 0}(x)\mathrm{\Theta }_{\nu 0}(x)}}=\frac{\mathrm{\Theta }^{\mu \rho }(x^{})h_\rho }{\sqrt{\mathrm{\Theta }^{\nu \alpha }(x^{})h_\alpha \mathrm{\Theta }_{\nu \beta }^{}(x^{})h^\beta }},$$
(3.50)
where in the last equality we have written the vector $`u^\mu (x)`$ with respect to the observer $`O^{}`$; i.e. if $`\mathrm{\Lambda }`$ denotes the passive Lorentz transformation from the frame of observer $`O`$ to that of observer $`O^{}`$, we have $`\mathrm{\Theta }^{\mu _1\mu _2}(x^{})=\mathrm{\Lambda }_{\nu _1}^{\mu _1}\mathrm{\Lambda }_{\nu _2}^{\mu _2}\mathrm{\Theta }^{\nu _1\nu _2}(x)`$ and $`h^\mu =\mathrm{\Lambda }_\nu ^\mu \delta _0^\nu `$. It is clear that this last expression does not equal the velocity field of the system relative to observer $`O^{}`$, which is
$$u^\mu (x^{})=\frac{\mathrm{\Theta }^{\mu 0}(x^{})}{\sqrt{\mathrm{\Theta }^{\nu 0}(x^{})\mathrm{\Theta }_{\nu 0}^{}(x^{})}}.$$
(3.51)
Similarly the two observers would also not agree on the particle probability densities if they described the system each relative to their rest frame. This implies that even at the standard quantum level the notion of the preferred observer would be present, which is in fact a sufficient reason to abandon the approach. A similar objection was raised by Bohm et al. in the context of the electromagnetic field .
This last objection could be removed by trying to find a ‘covariant’ determination for the preferred observer or the preferred frame, in a way similar as was done in the multi-particle Dirac case by Dürr et al. (cf. Section 3.2.2). However with the distinction that in the model of Dürr et al. the notion of a preferred foliation was only present at the subquantum level of the pilot-wave interpretation (at the level of the particle beables) and not at the statistical level of the theory. In the model presented above, the current $`j^\mu `$ is fully determined by the choice of the 4-vector $`n^\mu `$, and hence although there may be some ‘covariant law’ which determines this vector, it plays an important role in the statistical predictions of the model through the distribution $`j^0=\mathrm{\Theta }^{0\mu }n_\mu `$.
An example of such a covariant vector, is the total energy–momentum 4-vector
$$P^\mu =_\sigma 𝑑\sigma _\nu \mathrm{\Theta }^{\mu \nu }=d^3x\mathrm{\Theta }^{\mu 0},$$
(3.52)
with $`\sigma `$ an arbitrary spacelike hypersurface (the vector is independent of the choice of hypersurface due to the fact that the energy–momentum tensor is conserved). This vector transforms as a Lorentz 4-vector under Lorentz transformations and is conserved and future-causal (as the continuous sum over the future-causal vectors $`\mathrm{\Theta }^{\mu 0}`$). We could then use the normalized constant 4-vector $`n^\mu =P^\mu /\sqrt{P_\nu P^\nu }`$ to contract the energy–momentum tensor with. Although the vector $`P^\mu `$ is nonlocally determined, we can introduce some notion of covariance in this way.<sup>12</sup><sup>12</sup>12For the massive Klein-Gordon field, Dewdney, Horton and Nesteruk developed still another covariant model starting from the energy–momentum tensor . In their model the flowlines were defined as the integral curves of the 4-vector $`W^\mu `$, with $`W^\mu `$ a future-causal eigenvector of the energy–momentum tensor and the probability density is given by the intrinsic energy density.
Finally, there is the problem that if we consider a massive boson interacting with an electromagnetic field, then although the total energy–momentum tensor of the massive boson and the electromagnetic field is conserved, they are not conserved separately. We discuss this in detail in Section 3.8. Hence, if we would then construct the current $`j^\mu =\mathrm{\Theta }_{DKP}^{\mu \nu }n_\nu `$ just as in the free case ($`\mathrm{\Theta }_{DKP}^{\mu \nu }`$ is the energy–momentum tensor for the massive boson, which is unaltered in form after minimal coupling), then this current is no longer conserved. Although a particle interpretation of the current $`j^\mu `$ is still formally possible (the density is positive and the flowlines are future-causal), it would require the notion of particle creation and annihilation along the flowlines and hence this would contravene with the current experimental knowledge.<sup>13</sup><sup>13</sup>13This has of course nothing to do with pair creation, because even for energies below the threshold of pair creation, particle creation and annihilation as implied by a particle interpretation of the current $`j^\mu `$ would occur.
We can conclude that, although it is already hard to maintain a particle interpretation of $`j^\mu =\mathrm{\Theta }^{\mu \nu }n_\nu `$ even in the free case (due to the problems with Lorentz covariance and due to the fact that in general there does not correspond an empirically known quantity associated with the density $`j^0`$), it becomes in general untenable in the interacting case. A particle interpretation may perhaps only be maintained in the non-relativistic limit, with an electromagnetic field that is sufficiently weak. This is because in the non-relativistic limit, the current $`j^\mu `$ reduces to the same non-relativistic current as in the free current (see Section 3.6), and hence in the non-relativistic limit the current $`j^\mu `$ is conserved again. Of course this reasoning does not apply for the electromagnetic field (or the equivalent Harish-Chandra field). However, we hold the view that in general the trajectory model should rather be regarded as describing the tracks of energy flow, as proposed by Holland, and that all the results reported for this trajectory model should be reinterpreted as such. For example, the trajectories for the electromagnetic field drawn in for the double slit experiment<sup>14</sup><sup>14</sup>14Similar trajectories were presented in , but outside the context of pilot-wave theory. and for reflection through a glass slab should be seen as the lines of energy flow.
The rest of the chapter is organized as follows. In the following section we give some remarks on the definition of an inner product. In Section 3.6 we consider the non-relativistic limit of the trajectory model. In Section 3.7 we extend the trajectory model to many-particles. Finally, in Section 3.8 we consider the minimal coupling in the DKP theory.
### 3.5 Note on the definition of the inner product
For the moment we have mainly focussed on the possibility of constructing a conserved, future-causal current which could then be interpreted as a particle current. Related to the problem of finding such a current is the problem to define an inner product. Usually the inner product for the Kemmer theory is defined as
$$\psi _1|\psi _2=_\sigma 𝑑\sigma _\mu \overline{\psi }_1\beta ^\mu \psi _2=d^3x\overline{\psi }_1\beta ^0\psi _2,$$
(3.53)
where $`\sigma `$ is an arbitrary spacelike hypersurface. The quantity $`\psi _1|\psi _2`$ is inspired by the form of the charge current (3.16). This quantity is Lorentz invariant (it is hypersurface independent) and further satisfies all the requirements for an inner product, except for one, namely positivity. For an uncharged system the situation is somehow worse because then $`\psi |\psi =0`$ for all wavefunctions $`\psi `$. The problem in defining a positive inner product is another problem which dissolves when passing to quantum field theory (e.g. in the functional Schrödinger picture in quantum field theory, the inner product of two wavefunctionals is well defined).
In the previous section we have seen that the construction of a conserved, future-causal current from the energy–momentum tensor, as opposed to using the charge current, also has its problems. One can now consider the question whether one encounters the same problems when one tries to construct an inner product, starting from the energy–momentum tensor, i.e. inspired by the the current $`j^\mu =\mathrm{\Theta }^{\mu \nu }n_\nu `$. The answer to this question is yes. Suppose we would define
$$\psi _1|\psi _2=_\sigma 𝑑\sigma _\mu \overline{\psi }_1(\beta ^\mu \beta ^\nu +\beta ^\nu \beta ^\mu g^{\mu \nu })\psi _2n_\nu .$$
(3.54)
In the case $`n^\mu `$ represents the normal on some preferred foliation, then one can check that the quantity $`\psi _1|\psi _2`$ satisfies all the requirements for an inner product, however it is obviously not Lorentz-invariant because $`n^\mu `$ does not transform as a 4-vector under Lorentz transformations. Now if we would take a Lorentz 4-vector for $`n^\mu `$, e.g.
$$n^\mu _\sigma 𝑑\sigma _\nu \overline{\psi }_1(\beta ^\mu \beta ^\nu +\beta ^\nu \beta ^\mu g^{\mu \nu })\psi _2$$
(3.55)
then the quantity $`\psi _1|\psi _2`$ would be Lorentz invariant and it would satisfy all requirements for an inner product, except for linearity.
### 3.6 Non-relativistic limit of the trajectory model for massive spin-0 and spin-1 bosons
In this section we consider the non-relativistic limit of the trajectory model for massive bosons. Because the DKP energy–momentum tensor reduces to the Klein-Gordon (KG) energy–momentum tensor in the spin-0 representation and to the Proca energy–momentum tensor in the spin-1 representation (cf. Section 3.3.1), we can simply consider the non-relativistic limits of the KG and Proca energy–momentum tensor.
#### 3.6.1 The spin-0 case
In order to take the non-relativistic limit of the Klein-Gordon equation (3.25), we substitute $`\varphi =e^{imt}\psi ^{}/\sqrt{2}m`$, where the energy of the wavefunction $`\psi ^{}`$ is much smaller than the rest energy $`m`$, in the KG equation. As is well known this leads to the Schrödinger equation for $`\psi ^{}`$ (e.g. ). The KG energy–momentum tensor reduces to
$`\mathrm{\Theta }_{KG}^{00}`$ $`=`$ $`|\psi ^{}|^2,\mathrm{\Theta }_{KG}^{i0}={\displaystyle \frac{1}{m}}\text{Im}\left(\psi ^{}_i\psi ^{}\right),`$
$`\mathrm{\Theta }_{KG}^{ij}`$ $`=`$ $`{\displaystyle \frac{1}{m^2}}\text{Re}(_i\psi ^{}_j\psi ^{})+\delta _{ij}\left({\displaystyle \frac{\text{Im}(\psi ^{}_0\psi ^{})}{m}}{\displaystyle \frac{_k\psi ^{}_k\psi ^{}}{2m^2}}\right)`$ (3.56)
in the non-relativistic limit. Because $`\mathrm{\Theta }_{KG}^{00}`$, $`\mathrm{\Theta }_{KG}^{i0}`$, $`\mathrm{\Theta }_{KG}^{ij}`$ are respectively of zero order, of first order and of third order in $`p/m`$ and because the components $`n^i`$ are at least of the same order as $`n^0`$ ($`n^\mu `$ is future-causal), the current $`j^\mu =\mathrm{\Theta }_{KG}^{\mu \nu }n_\nu `$ reduces to the conventional Schrödinger current in the non-relativistic limit (up to the constant factor $`n^0`$, which can be removed by renormalizing $`\psi ^{}`$)
$`j^0`$ $`=`$ $`\mathrm{\Theta }_{KG}^{00}=|\psi ^{}|^2,`$
$`j^i`$ $`=`$ $`\mathrm{\Theta }_{KG}^{i0}={\displaystyle \frac{1}{m}}\text{Im}(\psi ^{}_i\psi ^{}).`$ (3.57)
In this way the trajectory model associated with the energy–momentum tensor reduces to the pilot-wave interpretation for the non-relativistic Schrödinger equation originally presented by de Broglie and Bohm (cf. Section 2.2).
#### 3.6.2 The spin-1 case
In order to take the non-relativistic limit in the massive spin-1 case, we write the Proca equations as a KG equation for each component of the field $`A^\mu `$
$$\mathrm{}A^\mu +m^2A^\mu =0$$
(3.58)
with subsidiary condition
$$_\mu A^\mu =0.$$
(3.59)
We again separate the rest energy by putting $`A^\mu =e^{imt}\varphi ^\mu /\sqrt{2}m`$. The condition $`_\mu A^\mu =0`$ then reduces to $`\varphi ^0=_i\varphi ^i/im`$, which implies that we can take $`\varphi ^0`$ as the small component of the wavefunction $`\varphi ^\mu `$. The non-relativistic limit of equation (3.58) results in the Schrödinger equation for each component $`\varphi ^\mu `$. If we define the wavefunction $`\mathrm{\Phi }=(\varphi ^1,\varphi ^2,\varphi ^3)^T`$, then the non-relativistic limit of the Proca equations can be written as
$$i\frac{\mathrm{\Phi }}{t}=\frac{^2\mathrm{\Phi }}{2m}.$$
(3.60)
With similar assumptions on the vector $`n^\mu `$ as in the spin-0 case, the spin-1 current $`j^\mu =\mathrm{\Theta }_P^{\mu \nu }n_\nu `$ reduces to the non-relativistic current for a spin-1 particle as given in (2.21) in the non-relativistic limit:
$`j^0`$ $`=`$ $`\mathrm{\Theta }_P^{00}=\mathrm{\Phi }^{}\mathrm{\Phi }`$
$`j^i`$ $`=`$ $`\mathrm{\Theta }_P^{i0}={\displaystyle \frac{1}{m}}\text{Im}\left(\mathrm{\Phi }^{}_i\mathrm{\Phi }+\varphi ^i_j\varphi ^j\varphi ^j_j\varphi ^i\right).`$ (3.61)
As a side remark, we note that in the non-relativistic limit, the DKP charge current also reduces to the non-relativistic currents which were presented in Section 2.2, in both the spin-0 case and the spin-1 case.
### 3.7 The many-particle Duffin-Kemmer-Petiau formalism
In this section we generalize the one-particle trajectory model to many particles. The trajectory model is constructed in close analogy with the pilot-wave model for the many-particle Dirac equation. The difference is that the trajectory model for bosons is constructed from the energy–momentum tensor and not from the multi-particle charge tensor. As indicated in the one-particle case, these boson trajectories should not be regarded as genuine particle trajectories but as representing flowlines of energy. The multi-particle generalization of the HC theory can be constructed analogously, but we will not consider this here.
The single-particle DKP theory can be extended to a $`N`$-particle system as follows. The wavefunction for the $`N`$-particle system is defined to be $`\psi _{r_1\mathrm{}r_N}(𝐱_1,\mathrm{},𝐱_N,t)`$, with the $`r_i`$ denote $`N`$ spin indices. The wavefunction is assumed to be symmetric under permutations of the particle labels. We also introduce the operators $`\beta _{(r)}^\mu `$ which operate only on the $`r^{\text{th}}`$ spin index. The wave equation for the $`N`$-particle wavefunction is defined to be
$`i_0\psi `$ $`=`$ $`{\displaystyle \underset{r=1}{\overset{N}{}}}\left(i\stackrel{~}{\beta }_{(r)}^i_i^{(r)}+\beta _0^{(r)}m_r\right)\psi ,`$ (3.62)
$`i\beta _{(r)}^i\left(\beta _0^{(r)}\right)^2_i^{(r)}\psi `$ $`=`$ $`m_r\left[1\left(\beta _0^{(r)}\right)^2\right]\psi ,`$ (3.63)
where $`_i^{(r)}=\frac{}{(𝐱_r)^i}`$ and $`m_r`$ is the mass of the $`r^{\text{th}}`$ particle.
The equations (3.62) and (3.63) are a straightforward generalization of the Schrödinger form of the one-particle DKP equation. It can easily be verified that a wavefunction $`\psi `$, constructed as an arbitrary superposition of direct products of one-particle DKP wavefunctions at equal time, obeys the many-particle DKP equations (3.62) and (3.63). Conversely, $`\psi `$ can only be written in such a form.
In order to construct the multi-particle generalization of the one-particle energy–momentum tensor (3.15) we first define the following operator:
$$\mathrm{\Gamma }_{(r)}^{\mu \nu }=m_r\eta _0^{(r)}(\beta _{(r)}^\mu \beta _{(r)}^\nu +\beta _{(r)}^\nu \beta _{(r)}^\mu g^{\mu \nu }).$$
(3.64)
The multi-particle energy–momentum tensor of rank $`2N`$ then reads (see also )
$$\mathrm{\Theta }_{DKP}^{\mu _1\mathrm{}\mu _{2N}}=\psi ^{}\mathrm{\Gamma }_{(1)}^{\mu _1\mu _2}\mathrm{}\mathrm{\Gamma }_{(N)}^{\mu _{2N1}\mu _{2N}}\psi $$
(3.65)
and satisfies the conservation equation
$$_0\mathrm{\Theta }_{DKP}^{0_1\nu _1\mathrm{}0_N\nu _N}+\underset{r=1}{\overset{N}{}}_{i_r}^{(r)}\mathrm{\Theta }_{DKP}^{0_1\nu _1\mathrm{}i_r\nu _r\mathrm{}0_N\nu _N}=0.$$
(3.66)
By contracting the energy–momentum tensor $`\mathrm{\Theta }^{\mu _1\mathrm{}\mu _{2N}}`$ with a constant tensor $`n^{\mu _1\mathrm{}\mu _N}`$ of rank $`N`$ we can construct a rank $`N`$ tensor current
$$j^{\mu _1\mathrm{}\mu _N}=\mathrm{\Theta }_{DKP}^{\mu _1\nu _1\mathrm{}\mu _N\nu _N}n_{\nu _1\mathrm{}\nu _N}$$
(3.67)
which satisfies the conservation equation
$$_0j^{0_1\mathrm{}0_N}+\underset{r=1}{\overset{N}{}}_{i_r}^{(r)}j^{0_1\mathrm{}i_r\mathrm{}0_N}=0.$$
(3.68)
This is the multi-particle generalization of the one-particle current $`j^\mu `$ defined in (3.48). In order to be able to define causal trajectories we require that the tensor $`n^{\mu _1\mathrm{}\mu _N}`$ is such that the vectors $`j^{0_1\mathrm{}\mu _r\mathrm{}0_N}`$ are future-causal for every $`r=1,\mathrm{},N`$. We will give examples below. The trajectories $`x_k^\mu (\tau )`$ $`(k=1,\mathrm{},N)`$ can then be defined as solutions to the ‘guidance equations’
$$\frac{dx_r^\mu }{d\tau }=\frac{j^{0_1\mathrm{}\mu _r\mathrm{}0_N}}{j^{0_1\mathrm{}0_N}}.$$
(3.69)
Because the vectors $`j^{0_1\mathrm{}\mu _r\mathrm{}0_N}`$ are assumed to be future-causal, the trajectories will be time-like or null. The density of crossings through constant time hyperplanes is defined to be the positive quantity $`j^{0_1\mathrm{}0_N}`$.
Examples for the tensor $`n^{\mu _1\mathrm{}\mu _N}`$ are generated by considering generalizations of the vector $`n^\mu `$ in the one-particle case. If $`a^\mu `$ is the constant future-causal 4-velocity of a particular observer, then we can take $`n^{\mu _1\mathrm{}\mu _N}=a^{\mu _1}\mathrm{}a^{\mu _N}`$. The proof that $`j^{0_1\mathrm{}\mu _r\mathrm{}0_N}`$ is future-causal for every $`r=1,\mathrm{},N`$ runs as follows. In Section 3.4 it was shown that the operator $`\mathrm{\Gamma }_{0\mu }^{(r)}a^\mu `$ is positive in spin space $`\mathrm{}^M`$, where $`M`$ is the dimension of the representation of the $`\beta `$ matrices. As a result the operator
$$\mathrm{\Gamma }=\mathrm{\Gamma }_{0\nu _1}^{(1)}a^{\nu _1}\mathrm{}\widehat{\mathrm{\Gamma }_{0\nu _r}^{(r)}a^{\nu _r}}\mathrm{}\mathrm{\Gamma }_{0\nu _N}^{(N)}a^{\nu _N},$$
(3.70)
where the hat indicates that the term should be omitted from the product, is a positive operator in $`N1`$ particle spin space $`(\mathrm{}^M)^{(N1)}`$. Consequently there exists an operator $`\mathrm{\Omega }`$ in $`N1`$ particle spin space such that $`\mathrm{\Gamma }=\mathrm{\Omega }^{}\mathrm{\Omega }`$. As a result one can write
$$j^{0_1\mathrm{}\mu _r\mathrm{}0_N}=(\mathrm{\Omega }\psi )^{}\mathrm{\Gamma }_{(r)}^{\mu _r\nu _r}a_{\nu _r}(\mathrm{\Omega }\psi ).$$
(3.71)
This shows that $`j^{0_1\mathrm{}\mu _r\mathrm{}0_N}`$ is the sum of $`M(N1)`$ (sum over all but the $`r^{\text{th}}`$ spin index of $`\mathrm{\Omega }\psi `$) vectors of the form $`\mathrm{\Psi }^{}\mathrm{\Gamma }_{(r)}^{\mu _r\nu _r}a_{\nu _r}\mathrm{\Psi }`$. Because each such term is future-causal the sum $`j^{0_1\mathrm{}\mu _r\mathrm{}0_N}`$ is also future-causal.
The one-particle vector $`P^\mu `$ defined in (3.52) is generalized to the many-particle case by
$$P^{\mu _1\mathrm{}\mu _N}=𝑑x_1^3\mathrm{}𝑑x_N^3\mathrm{\Theta }_{DKP}^{\mu _10_1\mathrm{}\mu _N0_N}.$$
(3.72)
The proof that $`P^{\mu _1\mathrm{}\mu _N}`$ leads to future-causal vectors $`j^{0_1\mathrm{}\mu _r\mathrm{}0_N}`$ proceeds essentially in the same way as in the first example, however we will not present it here.
One can easily show that for both definitions of $`n^{\mu _1\mathrm{}\mu _N}`$ the trajectory laws for a product state
$$\psi _{r_1\mathrm{}r_N}(𝐱_1,\mathrm{},𝐱_N,t)=\psi _{1,r_1}(𝐱_1,t)\mathrm{}\psi _{N,r_N}(𝐱_N,t)$$
(3.73)
reduce to the one-particle trajectory laws. I.e. for a product state, the density of crossings through constant time hyperplanes is given by the product of one-particle densities, i.e. $`j^{0_1\mathrm{}0_N}=j_1^{0_1}\mathrm{}j_N^{0_1}`$ with $`j_\alpha ^{0_\alpha }=\psi _\alpha ^{}\mathrm{\Gamma }^{0\mu }n_\mu \psi _\alpha `$, with $`n^\mu `$ the vector defined in the one-particle case. Second, the velocity field for each particle reduces to the one-particle velocity field defined in (3.49).
In the special case that $`n^{\mu _1\mathrm{}\mu _N}=\delta _0^{\mu _1}\mathrm{}\delta _0^{\mu _N}`$ the density of crossings becomes $`j^{0_1\mathrm{}0_N}=(\mathrm{\Pi }_rm_r)\psi ^{}\psi `$ and the guidance equations become
$$\frac{dx_r^i}{dt}=\frac{\psi ^{}\stackrel{~}{\beta }_{(r)}^i\psi }{\psi ^{}\psi }.$$
(3.74)
The resulting conservation equation (3.68), which can also be directly derived from (3.62), turns into
$$_0(\psi ^{}\psi )+\underset{r=1}{\overset{N}{}}_i^{(r)}(\psi ^{}\stackrel{~}{\beta }_{(r)}^i\psi )=0.$$
(3.75)
In this case the trajectory model displays a formal resemblance to the pilot-wave equations for the multi-particle Dirac equation as presented in Section 3.2.2.
We also consider the non-relativistic limit of the trajectory model. The non-relativistic limit of the multi-particle DKP equations reduces to the non-relativistic Schrödinger equation for $`N`$ spinless particles in the spin-0 representation and to the non-relativistic Schrödinger equation for $`N`$ spin-1 particles in the spin-1 representation. As in the one-particle case, the trajectory model constructed through the energy–momentum tensor reduces to the many-particle pilot-wave interpretations for spin-0 and spin-1 in the respective representations. In particular, in the spin-0 representation, the ‘guidance equations’ (3.69) reduce to the guidance equations originally presented by de Broglie and Bohm (cf. Section 2.2)
$$\frac{d𝐱_r}{dt}=\frac{\text{Im}\left(\psi ^{}\mathbf{}_r\psi ^{}\right)}{m|\psi ^{}|^2},$$
(3.76)
with $`\psi ^{}(𝐱_1,\mathrm{},𝐱_N)`$ the non-relativistic $`N`$-particle wavefunction. In the spin-1 representation the guidance equations become
$$\frac{d𝐱_r}{dt}=\frac{\text{Im}(\mathrm{\Phi }^{}\mathbf{}_r\mathrm{\Phi })}{m\mathrm{\Phi }^{}\mathrm{\Phi }}+\frac{\mathbf{}_r\times (\mathrm{\Phi }^{}𝐒_r^{(1)}\mathrm{\Phi })}{2m\mathrm{\Phi }^{}\mathrm{\Phi }},$$
(3.77)
with $`𝐒_r^{(1)}`$ operating only on the $`r^{\text{th}}`$ spin index of the non-relativistic multi-particle wavefunction $`\mathrm{\Phi }_{r_1,\mathrm{},r_N}(𝐱_1,\mathrm{},𝐱_N,t)`$, where each spin index $`r_i`$, $`i=1,\mathrm{},N`$, runs from 1 to 3. Note the spin contribution, which is similar to the spin contribution in the one-particle spin-1 guidance equation (2.22).
The multi-particle DKP charge tensor current is defined as
$$s_{DKP}^{\mu _1\mathrm{}\mu _N}=\psi ^{}\eta _{(1)}^0\beta _{(1)}^{\mu _1}\mathrm{}\eta _{(N)}^0\beta _{(N)}^{\mu _N}\psi $$
(3.78)
and satisfies the conservation equation
$$_0s_{DKP}^{0_1\mathrm{}0_N}+\underset{r=1}{\overset{N}{}}_i^{(r)}s_{DKP}^{0_1\mathrm{}i_r\mathrm{}0_N}=0.$$
(3.79)
Just as in the one-particle case, the vectors $`s_{DKP}^{0_1\mathrm{}\mu _r\mathrm{}0_N}`$ have the same non-relativistic limit as the vectors $`j^{0_1\mathrm{}\mu _r\mathrm{}0_N}`$.
### 3.8 The minimally coupled Duffin-Kemmer-Petiau theory
The minimally coupled DKP Lagrangian is obtained from the free DKP Lagrangian by applying the minimal coupling prescription $`_\mu D_\mu =_\mu +ieV_\mu `$, with $`V^\mu =(V_0,𝐕)`$, and yields
$$_{DKP}=\frac{i}{2}(\overline{\psi }\beta _\mu D^\mu \psi D_\mu ^{}\overline{\psi }\beta ^\mu \psi )m\overline{\psi }\psi .$$
(3.80)
The corresponding coupled wave equation reads
$$(i\beta ^\mu D_\mu m)\psi =0.$$
(3.81)
The symmetrized energy–momentum tensor is the same as in the free case
$$\mathrm{\Theta }_{DKP}^{\mu \nu }=m\psi ^{}\eta _0(\beta ^\mu \beta ^\nu +\beta ^\nu \beta ^\mu g^{\mu \nu })\psi .$$
(3.82)
However, the energy–momentum tensor is no longer conserved, but instead satisfies the equation
$$_\mu \mathrm{\Theta }_{DKP}^{\mu \nu }=F_\mu ^\nu s_{DKP}^\mu ,$$
(3.83)
with $`s_{DKP}^\mu `$ the charge current, which is formally the same as in the free case, and $`F^{\mu \nu }=^\mu V^\nu ^\nu V^\mu `$ the electromagnetic tensor. The right hand side of the equation (3.83) is recognized as the Lorentz force. Hence, if we would now construct the current $`j^\mu =\mathrm{\Theta }_{DKP}^{\mu \nu }n_\nu `$, as in the free case, then this current is no longer conserved and a particle interpretation becomes untenable.
It is interesting to note that Ghose et al. considered the energy–momentum tensor $`\mathrm{\Theta }_{DKP}^{\mu \nu }`$ to be conserved when an interaction with an electromagnetic field $`V^\mu `$ is introduced via minimal coupling. Therefore they thought that even in the case of an interaction a particle interpretation was possible. The reason for the discrepancy is that they introduced minimal coupling in the DKP theory in an other way than presented above.
Let us consider this in more detail. As explained above one can introduce minimal coupling at the level of the covariant form of the DKP equation (3.81). The covariant equation can also be written in the following Schrödinger form
$`iD_0\psi `$ $`=`$ $`(i\stackrel{~}{\beta }^iD_i+m\beta _0)\psi {\displaystyle \frac{ie}{2m}}F^{\mu \nu }(\beta _\nu \beta _0\beta _\mu +\beta _\nu g_{\mu 0})\psi ,`$ (3.84)
$`i\beta ^i\beta _0^2D_i\psi `$ $`=`$ $`m(1\beta _0^2)\psi .`$ (3.85)
On the other hand, as proposed by Ghose et al., one can introduce minimal coupling at the level of the Schrödinger form of the DKP theory (cf. (3.17) and (3.18)) which results in
$`iD_0\psi `$ $`=`$ $`(i\stackrel{~}{\beta }^iD_i+m\beta _0)\psi ,`$ (3.86)
$`i\beta ^i\beta _0^2D_i\psi `$ $`=`$ $`m(1\beta _0^2)\psi .`$ (3.87)
The difference is clear, the term involving $`F^{\mu \nu }`$ in the Schrödinger form (3.84) is not present in (3.86). This additional term has no equivalent in the spin-1/2 Dirac theory and is hard to interpret . It has recently been argued that this additional term is irrelevant . The argument is that when the DKP theory is reduced to its physical components, then the DKP theory reduces to the minimally coupled Klein-Gordon theory in the spin-0 case and to the minimally coupled Proca theory in the spin-1 case, where the anomalous term containing $`F^{\mu \nu }`$ disappears. However, this does not settle the question whether or not we can safely introduce minimal coupling at the level of the Schrödinger form in the DKP theory. We argue that in general the only correct way to introduce minimal coupling is at the level of the covariant form of the DKP equations.
The minimally coupled Schrödinger form implies the ‘covariant form’ (3.81), which can be seen by multiplying (3.86) with $`\beta _0`$ and by adding (3.87) to it . Hence, if the minimally coupled Schrödinger form is regarded as fundamental, then both the equations (3.84) and (3.86) should be valid. This implies that the additional term containing the tensor $`F^{\mu \nu }`$ should be zero. Without considering an explicit representation we can already show that this implies that the introduction of minimal coupling at the level of the Schrödinger form is in general untenable.
If minimal coupling is introduced at the level of the Schrödinger form of the DKP equation, we may derive from (3.86) (by multiplying (3.86) by $`\psi ^{}`$ from the right, multiplying the conjugate of (3.86) by $`\psi `$ from the left, and subtracting the two from each other ) that
$$_0(\psi ^{}\psi )+_i(\psi ^{}\stackrel{~}{\beta }^i\psi )=0$$
(3.88)
or differently written
$$_\mu \mathrm{\Theta }_{DKP}^{\mu 0}=0.$$
(3.89)
If we assume that the the Schrödinger form of the DKP equation is covariant,<sup>15</sup><sup>15</sup>15Although a pilot-wave model may be non-covariant at the subquantum level, it is desired that we have covariance at the quantum level, therefore the wave equation should be covariant. i.e. if we assume that the wave equation has the form (3.86) in every Lorentz frame, then we have $`_\mu \mathrm{\Theta }_{DKP}^{\mu 0}=0`$ in every Lorentz frame. Because $`_\mu \mathrm{\Theta }_{DKP}^{\mu \nu }`$ transforms as a 4-vector under Lorentz transformations, this implies that the time component of $`_\mu \mathrm{\Theta }_{DKP}^{\mu \nu }`$ is zero in every Lorentz frame. Because the only 4-vector which satisfies this property is the zero vector, we have
$$_\mu \mathrm{\Theta }_{DKP}^{\mu \nu }=0.$$
(3.90)
It was this equation which led Ghose et al. to conclude that a particle interpretation associated with the energy–momentum tensor was still possible in the presence of an external field. However, because the minimally coupled Schrödinger form implies the minimally coupled covariant form, both (3.83) and (3.90) should be valid when the minimally coupled Schrödinger form is regarded as fundamental. This implies however that the Lorentz force is zero, i.e.
$$F_\mu ^\nu s_{DKP}^\mu =0,$$
(3.91)
which is in general not the case. From this we may conclude that minimal coupling should be introduced at the level of the covariant form of the DKP equation.
The discrepancy between the introduction of minimal coupling at the level of the covariant form of the DKP equation and the introduction of minimal coupling at the level of the Schrödinger form of the DKP equation does not disappear when the theory is reduced to its physical components. Consider for example the case of spin-0. In the same way as was explained in Section 3.3.1, using the explicit representation in Appendix A, we can reduce the theory to its physical components. By using the explicit representation, it follows from both equations (3.84) and (3.86) that the DKP wavefunction $`\psi `$ can be written in terms of a remaining physical component $`\varphi `$:
$$\psi =\frac{1}{\sqrt{m}}\left(\begin{array}{c}D_\mu \varphi \\ m\varphi \end{array}\right)$$
(3.92)
where $`\varphi `$ satisfies the minimally coupled Klein-Gordon equation
$$(D_\mu D^\mu +m^2)\varphi =0.$$
(3.93)
Hence, both ways of introducing minimal coupling then lead to the minimally coupled Klein-Gordon theory . When the equation (3.83) is written in terms of the physical component $`\varphi `$ by making the substitution (3.92), then we obtain
$$_\mu \mathrm{\Theta }_{KG}^{\mu \nu }=F_\mu ^\nu s_{KG}^\mu ,$$
(3.94)
with
$$\mathrm{\Theta }_{KG}^{\mu \nu }=D^\mu \varphi (D^\nu \varphi )^{}+(D^\mu \varphi )^{}D^\nu \varphi g^{\mu \nu }\left(D_\alpha \varphi (D^\alpha \varphi )^{}m^2\varphi ^{}\varphi \right)$$
(3.95)
the Klein-Gordon energy–momentum tensor and
$$s_{KG}^\mu =ie\left(\varphi ^{}D^\mu \varphi (D^\mu \varphi )^{}\varphi \right)$$
(3.96)
the charge current in the minimally coupled Klein-Gordon theory. On the other hand, when (3.90) is written in terms of the physical component $`\varphi `$ we obtain
$$_\mu \mathrm{\Theta }_{KG}^{\mu \nu }=0.$$
(3.97)
Because only (3.94), and not (3.97), can be derived from the minimally coupled Klein-Gordon equation (3.93), we should introduce minimal coupling at the level of the covariant form of the DKP equation and not at the level of the Schrödinger form of the DKP equation.
Note that it is a general property for charged matter, that the matter energy–momentum tensor (even for fermions) is not conserved when an external electromagnetic field is introduced. This is because a charged particle exchanges energy and momentum when interacting with an electromagnetic field. On the other hand, the charge current is always conserved for bosons and for fermions, because the electromagnetic field does not carry charge. This is the reason why it presented no problem to maintain the particle interpretation in the Dirac theory, which was associated with the charge current and not with the energy–momentum tensor, after an electromagnetic field was introduced.
From our analysis it also follows that the physical interpretation of the term containing $`F^{\mu \nu }`$ in the Schrödinger form (3.84) may perhaps be sought in the fact that it contributes to the Lorentz force.
## Chapter 4 Field beables for bosonic quantum field theory
### 4.1 Introduction
The construction of a pilot-wave for bosonic quantum field theories has been discussed before in a number of papers. Already in his seminal paper in 1952, Bohm presented a pilot-wave interpretation for the free electromagnetic field. Later, in the 80’s and 90’s this pilot-wave interpretation was further elaborated on by Kaloyerou . In 1992, Valentini presented a different approach for the electromagnetic field. In the meanwhile the free massless scalar field was treated by Bohm and Hiley (later reviews of their treatment can be found in ). The extension to a free massive scalar field was given independently, and along the same lines, by Kaloyerou and Valentini . In , Valentini further considered the pilot-wave interpretation for the massive scalar field coupled to the electromagnetic field.
All these authors introduced fields as beables in the pilot-wave interpretation. As explained in the previous chapter, we also favour the field beable approach for quantum field theory. So it is along these lines that we try to develop the pilot-wave interpretation further.
We start with considering the pilot-wave interpretation for the quantized massive spin-0 and spin-1 field coupled to a non-quantized external electromagnetic field. Instead of simply presenting the pilot-wave interpretation for the quantized Klein-Gordon theory or the quantized Proca theory, we will take a different approach. We will start with the Duffin-Kemmer-Petiau (DKP) theory. First we will run through the quantization procedure, i.e. we will first apply the rules of canonical quantization, as set out by Dirac, in order to quantize the DKP theory. Then only afterwards we will present the corresponding pilot-wave theory. The reason to do so is twofold.
First, it was reported in the literature that the equivalence of the quantized DKP theory and the quantized Klein-Gordon theory in the spin-0 representation or the quantized Proca theory in the spin-1 representation is not that obvious, not even when only electromagnetic interaction is considered (although the equivalence as wave equations is well established). However, we show that, by using Dirac’s recipe of canonical quantization, it is straightforward to show the equivalence (in the case of spin-0 the equivalence was also shown by Fainberg and Pimentel by using similar arguments ). Once the DKP theory is quantized it is no problem to provide a pilot-wave interpretation in terms of field beables. But because of the equivalence, we could equally well have started from the quantized Klein-Gordon theory or quantized Proca theory from the start. In the case of a free spin-0 field the pilot-wave interpretation then of course reduces to the one originally presented by Kaloyerou and Valentini.
The second reason to go through the canonical quantization procedure for the DKP theory, is that it is a good example of how the quantization program of Dirac works. When we discuss the pilot-wave interpretation for the electromagnetic field, we will need to appeal to the canonical quantization program again. Only, in this case, slight complications arise due to the fact that we are dealing with a gauge theory. The canonical quantization of the electromagnetic field is of course well know, but it is instructive to recall it. This mean reason to do this is that the two existing approaches to a pilot-wave interpretation for the electromagnetic field, namely the one by Bohm and Kaloyerou, and the one by Valentini, find a natural home in two different ways of quantizing theories with gauge symmetries. Against the background of canonical quantization it is then easy to compare the two approaches.
A careful look at Valentini’s model reveals that it suffers from a problem. Namely the densities of field beables are non-normalizable. This problem is not present in the model by Bohm and Kaloyerou. The reason is that Valentini introduces beables corresponding to gauge degrees of freedom, whereas Bohm and Kaloyerou only introduce beables for gauge independent degrees of freedom.
After presenting the pilot-wave interpretation for the free electromagnetic field we then consider the pilot-wave interpretation for the quantized Klein-Gordon field coupled to quantized electromagnetic field (scalar quantum electrodynamics). Scalar quantum electrodynamics was first treated by Valentini , but in the model that we present here, beables are introduced only for gauge invariant degrees of freedom and hence our model does not suffer from the problem of non-normalizable field beable densities.
This is the organization of the chapter. We start with a review of Dirac’s procedure of canonical quantization in Section 4.2. In Sections 4.3 and 4.4, we consider the quantization of the quantized DKP field coupled to a non-quantized electromagnetic field in respectively the spin-0 and the spin-1 representation, together with the corresponding pilot-wave interpretation. In Section 4.5, we treat the quantization of the electromagnetic field and we discuss in detail the model of Bohm and Kaloyerou and the model of Valentini. Then, in Section 4.6, we consider the pilot-wave interpretation for scalar quantum electrodynamics. In Section 4.7 we discuss the possibility of constructing a pilot-wave interpretation for non-Abelian gauge theories. We end the chapter with a discussion on how the pilot-wave interpretation in terms of field beables solves the measurement problem.
The main conclusion of this chapter is that it presents no problem to construct a pilot-wave interpretation in terms of field beables for bosonic quantum field theory. Even for gauge theories we can develop a pilot-wave interpretation. This stands in sharp contrast with the problems we encountered when we tried to develop a pilot-wave interpretation for relativistic wave equations for bosons. However, in the following chapter, where we consider fermionic field theory, it will appear much more difficult to construct a field beable model for fermionic fields.
### 4.2 Canonical quantization of constrained systems
In this section we review Dirac’s procedure of canonical quantization of a constrained system. This review is mainly based on Dirac’s original presentation and the book by Henneaux and Teitelboim (which closely follows Dirac’s original presentation). We also consulted the book by Sundermeyer and the one by Gitman and Tyutin . A short introduction on canonical quantization can be found in Weinberg’s book \[112, pp. 325-330\]. Although this section is self-contained, we only review the essential ingredients of canonical quantization. We refer the interested reader to the aforementioned books for a more detailed treatment. The reader which is familiar with Dirac’s procedure of canonical quantization can skip this review.
#### 4.2.1 Hamiltonian formulation of a constrained system
For simplicity we present Dirac’s analysis for a system with a finite number of degrees of freedom. In Section 4.2.4 we indicate how the transition to a continuous number of degrees of freedom can be made.
We assume that the dynamics can be derived from the action
$$S=𝑑tL(q,\dot{q}).$$
(4.1)
Here $`L(q,\dot{q})`$ is a Lagrangian which is function of the coordinates $`q_n`$, where $`n=1,\mathrm{},N`$, and the corresponding velocities $`\dot{q}=dq/dt`$. By requiring that the action $`S`$ be stationary with respect to variations in the coordinates $`q_n`$ we obtain the Euler-Lagrange equations of motion
$$\frac{d}{dt}\frac{L}{\dot{q}_n}\frac{L}{q_n}=0.$$
(4.2)
In order to quantize the system, we have to move from the Lagrangian formulation, where the dynamical variables are the velocity phase-space variables $`q_n`$ and $`\dot{q}_n`$, to the Hamiltonian formulation, where the dynamical variables are the momentum phase-space variables $`q_n`$ and $`p_n`$. Canonical quantization then proceeds by associating operators with the momentum phase-space variables and by imposing certain commutation relations for these operators. Dirac’s procedure of canonical quantization provides a scheme for imposing these commutation relations in a way consistent with the various constraints that may arise in the Hamiltonian formulation.
In order to arrive at the Hamiltonian formulation we first need the momenta $`p_n`$ canonically conjugate to the coordinates $`q_n`$. These are defined as
$$p_n=\frac{L}{\dot{q}_n}(q,\dot{q}).$$
(4.3)
If the rank of the Hessian matrix
$$\frac{\delta ^2L}{\dot{q}_n^{}\dot{q}_n}(q,\dot{q})$$
(4.4)
is maximal at each point in velocity phase-space, then we can write all the velocities $`\dot{q}_n`$ as functions of the momenta and the coordinates. In this case the mapping from velocity phase-space variables to momentum phase-space variables is invertible. If the rank is not maximal, then this mapping is not invertible. In this case the momenta (4.4) are not all independent, but there are, rather, some relations
$$\chi _m(q,p)=0,m=1,\mathrm{},MN.$$
(4.5)
These relations are called the primary constraints.
For simplicity we have hereby assumed that the rank of the Hessian matrix is constant throughout velocity phase-space so that the constraints can be written in the particular form (4.5). For future convenience we also assume that the rank of $`\chi _m/(q_n,p_n^{})`$ is $`M`$ throughout velocity phase-space. This condition ensures that all primary constraints are lineary independent. It also excludes the use of equivalent sets of constraints, such as e.g. $`\chi _m^2(q,p)=0`$, $`m=1,\mathrm{},M`$.
We proceed by defining the canonical Hamiltonian
$$H_C=p_n\dot{q}_nL,$$
(4.6)
which is a function of the $`q_n`$ and the $`\dot{q}_n`$. By using the definition of the momenta and by using the Euler-Lagrange equations of motion, we find that the canonical Hamiltonian varies as
$$\delta H_C=\underset{n=1}{\overset{N}{}}(\delta p_n\dot{q}_n\delta q_n\dot{p}_n)$$
(4.7)
under infinitesimal variations of the coordinates and the velocities. Because $`\delta H_C`$ does not contain variations in the velocities $`\dot{q}_n`$, we can express $`H_C`$ solely in terms of the variables $`q_n`$ and $`p_n`$, i.e. $`H_C=H_C(q,p)`$.
The momentum phase-space variables cannot be varied independently because of the primary constraints. As a result the equations of motion which are derived from (4.7) contain arbitrary functions $`u_m`$ of the momentum phase-space variables:
$`\dot{q}_n`$ $`=`$ $`{\displaystyle \frac{H_C}{p_n}}+{\displaystyle \underset{m=1}{\overset{M}{}}}u_m{\displaystyle \frac{\chi _m}{p_n}},`$ (4.8)
$`\dot{p}_n`$ $`=`$ $`{\displaystyle \frac{H_C}{q_n}}{\displaystyle \underset{m=1}{\overset{M}{}}}u_m{\displaystyle \frac{\chi _m}{q_n}}.`$ (4.9)
At this point it is useful to introduce the Poisson bracket $`[F,G]_P`$ for two momentum phase-space functions $`F(q,p)`$ and $`G(q,p)`$:
$$[F,G]_P=\underset{n=1}{\overset{N}{}}\frac{F}{q_n}\frac{G}{p_n}\frac{F}{p_n}\frac{G}{q_n}.$$
(4.10)
Using the Poisson bracket, we can write the equation of motion for any function $`F(q,p)`$ as
$$\dot{F}=[F,H_C]_P+\underset{m=1}{\overset{M}{}}u_m[F,\chi _m]_P.$$
(4.11)
The equations of motion can be written in an even more concise form. In order to do so let us first introduce Dirac’s equality sign ‘$``$’ which is defined by
$$F(q,p)G(q,p)F(q,p)|_{\chi _m=0;m=1,\mathrm{},M}=G(q,p)|_{\chi _m=0;m=1,\mathrm{},M}.$$
(4.12)
If $`FG`$, one says that $`F`$ weakly equals $`G`$. By further introducing the total Hamiltonian
$$H_T=H_C+\underset{m=1}{\overset{M}{}}u_m\chi _m,$$
(4.13)
the equations of motion (4.11) can be written as
$$\dot{F}[F,H_T]_P.$$
(4.14)
Now we can further examine the consequences of these equations of motion. In the first place there will be some consistency conditions. The constraints $`\chi _m0`$ have to be weakly conserved in time. Hence we have the requirement
$$\dot{\chi }_m^{}[\chi _m^{},H_T]_P[\chi _m^{},H_C]_P+\underset{m=1}{\overset{M}{}}u_m[\chi _m^{},\chi _m]_P0$$
(4.15)
for $`m^{}=1,\mathrm{},M`$. Some of these equations may be equations which do not contain any $`u_m`$ and hence they may lead to new constraints that have to be satisfied. These constraints are called secondary constraints. In turn, the secondary constraints should also be conserved in time and hence may lead to further constraints. We can repeat this procedure until no further constraints are found. In each step of evaluating the consistency condition $`\dot{\chi }_m^{}0`$, the equality sign $``$ refers to the full set of constraints obtained at that stage. The newly obtained constraints are denoted by
$$\chi _k0,k=M+1,\mathrm{},M+K.$$
(4.16)
So that the full set of constraints can then be written as
$$\chi _j0,j=1,\mathrm{},J=M+K.$$
(4.17)
Not only may the consistency requirements that the constraints are conserved in time lead to new constraints, they may also determine some of the coefficients $`u_m`$. To see this consider again the consistency conditions
$$[\chi _j,H_T]_P[\chi _j,H_C]_P+\underset{m=1}{\overset{M}{}}u_m[\chi _j,\chi _m]_P0$$
(4.18)
where $`j`$ ranges from $`j=1,\mathrm{},J`$, and let us now regard them as equations for the $`u_m`$. Let $`u_m=U_m`$ be a particular solution for (4.18). The most general solution is then
$$u_m=U_m+\underset{a=1}{\overset{A}{}}v_aV_{am},$$
(4.19)
where the $`V_{am}`$, with $`a=1,\mathrm{},A`$, are $`A`$ independent solutions of the homogeneous equation
$$\underset{m=1}{\overset{M}{}}V_{am}[\chi _j,\chi _m]_P0$$
(4.20)
and the $`v_a`$ are arbitrary coefficients.
We may substitute these expression for the $`u_m`$ in the total Hamiltonian to obtain
$$H_T=H^{}+\underset{a=1}{\overset{A}{}}v_a\chi _a$$
(4.21)
where
$`H^{}`$ $`=`$ $`H_C+{\displaystyle \underset{m=1}{\overset{M}{}}}U_m\chi _m,`$ (4.22)
$`\chi _a`$ $`=`$ $`{\displaystyle \underset{m=1}{\overset{M}{}}}V_{am}\chi _m.`$ (4.23)
We see that the equation of motion for an arbitrary function $`F(q,p)`$ now reads
$$\dot{F}[F,H_T]_P[F,H^{}]_P+\underset{a=1}{\overset{A}{}}v_a[F,\chi _a]_P$$
(4.24)
Notice that the equation of motion for $`F`$ may depend on the arbitrary coefficients $`v_a`$. This remaining arbitrariness denotes the presence of some gauge invariance. In this context it is said that the $`\chi _a`$, $`a=1,\mathrm{},A`$, generate infinitesimal gauge transformations. This can be seen as follows. Let $`F_v`$ and $`F_v^{}`$ be phase-space functions evolving form the initial value $`F_0`$ with two different sets of coefficients, $`v_a`$ and $`v_a^{}`$. To first order we have
$$F_v(\delta t)=F_0+[F_0,H_T]_P\delta tF_0+[F_0,H^{}]_P\delta t+\underset{a=1}{\overset{A}{}}v_a[F_0,\chi _a]_P\delta t$$
(4.25)
and hence
$$\delta F(\delta t)=F_v(\delta t)F_v^{}(\delta t)\underset{a=1}{\overset{A}{}}(v_av_a^{})[F_0,\chi _a]_P\delta t.$$
(4.26)
So the functions $`F_v(\delta t)`$ and $`F_v^{}(\delta t)`$ are related by an infinitesimal canonical transformation generated by $`_{a=1}^A(v_av_a^{})\chi _a\delta t`$.<sup>1</sup><sup>1</sup>1A transformation $`Q(q,p,t)`$, $`P(q,p,t)`$ is canonical if
$$[Q_n,Q_n^{}]_P=[P_n,P_n^{}]_P=0,[Q_n,P_n^{}]_P=\delta _{nn^{}},$$
(4.27) where the Poisson bracket is calculated with respect to the variables $`q`$ and $`p`$. This means that in the case of a canonical transformation, the Poisson brackets of phase-space functionals are the same, whether calculated with old or new phase-space variables. The new Hamiltonian $`H^{}(Q,P)`$ corresponding to the transformed system is
$$H^{}=\underset{n=1}{\overset{N}{}}P_n\dot{Q}_nL+\frac{dF}{dt},$$
(4.28) where the Lagrangian $`L`$ is expressed in terms of the primed variables. The function $`F(Q,P)`$ is the generating function, which is determined by $`{\displaystyle \frac{\delta F}{\delta Q_n}}`$ $`=`$ $`{\displaystyle \underset{n^{}=1}{\overset{N}{}}}p_n^{}{\displaystyle \frac{\delta q_n^{}}{\delta Q_n}}P_n,`$ $`{\displaystyle \frac{\delta F}{\delta P_n}}`$ $`=`$ $`{\displaystyle \underset{n^{}=1}{\overset{N}{}}}p_n^{}{\displaystyle \frac{\delta q_n^{}}{\delta P_n}}.`$ (4.29) The constraints $`\chi _a`$, with $`a=1,\mathrm{},A`$, are not the only generators of infinitesimal gauge transformations. One can for example show that $`[\chi _a,\chi _a^{}]_P`$ and $`[H^{},\chi _a^{}]_P`$, with $`a,a^{}=1,\mathrm{},A`$, also generate gauge transformations. In practice this means that some of the non-primary constraints will also generate infinitesimal gauge transformations (one can show that $`[\chi _a,\chi _a^{}]_P`$ and $`[H^{},\chi _a]_P`$ weakly vanish and hence they are linear combinations of the constraints $`\chi _j`$, $`j=1,\mathrm{},J`$).
Suppose we have some additional lineary independent constraints $`\chi _a`$, with $`a=A+1,\mathrm{},A^{}`$, which generate infinitesimal gauge transformations. One can then make the gauge invariance of the dynamics explicit by using the extended Hamiltonian
$$H_E=H_T+\underset{a=A+1}{\overset{A^{}}{}}v_a\chi _a.$$
(4.30)
The corresponding equation of motion for a function $`F`$ reads
$$\dot{F}[F,H_E]_P[F,H^{}]_P+\underset{a=1}{\overset{A^{}}{}}v_a[F,\chi _a]_P.$$
(4.31)
It is clear that a function $`F`$ for which the equation of motion depends on the arbitrary functions $`v_a`$, $`a=1,\mathrm{},A^{}`$ can not be an observable quantity. Hence an observable function should have Poisson brackets zero with all the generators of gauge transformations $`\chi _a`$, $`a=1,\mathrm{},A^{}`$; in other words an observable function is gauge invariant.
The constraints $`\chi _b`$ which will be added to the Hamiltonian are the first class constraints. A function $`F`$ of the momentum phase-space variables is called first class if it has zero Poisson brackets with all the constraints, i.e.
$$[F,\chi _j]_P0,j=1,\mathrm{},J.$$
(4.32)
Using their definition (4.23), one can check that the constraints $`\chi _a`$, $`a=1,\mathrm{},A`$, which were defined in (4.23), are first class, as well as $`H^{}`$. One can also show that the Poisson bracket of two first class functions is also first class. Hence $`[\chi _a,\chi _a^{}]_P`$ and $`[H^{},\chi _a]_P`$ are also first class. So it seems that the set of first class constraints corresponds to the set of generators of infinitesimal gauge transformations. This was indeed conjectured by Dirac. However, later, counter examples of this conjecture were presented (although these had no physical relevance). Although the set of first class constraints does hence not correspond to the set of infinitesimal gauge transformations, we will have to treat the first class as such when we try to quantize the system. I.e. the first class constraints should be added to the total Hamiltonian to yield the extended Hamiltonian.
The constraints which are not first class are called second class constraints. The distinction between first class constraints and second class constraints is important if we want to quantize the system. We now discuss two distinct cases separately. In the first case we assume a system which only has second class constraints and in the second case we assume a system which has only first class constraints.
#### 4.2.2 Canonical quantization of systems with second class constraints
Assume that all the constraints are second class constraints. This implies that the matrix
$$C_{jj^{}}=[\chi _j,\chi _j^{}]_P,$$
(4.33)
with $`j,j^{}=1,\mathrm{},J`$, is non-singular (for every point in momentum phase-space). Otherwise there would exist a linear combination of the constraints $`\chi _j`$ which has Poisson brackets zero with every constraint and hence this linear combination would constitute a first class constraint. Because $`C`$ is anti-symmetric the number of second class constraints must necessarily be even, otherwise the determinant of $`C`$ would be zero. The inverse of $`C`$ is denoted by $`C^1`$.
If we would now quantize the system, by associating the operators $`\widehat{F}`$ and $`\widehat{G}`$ with some momentum phase-space functions $`F(q,p)`$ and $`G(q,p)`$, and by imposing the commutation relations<sup>2</sup><sup>2</sup>2One cannot associate quantum operators with all momentum phase-space functions, because some phase-space functions may correspond to more than one quantum operator. Therefore applying the prescription (4.34) to any momentum phase-space function would lead to contradictions. This is the operator ordering ambiguity. Different choices for operator orderings may correspond to different quantum theories.
$$[\widehat{F},\widehat{G}]=i\widehat{[F,G]_P},$$
(4.34)
then these commutation relations would not always be consistent with the constraints.<sup>3</sup><sup>3</sup>3For fermionic theories anti-commutation relations can be imposed. I.e. the operations of imposing the constraints and taking the commutator of the operators would not always commute. For this reason Dirac introduced the Dirac bracket
$$[F,G]_D=[F,G]_P[F,\chi _j]_PC_{jj^{}}^1[\chi _j^{},G]_P,$$
(4.35)
which is defined for any two phase-space functions $`F(q,p)`$ and $`G(q,p)`$. The basic properties of the Dirac bracket are the same as for the Poisson bracket: linearity, antisymmetry, Leibnitz rule and Jacobi identity. In addition, if $`F`$ or $`G`$ is a linear combination of constraints then $`[F,G]_D=0`$. Hence the operations of imposing the constraints and taking the Dirac bracket commute. The system can now be quantized by imposing the following equal-time commutation relations for the operators $`\widehat{F}`$ and $`\widehat{G}`$
$$[\widehat{F},\widehat{G}]=i\widehat{[F,G]_D}.$$
(4.36)
One can easily verify that equations of motion for any phase-space function $`F`$ can be written as
$$\dot{F}[F,H_C]_D[F,H]_D,$$
(4.37)
with
$$H=H_C|_{\chi _j=0;j=1,\mathrm{},J}.$$
(4.38)
Hence, once the Dirac bracket is found, the total Hamiltonian is of no further use. By using the Dirac bracket we can calculate the equations of motion using the Hamiltonian $`H`$. So the total Hamiltonian is only needed to find the secondary and further $`n`$-ary constraints.
##### The true degrees of freedom
By using the Dirac bracket we can quantize a system with second class constraints. However, because of the constraints the phase-space variables are not all independent; we are in fact dealing with too many phase-space variables. Nevertheless it is in principle possible to reduce the number of phase-space variables by isolating the true degrees of freedom. This is due to a theorem by Maskawa and Nakajima \[113, 112, pp. 329-330\].<sup>4</sup><sup>4</sup>4The same result was also presented in \[110, pp. 82-85\] and in \[111, p. 30\]. It will appear crucial for the construction of a pilot-wave interpretation, at least for the field theories we consider, to separate out the true degrees of freedom.
The Maskawa-Nakajima theorem states that if there are $`J=2R`$ second class constraints (and no first class constraints), then we can, at least locally, perform a canonical transformation such that the new canonical variables can be written in terms of two sets $`Q_l`$ and $`\overline{Q}_k`$ and their respective conjugate momenta $`P_l`$ and $`\overline{P}_k`$, with $`l=1,\mathrm{},NR`$ and $`k=NR+1,\mathrm{},N`$, such that the constraints in terms of the new variables read $`\overline{Q}_k=\overline{P}_k=0`$ for $`k=NR+1,\mathrm{},N`$. The theorem further states that
$$[F,G]_D|_{\chi _j=0}=\underset{l=1}{\overset{nr}{}}\frac{F^{}}{Q_l}\frac{G^{}}{P_l}\frac{G^{}}{Q_l}\frac{F^{}}{P_l},$$
(4.39)
with $`F^{}(Q,P)=F(q(Q,P,\overline{Q},\overline{P}),q(Q,P,\overline{Q},\overline{P}))|_{\overline{Q}_k=\overline{P}_k=0}`$ and a similar definition for $`G^{}`$. This means that the Dirac bracket equals the Poisson bracket ‘restricted to the unconstrained variables’, when the constraints are imposed. The Hamiltonian $`H`$ defined (4.38) reduces to the following Hamiltonian for the variables $`Q_l`$ and $`P_l`$, with $`l=1,\mathrm{},nr`$:
$$H(P,Q)=\underset{l=1}{\overset{nr}{}}P_l\dot{Q}_lL|_{\overline{Q}_k=\overline{P}_k=0}+\frac{dF}{dt}|_{\overline{Q}_k=\overline{P}_k=0}.$$
(4.40)
$`F`$ is the generation function of the canonical transformation (cf. the footnote on p. 1). Gitman and Tyutin call the Hamiltonian $`H(P,Q)`$ the physical Hamiltonian \[111, p. 31\].
In this way, the theory can at least locally be recast in terms of unconstrained variables (the true degrees of freedom) $`Q_l`$ and $`P_l`$ with $`l=1,\mathrm{},nr`$, for which the Dirac bracket equals the Poisson bracket, and for which the dynamics is governed by the Hamiltonian (4.40). The canonical variables $`\overline{Q}_k`$ and $`\overline{P}_k`$ ($`k=nr+1,\mathrm{},n`$) are the constraints and can hence be omitted in the description of the system. In this way, the dimension of the phase-space is reduced from $`2n`$ to $`2n2r`$, this is the number of phase-space variables we started with, minus the number of constraints.
One can also show that the true degrees of freedom are, at least locally, unique up to a canonical transformation \[111, p. 31\]. I.e. given another set of true degrees of freedom $`(\stackrel{~}{Q}_l,\stackrel{~}{P}_l)`$, then there exists a canonical transformation from $`(Q_l,P_l)`$ to $`(\stackrel{~}{Q}_l,\stackrel{~}{P}_l)`$, which maps the physical Hamiltonian $`H(P,Q)`$ to the physical Hamiltonian $`\stackrel{~}{H}(\stackrel{~}{Q},\stackrel{~}{P})`$ for the true degrees of freedom $`(\stackrel{~}{Q}_l,\stackrel{~}{P}_l)`$.
If we now quantize the system, by associating operators with the remaining $`2n2r`$ canonical variables, the commutation relations (4.36) for these operators become
$$[\widehat{Q}_{l_1},\widehat{Q}_{l_2}]=[\widehat{P}_{l_1},\widehat{P}_{l_2}]=0,[\widehat{Q}_{l_1},\widehat{P}_{l_2}]=i\delta _{l_1l_2},$$
(4.41)
with $`l_1,l_2=1,\mathrm{},nr`$.
The quantum description of the system then runs as follows. A quantum system is described by a vector $`|\psi `$ (the state vector) in a Hilbert space, with inner product $`\psi _2|\psi _1`$. The operators $`\widehat{Q}_l`$ and $`\widehat{P}_l`$ $`(l=1,\mathrm{},nr)`$ now act on these state vectors. In the Heisenberg picture, the operators $`\widehat{Q}_l`$ and $`\widehat{P}_l`$ are the dynamical objects and the states are time independent. However, in order to construct a pilot-wave interpretation we will not need the Heisenberg picture, but the Schrödinger picture. In the Schrödinger picture, the states and not the operators, are the dynamical objects. Because we have the standard canonical commutation relations for the operators $`\widehat{Q}_l`$ and $`\widehat{P}_l`$, we can use the standard representation
$$\widehat{Q}_l=Q_l,\widehat{P}_l=i\frac{}{Q_l},$$
(4.42)
for the operators in the Schrödinger picture. In this representation, the operators act on the wavefunction $`\mathrm{\Psi }(Q_1,\mathrm{},Q_{nr},t)=Q_1,\mathrm{},Q_{nr}|\mathrm{\Psi }(t)`$, with $`|Q_1,\mathrm{},Q_{nr}`$ the simultaneous eigenstates of the operators $`\widehat{Q}_l`$. In the representation (4.42) the Hamiltonian operator is written as $`\widehat{H}\left(Q,i/Q\right)`$ and the dynamical evolution of $`\mathrm{\Psi }(Q_1,\mathrm{},Q_{nr},t)`$ is given by the Schrödinger equation
$$i\frac{}{t}\mathrm{\Psi }(Q_1,\mathrm{},Q_{nr},t)=\widehat{H}(Q,i/Q)\mathrm{\Psi }(Q_1,\mathrm{},Q_{nr},t).$$
(4.43)
A pilot-wave interpretation can then be devised by looking at the conservation equation for the probability density $`|\mathrm{\Psi }(Q_1,\mathrm{},Q_{nr},t)|^2`$.
For a system described by a continuum number of degrees of freedom, the scheme to construct a pilot-wave interpretation proceeds along the same lines. The transition to a continuum number of degrees of freedom is discussed in Section 4.2.4. In the case of fermionic degrees of freedom we have to use a different representation from the one above. We discuss this in the next chapter.
#### 4.2.3 Canonical quantization of a system with first class constraints
Suppose now that all constraints $`\chi _j`$ are first class constraints. If there were also second class constraints, these could be dealt with separately, in the way described in the previous section. In the case a system has first class constraints, the matrix $`C_{jj^{}}`$ is singular and Dirac’s method of quantization of systems with only second class constraints cannot be applied. There are two ways to proceed<sup>5</sup><sup>5</sup>5We mention here only two methods of dealing with first class constraints, there are still other methods \[110, p. 110\]. However, these are less frequently used and, moreover, it is unclear whether these approaches may lead to a pilot-wave interpretation.:
* Constraints as operator identities: We have seen that the presence of first class constraints indicates the presence of some gauge invariance. This gauge invariance means that the evolution of the coordinates $`q`$ and $`p`$ is not uniquely fixed by their initial values; at each time one can perform a gauge transformation which yields a physical equivalent state. Only phase-space functions $`F(q,p)`$ which are gauge invariant are physically observable.
One can eliminate the gauge variables by adding further restrictions on the canonical variables. This is done by imposing further constraints, called gauge constraints. It is permissible to bring in these further constraints because they merely remove the arbitrary elements in the theory and do not affect the gauge invariant quantities.
A good set of gauge constraints
$$C_j^{}(q,p)0,$$
(4.44)
called a canonical gauge by Henneaux and Teitelboim \[109, p. 27\] and an admissible gauge by Sundermeyer \[110, p. 102\], satisfies the following two properties:
+ The gauge must be attainable. I.e. given a set of canonical variables $`q`$ and $`p`$ there exists a gauge transformation which brings the given set into one which satisfies (4.44). The transformation must be obtained by iteration of infinitesimal transformations of the form $`ϵ_j[F,\chi _j]`$, where $`F`$ represents the canonical variables $`q`$ and $`p`$.
+ Second, the conditions (4.44) must fix the gauge completely. As long as there would be a residual gauge freedom, the initial conditions on the canonical variables would be insufficient to uniquely determine their future evolution.
This condition implies that no gauge transformations but the identity preserve (4.44) or in other words that the equations $`_jϵ_j[C_j^{},\chi _j]0`$ must imply that $`ϵ_j=0`$.
One can show that the two conditions taken together imply that the number of gauge constraints must equal the number of first class constraints \[109, p. 27\]. The second condition further implies that the set of constraints $`\{C_j^{},\chi _j\}`$ is second class. This means that by adding the gauge constraints, the first class constraints turned second class so that we can quantize the system by using the Dirac bracket.
Once we have a suitable set of gauge constraints, we can, at least locally, perform a Maskawa-Nakajima canonical transformation to new canonical coordinates. In terms of these new coordinates the constraints form a set of canonical pairs. The other canonical pairs, which are unconstrained, then form the true degrees of freedom and the system can be expressed solely in terms of these true degrees of freedom. By definition, the true degrees of freedom have zero Poisson brackets with the constraints. In particular, the true degrees of freedom will have zero Poisson brackets with the first class constraints (the first class constraints in the new coordinates are found by performing the canonical transformation to the constraints $`\chi _j`$). This means that the true degrees of freedom are gauge independent variables. Therefore we will often refer to the true degrees of freedom as the gauge independent variables. The true degrees of freedom are unique up to a canonical transformation; in particular they are independent of the particular choice of admissible gauge (an extensive discussion of this can be found in \[111, pp. 36-60\]).
* Constraints as conditions on states: We can also quantize the canonical variables as if there were no constraints. The commutation relations for the operators associated with the canonical variables, are determined by the Poisson bracket as in equation (4.34). The constraints are reintroduced by demanding that physical states $`|\mathrm{\Psi }`$ satisfy
$$\widehat{\chi }_j|\mathrm{\Psi }=0.$$
(4.45)
The advantage of this method of dealing with constraints is that, because the commutation relations are simply the standard commutation relations, we can use the standard representation for the operators. If we quantize by treating the constraints as operators, the commutation relations are derived from the Dirac bracket. This can make it more complicated to find a suitable representation, because it requires the identification of the true degrees of freedom.
The disadvantage of dealing with constraints as conditions on states is that in general we will have to introduce a non-trivial measure on the configuration space in order to construct an inner product which yields finite numbers. A suitable measure can be found by applying the Faddeev-Popov formalism.
If we want to construct a pilot-wave model in the context of this method of dealing with constraints, we encounter a similar problem. The density of field beables will be non-normalizable. We will illustrate this explicitly in Section 4.5, where we will try to construct a pilot-wave theory for the quantized electromagnetic field by starting from this scheme. It will turn out that when we try to solve this problem, we are naturally led to the pilot-wave interpretation which may be obtained by quantizing the electromagnetic field by treating constraints as operator identities. We will discuss these issues in more detail in Section 4.5.4.
Finally, we want to note that these two schemes of dealing with first class constraints are very well suited for the quantization of the electromagnetic field because it presents no problem to fix the gauge globally. On the other hand, if we consider the quantization of non-Abelian gauge theories (Yang-Mills theories), it is more difficult to find a suitable gauge. This will lead to difficulties in both schemes of dealing with first class constraints. In the first scheme this is because a gauge is imposed from the start (as additional constraints). In the second scheme, this is because a gauge is needed to perform the Faddeev-Popov trick.
#### 4.2.4 Quantization of a field theory
So far we only have considered system which can be described by a finite number of degrees of freedom. The transition to a system which is described by a continuum number of degrees of freedom is straightforward. One can think of the transition as a replacement of the discrete label $`n`$ of the coordinates $`q_n`$ by a continuum label $`𝐱`$, i.e. $`q_n(t)=q(t,n)\psi _𝐱(t)=\psi (t,𝐱)`$. Usually the fields also carry an additional discrete label. Throughout this chapter, we will assume that the fields $`\psi _i(𝐱)`$ and their derivatives vanish sufficiently fast at spatial infinity. In this way possible boundary terms that arise when performing partial integration may be omitted.
Sums which appeared for systems with a finite number of degrees of freedom change to integrals. Derivatives with respect to the canonical coordinates $`/q_n`$ are replaced by functional derivatives $`\delta /\delta \psi _i(𝐱)`$.
As an example we can consider the definition of the Poisson bracket for fields. Let $`\psi _i(t,𝐱)`$ be the fields with canonically conjugate momenta $`\mathrm{\Pi }_{\psi _i}(t,𝐱)`$. The Poisson bracket for two functions $`F(\psi _i,\mathrm{\Pi }_{\psi _i})`$ and $`G(\psi _i,\mathrm{\Pi }_{\psi _i})`$ then reads
$$[F,G]_P=\underset{i}{}d^3x\left(\frac{\delta F}{\delta \psi _i(t,𝐱)}\frac{\delta G}{\delta \mathrm{\Pi }_{\psi _i}(t,𝐱)}\frac{\delta G}{\delta \psi _i(t,𝐱)}\frac{\delta F}{\delta \mathrm{\Pi }_{\psi _i}(t,𝐱)}\right).$$
(4.46)
In the expression all the fields are considered at the time $`t`$.
In order to construct a pilot-wave theory we will work in the functional Schrödinger picture. This is realized by using the representation
$$\widehat{\psi }_i(𝐱)=\psi _i(𝐱),\widehat{\mathrm{\Pi }}_{\psi _i}(𝐱)=i\frac{\delta }{\delta \psi _i(𝐱)}.$$
(4.47)
The Schrödinger equation is then a functional differential equation for the wavefunctional $`\mathrm{\Psi }(\psi _i(𝐱),t)`$.<sup>6</sup><sup>6</sup>6The wavefunctional is also called the super-wavefunction and the corresponding Schrödinger equation is called the super-Schrödinger equation .
If the fields $`\psi _i`$ represent the true degrees of freedom then we can define the inner product of two wavefunctionals as
$$\mathrm{\Psi }_1|\mathrm{\Psi }_2=\left(\mathrm{\Pi }_j𝒟\psi _j\right)\mathrm{\Psi }_1^{}(\psi _i,t)\mathrm{\Psi }_2(\psi _i,t),$$
(4.48)
with $`𝒟\psi _i=\mathrm{\Pi }_𝐱d\psi _i(𝐱)`$.
The wavefunctional can be written as $`\mathrm{\Psi }(\psi _i(𝐱),t)=\psi _i(𝐱)|\mathrm{\Psi }`$ where the $`|\psi _i(𝐱)`$ form a basis of the Hilbert space. They are the joint eigenstates of the operators $`\widehat{\psi }_i(𝐱)`$, i.e. $`\widehat{\psi }_j(𝐱)|\psi _i(𝐱)=\psi _j(𝐱)|\psi _i(𝐱)`$. In the standard quantum mechanical interpretation, the quantity
$$|\mathrm{\Psi }(\psi _i(𝐱),t)|^2=|\psi _i(𝐱)|\mathrm{\Psi }(t)|^2$$
(4.49)
can be interpreted as the probability density to find the system with wavefunctional $`\mathrm{\Psi }`$ in the field configuration $`(\psi _i(𝐱))`$. In the pilot-wave interpretation, we will introduce field beables $`\psi _i(t,𝐱)`$ which are distributed according to this density $`|\mathrm{\Psi }(\psi _i(𝐱),t)|^2`$ and for which the dynamics is governed by the guidance equations. As in the preceding chapters, the guidance equations will be derived from the continuity equation for $`|\mathrm{\Psi }(\psi _i(𝐱),t)|^2`$, by considering the analogy with the continuity equation in hydrodynamics.
We want to stress that in using the functional Schrödinger picture and in deriving the pilot-wave interpretation we will not adopt the greatest possible mathematical rigour. There are several problems associated to dealing with an infinite number of degrees of freedom, which will not be addressed in this thesis. First, there is the problem of infinities which plagues quantum field theory. Here, we will not make an effort to incorporate some renormalization scheme in devising the pilot-wave interpretations. Second there is the problem how to make mathematical sense out of the measure $`𝒟\psi _i`$; because it is a measure on an infinite dimensional configuration space it cannot be a Lebesgue measure. Third, the fields $`\psi _i(𝐱)`$ in the representation (4.47) can in fact not be treated as smooth functions, but should be distributions and accordingly they should be smeared \[114, p. 56\].
Although we do not elaborate on these problems, they certainly need attention in the future. The reader which is uncomfortable with our neglection of these problems can implicitly assume that we describe fields confined to a box of finite volume and with periodic boundary conditions, for which a cutoff is introduced for large momenta. Under these assumptions the fields are described by a finite number of degrees of freedom in momentum space, so that the above problems dissolve.
We can also make some general notes on locality and covariance of the pilot-wave field models. In Section 1.1 we noted the pilot-wave interpretation for non-relativistic quantum theory is nonlocal at the subquantum level. I.e. the motions of the particle beables are nonlocally correlated. But the nonlocality can not be used for superluminal signaling (at least not if the particles are distributed according to the quantum equilibrium hypothesis, which we assume in this thesis). Then in Section 3.2.3 we noted that the pilot-wave interpretation is not Lorentz invariant either. But, as with the locality, Lorentz invariance is satisfied on the empirical level. These properties also apply to the field beable approach. I.e. on the subquantum level the pilot-wave interpretations for quantum field theories, will be nonlocal and not Lorentz invariant. But at the empirical level the pilot-wave interpretation makes the same (statistical) predictions as the standard interpretation, so that at this level we regain locality and Lorentz invariance. These issues are thoroughly discussed in and we will not re-address them here.
### 4.3 The Duffin-Kemmer-Petiau theory in the spin-0 representation
In this section, we consider the construction of a pilot-wave interpretation for the Duffin-Kemmer-Petiau (DKP) theory in the spin-0 representation. First we show the equivalence of the quantized DKP theory and the quantized Klein-Gordon theory. Then afterwards we present the corresponding pilot-wave interpretation.
#### 4.3.1 Equivalence with the canonically quantized Klein-Gordon theory
We start with quantizing the DKP theory by applying Dirac’s canonical quantization procedure. We consider the DKP field coupled to a non-quantized electromagnetic field $`V^\mu =(V_0,𝐕)`$, which is introduced via the minimal coupling prescription $`_\mu D_\mu =_\mu +ieV_\mu `$.<sup>7</sup><sup>7</sup>7Lorentzian indices will be denoted by Greek letters $`\mu ,\nu ,\mathrm{}`$ and Euclidean indices will be denoted by Latin letters $`i,j,\mathrm{}`$ The Lorentzian indices are raised and lowered by the metric $`g_{\mu \nu }=\text{diag}(1,1,1,1)`$ and the Euclidean indices are raised and lowered by the metric $`\delta _{ij}=\text{diag}(1,1,1)`$. In this chapter we will often start from a Lagrangian density written in terms of Lorentzian vectors, but when we pass to the Hamiltonian formulation we will use Euclidean vectors. Note that in the previous chapter, the indices $`i,j,\mathrm{}`$ were used to denote the space index of tensorial objects. The equivalence of the quantized DKP theory in the spin-0 representation, coupled to a quantized electromagnetic field, and the quantized Klein-Gordon theory, coupled to a quantized electromagnetic field, can be shown in the same way. We will not do this explicitly here. In Section 4.6, where we discuss the pilot-wave interpretation for a massive bosonic field interacting with a quantized electromagnetic field, we will start from the coupled Klein-Gordon theory instead of from the coupled DKP theory.
If we write the five component DKP field as $`\psi =m^{1/2}(\varphi _\mu ,m\varphi )^T`$ with $`\varphi ^\mu =(\varphi _0,\mathit{\varphi })`$, then in the spin-0 representation given in Appendix A, the minimally coupled DKP Lagrangian (which is obtained from the Lagrangian density (3.80)) reads
$`L_K={\displaystyle d^3x_K}`$ $`=`$ $`{\displaystyle }d^3x({\displaystyle \frac{1}{2}}(\varphi _\mu ^{}D^\mu \varphi \varphi ^{}D_\mu \varphi ^\mu +(D_\mu \varphi )^{}\varphi ^\mu (D_\mu \varphi ^\mu )^{}\varphi )`$ (4.50)
$`m^2\varphi ^{}\varphi \varphi _\mu ^{}\varphi ^\mu ).`$
The equations of motion are
$`D_\mu \varphi ^\mu +m^2\varphi `$ $`=0,`$ $`\left(D_\mu \varphi ^\mu \right)^{}+m^2\varphi ^{}=0,`$
$`D_\mu \varphi \varphi _\mu `$ $`=0,`$ $`\left(D_\mu \varphi \right)^{}\varphi _\mu ^{}=0.`$ (4.51)
The canonically conjugate momenta are<sup>8</sup><sup>8</sup>8Often the space dependence of the fields will not be written explicitly. E.g. instead of writing $`\mathrm{\Pi }_\psi (𝐱)=\frac{\delta L}{\delta \dot{\psi }(𝐱)}`$ we write $`\mathrm{\Pi }_\psi =\frac{\delta L}{\delta \dot{\psi }}`$. Similarly instead of writing the constraints as $`\chi _\psi (𝐱)=0`$ we write them as $`\chi _\psi =0`$. In this way the constraint $`\chi _\psi =0`$ represents in fact an infinite number of constraints; one corresponding to each point in space.
$`\mathrm{\Pi }_{\varphi _0}=`$ $`{\displaystyle \frac{\delta L}{\delta \dot{\varphi }_0}}={\displaystyle \frac{\varphi ^{}}{2}},\mathrm{\Pi }_{\varphi _0^{}}`$ $`={\displaystyle \frac{\delta L}{\delta \dot{\varphi }_0^{}}}={\displaystyle \frac{\varphi }{2}},`$
$`\mathrm{\Pi }_{\varphi _i}=`$ $`{\displaystyle \frac{\delta L}{\delta \dot{\varphi }_i}}=0,\mathrm{\Pi }_{\varphi _i^{}}`$ $`={\displaystyle \frac{\delta L}{\delta \dot{\varphi }_i^{}}}=0,`$
$`\mathrm{\Pi }_\varphi =`$ $`{\displaystyle \frac{\delta L}{\delta \dot{\varphi }}}={\displaystyle \frac{\varphi _0^{}}{2}},\mathrm{\Pi }_\varphi ^{}`$ $`={\displaystyle \frac{\delta L}{\delta \dot{\varphi }^{}}}={\displaystyle \frac{\varphi _0}{2}}.`$ (4.52)
We can immediately identify the primary constraints
$`\chi _{\varphi _0}=`$ $`\mathrm{\Pi }_{\varphi _0}+{\displaystyle \frac{\varphi ^{}}{2}},\chi _{\varphi _0^{}}`$ $`=\mathrm{\Pi }_{\varphi _0^{}}+{\displaystyle \frac{\varphi }{2}},`$
$`\chi _{\varphi _i}=`$ $`\mathrm{\Pi }_{\varphi _i},\chi _{\varphi _i^{}}`$ $`=\mathrm{\Pi }_{\varphi _i^{}},`$
$`\chi _\varphi =`$ $`\mathrm{\Pi }_\varphi {\displaystyle \frac{\varphi _0^{}}{2}},\chi _\varphi ^{}`$ $`=\mathrm{\Pi }_\varphi ^{}{\displaystyle \frac{\varphi _0}{2}}.`$ (4.53)
The corresponding canonical Hamiltonian reads
$$H_C=d^3x\left(\varphi ^{}D_i\varphi _i+\varphi D_i^{}\varphi _i^{}+m^2\varphi ^{}\varphi +\varphi _0^{}\varphi _0\varphi _i^{}\varphi _i+ieV_0(\varphi ^{}\varphi _0\varphi \varphi _0^{})\right).$$
(4.54)
The total Hamiltonian reads
$$H_T=H_C+\underset{\gamma }{}d^3xu_\gamma (𝐱)\chi _\gamma (𝐱),$$
(4.55)
where the label $`\gamma `$ takes the values $`\varphi _0,\varphi _0^{},\varphi ,\varphi ^{},\varphi _i,\varphi _i^{}`$ and the $`u_\gamma `$ are arbitrary fields. In order to find out whether there are secondary constraints we impose the consistency conditions that the primary constraints are weakly conserved in time, i.e. $`[\chi _\gamma ,H_T]_P0`$. The conditions that the constraints $`\chi _{\varphi _0},\chi _{\varphi _0^{}},\chi _\varphi ,\chi _\varphi ^{}`$ are conserved yield respectively
$`u_\varphi ^{}`$ $`=`$ $`\varphi _0^{}+ieV_0\varphi ^{},`$
$`u_\varphi `$ $`=`$ $`\varphi _0ieV_0\varphi ,`$
$`u_{\varphi _0^{}}`$ $`=`$ $`D_i^{}\varphi _i^{}m^2\varphi ^{}+ieV_0\varphi _0^{},`$
$`u_{\varphi _0}`$ $`=`$ $`D_i\varphi _im^2\varphi ieV_0\varphi _0.`$ (4.56)
Hence these consistency conditions determine some of the arbitrary fields $`u_\gamma `$ and do not lead to further constraints. The conditions that $`\varphi _i`$ and $`\varphi _i^{}`$ are weakly conserved lead to the secondary constraints
$$\chi _{s\varphi _i}=D_i\varphi +\varphi _i,\chi _{s\varphi _i^{}}=D_i^{}\varphi ^{}+\varphi _i^{}.$$
(4.57)
In turn, the requirement that the secondary constraints $`\chi _{s\varphi _i}`$ and $`\chi _{s\varphi _i^{}}`$ are conserved, determine the fields $`u_{\varphi _i^{}}`$ and $`u_{\varphi _i}`$
$$u_{\varphi _i^{}}=D_i^{}u_\varphi ^{},u_{\varphi _i}=D_iu_\varphi .$$
(4.58)
Hence we see that all the fields $`u_\gamma `$ are determined. This means that all the constraints are second class constraints and the system can be quantized by using the Dirac bracket.
In order to construct the Dirac bracket we need the inverse of the matrix (cf. Section 4.2.2)
$$C_{\gamma \gamma ^{}}(𝐱,𝐲)=[\chi _\gamma (𝐱),\chi _\gamma ^{}(𝐲)]_P$$
(4.59)
where the labels $`\gamma `$ and $`\gamma ^{}`$ take the values $`\varphi _0,\varphi _0^{},\varphi _i,\varphi _i^{},s\varphi _i,s\varphi _i^{}`$. This matrix has nonzero components
$`C_{\varphi _0^{},\varphi }(𝐱,𝐲)`$ $`=`$ $`C_{\varphi ,\varphi _0^{}}(𝐲,𝐱)=\delta (𝐱𝐲),`$
$`C_{\varphi _0,\varphi ^{}}(𝐱,𝐲)`$ $`=`$ $`C_{\varphi ^{},\varphi _0}(𝐲,𝐱)=\delta (𝐱𝐲),`$
$`C_{\varphi ,s\varphi _i}(𝐱,𝐲)`$ $`=`$ $`C_{s\varphi _i,\varphi }(𝐲,𝐱)=D_{x_i}^{}\delta (𝐱𝐲),`$
$`C_{\varphi ^{},s\varphi _i^{}}(𝐱,𝐲)`$ $`=`$ $`C_{s\varphi _i^{},\varphi ^{}}(𝐲,𝐱)=D_{x_i}\delta (𝐱𝐲),`$
$`C_{s\varphi _i,\varphi _j}(𝐱,𝐲)`$ $`=`$ $`C_{\varphi _i,s\varphi _j}(𝐲,𝐱)=\delta _{ij}\delta (𝐱𝐲),`$
$`C_{s\varphi _i^{},\varphi _j^{}}(𝐱,𝐲)`$ $`=`$ $`C_{\varphi _i^{},s\varphi _j^{}}(𝐲,𝐱)=\delta _{ij}\delta (𝐱𝐲).`$ (4.60)
The inverse $`C_{\gamma \gamma ^{}}^1`$ has the following nonzero components
$`C_{\varphi ,\varphi _0^{}}^1(𝐱,𝐲)`$ $`=`$ $`C_{\varphi _0^{},\varphi }^1(𝐲,𝐱)=\delta (𝐱𝐲),`$
$`C_{\varphi ^{},\varphi _0}^1(𝐱,𝐲)`$ $`=`$ $`C_{\varphi _0,\varphi ^{}}^1(𝐲,𝐱)=\delta (𝐱𝐲),`$
$`C_{\varphi _i^{},\varphi _0}^1(𝐱,𝐲)`$ $`=`$ $`C_{\varphi _0,\varphi _i^{}}^1(𝐲,𝐱)=D_{x_i}\delta (𝐱𝐲),`$
$`C_{\varphi _0^{},\varphi _i}^1(𝐱,𝐲)`$ $`=`$ $`C_{\varphi _i,\varphi _0^{}}^1(𝐲,𝐱)=D_{x_i}\delta (𝐱𝐲),`$
$`C_{\varphi _i,s\varphi _j}^1(𝐱,𝐲)`$ $`=`$ $`C_{s\varphi _i,\varphi _j}^1(𝐲,𝐱)=\delta _{ij}\delta (𝐱𝐲),`$
$`C_{\varphi _i^{},s\varphi _j^{}}^1(𝐱,𝐲)`$ $`=`$ $`C_{s\varphi _i^{},\varphi _j^{}}^1(𝐲,𝐱)=\delta _{ij}\delta (𝐱𝐲).`$ (4.61)
The Dirac bracket for the fields $`\varphi ,\varphi ^{},\varphi _0,\varphi _0^{}`$ now reads
$`[\varphi (𝐱),\varphi _0^{}(𝐲)]_D`$ $`=`$ $`\delta (𝐱𝐲),`$
$`[\varphi ^{}(𝐱),\varphi _0(𝐲)]_D`$ $`=`$ $`\delta (𝐱𝐲)`$ (4.62)
and all the other Dirac brackets involving the fields $`\varphi ,\varphi ^{},\varphi _0,\varphi _0^{}`$ are zero. Because it is a property of the Dirac bracket that one can impose the constraints before evaluating the Dirac bracket, the commutation relations involving other fields can be derived from the commutation relations of $`\varphi ,\varphi ^{},\varphi _0,\varphi _0^{}`$, by using the constraints (4.53) and (4.57). Therefore there is no need to give them explicitly.
In order to find the true degrees of freedom, we perform the Maskawa-Nakajima transformation to new canonical variables, which we denote with an additional twidle:
$`\stackrel{~}{\varphi }=`$ $`{\displaystyle \frac{\varphi }{2}}\mathrm{\Pi }_{\varphi _0^{}},\stackrel{~}{\varphi }^{}`$ $`={\displaystyle \frac{\varphi ^{}}{2}}\mathrm{\Pi }_{\varphi _0},`$
$`\mathrm{\Pi }_{\stackrel{~}{\varphi }}=`$ $`{\displaystyle \frac{\varphi _0^{}}{2}}+\mathrm{\Pi }_\varphi +D_i^{}\mathrm{\Pi }_{\varphi _i},\mathrm{\Pi }_{\stackrel{~}{\varphi }^{}}`$ $`={\displaystyle \frac{\varphi _0}{2}}+\mathrm{\Pi }_\varphi ^{}+D_i\mathrm{\Pi }_{\varphi _i^{}},`$
$`\stackrel{~}{\varphi }_0=`$ $`{\displaystyle \frac{\varphi _0}{2}}\mathrm{\Pi }_\varphi ^{}D_i\mathrm{\Pi }_{\varphi _i^{}},\stackrel{~}{\varphi }_0^{}`$ $`={\displaystyle \frac{\varphi _0^{}}{2}}\mathrm{\Pi }_\varphi D_i^{}\mathrm{\Pi }_{\varphi _i},`$
$`\mathrm{\Pi }_{\stackrel{~}{\varphi }_0}=`$ $`\mathrm{\Pi }_{\varphi _0}+{\displaystyle \frac{\varphi ^{}}{2}},\mathrm{\Pi }_{\stackrel{~}{\varphi }_0^{}}`$ $`=\mathrm{\Pi }_{\varphi _0^{}}+{\displaystyle \frac{\varphi }{2}},`$
$`\stackrel{~}{\varphi }_i=`$ $`\varphi _i+D_i\varphi ,\stackrel{~}{\varphi }_i^{}`$ $`=\varphi _i^{}+D_i^{}\varphi ^{},`$
$`\mathrm{\Pi }_{\stackrel{~}{\varphi }_i}=`$ $`\mathrm{\Pi }_{\varphi _i},\mathrm{\Pi }_{\stackrel{~}{\varphi }_i^{}}`$ $`=\mathrm{\Pi }_{\varphi _i^{}}.`$ (4.63)
For the new variables the constraints read
$$\stackrel{~}{\varphi }_0=\mathrm{\Pi }_{\stackrel{~}{\varphi }_0}=\stackrel{~}{\varphi }_0^{}=\mathrm{\Pi }_{\stackrel{~}{\varphi }_0^{}}=\stackrel{~}{\varphi }_i=\mathrm{\Pi }_{\stackrel{~}{\varphi }_i}=\stackrel{~}{\varphi }_i^{}=\mathrm{\Pi }_{\stackrel{~}{\varphi }_i^{}}=0,$$
(4.64)
so that the true degrees of freedom are $`\stackrel{~}{\varphi },\stackrel{~}{\varphi }^{},\mathrm{\Pi }_{\stackrel{~}{\varphi }},\mathrm{\Pi }_{\stackrel{~}{\varphi }^{}}`$. The physical Hamiltonian $`H`$ is found by performing the canonical transformation (4.63) on the canonical Hamiltonian $`H_C`$ and by imposing the constraints:
$$H=d^3x\left(\mathrm{\Pi }_{\stackrel{~}{\varphi }^{}}\mathrm{\Pi }_{\stackrel{~}{\varphi }}+\left(D_i^{}\stackrel{~}{\varphi }^{}\right)D_i\stackrel{~}{\varphi }+m^2\stackrel{~}{\varphi }^{}\stackrel{~}{\varphi }+ieV_0\left(\stackrel{~}{\varphi }^{}\mathrm{\Pi }_{\stackrel{~}{\varphi }^{}}\stackrel{~}{\varphi }\mathrm{\Pi }_{\stackrel{~}{\varphi }}\right)\right).$$
(4.65)
This Hamiltonian is recognized as the Klein-Gordon Hamiltonian. Because the Dirac bracket for the true degrees of freedom is simply the Poisson bracket, the theory is quantized by imposing the commutation relations
$`[\widehat{\stackrel{~}{\varphi }}(𝐱),\widehat{\mathrm{\Pi }}_{\stackrel{~}{\varphi }}(𝐲)]`$ $`=`$ $`i[\stackrel{~}{\varphi }(𝐱),\mathrm{\Pi }_{\stackrel{~}{\varphi }}(𝐲)]_P=i\delta (𝐱𝐲),`$
$`[\widehat{\stackrel{~}{\varphi }^{}}(𝐱),\widehat{\mathrm{\Pi }}_{\stackrel{~}{\varphi }^{}}(𝐲)]`$ $`=`$ $`i[\stackrel{~}{\varphi }^{}(𝐱),\mathrm{\Pi }_{\stackrel{~}{\varphi }^{}}(𝐲)]_P=i\delta (𝐱𝐲).`$ (4.66)
The other fundamental commutation relations involving the operators $`\widehat{\stackrel{~}{\varphi }},\widehat{\mathrm{\Pi }}_{\stackrel{~}{\varphi }},\widehat{\stackrel{~}{\varphi }^{}}`$ and $`\widehat{\mathrm{\Pi }}_{\stackrel{~}{\varphi }^{}}`$ are zero. The commutation relations are realized by the representation
$`\widehat{\stackrel{~}{\varphi }}(𝐱)`$ $`=`$ $`\varphi (𝐱),\widehat{\mathrm{\Pi }}_{\stackrel{~}{\varphi }}(𝐱)=i{\displaystyle \frac{\delta }{\delta \varphi (𝐱)}},`$
$`\widehat{\stackrel{~}{\varphi }^{}}(𝐱)`$ $`=`$ $`\varphi ^{}(𝐱),\widehat{\mathrm{\Pi }}_{\stackrel{~}{\varphi }^{}}(𝐱)=i{\displaystyle \frac{\delta }{\delta \varphi ^{}(𝐱)}}.`$ (4.67)
In this representation the Hamiltonian (4.65) reads<sup>9</sup><sup>9</sup>9In fact there appears an operator ordering ambiguity at this point. The term proportional to $`V_0`$ in the Hamiltonian arises from associating operators to $`\stackrel{~}{\varphi }^{}\mathrm{\Pi }_{\stackrel{~}{\varphi }^{}}\stackrel{~}{\varphi }\mathrm{\Pi }_{\stackrel{~}{\varphi }}`$ (this quantity is proportional to the charge density). We have chosen $`\widehat{\stackrel{~}{\varphi }^{}}\widehat{\mathrm{\Pi }}_{\stackrel{~}{\varphi }^{}}\widehat{\stackrel{~}{\varphi }}\widehat{\mathrm{\Pi }}_{\stackrel{~}{\varphi }}`$. This operator ordering is also the Weyl ordering \[115, p. 347\]. Another operator ordering choice is $`\widehat{\stackrel{~}{\varphi }^{}}\widehat{\mathrm{\Pi }}_{\stackrel{~}{\varphi }^{}}\widehat{\mathrm{\Pi }}_{\stackrel{~}{\varphi }}\widehat{\stackrel{~}{\varphi }}`$. Although this choice is not Weyl ordered, it is Hermitian. The same operator ordering ambiguity arises when we quantize the coupled Klein-Gordon theory. So one can always choose operator orderings so that the quantized DKP theory is equivalent with the quantized Klein-Gordon theory. A similar remark will apply in the spin-1 case.
$$\widehat{H}=d^3x\left(\frac{\delta }{\delta \varphi ^{}}\frac{\delta }{\delta \varphi }+|D_i\varphi |^2+m^2|\varphi |^2+eV_0\left(\varphi ^{}\frac{\delta }{\delta \varphi ^{}}\varphi \frac{\delta }{\delta \varphi }\right)\right).$$
(4.68)
The operators act on wavefunctionals $`\mathrm{\Psi }(\varphi ,\varphi ^{},t)=\varphi ,\varphi ^{}|\mathrm{\Psi }(t)`$, so that we obtain the functional Schrödinger equation
$$i\frac{\mathrm{\Psi }(\varphi ,\varphi ^{},t)}{t}=\widehat{H}\mathrm{\Psi }(\varphi ,\varphi ^{},t).$$
(4.69)
This quantum field theory can also be derived from quantizing the Klein-Gordon theory.<sup>10</sup><sup>10</sup>10Fainberg and Pimentel also established the result that the canonical quantization of the DKP theory in the spin-0 representation leads to the canonically quantized Klein-Gordon theory. They also give a strict proof of equivalence of the theories for the method of path-integral quantization.
In fact there was also a shorter route to get to the quantized theory. From the constraints we can already read of that the fields $`\varphi ,\varphi ^{},\varphi _0,\varphi _0^{}`$ can be taken as independent degrees of freedom, because all the other canonical variables can be expressed in terms of these fields using the constraints. We can then quantize the theory by imposing the commutation relations
$`[\widehat{\varphi }(𝐱),\widehat{\varphi }_0^{}(𝐲)]`$ $`=`$ $`i[\varphi (𝐱),\varphi _0^{}(𝐲)]_D=i\delta (𝐱𝐲),`$
$`[\widehat{\varphi }^{}(𝐱),\widehat{\varphi }_0(𝐲)]`$ $`=`$ $`i[\varphi ^{}(𝐱),\varphi _0(𝐲)]_D=i\delta (𝐱𝐲).`$ (4.70)
Operators corresponding to the other degrees of freedom can be expressed in terms of the operators $`\widehat{\varphi },\widehat{\varphi }^{},\widehat{\varphi }_0`$ and $`\widehat{\varphi }_0^{}`$. If we then realize the commutation relations by the representation
$`\widehat{\varphi }(𝐱)`$ $`=`$ $`\varphi (𝐱),\widehat{\varphi }_0^{}(𝐱)=i{\displaystyle \frac{\delta }{\delta \varphi (𝐱)}},`$
$`\widehat{\varphi }^{}(𝐱)`$ $`=`$ $`\varphi ^{}(𝐱),\widehat{\varphi }_0(𝐱)=i{\displaystyle \frac{\delta }{\delta \varphi ^{}(𝐱)}},`$ (4.71)
we obtain the same quantum theory as the one presented above. If we look at the Maskawa-Nakajima canonical transformation (4.63), then we see that on the constraint space, the true degrees of freedom read
$`\stackrel{~}{\varphi }`$ $`=`$ $`\varphi ,\stackrel{~}{\varphi }^{}=\varphi ^{},`$
$`\mathrm{\Pi }_{\stackrel{~}{\varphi }}`$ $`=`$ $`\varphi _0^{},\mathrm{\Pi }_{\stackrel{~}{\varphi }^{}}=\varphi _0.`$ (4.72)
Hence the fields $`\varphi ,\varphi ^{},\varphi _0,\varphi _0^{}`$ could be regarded as the true degrees of freedom from the start. In the rest of the chapter we will often take this shortcut to obtain the quantized theory. The independent degrees of freedom will be identified as the true degrees of freedom, and the commutation relations of the corresponding operators will be derived from the Dirac bracket.
#### 4.3.2 Pilot-wave interpretation
Now that we have the quantum field theory, we can construct a pilot-wave interpretation. The conservation equation corresponding to the functional Schrödinger equation (4.69) reads
$$\frac{|\mathrm{\Psi }|^2}{t}+d^3x\left(\frac{\delta J_\varphi }{\delta \varphi }+\frac{\delta J_\varphi ^{}}{\delta \varphi ^{}}\right)=0,$$
(4.73)
with
$`J_\varphi `$ $`=`$ $`{\displaystyle \frac{1}{2i}}\left(\mathrm{\Psi }^{}{\displaystyle \frac{\delta }{\delta \varphi ^{}}}\mathrm{\Psi }\mathrm{\Psi }{\displaystyle \frac{\delta }{\delta \varphi ^{}}}\mathrm{\Psi }^{}\right)ieV_0|\mathrm{\Psi }|^2\varphi ,`$
$`J_\varphi ^{}`$ $`=`$ $`{\displaystyle \frac{1}{2i}}\left(\mathrm{\Psi }^{}{\displaystyle \frac{\delta }{\delta \varphi }}\mathrm{\Psi }\mathrm{\Psi }{\displaystyle \frac{\delta }{\delta \varphi }}\mathrm{\Psi }^{}\right)+ieV_0|\mathrm{\Psi }|^2\varphi ^{}.`$ (4.74)
In the pilot-wave interpretation we introduce the field beables $`\varphi `$ and $`\varphi ^{}`$ whose motion is governed by the wavefunctional via the guidance equations
$$\dot{\varphi }=J_\varphi /|\mathrm{\Psi }|^2,\dot{\varphi }^{}=J_\varphi ^{}/|\mathrm{\Psi }|^2.$$
(4.75)
The density of field beables is given by the equilibrium density $`|\mathrm{\Psi }|^2`$. In the free case this pilot-wave interpretation reduces to the one originally presented by Kaloyerou and by Valentini .
### 4.4 The Duffin-Kemmer-Petiau theory in the spin-1 representation
In this section, we consider the construction of a pilot-wave interpretation for the Duffin-Kemmer-Petiau (DKP) theory in the spin-1 representation. First we show the equivalence of the quantized DKP theory and the quantized Proca theory. Then afterwards we present the corresponding pilot-wave interpretation.
#### 4.4.1 Equivalence with the canonically quantized Proca theory
As in the spin-0 case we consider the DKP field coupled to a non-quantized electromagnetic field $`V^\mu `$. The treatment of the DKP field coupled to a quantized electromagnetic field is completely analogous and will not be treated explicitly.
If we write the ten component DKP field $`\psi `$ as
$$\psi =m^{1/2}(𝐄,𝐁,m𝐀,mA_0)^T,$$
(4.76)
then in the spin-1 representation given in Appendix A, the minimally coupled DKP Lagrangian reads
$`L_K={\displaystyle d^3x_K}`$ $`=`$ $`{\displaystyle }d^3x({\displaystyle \frac{1}{2}}(A_i^{}D_0E_i+A_iD_0^{}E_i^{}E_i^{}D_0A_iE_iD_0^{}A_i^{})`$ (4.77)
$`+A_0^{}D_iE_iE_i^{}D_iA_0\epsilon _{ijk}(B_i^{}D_jA_k+A_i^{}D_jB_k)`$
$`E_i^{}E_i+B_i^{}B_i+m^2A_0^{}A_0m^2A_i^{}A_i)`$
The equations of motion are
$`D_0E_i=`$ $`\epsilon _{ijk}D_jB_k+m^2A_i,D_0^{}E_i^{}`$ $`=\epsilon _{ijk}D_j^{}B_k^{}+m^2A_i^{},`$
$`D_0A_i=`$ $`E_iD_iA_0,D_0^{}A_i^{}`$ $`=E_i^{}D_i^{}A_0^{},`$
$`D_iE_i=`$ $`m^2A_0,D_i^{}E_i^{}`$ $`=m^2A_0^{},`$
$`B_i=`$ $`\epsilon _{ijk}D_jA_k,B_i^{}`$ $`=\epsilon _{ijk}D_j^{}A_k^{}.`$ (4.78)
The canonically conjugate momenta are
$`\mathrm{\Pi }_{A_i}=`$ $`{\displaystyle \frac{\delta L}{\delta \dot{A}_i}}={\displaystyle \frac{1}{2}}E_i^{},\mathrm{\Pi }_{A_i^{}}=`$ $`{\displaystyle \frac{\delta L}{\delta \dot{A}_i^{}}}={\displaystyle \frac{1}{2}}E_i,`$
$`\mathrm{\Pi }_{A_0}=`$ $`{\displaystyle \frac{\delta L}{\delta \dot{A}_0}}=0,\mathrm{\Pi }_{A_0^{}}=`$ $`{\displaystyle \frac{\delta L}{\delta \dot{A}_0^{}}}=0,`$
$`\mathrm{\Pi }_{E_i}=`$ $`{\displaystyle \frac{\delta L}{\delta \dot{E}_i}}={\displaystyle \frac{1}{2}}A_i^{},\mathrm{\Pi }_{E_i^{}}=`$ $`{\displaystyle \frac{\delta L}{\delta \dot{E}_i^{}}}={\displaystyle \frac{1}{2}}A_i,`$
$`\mathrm{\Pi }_{B_i}=`$ $`{\displaystyle \frac{\delta L}{\delta \dot{B}_i}}=0,\mathrm{\Pi }_{B_i^{}}=`$ $`{\displaystyle \frac{\delta L}{\delta \dot{B}_i^{}}}=0.`$ (4.79)
Hence the primary constraints are
$`\chi _{A_i}=`$ $`\mathrm{\Pi }_{A_i}+{\displaystyle \frac{1}{2}}E_i^{},\chi _{A_i^{}}`$ $`=\mathrm{\Pi }_{A_i^{}}+{\displaystyle \frac{1}{2}}E_i,`$
$`\chi _{A_0}=`$ $`\mathrm{\Pi }_{A_0},\chi _{A_0^{}}`$ $`=\mathrm{\Pi }_{A_0^{}},`$
$`\chi _{E_i}=`$ $`\mathrm{\Pi }_{E_i}{\displaystyle \frac{1}{2}}A_i^{},\chi _{E_i^{}}`$ $`=\mathrm{\Pi }_{E_i^{}}{\displaystyle \frac{1}{2}}A_i,`$
$`\chi _{B_i}=`$ $`\mathrm{\Pi }_{B_i},\chi _{B_i^{}}`$ $`=\mathrm{\Pi }_{B_i^{}}.`$ (4.80)
The canonical Hamiltonian reads
$`H_C`$ $`=`$ $`{\displaystyle }d^3x(E_i^{}D_iA_0+\epsilon _{ijk}B_i^{}D_jA_k+E_iD_i^{}A_0^{}+\epsilon _{ijk}B_iD_j^{}A_k^{}`$ (4.81)
$`+E_i^{}E_iB_i^{}B_i+m^2(A_i^{}A_iA_0^{}A_0)+ieV_0(A_iE_i^{}A_i^{}E_i)).`$
The total Hamiltonian reads
$$H_T=H_C+\underset{\gamma }{}d^3xu_\gamma (𝐱)\chi _\gamma (𝐱),$$
(4.82)
where the label $`\gamma `$ takes the values $`A_i,A_i^{},A_0,A_0^{},E_i,E_i^{},B_i,B_i^{}`$ and the $`u_\gamma `$ are arbitrary fields. The requirement that the constraints $`\chi _{A_i},\chi _{A_i^{}},\chi _{E_i}`$ and $`\chi _{E_i^{}}`$ are weakly conserved in time respectively yield
$`u_{E_i^{}}`$ $`=`$ $`\epsilon _{ijk}D_k^{}B_j^{}+m^2A_i^{}+ieV_0E_i^{},`$
$`u_{E_i}`$ $`=`$ $`\epsilon _{ijk}D_kB_j+m^2A_iieV_0E_i,`$
$`u_{A_i^{}}`$ $`=`$ $`D_i^{}A_0^{}+ieV_0A_i^{}E_i^{},`$
$`u_{A_i}`$ $`=`$ $`D_iA_0ieV_0A_iE_i.`$ (4.83)
The consistency requirements that the constraints $`\chi _{A_0},\chi _{A_0^{}},\chi _{B_i}`$ and $`\chi _{B_i^{}}`$ are weakly conserved in time yield the secondary constraints
$`\chi _{sA_0^{}}`$ $`=`$ $`D_i^{}E_i^{}m^2A_0^{},\chi _{sA_0}=D_iE_im^2A_0,`$
$`\chi _{sB_i^{}}`$ $`=`$ $`\epsilon _{ijk}D_j^{}A_k^{}B_i^{},\chi _{sB_i}=\epsilon _{ijk}D_jA_kB_i.`$ (4.84)
The requirement that the constraints $`\chi _{sA_0},\chi _{sA_0^{}},\chi _{sB_i}`$ and $`\chi _{sB_i^{}}`$ are weakly conserved in time fix the the remaining fields $`u_{A_0},u_{A_0^{}},u_{B_i}`$ and $`u_{B_i^{}}`$
$`u_{A_0}`$ $`=`$ $`D_i^{}u_{E_i^{}}/m^2,u_{A_0^{}}=D_iu_{E_i}/m^2,`$
$`u_{B_i}`$ $`=`$ $`\epsilon _{ijk}D_ju_{A_k},u_{B_i^{}}=\epsilon _{ijk}D_j^{}u_{A_k^{}}.`$ (4.85)
Because all the fields $`u_\gamma `$ are determined by the consistency conditions, all constraints are second class constraints.
At this stage we can already identify the true degrees of freedom. From the constraints we see that we can take the fields $`A_i,A_i^{},E_i,E_i^{}`$ as the true degrees of freedom. When the constraints are imposed, the canonical Hamiltonian can be written in terms of these fields<sup>11</sup><sup>11</sup>11The distribution $`h_{ij}`$ is introduced for notational convenience. By using this distribution the Hamiltonian seems to depend nonlocally on the fields. However, this is only apparently because one can explicitly perform the integration over $`h_{ij}(𝐲,𝐳)`$ in order to obtain a local Hamiltonian.
$`H`$ $`=`$ $`{\displaystyle }d^3yd^3zh_{ij}(𝐲,𝐳)E_i^{}(𝐲)E_j(𝐳)+{\displaystyle }d^3x({\displaystyle \frac{1}{2}}G_{ij}^{}(𝐱)G_{ij}(𝐱)`$ (4.86)
$`+m^2A_i^{}(𝐱)A_i(𝐱)+ieV_0(A_i(𝐱)E_i^{}(𝐱)A_i^{}(𝐱)E_i(𝐱))),`$
with
$`G_{ij}`$ $`=`$ $`D_iA_jD_jA_i,`$
$`h_{ij}(𝐲,𝐳)`$ $`=`$ $`\left({\displaystyle \frac{1}{m^2}}D_{y_i}D_{y_j}+\delta _{ij}\right)\delta (𝐲𝐳).`$ (4.87)
In order to quantize the system we need the Dirac bracket. The matrix $`C_{\gamma \gamma ^{}}(𝐱,𝐲)`$, where the labels $`\gamma `$ and $`\gamma ^{}`$ take the values $`A_0,A_0^{},A_i,A_i^{},E_i,E_i^{},B_i,B_i^{},`$ $`sA_0,sA_0^{},sB_i,sB_i^{}`$, has the following nonzero components
$`C_{A_i,E_j^{}}(𝐱,𝐲)`$ $`=`$ $`C_{E_j^{},A_i}(𝐲,𝐱)=\delta _{ij}\delta (𝐱𝐲),`$
$`C_{A_i^{},E_j}(𝐱,𝐲)`$ $`=`$ $`C_{E_j,A_i^{}}(𝐲,𝐱)=\delta _{ij}\delta (𝐱𝐲),`$
$`C_{A_0,sA_0}(𝐱,𝐲)`$ $`=`$ $`C_{sA_0,A_0}(𝐲,𝐱)=m^2\delta (𝐱𝐲),`$
$`C_{A_0^{},sA_0^{}}(𝐱,𝐲)`$ $`=`$ $`C_{sA_0^{},A_0^{}}(𝐲,𝐱)=m^2\delta (𝐱𝐲),`$
$`C_{B_j,sB_i}(𝐱,𝐲)`$ $`=`$ $`C_{sB_i,B_j}(𝐲,𝐱)=\delta _{ij}\delta (𝐱𝐲),`$
$`C_{B_j^{},sB_i^{}}(𝐱,𝐲)`$ $`=`$ $`C_{sB_i^{},B_j^{}}(𝐲,𝐱)=\delta _{ij}\delta (𝐱𝐲),`$
$`C_{E_i,sA_0}(𝐱,𝐲)`$ $`=`$ $`C_{sA_0,E_i}(𝐲,𝐱)=D_{x_i}^{}\delta (𝐱𝐲),`$
$`C_{E_i^{},sA_0^{}}(𝐱,𝐲)`$ $`=`$ $`C_{sA_0^{},E_i^{}}(𝐲,𝐱)=D_{x_i}\delta (𝐱𝐲),`$
$`C_{sB_j,A_i}(𝐱,𝐲)`$ $`=`$ $`C_{A_i,sB_j}(𝐲,𝐱)=\epsilon _{ijk}D_{x_k}^{}\delta (𝐱𝐲),`$
$`C_{A_i^{},sB_j^{}}(𝐱,𝐲)`$ $`=`$ $`C_{sB_j^{},A_i^{}}(𝐲,𝐱)=\epsilon _{ijk}D_{x_k}\delta (𝐱𝐲).`$ (4.88)
Its inverse $`C_{\gamma \gamma ^{}}^1`$ has the following nonzero components
$`C_{E_j^{},A_i}^1(𝐱,𝐲)`$ $`=`$ $`C_{A_i,E_j^{}}^1(𝐲,𝐱)=\delta _{ij}\delta (𝐱𝐲),`$
$`C_{E_i,A_j^{}}^1(𝐱,𝐲)`$ $`=`$ $`C_{A_j^{},E_i}^1(𝐲,𝐱)=\delta _{ij}\delta (𝐱𝐲),`$
$`C_{sB_j,B_i}^1(𝐱,𝐲)`$ $`=`$ $`C_{B_i,sB_j}^1(𝐲,𝐱)=\delta _{ij}\delta (𝐱𝐲),`$
$`C_{sB_j^{},B_i^{}}^1(𝐱,𝐲)`$ $`=`$ $`C_{B_i^{},sB_j^{}}^1(𝐲,𝐱)=\delta _{ij}\delta (𝐱𝐲),`$
$`C_{sA_0,A_0}^1(𝐱,𝐲)`$ $`=`$ $`C_{A_0,sA_0}^1(𝐲,𝐱)={\displaystyle \frac{1}{m^2}}\delta (𝐱𝐲),`$
$`C_{sA_0^{},A_0^{}}^1(𝐱,𝐲)`$ $`=`$ $`C_{A_0^{},sA_0^{}}^1(𝐲,𝐱)={\displaystyle \frac{1}{m^2}}\delta (𝐱𝐲),`$
$`C_{A_i,A_0^{}}^1(𝐱,𝐲)`$ $`=`$ $`C_{A_0^{},A_i}^1(𝐲,𝐱)={\displaystyle \frac{1}{m^2}}D_i\delta (𝐱𝐲),`$
$`C_{A_0,A_i^{}}^1(𝐱,𝐲)`$ $`=`$ $`C_{A_i^{},A_0}^1(𝐲,𝐱)={\displaystyle \frac{1}{m^2}}D_i\delta (𝐱𝐲),`$
$`C_{E_i,B_j^{}}^1(𝐱,𝐲)`$ $`=`$ $`C_{B_j^{},E_i}^1(𝐲,𝐱)=\epsilon _{ijk}D_k\delta (𝐱𝐲),`$
$`C_{B_i,E_j^{}}^1(𝐱,𝐲)`$ $`=`$ $`C_{E_j^{},B_i}^1(𝐲,𝐱)=\epsilon _{ijk}D_k\delta (𝐱𝐲).`$ (4.89)
The Dirac bracket for the fields $`A_i,A_i^{},E_i,E_i^{}`$ reads
$`[A_i(𝐱),E_j^{}(𝐲)]_D`$ $`=`$ $`\delta _{ij}\delta (𝐱𝐲),`$
$`[A_i^{}(𝐱),E_j(𝐲)]_D`$ $`=`$ $`\delta _{ij}\delta (𝐱𝐲)`$ (4.90)
and all the other Dirac brackets involving the fields $`A_i,A_i^{},E_i,E_i^{}`$ are zero. The Dirac bracket for the other fields can be derived from these by using the constraints (4.80) and (4.84).
Because we already have identified the true degrees of freedom, there is no need to perform the Maskawa and Nakajima canonical transformation explicitly. The theory is quantized by imposing the following commutation relations for the true degrees of freedom
$`[\widehat{A}_i(𝐱),\widehat{E}_j^{}(𝐲)]`$ $`=`$ $`i[A_i(𝐱),E_j^{}(𝐲)]_D=i\delta _{ij}\delta (𝐱𝐲),`$
$`[\widehat{A}_i^{}(𝐱),\widehat{E}_j(𝐲)]`$ $`=`$ $`i[A_i^{}(𝐱),E_j(𝐲)]_D=i\delta _{ij}\delta (𝐱𝐲).`$ (4.91)
All the other fundamental commutation relations involving the operators $`\widehat{A}_i,\widehat{E}_i^{},\widehat{A}_i^{},\widehat{E}_i`$ and $`A_i^{}(𝐱)`$ are zero. The commutation relations are realized by the representation
$`\widehat{A}_i(𝐱)`$ $`=`$ $`A_i(𝐱),\widehat{E}_j^{}(𝐱)=i{\displaystyle \frac{\delta }{\delta A_j(𝐱)}},`$
$`\widehat{A}_i^{}(𝐱)`$ $`=`$ $`A_i^{}(𝐱),\widehat{E}_j(𝐱)=i{\displaystyle \frac{\delta }{\delta A_j^{}(𝐱)}}.`$ (4.92)
In this representation the Hamiltonian operator reads
$`\widehat{H}`$ $`=`$ $`{\displaystyle d^3yd^3zh_{ij}(𝐲,𝐳)\frac{\delta }{\delta A_i(𝐲)}\frac{\delta }{\delta A_j^{}(𝐳)}}`$ (4.93)
$`+{\displaystyle }d^3x[eV_0(𝐱)(A_i^{}(𝐱){\displaystyle \frac{\delta }{\delta A_i^{}(𝐱)}}A_i(𝐱){\displaystyle \frac{\delta }{\delta A_i(𝐱)}})`$
$`+{\displaystyle \frac{1}{2}}G_{ij}^{}(𝐱)G_{ij}(𝐱)+m^2A_i^{}(𝐱)A_i(𝐱)].`$
The operators act on wavefunctionals $`\mathrm{\Psi }(A_i,A_i^{},t)=A_i,A_i^{}|\mathrm{\Psi }(t)`$. The functional Schrödinger equation for the wavefunctional is given by
$$i\frac{\mathrm{\Psi }(A_i,A_i^{},t)}{t}=\widehat{H}\mathrm{\Psi }(A_i,A_i^{},t).$$
(4.94)
This is the same quantum field theory that can be derived from quantizing the Proca theory. This shows the equivalence of the quantized DKP theory in the spin-1 representation and the quantized Proca theory. The same pilot-wave interpretation could hence be obtained by considering the quantized Proca theory from the outset.
#### 4.4.2 Pilot-wave interpretation
In order to obtain the pilot-wave interpretation we consider the corresponding conservation equation
$$\frac{|\mathrm{\Psi }|^2}{t}+d^3x\left(\frac{\delta J_{A_i}(𝐱)}{\delta A_i(𝐱)}+\frac{\delta J_{A_i^{}}(𝐱)}{\delta A_i^{}(𝐱)}\right)=0,$$
(4.95)
with
$`J_{A_i}(𝐱)`$ $`=`$ $`{\displaystyle \frac{1}{2i}}{\displaystyle d^3yh_{ij}(𝐱,𝐲)\left(\mathrm{\Psi }^{}\frac{\delta }{\delta A_j^{}(𝐲)}\mathrm{\Psi }\mathrm{\Psi }\frac{\delta }{\delta A_j^{}(𝐲)}\mathrm{\Psi }^{}\right)}`$
$`ieV_0(𝐱)|\mathrm{\Psi }|^2A_i(𝐱),`$
$`J_{A_i^{}}(𝐱)`$ $`=`$ $`{\displaystyle \frac{1}{2i}}{\displaystyle d^3yh_{ij}^{}(𝐱,𝐲)\left(\mathrm{\Psi }^{}\frac{\delta }{\delta A_j(𝐲)}\mathrm{\Psi }\mathrm{\Psi }\frac{\delta }{\delta A_j(𝐲)}\mathrm{\Psi }^{}\right)}`$ (4.96)
$`+ieV_0(𝐱)|\mathrm{\Psi }|^2A_i^{}(𝐱).`$
In the pilot-wave interpretation, the conserved density of the field beables $`A_i,A_i^{}`$ is given by $`|\mathrm{\Psi }|^2`$ and the guidance equations for the fields are
$$\dot{A}_i=J_{A_i}/|\mathrm{\Psi }|^2,\dot{A}_i^{}=J_{A_i^{}}/|\mathrm{\Psi }|^2.$$
(4.97)
### 4.5 The electromagnetic field
In this section we consider the pilot-wave interpretation of the quantized the electromagnetic field. We start with reconsidering the quantization of the electromagnetic field. The reason to do is because the two existing pilot-wave approaches, the one originated by Bohm and the one by Valentini , find a natural home in two different ways of quantizing the electromagnetic field. Although these different ways of quantizing the electromagnetic field yield equivalent quantum theories, the corresponding pilot-wave interpretations are not equivalent. We will indicate some problems in the approach by Valentini, which, in our opinion, makes the original approach by Bohm favourable.
This is the organization of the section. First, in Section 4.5.1, we recall the Hamiltonian formulation of Maxwell’s theory for the electromagnetic field. In order not to obscure the issue, we will not start from the Harish-Chandra theory (cf. Section 3.3.2), which is equivalent to Maxwell’s theory at the level of classical field equations, but from the Maxwell form straight away. Most probably it presents no problem to show the equivalence for the quantized theories. The Hamiltonian formulation of Maxwell’s theory can also be found in . In Section 4.5.2 we then consider the quantization of the electromagnetic field by imposing the constraints as operator identities. We will work with the Coulomb gauge and then later in Section 4.5.3 we will discuss some other gauges. This approach of quantizing the electromagnetic field will lead to Bohm’s original pilot-wave interpretation. In Section 4.5.4 we then consider the quantization of the electromagnetic field by imposing constraints as conditions on states. This approach will lead to Valentini’s pilot-wave approach.
#### 4.5.1 Hamiltonian formulation of the electromagnetic field
The free Lagrangian density for the electromagnetic vector potential $`V^\mu =(V_0,𝐕)`$ is given by
$`L_M`$ $`=`$ $`{\displaystyle d^3x_M}={\displaystyle \frac{1}{4}}{\displaystyle d^3xF^{\mu \nu }F_{\mu \nu }}`$ (4.98)
$`=`$ $`{\displaystyle d^3x\left(\frac{1}{2}(_0V_i+_iV_0)(_0V_i+_iV_0)\frac{1}{4}F_{ij}F_{ij}\right)}`$
with $`F^{\mu \nu }=^\mu V^\nu ^\nu V^\mu `$ and $`F_{ij}=_iV_j_jV_i`$. The Lagrangian is invariant under gauge transformations
$$V^\mu V^\mu ^\mu \theta .$$
(4.99)
The equations of motion are
$$_\mu F^{\mu \nu }=0.$$
(4.100)
The canonically conjugate momenta of the fields are
$`\mathrm{\Pi }_{V_0}`$ $`=`$ $`{\displaystyle \frac{\delta L}{\delta \dot{V}_0}}=0,`$
$`\mathrm{\Pi }_{V_i}`$ $`=`$ $`{\displaystyle \frac{\delta L}{\delta \dot{V}_i}}=(_0V_i+_iV_0).`$ (4.101)
Because $`\mathrm{\Pi }_{V_0}=0`$, we have a primary constraint
$$\chi _1=\mathrm{\Pi }_{V_0}.$$
(4.102)
The canonical Hamiltonian reads
$$H_C=d^3x\left(\frac{1}{2}\mathrm{\Pi }_{V_i}\mathrm{\Pi }_{V_i}+\frac{1}{4}F_{ij}F_{ij}\mathrm{\Pi }_{V_i}_iV_0\right).$$
(4.103)
The total Hamiltonian reads
$$H_T=H_C+d^3xu_1\chi _1.$$
(4.104)
The consistency requirement that the primary constraint $`\chi _1`$ is conserved in time leads to the secondary constraint
$$\chi _2=_i\mathrm{\Pi }_{V_i},$$
(4.105)
which is called the Gauss constraint. There are no further constraints and the field $`u_1`$ remains undetermined. Hence both constraints are first class constraints. The constraint $`\chi _1`$ should be included in the Hamiltonian to yield the extended Hamiltonian
$$H_E=H_C+d^3xu_1\chi _1+d^3xu_2\chi _2.$$
(4.106)
In the next section, we will apply both approaches of dealing with first class constraints, as described in Section 4.2.3. Let us first identify the true degrees of freedom. In order to do so we make use of the following property. If $`F_i`$ is a vector field (which, together its its spatial derivatives, vanishes sufficiently fast at spatial infinity), one can uniquely decompose it into a transversal part and a longitudinal part:
$$F_i=F_i^T+F_i^L,$$
(4.107)
with
$`F_i^T`$ $`=`$ $`\left(\delta _{ij}{\displaystyle \frac{_i_j}{^2}}\right)F_j,`$
$`F_i^L`$ $`=`$ $`{\displaystyle \frac{_i_j}{^2}}F_j`$ (4.108)
and where the operator $`^2`$ acts as
$$\frac{1}{^2}f(𝐱)=d^3y\frac{f(𝐲)}{4\pi |𝐱𝐲|}.$$
(4.109)
By using this decomposition for the fields $`V_i`$ and $`\mathrm{\Pi }_{V_i}`$, one sees that transversal components $`V_i^T`$ and $`\mathrm{\Pi }_{V_i}^T`$ live on the constraint space $`\mathrm{\Pi }_{V_0}=_i\mathrm{\Pi }_{V_i}=0`$. One can also check that $`V_i^T`$ and $`\mathrm{\Pi }_{V_i}^T`$ have Poisson brackets zero with the constraints $`\chi _1`$ and $`\chi _2`$. Because $`\chi _1`$ and $`\chi _2`$ were the generators of gauge transformations this means that the transversal fields $`V_i^T`$ and $`\mathrm{\Pi }_{V_i}^T`$ are gauge invariant and hence they represent the true degrees of freedom. We can do a counting to see that the transversal fields represent all the true degrees of freedom. We have 8 canonical variables at each point in space and 4 constraints. That leaves us with 4 true degrees of freedom at each point in space and this is exactly the number of degrees of freedom of the transversal fields. The Hamiltonian $`H`$ on the constraint space can also be expressed in terms of $`V_i^T`$ and $`\mathrm{\Pi }_{V_i}^T`$:
$$H=\frac{1}{2}d^3x\left(\mathrm{\Pi }_{V_i}^T\mathrm{\Pi }_{V_i}^TV_i^T^2V_i^T\right).$$
(4.110)
#### 4.5.2 Constraints as operator identities: The Coulomb gauge
##### Commutation relations in the Coulomb gauge
The two first class constraints $`\chi _1`$ and $`\chi _2`$ are generators of infinitesimal gauge transformations. In particular, the infinitesimal gauge transformations $`V^\mu V^\mu +\delta V^\mu =V^\mu ^\mu \epsilon `$, which are symmetry transformations of the Lagrangian equations of motion (4.100), are generated by the linear combination $`d^3y\left(\chi _1(y)\epsilon _1(y)+\chi _2(y)\epsilon _2(y)\right)`$ with $`\epsilon _1=_0\epsilon `$ and $`\epsilon _2=\epsilon `$, i.e.
$$\delta V^\mu (x)=[V^\mu (x),d^3y\left(\chi _1(y)\epsilon _1(y)+\chi _2(y)\epsilon _2(y)\right)]_P=^\mu \epsilon (x),$$
(4.111)
where the fields are taken at equal time, i.e. $`x_0=y_0`$. In fact there still exist gauge transformations generated by the constraints, which are not gauge symmetries of the Lagrangian equations of motion, see \[110, p. 134\]. This indicates that the set of infinitesimal gauge transformations of the Hamiltonian equations of motion (which are generated by the first class constraints) does not necessarily correspond to the set of infinitesimal gauge transformation of the Lagrangian equations of motion.
As indicated in Section 4.2.3 we can impose additional constraints, i.e. gauge constraints, so that the full set of constraints becomes second class. A suitable set of constraints is given by the Coulomb gauge
$`\chi _3`$ $`=`$ $`_iV_i,`$ (4.112)
$`\chi _4`$ $`=`$ $`V_0.`$ (4.113)
The Coulomb gauge satisfies the requirements of an admissible gauge (cf. Section 4.2.3). First, the Coulomb gauge can be attained with the gauge transformation (4.99) with
$$\theta =\frac{1}{^2}_iV_i$$
(4.114)
Second, if we restrict ourself to gauge transformations for which the function $`\theta `$ vanishes at spatial infinity, then the Coulomb gauge fixes the gauge uniquely.
The quantization of the electromagnetic field in the Coulomb gauge, by imposing constraints as operator equations, is a textbook example of dealing with constraints. The treatment can for example be found in \[110, pp. 123-140\] and \[112, pp. 339-350\]. Here we repeat only the basic elements. In Section 4.5.3, we will consider quantization in other gauges.
The Dirac bracket can be calculated by using the inverse of the matrix $`C_{NM}(𝐱,𝐲)=[\chi _N(𝐱),\chi _M(𝐲)]_P`$, with $`M,N=1,\mathrm{},4`$. The matrix $`C`$ has the following nonzero components
$`C_{23}(𝐱,𝐲)`$ $`=`$ $`C_{32}(𝐱,𝐲)=^2\delta (𝐱𝐲),`$
$`C_{14}(𝐱,𝐲)`$ $`=`$ $`C_{41}(𝐱,𝐲)=\delta (𝐱𝐲).`$ (4.115)
One can construct an inverse matrix $`C_{NM}^1(𝐱,𝐲)`$, which has the following nonzero components
$`C_{14}^1(𝐱,𝐲)`$ $`=`$ $`C_{41}^1(𝐱,𝐲)=\delta (𝐱𝐲),`$
$`C_{23}^1(𝐱,𝐲)`$ $`=`$ $`C_{32}^1(𝐱,𝐲)={\displaystyle \frac{1}{^2}}\delta (𝐱𝐲).`$ (4.116)
Note that this inverse is not unique. There is an ambiguity in the matrix elements $`C_{23}^1(𝐱,𝐲)`$ and $`C_{32}^1(𝐱,𝐲)`$. They are both determined up to a function $`g(𝐱,𝐲)`$ which satisfies $`_x^2g(𝐱,𝐲)=_y^2g(𝐱,𝐲)=0`$. This ambiguity in the matrix $`C^1`$ may lead to an ambiguity in the Dirac bracket and hence in the field commutators. However, the ambiguity for the inverse matrix $`C^1`$ can be removed by considering the boundary conditions for the fields (they vanish sufficiently fast at spatial infinity). We shall not do this analysis here, but we refer the reader to \[110, pp. 65-72\] and where the same issue is treated in the context of light-cone quantization.
Note that there was no such an ambiguity in the Duffin-Kemmer-Petiau theory. In the Duffin-Kemmer-Petiau theory, the inverse of the matrix $`C`$ was always uniquely determined.
Using this inverse matrix $`C^1`$ we obtain the following Dirac bracket for the fields
$`[V_i(𝐱),\mathrm{\Pi }_{V_j}(𝐲)]_D`$ $`=`$ $`\left(\delta _{ij}{\displaystyle \frac{_i_j}{^2}}\right)\delta (𝐱𝐲),`$ (4.117)
$`[V_i(𝐱),V_j(𝐲)]_D`$ $`=`$ $`[\mathrm{\Pi }_{V_i}(𝐱),\mathrm{\Pi }_{V_j}(𝐲)]_D=0.`$ (4.118)
The Dirac brackets involving the fields $`V_0`$ and $`\mathrm{\Pi }_{V_0}`$ are zero.
Quantization proceeds by imposing the following commutation relations for the operators
$`[\widehat{V}_i(𝐱),\widehat{\mathrm{\Pi }}_{V_j}(𝐲)]`$ $`=`$ $`i\left(\delta _{ij}{\displaystyle \frac{_i_j}{^2}}\right)\delta (𝐱𝐲),`$ (4.119)
$`[\widehat{V}_i(𝐱),\widehat{V}_j(𝐲)]`$ $`=`$ $`[\widehat{\mathrm{\Pi }}_{V_i}(𝐱),\widehat{\mathrm{\Pi }}_{V_j}(𝐲)]=0.`$ (4.120)
The commutation relations involving the operators $`\widehat{V}_0`$ and $`\widehat{\mathrm{\Pi }}_{V_0}`$ are zero.
The rest of this section is divided in three parts:
* First, we consider in detail the construction of a class of Maskawa and Nakajima canonical transformations which enable us to separate the true degrees of freedom from the constraints.
* Second, we consider the functional Schrödinger equation in terms of the true degrees of freedom, together with the pilot-wave interpretation.
* Finally, we consider some explicit examples of Maskawa and Nakajima canonical transformations.
##### The true degrees of freedom
We have to find a representation for the operators $`\widehat{V}_0,\widehat{V}_i,\widehat{\mathrm{\Pi }}_{V_0}`$ and $`\widehat{\mathrm{\Pi }}_{V_i}`$ so that the constraints and the commutation relations are satisfied. The operators $`\widehat{V}_0`$ and $`\widehat{\mathrm{\Pi }}_{V_0}`$ are zero as constraints, so we only have to consider a representation for the operators $`\widehat{V}_i`$ and $`\widehat{\mathrm{\Pi }}_{V_i}`$.
Inspired by the decomposition (4.107) into longitudinal and transversal part of fields, it would be tempting to use the representation
$`\widehat{V}_i(𝐱)`$ $`=`$ $`\left(\delta _{ij}{\displaystyle \frac{_i_j}{^2}}\right)V_j(𝐱),`$
$`\widehat{\mathrm{\Pi }}_{V_i}(𝐱)`$ $`=`$ $`i\left(\delta _{ij}{\displaystyle \frac{_i_j}{^2}}\right){\displaystyle \frac{\delta }{\delta V_j(𝐱)}},`$ (4.121)
where $`V_j(𝐱)`$ is a real-valued three component field. This representation satisfies the constraints and the commutation relations (4.120). However, with this representation we would introduce superfluous degrees of freedom. This is because this representation keeps 6 degrees of freedom at each point in space, namely $`V_i(𝐱)`$ and $`\delta /\delta V_i(𝐱)`$. This means that the representation has two degrees of freedom in excess. The reason for the superfluous degrees of freedom is that the representation (4.121) is invariant under the transformations $`V_i(𝐱)V_i(𝐱)+_i\theta (𝐱)`$ and is therefore not one-to-one.
In order to construct a correct representation, we will first explicitly perform the Maskawa-Nakajima (MN) canonical transformation. This canonical transformation will separate the true degrees of freedom from the constraints. Then we can use the standard representation for the operators corresponding to the true degrees of freedom. In order to find the MN canonical transformation we will work constructively. We will also try to keep some generality. In this way we will in fact end up with a class of possible MN canonical transformations.
Because the constraints $`\chi _1=\mathrm{\Pi }_{V_0}`$ and $`\chi _4=V_0`$ already form a canonical pair they should not be involved in the MN canonical transformation of the other canonical variables $`V_i`$ and $`\mathrm{\Pi }_{V_i}`$. Because the Coulomb gauge is linear in the fields $`V_i`$, it is most simple to try to find a canonical transformation which is also linear in the fields. With $`\stackrel{~}{V}_i`$ the new field variables, with corresponding momenta $`\mathrm{\Pi }_{\stackrel{~}{V}_i}`$, we then look for a transformation of the form
$`V_i(𝐱)`$ $`=`$ $`{\displaystyle d^3yK_{ij}(𝐱,𝐲)\stackrel{~}{V}_j(𝐲)},`$
$`\stackrel{~}{V}_i(𝐱)`$ $`=`$ $`{\displaystyle d^3yK_{ij}^1(𝐱,𝐲)V_j(𝐲)},`$ (4.122)
with $`K`$ a matrix with inverse $`K^1`$
$`{\displaystyle d^3xK_{ki}^1(𝐳,𝐱)K_{ij}(𝐱,𝐲)}=\delta _{kj}\delta (𝐳𝐲),`$ (4.123)
$`{\displaystyle d^3xK_{ki}(𝐳,𝐱)K_{ij}^1(𝐱,𝐲)}=\delta _{kj}\delta (𝐳𝐲).`$ (4.124)
We further need a transformation of the momenta $`\mathrm{\Pi }_{V_i}`$. For simplicity we construct a canonical transformation with corresponding generating function $`F=0`$. With $`F=0`$, it follows from (4.29) that
$$\mathrm{\Pi }_{\stackrel{~}{V}_i}(𝐱)=d^3y\mathrm{\Pi }_{V_j}(𝐲)\frac{\delta V_j(𝐲)}{\delta \stackrel{~}{V}_i(𝐱)}=d^3y\mathrm{\Pi }_{V_j}(𝐲)K_{ji}(𝐲,𝐱).$$
(4.125)
The inverse transformation reads
$$\mathrm{\Pi }_{V_i}(𝐱)=d^3y\mathrm{\Pi }_{\stackrel{~}{V}_j}(𝐲)K_{ji}^1(𝐲,𝐱).$$
(4.126)
In this way, the transformation determined by (4.122), (4.125) and (4.126) is a canonical transformation (and even a point transformation).
Now we further need to assure the the canonical transformation is a MN canonical transformation. We will construct a canonical transformation such that the remaining constraints $`_iV_i=_i\mathrm{\Pi }_{V_i}=0`$ read $`\stackrel{~}{V}_3=\mathrm{\Pi }_{\stackrel{~}{V}_3}=0`$ in terms of the new variables. The true degrees of freedom then will be $`\stackrel{~}{V}_1,\stackrel{~}{V}_2,\mathrm{\Pi }_{\stackrel{~}{V}_1}`$ and $`\mathrm{\Pi }_{\stackrel{~}{V}_2}`$. Therefore we further require that the transformations $`K`$ and $`K^1`$ satisfy
$`_{x_i}K_{ij}(𝐱,𝐲)`$ $`=`$ $`0,\text{for }j=1,2,`$ (4.127)
$`\epsilon _{ilk}_{x_l}K_{kj}(𝐱,𝐲)`$ $`=`$ $`0,\text{for }j=3,`$ (4.128)
$`_{x_i}K_{ji}^1(𝐲,𝐱)`$ $`=`$ $`0,\text{for }j=1,2,`$ (4.129)
$`\epsilon _{ilk}_{x_l}K_{jk}^1(𝐲,𝐱)`$ $`=`$ $`0,\text{for }j=3.`$ (4.130)
This implies that the transversal parts of the canonical variables $`V_i`$ and $`\mathrm{\Pi }_{V_i}`$ can be written respectively in terms of $`\stackrel{~}{V}_1`$ and $`\stackrel{~}{V}_2`$, and $`\mathrm{\Pi }_{\stackrel{~}{V}_1}`$ and $`\mathrm{\Pi }_{\stackrel{~}{V}_2}`$. The longitudinal parts of the canonical variables $`V_i`$ and $`\mathrm{\Pi }_{V_i}`$ can be written respectively in terms of $`\stackrel{~}{V}_3`$ and $`\mathrm{\Pi }_{\stackrel{~}{V}_3}`$. I.e. we have
$`V_i^T(𝐱)`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{2}{}}}{\displaystyle d^3yK_{ij}(𝐱,𝐲)\stackrel{~}{V}_j(𝐲)},`$
$`V_i^L(𝐱)`$ $`=`$ $`{\displaystyle d^3yK_{i3}(𝐱,𝐲)\stackrel{~}{V}_3(𝐲)},`$
$`\mathrm{\Pi }_{V_i}^T(𝐱)`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{2}{}}}{\displaystyle d^3y\mathrm{\Pi }_{\stackrel{~}{V}_j}(𝐲)K_{ji}^1(𝐲,𝐱)},`$
$`\mathrm{\Pi }_{V_i}^L(𝐱)`$ $`=`$ $`{\displaystyle d^3y\mathrm{\Pi }_{\stackrel{~}{V}_3}(𝐲)K_{3i}^1(𝐲,𝐱)}.`$ (4.131)
Because $`K_{i3}`$ and $`K_{3i}^1`$ are irrotational, cf. (4.128) and (4.130), we can write
$`K_{i3}(𝐱,𝐲)`$ $`=`$ $`_{x_i}U(𝐱,𝐲),`$ (4.132)
$`K_{3i}^1(𝐲,𝐱)`$ $`=`$ $`_{x_i}\overline{U}(𝐲,𝐱).`$ (4.133)
From the expression (4.123) for $`j=k=3`$ it then follows that
$$d^3xK_{3i}^1(𝐳,𝐱)K_{i3}(𝐱,𝐲)=d^3x\left(_{x_i}\overline{U}(𝐳,𝐱)\right)\left(_{x_i}U(𝐱,𝐲)\right)=\delta (𝐳𝐲).$$
(4.134)
Because $`U`$ and $`\overline{U}`$ will make no appearance in the theory when the constraints are imposed as operator identities (because $`U`$ and $`\overline{U}`$ are only present in the longitudinal components of the fields), we can make some more assumptions on these matrices. We will assume that
$$\overline{U}(𝐳,𝐱)=\frac{1}{_x^2}U^1(𝐳,𝐱),$$
(4.135)
with $`U^1`$ the inverse of $`U`$, and that the boundary terms in (4.134) vanish after partial integration. Under these circumstances, the relation (4.134) is then satisfied.
From
$$_iV_i(𝐱)=_iV_i^L(𝐱)=d^3y_x^2U(𝐱,𝐲)\stackrel{~}{V}_3(𝐲)$$
(4.136)
it follows that $`\stackrel{~}{V}_3=0_iV_i=0`$. On the other hand, from
$$\stackrel{~}{V}_3(𝐱)=d^3y\frac{_{y_i}}{_y^2}U^1(𝐱,𝐲)V_i(𝐲)=d^3y\frac{1}{_y^2}U^1(𝐱,𝐲)_{y_i}V_i(𝐲)$$
(4.137)
(for the last equality, we used the fact that the fields $`V_i`$ vanish sufficiently fast at spatial infinity), it follows that $`_iV_i=0\stackrel{~}{V}_3=0`$. Similarly, we have that $`_i\mathrm{\Pi }_{V_i}=0\mathrm{\Pi }_{\stackrel{~}{V}_3}=0`$. Hence, in terms of the new variables the constraints read $`\stackrel{~}{V}_3=\mathrm{\Pi }_{\stackrel{~}{V}_3}=0`$. The true degrees of freedom are $`\stackrel{~}{V}_1,\stackrel{~}{V}_2,\mathrm{\Pi }_{\stackrel{~}{V}_1}`$ and $`\mathrm{\Pi }_{\stackrel{~}{V}_2}`$. As mentioned before the true degrees of freedom are gauge independent degrees of freedom. In this case this follows from the fact that $`\stackrel{~}{V}_1,\stackrel{~}{V}_2,\mathrm{\Pi }_{\stackrel{~}{V}_1}`$ and $`\mathrm{\Pi }_{\stackrel{~}{V}_2}`$ have zero Poisson brackets with $`\mathrm{\Pi }_{\stackrel{~}{V}_3}`$ and $`\mathrm{\Pi }_{\stackrel{~}{V}_3}`$ is the generator of gauge transformations in the new coordinates (recall that the first class constraint $`_i\mathrm{\Pi }_{V_i}=0`$ reads $`\mathrm{\Pi }_{\stackrel{~}{V}_3}=0`$ in the new coordinates).
We can conclude that the transformation determined by (4.122), (4.125) and (4.126) satisfies the MN theorem.
One can also explicitly check that the Dirac brackets (4.120) of the canonical variables $`V_i(𝐱)`$ and $`\mathrm{\Pi }_{V_j}(𝐲)`$ equal the Poisson brackets restricted to the unconstrained variables $`\stackrel{~}{V}_i`$ and $`\mathrm{\Pi }_{\stackrel{~}{V}_i}`$ with $`i=1,2`$. This is done by using the relation
$`{\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle d^3xK_{ki}(𝐳,𝐱)K_{ij}^1(𝐱,𝐲)}`$ $`=`$ $`\delta _{kj}\delta (𝐳𝐲){\displaystyle d^3xK_{k3}(𝐳,𝐱)K_{3j}^1(𝐱,𝐲)}`$ (4.138)
$`=`$ $`\left(\delta _{kj}{\displaystyle \frac{_{y_k}_{y_j}}{^2}}\right)\delta (𝐳𝐲),`$
which is found from (4.123), (4.132), (4.133) and (4.135).
##### Functional Schrödinger equation and pilot-wave interpretation
By applying the canonical transformation given by (4.122) and (4.126) to the Hamiltonian (4.106), and by using the constraints $`\stackrel{~}{V}_3=0`$ and $`\mathrm{\Pi }_{\stackrel{~}{V}_3}=0`$, or by directly applying the canonical transformation to the Hamiltonian (4.110), and by using the fact that the generating function is zero, we obtain the following Hamiltonian for the unconstrained variables $`\stackrel{~}{V}_i`$ and $`\mathrm{\Pi }_{\stackrel{~}{V}_i}`$ $`(i=1,2)`$:
$$H=\underset{k,l=1}{\overset{2}{}}d^3yd^3z(h_{kl}(𝐲,𝐳)\mathrm{\Pi }_{\stackrel{~}{V}_k}(𝐲)\mathrm{\Pi }_{\stackrel{~}{V}_l}(𝐳)+\overline{h}_{kl}(𝐲,𝐳)\stackrel{~}{V}_k(𝐲)\stackrel{~}{V}_l(𝐳)),$$
(4.139)
with
$`h_{kl}(𝐲,𝐳)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{3}{}}}{\displaystyle d^3xK_{ki}^1(𝐲,𝐱)K_{li}^1(𝐳,𝐱)},`$ (4.140)
$`\overline{h}_{kl}(𝐲,𝐳)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{3}{}}}{\displaystyle d^3x\epsilon _{imn}_{x_m}K_{nk}(𝐱,𝐲)\epsilon _{irs}_{x_r}K_{sl}(𝐱,𝐳)}.`$ (4.141)
Because the fields $`\stackrel{~}{V}_k`$ and $`\mathrm{\Pi }_{\stackrel{~}{V}_k}`$ $`(k=1,2)`$ are unconstrained, the theory is quantized by using the standard commutation relation for the corresponding operators and hence we can use the standard representation
$$\widehat{\stackrel{~}{V}}_k(𝐱)=\stackrel{~}{V}_k(𝐱),\widehat{\mathrm{\Pi }}_{\stackrel{~}{V}_k}(𝐱)=i\frac{\delta }{\delta \stackrel{~}{V}_k(𝐱)},\text{for }k=1,2.$$
(4.142)
In this representation, the functional Schrödinger equation for the wavefunctional $`\mathrm{\Psi }(\stackrel{~}{V}_1,\stackrel{~}{V}_2,t)=\stackrel{~}{V}_1,\stackrel{~}{V}_2|\mathrm{\Psi }(t)`$ reads
$`i{\displaystyle \frac{\mathrm{\Psi }(\stackrel{~}{V}_1,\stackrel{~}{V}_2,t)}{t}}`$ $`=`$ $`{\displaystyle \underset{k,l=1}{\overset{2}{}}}{\displaystyle }d^3yd^3z(h_{kl}(𝐲,𝐳){\displaystyle \frac{\delta }{\delta \stackrel{~}{V}_k(𝐲)}}{\displaystyle \frac{\delta }{\delta \stackrel{~}{V}_l(𝐳)}}+`$ (4.143)
$`\overline{h}_{kl}(𝐲,𝐳)\stackrel{~}{V}_k(𝐲)\stackrel{~}{V}_l(𝐳))\mathrm{\Psi }(\stackrel{~}{V}_1,\stackrel{~}{V}_2,t).`$
The corresponding continuity equation for the density $`|\mathrm{\Psi }(\stackrel{~}{V}_1,\stackrel{~}{V}_2,t)|^2`$ reads
$$\frac{|\mathrm{\Psi }|^2}{t}+\underset{k=1}{\overset{2}{}}d^3x\frac{\delta J_{\stackrel{~}{V}_k}(𝐱)}{\delta \stackrel{~}{V}_k(𝐱)}=0,$$
(4.144)
with
$$J_{\stackrel{~}{V}_k}(𝐱)=\frac{1}{2i}\underset{l=1}{\overset{2}{}}d^3yh_{kl}(𝐱,𝐲)\left(\mathrm{\Psi }^{}\frac{\delta }{\delta \stackrel{~}{V}_l(𝐲)}\mathrm{\Psi }\mathrm{\Psi }\frac{\delta }{\delta \stackrel{~}{V}_l(𝐲)}\mathrm{\Psi }^{}\right).$$
(4.145)
The pilot-wave interpretation is straightforward. The conserved density of the field beables $`\stackrel{~}{V}_1`$ and $`\stackrel{~}{V}_2`$ is given by $`|\mathrm{\Psi }|^2`$ and the guidance equations for the fields are
$$\dot{\stackrel{~}{V}}_k(𝐱)=J_{\stackrel{~}{V}_k}(𝐱)/|\mathrm{\Psi }|^2,\text{for }k=1,2.$$
(4.146)
It is important to note that the fields $`\stackrel{~}{V}_k(𝐱)`$ ($`k=1,2`$) do not necessary live in physical 3-space,<sup>12</sup><sup>12</sup>12This could perhaps be indicated more explicitly in the notation, e.g. by replacing the variable $`𝐱`$ by some other variable. but in some abstract space defined by the transformation (4.122). The transversal part of the electromagnetic field potential which lives in physical 3-space can always be obtained by considering (4.131). For the quantized Duffin-Kemmer-Petiau field, it was no problem to find true degrees of freedom which live in physical 3-space. This will be difficult to do for the electromagnetic field. In the following paragraph we consider some examples of possible representations. In the first example the fields $`\stackrel{~}{V}_k(𝐱)`$ ($`k=1,2`$) will live in physical 3-space, in the second example not.
##### Explicit examples of Maskawa and Nakajima canonical transformations
Example one: Assume the transformation matrix
$$K_{ij}(𝐱,𝐲)=\{\begin{array}{cc}\delta _{ij}\delta (𝐱𝐲)\frac{1}{2}\delta _{i3}_{x_j}\delta (x_1y_1)\delta (x_2y_2)\text{sgn}(x_3y_3)\hfill & \text{if }j=1,2\hfill \\ _{x_i}U(𝐱,𝐲)\hfill & \text{if }j=3\hfill \end{array}$$
(4.147)
where $`U`$ is an arbitrary non-singular matrix which has the properties discussed in the previous section (cf. the paragraph containing equation (4.134)) and ‘sgn’ the sign function. Note that $`K`$ satisfies the requirements (4.127) and (4.128). The inverse of $`K`$ is given by
$$K_{ij}^1(𝐱,𝐲)=\{\begin{array}{cc}\left(\delta _{ij}\frac{_{x_i}_{x_j}}{^2}\right)\delta (𝐱𝐲)\hfill & \text{if }i=1,2\hfill \\ \frac{_{y_j}}{_y^2}U^1(𝐱,𝐲)\hfill & \text{if }i=3\hfill \end{array}.$$
(4.148)
With this transformation the transversal part of $`V_i`$ reads
$`V_1^T(𝐱)`$ $`=`$ $`\stackrel{~}{V}_1(𝐱),`$ (4.149)
$`V_2^T(𝐱)`$ $`=`$ $`\stackrel{~}{V}_2(𝐱),`$ (4.150)
$`V_3^T(𝐱)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left({\displaystyle _{x_3}^+\mathrm{}}{\displaystyle _{\mathrm{}}^{x_3}}\right)\left(ds{\displaystyle \underset{i=1}{\overset{2}{}}}_{x_i}\stackrel{~}{V}_i(x_1,x_2,s)\right).`$ (4.151)
The transversal momentum field reads
$$\mathrm{\Pi }_{V_i}^T(𝐱)=\underset{k=1}{\overset{2}{}}\left(\delta _{ik}\frac{_{x_i}_{x_k}}{^2}\right)\mathrm{\Pi }_{\stackrel{~}{V}_k}(𝐱).$$
(4.152)
The expression of the fields $`V_i`$ in terms of unconstrained variables given by (4.149)-(4.151) was suggested by Weinberg . As shown in the preceding section we can write the Schrödinger equation solely in terms of the unconstrained variables $`\stackrel{~}{V}_1`$ and $`\stackrel{~}{V}_2`$ with the representation (4.142) and we can construct the corresponding pilot-wave interpretation. A pleasant feature is that the unconstrained variables $`\stackrel{~}{V}_1`$ and $`\stackrel{~}{V}_2`$ live in physical 3-space. But the Hamiltonian will display a highly nonlocal dependence on the unconstrained variables because the field $`V_3^T(𝐱)`$ depends nonlocally on $`\stackrel{~}{V}_1`$ and $`\stackrel{~}{V}_2`$.
However, there is a problem with this transformation which is left unmentioned by Weinberg and which prevents us from using it. Because
$$\underset{x_3\pm \mathrm{}}{lim}V_3^T(𝐱)=\frac{1}{2}_{\mathrm{}}^+\mathrm{}𝑑s\left(\underset{k=1}{\overset{2}{}}_{x_k}\stackrel{~}{V}_k(x_1,x_2,s)\right)$$
(4.153)
is not zero (unless further constraints are brought into play), our assumption that the fields vanish at infinity are not met. This in itself is not a problem, we could do with different boundary conditions. However, because $`V_3^T(𝐱)`$ does not vanish at spatial infinity there appears an explicit infinity in the Hamiltonian (4.139) and hence in the equations of motion, which is intolerable. The infinity appears explicitly because for $`k,l=1,2`$, $`\overline{h}_{kl}(𝐲,𝐳)`$ contains the term $`\frac{1}{2}_{i=1}^2d^3x_{x_i}K_{3k}(𝐱,𝐲)_{x_i}K_{3l}(𝐱,𝐳)`$ for which
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle d^3x_{x_i}K_{3k}(𝐱,𝐲)_{x_i}K_{3l}(𝐱,𝐳)}`$ $``$ $`{\displaystyle 𝑑x_3\text{sgn}(x_3y_3)\text{sgn}(x_3z_3)}`$ (4.154)
$`=`$ $`\left({\displaystyle _{\mathrm{}}^+\mathrm{}}2{\displaystyle _{\text{min}(y_3,z_3)}^{\text{max}(y_3,z_3)}}\right)dx_3`$
$`=`$ $`+\mathrm{}.`$
This problem can not be solved by merely adding a function independent of $`x_3`$ to the right hand side of (4.151). It can only be solved by adding further constraints or by introducing a cut-off in the Hamiltonian.
Example two: Let us first introduce two 3-vectors $`𝜺^j(𝐤)`$, $`j=1,2`$, for each 3-vector $`𝐤`$, such that the vectors $`𝜺^1(𝐤),𝜺^2(𝐤),𝐤/k`$ forms a orthonormal triad. In other words the following conditions should be satisfied: orthogonality
$$\underset{m=1}{\overset{3}{}}\epsilon _m^j(𝐤)\epsilon _m^i(𝐤)=\delta _{ji},i,j=1,2,$$
(4.155)
the transversality condition
$$𝐤𝜺^j(𝐤)=0$$
(4.156)
and the completeness relation
$$\underset{j=1}{\overset{2}{}}\epsilon _m^j(𝐤)\epsilon _n^j(𝐤)=\delta _{mn}\frac{k_mk_n}{k^2}.$$
(4.157)
We also demand that
$$𝜺^j(𝐤)=𝜺^j(𝐤).$$
(4.158)
For example, we could take the following choice for the vectors $`𝜺^1(𝐤)`$ and $`𝜺^2(𝐤)`$. For each $`k_30`$ we take
$$𝜺^j(𝐤)=R(𝐤/k,𝐞_1)\left(\begin{array}{c}0\\ \frac{1}{\sqrt{2}}\\ (1)^j\frac{1}{\sqrt{2}}\end{array}\right)\text{for }k_30,$$
(4.159)
where $`R(𝐤/k,𝐞_1)`$ is the rotation matrix that carries the unit vector $`𝐞_1=(1,0,0)^T`$ to the unit vector $`𝐤/k`$. For $`k_3<0`$ the vectors are defined by imposing the condition $`𝜺^j(𝐤)=𝜺^j(𝐤)`$.
We are now ready to define the transformation matrix as
$$K_{ij}(𝐱,𝐲)=\{\begin{array}{cc}\frac{1}{(2\pi )^3}d^3ke^{i𝐤(𝐱𝐲)}\epsilon _i^j(𝐤)\hfill & \text{if }j=1,2\hfill \\ _{x_i}U(𝐱,𝐲)\hfill & \text{if }j=3\hfill \end{array}$$
(4.160)
where $`U`$ is an arbitrary non-singular matrix which has the properties discussed in the previous section (cf. the paragraph containing equation (4.134)). Note that $`K`$ is a real matrix because of (4.158) and that $`K`$ satisfies the requirements (4.127) and (4.128). The inverse of $`K`$ is given by
$$K_{ij}^1(𝐱,𝐲)=\{\begin{array}{cc}\frac{1}{(2\pi )^3}d^3ke^{i𝐤(𝐱𝐲)}\epsilon _j^i(𝐤)\hfill & \text{if }i=1,2\hfill \\ \frac{_{y_j}}{_y^2}U^1(𝐱,𝐲)\hfill & \text{if }i=3\hfill \end{array}.$$
(4.161)
In this representation we have the pleasant feature that
$`h_{kl}(𝐲,𝐳)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\delta (𝐲𝐳),`$
$`\overline{h}_{kl}(𝐲,𝐳)`$ $`=`$ $`{\displaystyle \frac{1}{2}}^2(𝐲𝐳).`$ (4.162)
Hence the Schrödinger equation (4.143) reads
$$i\frac{\mathrm{\Psi }(\stackrel{~}{V}_1,\stackrel{~}{V}_2,t)}{t}=\frac{1}{2}\underset{l=1}{\overset{2}{}}d^3x\left(\frac{\delta ^2}{\delta \stackrel{~}{V}_l(𝐱)^2}\stackrel{~}{V}_l(𝐱)^2\stackrel{~}{V}_l(𝐱)\right)\mathrm{\Psi }(\stackrel{~}{V}_1,\stackrel{~}{V}_2,t).$$
(4.163)
This is the same Schrödinger equation that would be obtained for two noninteracting massless uncharged spin-0 particles.<sup>13</sup><sup>13</sup>13Although we did not treat the massless spin-0 case explicitly, it can be obtained form the massive spin-0 theory (which was considered in Section 4.3) simply by putting the mass $`m`$ equal to zero. Note that the massless spin-1 case cannot be obtained by doing this in the Proca theory. The only important difference is that here, the fields $`\stackrel{~}{V}_k(𝐱)`$ ($`k=1,2`$) do not live in physical 3-space. The corresponding guidance equations (4.146) are
$$\dot{\stackrel{~}{V}}_l=\frac{1}{2i|\mathrm{\Psi }|^2}\left(\mathrm{\Psi }^{}\frac{\delta }{\delta \stackrel{~}{V}_l}\mathrm{\Psi }\mathrm{\Psi }\frac{\delta }{\delta \stackrel{~}{V}_l}\mathrm{\Psi }^{}\right),\text{with }l=1,2.$$
(4.164)
We can rewrite the Schrödinger equation and the corresponding pilot-wave interpretation by considering the Fourier expansion
$$\stackrel{~}{V}_l(𝐱)=\frac{1}{\sqrt{(2\pi )^3}}d^3kq_l(𝐤)e^{i𝐤𝐱},\text{with }l=1,2$$
(4.165)
and with $`q_l(𝐤)=q_l^{}(𝐤)`$, because the fields $`\stackrel{~}{V}_l`$ are real. In terms of the Fourier modes $`q_l(𝐤)`$, the Schrödinger equation (4.163) becomes the Schrödinger equation
$$i\frac{\mathrm{\Phi }(q_1,q_2,t)}{t}=\frac{1}{2}\underset{l=1}{\overset{2}{}}d^3k\left(\frac{\delta ^2}{\delta q_l(𝐤)\delta q_l^{}(𝐤)}+k^2q_l(𝐤)q_l^{}(𝐤)\right)\mathrm{\Phi }(q_1,q_2,t)$$
(4.166)
for the wavefunctional $`\mathrm{\Phi }(q_l,t)`$. If $`\mathrm{\Psi }(\stackrel{~}{V}_1,\stackrel{~}{V}_2,t)`$ is a solution of the Schrödinger equation (4.163) then we can construct the following solution
$$\mathrm{\Phi }(q_l,t)\mathrm{\Psi }((2\pi )^{3/2}d^3kq_l(𝐤)e^{i𝐤𝐱},t)$$
(4.167)
for the Schrödinger equation (4.166).
We can also write $`V_i^T(𝐱)`$, the transversal component of the vector-potential, directly in terms of the modes $`q_l(𝐤)`$:
$$V_i^T(𝐱)=\frac{1}{\sqrt{(2\pi )^3}}\underset{l=1}{\overset{2}{}}d^3kq_l(𝐤)\epsilon _i^l(𝐤)e^{i𝐤𝐱}.$$
(4.168)
This relation could of course be used to construct a Maskawa and Nakajima canonical transformation which directly yields us the true degrees of freedom $`q_l(𝐤)`$, without the need to introduce the variables $`\stackrel{~}{V}_l(𝐱)`$. This is in fact what most textbooks on quantum field theory do.
The representation in terms of the $`q_l(𝐤)`$ is certainly the most transparent picture one in the free case. But when the electromagnetic field is coupled to a charged matter field then the Hamiltonian density becomes nonlocal. I.e. the value of the Hamiltonian density at the momentum $`𝐤`$ will depend on the value of fields at other momenta $`𝐥𝐤`$. Therefore in the coupled case, the representation loses some of its attractiveness.
The corresponding guidance equations in terms of the Fourier modes read
$$\dot{q}_l(𝐤)=\frac{1}{2i|\mathrm{\Phi }|^2}\left(\mathrm{\Phi }^{}\frac{\delta }{\delta \dot{q}_l(𝐤)}\mathrm{\Phi }\mathrm{\Phi }\frac{\delta }{\delta \dot{q}_l(𝐤)}\mathrm{\Phi }^{}\right),\text{with }l=1,2.$$
(4.169)
Here, we recognize the pilot-wave interpretation for the electromagnetic field that was originally presented by Bohm and which was further developed by Kaloyerou .<sup>14</sup><sup>14</sup>14In fact there is a slight difference between the model presented in and the model presented in . In the model in the fields satisfy $`q_1(𝐤)=q_1^{}(𝐤)`$ and $`q_2(𝐤)=q_2^{}(𝐤)`$, whereas in the model presented in the fields satisfy $`q_l(𝐤)=q_l^{}(𝐤)`$. The model presented here hence corresponds with the model of Bohm. Of course the transition from one model to the other is a triviality.
Bohm and Kaloyerou directly used the expansion in terms of the transversal Fourier modes of the electromagnetic potential and a similar expansion for the transversal part of the momentum field, in order for the fields to satisfy respectively the Coulomb gauge constraint and the Gauss constraint. Only afterwards, the commutation relations for the transversal Fourier modes are then imposed. This is in accordance with the treatment that can be found in many textbooks. The Dirac approach of dealing with constraints justifies these arguments.
In Kaloyerou also took a different approach to a pilot-wave interpretation. Instead of passing to Fourier modes, Kaloyerou devised a pilot-wave interpretation in which he introduced field beables corresponding to the vector potential $`𝐕(𝐱)`$. However, we think this approach is not justified (or at least incomplete). The reason is the following. In order to find a representation for the field operators $`\widehat{V}_i`$ and $`\widehat{\mathrm{\Pi }}_{V_i}`$ he replaced the transversal $`\delta `$-function, i.e. $`(\delta _{ij}_i_j/^2)\delta (𝐱𝐲)`$, on the right hand side of equation (4.119) by $`\delta _{ij}\delta (𝐱𝐲)`$. Kaloyerou argued that this replacement is justified because the same equations are obtained when the vector potential is expressed in terms of normal modes. By replacing the transverse $`\delta `$-function, the commutation relations for the operators become indeed simple and a representation is easily found. However, Kaloyerou’s argument is not correct. It is a direct consequence of the Gauss constraint and the Coulomb gauge constraint that we need the transverse $`\delta `$-function in (4.119). By replacing the transverse $`\delta `$-function by $`\delta (𝐱𝐲)`$, it is in fact implicitly assumed that the Gauss constraint and the Coulomb gauge do not apply. But without the Gauss constraint and the Coulomb gauge constraint we cannot dismiss the longitudinal modes in the Fourier expansions of the field operators $`\widehat{V}_i`$ and $`\widehat{\mathrm{\Pi }}_{V_i}`$, whereas we can do this if we have the Gauss constraint and the Coulomb gauge constraint. Nevertheless, when Kaloyerou makes the Fourier expansion of the fields $`V_i`$ and $`\mathrm{\Pi }_{V_i}`$, as presented above, he correctly dismisses the longitudinal components.
Kaloyerou also discusses the gauge invariance of the pilot-wave model. We will turn to this issue in the next section.
In , Kaloyerou further discusses in detail the application of this pilot-wave interpretation to various typical quantum phenomena involving the electromagnetic field. In particular a detailed account was given for the photo-electric and Compton effect, along the lines of Bohm . Recently Kaloyerou also applied the pilot-wave interpretation to describe the Mach-Zehnder Wheeler delayed-choice experiment .
#### 4.5.3 Quantization in other gauges
##### Admissible gauges
The Coulomb gauge is a natural gauge for the quantization of the electromagnetic field. It allows us to identify the transversal components of the electromagnetic vector potential as the true degrees of freedom and longitudinal component of the vector potential as the gauge degree of freedom. There also exist other gauges which allow for the quantization of the electromagnetic field by means of Dirac’s prescription. These gauges were called the admissible gauges in Section 4.2.3. A gauge is admissible if it is attainable by a sequence of infinitesimal gauge transformations and if it uniquely fixes the gauge. Sometimes remaining gauge invariance can be removed restricting the possible gauge transformations (4.99) by requiring boundary condition for the function $`\theta (x)`$ (for example in the Coulomb gauge we require that $`\theta (x)`$ vanishes at spatial infinity).
If a gauge is admissible then we can perform the Dirac method of quantization by adding the gauge to the set of first class constraints. For a finite number of degrees of freedom, the true degrees of freedom that may be found by performing a Maskawa and Nakajima canonical transformation are unique up to a canonical transformation. As a result, for a finite number of degrees of freedom, the Hamiltonian formulation that is obtained by imposing gauge constraints does not depend on the particular choice of gauge and hence there might only be one ambiguity in the canonical quantization procedure, which is the operator ordering ambiguity. Although it is to be expected that the true degrees of freedom are also unique for a system described by an infinite number of degrees of freedom, we have not seen an explicit statement or proof of this. But even if we take it for granted that the Hamiltonian formulation for a system described by an infinite number of degrees of freedom does not depend on the particular choice of gauge, there still may appear other ambiguities when applying the canonical quantization procedure. Apart from the operator ordering problem there is also the problem that one can use representations which yield unitarily inequivalent quantum theories \[114, pp. 53-55\]. This is a peculiarity of quantum field theory which is not present when quantizing a system with a finite number of degrees of freedom.
Different quantum field theories will of course lead to different pilot-wave interpretations. However, even equivalent quantum field theories may lead to different pilot-wave interpretations. We have already seen that there is an ambiguity in identifying the guidance equations. But also, different representations, which may yield equivalent quantum theories, may lead to pilot-wave interpretations with inequivalent ontologies. For example, in non-relativistic quantum theory one can use the configuration representation or the momentum representation, which are equivalent at the quantum level. But, as shown by Brown and Hiley, the corresponding pilot-wave interpretations are not equivalent .
For certain classes of representations, may we explicitly show that the corresponding quantum theories and the corresponding pilot-wave interpretations are equivalent. An example is given in Appendix B, where we consider the quantum theories which arise by using different transformation matrices $`K`$, which were used in the Maskawa and Nakajima canonical transformation (4.122), (4.125) and (4.126), and we show that they are equivalent on the level of quantum theory and on the level of the pilot-wave interpretation.
In the rest of this section, we will leave this issue of uniqueness aside. In the next paragraph we will consider some frequently used gauges and look whether they may naturally lead to a pilot-wave interpretation. We will see that none of the discussed gauges is straightforwardly amenable for developing a pilot-wave interpretation. In the second to next paragraph, we make a note on the gauge invariance of the pilot-wave model.
##### Some examples of frequently used gauges
A first example is the Lorentz gauge: $`_\mu V^\mu =0`$. This gauge is often used because it is explicitly Lorentz covariant. However, the Lorentz gauge contains $`_0V_0`$ which is not expressible in terms of the conjugate momenta and hence this gauge is not suitable for the Dirac procedure of quantization. Nevertheless, the Dirac procedure could be maintained in this particular case by introducing fermionic fields (the ghost fields) \[110, p. 119\]. However, because it is difficult at present to construct a pilot-wave interpretation for fermionic fields (see following chapter), we will not pursue this approach.
A second example is the axial gauge: $`V_3=\mathrm{\Pi }_{V_3}+_3V_0=0`$. This is an example of an admissible gauge. The electromagnetic potential is uniquely fixed by this gauge if we restrict ourself to gauge transformations (4.99) for which the function $`\theta `$ vanishes at spatial infinity. The Dirac brackets are very simple in this case; for the variables $`V_1,V_2,\mathrm{\Pi }_{V_1}`$ and $`\mathrm{\Pi }_{V_2}`$, the Dirac bracket equals the Poisson bracket \[110, p. 142\]. However, using this gauge leads to explicit infinities in the Hamiltonian ,<sup>15</sup><sup>15</sup>15Although these papers concern non-Abelian gauge theories, some of the content can be applied to the Abelian case as well. in a similar way as we encountered in example one in previous section (on p. 4.5.2). Various suggestions of how these problems may be overcome can be found in . In , the problem is treated in the context of non-Abelian gauge theories ($`SU(N)`$, $`N>1`$, Yang-Mills theories).
Another example of an admissible gauge is the superaxial gauge. This gauge was presented by Girotti and Rothe as a solution for the infinities appearing in the Hamiltonian in the axial gauge. The superaxial gauge reads
$`V_1(x_0,x_1,x_2,x_3^{(0)})`$ $`=`$ $`V_2(x_0,x_1^{(0)},x_2,x_3^{(0)})=V_3(x_0,x_1,x_2,x_3)=0,`$
$`V_0(x_0,x_1,x_2,x_3)`$ $`=`$ $`{\displaystyle _{x_1^{(0)}}^{x_1}}𝑑x_1^{}\mathrm{\Pi }_{V_1}(x_0,x_1^{},x_2,x_3^{(0)})`$ (4.170)
$`+{\displaystyle _{x_1^{(0)}}^{x_2}}𝑑x_2^{}\mathrm{\Pi }_{V_2}(x_0,x_1^{(0)},x_2^{},x_3^{(0)})`$
$`+{\displaystyle _{x_3^{(0)}}^{x_3}}𝑑x_3^{}\mathrm{\Pi }_{V_3}(x_0,x_1,x_2,x_3^{}),`$
where $`𝐱^{(0)}`$ is some fixed point. This gauge picks a unique representative out of the equivalence class of gauge equivalent fields and can be attained with the gauge transformation (4.99) with
$`\theta `$ $`=`$ $`{\displaystyle _{x_1^{(0)}}^{x_1}}𝑑x_1^{}V_1(x_0,x_1^{},x_2,x_3^{(0)})+{\displaystyle _{x_1^{(0)}}^{x_2}}𝑑x_2^{}V_2(x_0,x_1^{(0)},x_2^{},x_3^{(0)})`$ (4.171)
$`+`$ $`{\displaystyle _{x_3^{(0)}}^{x_3}}𝑑x_3^{}V_3(x_0,x_1,x_2,x_3^{}){\displaystyle _{x_0^{(0)}}^{x_0}}𝑑x_0^{}V_0(x_0^{},x_1^{(0)},x_2^{(0)},x_3^{(0)}).`$
The resulting equal-time commutation relations for the field operators are
$$[\widehat{V}_i(𝐱),\widehat{\mathrm{\Pi }}_{V_j}(𝐲)]=i\delta _{ij}\delta (𝐱𝐲)i_{x_i}r_j(𝐱,𝐲)$$
(4.172)
with
$`r_j(𝐱,𝐲)`$ $`=`$ $`\delta _{1j}\mathrm{\Delta }(x_1,x_1^{(0)};y_1)\delta (x_2y_2)\delta (x_3^{(0)}y_3)`$
$`+\delta _{2j}\delta (x_1^{(0)}y_1)\text{sgn}(x_2y_2)\delta (x_3^{(0)}y_3)/2`$
$`+\delta _{3j}\delta (x_1y_1)\delta (x_2y_2)\mathrm{\Delta }(x_3,x_3^{(0)};y_3),`$
$`\mathrm{\Delta }(x,x^{(0)};y)`$ $`=`$ $`{\displaystyle _{x^{(0)}}^x}𝑑x^{}\delta (x^{}y).`$ (4.173)
The commutation relations involving $`\widehat{\mathrm{\Pi }}_{V_0}`$ are zero and the commutation relations involving $`\widehat{V}_0`$ can be obtained from the commutation relations above by using the operator equivalents of the gauge constraints (4.170). For the superaxial gauge it is suggestive to take $`V_1`$ and $`V_2`$, with $`V_1(x_1,x_2,x_3^{(0)})=V_2(x_1^{(0)},x_2,x_3^{(0)})=0`$, as the true degrees of freedom. This leads to the following natural representation. Take
$`\widehat{V}_k(𝐱)`$ $`=`$ $`V_k(𝐱)k=1,2,`$
$`\widehat{V}_3(𝐱)`$ $`=`$ $`0,`$ (4.174)
where $`V_1(x_1,x_2,x_3^{(0)})=V_2(x_1^{(0)},x_2,x_3^{(0)})=0`$. Because the representation should be compatible with the commutation relations (4.172), we find
$`\widehat{\mathrm{\Pi }}_1(𝐱)`$ $`=`$ $`i{\displaystyle \frac{\delta }{\delta V_1(𝐱)}}+i\delta (x_3^{(0)}x_3){\displaystyle 𝑑x_3^{}\frac{\delta }{\delta V_1(x_1,x_2,x_3^{})}}`$
$`i_{x_2}\delta (x_3^{(0)}x_3){\displaystyle 𝑑x_1^{}𝑑x_3^{}\mathrm{\Delta }(x_1^{},x_1^{(0)};x_1)\frac{\delta }{\delta V_2(x_1^{},x_2,x_3^{})}}`$
$`\widehat{\mathrm{\Pi }}_2(𝐱)`$ $`=`$ $`i{\displaystyle \frac{\delta }{\delta V_2(𝐱)}}+i\delta (x_1^{(0)}x_1)\delta (x_3^{(0)}x_3){\displaystyle 𝑑x_1^{}𝑑x_3^{}\frac{\delta }{\delta V_2(x_1^{},x_2,x_3^{})}},`$
$`\widehat{\mathrm{\Pi }}_3(𝐱)`$ $`=`$ $`i{\displaystyle \underset{k=1}{\overset{2}{}}}{\displaystyle 𝑑x_3^{}_{x_k}\mathrm{\Delta }(x_3^{},x_3^{(0)};x_3)\frac{\delta }{\delta V_1(x_1,x_2,x_3^{})}},`$ (4.175)
where it is understood that functional derivatives with respect to the fields $`V_1(x_1,x_2,x_3^{(0)})`$ and $`V_2(x_1^{(0)},x_2,x_3^{(0)})`$ are put zero.
Although the commutation relations (4.172) and the constraints are satisfied in this representation, it is not possible to use it, because this representation leads to explicit infinities in the Hamiltonian (this is a direct consequence of the $`\delta `$-functions appearing in the conjugate momenta).
Recall that the superaxial gauge was introduced to deal with the infinities in the Hamiltonian for the axial gauge. Although the superaxial gauge does not lead to explicit infinities in the Hamiltonian operator, the infinities in the Hamiltonian reappear if we take the most obvious representation given by (4.174) and (4.175).
We can conclude that admissible gauges are not always suitable to construct a quantum theory or they may have associated representation that are not suitable. Without a suitable quantum theory it is then of course not possible to devise a pilot-wave interpretation.
##### Note on gauge invariance of the pilot-wave interpretation
In , Kaloyerou made some comments on the gauge invariance of the pilot-wave interpretation. Kaloyerou posed two questions:<sup>16</sup><sup>16</sup>16These questions are cited from .
* According to our ontology, are physical results, i.e. expectation values of field observables gauge invariant?
* Does the gauge freedom conflict with attributing ontological significance to the potentials?
Kaloyerou answers the first question affirmative by recalling the equivalence, at the standard quantum mechanical level, of the Gupta-Bleuler formalism and the formalism that arises in the Coulomb gauge. Because the Gupta-Bleuler formalism leads to gauge invariant expectation values, then also the formalism in the Coulomb gauge leads to gauge invariant expectation values. Because the pilot-wave model, devised for the quantum formalism in the Coulomb gauge, produces the same statistics as the standard quantum interpretation, the pilot-wave model yields gauge invariant expectation values of field observables. However, there is in fact no need to compare the quantum theory in the Coulomb gauge to the Gupta-Bleuler formalism. From our analysis it follows that the quantum theory, although it is devised by using the Coulomb gauge, was formulated solely in terms of gauge invariant degrees of freedom, i.e. degrees of freedom which commute with the generators of gauge transformations, and hence all the predictions of this quantum theory are gauge invariant. Because the pilot-wave model is equivalent to the standard interpretation at the empirical level, its empirical predictions are also gauge invariant.
The second question was inspired by the situation for classical electromagnetism. In classical electromagnetism different potentials may correspond to the same physical situation. Therefore one can either adopt the position to attach an ontological status to all the potentials, but in physical measurements the potential may only be revealed up to a gauge transformation. Or one could adopt the position to attach an ontological status only to the potentials that satisfy a particular gauge. Now the second question does in fact not apply to the pilot-wave model presented above, because in this model beables were only introduced for gauge invariant variables. The question is meaningful for pilot-wave models where beables are introduced for gauge variables as well, as in Valentini’s model, which is to be discussed in the next section. However, we do not favor such models because, as we will see in the following section, these models tend to lead to non-normalizable densities for the field beables.
#### 4.5.4 Quantization with constraints as conditions on the state
##### Quantization and functional Schrödinger equation
Instead of adding further gauge constraints to the set of first class constraints $`\{\chi _1,\chi _2\}`$ in order to quantize the electromagnetic field, we can also proceed another way. As explained in Section 4.2.3, we can use the standard commutation relations for the fields, i.e.
$`[\widehat{V}_0(𝐱),\widehat{\mathrm{\Pi }}_{V_0}(𝐲)]`$ $`=`$ $`i[V_0(𝐱),\mathrm{\Pi }_{V_0}(𝐲)]_P=i\delta (𝐱𝐲)`$
$`[\widehat{V}_i(𝐱),\widehat{\mathrm{\Pi }}_{V_j}(𝐲)]`$ $`=`$ $`i[V_i(𝐱),\mathrm{\Pi }_{V_j}(𝐲)]_P=i\delta _{ij}\delta (𝐱𝐲)`$ (4.176)
and impose the constraints as conditions on states,
$$\widehat{\chi }_1|\mathrm{\Psi }=\widehat{\mathrm{\Pi }}_{V_0}|\mathrm{\Psi }=0,\widehat{\chi }_2|\mathrm{\Psi }=_i\widehat{\mathrm{\Pi }}_{V_i}|\mathrm{\Psi }=0.$$
(4.177)
The states $`|\mathrm{\Psi }`$ which satisfy these constraint equations are then the physical states. Because the constraints $`\chi _1`$ and $`\chi _2`$ were generators of gauge transformations, the conditions (4.177) mean that physical states are gauge invariant.
Because we have the commutation relations (4.176), we can use the standard representation for the field operators in the Schrödinger picture, i.e.
$`\widehat{V}_0(𝐱)`$ $`=`$ $`V_0(𝐱),\widehat{\mathrm{\Pi }}_{V_0}(𝐱)=i{\displaystyle \frac{\delta }{\delta V_0(𝐱)}},`$
$`\widehat{V}_i(𝐱)`$ $`=`$ $`V_i(𝐱),\widehat{\mathrm{\Pi }}_{V_i}(𝐱)=i{\displaystyle \frac{\delta }{\delta V_i(𝐱)}}.`$ (4.178)
By using this representation for the extended Hamiltonian, given in (4.106), we obtain the following functional Schrödinger equation for the wavefunctional $`\mathrm{\Psi }(V_0,V_i,t)=V_0,V_i|\mathrm{\Psi }(t)`$
$$i\frac{\mathrm{\Psi }}{t}=d^3x\left(\frac{1}{2}\frac{\delta ^2}{\delta V_i\delta V_i}+\frac{1}{4}F_{ij}F_{ij}iV_0_j\frac{\delta }{\delta V_j}iu_1\frac{\delta }{\delta V_0}iu_2_j\frac{\delta }{\delta V_j}\right)\mathrm{\Psi }.$$
(4.179)
The physical states further have to satisfy (4.177), i.e.
$`{\displaystyle \frac{\delta }{\delta V_0}}\mathrm{\Psi }`$ $`=`$ $`0,`$ (4.180)
$`_i{\displaystyle \frac{\delta }{\delta V_i}}\mathrm{\Psi }`$ $`=`$ $`0.`$ (4.181)
The first constraint implies that $`\mathrm{\Psi }`$ does not depend on $`V_0`$. This means that $`\mathrm{\Psi }`$ is invariant under transformations of the form
$$V_0(𝐱)V_0^{\theta _0}(𝐱)=V_0(𝐱)+\theta _0(𝐱)$$
(4.182)
with $`\theta _0`$ an arbitrary space-dependent function. The second constraint implies that $`\mathrm{\Psi }`$ is invariant under time independent gauge transformations, i.e. that $`\mathrm{\Psi }`$ is invariant under transformations
$$V_i(𝐱)V_i^\theta (𝐱)=V_i(𝐱)+_i\theta (𝐱),$$
(4.183)
with $`\theta `$ an arbitrary space dependent function.
Hence, for a physical state $`\mathrm{\Psi }`$, the functional Schrödinger equation (4.179) can be written as
$$i\frac{\mathrm{\Psi }}{t}=d^3x\left(\frac{1}{2}\frac{\delta ^2}{\delta V_i\delta V_i}+\frac{1}{4}F_{ij}F_{ij}\right)\mathrm{\Psi }.$$
(4.184)
##### Pilot-wave interpretation
There is a corresponding conservation equation
$$\frac{|\mathrm{\Psi }|^2}{t}+d^3x\left(\frac{\delta J_{V_0}(𝐱)}{\delta V_0(𝐱)}+\frac{\delta J_{V_i}(𝐱)}{\delta V_i(𝐱)}\right)=0,$$
(4.185)
with
$`J_{V_0}(𝐱)`$ $`=`$ $`0,`$ (4.186)
$`J_{V_j}(𝐱)`$ $`=`$ $`{\displaystyle \frac{1}{2i}}\left(\mathrm{\Psi }^{}{\displaystyle \frac{\delta }{\delta V_j(𝐱)}}\mathrm{\Psi }\mathrm{\Psi }{\displaystyle \frac{\delta }{\delta V_j(𝐱)}}\mathrm{\Psi }^{}\right).`$ (4.187)
It would now be tempting to construct a pilot-wave interpretation, inspired by this conservation equation, by taking the guidance equations for the field beables $`V_0`$ and $`V_i`$ as
$$\dot{V}_0(𝐱)=J_{V_0}(𝐱)/|\mathrm{\Psi }|^2,\dot{V}_i(𝐱)=J_{V_i}(𝐱)/|\mathrm{\Psi }|^2.$$
(4.188)
If the field beable $`V_0`$ is discarded, it is in fact simply constant in time by (4.186), this is exactly the approach to a pilot-wave interpretation by Valentini . Though, Valentini derived his pilot-wave interpretation by considering different starting principles (this is the ‘3+1’ view). To be more precise, from the constrained dynamics point of view, the quantization scheme implicitly used by Valentini is a mixture of the two schemes described in Section 4.2.3. First, the temporal gauge $`V_0=0`$ is chosen. In accordance with the first scheme explained in Section 4.2.3, the constraints $`V_0=\mathrm{\Pi }_{V_0}=0`$ are then treated as operator identities, so that the fields $`V_0`$ and $`\mathrm{\Pi }_{V_0}`$ will make no further appearance in the theory. Second, the unconstrained canonical commutation relations are imposed on the fields $`V_i`$ and $`\mathrm{\Pi }_{V_i}`$ and the remaining constraint $`\chi _2=_iV_i=0`$ is then imposed as a condition on states, in accordance with the second scheme explained in Section 4.2.3.
##### Problem with non-normalizable field beable densities
There is a problem with this pilot-wave approach. The density $`|\mathrm{\Psi }(V_0,V_i,t)|^2`$ of field beables is not normalizable with respect to the variables $`V_0`$ and $`V_i`$. This is because $`\mathrm{\Psi }`$ does not depend on $`V_0`$ and is further invariant under time independent gauge transformations (4.183). Hence the integral is proportional to the volume of the gauge group (by the gauge group we mean the group of transformations determined by (4.182) and (4.183)), which is infinite. I.e. we have
$$𝒟V_0\left(\mathrm{\Pi }_{j=1}^3𝒟V_j\right)|\mathrm{\Psi }(V_0,V_i,t)|^2𝒟V_0𝒟\theta 𝒟\theta _0𝒟\theta \mathrm{}.$$
(4.189)
Nevertheless, there is an easy way out to this problem. Let us make a change in field variables from $`V_i`$ to $`\stackrel{~}{V}_i`$ by the transformation
$$V_i(𝐱)=d^3yK_{ij}(𝐱,𝐲)\stackrel{~}{V}_j(𝐲),$$
(4.190)
with $`K`$ a matrix which satisfies the properties discussed in Section 4.5.2. We can use this transformation to rewrite the functional Schrödinger equation (4.184) and the constraints (4.182) and (4.183) as equations for a wavefunctional $`\mathrm{\Psi }^{}(V_0,\stackrel{~}{V}_i)`$.
Let us first look at the constraints. As before, the constraint (4.182) implies that $`\mathrm{\Psi }^{}`$ is independent of $`V_0`$. For the second constraint, we use the relation
$$\frac{\delta }{\delta V_i}(𝐱)=d^3yK_{ji}^1(𝐲,𝐱)\frac{\delta }{\delta \stackrel{~}{V}_j(𝐲)}$$
(4.191)
to obtain
$`_{x_i}{\displaystyle \frac{\delta }{\delta V_i(𝐱)}}\mathrm{\Psi }=0`$ $``$ $`{\displaystyle d^3y_{x_i}K_{ji}^1(𝐲,𝐱)\frac{\delta }{\delta \stackrel{~}{V}_j(𝐲)}\mathrm{\Psi }^{}}=0`$ (4.192)
$``$ $`{\displaystyle d^3yU^1(𝐲,𝐱)\frac{\delta }{\delta \stackrel{~}{V}_3(𝐲)}\mathrm{\Psi }^{}}=0`$
$``$ $`{\displaystyle \frac{\delta }{\delta \stackrel{~}{V}_3(𝐲)}}\mathrm{\Psi }^{}=0.`$
Hence we find the constraints imply that $`\mathrm{\Psi }^{}`$ is independent of $`V_0`$ and $`\stackrel{~}{V}_3`$, i.e. $`\mathrm{\Psi }^{}=\mathrm{\Psi }^{}(\stackrel{~}{V}_1,\stackrel{~}{V}_2,t)`$.
If we rewrite the functional Schrödinger equation (4.184) as an equation for $`\mathrm{\Psi }^{}(\stackrel{~}{V}_1,\stackrel{~}{V}_2,t)`$ we obtain
$`i{\displaystyle \frac{\mathrm{\Psi }^{}(\stackrel{~}{V}_1,\stackrel{~}{V}_2,t)}{t}}`$ $`=`$ $`{\displaystyle \underset{k,l=1}{\overset{2}{}}}{\displaystyle }d^3yd^3z(h_{kl}(𝐲,𝐳){\displaystyle \frac{\delta }{\delta \stackrel{~}{V}_k(𝐲)}}{\displaystyle \frac{\delta }{\delta \stackrel{~}{V}_l(𝐳)}}+`$ (4.193)
$`\overline{h}_{kl}(𝐲,𝐳)\stackrel{~}{V}_k(𝐲)\stackrel{~}{V}_l(𝐳))\mathrm{\Psi }^{}(\stackrel{~}{V}_1,\stackrel{~}{V}_2,t),`$
with $`h_{kl}(𝐲,𝐳)`$ and $`\overline{h}_{kl}(𝐲,𝐳)`$ as defined in (4.140) and (4.141). In other words we obtain the functional Schrödinger equation (4.143), which was the functional Schrödinger equation in terms of unconstrained variables when we quantized the electromagnetic field in the Coulomb gauge.
Because the fields $`V_0`$ and $`\stackrel{~}{V}_3`$ make no appearance in the theory anymore we can in fact safely forget about them. We should only introduce beables corresponding to the fields $`\stackrel{~}{V}_1`$ and $`\stackrel{~}{V}_2`$ and not to the fields $`V_0`$ and $`\stackrel{~}{V}_3`$. It would not even be very meaningful to introduce beables for the fields $`V_0`$ and $`\stackrel{~}{V}_3`$. They would just remain constant with time. In addition, by dismissing the fields $`V_0`$ and $`\stackrel{~}{V}_3`$ we do not encounter a problem anymore with infinities when normalizing the density of field beables $`|\mathrm{\Psi }^{}|^2`$. The infinity would appear if we would integrate the density over the fields $`V_0`$ and $`\stackrel{~}{V}_3`$.
Of course the pilot-wave theory that we obtain by dismissing the fields $`V_0`$ and $`\stackrel{~}{V}_3`$ is exactly the pilot-wave theory that was obtained in the previous section. We think this is the most natural approach to a pilot-wave interpretation for the electromagnetic field. By getting rid of the gauge degrees of freedom we do not have a problem with normalizing the densities of field beables. In addition, we saw that introducing beables for gauge degrees of freedom is rather meaningless because the constraints imply that these beables remain constant in time.
##### Note on the definition of the inner product
We saw that keeping gauge degrees of freedom when developing a pilot-wave interpretation leads to non-normalizable densities of field beables. A related problem arises in standard quantum field theory. When the inner product of two gauge invariant states $`|\mathrm{\Psi }_1`$ and $`|\mathrm{\Psi }_2`$ were to be defined as
$$\mathrm{\Psi }_1|\mathrm{\Psi }_2=𝒟V_0\left(\mathrm{\Pi }_{j=1}^3𝒟V_j\right)\mathrm{\Psi }_1^{}\mathrm{\Psi }_2,$$
(4.194)
then this product would be infinite for the same reasons. In addition, if expectation values of operators would be defined in a similar way, one would encounter ambiguities . This would for example be the case for the following expectation value
$$\mathrm{\Psi }|[\widehat{V}_i(𝐱),_j\widehat{\mathrm{\Pi }}_{V_j}(𝐲)]|\mathrm{\Psi }.$$
(4.195)
By using (4.177) one finds that the expectation value is zero. If on the other hand the expectation value is calculated with (4.176) then one finds that this quantity is different from zero.
In standard quantum field theory, these problems are solved by introducing a measure $`\mu (V)`$ on the fields . The measure is found by applying the Faddeev-Popov trick. With the Faddeev-Popov trick, the gauge volume can explicitly be factored out from the integral in (4.194). The remaining part then represent the integration of $`\mathrm{\Psi }_1^{}\mathrm{\Psi }_2`$ over gauge independent variables, which yields a finite number.
It is instructive to apply the Faddeev-Popov formalism explicitly. Suppose a gauge $`\overline{\chi }_3(V_i)=0`$ which picks a unique representative from each equivalence class of fields that are connected by time independent gauge transformations. Then
$$1=\mathrm{\Delta }(\overline{\chi }_3(V_j))𝒟\theta \delta \left(\overline{\chi }_3\left(V_i^\theta \right)\right)$$
(4.196)
with
$$\mathrm{\Delta }(\overline{\chi }_3(V_j))=\left|det\left(\frac{\delta \overline{\chi }_3(V_j^\theta )(𝐱)}{\delta \theta (𝐲)}\right)\right||_{\overline{\chi }_3(V_i^\theta )=0}$$
(4.197)
the Faddeev-Popov determinant, which is gauge invariant. Suppose similarly a gauge $`\overline{\chi }_4(V_0)=0`$ which picks a unique representative from each equivalence class of fields that are connected by the transformations (4.182), then
$$1=\mathrm{\Delta }(\overline{\chi }_4(V_0))𝒟\theta _0\delta \left(\overline{\chi }_4\left(V_0^{\theta _0}\right)\right),$$
(4.198)
with
$$\mathrm{\Delta }(\overline{\chi }_4(V_0))=\left|det\left(\frac{\delta \overline{\chi }_4(V_0^{\theta _0})(𝐱)}{\delta \theta _0(𝐲)}\right)\right||_{\overline{\chi }_4(V_0^{\theta _0})=0}.$$
(4.199)
By substituting (4.196) and (4.198) in the inner product (4.207), we can write
$`{\displaystyle 𝒟V_0\left(\mathrm{\Pi }_{j=1}^3𝒟V_j\right)\mathrm{\Psi }_1^{}\mathrm{\Psi }_2}`$
$`={\displaystyle 𝒟\theta 𝒟\theta _0𝒟V_0\left(\mathrm{\Pi }_{j=1}^3𝒟V_j\right)\mathrm{\Psi }_1^{}\mathrm{\Psi }_2\mathrm{\Delta }(\overline{\chi }_3(V_i))\mathrm{\Delta }(\overline{\chi }_4(V_0))\delta \left(\overline{\chi }_3\left(V_i\right)\right)\delta \left(\overline{\chi }_4\left(V_0\right)\right)}`$
$`=\left({\displaystyle 𝒟\theta 𝒟\theta _0}\right){\displaystyle 𝒟V_0\left(\mathrm{\Pi }_{j=1}^3𝒟V_j\right)\mathrm{\Psi }_1^{}\mathrm{\Psi }_2\mathrm{\Delta }(\overline{\chi }_3(V_i))\mathrm{\Delta }(\overline{\chi }_4(V_0))}`$
$`\times \delta \left(\overline{\chi }_3\left(V_i\right)\right)\delta \left(\overline{\chi }_4\left(V_0\right)\right).`$ (4.200)
The last equality arises because $`\mathrm{\Psi }_1`$, $`\mathrm{\Psi }_2`$, the measure $`𝒟V_0\left(\mathrm{\Pi }_{j=1}^3𝒟V_j\right)`$ and the Faddeev-Popov determinants are invariant under the gauge transformations (4.182) and (4.183). In this way, we have been able to separate the infinite gauge part from the integral. We can now define the measure
$$\mu (V)=\frac{\delta \left(\overline{\chi }_3\left(V_i\right)\right)\delta \left(\overline{\chi }_4\left(V_0\right)\right)}{𝒟\theta 𝒟\theta _0}$$
(4.201)
and a new inner product
$`\mathrm{\Psi }_1|\mathrm{\Psi }_2`$ $`=`$ $`{\displaystyle 𝒟V_0\left(\mathrm{\Pi }_{j=1}^3𝒟V_j\right)\mu (V)\mathrm{\Psi }_1^{}\mathrm{\Psi }_2}`$ (4.202)
$`=`$ $`{\displaystyle 𝒟V_0\left(\mathrm{\Pi }_{j=1}^3𝒟V_j\right)\mathrm{\Psi }_1^{}\mathrm{\Psi }_2\mathrm{\Delta }(\overline{\chi }_3(V_i))\mathrm{\Delta }(\overline{\chi }_4(V_0))\delta \left(\overline{\chi }_3\left(V_i\right)\right)\delta \left(\overline{\chi }_4\left(V_0\right)\right)}`$
$`=`$ $`{\displaystyle \left(\mathrm{\Pi }_{j=1}^3𝒟V_j\right)\mathrm{\Psi }_1^{}\mathrm{\Psi }_2\mathrm{\Delta }(\overline{\chi }_3(V_i))\delta \left(\overline{\chi }_3\left(V_i\right)\right)}`$
which is finite. With the introduction of the measure $`\mu (V)`$ also the ambiguities with expectation values such as (4.195) are removed .
We can for example use the Coulomb gauge for $`\overline{\chi }_3`$, i.e. $`\overline{\chi }_3=_iV_i`$. In order to perform the integral in the inner product in (4.194) explicitly, we can make a transition to the new variables $`\stackrel{~}{V}_i`$, $`i=1,2,3`$, by performing the transformation (4.190). First we define
$$N_k\mathrm{\Psi }_k^{}(V_0,\stackrel{~}{V}_i,t)=\mathrm{\Psi }_k(V_0,\underset{j=1}{\overset{3}{}}d^3yK_{ij}(𝐱,𝐲)\stackrel{~}{V}_j(𝐲),t)$$
(4.203)
for $`k=1,2`$. The $`N_k`$ are normalization constants which will be determined later. With $`\mathrm{\Psi }_1`$ and $`\mathrm{\Psi }_2`$ being gauge invariant states we know that $`\mathrm{\Psi }_1^{}`$ and $`\mathrm{\Psi }_2^{}`$ will not depend on $`V_0`$ and $`\stackrel{~}{V}_3`$.
If we now make use of the identities
$$\mathrm{\Delta }(\overline{\chi }_3(V_i))=\left|det\left(^2\delta (𝐱𝐲)\right)\right|,$$
(4.204)
and
$$\delta (_iV_i)=\frac{\delta (\stackrel{~}{V}_3)}{\left|det\left(_x^2U(𝐱,𝐲)\right)\right|},$$
(4.205)
we can write the inner product in (4.195) as
$$\mathrm{\Psi }_1|\mathrm{\Psi }_2=\frac{N_1N_2\left|detK\right|}{\left|detU\right|}𝒟\stackrel{~}{V}_1𝒟\stackrel{~}{V}_2\mathrm{\Psi }_1^{}(\stackrel{~}{V}_1,\stackrel{~}{V}_2,t)^{}\mathrm{\Psi }_2^{}(\stackrel{~}{V}_1,\stackrel{~}{V}_2,t).$$
(4.206)
In particular the norm of gauge invariant wavefunctionals $`\mathrm{\Psi }`$ reads
$$\mathrm{\Psi }|\mathrm{\Psi }=\frac{N^2\left|detK\right|}{\left|detU\right|}𝒟\stackrel{~}{V}_1𝒟\stackrel{~}{V}_2|\mathrm{\Psi }^{}(\stackrel{~}{V}_1,\stackrel{~}{V}_2,t)|^2.$$
(4.207)
with $`\mathrm{\Psi }^{}`$ defined similarly as in (4.195). If we take $`N_k^2=\left|detK\right|/\left|detU\right|`$, then the wavefunctionals $`\mathrm{\Psi }^{}(\stackrel{~}{V}_1,\stackrel{~}{V}_2,t)`$ are normalized with respect to the variables $`\stackrel{~}{V}_1`$ and $`\stackrel{~}{V}_2`$. From this expression for the norm of a state, we see that the natural definition for the density of field beables $`\stackrel{~}{V}_1`$ and $`\stackrel{~}{V}_2`$ is $`|\mathrm{\Psi }^{}(\stackrel{~}{V}_1,\stackrel{~}{V}_2,t)|^2`$. The pilot-wave theory that follows by considering the continuity equation for this density is of course the one we presented before.
### 4.6 Scalar quantum electrodynamics
In this section we present a pilot-wave field interpretation for a quantized bosonic field interacting with a quantized electromagnetic field. In Sections 4.3 and 4.4, we have shown that the quantized DKP theory coupled to a non-quantized electromagnetic field is the same as the quantized Klein-Gordon or quantized Proca theory coupled to a non-quantized electromagnetic field. This equivalence is also true if the electromagnetic field is quantized. Therefore we do not not bother to start from the coupled DKP theory here. Instead we start with the coupled Klein-Gordon theory (called scalar quantum electrodynamics). The Proca theory can be treated completely analogously and will therefore not be discussed here.
For the electromagnetic field the quantization is only slightly different compared to the free case. This is because the constraints for the electromagnetic field differ. After the discussion on the quantization of the theory we present the corresponding pilot-wave interpretation.
#### 4.6.1 Quantization in the Coulomb gauge
The Lagrangian is given by the sum of the minimally coupled Klein-Gordon Lagrangian and the free Maxwell Lagrangian
$$L=d^3x=d^3x\left(D_\mu ^{}\varphi ^{}D^\mu \varphi m\varphi ^{}\varphi \frac{1}{4}F^{\mu \nu }F_{\mu \nu }\right).$$
(4.208)
The equations of motion are
$`_\mu F^{\mu \nu }=s_{KG}^\nu ,`$ (4.209)
$`D_\mu D^\mu \varphi +m^2\varphi =0,D_\mu ^{}D^\mu \varphi ^{}+m^2\varphi ^{}=0`$ (4.210)
with $`s_{KG}^\mu =ie\left(\varphi ^{}D^\mu \varphi \varphi D^\mu \varphi ^{}\right)`$ the Klein-Gordon charge current. In the following we will not write the subscript ‘$`KG`$’ anymore. The Lagrangian and hence the equations of motion are invariant under the gauge transformations
$`\varphi e^{ie\theta }\varphi ,\varphi ^{}e^{ie\theta }\varphi ^{},`$
$`V^\mu V^\mu ^\mu \theta .`$ (4.211)
The canonically conjugate momenta read
$`\mathrm{\Pi }_\varphi ={\displaystyle \frac{\delta L}{\delta \dot{\varphi }}}=D_0^{}\varphi ^{},\mathrm{\Pi }_\varphi ^{}={\displaystyle \frac{\delta L}{\delta \dot{\varphi }^{}}}=D_0\varphi `$
$`\mathrm{\Pi }_{V_0}={\displaystyle \frac{\delta L}{\delta \dot{V}_0}}=0,\mathrm{\Pi }_{V_i}={\displaystyle \frac{\delta L}{\delta \dot{V}_i}}=(_0V_i+_iV_0).`$ (4.212)
We can read of that we have one primary constraint, i.e.
$$\chi _1=\mathrm{\Pi }_{V_0}.$$
(4.213)
The canonical Hamiltonian is given by
$`H_C`$ $`=`$ $`{\displaystyle d^3x\left(\mathrm{\Pi }_\varphi \mathrm{\Pi }_\varphi ^{}+\left(D_i^{}\varphi ^{}\right)D_i\varphi +m^2\varphi ^{}\varphi +ieV_0\left(\varphi ^{}\mathrm{\Pi }_\varphi ^{}\varphi \mathrm{\Pi }_\varphi \right)\right)}`$ (4.214)
$`+{\displaystyle d^3x\left(\frac{1}{2}\mathrm{\Pi }_{V_i}\mathrm{\Pi }_{V_i}+\frac{1}{4}F_{ij}F_{ij}\mathrm{\Pi }_{V_i}_iV_0\right)}.`$
The total Hamiltonian reads
$$H_T=H_C+d^3xu_1\chi _1.$$
(4.215)
The requirement that $`\chi _1`$ is conserved leads to the secondary constraint
$$\chi _2=_i\mathrm{\Pi }_{V_i}+s_0.$$
(4.216)
In the constraint the charge density is to be considered in terms of momentum phase-space variables, i.e. $`s_0=ie\left(\varphi ^{}\mathrm{\Pi }_\varphi ^{}\varphi \mathrm{\Pi }_\varphi \right)`$. There are no further constraints. The requirement that $`\chi _2`$ is conserved in time is identically fulfilled and hence does not determine the field $`u_1(𝐱)`$. Hence we are left with two first class constraints and the constraint $`\chi _2`$ can be added to the total Hamilonian to yield the extended Hamiltonian
$$H_E=H_C+d^3xu_1\chi _1+d^3xu_2\chi _2.$$
(4.217)
We can use the Coulomb gauge to quantize the system. In this case the Coulomb gauge reads \[110, p. 151\]
$`\chi _3`$ $`=`$ $`_iV_i,`$ (4.218)
$`\chi _4`$ $`=`$ $`V_0+{\displaystyle \frac{1}{^2}}s_0.`$ (4.219)
As in the case of the free electromagnetic field this is an admissible gauge. It can be obtained by the transformation (4.99) with
$$\theta =\frac{1}{^2}_iV_i.$$
(4.220)
The Dirac bracket can be calculated by using the inverse of the matrix $`C_{NM}(𝐱,𝐲)=[\chi _N(𝐱),\chi _M(𝐲)]_P`$, with $`M,N=1,\mathrm{},4`$. The matrix $`C`$ has the following nonzero components
$`C_{14}(𝐱,𝐲)`$ $`=`$ $`C_{41}(𝐱,𝐲)=\delta (𝐱𝐲),`$
$`C_{23}(𝐱,𝐲)`$ $`=`$ $`C_{32}(𝐱,𝐲)=^2\delta (𝐱𝐲).`$ (4.221)
The inverse reads<sup>17</sup><sup>17</sup>17As in the case of the free electromagnetic field, the inverse is unique by considering the boundary conditions for the fields.
$`C_{14}^1(𝐱,𝐲)`$ $`=`$ $`C_{41}^1(𝐱,𝐲)=\delta (𝐱𝐲),`$
$`C_{23}^1(𝐱,𝐲)`$ $`=`$ $`C_{32}^1(𝐱,𝐲)={\displaystyle \frac{1}{^2}}\delta (𝐱𝐲).`$ (4.222)
We do not give the Dirac brackets explicitly, instead we will directly present the commutation relations for the operators below.
The Hamiltonian which will generate the dynamics of the fields is derived from the canonical Hamiltonian (4.214) by imposing the constraints:
$`H`$ $`=`$ $`{\displaystyle d^3x\left(\mathrm{\Pi }_\varphi ^{}\mathrm{\Pi }_\varphi +\left(D_i^T\varphi ^{}\right)D_i^T\varphi +m^2\varphi ^{}\varphi +\frac{1}{2}\mathrm{\Pi }_{V_i}^T\mathrm{\Pi }_{V_i}^T\frac{1}{2}V_i^T^2V_i^T\right)}`$ (4.223)
$`+{\displaystyle \frac{1}{2}}{\displaystyle d^3xd^3y\frac{s_0(𝐱)s_0(𝐲)}{4\pi |𝐱𝐲|}}.`$
where
$$D_i^T=_iieV_i^T.$$
(4.224)
We have hereby used the decomposition of the fields in longitudinal and transversal components as defined in (4.107). The longitudinal part of the field $`V_i`$ is zero because of the Coulomb constraint $`_iV_i=0`$ and the longitudinal part of the field $`\mathrm{\Pi }_{V_i}`$ can be expressed in terms of the charge density as
$$\mathrm{\Pi }_{V_i}^L=\frac{_i_j}{^2}\mathrm{\Pi }_{V_j}=\frac{_i}{^2}s_0.$$
(4.225)
Let us now quantize the theory by associating operators to the field variables. The equal time commutation relations for the field operators are found using the Dirac bracket. The commutation relations for the operators corresponding to the matter field read
$$[\widehat{\varphi }(𝐱),\widehat{\mathrm{\Pi }}_\varphi (𝐲)]=[\widehat{\varphi }^{}(𝐱),\widehat{\mathrm{\Pi }}_\varphi ^{}(𝐲)]=i\delta (𝐱𝐲)$$
(4.226)
The other commutation relations between the field operators $`\widehat{\varphi },\widehat{\mathrm{\Pi }}_\varphi ,\widehat{\varphi }^{},\widehat{\mathrm{\Pi }}_\varphi ^{}`$ are zero. The commutation relations for the operators corresponding to the electromagnetic field read
$`[\widehat{V}_i(𝐱),\widehat{\mathrm{\Pi }}_{V_j}(𝐲)]=i\left(\delta _{ij}{\displaystyle \frac{_i_j}{^2}}\right)\delta (𝐱𝐲),`$
$`[\widehat{V}_i(𝐱),\widehat{V}_j(𝐲)]=[\widehat{\mathrm{\Pi }}_{V_i}(𝐱),\widehat{\mathrm{\Pi }}_{V_j}(𝐲)]=0.`$ (4.227)
The commutation relations involving the operator $`\widehat{\mathrm{\Pi }}_{V_0}`$ are zero. The commutation relations of the operator $`\widehat{V}_0`$ and other operators corresponding to the electromagnetic field are also zero. If $`R`$ is a functional of the canonical variables of the matter field, then we have the following commutation relations
$`[\widehat{R},\widehat{V}_i(𝐱)]=0,[\widehat{R},\widehat{V}_0(𝐱)]=[\widehat{R},{\displaystyle \frac{1}{^2}}\widehat{s}_0(𝐱)],`$
$`[\widehat{R},\widehat{\mathrm{\Pi }}_{V_i}(𝐱)]=[\widehat{R},_i\widehat{V}_0(𝐱)]=[\widehat{R},{\displaystyle \frac{_i}{^2}}\widehat{s}_0(𝐱)].`$ (4.228)
If we make use of the transformation matrix $`K`$ defined in Section 4.5.2, a general representation for the field operators, consistent with the commutation relations and the constraints, is given by
$`\widehat{\varphi }(𝐱)`$ $`=`$ $`\varphi (𝐱),\widehat{\mathrm{\Pi }}_\varphi (𝐱)=i{\displaystyle \frac{\delta }{\delta \varphi (𝐱)}},`$
$`\widehat{\varphi }^{}(𝐱)`$ $`=`$ $`\varphi ^{}(𝐱),\widehat{\mathrm{\Pi }}_\varphi ^{}(𝐱)=i{\displaystyle \frac{\delta }{\delta \varphi ^{}(𝐱)}},`$
$`\widehat{V}_i(𝐱)`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{2}{}}}{\displaystyle d^3yK_{ik}(𝐱,𝐲)\stackrel{~}{V}_k(𝐲)},`$
$`\widehat{\mathrm{\Pi }}_{V_i}(𝐱)`$ $`=`$ $`\widehat{\mathrm{\Pi }}_{V_i}^T(𝐱)+\widehat{\mathrm{\Pi }}_{V_i}^L(𝐱),`$
$`\widehat{\mathrm{\Pi }}_{V_i}^T(𝐱)`$ $`=`$ $`i{\displaystyle \underset{k=1}{\overset{2}{}}}{\displaystyle d^3yK_{ki}^1(𝐲,𝐱)\frac{\delta }{\delta \stackrel{~}{V}_k(𝐲)}},`$
$`\widehat{\mathrm{\Pi }}_{V_i}^L(𝐱)`$ $`=`$ $`_i\widehat{V}_0(𝐱)={\displaystyle \frac{_i}{^2}}\widehat{s}_0(𝐱),`$
$`\widehat{s}_0(𝐱)`$ $`=`$ $`e\left(\varphi ^{}(𝐱){\displaystyle \frac{\delta }{\delta \varphi ^{}(𝐱)}}\varphi (𝐱){\displaystyle \frac{\delta }{\delta \varphi (𝐱)}}\right).`$ (4.229)
With this representation, the operator Hamiltonian becomes
$`\widehat{H}`$ $`=`$ $`{\displaystyle d^3x\left(\frac{\delta }{\delta \varphi ^{}}\frac{\delta }{\delta \varphi }+\left(D_i^T\varphi ^{}\right)D_i^T\varphi +m^2|\varphi |^2\right)}`$ (4.230)
$`+{\displaystyle \underset{k,l=1}{\overset{2}{}}}{\displaystyle d^3yd^3z\left(h_{kl}(𝐲,𝐳)\frac{\delta }{\delta \stackrel{~}{V}_k(𝐲)}\frac{\delta }{\delta \stackrel{~}{V}_l(𝐳)}+\overline{h}_{kl}(𝐲,𝐳)\stackrel{~}{V}_k(𝐲)\stackrel{~}{V}_l(𝐳)\right)}`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle d^3xd^3y\frac{\widehat{s}_0(𝐱)\widehat{s}_0(𝐲)}{4\pi |𝐱𝐲|}}`$
with
$$D_{x_i}^T=_{x_i}ie\underset{k=1}{\overset{2}{}}d^3yK_{ik}(𝐱,𝐲)\stackrel{~}{V}_k(𝐲).$$
(4.231)
The Schrödinger equation for the wavefunctional $`\mathrm{\Psi }(\varphi ,\varphi ^{},\stackrel{~}{V}_1,\stackrel{~}{V}_2,t)`$ reads
$$i\frac{\mathrm{\Psi }}{t}=\widehat{H}\mathrm{\Psi }.$$
(4.232)
Note that the Hamiltonian density depends nonlocally on the true degrees of freedom, even when the theory is written in the momentum representation. This is not only due to the Coulomb interaction (this is the last term in the Hamiltonian), but also due to the presence of the term $`\left(D_i^T\varphi ^{}\right)D_i^T\varphi `$ in the Hamiltonian.
#### 4.6.2 Pilot-wave interpretation
Now we are ready to present the pilot-wave interpretation. The conservation equation corresponding to the functional Schrödinger equation (4.232) reads
$$\frac{|\mathrm{\Psi }|^2}{t}+d^3x\left(\frac{\delta J_\varphi (𝐱)}{\delta \varphi (𝐱)}+\frac{\delta J_\varphi ^{}(𝐱)}{\delta \varphi ^{}(𝐱)}+\underset{k=1}{\overset{2}{}}\frac{\delta J_{\stackrel{~}{V}_k}(𝐱)}{\delta \stackrel{~}{V}_k(𝐱)}\right)=0,$$
(4.233)
with
$`J_\varphi (𝐱)`$ $`=`$ $`{\displaystyle \frac{1}{2i}}\left(\mathrm{\Psi }^{}{\displaystyle \frac{\delta }{\delta \varphi ^{}(𝐱)}}\mathrm{\Psi }\mathrm{\Psi }{\displaystyle \frac{\delta }{\delta \varphi ^{}(𝐱)}}\mathrm{\Psi }^{}\right),`$
$`+{\displaystyle \frac{e^2}{2i}}{\displaystyle }d^3y{\displaystyle \frac{1}{4\pi |𝐱𝐲|}}({\displaystyle \frac{\delta \mathrm{\Psi }^{}}{\delta \varphi (𝐲)}}\varphi (𝐱)\varphi (𝐲)\mathrm{\Psi }+\mathrm{\Psi }^{}\varphi (𝐱)\varphi ^{}(𝐲){\displaystyle \frac{\delta \mathrm{\Psi }}{\delta \varphi ^{}(𝐲)}}`$
$`\mathrm{\Psi }^{}\varphi (𝐱)\varphi (𝐲){\displaystyle \frac{\delta \mathrm{\Psi }}{\delta \varphi (𝐲)}}{\displaystyle \frac{\delta \mathrm{\Psi }^{}}{\delta \varphi ^{}(𝐲)}}\varphi (𝐱)\varphi ^{}(𝐲)\mathrm{\Psi }),`$
$`J_\varphi ^{}(𝐱)`$ $`=`$ $`J_\varphi ^{}(𝐱),`$
$`J_{\stackrel{~}{V}_k}(𝐱)`$ $`=`$ $`{\displaystyle \frac{1}{2i}}{\displaystyle \underset{l=1}{\overset{2}{}}}{\displaystyle d^3yh_{kl}(𝐱,𝐲)\left(\mathrm{\Psi }^{}\frac{\delta }{\delta \stackrel{~}{V}_l(𝐲)}\mathrm{\Psi }\mathrm{\Psi }\frac{\delta }{\delta \stackrel{~}{V}_l(𝐲)}\mathrm{\Psi }^{}\right)},`$ (4.234)
with $`k=1,2`$.
In the pilot-wave interpretation the conserved density of the field beables $`\varphi ,\varphi ^{},\stackrel{~}{V}_1,\stackrel{~}{V}_2`$ is given by $`|\mathrm{\Psi }|^2`$ and the guidance equations for the fields are
$$\dot{\varphi }=J_\varphi /|\mathrm{\Psi }|^2,\dot{\varphi }^{}=J_\varphi ^{}/|\mathrm{\Psi }|^2,$$
(4.235)
$$\dot{\stackrel{~}{V}}_k=J_{\stackrel{~}{V}_k}/|\mathrm{\Psi }|^2,\text{for }k=1,2.$$
(4.236)
#### 4.6.3 Conclusion
We see that we can give a pilot-wave interpretation for a quantized scalar field coupled to a quantized electromagnetic field. In the pilot-wave interpretation we only have introduced beables corresponding to true degrees of freedom. This approach has to be contrasted with Valentini’s approach , where also beables are introduced corresponding to gauge degrees of freedom. However, just as in the free case this leads to densities of the field beables which are not normalizable.
We have used the Coulomb gauge to quantize the Maxwell field, but of course, other admissible gauges could equally well be used. An example of an admissible gauge is the superaxial gauge . However, as we have seen already in the free case, the most obvious representation for the superaxial gauge leads to infinities in the Hamiltonian. Another interesting gauge is the unitary gauge, which can be used for the treatment of the Abelian Higgs model. We could in fact give a pilot-wave account for spontaneous symmetry breaking by adding the Higgs potential to the Lagrangian of the scalar field and by quantizing the system in unitary gauge. But we shall not do it here.
### 4.7 A note on the quantization of non-Abelian gauge theories
We have seen that in order to construct a pilot-wave interpretation for a constrained system it seems essential to isolate the true degrees of freedom. This presented no problem for the electromagnetic field. It is no problem to find an admissible gauge for the electromagnetic field. These admissible gauges then in turn suggest the use of some particular set of true degrees of freedom. The situation is different for non-Abelian gauge theories ($`SU(N)`$, $`N>1`$, Yang-Mills theories).
It seems difficult to find a admissible gauge for non-Abelian gauge theories. The Coulomb gauge for example does not uniquely fix the gauge. As shown by Gribov , there remain gauge equivalent fields which satisfy the Coulomb gauge (and which are non-perturbative of nature). As was shown later this is not a problem particular for the Coulomb gauge. It was namely shown that there exists no continuous gauge for (non-Abelian) Yang-Mills theories on a compactified space or space-time which uniquely fixes the gauge . Although we are not dealing with a compact space or space-time this theorem is powerful because space could be treated as compactified if the fields have suitable boundary conditions, for example if all the fields vanish at infinity. Hence if we want a continuous gauge, which uniquely fixes the gauge, then the fields (i.e. the vector potentials and the corresponding momenta) behave non-trivially at spatial infinity.
Because some gauges, such as the Coulomb gauge, can be used to fix the gauge locally, a possible solution could be to restrict the configuration space of the fields to a certain subset called the fundamental modular region such that the gauge picks a unique representative in this subset .
The axial gauge is suitable for the Dirac procedure. However, this gauge leads to explicit infinities in the Hamiltonian. The superaxial gauge , which does not yield infinities in the Hamiltonian, does not bring us any further either, because the corresponding natural representation leads, make the infinities reappear in the Hamiltonian, just as in the case of Maxwell’s theory (cf. Section 4.5.3).
The difficulty in finding an admissible gauge now leads to problems if we want to quantize the theory either by imposing constraints as operator identities or by imposing constraints as conditions on states. In the first method this is because the matrix which components are the Poisson brackets of the first class constraints $`[\chi _i,\chi _j]_P`$ is then not invertible. This blocks the construction of the Dirac bracket and hence it is unclear what commutation relations should be used for the operators. In the second method one can impose commutation relations for the operators, by taking the constraints as conditions on the states. The functional Schrödinger equation is then easily found. However, the difficulties arise with the construction of an inner product. In order to render the inner product finite, the gauge volume should be separated out of the functional integral and in order to accomplish this by performing the Faddeev-Popov trick, a admissible gauge is needed.
The issue of finding an admissible gauge and hence of finding the true degrees of freedom, is already a problem for the standard interpretation. Of course these problems persist when we want to construct a pilot-wave approach. Valentini, for example, considered an approach where not only beables are introduced for the true degrees of freedom, but also for gauge variables . However, just as in the Abelian case, cf. Section 4.5.4, the density of field beables is not normalizable.
### 4.8 The measurement process in terms of field beables
In Section 1.2 we gave the pilot-wave description of a measurement process in terms of particle beables. The description in terms of field beables proceeds along similar lines. In the field interpretation, the wavefunctionals develop non-overlapping branches in the configuration space of fields. The field beables then enter one of the branches, and under suitable conditions (the branches should not overlap again at a later time) the empty branch may be dismissed for the future description of the system. In the standard interpretation, this would then be referred to as the collapse of the wavefunctional.
Let us consider this in some more detail. Suppose we have a system described by the wavefunctional $`\mathrm{\Psi }^{(s)}(\varphi )`$. In a measurement situation the system couples to a measurement apparatus. We write the wavefunctional of the apparatus as $`\chi ^{(a)}(\stackrel{~}{\varphi })`$, where the argument $`\stackrel{~}{\varphi }`$ represents all the field degrees of freedom of the apparatus. During the measurement, the total wavefunctional, of the system and apparatus, then evolves as
$$\mathrm{\Psi }^{(s)}(\varphi )\chi ^{(a)}(\stackrel{~}{\varphi })\underset{i}{}\mathrm{\Psi }_i^{(s)}(\varphi )\chi _i^{(a)}(\stackrel{~}{\varphi }).$$
(4.237)
If the different terms $`\mathrm{\Psi }_i^{(s)}(\varphi )\chi _i^{(a)}(\stackrel{~}{\varphi })`$ are non-overlapping in the configuration space of fields $`(\varphi ,\stackrel{~}{\varphi })`$ and if they remain non-overlapping in the future, the field beables are effectively guided by one of the wavefunctionals, say $`\mathrm{\Psi }_k^{(s)}(\varphi )\chi _k^{(a)}(\stackrel{~}{\varphi })`$. The empty wavefunctionals can then be dismissed from the future description of the field beables and we have an effective collapse. This is completely analogous to the situation in non-relativistic quantum theory. There is, however, one issue that needs to be addressed. In non-relativistic quantum theory it was guaranteed, for a general measurement situation, that different terms in a macroscopic superposition were non-overlapping in the configuration space, because the different macroscopic states generally correspond to systems localized at different regions in physical space (you can for example think of states which correspond to a macroscopic needle pointing in different directions). This was straightforward to show. Now with a field ontology this is not so straightforward anymore. Are states corresponding corresponding to macroscopic systems non-overlapping in the configuration space of fields?
A natural approach would be to consider quantum states which correspond to systems which are localized at distinct regions in physical 3-space and to look whether these states are non-overlapping in the configuration space of fields. Valentini addressed this question for non-relativistic one-particle states . It is interesting to consider his analysis here.
Valentini considered the real Klein-Gordon field. By letting the Klein-Gordon field operator $`\widehat{\varphi }(𝐱)`$ act on the ground state $`|0`$, one-particle states $`|𝐱\widehat{\varphi }(𝐱)|0`$ are constructed. These states obey the Klein-Gordon equation (because the field operator $`\widehat{\varphi }(𝐱)`$ obeys the Klein-Gordon equation). But only in the non-relativistic limit these states represent strictly localized particles. This is because only in the non-relativistic limit the different states $`|𝐱`$ become orthogonal. In the field basis $`|\varphi `$, we have $`\varphi |𝐱=\varphi (𝐱)\varphi |0`$ with $`\varphi |0`$ the wavefunctional of the vacuum.
If we now consider a low energy state $`|\mathrm{\Psi }`$, which contains only one particle, then we can expand the corresponding wavefunctional as
$$\mathrm{\Psi }(\varphi )=\varphi |\mathrm{\Psi }=d^3x\varphi |𝐱𝐱|\mathrm{\Psi }=\varphi |0d^3x\varphi (𝐱)\mathrm{\Psi }(𝐱),$$
(4.238)
with $`\mathrm{\Psi }(𝐱)=𝐱|\mathrm{\Psi }`$ an amplitude which obeys the non-relativistic Schrödinger equation. So, a non-relativistic particle which would be described by the wavefunction $`\mathrm{\Psi }(𝐱)`$ in non-relativistic quantum theory, is described by the wavefunctional $`\mathrm{\Psi }(\varphi )`$ given in (4.238) in quantum field theory.
Valentini showed that the probability density $`|\mathrm{\Psi }(\varphi )|^2`$ reaches its maximum for a field $`\varphi (𝐱)`$ which mimics the non-relativistic wavefunction $`\mathrm{\Psi }(𝐱)`$. For example if $`\mathrm{\Psi }(𝐱)`$ happens to be real then $`\varphi (𝐱)`$ is proportional to $`\mathrm{\Psi }(𝐱)`$. Suppose now that the probability density $`|\mathrm{\Psi }(\varphi )|^2`$ is sharply peaked around the field configuration $`\varphi (𝐱)=\mathrm{\Psi }(𝐱)`$. Let us further consider two wavefunctionals $`\mathrm{\Psi }_1(\varphi )`$ and $`\mathrm{\Psi }_2(\varphi )`$ which describe a non-relativistic particle localized at different regions in physical 3-space, i.e.
$$\mathrm{\Psi }_i(\varphi )=\varphi |0d^3x\varphi (𝐱)\mathrm{\Psi }_i(𝐱),i=1,2,$$
(4.239)
where the non-relativistic wavefunctions $`\mathrm{\Psi }_1(𝐱)`$ and $`\mathrm{\Psi }_2(𝐱)`$ have a support at distinct regions in physical 3-space. Now, with the assumption that the probability densities $`|\mathrm{\Psi }_i(\varphi )|^2`$ are sharply peaked around the field configurations $`\varphi (𝐱)=\mathrm{\Psi }_i(𝐱)`$ it is clear that the wavefunctionals $`\mathrm{\Psi }_i(\varphi )`$ are non-overlapping in the configuration space of fields.
So, if it could be shown that the probability densities $`|\mathrm{\Psi }_i(\varphi )|^2`$ are sharply peaked around the field with maximum probability, wavefunctionals which correspond to particles which are localized at different regions in physical 3-space would be non-overlapping in the configuration space of fields. And under these circumstances it is to be expected that wavefunctionals corresponding to different macroscopic systems, such as macroscopic pointer needles, are also non-overlapping. However, this need not to be the case at all. I.e. the different wavefunctionals $`\chi _i^{(a)}`$ of the apparatus need not to correspond with to different configurations in physical 3-space in order to be non-overlapping. It is in for example sufficient to have non-overlap if the wavefunctionals $`\chi _i`$ are peaked at different values for the electromagnetic field. We think it might be interesting to consider coherent states for the electromagnetic field in this respect.
## Chapter 5 Field beables for fermionic quantum field theory
### 5.1 Introduction
Little work has yet appeared on the construction of a pilot-wave interpretation for fermionic field theory with fields as beables. Bohm, Hiley and Kaloyerou argued that the anti-commutation relations of fermionic fields do not permit a pilot-wave approach in which the beables are continuous fields . According to their view, the field beable interpretation should only be adopted for bosons. For fermions Bohm et al. prefer the particle beable approach. However, in the previous chapter, we gave arguments why we do not favour this approach. One of the main reasons is that the model of Bohm et al. requires the notion of a Dirac sea, which makes the model only suitable for quantum electrodynamics.
In 1988, shortly after Bohm et al. argued against a field beable approach to fermionic quantum fields, Holland presented such a model . Holland presented his model for the quantized non-relativistic Schrödinger field, but the model can be straightforwardly extended to any fermionic field theory. The beables in Holland’s model are the Euler angles at each point in momentum space.
Later, in 1992, Valentini presented a pilot-wave interpretation for the quantized Van der Waerden theory (which describes relativistic spin-$`1/2`$ fields) . In Valentini’s model the beables are anti-commuting fields, also called Grassmann fields. It is on this last approach by Valentini that we will focus in this chapter. We will argue that this approach is untenable. We will see that it is no problem to write a fermionic field theory in the functional Schrödinger picture. However, the quantity which would be identified as the probability of the Grassmann fields, is itself an element of the Grassmann algebra and hence cannot be interpreted as a probability. In addition, it is hard to introduce meaningful guidance equations for Grassmann fields.
Because Valentini’s pilot-wave approach to fermionic fields is often quoted as a valid alternative (for more than ten years by now) we devote a chapter to it, explaining in detail where the problems are situated. Instead of using the quantized Dirac theory or the equivalent Van der Waerden theory to treat relativistic spin-$`1/2`$ fields, we will start with the quantized non-relativistic Schrödinger theory. Apart from notational simplicity, it has the additional advantage that it can be quantized using both Fermi-Dirac statistics and Bose-Einstein statistics. We can then clearly indicate the analogies and differences between bosonic and fermionic quantization. Our approach to introduce a pilot-wave model with Grassmann beables slightly differs from Valentini’s original approach (nevertheless we face the same problems). In Section 5.4.4, we compare our approach to Valentini’s one. Although the model of Holland hence seems to be the only alternative for a field beable approach for the moment, we do not elaborate on this model.
In the previous chapter we also expressed the opinion that a field beable approach to quantum field theory seems the most natural approach, but this by no means excludes the possibility of a particle beable approach. In fact such an approach was even presented by Bell who treated fermionic field theory on a lattice . Later Dürr et al. constructed a continuum version and applied it to quantum electrodynamics (where particle beables are only introduced for the fermions), see and references therein. These models differ form ‘ordinary’ pilot-wave models in the sense that they also include an element of stochasticity. These models were therefore termed Bell-type models by Dürr et al. We do not consider these Bell-type models further either in this thesis.
There is also a particle model by Colin , who also tried to find a continuum version of the Bell model. Colin ended up with a deterministic pilot-wave type model, instead of a stochastic one. This model is in fact the same as the one originally presented by Bohm for the Dirac equation, in which beables are introduced for all the particles in the Dirac sea. Contrary to the model of Dürr et al., Colin’s model relies hence on the existence of a Dirac sea and therefore this model might perhaps not be extendible to other type of interactions, such as weak interaction.
### 5.2 The quantized non-relativistic Schrödinger theory
The Lagrangian for the non-relativistic Schrödinger theory in one spatial dimension reads
$$L=𝑑x\left(\frac{i\mathrm{}}{2}\left(\psi ^{}_t\psi \psi _t\psi ^{}\right)\frac{\mathrm{}^2}{2m}_x\psi ^{}_x\psi \right).$$
(5.1)
Dirac’s method of quantization is simple in this case (one only encounters second class constraints) and the resulting quantum field theory can be found in many textbooks. Therefore there is no need to repeat the analysis here. We directly present the well-known expression for the Hamiltonian operator
$$\widehat{H}=\frac{\mathrm{}^2}{2m}𝑑x\widehat{\psi }^{}_x^2\widehat{\psi }$$
(5.2)
and the bosonic and fermionic commutation relations for the field operators $`\widehat{\psi }`$ and $`\widehat{\psi }^{}`$:
$$[\widehat{\psi }(x),\widehat{\psi }^{}(y)]_\pm =\delta (xy),[\widehat{\psi }(x),\widehat{\psi }(y)]_\pm =[\widehat{\psi }^{}(x),\widehat{\psi }^{}(y)]_\pm =0.$$
(5.3)
Here $`[.,.]_{}`$ denotes the commutator (bosonic quantization) and $`[.,.]_+`$ the anti-commutator (fermionic quantization). For notational convenience, we write the theory in terms of Fourier components
$`\widehat{\psi }(x)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2\pi }}}{\displaystyle 𝑑k\widehat{a}(k)e^{ikx}},`$
$`\widehat{\psi }^{}(x)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2\pi }}}{\displaystyle 𝑑k\widehat{a}^{}(k)e^{ikx}}.`$ (5.4)
The Hamiltonian operator (5.2) then becomes
$$\widehat{H}=𝑑kE(k)\widehat{a}^{}(k)\widehat{a}(k)$$
(5.5)
with $`E(k)=\mathrm{}^2k^2/2m`$ and the commutation relations (5.3) then imply
$$[\widehat{a}(k),\widehat{a}^{}(k^{})]_\pm =\delta (kk^{}),[\widehat{a}(k),\widehat{a}(k^{})]_\pm =[\widehat{a}^{}(k),\widehat{a}^{}(k^{})]_\pm =0.$$
(5.6)
We now deal with the bosonic and fermionic case separately.
### 5.3 Bosonic quantization of the non-relativistic <br>Schrödinger equation
As shown in the previous chapter, the construction of a pilot-wave interpretation presents no difficulty in the case of bosonic quantization. The quantized non-relativistic Schrödinger field with Bose-Einstein statistics was already treated by Holland . It is this model that we recall here, the only difference is that we use take a continuous momentum space instead of a discretized one.
The bosonic commutation relations $`[\widehat{a}(k),\widehat{a}^{}(k^{})]_{}=\delta (kk^{})`$ can be realized with the representation
$`\widehat{a}(k)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(q(k)+{\displaystyle \frac{\delta }{\delta q(k)}}\right),`$
$`\widehat{a}^{}(k)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(q(k){\displaystyle \frac{\delta }{\delta q(k)}}\right).`$ (5.7)
By substituting these relations for $`\widehat{a}(k)`$ and $`\widehat{a}^{}(k)`$ in the Hamiltonian (5.5) we obtain<sup>1</sup><sup>1</sup>1The irrelevant infinite $`c`$-number term could be omitted in the Hamiltonian.
$$H=\frac{1}{2}𝑑kE(k)\left(\frac{\delta ^2}{\delta q^2(k)}+q^2(k)\delta (0)\right).$$
(5.8)
The Schrödinger equation for the wavefunctional $`\mathrm{\Psi }(q(k),t)`$ (which is complex valued) then reads
$$i\mathrm{}\frac{\mathrm{\Psi }}{t}=\frac{1}{2}𝑑kE(k)\left(\frac{\delta ^2}{\delta q^2(k)}+q^2(k)\delta (0)\right)\mathrm{\Psi }.$$
(5.9)
The corresponding conservation equation is
$`{\displaystyle \frac{|\mathrm{\Psi }|^2}{t}}`$ $`+`$ $`{\displaystyle 𝑑k\frac{\delta J_q(k)}{\delta q(k)}}=0,`$
$`J_q(k)`$ $`=`$ $`{\displaystyle \frac{E(k)}{2i\mathrm{}}}\left(\mathrm{\Psi }^{}{\displaystyle \frac{\delta }{\delta q(k)}}\mathrm{\Psi }\mathrm{\Psi }{\displaystyle \frac{\delta }{\delta q(k)}}\mathrm{\Psi }^{}\right).`$ (5.10)
In the pilot-wave interpretation the density of the fields $`q(k)`$ is given by $`|\mathrm{\Psi }|^2`$ and the guidance equation reads
$$\dot{q}(k)=J_q(k)/|\mathrm{\Psi }|^2.$$
(5.11)
### 5.4 Fermionic quantization of the non-relativistic Schrödinger equation
#### 5.4.1 Functional Schrödinger representation
Let us now try to apply the same scheme in the fermionic case. The first thing to do, is to find a representation for $`\widehat{a}(k)`$ and $`\widehat{a}^{}(k)`$ in terms of certain variables and differential operators with respect to these variables, such that the anti-commutation relations $`[\widehat{a}(k),\widehat{a}^{}(k^{})]_+=\delta (kk^{})`$ are satisfied. With such a representation, we can then obtain a functional Schrödinger equation. A possible representation is the one in terms of Euler angles which was used by Holland . Other possible representations which are more commonly used nowadays are written in terms of Grassmann numbers . In this thesis we only focus on this representation.
For each wavenumber $`k`$, we introduce a Grassmann number<sup>2</sup><sup>2</sup>2For the definitions and properties concerning Grassmann numbers we refer to Appendix C. $`\eta (k)`$ and its conjugate $`\eta ^{}(k)`$, which satisfy
$$[\eta (k),\eta (l)]_+=[\eta (k),\eta ^{}(l)]_+=[\eta ^{}(k),\eta ^{}(l)]_+=0.$$
(5.12)
The anti-commutation relations for the creation and annihilation operators are then realized in the representation <sup>3</sup><sup>3</sup>3Other representations can be used as well. An alternative representation is for example $`\widehat{a}(k)=\eta (k)`$, $`a^{}(k)=\stackrel{}{\delta }/\delta \eta (k)`$ . Note that this representation is one-to-one, whereas the representation (5.13) is not. However, because we will encounter more severe obstructions when trying to formulate a pilot-wave interpretation using a Grassmann representation, we will not further consider this fact.
$`\widehat{a}(k)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(\eta (k)+{\displaystyle \frac{\stackrel{}{\delta }}{\delta \eta ^{}(k)}}\right),`$
$`\widehat{a}^{}(k)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(\eta ^{}(k)+{\displaystyle \frac{\stackrel{}{\delta }}{\delta \eta (k)}}\right).`$ (5.13)
By substituting these relations in the second quantized Hamiltonian (5.5), we obtain the Hamiltonian
$$H=\frac{1}{2}𝑑kE(k)(\frac{\stackrel{}{\delta }}{\delta \eta (k)}\frac{\stackrel{}{\delta }}{\delta \eta ^{}(k)}+\eta ^{}(k)\frac{\stackrel{}{\delta }}{\delta \eta ^{}(k)}\eta (k)\frac{\stackrel{}{\delta }}{\delta \eta (k)}+\eta ^{}(k)\eta (k)+\delta (0))$$
(5.14)
which leads to the following Schrödinger equation for the wavefunctional $`\mathrm{\Psi }(\eta ,\eta ^{},t)`$
$`i\mathrm{}{\displaystyle \frac{\mathrm{\Psi }}{t}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle }dkE(k)({\displaystyle \frac{\stackrel{}{\delta }}{\delta \eta (k)}}{\displaystyle \frac{\stackrel{}{\delta }}{\delta \eta ^{}(k)}}+`$ (5.15)
$`\eta ^{}(k){\displaystyle \frac{\stackrel{}{\delta }}{\delta \eta ^{}(k)}}\eta (k){\displaystyle \frac{\stackrel{}{\delta }}{\delta \eta (k)}}+\eta ^{}(k)\eta (k)+\delta (0))\mathrm{\Psi }.`$
We want to stress the fact that the wavefunctional is an element of the Grassmann algebra with generators $`\eta (k)`$ and $`\eta ^{}(k)`$, and that hence the wavefunctional is not a complex valued functional. This fact is a direct consequence of the representation in terms of Grassmann numbers for the creation and annihilation operators.
The conventional inner product of two Grassmann valued functionals $`\mathrm{\Psi }_1`$ and $`\mathrm{\Psi }_2`$ is defined by <sup>4</sup><sup>4</sup>4In the literature other definitions can be encountered, e.g. Hallin and Liljenberg use a Grassmann valued inner product instead of a complex valued inner product. However, all these definitions boil down to the same expressions for probability amplitudes at the end.
$$\mathrm{\Psi }_1|\mathrm{\Psi }_2=𝒟\eta ^{}𝒟\eta \mathrm{\Psi }_1^{}\mathrm{\Psi }_2=\mathrm{\Psi }_2|\mathrm{\Psi }_1^{}$$
(5.16)
with $`\mathrm{\Psi }^{}`$ the dual of $`\mathrm{\Psi }`$ (and not the Hermitian conjugate of $`\mathrm{\Psi }`$ which is denoted by a dagger) given by
$$\mathrm{\Psi }^{}(\eta ,\eta ^{},t)=𝒟\overline{\eta }^{}𝒟\overline{\eta }\mathrm{exp}(\overline{\eta }\eta ^{}+\overline{\eta }^{}\eta )\mathrm{\Psi }^{}(\overline{\eta },\overline{\eta }^{},t),$$
(5.17)
with $`𝒟\overline{\eta }=_kd\overline{\eta }(k)`$ and $`\mathrm{\Psi }^{}`$ the Hermitian conjugate of $`\mathrm{\Psi }`$. We used the notation $`\overline{\eta }\eta ^{}=𝑑k\overline{\eta }(k)\eta ^{}(k)`$.
At this stage, it is important not to confuse $`\psi (\eta ,\eta ^{},t)`$ with $`\eta \eta ^{}|\psi (t)`$. The wavefunctional $`\psi (\eta ,\eta ^{},t)`$ is an element of the Grassmann algebra (in order to satisfy the Schrödinger equation (5.15)) and because $`\eta \eta ^{}|\psi (t)`$ is an inner product it is a complex number. This difference was also stressed in .
#### 5.4.2 Problem to identify a suitable density of field beables
We can now turn to the question whether this functional Schrödinger picture admits for a pilot-wave interpretation. In order to identify a candidate for the probability of the beables we can consider the norm of a wavefunctional. Such an approach is likely to guarantee equivalence between standard fermionic quantum field theory and a possible pilot-wave interpretation. The norm of a wavefunctional reads
$`\mathrm{\Psi }|\mathrm{\Psi }`$ $`=`$ $`{\displaystyle 𝒟\eta ^{}𝒟\eta \mathrm{\Psi }^{}\mathrm{\Psi }}`$ (5.18)
$`=`$ $`{\displaystyle 𝒟\eta ^{}𝒟\eta \frac{1}{2}\left(\mathrm{\Psi }^{}\mathrm{\Psi }+\mathrm{\Psi }^{}(\mathrm{\Psi }^{})^{}\right)}`$
$`=`$ $`{\displaystyle 𝒟\eta ^{}𝒟\eta P(\eta ,\eta ^{},t)}`$
where
$$P(\eta ,\eta ^{},t)=\frac{1}{2}(\mathrm{\Psi }^{}\mathrm{\Psi }+\mathrm{\Psi }^{}(\mathrm{\Psi }^{})^{})=P(\eta ,\eta ^{},t)^{}.$$
(5.19)
Hence, $`P(\eta ,\eta ^{},t)`$ would be the most natural candidate for the probability density. However, $`P`$ is an element of the Grassmann algebra and hence not a real positive number. This means that $`P(\eta ,\eta ^{},t)`$ can not be interpreted as a probability density of field beables. This is a first problem we encounter when we try to construct a pilot-wave interpretation. There is no clear candidate for the probability density of field beables.
The core of the problem is that the configuration space of Grassmann fields is trivial, it consist of just one configuration $`(\eta (k),\eta ^{}(k))`$. The wavefunctional $`\mathrm{\Psi }(\eta ,\eta ^{},t)`$ is a mapping this one configuration $`(\eta (k),\eta ^{}(k))`$ to the Grassmann algebra. This situation has to be contrasted with the situation in the bosonic case. In the bosonic case the configuration space of fields is the space of smooth functions which consists hence of more than one configuration. The wavefunctional $`\mathrm{\Psi }(q(k),t)`$ for a bosonic system is a mapping from this configuration space of smooth functions $`(q(k))`$ to the complex numbers. From this point of view it is clear that is futile to introduce a probability distribution on the the configuration space of Grassmann fields.
For the same reason it is unclear how to introduce a meaningful guidance equation for the Grassmann fields. Because the configuration space of Grassmann fields consists only of one configuration, trajectories in this configuration space are trivial.
#### 5.4.3 Problem to construct a well defined guidance equation
Although we already anticipated the problems in constructing a guidance equation for the Grassmann fields in the previous section, it is still instructive to make an explicit attempt. For this purpose we will treat the quantity $`P`$ formally as a density of Grassmann field beables. As usual we try to identify a guidance equation by considering the continuity equation for $`P`$.
After a rather tedious calculation, one can obtain the following conservation equation for $`\mathrm{\Psi }^{}(\eta ,\eta ^{},t)\mathrm{\Psi }(\eta ,\eta ^{},t)`$:
$`{\displaystyle \frac{\mathrm{\Psi }^{}\mathrm{\Psi }}{t}}`$ $`+`$ $`{\displaystyle 𝑑k\left(\frac{I_\eta (k)\stackrel{}{\delta }}{\delta \eta (k)}+\frac{\stackrel{}{\delta }I_\eta ^{}(k)}{\delta \eta ^{}(k)}\right)}=0,`$
$`I_\eta (k)`$ $`=`$ $`{\displaystyle \frac{iE(k)}{2\mathrm{}}}\left(\mathrm{\Psi }^{}\mathrm{\Psi }\eta (k){\displaystyle \frac{1}{2}}{\displaystyle \frac{\stackrel{}{\delta }\overline{\mathrm{\Psi }}^{}}{\delta \eta ^{}(k)}}\overline{\mathrm{\Psi }}+{\displaystyle \frac{1}{2}}\mathrm{\Psi }^{}{\displaystyle \frac{\stackrel{}{\delta }\overline{\mathrm{\Psi }}}{\delta \eta ^{}(k)}}\right),`$
$`I_\eta ^{}(k)`$ $`=`$ $`{\displaystyle \frac{iE(k)}{2\mathrm{}}}\left(\eta ^{}(k)\mathrm{\Psi }^{}\mathrm{\Psi }{\displaystyle \frac{1}{2}}{\displaystyle \frac{\overline{\mathrm{\Psi }}^{}\stackrel{}{\delta }}{\delta \eta (k)}}\mathrm{\Psi }+{\displaystyle \frac{1}{2}}\overline{\mathrm{\Psi }}^{}{\displaystyle \frac{\overline{\mathrm{\Psi }}\stackrel{}{\delta }}{\delta \eta (k)}}\right),`$ (5.20)
where $`\overline{\mathrm{\Psi }}=\mathrm{\Psi }_e\mathrm{\Psi }_o`$ with $`\mathrm{\Psi }_e`$ and $`\mathrm{\Psi }_o`$ respectively the even part and the odd part of $`\mathrm{\Psi }=\mathrm{\Psi }_e+\mathrm{\Psi }_o`$. Hence the conservation equation for $`P(\eta ,\eta ^{},t)`$ is
$$\frac{P(\eta ,\eta ^{},t)}{t}+𝑑k\left(\frac{\stackrel{}{\delta }}{\delta \eta ^{}(k)}\left(I_{\eta (k)}^{}+I_{\eta ^{}(k)}\right)+\left(I_{\eta (k)}+I_{\eta ^{}(k)}^{}\right)\frac{\stackrel{}{\delta }}{\delta \eta (k)}\right)=0.$$
(5.21)
With $`P`$ formally interpreted as a ‘probability density’, the guidance equations for the fields $`\eta (k),\eta ^{}(k)`$ should look something like
$`\dot{\eta }(k)`$ $`=`$ $`P(\eta ,\eta ^{},t)^1\left(I_{\eta (k)}+I_{\eta ^{}(k)}^{}\right),`$
$`\dot{\eta }^{}(k)`$ $`=`$ $`\left(I_{\eta (k)}^{}+I_{\eta ^{}(k)}\right)P(\eta ,\eta ^{},t)^1.`$ (5.22)
In order for these guidance equations to be well defined, we have to address some problems:
1. Not every element of the Grassmann algebra has an inverse for the multiplication (see Appendix C) and hence we must ensure that $`P^1`$ is always well defined. This is the case if and only if $`\mathrm{\Psi }^1`$ is well defined. For example the ground state of the Schrödinger wave equation $`\mathrm{\Psi }_v`$ is given by
$$\mathrm{\Psi }_v=N\mathrm{exp}\left(𝑑k\eta (k)\eta ^{}(k)\right).$$
(5.23)
Because $`𝒟\eta ^{}𝒟\eta _k\left(\eta (k)\eta ^{}(k)\right)\mathrm{\Psi }_v=N0`$ the ground state has an inverse (cf. Appendix C). This inverse is given by
$$\mathrm{\Psi }_v^1=N^1\mathrm{exp}\left(𝑑k\eta (k)\eta ^{}(k)\right).$$
(5.24)
The first excited state state, describing one particle with energy $`E_l`$, has the form
$$\mathrm{\Psi }_l=a_l^{}\mathrm{\Psi }_ve^{iE_lt/\mathrm{}}=\sqrt{2}\eta _l^{}\mathrm{\Psi }_ve^{iE_lt/\mathrm{}}.$$
(5.25)
But because this state is proportional to $`\eta _l^{}`$ it has no inverse. Further excited states will exhibit the same problem.
A possible way to regularize these excited states is by taking a suitable superposition with the ground state. Consider for example the first excited state $`\mathrm{\Psi }_l`$, which we superpose with the ground state as follows
$$\mathrm{\Psi }_l^ϵ=\sqrt{1ϵ}\mathrm{\Psi }_l+\sqrt{ϵ}\mathrm{\Psi }_v.$$
(5.26)
The inverse of the state $`\mathrm{\Psi }_l^ϵ`$ is now well defined because
$$𝒟\eta ^{}𝒟\eta \underset{k}{}\eta (k)\eta ^{}(k)\mathrm{\Psi }_l^ϵ=\sqrt{ϵ}𝒟\eta ^{}𝒟\eta \underset{k}{}\eta (k)\eta ^{}(k)\mathrm{\Psi }_v=\sqrt{ϵ}0.$$
(5.27)
In the limit $`ϵ0`$ the state $`\mathrm{\Psi }_l^ϵ`$ will approach $`\mathrm{\Psi }_l`$. Because $`\mathrm{\Psi }_l`$ and $`\mathrm{\Psi }_v`$ are orthogonal with respect to the inner product (5.16), the probability of finding the state $`\mathrm{\Psi }_l`$ is given by $`1ϵ`$ and the probability of finding the ground state is $`ϵ`$.
2. Because $`\dot{\eta }(k)`$ and $`\dot{\eta }^{}(k)`$ should be anti-commuting variables, the right hand sides of (5.22) should be odd elements of the Grassmann algebra. This can only be accomplished if the wavefunctional $`\mathrm{\Psi }`$ is an even element of the Grassmann algebra. This implies that $`\mathrm{\Psi }`$ is a superposition of wavefunctionals which describe an even number of particles. This is because the ground state of the Schrödinger wave equation $`\mathrm{\Psi }_v`$ is even. The state $`\mathrm{\Psi }_l`$ which is the first excited state with energy $`E_l`$ describes one particle and is an odd element of the Grassmann algebra. Generally, applying $`n`$ different creation operators on the ground state, multiplies the groundstate with $`n`$ Grassmann numbers. Hence, only states describing an even number of particles are even and only these states will lead to well defined guidance equations.
3. The third and most important problem arises when we try to make sense of the ‘guidance equations’ (5.22) as differential equations. This is basically due to the problem that the configuration space of Grassmann fields only contains one configuration $`(\eta (k),\eta ^{}(k))`$.
A possible way to make sense of the guidance equations as differential equations could be by introducing a Grassmann algebra which is generated by $`u(k),u^{}(k)`$, where there are as many $`u(k)`$’s as there are $`\eta (k)`$’s, and by expressing the $`\eta (k)`$’s as time dependent superpositions of odd elements of the basis:
$`\eta (k)`$ $`=`$ $`{\displaystyle 𝑑l_1f^{10}(k;l_1;t)u(l_1)}+{\displaystyle 𝑑l_1f^{01}(k;l_1;t)u^{}(l_1)}`$ (5.28)
$`+`$ $`{\displaystyle 𝑑l_1𝑑l_2𝑑l_3f^{21}(k;l_1,l_2,l_3;t)u(l_1)u(l_2)u^{}(l_3)}+\mathrm{}`$
where every $`f^{mn}(k;l_1,\mathrm{},l_{m+n})`$ is, for every $`k`$, a time dependent distributions. The guidance equations (5.22) are then well defined as differential equations for the coefficients $`f^{mn}`$. By this construction we have in fact introduced a nontrivial configuration space of fields. The fields are now the coefficients $`f^{mn}`$ which are distributions. But it is unclear how to proceed from this point. It is unclear what the probability distribution on the configuration space of fields $`f^{mn}`$ should look like.
In conclusion, we see that we encounter problems with the construction of a pilot-wave interpretation for the functional Schrödinger equation in the Grassmann representation. The main problems have to do with the fact that the configuration space of Grassmann fields consists only of a single configuration. Hence it makes no sense to introduce a probability density on this configuration space nor does it make sense to introduce dynamics on this configuration space.
#### 5.4.4 Relation to Valentini’s work
It was originally believed by Valentini that a pilot-wave interpretation could only be given for field theories for which the field equations contain second-order time derivatives, such as the Van der Waerden spin-$`1/2`$ equation . However, as shown in the preceding chapter, we could give a pilot-wave interpretation for the Duffin-Kemmer-Petiau theory which is first order in time. Also the non-relativistic Schrödinger equation is first-order in time and as shown above, the pilot-wave interpretation could be adopted equally well for this first-order theory when using the Bose-Einstein statistics. In fact, one can even show that the quantized Van der Waerden theory and the quantized Dirac theory can be transformed into one another by a canonical transformation and are hence equivalent (as should be expected because they are equivalent already on the first quantized level).
Valentini did not start with the inner product (5.16) either to arrive at the ‘probability density’. He considered the conservation equation for $`\mathrm{\Psi }^{}\mathrm{\Psi }`$, which has the following form in the case of the non-relativistic Schrödinger theory
$`{\displaystyle \frac{\mathrm{\Psi }^{}\mathrm{\Psi }}{t}}`$ $`+`$ $`{\displaystyle 𝑑k\left(\frac{C_\eta (k)\stackrel{}{\delta }}{\delta \eta (k)}+\frac{\stackrel{}{\delta }C_\eta ^{}(k)}{\delta \eta ^{}(k)}\right)}=0,`$
$`C_\eta (k)`$ $`=`$ $`{\displaystyle \frac{iE(k)}{2\mathrm{}}}\left(\mathrm{\Psi }^{}\mathrm{\Psi }\eta (k){\displaystyle \frac{\stackrel{}{\delta }\overline{\mathrm{\Psi }}^{}}{\delta \eta ^{}(k)}}\overline{\mathrm{\Psi }}+\mathrm{\Psi }^{}{\displaystyle \frac{\stackrel{}{\delta }\overline{\mathrm{\Psi }}}{\delta \eta ^{}(k)}}\right).`$ (5.29)
In the pilot-wave interpretation the conserved quantity $`\mathrm{\Psi }^{}\mathrm{\Psi }`$ would then be the probability density of the field beables and the guidance equations for these field beables would then read
$`\dot{\eta }(k)`$ $`=`$ $`(\mathrm{\Psi }^{}\mathrm{\Psi })^1C_\eta (k),`$
$`\dot{\eta }^{}(k)`$ $`=`$ $`C_\eta ^{}(k)(\mathrm{\Psi }^{}\mathrm{\Psi })^1.`$ (5.30)
However, apart from the fact that $`\mathrm{\Psi }^{}\mathrm{\Psi }`$ is also an element of the Grassmann algebra, it is in general not normalizable. Integrating $`\mathrm{\Psi }^{}\mathrm{\Psi }`$ over the configuration space of Grassmann fields can yield zero for a nonzero wavefunctional. Even the quantity $`\mathrm{\Psi }^{}\mathrm{\Psi }`$ can be zero for a nonzero wavefunctional. These problems are not present for the quantity $`P`$ defined in (5.19). This is also the reason why the inner product was defined as (5.16).
## Chapter 6 On Peres’ statement “opposite momenta lead to opposite directions”, decaying systems and optical imaging
### 6.1 Introduction
In this chapter we will consider two things.<sup>1</sup><sup>1</sup>1The results of this chapter are published in . First, we consider the question to which extent opposite momenta lead to opposite directions for the fragments of a decaying quantum system. We hereby improve an analysis by Peres . According to Peres, there are only two sources for deviation from perfect angular alignment for a two-particle system with total momentum zero. We will argue that there is also another contribution to the deviation from angular alignment, which is due to the uncertainty of the location of the source. It will appear that Peres’ estimation for the angular deviation only applies in, what we will call, the large time or large distance regime. This is the regime where the two particles have traveled a large distance from the source. In the small time or small distance regime, the other contribution to the deviation becomes dominant. Peres applied his analysis to two different experiments , the thought experiment of Popper and the optical imaging experiment reported by Pittman et al. . We will argue that Popper’s thought experiment occurs in the small time regime and hence Peres’ analysis should not be applied there. On the other hand, Peres’ analysis can be correctly applied to the experiment of Pittman et al., which can be seen as occurring in the large time regime.
Second, we will reconsider the experiment by Pittman et al., from another point of view, namely pilot-wave theory. It is clear that one of the main merits of the pilot-wave interpretation is that it provides an observer independent description of quantum phenomena. Another merit of the pilot-wave interpretation is that it can be used to visualize quantum processes (just consider the many examples in Holland’s book ). Now, the paper by Pittman et al. contains drawings of conceptual photon trajectories. These trajectories, which are derived with ‘usual’ geometrical optics, are not the real paths of the photons (because, according to quantum theory, photons do not exist as localized entities between two measurements), but merely serve as a tool to visualize the experiment. However, when we calculate the trajectories predicted by pilot-wave theory they coincide with these conceptual trajectories. In this way the trajectories predicted by pilot-wave theory can serve as a theoretical basis for otherwise rather ad hoc drawings. Note however that we derive the trajectories for the massive particle equivalent of the experiment of Pittman et al.; it would take us too far to derive them for example from the Harish-Chandra theory (which is equivalent to Maxwell’s theory).<sup>2</sup><sup>2</sup>2Although the trajectory model that can be derived for the Harish-Chandra theory cannot serve as a valid interpretation for photons, it serves well for illustrative purposes as indicated in Chapter 3.
In the following section we start with recalling Popper’s thought experiment and the optical imaging experiment of Pittman et al.. In Section 6.3 we consider the question to which extent ‘opposite momenta’ lead to ‘opposite directions’. In section 6.4 we first consider a simplified pilot-wave description of a decaying system and then in Section 6.5 we consider the pilot-wave description of the massive particle equivalent to the experiment of Pittman et al.
### 6.2 On Popper’s experiment
In 1934, Popper proposed his experiment which aimed to test the general validity of quantum mechanics . Popper assumes a source $`S`$ from which pairs of particles are emitted in opposite directions. Two observers, say Alice and Bob, are located at opposite sides of the source, both equipped with an array of detectors. If Alice puts a screen with a slit in her way of the particles, she will observe a diffraction pattern behind the screen. According to Popper, quantum mechanics will also predict a diffraction pattern on the other side of the source, where Bob is located, when coincidence counts are considered. This is because every measurement by Alice is in fact a virtual position measurement of the correlated particle on Bob’s side, leading to an increased momentum uncertainty for Bob’s particle as well. This is the same diffraction pattern that would be observed when a physical slit was placed on Bob’s side. Popper, who declared himself a metaphysical realist, found this idea of ‘virtual scattering’ absurd and predicted no increased momentum uncertainty for Bob’s measurement due to Alice’s position measurement. He therefore saw his proposed experiment as a possible test against quantum mechanics and in favor of his realist vision in which particles have at each time well defined positions and momenta; particles for which the Heisenberg uncertainty, for example, is only a lower, statistical limit of scatter.
Unfortunately, to describe the setup of the experiment, Popper occasionally invoked classical language, which veiled some severe problems which could obstruct a practical realization of his experiment (for an extensive discussion see Peres ). For example, Popper writes: “We have a source $`S`$ (positronium, say) from which pairs of particles that have interacted are emitted in opposite directions. We consider pairs of particles that move in opposite directions …”. It is the validity of this statement, which appears to be a very delicate issue, that is one of the topics we shall deal with in this chapter.
Of course, when we consider a decaying system at rest in classical mechanics, the fragments will have opposite momenta
$$𝐩_1+𝐩_2=\mathrm{𝟎}$$
(6.1)
and if we take the place of decay of the system as the centre of our coordinate system the positions of the two fragments will satisfy $`m_1𝐱_1+m_2𝐱_2=\mathrm{𝟎}`$, so the fragments will be found in opposite, isotropically distributed directions. If the fragments have equal masses, then they will be found at opposite places, relative to the centre of the coordinate system.
But these properties do not hold in quantum mechanics. Suppose that a system has a two-particle wavefunction $`\psi `$ which has a sharp distribution at $`𝐩_1+𝐩_2=\mathrm{𝟎}`$, i.e. both $`|\widehat{p}_{1j}+\widehat{p}_{2j}|`$ and $`\mathrm{\Delta }(\widehat{p}_{1j}+\widehat{p}_{2j})`$ are small for every component $`j`$ of the momentum vectors $`\widehat{𝐩}_1+\widehat{𝐩}_2`$. Then according to the uncertainty relations
$$\mathrm{\Delta }(\widehat{p}_{1j}+\widehat{p}_{2j})\mathrm{\Delta }(m_1\widehat{x}_{1j}+m_2\widehat{x}_{2j})\frac{\mathrm{}}{2}(m_1+m_2),$$
(6.2)
the distribution of $`m_1x_{1j}+m_2x_{2j}`$ will be broad for every component $`j`$. In the case of equal masses, the inequalities in (6.2) imply that opposite momenta are incompatible with opposite positions. In particular, at the moment of decay, the inequalities imply that opposite momenta of the particles are incompatible with the latter being both located at the origin of the coordinate system. It was even shown by Collett and Loudon that this initial uncertainty on the location of the source implies that Popper’s experiment is inconclusive.
Although the original proposal of Popper’s experiment can hence not be performed practically, due to the fact that opposite momenta are incompatible with opposite positions, the intention of Popper’s proposal can be maintained if we have a two-particle system which displays some form of entanglement in the position coordinates.<sup>3</sup><sup>3</sup>3The position entanglement should not be exact, otherwise, as shown by Short , both of the observers would observe an infinite momentum spread. One is then in principle able to test experimentally whether one of the particles will experience an increased momentum spread due to a position measurement (within a slit width) of the correlated particle, i.e. we would be able to test a possible ‘virtual scattering’. Such a form of position entanglement was obtained with the phenomenon of optical imaging . By making use of optical imaging, Kim and Shih were able to perform an experiment in the spirit of Popper’s original proposal.
Let us briefly review this experiment by Pittman et al.. The experiment uses momentum correlated photons resulting from spontaneous parametric down conversion (see Fig. 6.1).
In this process a pump photon incident on a nonlinear beta barium borate (BBO) crystal leads to the creation of a signal and an idler photon. Due to momentum conservation, the sum of the momenta of these photons has to equal the momentum of the pump photon. This results in the momentum entanglement of the two photons, because the momenta of the idler and signal photon can be combined in an infinite number of ways to equal the momentum of the pump photon. In order to avoid a momentum spread due to spatial confinement, the width of the pump beam is sufficiently large. Hence, the exact place of creation within the BBO crystal is unknown. In the same way, the energies of the created photons add up to the energy of the pump photon. In the experiment, the signal and idler photons are sent in two different directions where coincidence records may be performed by two photon counting detectors. A convex lens, with focal length $`f`$, is placed in the signal beam in order to turn the momentum correlation of the created photons into spatial correlation. In front of the detector for the signal beam an aperture is placed at a distance $`S`$ from the lens. By placing the detector for the idler beam at a distance $`S^{}`$ from the lens, prescribed by the Gaussian thin lens equation, i.e.
$$\frac{1}{S}+\frac{1}{S^{}}=\frac{1}{f}$$
(6.3)
and scanning in the transverse plane of the idler beam, an image of this aperture is observed in the coincidence counts. This image obeys the classical lens equations in the following sense. If a classical point-like light source would be placed in the plane of the aperture, where the signal photon was detected, it would have an image where the idler photon was detected. This is the spatial correlation of the photons.
Finally, although we will not further deal with the experiment of Kim and Shih, it is interesting to note that unfortunately they failed in their original intention to perform Popper’s gedankenexperiment. As was shown by Short , the diameter of the incoming beam of pump photons was still to small to guarantee perfect momentum entanglement of the parametric down converted photons and this blurred the predicted results.
### 6.3 Opposite momenta and opposite directions
Peres gave an analysis of the extent to which opposite momenta lead to opposite directions . He argues that the inequalities in (6.2) do not exclude a priori the possibility that opposite momenta of particles lead to opposite directions (instead of opposite positions) where the particles will be found when a measurement is performed; on the contrary, the operator equivalent of (6.1) would even lead to an observable alignment of the detection points of the two particles. We will indicate that Peres’ analysis is correct in the large time regime, but in the small time regime there is an additional source for the deviation from perfect alignment which is not mentioned by Peres.
Let us recall Peres’ arguments. To discuss the possible angular alignment of momentum entangled particles, Peres considers a non-relativistic wavefunction describing massive particles. According to Peres, the reason for angular correlation of momentum entangled photons (as in the experiment of Pittman et al.) is the same as in the considered massive case.
The momentum correlated particles can be assumed to result from a decaying system at rest. The decaying system can then be described by the wavefunction
$$\psi (𝐱_1,𝐱_2,t)=d^3p_1d^3p_2F(𝐩_1,𝐩_2)e^{i(𝐩_1𝐱_1+𝐩_2𝐱_2Et)/\mathrm{}}$$
(6.4)
where the momentum distribution $`F`$ is peaked around $`𝐩_1+𝐩_2=\mathrm{𝟎}`$ and around the rest energy $`E_0`$ of the decaying system. According to Peres the opposite momenta of the particles lead to opposite directions where the particles will be found when a measurement is performed. He shows this by applying the stationary phase method. The main contribution to the integral in (6.4) comes from values $`𝐩_1`$ and $`𝐩_2`$ for which $`𝐩_1+𝐩_2\mathrm{𝟎}`$. Because of the rapid oscillations of the phase in the integrand in (6.4), the integral will be appreciably different from zero only if the phase is stationary with respect to the six integration variables $`𝐩_1`$ and $`𝐩_2`$ in the vicinity of $`𝐩_1+𝐩_2=\mathrm{𝟎}`$, i.e.
$$\frac{S}{𝐩_i}+𝐱_i\frac{E}{𝐩_i}t=\mathrm{𝟎},i=1,2$$
(6.5)
where $`S`$ is the phase of $`F(𝐩_1,𝐩_2)`$ measured in units of $`\mathrm{}`$, i.e.
$$F(𝐩_1,𝐩_2)=|F(𝐩_1,𝐩_2)|e^{iS(𝐩_1,𝐩_2)/\mathrm{}}$$
(6.6)
and the equations (6.5) have to be evaluated for $`𝐩_1+𝐩_2=\mathrm{𝟎}`$. The equations (6.5) then determine the conditions on $`𝐱_i`$ in order to have a non-zero $`\psi `$ (and $`|\psi |^2`$).
Peres then introduces spherical coordinates to describe $`𝐩_i`$ and $`𝐱_i`$, and varies the phase $`S`$ with respect to the six spherical variables of $`𝐩_i`$: ($`p_i=|𝐩_i|,\varphi _i,\theta _i)`$. Peres further assumes that the phase $`S`$ obeys
$$S/𝐩_i=\mathrm{𝟎}$$
(6.7)
in the vicinity of $`𝐩_1+𝐩_2=\mathrm{𝟎}`$. This would restrict the place of decay near the origin of the coordinate system, because of (6.5). By varying with respect to the momentum angles, Peres obtains that the phase is stationary if $`𝐩_i`$ and $`𝐱_i`$ have the same direction. Because $`F`$ is peaked at $`𝐩_1+𝐩_2=\mathrm{𝟎}`$, this results in
$`\theta _1^{}+\theta _2^{}`$ $`=`$ $`\pi `$
$`|\varphi _1^{}\varphi _2^{}|`$ $`=`$ $`\pi `$ (6.8)
where $`(x_i=|𝐱_i|,\varphi _i^{},\theta _i^{})`$ are the spherical coordinates of $`𝐱_i`$. These equations show that the two particles can only be detected at opposite directions relative to the centre of our coordinate system. This is because the integral in (6.4) (and hence $`|\psi |^2`$) would only be appreciably different from zero if the vectors $`𝐱_i`$ obey (6.8).
Peres mentions two causes for deviation from perfect angular alignment. The first is a transversal deviation of the order $`\sqrt{ht/m}`$ due to the spreading of the wavefunction, which was recognized as the standard quantum limit . The second is an angular spread of the order $`\mathrm{\Delta }(\widehat{p}_{1j}+\widehat{p}_{2j})/p_i`$. Below, we show that there is another cause for deviation which arises from the uncertainty on the source and which is particularly important in the ‘small distance’ or ‘small time’ regime.
First, we want to note that there is, apart from the conditions on $`\theta _i^{}`$ and $`\varphi _i^{}`$, also a condition on the variables $`x_i`$, which is not mentioned by Peres. This condition is obtained by varying the phase of the integrand with respect to $`p_i`$, having in mind the previous result that $`𝐩_i`$ and $`𝐱_i`$ have the same direction. If we define $`v_i=dE/dp_i`$, then the additional condition reads
$$x_i=v_it,$$
(6.9)
where the $`v_i`$ have to be evaluated for $`𝐩_1+𝐩_2=\mathrm{𝟎}`$. Because $`\psi `$ obeys the non-relativistic Schrödinger equation, $`E=\frac{p_1^2}{2m_1}+\frac{p_2^2}{2m_2}`$, with $`m_i`$ the masses of the particles. In this way (6.9) becomes
$$x_1=\frac{p_1}{m_1}t,x_2=\frac{p_2}{m_2}t.$$
(6.10)
Using $`𝐩_1+𝐩_2=\mathrm{𝟎}`$ one obtains
$$x_1m_1=x_2m_2.$$
(6.11)
Note that we have not yet used the fact that $`F`$ is peaked around a certain energy $`E_0`$, as is required in the case of a decaying system at rest. As Peres notes in his paper, a restriction of the energy to $`E_0`$ further restricts the momenta of the particles to satisfy $`p_1^2=p_2^2=2E_0m_1m_2/(m_1+m_2)`$. By combining (6.8) and (6.11), we obtain that the joint detection probability has a maximum for the classically expected relation $`m_1𝐱_1+m_2𝐱_2=\mathrm{𝟎}`$. Hence, Peres’ statement ‘opposite momenta lead to opposite directions’ may be replaced by a stronger statement, namely that the opposite momenta lead to a maximum detection probability for $`m_1𝐱_1+m_2𝐱_2=\mathrm{𝟎}`$.
Let us now consider the possible sources for deviation from this classical relation. Classically one can, in theory, make both quantities $`\mathrm{\Delta }(\widehat{p}_{1j}+\widehat{p}_{2j})`$ and $`\mathrm{\Delta }(m_1\widehat{x}_{1j}+m_2\widehat{x}_{2j})`$ as small as wanted. Quantum mechanically one can at best prepare the system, such that initially the equality in
$$\mathrm{\Delta }(\widehat{p}_{1j}+\widehat{p}_{2j})\mathrm{\Delta }(m_1\widehat{x}_{1j}+m_2\widehat{x}_{2j})\frac{\mathrm{}}{2}(m_1+m_2)$$
(6.12)
is reached. Note that this equation implies that in case of opposite momenta, the particles can not depart from a confined, fixed source.
Because the operator $`\widehat{p}_{1j}+\widehat{p}_{2j}`$ commutes with the free Hamiltonian, the variance of the momentum operator $`\widehat{p}_{1j}+\widehat{p}_{2j}`$ is stationary. The variance of $`\widehat{x}_{1j}+\widehat{x}_{2j}`$ however, will in general increase with time due to the spreading of the wavefunction. This can be seen if we write down the expression for the free evolution of the operator $`m_1\widehat{𝐱}_1+m_2\widehat{𝐱}_2`$ in the Heisenberg picture
$$m_1\widehat{𝐱}_1(t)+m_2\widehat{𝐱}_2(t)=m_1\widehat{𝐱}_1(0)+\widehat{𝐩}_1(0)t+m_2\widehat{𝐱}_2(0)+\widehat{𝐩}_2(0)t$$
(6.13)
The variance of this operator for an arbitrary component $`j`$ is
$`\mathrm{\Delta }\left(m_1\widehat{x}_{1j}(t)+m_2\widehat{x}_{2j}(t)\right)^2`$
$`=\mathrm{\Delta }\left(m_1\widehat{x}_{1j}(0)+m_2\widehat{x}_{2j}(0)\right)^2+\mathrm{\Delta }\left(\widehat{p}_{1j}(0)+\widehat{p}_{2j}(0)\right)^2t^2`$
$`+<[\left(m_1\widehat{x}_{1j}(0)+m_2\widehat{x}_{2j}(0)\right),\left(\widehat{p}_{1j}(0)+\widehat{p}_{2j}(0)\right)]_+>t`$
$`2<m_1\widehat{x}_{1j}(0)+m_2\widehat{x}_{2j}(0)><\widehat{p}_{1j}(0)+\widehat{p}_{2j}(0)>t`$ (6.14)
where the brackets $`[,]_+`$ denote the anti-commutator. If we assume a distribution $`F`$ which is real and symmetric, i.e. $`F(𝐩_1,𝐩_2)=F(𝐩_1,𝐩_2)`$ then the last two terms in (6.14) are both zero. So the variance of $`m_1\widehat{x}_{1j}(t)+m_2\widehat{x}_{2j}(t)`$ increases with time
$$\mathrm{\Delta }\left(m_1\widehat{x}_{1j}(t)+m_2\widehat{x}_{2j}(t)\right)^2=\mathrm{\Delta }\left(m_1\widehat{x}_{1j}(0)+m_2\widehat{x}_{2j}(0)\right)^2+\mathrm{\Delta }\left(\widehat{p}_{1j}(0)+\widehat{p}_{2j}(0)\right)^2t^2.$$
(6.15)
This leads to an increasing deviation from the relation $`m_1𝐱_1+m_2𝐱_2=\mathrm{𝟎}`$.
Thus there is always an interplay between opposite momenta and opposite directions which is expressed in (6.12) and (6.15). We argue that one should study (6.15), where the variances at $`t=0`$ in the right hand side of the expression are limited by the Heisenberg uncertainty in (6.12), to determine to which extent we can speak of possible angular alignment.
Let us now see how Peres’ estimates for deviation from perfect alignment come about and in which regime they are important. Assume for convenience that $`m_1=m_2=m`$. The transversal deviation $`L(t)`$ can be taken of the order $`\mathrm{\Delta }\left(\widehat{x}_{1j}(t)+\widehat{x}_{2j}(t)\right)`$. There are two contributions to this transversal deviation. The first is $`\mathrm{\Delta }\left(\widehat{x}_{1j}(0)+\widehat{x}_{2j}(0)\right)=L(0)`$ and is important for small times. The second contribution is $`\mathrm{\Delta }\left(\widehat{p}_{1j}+\widehat{p}_{2j}\right)t/m`$ which becomes important for larger times. The angular deviation $`\theta `$ may be derived from $`\mathrm{tan}(\theta )=L(t)/R(t)`$, where $`R(t)=pt/m`$ is the distance that both particles have traveled. For small times one has $`\mathrm{tan}(\theta )\mathrm{\Delta }\left(\widehat{x}_{1j}(0)+\widehat{x}_{2j}(0)\right)m/pt`$ and for large times one has $`\mathrm{tan}(\theta )\mathrm{\Delta }\left(\widehat{p}_{1j}+\widehat{p}_{2j}\right)/p`$. Hence, for large times we obtain the estimate of deviation mentioned by Peres. From the relations (6.12) and (6.15) one can also easily derive the standard quantum limit
$`\mathrm{\Delta }\left(\widehat{x}_{1j}(t)+\widehat{x}_{2j}(t)\right)^2`$ $``$ $`2\mathrm{\Delta }\left(\widehat{x}_{1j}(0)+\widehat{x}_{2j}(0)\right)\mathrm{\Delta }\left(\widehat{p}_{1j}(0)+\widehat{p}_{2j}(0)\right)t/m`$ (6.16)
$``$ $`2\mathrm{}t/m.`$
However, this uncertainty is misleading because for small times it neglects the contribution arising from the uncertainty on the source $`\mathrm{\Delta }\left(m_1\widehat{x}_{1j}(0)+m_2\widehat{x}_{2j}(0)\right)`$. Especially in the considered case of nearly opposite momenta, this contribution will be large because $`\mathrm{\Delta }\left(\widehat{p}_{1j}(0)+\widehat{p}_{2j}(0)\right)`$ is small.
In summary, we see that Peres gave causes for deviation which only apply in the large time regime. These causes are in perfect agreement with the ‘scattering into cones’ theorem which states that for every cone $`C`$ in $`\mathrm{}^m`$ with apex in the origin
$$\underset{t\mathrm{}}{lim}_Cd^mx|\psi (x,t)|^2=_Cd^mp|\varphi (p)|^2$$
(6.17)
with $`\varphi (p)`$ the momentum wave function . This means that the probability that in the infinite future the particles will be found in the cone $`C`$ is equal to the probability that their momenta lie in the same cone. In the small time regime the uncertainty on the source gives the major contribution to deviation.
Let us give a simple example. We can use a Gaussian distribution to represent the momentum correlation
$$F(𝐩_1,𝐩_2)e^{\frac{(𝐩_1+𝐩_2)^2}{\sigma }}.$$
(6.18)
The smaller the value of $`\sigma `$, the better the momentum correlation between the two fragments. In the limit $`\sigma 0`$ this distribution approaches the Dirac $`\delta `$-distribution. Note that this distribution is not peaked around a certain energy $`E_0`$ as should be required for a decaying system at rest. In Appendix D it is explained why we can leave this restriction on the energy aside without changing the main result. It will follow that a reasonable energy width, peaked around $`E_0`$, will imply only a minor broadening of the wavefunction. The wavefunction at $`t=0`$ is
$$\psi (𝐱_1,𝐱_2,0)\delta (𝐱_1𝐱_2)e^{𝐱_1^2\sigma /4\mathrm{}^2}.$$
(6.19)
Thus clearly $`\psi `$ represents a decaying system because initially $`𝐱_1=𝐱_2`$. But for small values of $`\sigma `$ (when $`F`$ is peaked around $`𝐩_1+𝐩_2=\mathrm{𝟎}`$) the probability of finding the particles at $`t=0`$ at some configuration is totally smeared out, though the probability has a maximum at the origin of the coordinate system. Note that although the relation $`S/𝐩_i=\mathrm{𝟎}`$ is satisfied, this condition does not restrict the place of decay near the origin of the coordinate system, as was assumed by Peres. So, in the small time regime, it may be hard to speak of possible opposite movements of particles relative to the origin because we cannot exactly say (at least without measurement) where the decay of the system took place. As time increases, the deviation from the detection probability peak at $`m_1𝐱_1+m_2𝐱_2=\mathrm{𝟎}`$ will even increase with time as was shown above. However, because the distance from the particles to the source increases, the uncertainty of the source will become less important for the angular deviation; in the large time regime the angular deviation will then be dominated by the momentum uncertainty.
In the following section, we will show that we can retain the classical picture of a decaying system in the previous example, in both the small time and the large time regime, when it is described by pilot-wave theory. For example, in the case of the wavefunction considered above, the particle beables will depart near each other and will move along opposite directions. However, the place of departure will vary from pair to pair over an extended region. The more this initial region is confined, the less perfect the momentum entanglement will be, and the less perfect the opposite movements of the pilot-wave particles will be.
When do we have a transition between the small time and the large time regime? We could say that the transition between the small time regime and the large time regime occurs at time $`T`$ when both contributions to $`L(t)`$ are equally large, i.e. $`L(0)=\mathrm{\Delta }\left(\widehat{p}_{1j}(0)+\widehat{p}_{2j}(0)\right)T/m`$ (for fragments with equal masses). If we have $`\mathrm{\Delta }\left(\widehat{p}_{1j}(0)+\widehat{p}_{2j}(0)\right)\mathrm{\Delta }\left(\widehat{x}_{1j}(0)+\widehat{p}_{xj}(0)\right)\mathrm{}`$ then the transition time is given by $`TL(0)^2k_c/c`$, with $`k_c`$ the wavenumber corresponding to the Compton wavelength of the particles.
If we now consider the experiment performed by Pittman et al., then Peres’ analysis can be applied. This is because the experiment can be seen to occur in the large distance regime because of presence of the lens. In some sense the lens can be seen as projecting the angular correlation at infinity, to finite distances (at distances $`2f`$ from the lens). The better the momentum correlation, the better the angular correlation is at infinity or the better the optical imaging is.
It is interesting to use the data of this experiment to give an example where the transition can be situated between the large distance and the small distance regime for a system which displays strong momentum correlation. The wavelength of the pump photon is $`351.1nm`$ and the width of the pump beam is $`L(0)=2mm`$. If we now assume that the equality is reached for the Heisenberg uncertainty then, in the absence of a lens, the transition time is of the order $`TL(0)^2k_c/c`$. The distance traveled by the photons at that time would be $`RL(0)^2k_c70m`$. This means that if we would create momentum entangled photons via spontaneous parametric down conversion, then the angular deviation would be dominated by the error arising from the uncertainty of the source within a distance $`R70m`$.
The deviation from perfect alignment resulting from the uncertainty on the source is also important in Popper’s experimental proposal. It can easily be seen that Popper’s experiment cannot occur in the large time regime. If the particles would travel large distances (and hence obtain good angular correlation), the virtual slit (which is of the order of the transversal deviation) would be too large to have virtual diffraction of the particles on Bob’s side. Hence, Popper’s original experimental proposal should be considered in the small time regime. However, in this regime the uncertainty on the source becomes the most important contribution to the angular deviation and as follows from the analysis of Collett and Loudon, this uncertainty makes a detectable virtual diffraction impossible. This also implies that in discussing Popper’s experiment one should be careful with statements such as “…the allowed deviation from perfect alignment is of the order of $`\mathrm{\Delta }|p_1+p_2|/|p_1p_2|`$, which is much too small to be of any consequence in the present discussion.’ and “…nearly perfect alignment can be taken for granted, …” .
### 6.4 Pilot-wave description of a decaying system
In this section, we give a simplified description of a system consisting of two particles with opposite momenta, resulting from a decaying system at rest, in terms of the pilot-wave interpretation.<sup>4</sup><sup>4</sup>4Our simplified pilot-wave approach to a decaying system is to be distinguished from the one studied by Y. Nogami et al. , where a decaying system is represented by a particle that leaks out from a region surrounded by a repulsive potential barrier. We use the wavefunction in (6.4) with momentum distribution
$$F(𝐩_1,𝐩_2)=N\delta (𝐩_1+𝐩_2)e^{\alpha p_1^2/\mathrm{}}$$
(6.20)
where $`N`$ is a normalization factor. The parameter $`\alpha `$ sets the scale of the initial separation of the two particles, as will be seen soon. A small $`\alpha `$ will correspond to the considered physical situation of a decaying system. The parameter is introduced in order to avoid singularities arising from the $`\delta `$-distribution when calculating the pilot-wave trajectories later on. Although a system that decays from rest has a certain fixed total energy, we omit this energy restriction, just as in Section 6.3, for reasons explained in Appendix D. Note that the exponential factor in (6.20) does not restrict the value of the total energy for a small $`\alpha `$.
The wavefunction corresponding to the distribution $`F`$ is
$$\psi (𝐱_1,𝐱_2,t)=N\left(\frac{\pi \mathrm{}}{\alpha +it/2\mu }\right)^{3/2}e^{(𝐱_1𝐱_2)^2/4\mathrm{}(\alpha +\frac{it}{2\mu })}$$
(6.21)
with $`\mu `$ the reduced mass of the fragments: $`\frac{1}{\mu }=\frac{1}{m_1}+\frac{1}{m_2}`$. At $`t=0`$ the probability distribution is
$$|\psi (𝐱_1,𝐱_2,0)|^2=N^2\left(\frac{\pi \mathrm{}}{\alpha }\right)^3e^{(𝐱_1𝐱_2)^2/2\mathrm{}\alpha }.$$
(6.22)
It follows that a small value for $`\alpha `$ corresponds to the considered physical situation of a decaying system. However, the place of decay is unknown. This is a consequence of the Heisenberg uncertainty, as explained in the Section 6.3.
The trajectories $`𝐱_j(t)`$ of the particle beables are found by solving the guidance equations
$$\frac{d𝐱_j}{dt}=\frac{1}{m_j}\frac{\text{Re}\left(\psi ^{}(𝐱_1,𝐱_2,t)\widehat{𝐩}_j\psi (𝐱_1,𝐱_2,t)\right)}{|\psi (𝐱_1,𝐱_\mathrm{𝟐},t)|^2}.$$
(6.23)
Because $`(\widehat{𝐩}_1+\widehat{𝐩}_2)\psi =0`$ the trajectories of the particles satisfy
$$\frac{d}{dt}(m_1𝐱_1+m_2𝐱_2)=\mathrm{𝟎}.$$
(6.24)
This shows that the particles have opposite speeds and thus move in opposite directions. Integration of the differential equations (6.23) leads to
$`𝐱_1(t)`$ $`=`$ $`𝐜_1+𝐜_2\sqrt{t^2/4\mu ^2+\alpha ^2}`$
$`𝐱_2(t)`$ $`=`$ $`𝐜_1𝐜_2\sqrt{t^2/4\mu ^2+\alpha ^2}`$ (6.25)
where $`𝐜_1`$ and $`𝐜_2`$ are arbitrary constant vectors. It follows that the particles also move along straight lines. At $`t=0`$ the probability distribution $`|\psi |^2`$ is sharply peaked at $`𝐱_1=𝐱_2`$ (for small values of $`\alpha `$) and hence the particle beables will depart near each other. As follows from (6.25), their further propagation proceeds along straight lines, in the direction of their connecting line. Thus opposite momenta lead to opposite directions of movement for pilot-wave particles. But their place of departure is located within an extended area, in order to preserve momentum correlation.
By using pilot-wave theory, we are thus able to retain part of the classical picture of a decaying system at rest. Note the similarity in language with the one used by Popper to describe his experiment. The difference is that Popper assumed the particles to depart from a confined region (which is however incompatible with opposite momenta in quantum mechanics).
### 6.5 Pilot-wave description of the experiment of Pittman et al.
Although the experiment of Pittman et al. can be correctly explained with quantum optics, we will provide a pilot-wave account of the experiment, when it is ‘translated’ into its massive particle equivalent. One of the reasons to use pilot-wave theory is that it justifies the conceptual photon trajectories drawn by Pittman et al. . I.e. the photon trajectories coincide with the trajectories of pilot-wave particles in the massive particle equivalent of the experiment. This pilot-wave approach is to be contrasted with the explanation in terms of ‘usual’ geometrical optics used by Pittman et al.. In quantum optics, these paths are usually regarded as a visualization of the different contributions to the detection probabilities.
Because there is at present no satisfactory pilot-wave interpretation in terms of particle beables for photons, see Chapter 3, we will follow Peres’ point of departure and we will consider the non-relativistic massive particle wavefunction in (6.4) which can then be seen as describing the massive particle equivalent of the experiment by Pittman et al. The spontaneous parametric down conversion source then corresponds to a decaying system at rest, resulting in two energy and momentum correlated fragments. We will assume the total momentum of the fragments to be zero, instead of some fixed value corresponding to the initial momentum of the total system (which would represent the momentum of the pump photon). This assumption corresponds to the ‘unfolded’ schematic introduced by Pittman et al. (which is displayed in Fig. 6.1). In this way we can use the momentum distribution $`F`$ defined in the previous section
$$F(𝐩_1,𝐩_2)=N\delta (𝐩_1+𝐩_2)e^{\alpha p_1^2/\mathrm{}}.$$
(6.26)
In the preceding section, we described the free evolution after the decay of the system. The unknown place of decay in the massive particle case corresponds to the unknown place of creation within the BBO crystal in the photon case, due to the width of the pump beam. To complete the pilot-wave description of the massive particle equivalent of optical imaging, we just have to describe the system’s interaction with the lens. In classical optics we can use ray optics to describe the action of the lens on an impinging light beam . The rays are such that the Gaussian thin lens equations are satisfied. Two generic examples, which we will need later on, are the following. The effect of the lens on a plane wave is to turn it into a converging wave, with focus in the focal plane, such that the corresponding rays obey the lens equations. The characteristics of the converging wave are then determined by the momentum of the incoming plane wave and the focal length. A second example is a spherical wave, representing a point source. If we assume that the light source is located in a plane at a distance $`S`$ from the lens, then the spherical wave will turn into a converging wave with focus in the plane at a distance $`S^{}`$ from the lens so that $`1/S+1/S^{}=1/f`$ and the source, the image and the centre of the lens will be aligned. In massive particle quantum physics the equivalent of optical lenses are electrostatic or magnetic lenses. These electromagnetic lenses are generally used to collimate or focus beams of charged particles. This field of research is usually called optics of charged-particle beams or the theory of charged-particle beams through electromagnetic systems. Most of the literature deals with the classical description of the particles and only recently the quantum mechanical approach has been studied, see for example Hawkes and Kasper , and Khan and Jagannathan and references therein. Here, we will not consider the detailed analysis of particles passing through such electromagnetic lenses, and use directly, in the spirit of de Broglie, the analogy with classical optics. For example, we can describe the action of an electromagnetic lens as turning a quantum mechanical plane wave into a Gaussian wave (we can take this as the analogue of the converging wave in classical optics, because a Gaussian wave is contracting before expanding), determined by the momentum of the incoming wave and the focal length. This analogy is very appealing because the rays in classical optics can be ‘identified’ with the pilot-wave trajectories. This is because in the one-particle case, the curves determined by the normals of the wavefronts of the quantum mechanical wavefunction are just the possible trajectories of the particle beables.<sup>5</sup><sup>5</sup>5For a non-relativistic particle the guidance equation can be written as $`d𝐱/dt=\mathbf{}S/m`$, where $`S`$ is the phase of the wavefunction measured in units $`\mathrm{}`$, i.e. $`\psi =|\psi |\mathrm{exp}(iS/\mathrm{})`$ . If we apply this to our decaying system, then every plane wave of the particle impinging on the lens, say particle two, in the integral in (6.4) is turned into a particular Gaussian wave. The resulting wave is then
$$\psi ^{}(𝐱_1,𝐱_2,t)=d^3p_1d^3p_2F(𝐩_1,𝐩_2)e^{i(𝐩_1𝐱_1p_1^2t/2m_1)/\mathrm{}}G(𝐱_2,𝐩_2,f)$$
(6.27)
where $`G`$ represents the Gaussian wave. This wave is guiding the particle beables after particle beable two passed the lens. To avoid unnecessary mathematical complications when calculating the trajectories implied by the wave (6.27), we assume that the place of decay of the system is somewhere in the middle between the lens and the detector on the right (where the idler photon arrives in the experiment of Pittman et al.). This corresponds to a BBO crystal placed in the middle instead of it placed near the lens, as in the experiment. When particle beable two arrives in the vicinity of the lens, particle beable one will arrive in the vicinity of the transversal plane on the right, where the detectors are placed. Suppose that the particle beable is detected in the transversal plane at the position $`𝐚`$. Due to the correlation of the detector and the particle on the right, particle beable two will be effectively guided by the wavefunction<sup>6</sup><sup>6</sup>6In the standard interpretation of quantum mechanics, this process is the collapse of the wavefunction. The pilot-wave description of this process was given in Section 1.2.
$`\psi _2(𝐱_2,t)`$ $``$ $`{\displaystyle d^3pe^{\frac{i}{\mathrm{}}\left(𝐩(𝐚𝐱_2)p^2t/2m_2\right)\alpha p^2/\mathrm{}}}`$ (6.28)
$``$ $`\left({\displaystyle \frac{\pi \mathrm{}}{\alpha +it/2m_2}}\right)^{3/2}e^{(𝐚𝐱_2)^2/4\mathrm{}(\alpha +\frac{it}{2m_2})}.`$
The phase of this wave is
$$S(𝐱_2,t)=\frac{t(𝐚𝐱_2)^2}{8m_2\alpha ^2+t^2/m_2}\frac{3\mathrm{}}{2}\mathrm{tan}^1(t/2m_2\alpha ).$$
(6.29)
So, the wavefronts of the guiding wave of particle two are spheres with centre in $`𝐚`$. Because the detectors are placed in transversal planes at distances $`S`$ and $`S^{}`$ from the lens, with $`S`$ and $`S^{}`$ obeying the Gaussian lens equation (6.3), this wave will result, after propagation through the lens, in a converging wave with focus in the plane at a distance $`S`$ from the lens and where the focus is determined by the Gaussian thin lens equations. Hereby we used again the analogy with classical Gaussian optics. If for example $`S=S^{}`$ and if the centre of the lens is taken as the origin of our coordinate system, then the focus will be at $`𝐚`$ (see Fig. 1). As a result, particle beable two will be detected in the focus of the wave. Because we used a Gaussian to describe the converging wave, the trajectories will not be straight lines, but will be curved (for images see Holland p162). The curvature will depend on the width in the focus of the Gaussian. In the limit of a zero width, however, the trajectories will approach straight lines, directed from the lens towards the focus of the wave. When the coincidence detections are considered, it will appear as if particle beable two departs from the place of detection of particle beable one.
This completes the analysis in terms of the pilot-wave interpretation of the phenomenon of optical imaging. Before the fragments reach the lens, they move along straight lines from the place of decay. Note that this place of decay is not fixed, in order to guarantee the momentum correlation $`𝐩_1+𝐩_2=\mathrm{𝟎}`$. When one of the particles reaches the lens, its direction of movement will change in accordance with the classical thin lens approximations. We assumed hereby that the place of decay is centred between the right detector and the lens. It can be expected, although it is not proven, that a random place of decay (for example near the lens) will lead to the same results in the pilot-wave description of the experiment.
### 6.6 Conclusion
In conclusion, we showed that Peres’ analysis concerning the question to what extent opposite momenta lead to opposite directions, is only valid in the large distance regime. In the small time regime there is an additional source of angular deviation. On the other hand the statement ‘opposite momenta lead to opposite directions’ is true in the language of pilot-wave theory. I.e. the particle beables travel in opposite directions when the wavefunction has eigenvalue zero for the total momentum operator (however from an unknown place of departure). We also showed that pilot-wave trajectories coincide with the pictures of conceptual trajectories present in the paper Pittman et al. Hence, pilot-wave theory could be used to mathematically underpin these conceptual trajectories.
Note that the experiment of Kim and Shih is very illustrative for the need for perfect momentum correlation of the photons, or equivalently that there must be very little restriction on the place of creation of the photons to create a perfect image. This is because Kim and Shih failed in their original intention to perform Popper’s gedankenexperiment, due to the restricted diameter of the pump beam used in the experiment. The imperfect momentum correlation then led to an imperfect optical image . This is immediately obvious when we consider our pilot-wave description of optical imaging, because if the momentum correlation is imperfect, the pilot-wave particles will not move in opposite directions before the system reaches the lens.
## Chapter 7 Conclusion and outlook
In this thesis we studied the de Broglie-Bohm pilot-wave interpretation of quantum theory. We mainly focussed on the question to which extent it is possible to provide a pilot-wave interpretation for quantum theory.
In the context of non-relativistic quantum theory it is no problem to construct a pilot-wave interpretation. It is straightforward to extend the pilot-wave interpretation of de Broglie and Bohm for a spinless particle to arbitrary spin. In fact there even exists more than one approach. In the model we presented here, the beables are point-particles without any spin properties.
Problems arise when we try to transcript this pilot-wave model to relativistic wave equations. In fact these problems do not only arise when we try to develop a pilot-wave interpretation. Problems already arise when we try to transcript the quantum mechanical interpretation of non-relativistic quantum theory to relativistic wave equations. In the first place, there is the problem of identifying a future-causal current which can be interpreted as a particle probability current. Closely related to this problem is the problem of defining a positive definite inner product. Then there is also the problem of interpreting the negative energy states. Only in restricted cases these problems can be solved (e.g. as in the spin-1/2 Dirac theory coupled to the electromagnetic field). As is well known these problems led to the conception of quantum field theory. The problems which blocked a quantum mechanical interpretation of relativistic wave equations are not present in quantum field theory. One has a positive definite inner product and there are no problems with the interpretation of negative energy states. Particles and anti-particles can be freely created and annihilated without the need of considering negative energy states. Hence in the pilot-wave approach we should not stick to relativistic wave equations either, but focus on quantum field theory instead.
In quantum field theory, field operators take over the role of the particle operators in non-relativistic quantum theory. This suggests that the notion of fields is more fundamental than the notion of particles for high energy physics. Therefore, fields seem to be the most natural beables in the pilot-wave approach. This seems to be confirmed at least in the case of bosonic quantum field theory. For bosonic quantum field theory it was straightforward to construct a pilot-wave interpretation in terms of field beables. We could construct a pilot-wave theory for massive spin-0 fields and spin-1 fields. Even for gauge theories we could find a pilot-wave interpretation. We took an approach in which only beables are introduced for gauge independent variables, as in the model by Bohm and Kaloyerou. Valentini presented an other approach in which beables are also introduced for gauge variables. But we argued that this last approach leads to non-normalizable densities of field beables. In addition the guidance equations for the gauge variables are rather meaningless because they just express that these variables are stationary. We also indicated that for non-Abelian gauge theories a pilot-wave interpretation is in principle possible along the same lines as in the Abelian case (i.e. the electromagnetic field case), the only problem is the identification of the gauge invariant degrees of freedom.
In the case of fermionic quantum field theory the construction of a pilot-wave interpretation is not so straightforward. We elaborated on an idea of Valentini to use a representation for the field operators in terms of Grassmann fields. In this representation a functional Schrödinger picture can be obtained. However, we found that it was not possible to devise a pilot-model where the beables are Grassmann fields, at least not in the way Valentini originally suggested.
Is a pilot-wave interpretation in terms of field beables therefore impossible? It seems not, because there is still a model by Holland, where the beables correspond to rotators in each point of the configuration space (which can be physical 3-space or momentum space). However, this model brings with it new questions which have yet to be answered. Is the ontology for fermionic fields in terms of rotator beables compatible with the ontology for bosonic fields in terms of ‘ordinary’ field beables? Or can one construct a rotator ontology for bosonic fields as well? Another question is whether wavefunctionals which correspond to distinct macroscopic states non-overlapping in the configuration space of rotator fields? Admittedly, this last question is in fact not yet sufficiently addressed either in the context of ‘ordinary’ field beables.
Another possibility is that we should go along with a particle ontology, instead of with a field ontology, for fermionic field theory. This is an approach taken by Dürr, Goldstein, Tumulka and Zanghì. The Bell-type model they introduced for quantum electrodynamics introduces point particles as beables, but only for the fermions and not for the electromagnetic field. It is the question whether this type of model can be extended so that particle beables are introduced for the bosonic fields as well. It could also be that beables need not be introduced for bosons at all. Perhaps introducing beables corresponding to fermions is sufficient to solve the measurement problem. It would then be interesting to see the extension of this model to account for weak and strong interactions.
We think that the construction of a coherent pilot-wave interpretation for fermionic field theory is the main issue to be addressed in pilot-wave theory. This could be done by further elaborating on the model of Holland. Or if one allows stochasticity in the model, one could start from the model by Dürr et al. In any case, we think that a coherent pilot-wave interpretation for fermionic fields must be possible. The reason is that the standard interpretation for fermionic fields does not encounters problems, i.e. a Hilbert space can be set up together with an associated operator formalism. This is contrary to the case of relativistic wave equations where the problems in providing a pilot-wave interpretation were directly related to problems already rooted in the standard interpretation. With a coherent pilot-wave interpretation for fermionic fields and together with the pilot-wave interpretation for massive spin-0 and massive spin-1 presented in this thesis and with a pilot-wave interpretation for non-Abelian gauge theories (which presents no problem in principle), we could in principle be given a pilot-wave interpretaton for the standard model, which represents to this day the credo of high energy quantum physics.
## Appendix A Representation of the Kemmer-Duffin-Petiau matrices
We adopt the matrix representation used in , where the $`\beta ^i`$ correspond to $`i\beta ^i`$ in . For spin-0 the $`\beta ^\mu `$ matrices read
$`\beta ^0=\left(\begin{array}{c}0\\ 0\\ 0\\ 0\\ i\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}i\\ 0\\ 0\\ 0\\ 0\end{array}\right),\beta ^1=\left(\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ i\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ i\\ 0\\ 0\\ 0\end{array}\right),`$ (A.51)
$`\beta ^2=\left(\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ i\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ i\\ 0\\ 0\end{array}\right),\beta ^3=\left(\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ i\end{array}\begin{array}{c}0\\ 0\\ 0\\ i\\ 0\end{array}\right).`$ (A.102)
The $`\gamma `$ matrix used in the massless spin-0 theory reads
$`\gamma =\left(\begin{array}{c}1\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 1\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 1\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 1\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\end{array}\right).`$ (A.128)
For spin-1 the $`\beta ^\mu `$ matrices read
$`\beta ^0=\left(\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\\ 0\\ i\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ i\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ i\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}i\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ i\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ i\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\right),\beta ^1=\left(\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ i\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ i\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ i\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\\ i\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ i\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}i\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\right),`$ (A.329)
$`\beta ^2=\left(\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ i\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ i\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\\ 0\\ i\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\\ i\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ i\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ i\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\right),\beta ^3=\left(\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ i\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ i\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\\ 0\\ i\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ i\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ i\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ i\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\right).`$ (A.530)
The $`\gamma `$ matrix used in the massless spin-0 theory reads
$`\gamma =\left(\begin{array}{c}1\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 1\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 1\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 1\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 1\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\\ 1\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\right).`$ (A.631)
## Appendix B Class of representations which leads to equivalent pilot-wave interpretations for the electromagnetic field
Suppose we have two different canonical transformations represented by $`K`$ and $`K^{}`$ (cf. Section 4.5.2) which allow for a separation of the true degrees of freedom from the constraints in the Coulomb gauge:
$`V_i(𝐱)`$ $`=`$ $`{\displaystyle d^3yK_{ij}(𝐱,𝐲)\stackrel{~}{V}_j(𝐲)}={\displaystyle d^3yK_{ij}^{}(𝐱,𝐲)\stackrel{~}{V}_j^{}(𝐲)},`$
$`\mathrm{\Pi }_{V_i}(𝐱)`$ $`=`$ $`{\displaystyle d^3y\mathrm{\Pi }_{\stackrel{~}{V}_j}(𝐲)K_{ji}^1(𝐲,𝐱)}={\displaystyle d^3y\mathrm{\Pi }_{\stackrel{~}{V}_j^{}}(𝐲)K_{ji}^1(𝐲,𝐱)}.`$ (B.1)
We assume that the constraints are $`\stackrel{~}{V}_3=\mathrm{\Pi }_{\stackrel{~}{V}_3}=0`$ for the set of unprimed variables and $`\stackrel{~}{V}_3^{}=\mathrm{\Pi }_{\stackrel{~}{V}_3^{}}=0`$ for the set of primed variables. Quantization then proceeds by imposing the standard commutation relations for the true degrees of freedom $`\stackrel{~}{V}_k,\mathrm{\Pi }_{\stackrel{~}{V}_k}`$, $`k=1,2`$ for the set of unprimed variables or $`\stackrel{~}{V}_k^{},\mathrm{\Pi }_{\stackrel{~}{V}_k^{}}`$, $`k=1,2`$ for the set of primed variables. These commutation relations can then be realized by the standard representation
$$\widehat{\stackrel{~}{V}}_k(𝐱)=\stackrel{~}{V}_k(𝐱),\widehat{\mathrm{\Pi }}_{\stackrel{~}{V}_k}(𝐱)=i\frac{\delta }{\delta \stackrel{~}{V}_k(𝐱)},\text{for }k=1,2,$$
(B.2)
and similarly for the primed variables.
In terms of respectively the unprimed and the primed variables, the matrices $`h_{kl}`$ and $`\overline{h}_{kl}`$ in the Hamiltonian (4.141) will respectively depend on the unprimed and primed transformations $`K`$ and $`K^1`$. The different Hamiltonians can then be transformed into each other. This can be done by using the transformation (B.1) but then applied to the fields which are used in the representation (B.5).<sup>1</sup><sup>1</sup>1Note that we have used the same notation for the classical fields (i.e. the unquantized fields) and the fields in the representation (B.5). But despite the same notation, these fields have a different meaning.
The transformation between the primed and unprimed fields which are used in the above representation read
$`\stackrel{~}{V}_k^{}(𝐱)`$ $`=`$ $`{\displaystyle \underset{l=1}{\overset{2}{}}}{\displaystyle d^3yT_{kl}(𝐱,𝐲)\stackrel{~}{V}_l(𝐲)},`$
$`\stackrel{~}{V}_k(𝐱)`$ $`=`$ $`{\displaystyle \underset{l=1}{\overset{2}{}}}{\displaystyle d^3yT_{kl}^1(𝐱,𝐲)\stackrel{~}{V}_l^{}(𝐲)},`$
$`{\displaystyle \frac{\delta }{\delta \stackrel{~}{V}_k(𝐱)}}`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{2}{}}}{\displaystyle d^3yT_{kl}(𝐲,𝐱)\frac{\delta }{\delta \stackrel{~}{V}_l^{}(𝐲)}},`$
$`{\displaystyle \frac{\delta }{\delta \stackrel{~}{V}_{}^{}{}_{k}{}^{}(𝐱)}}`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{2}{}}}{\displaystyle d^3yT_{kl}^1(𝐲,𝐱)\frac{\delta }{\delta \stackrel{~}{V}_l(𝐲)}}`$ (B.3)
with $`k=1,2`$ and
$`T_{kl}(𝐱,𝐲)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{3}{}}}{\displaystyle d^3zK_{ki}^1(𝐱,𝐳)K_{il}(𝐳,𝐲)},`$
$`T_{kl}^1(𝐱,𝐲)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{3}{}}}{\displaystyle d^3zK_{ki}^1(𝐱,𝐳)K_{il}^{}(𝐳,𝐲)}.`$ (B.4)
We have hereby used the properties (4.127)-(4.130) of $`K`$ and $`K^{}`$ to show identities like
$`{\displaystyle \underset{i=1}{\overset{3}{}}}{\displaystyle d^3zK_{ki}^1(𝐱,𝐳)K_{i3}(𝐳,𝐲)}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{3}{}}}{\displaystyle d^3zK_{ki}^1(𝐱,𝐳)_{z_i}U(𝐳,𝐲)}`$ (B.5)
$`=`$ $`{\displaystyle \underset{i=1}{\overset{3}{}}}{\displaystyle d^3z\left(_{z_i}K_{ki}^1(𝐱,𝐳)\right)U(𝐳,𝐲)}`$
$`=`$ $`0,\text{for }k=1,2.`$
We have hereby assumed that the boundary terms vanishes after the partial integration. Identities like these are the reason why the transformations are independent of the constraint variables $`\stackrel{~}{V}_3,\mathrm{\Pi }_{\stackrel{~}{V}_3}`$ and $`\stackrel{~}{V}_3^{},\mathrm{\Pi }_{\stackrel{~}{V}_3^{}}`$, as should be.
If the normalized wavefunction $`\mathrm{\Psi }(\stackrel{~}{V}_i)`$ is a solution to the functional Schrödinger equation (4.143) where the functions $`h_{kl}`$ and $`\overline{h}_{kl}`$ depend on the unprimed transformations $`K`$ and $`K^1`$, then the wavefunctional
$$\mathrm{\Psi }^{}(\stackrel{~}{V}_j^{})=N\mathrm{\Psi }\left(\underset{j=1}{\overset{2}{}}d^3yT_{ij}^1(𝐱,𝐲)\stackrel{~}{V}_j^{}(𝐲)\right)=N\mathrm{\Psi }(\stackrel{~}{V}_i)$$
(B.6)
satisfies the functional Schrödinger equation (4.143) where the Hamiltonian depends on the primed transformations. The normalization constant is determined by the Jacobian of the transformation, i.e. $`|N|^2=detT^2`$.
Hence the quantum theories in terms of primed and unprimed fields are equivalent. The pilot-wave interpretations that may be constructed in terms of the two sets of unconstrained variables are also equivalent; the field beables $`\stackrel{~}{V}_l(𝐱,t)`$ and $`\stackrel{~}{V}_l^{}(𝐱,t)`$ can be transformed to each other by application of the same transformation (B.3) as the classical fields.
## Appendix C The Grassmann algebra: definitions and properties
### C.1 The Grassmann algebra
An algebra over the complex numbers whose generators $`\eta (k)`$, labeled by $`k`$, satisfy
$$[\eta (k),\eta (l)]_+=0,$$
(C.1)
is called a Grassmann algebra . The generators are called Grassmann numbers. An element of the Grassmann algebra which respectively commutes or anti-commutes with all the other elements of the Grassmann algebra is respectively called an even or odd element of the Grassmann algebra.
The left derivatives $`\frac{\stackrel{}{\delta }}{\delta \eta (k)}`$ and right derivatives $`\frac{\stackrel{}{\delta }}{\delta \eta (k)}`$ are defined by
$`[\eta (k),{\displaystyle \frac{\stackrel{}{\delta }}{\delta \eta (l)}}]_+=[\eta (k),{\displaystyle \frac{\stackrel{}{\delta }}{\delta \eta (l)}}]_+=\delta (kl),`$
$`[{\displaystyle \frac{\stackrel{}{\delta }}{\delta \eta (k)}},{\displaystyle \frac{\stackrel{}{\delta }}{\delta \eta (l)}}]_+=[{\displaystyle \frac{\stackrel{}{\delta }}{\delta \eta (k)}},{\displaystyle \frac{\stackrel{}{\delta }}{\delta \eta (l)}}]_+=[{\displaystyle \frac{\stackrel{}{\delta }}{\delta \eta (k)}},{\displaystyle \frac{\stackrel{}{\delta }}{\delta \eta (l)}}]_+=0.`$ (C.2)
With this definition we have
$`{\displaystyle \frac{\stackrel{}{\delta }}{\delta \eta (k)}}\left(AB\right)`$ $`=`$ $`{\displaystyle \frac{\stackrel{}{\delta }A}{\delta \eta (k)}}B+(1)^{P_A}A{\displaystyle \frac{\stackrel{}{\delta }B}{\delta \eta (k)}},`$
$`\left(AB\right){\displaystyle \frac{\stackrel{}{\delta }}{\delta \eta (k)}}`$ $`=`$ $`(1)^{P_B}{\displaystyle \frac{A\stackrel{}{\delta }}{\delta \eta (k)}}B+A{\displaystyle \frac{B\stackrel{}{\delta }}{\delta \eta (k)}},`$ (C.3)
where $`P_A`$ equals $`0`$ or $`1`$, depending on whether the element $`A`$ is respectively even or odd.
With the symbols $`d\eta (k)`$ which satisfy
$$[\eta (k),d\eta (l)]_+=[d\eta (k),d\eta (l)]_+=0$$
(C.4)
we may define the integrals
$$𝑑\eta (k)=0,\eta (k)𝑑\eta (k)=1.$$
(C.5)
The integral over an arbitrary element of the Grassmann algebra is defined by linearly extending the definitions (C.5).
For each element of the Grassmann algebra $`f`$ we may also introduce a conjugate $`f^{}`$ which has the properties
$$(f_1f_2)^{}=f_2^{}f_1^{},\left(f^{}\right)^{}=f,[\eta (k),\eta ^{}(l)]_+=[\eta ^{}(k),\eta ^{}(l)]_+=0.$$
(C.6)
and which is such that when the conjugate is restricted to the subspace of complex numbers, it is simply the complex conjugate. As a corollary we have
$$\left(\frac{\stackrel{}{\delta }}{\delta \eta (k)}f\right)^{}=f^{}\frac{\stackrel{}{\delta }}{\delta \eta ^{}(k)},\left(f\frac{\stackrel{}{\delta }}{\delta \eta (k)}\right)^{}=\frac{\stackrel{}{\delta }}{\delta \eta ^{}(k)}f^{}$$
(C.7)
and similarly for derivatives with respect to the fields $`\eta ^{}(k)`$.
### C.2 The inverse of an element of the Grassmann algebra
Suppose an infinite dimensional complex Grassmann algebra generated by the Grassmann numbers $`\{\eta (k),\eta ^{}(k)\}`$ with $`\eta ^{}(k)`$ the conjugate of $`\eta (k)`$, which satisfy the anti-commutation relations
$$[\eta (k),\eta (l)]_+=[\eta (k),\eta ^{}(l)]_+=[\eta ^{}(k),\eta ^{}(l)]_+=0.$$
(C.8)
The inverse element for the multiplication can not be defined for any element of the Grassmann algebra, e.g. if the Grassmann number $`\eta (k)`$ would have an inverse then this would be in contradiction with the property $`\eta (k)^2=0`$. However, there is a class of elements for which we can define an inverse. Consider therefore the following decomposition of an arbitrary element $`f`$ of the Grassmann algebra as $`f=c+g`$, where $`c`$ is a complex number with
$$𝒟\eta ^{}𝒟\eta \underset{k}{}\eta (k)\eta ^{}(k)f=c.$$
(C.9)
This last equation implies that $`g`$ can be written as $`g=_ig_i`$ with $`g_i^2=0`$. For every element $`f`$ of the Grassmann algebra for which $`c0`$, one can define the unique inverse element $`f^1`$ as
$$f^1=\frac{1}{c}\underset{l=0}{\overset{+\mathrm{}}{}}\left(\frac{g}{c}\right)^l,$$
(C.10)
for which $`f^1f=ff^1=1`$. The inverse of a product satisfies
$$(fg)^1=g^1f^1.$$
(C.11)
For an element $`f`$ for which $`c=0`$, no such elements $`f^1`$ exist.
## Appendix D Note on the energy restriction for a decaying system
Conservation of energy requires that the energy of the total system equals the energy of the decaying system $`E_0`$. If this decaying system is initially at rest, $`E_0`$ will be the rest energy of the system. Suppose now that we take a $`\delta `$-distribution for this energy restriction i.e.
$$F(𝐩_1,𝐩_2)=f(𝐩_1,𝐩_2)\delta (EE_0)$$
(D.1)
where $`f`$ determines the momentum correlation (this is for example the distribution in (6.18) or (6.20)). If we take a distribution $`f`$ which satisfies $`f(𝐩_1,𝐩_2)=f^{}(𝐩_1,𝐩_2)`$, then the probability currents of the two particles are both zero for all times, i.e.
$$𝐣_i=\frac{\mathrm{}}{m_i}\text{Im}\left(\psi ^{}_{𝐱_i}\psi \right)=\mathrm{𝟎},i=1,2.$$
(D.2)
This implies that the particles show no evolution. In pilot-wave theory this corresponds with particle beables that stand still, because the speeds are defined as $`d𝐱_i/dt=𝐣_i/|\psi |^2`$. As a result, the considered distribution in (D.1) does not actually represent a decaying system. We can resolve this problem by allowing a finite energy width centred around $`E_0`$. Nevertheless, it will follow that, if the wavefunction displays strong momentum correlation in the sense that $`𝐩_1+𝐩_2=\mathrm{𝟎}`$, then the restriction to a small energy width only involves a minor broadening of the wavefunction, which implies that we can leave the restriction on the total energy aside for our qualitative analysis.
For the momentum distribution $`f`$ we will take the distribution in (6.20), i.e. $`f(𝐩_1,𝐩_2)\delta (𝐩_1+𝐩_2)e^{\alpha p_1^2/\mathrm{}}`$. The restriction on the total energy is accomplished by integrating over values for $`(𝐩_1,𝐩_2)`$ for which $`E_{}EE_+`$, for a certain minimum energy value $`E_{}`$ and a certain maximum energy value $`E_+`$. We therefore define the following function
$$\text{disc}(E_\pm )(𝐩_1,𝐩_2)=\{\begin{array}{cc}1\hfill & \text{if}\frac{p_1^2}{2m_1}+\frac{p_2^2}{2m_2}E_\pm \hfill \\ 0\hfill & \text{otherwise}\hfill \end{array}.$$
The momentum distribution then becomes
$$F(𝐩_1,𝐩_2)=f(𝐩_1,𝐩_2)\left[\text{disc}(E_+)\text{disc}(E_{})\right].$$
(D.3)
The resulting wavefunction of the system is then
$`\psi (𝐱_1,𝐱_2,t)`$ $`=`$ $`{\displaystyle f(𝐩_1,𝐩_2)\left[\text{disc}(E_+)\text{disc}(E_{})\right]e^{i(𝐩_1𝐱_1+𝐩_2𝐱_2Et)/\mathrm{}}𝑑𝐩_1𝑑𝐩_2}`$ (D.4)
$``$ $`{\displaystyle e^{i𝐩(𝐱_1𝐱_2)/\mathrm{}(it/2\mu +\alpha )p^2/\mathrm{}}\left[\text{disc}^{}(E_+)\text{disc}^{}(E_{})\right]𝑑𝐩}`$
where
$$\text{disc}^{}(E_\pm )(𝐩)=\{\begin{array}{cc}1\hfill & \text{if}p^2/2\mu E_\pm \hfill \\ 0\hfill & \text{otherwise}\hfill \end{array}.$$
If we write (D.4) as a Fourier transform then we can apply the convolution theorem
$`\psi (𝐱_1,𝐱_2,t)`$
$`_{\{(𝐱_1𝐱_2)/\mathrm{}\}}^+\left(e^{(it/2\mu +\alpha )p^2/\mathrm{}}\right)_{\{(𝐱_1𝐱_2)/\mathrm{}\}}^+\left(\text{disc}^{}(E_+)\text{disc}^{}(E_{})\right)`$
$`h(x,t)\overline{g}(x),`$ (D.5)
where
$`h(x,t)`$ $`=`$ $`_{\{(𝐱_1𝐱_2)/\mathrm{}\}}^+\left(e^{(it/2\mu +\alpha )p^2/\mathrm{}}\right),`$
$`\overline{g}(x)`$ $`=`$ $`_{\{(𝐱_1𝐱_2)/\mathrm{}\}}^+\left(\text{disc}^{}(E_+)\text{disc}^{}(E_{})\right),`$
$`x`$ $`=`$ $`|𝐱_1𝐱_2|.`$ (D.6)
The first function in the convolution in (D.5) is just the wavefunction in (6.21),
$$h(x,t)\left(\frac{\pi \mathrm{}}{\alpha +it/2\mu }\right)^{3/2}e^{(𝐱_1𝐱_2)^2/4\mathrm{}(\alpha +\frac{it}{2\mu })}.$$
(D.7)
If we define $`a_\pm =2\pi \sqrt{E_\pm 2\mu }/\mathrm{}`$, then the second function in the convolution in (D.5) becomes
$$\overline{g}(x)g(x)=\left[a_+J_1(a_+x)a_{}J_1(a_{}x)\right]/x$$
(D.8)
where $`J_1`$ is the first order spherical Bessel function.
We present now two ways to show that the restriction on the energy can be relinquished, without changing the qualitative analysis. The first way proceeds as follows. In order to evaluate $`g(x)`$, we will substitute some reasonable values for $`E_+`$ and $`E_{}`$ in (D.8). In addition we will assume that the fragments have equal masses so that we can put $`2\mu =m`$, with $`m`$ the mass of one fragment. For $`E_+`$ we will take one percent of the rest mass of the total system in order to avoid the relativistic regime, $`E_+=0.02mc^2`$. We will take an energy gap of $`0.001E_+`$, so that $`E_{}=0.999E_+`$. In this way $`a_+5.58309/\lambda _c`$ and $`a_{}5.58023/\lambda _c`$, where $`\lambda _c`$ is the Compton wavelength of the fragments. In Fig. D.1 the function $`g(x)=\left[a_+J_1(a_+x)a_{}J_1(a_{}x)\right]/x`$ is plotted for $`x`$ in units of the Compton wavelength $`\lambda _c`$.
The figure shows that $`g(x)`$ is a rapidly oscillating function with a peak at $`x=0`$. We now give the distribution of $`h(x,t)`$ at $`t=0`$ a width of the order of twenty times the Compton wavelength, which can be done by adjusting $`\alpha `$. Recall that the function $`h(x,t)`$ was in fact the wavefunction of the system if we did not restrict the energy. So, the width of $`h(x,t)`$ is in fact the measure of the initial nearness of the fragments, which is then of the order of twenty times the Compton wavelength. Then due to the rapid oscillation, the main contribution in the convolution will arise only from the peak in $`g(x)`$ at $`x=0`$. This peak will result in only a small broadening of $`h(x,0)`$ so that $`h(x,0)g(x)h(x,0)`$.
Because the width of $`h(x,t)`$ only increases with time, this approximation will become more precise with time. In conclusion, we can put the unnormalized wavefunction equal to
$$\psi (𝐱_1,𝐱_2,t)\left(\frac{\pi \mathrm{}}{\alpha +it/2\mu }\right)^{3/2}e^{(𝐱_1𝐱_2)^2/4\mathrm{}(\alpha +\frac{it}{2\mu })}.$$
(D.9)
The larger the energy gap or the larger the rest energy, the more rapid the oscillation of $`g(x)`$ will be and the more peaked $`g(x)`$ will be at $`x=0`$, and as a result the more the approximation is valid.
A second way to achieve this result is to assume that $`h(x,t)`$ is very narrowly peaked at $`x=0`$ for $`t=0`$, so that $`h(x,0)\delta (x)`$. This is the case if $`\alpha `$ approaches zero. As a result $`h(x,0)g(x)g(x)`$. Thus in this case the non-normalized initial probability distribution is $`g(x)^2`$. This distribution is plotted in Fig. D.2 with again $`x`$ in units of the Compton wavelength. This figure shows that most of the probability is concentrated within a few times the Compton wavelength.
This implies that the particle beables depart from a very narrow region, only a few times the Compton wavelength in diameter, from each other. Because the wavefunction $`\psi (𝐱_1,𝐱_2,t)`$ in (D.5) only depends on the difference $`|𝐱_1𝐱_2|`$, the velocities (as defined in (6.23)) will be opposite and the particle beables will travel along straight lines in opposite directions.
In conclusion, we have shown that the energy restriction does not put a restriction on the pilot-wave description of the system. The only effect of the energy restriction is a minor broadening of the probability density. |
warning/0506/hep-ph0506188.html | ar5iv | text | # 1 Introduction
## 1 Introduction
One of the key ingredients in the Standard Model (SM) is CP violation, which can be described by the Cabibbo-Kobayashi-Maskawa (CKM) matrix. However, even with this description we still have an incomplete picture concerning the origin of CP violation in the SM. The exploitation of CP violation from the theoretical and experimental sides of physics is very exciting, as it may open a window to the existence of new physics beyond the SM. Note that the existence of CP violation is a well established fact in $`K`$ and $`B`$ meson systems.
In order to study the sources of CP violation it is promising to consider those observables which are sensitive to the possible CP phases. For example, CP asymmetries in decay widths and lepton polarization asymmetries, such as explored in references .
One of the promising directions for measuring CP violation is the analysis of rare semi-leptonic decays. From the experimental perspective the exclusive decay modes, such as $`BK\mathrm{}^+\mathrm{}^{}`$ and $`BK^{}\mathrm{}^+\mathrm{}^{}`$, are easy to measure. Two years ago the Belle and BaBar collaborations announced the following results for the branching ratios for the $`BK\mathrm{}^+\mathrm{}^{}`$ and $`BK^{}\mathrm{}^+\mathrm{}^{}`$ decays;
$`Br(BK\mathrm{}^+\mathrm{}^{})=\{\begin{array}{cc}\left(4.8_{0.9}^{+1.0}\pm 0.3\pm 0.1\right)\times 10^7\hfill & \text{[8]},\\ & \\ \left(0.65_{0.13}^{+0.14}\pm 0.04\right)\times 10^6\hfill & \text{[9]},\end{array}`$ (4)
$`Br(BK^{}\mathrm{}^+\mathrm{}^{})=\{\begin{array}{cc}\left(11.5_{3.4}^{+2.6}\pm 0.8\pm 0.2\right)\times 10^7\hfill & \text{[8]},\\ & \\ \left(0.88_{0.29}^{+0.23}\right)\times 10^6\hfill & \text{[9]}.\end{array}`$ (8)
The analysis for study of possible CP violation in $`BK^{}\mathrm{}^+\mathrm{}^{}`$ was done in earlier works . The goal of our present work is to similarly study the possible CP violation asymmetry in the exclusive $`BK\mathrm{}^+\mathrm{}^{}`$ decay using the most general form of the effective Hamiltonian, including all possible forms of interactions. Such an analysis will be useful for comparisons with experimental results, as the inclusive modes are generally hard to measure. Note that the CP violation in the decay $`BK\mathrm{}^+\mathrm{}^{}`$ is induced by the $`bs\mathrm{}^+\mathrm{}^{}`$ transition, which in the SM is practically equal to zero. This is due to the CKM factors $`V_{ub}V_{us}^{}`$ being negligible, with the result that the unitarity condition produces only an overall phase factor in the matrix element. Therefore the CP asymmetry is strongly suppressed. As such, any deviation from zero for the CP asymmetry would be an indication of new physics.
This paper is organized as follows. In section 2, using the most general form of the effective Hamiltonian, we derive the matrix element of the $`BK\mathrm{}^+\mathrm{}^{}`$ decay in terms of the $`BK`$ transition form-factors. We also derive in this section the general analytic expression for the CP violating asymmetry. Section 3 contains our numerical analysis of the CP violating asymmetries together with our conclusions.
## 2 The matrix element for the $`BK\mathrm{}^+\mathrm{}^{}`$ decay
In this section we calculate the matrix element for the $`BK\mathrm{}^+\mathrm{}^{}`$ decay, which is governed by the $`bs\mathrm{}^+\mathrm{}^{}`$ transition at the quark level. The most general form of the effective Hamiltonian, for the $`bs\mathrm{}^+\mathrm{}^{}`$ transition (in terms of the twelve model independent four-Fermi interactions) can be written in the following form;
$`_{eff}`$ $`=`$ $`{\displaystyle \frac{\alpha G_F}{\sqrt{2}\pi }}V_{tb}V_{ts}^{}[C_{SL}\left(\overline{s}i\sigma _{\mu \nu }{\displaystyle \frac{q^\nu }{q^2}}Lb\right)\overline{\mathrm{}}\gamma ^\mu \mathrm{}+C_{BR}\left(\overline{s}i\sigma _{\mu \nu }{\displaystyle \frac{q^\nu }{q^2}}Rb\right)\overline{\mathrm{}}\gamma ^\mu \mathrm{}`$ (9)
$`+C_{LL}^{tot}\left(\overline{s}_L\gamma _\mu b_L\right)\overline{\mathrm{}}_L\gamma ^\mu \mathrm{}_L+C_{LR}^{tot}\left(\overline{s}_L\gamma _\mu b_L\right)\overline{\mathrm{}}_R\gamma ^\mu \mathrm{}_R+C_{RL}\left(\overline{s}_R\gamma _\mu b_R\right)\overline{\mathrm{}}_L\gamma ^\mu \mathrm{}_L`$
$`+C_{RR}\left(\overline{s}_R\gamma _\mu b_R\right)\overline{\mathrm{}}_R\gamma ^\mu \mathrm{}_R+C_{LRLR}\left(\overline{s}_Lb_R\right)\overline{\mathrm{}}_L\mathrm{}_R+C_{RLLR}\left(\overline{s}_Rb_L\right)\overline{\mathrm{}}_L\mathrm{}_R`$
$`+C_{LRRL}\left(\overline{s}_Lb_R\right)\overline{\mathrm{}}_R\mathrm{}_L+C_{RLRL}\left(\overline{s}_Rb_L\right)\overline{\mathrm{}}_R\mathrm{}_L+C_T\overline{s}\sigma _{\mu \nu }b\overline{\mathrm{}}\sigma ^{\mu \nu }\mathrm{}`$
$`+iC_{TE}ϵ^{\mu \nu \alpha \beta }\overline{s}\sigma _{\mu \nu }b\overline{\mathrm{}}\sigma _{\alpha \beta }\mathrm{}],`$
where $`L/R=\frac{1}{2}(1\gamma _5)`$, the $`C_X`$’s are the Wilson coefficients of the four-Fermi interactions and $`q_\mu =(p_Bp_K)_\mu =(p_++p_{})_\mu `$ is the momentum transfer. Among the twelve Wilson coefficients several already exist in the SM. For example, the first two terms with coefficients $`C_{SL}`$ and $`C_{BR}`$ describe the penguin operators, where in the SM these coefficients are equal to $`2m_sC_7^{eff}`$ and $`2m_bC_7^{eff}`$. The next four terms in Eq.(9) are the vector type interactions with coefficients $`C_{LL}^{tot}`$, $`C_{LR}^{tot}`$, $`C_{RL}`$ and $`C_{RR}`$. Two of these vector interactions, $`C_{LL}^{tot}`$ and $`C_{LR}^{tot}`$, also exist in the SM with the form $`(C_9^{eff}C_{10})`$ and $`(C_9^{eff}+C_{10})`$. Therefore we can say that the coefficients $`C_{LL}^{tot}`$ and $`C_{LR}^{tot}`$ describe the sum of the contributions from the SM and the new physics, where they can be written as;
$`C_{LL}^{tot}`$ $`=`$ $`C_9^{eff}C_{10}+C_{LL},`$
$`C_{LR}^{tot}`$ $`=`$ $`C_9^{eff}+C_{10}+C_{LR}.`$
The terms with coefficients $`C_{LRLR}`$, $`C_{RLLR}`$, $`C_{LRRL}`$ and $`C_{RLRL}`$ describe the scalar type interactions. The last two terms, with the coefficients $`C_T`$ and $`C_{TE}`$, describe the tensor type interactions.
Now that we have the effective Hamiltonian, describing the $`bs\mathrm{}^+\mathrm{}^{}`$ decay at a scale $`\mu m_B`$, we can write down the matrix elements for the $`BK\mathrm{}^+\mathrm{}^{}`$ decay. The matrix element for this decay can be obtained by sandwiching the effective Hamiltonian between $`B`$ and $`K`$ meson states; which are parameterized in terms of form-factors which depend on the momentum transfer squared, $`q^2=(p_Bp_K)^2=(p_+p_{})^2`$. It follows from Eq.(9) that in order to calculate the amplitude of the $`BK\mathrm{}^+\mathrm{}^{}`$ decay the following matrix elements are required;
$$K\left|\overline{s}\gamma _\mu b\right|B,K\left|\overline{s}i\sigma _{\mu \nu }q^\nu b\right|B,K\left|\overline{s}b\right|B,K\left|\overline{s}\sigma _{\mu \nu }b\right|B.$$
These matrix elements are defined as follows ;
$`K(p_K)\left|\overline{s}\gamma _\mu b\right|B(p_B)=f_+\left[(p_B+p_K)_\mu {\displaystyle \frac{m_B^2m_K^2}{q^2}}q_\mu \right]+f_0{\displaystyle \frac{m_B^2m_K^2}{q^2}}q_\mu ,`$ (10)
$`K(p_K)\left|\overline{s}\sigma _{\mu \nu }b\right|B(p_B)=i{\displaystyle \frac{f_T}{m_B+m_K}}\left[(p_B+p_K)_\mu q_\nu q_\mu (p_B+p_K)_\nu \right].`$ (11)
Note that the finiteness of Eq.(9) at $`q^2=0`$ is guaranteed by assuming that $`f_+(0)=f_0(0)`$.
The matrix elements $`K(p_K)\left|\overline{s}i\sigma _{\mu \nu }q^\nu b\right|B(p_B)`$ and $`K\left|\overline{s}b\right|B`$ can be obtained from Eqs.(10) and (11) by multiplying both sides of these equations by $`q^\mu `$ and using the equations of motion, we get;
$`K(p_K)\left|\overline{s}b\right|B(p_B)`$ $`=`$ $`f_0{\displaystyle \frac{m_B^2m_K^2}{m_bm_s}},`$ (12)
$`K(p_K)\left|\overline{s}i\sigma _{\mu \nu }q^\nu b\right|B(p_B)`$ $`=`$ $`{\displaystyle \frac{f_T}{m_B+m_K}}\left[(p_B+p_K)_\mu q^2q_\mu (m_B^2m_K^2)\right].`$ (13)
As we have already mentioned the form-factors entering Eqs.(10)-(13) represent the hadronization process, where in order to calculate these form-factors information about the nonperturbative region of QCD is required. Therefore for the estimation of the form-factors to be reliable a nonperturbative approach is needed. Among the nonperturbative approaches the QCD sum rule is more predictive in studying the properties of hadrons. The form-factors appearing in the $`BK`$ transition are computed in the framework of the three point QCD sum rules and in the light cone QCD sum rules. We will use the result of the work in where radiative corrections to the leading twist wave functions and $`SU(3)`$ breaking effects are taken into account. As a result the form-factors are parameterized in the following way ;
$$f_i(q^2)=\frac{r_1}{1q^2/m_1^2}+\frac{r_2}{(1q^2/m_1^2)^2},$$
(14)
where $`1=+`$ or $`T`$, and
$$f_0(q^2)=\frac{r_2}{1q^2/m_{fit}^2},$$
(15)
with $`m_1=5.41`$GeV and the other parameters as given in Table 1.
Using the definition of the form factors given in Eqs.(10)-(13) we arrive at the following matrix element for the $`BK\mathrm{}^+\mathrm{}^{}`$ decay;
$`(BK\mathrm{}^+\mathrm{}^{})`$ $`=`$ $`{\displaystyle \frac{G_F\alpha }{4\sqrt{2}\pi }}V_{tb}V_{ts}^{}\{\overline{\mathrm{}}\gamma ^\mu \mathrm{}[A(p_B+p_K)_\mu +Bq_\mu ]`$ (16)
$`+\overline{\mathrm{}}\gamma ^\mu \gamma _5\mathrm{}\left[C(p_B+p_K)_\mu +Dq_\mu \right]+\overline{\mathrm{}}\mathrm{}Q+\overline{\mathrm{}}\gamma _5\mathrm{}N`$
$`+4\overline{\mathrm{}}\sigma ^{\mu \nu }\mathrm{}(iG)\left[(p_B+p_K)_\mu q_\nu (p_B+p_K)_\nu q_\mu \right]`$
$`+4\overline{\mathrm{}}\sigma ^{\alpha \beta }\mathrm{}ϵ_{\mu \nu \alpha \beta }H[(p_B+p_K)_\mu q_\nu (p_B+p_K)_\nu q_\mu ]\}.`$
The functions entering Eq.(16) are defined as;
$`A`$ $`=`$ $`(C_{LL}^{tot}+C_{LR}^{tot}+C_{RL}+C_{RR})f_++2(C_{BR}+C_{SL}){\displaystyle \frac{f_T}{m_B+m_K}},`$
$`B`$ $`=`$ $`(C_{LL}^{tot}+C_{LR}^{tot}+C_{RL}+C_{RR})f_{}2(C_{BR}+C_{SL}){\displaystyle \frac{f_T}{(m_B+m_K)q^2}}(m_B^2m_K^2),`$
$`C`$ $`=`$ $`(C_{LR}^{tot}+C_{RR}C_{LL}^{tot}C_{RL})f_+,`$
$`D`$ $`=`$ $`(C_{LR}^{tot}+C_{RR}C_{LL}^{tot}C_{RL})f_{},`$
$`Q`$ $`=`$ $`f_0{\displaystyle \frac{m_B^2m_K^2}{m_bm_s}}(C_{LRLR}+C_{RLLR}+C_{LRRL}+C_{RLRL}),`$
$`N`$ $`=`$ $`f_0{\displaystyle \frac{m_B^2m_K^2}{m_bm_s}}(C_{LRLR}+C_{RLLR}C_{LRRL}C_{RLRL}),`$
$`G`$ $`=`$ $`{\displaystyle \frac{C_T}{m_B+m_K}}f_T,`$
$`H`$ $`=`$ $`{\displaystyle \frac{C_{TE}}{m_B+m_K}}f_T,`$ (17)
where
$`f_{}=(f_0f_+){\displaystyle \frac{m_B^2m_K^2}{q^2}}.`$
From Eq.(16) it follows that the difference from the SM is due to the last four terms only, namely the scalar and tensor type interactions. For an analysis of the CP asymmetry it is necessary to compute the differential decay width for $`BK\mathrm{}^+\mathrm{}^{}`$. From the expression of the matrix element given in Eq.(16) we calculate the following result for the dilepton invariant mass spectrum;
$`{\displaystyle \frac{d\mathrm{\Gamma }}{d\widehat{s}}}(BK\mathrm{}^+\mathrm{}^{})={\displaystyle \frac{G^2\alpha ^2m_B}{2^{14}\pi ^5}}\left|V_{tb}V_{ts}^{}\right|^2\lambda ^{1/2}(1,\widehat{r}_K,\widehat{s})v\mathrm{\Delta }(\widehat{s}),`$ (18)
where $`\lambda (1,\widehat{r}_K,\widehat{s})=1+\widehat{r}_K^2+\widehat{s}^22\widehat{r}_K2\widehat{s}2\widehat{r}_K\widehat{s}`$, $`\widehat{s}=q^2/m_B^2`$, $`\widehat{r}_K=m_K^2/m_B^2`$, $`\widehat{m}_{\mathrm{}}=m_{\mathrm{}}/m_B`$, $`v=\sqrt{14\widehat{m}_{\mathrm{}}^2/\widehat{s}}`$ is the final lepton velocity, and $`\mathrm{\Delta }(\widehat{s})`$ is;
$`\mathrm{\Delta }`$ $`=`$ $`{\displaystyle \frac{4m_B^2}{3}}\text{Re}[96\lambda m_B^3\widehat{m}_{\mathrm{}}(AG^{})+24m_B^2\widehat{m}_{\mathrm{}}^2(1\widehat{r}_K)(CD^{})+12m_B\widehat{m}_{\mathrm{}}(1\widehat{r}_K)(CN^{})`$ (19)
$`+12m_B^2\widehat{m}_{\mathrm{}}^2\widehat{s}\left|D\right|^2+3\widehat{s}\left|N\right|^2+12m_B\widehat{m}_{\mathrm{}}\widehat{s}(DN^{})+256\lambda m_B^4\widehat{s}v^2\left|H\right|^2+\lambda m_B^2(3v^2)\left|A\right|^2`$
$`+s3\widehat{s}v^2|Q|^2+64\lambda m_B^4\widehat{s}(32v^2)|G|^2+m_B^2\{2\lambda (1v^2)[2\lambda 3(1\widehat{r}_K)^2]\}\left|C|^2\right].`$
As we have already mentioned, our goal in this work is the study of possible CP violating asymmetries beyond the SM in the $`BK\mathrm{}^+\mathrm{}^{}`$ decay; at this point we shall briefly remind the reader of the situation in the SM. In the SM the $`C_9`$ Wilson coefficient is the only one to have strong and weak phases. Strong phases arise from the short distance effects and resonances whereas the weak phase comes from the CKM elements. The remaining two coefficients, $`C_7`$ and $`C_{10}`$, are strictly real within the SM. From the parameterization of the form-factors it follows that they are inherently real and thus the imaginary parts in the functions in Eq.(19) can come only from the Wilson coefficients in Eq.(9). By strong and weak phases we mean the phases which are CP even and odd respectively. In other words we shall consider the picture where CP violating effects due to the short distance dynamics are parameterized by the Wilson coefficients. In principle all Wilson coefficients can have nonzero strong and weak phases .
In general the amplitude for $`\overline{B}K`$ has the general form ;
$$A(\overline{B}K)=e^{i\varphi _1}A_1e^{i\delta _1}+e^{i\varphi _2}A_2e^{i\delta _2},$$
(20)
where the strong phases are labeled as $`\delta `$’s and the weak phases by $`\varphi `$’s. As noted above the strong phases are CP even, whereas weak phases are odd under CP. Thus we arrive at an amplitude for the conjugated process, $`B\overline{K}`$, from Eq.(20);
$$\overline{A}(B\overline{K})=e^{i\varphi _1}A_1e^{i\delta _1}+e^{i\varphi _2}A_2e^{i\delta _2},$$
(21)
where the amplitudes of the decay rate of particle and anti-particle can be defined by the CP asymmetry (in the decay rate) as;
$$A_{CP}=\frac{|A|^2|\overline{A}|^2}{|A|^2+|\overline{A}|^2}=\frac{2A_1A_2Sin(\varphi _1\varphi _2)Sin(\delta _1\delta _2)}{A_1^2+2A_1A_2Cos(\varphi _1\varphi _2)Cos(\delta _1\delta _2)+A_2^2}.$$
(22)
Note that from the above expression we observe that in order to have CP asymmetry we should have both strong and weak phases in the amplitude; where the strong phases are provided by $`C_9^{eff}`$. In the SM the weak phases for the $`bs\mathrm{}^+\mathrm{}^{}`$ transition are negligible and hence the CP asymmetry for processes based on the quark level transitions, $`bs\mathrm{}^+\mathrm{}^{}`$, are highly suppressed. We will now consider the CP asymmetry in the decay width which is defined as;
$$A_{CP}(q^2)=\frac{{\displaystyle \frac{d\mathrm{\Gamma }}{d\widehat{s}}}(\overline{B}K\mathrm{}^+\mathrm{}^{}){\displaystyle \frac{d\mathrm{\Gamma }}{d\widehat{s}}}(B\overline{K}\mathrm{}^+\mathrm{}^{})}{{\displaystyle \frac{d\mathrm{\Gamma }}{d\widehat{s}}}(\overline{B}K\mathrm{}^+\mathrm{}^{})+{\displaystyle \frac{d\mathrm{\Gamma }}{d\widehat{s}}}(B\overline{K}\mathrm{}^+\mathrm{}^{})}.$$
(23)
Note that one can also have a CP asymmetry from the Forward-Backward (FB) asymmetry. However, in our present case the FB asymmetry for $`BK\mathrm{}^+\mathrm{}^{}`$ vanishes within the SM.
We shall now consider the minimal extension of these Wilson coefficients. In this approach we shall assume that the Wilson coefficients corresponding to scalar and tensor type interactions vanish identically (of course in the general case we can consider all Wilson coefficients with an arbitrary weak phase). For scalar type operators which emerge in Supersymmetric (SUSY) models and two Higgs Doublet Models (2HDM) this assumption is justified when we have electrons or muons in the final state. The reason being that in SUSY and 2HDM these operators originate from an Higgs exchange which results in Wilson coefficients which are proportional to $`m_{\mathrm{}}`$, and hence negligible for $`\mathrm{}=e,\mu `$.
The Wilson coefficients for the dipole operator obeys;
$$C_{BR}=2C_7^{eff}m_b,C_{SL}=2C_7^{eff}m_s,$$
(24)
with
$$C_7^{eff}=|C_7^{eff}|exp(i\varphi _7),$$
where $`\varphi _7`$ is an arbitrary phase and it is not constrained by the already observed branching ratio $`Br(BK^{}\gamma )`$.
Regarding the appearance of the new weak phase in $`C_{10}`$ we feel that a few words are in order. One of the possible discrepancies between the experimental results and the theoretical prediction for $`B\pi K`$ (from the $`B\pi \pi `$ data) can be resolved, as proposed in , by introducing a complex phase in the Wilson coefficient $`C_{10}=C_{10}^{SM}\mathrm{exp}(i\varphi _{10})`$. In this prescription the weak phase given to $`C_{10}`$ does not effect the CP asymmetry in $`BK\mathrm{}^+\mathrm{}^{}`$.
We will assume that the Wilson coefficients $`C_{RL}`$ and $`C_{RR}`$ also have weak phases, that is;
$`C_{RL}`$ $`=`$ $`|C_{RL}|\mathrm{exp}(i\varphi _{RL}),`$
$`C_{RR}`$ $`=`$ $`|C_{RR}|\mathrm{exp}(i\varphi _{RR}).`$ (25)
The Wilson coefficient $`C_9^{eff}(m_b,q^2)`$ has a finite phase, where, in order to better appreciate this, we write its explicit phase content as;
$`C_9^{eff}(m_b)=C_9(m_b)\left\{1+{\displaystyle \frac{\alpha _s\left(\mu \right)}{\pi }}\omega \left(\widehat{s}\right)\right\}+Y_{SD}(m_b,\widehat{s})+Y_{LD}(m_b,\widehat{s}),`$ (26)
where $`C_9(m_b)=4.334`$. Here $`\omega \left(\widehat{s}\right)`$ represents the $`𝒪(\alpha _s)`$ corrections coming from the four quark operator $`𝒪_9`$ ;
$`\omega \left(\widehat{s}\right)`$ $`=`$ $`{\displaystyle \frac{2}{9}}\pi ^2{\displaystyle \frac{4}{3}}Li_2\left(\widehat{s}\right){\displaystyle \frac{2}{3}}\mathrm{ln}\left(\widehat{s}\right)\mathrm{ln}\left(1\widehat{s}\right){\displaystyle \frac{5+4\widehat{s}}{3\left(1+2\widehat{s}\right)}}\mathrm{ln}\left(1\widehat{s}\right)`$ (27)
$`{\displaystyle \frac{2\widehat{s}\left(1+\widehat{s}\right)\left(12\widehat{s}\right)}{3\left(1\widehat{s}\right)^2\left(1+2\widehat{s}\right)}}\mathrm{ln}\left(\widehat{s}\right)+{\displaystyle \frac{5+9\widehat{s}6\widehat{s}^2}{3\left(1\widehat{s}\right)\left(1+2\widehat{s}\right)}}.`$
In Eq.(26) $`Y_{SD}`$ and $`Y_{LD}`$ represent, respectively, the short and long distance contributions to the four quark operators $`𝒪_{i=1,\mathrm{},6}`$ . Here $`Y_{SD}`$ can be obtained by a perturbative calculation;
$`Y_{SD}(m_b,\widehat{s})`$ $`=`$ $`g(\widehat{m}_c,\widehat{s})\left[3C_1+C_2+3C_3+C_4+3C_5+C_6\right]`$ (28)
$`{\displaystyle \frac{1}{2}}g(1,\widehat{s})\left[4C_3+4C_4+3C_5+C_6\right]`$
$`{\displaystyle \frac{1}{2}}g(0,\widehat{s})\left[C_3+3C_4\right]+{\displaystyle \frac{2}{9}}\left[3C_3+C_4+3C_5+C_6\right]`$
$`{\displaystyle \frac{V_{us}^{}V_{ub}}{V_{ts}^{}V_{tb}}}\left[3C_1+C_2\right]\left[g(0,\widehat{s})g(\widehat{m}_c,\widehat{s})\right],`$
where the loop function $`g(m_q,s)`$ represents the loops of quarks with mass $`m_q`$ at the dilepton invariant mass $`s`$. This function develops absorptive parts for dilepton energies $`s=4m_q^2`$;
$`g(\widehat{m}_q,\widehat{s})`$ $`=`$ $`{\displaystyle \frac{8}{9}}\mathrm{ln}\widehat{m}_q+{\displaystyle \frac{8}{27}}+{\displaystyle \frac{4}{9}}y_q{\displaystyle \frac{2}{9}}\left(2+y_q\right)\sqrt{\left|1y_q\right|}`$ (29)
$`\times \left\{\mathrm{\Theta }(1y_q)\left(\mathrm{ln}{\displaystyle \frac{1+\sqrt{1y_q}}{1\sqrt{1y_q}}}i\pi \right)+\mathrm{\Theta }(y_q1)2\mathrm{arctan}{\displaystyle \frac{1}{\sqrt{y_q1}}}\right\},`$
where $`\widehat{m}_q=m_q/m_b`$ and $`y_q=4\widehat{m}_q^2/\widehat{s}`$. Therefore, due to the extension of the absorptive parts of $`g(\widehat{m}_q,\widehat{s})`$ we see that the strong phases come from $`Y_{SD}`$. In particular one notices that the terms proportional to $`g(0,\widehat{s})`$ have a non-vanishing imaginary part independent of the dilepton invariant mass.
In addition to these perturbative contributions the $`\overline{c}c`$ loops can excite low-lying charmonium states $`\psi (1s),\mathrm{},\psi (6s)`$ whose contributions are represented by $`Y_{LD}`$ ;
$`Y_{LD}(m_b,\widehat{s})`$ $`=`$ $`{\displaystyle \frac{3}{\alpha ^2}}\left\{{\displaystyle \frac{V_{cs}^{}V_{cb}}{V_{ts}^{}V_{tb}}}C^{\left(0\right)}{\displaystyle \frac{V_{us}^{}V_{ub}}{V_{ts}^{}V_{tb}}}\left[3C_3+C_4+3C_5+C_6\right]\right\}`$ (30)
$`\times {\displaystyle \underset{V_i=\psi \left(1s\right),\mathrm{},\psi \left(6s\right)}{}}{\displaystyle \frac{\pi \kappa _i\mathrm{\Gamma }\left(V_i\mathrm{}^+\mathrm{}^{}\right)M_{V_i}}{\left(M_{V_i}^2\widehat{s}m_b^2iM_{V_i}\mathrm{\Gamma }_{V_i}\right)}},`$
where $`\kappa _i`$ is a phenomenological parameter taken here to be 2.3 so as to produce the correct branching ratio of $`Br(BJ/\psi K^{}K^{}\mathrm{}\mathrm{})=Br(BJ/\psi K^{})Br(J/\psi \mathrm{}\mathrm{})`$ , and $`C^{\left(0\right)}3C_1+C_2+3C_3+C_4+3C_5+C_6=0.362`$. Contrary to $`Y_{SD}`$ the long-distance contribution in $`Y_{LD}`$ has both weak and strong phases. The weak phases follow from the CKM elements whereas the strong phases come from the $`\widehat{s}`$ values for which the $`i`$-th charmonium states are on shell. Therefore, the Wilson coefficient $`C_9^{eff}(m_b)`$ has both weak and strong phases already in the SM.
In this sense the Wilson coefficients $`C_7^{eff}(m_b)`$ and $`C_{10}(m_b)`$ can not develop any strong phase, and thus, $`\varphi _7`$ and $`\varphi _{10}`$ should necessarily originate from physics beyond the SM. As such the phases of $`\varphi _7`$ and $`\varphi _{10}`$ can be chosen to have a purely weak character.
## 3 Numerical analysis
In this section we present our numerical results for the asymmetries $`A_{CP}`$ for the $`BK\mu ^+\mu ^{}`$ decay. Note that the parameters for the hadronic form-factors are taken from Table I. For values of the Wilson coefficients in the SM we have used $`C_3=0.011`$, $`C_4=0.026`$, $`C_5=0.007`$, $`C_6=0.031`$, $`C_7^{eff}=0.313`$, $`C_9=4.344`$, and $`C_{10}=4.664`$. For further numerical analysis the values of the new Wilson coefficients are needed, where we have varied them in the range $`|C_{10}|<C_X<|C_{10}|`$. The experimental value of the branching ratio of the $`BK(K^{})\mathrm{}^+\mathrm{}^{}`$ decays and the bound on $`Br(B\mu ^+\mu ^{})`$ suggest that this is the right order of magnitude. It should be noted that the experimental results lead to strong restrictions on some of the Wilson coefficients, namely $`2C_{LL}`$ and $`C_{RL}2.3`$, while the remaining coefficients vary in the range $`|C_{10}|<C_X<|C_{10}|`$. For the remaining parameters we take $`m_b=4.8`$GeV, $`m_c=1.35`$GeV, $`m_B=5.28`$GeV and $`m_K=0.496`$GeV.
For the kinematical interval the dilepton invariant mass is $`4m_{\mathrm{}}^2q^2(m_Bm_K)^2`$ where the $`J/\psi `$ family of resonances can be excited. The dominant contribution comes from the three low-lying resonances $`J/\psi ,\psi ^{^{}},\psi ^{^{\prime \prime }}`$ in the interval $`8`$GeV$`{}_{}{}^{2}\stackrel{<}{_{}}q^2\stackrel{<}{_{}}14.5`$GeV<sup>2</sup>. In order to minimize the hadronic uncertainties we will discard this subinterval in the analysis below by dividing the $`q^2`$ region in to $`low`$ and $`high`$ dilepton mass intervals;
$`\text{Region I}:4m_{\mathrm{}}^2q^28\mathrm{GeV}^2,`$
$`\text{Region II}:14.5\mathrm{GeV}^2q^2(m_Bm_K)^2,`$ (31)
where the contribution of the higher resonances do still exist in the second region.
As mentioned previously, we have analyzed the case where there are four weak phases: $`\varphi _7`$, $`\varphi _{10}`$, $`\varphi _{RL}`$ and $`\varphi _{RR}`$.
In fig. 1 we have presented $`A_{CP}`$ in the $`\varphi _7`$$`q^2`$ plane for the $`BK\mu ^+\mu ^{}`$ decay for Region I and Region II respectively. In Region I the CP asymmetry is practically independent of $`q^2`$, becoming maximal in the value for CP violation at $`\varphi _7=\pi /2`$. In Region II, however, the $`q^2`$ dependence is comparatively enhanced as the dominance of the dipole coefficient is now reduced. Aside from this our figures suggest that the CP asymmetry in Region II is four times larger than in Region I, and this confirms our earlier expectation.
Since the CP asymmetry is dependent on $`q^2`$ and the new weak phases there can appear some difficulties. The dependence of one of the variables, for example $`q^2`$, can be removed by integrating over $`q^2`$ in the allowed practical kinematical region, where the averaged asymmetries could be measured more easily experimentally. Therefore we shall now discuss only averaged CP asymmetries, which we define in the following way. That is, our averaging procedure is defined by;
$$A_{CP}=\frac{{\displaystyle _{R_i}}A_{CP}{\displaystyle \frac{d\mathrm{\Gamma }}{dq^2}}𝑑q^2}{{\displaystyle _{R_i}}{\displaystyle \frac{d\mathrm{\Gamma }}{dq^2}}𝑑q^2}.$$
(32)
where $`R_i`$ means Region I or II.
We now depict in fig. 2 the $`\varphi _7`$ dependence of the averaged asymmetries $`A_{CP}`$. From this figure it can be observed that the average $`CP`$ asymmetry can attain values of $`3\%`$. Differences from zero of any value of $`A_{CP}`$ would be an unambiguous indication of the existence physics beyond the SM.
In fig. 3 we have plotted the dependence of CP asymmetry on the dilepton invariant mass and $`\varphi _{RL}`$. In fig. 5 we have shown the same kind of plot but for $`\varphi _{RR}`$. We have also shown the correlation of averaged CP asymmetry and the integrated branching ratios. In fig. 4 the variation of $`A_{CP}`$ with integrated branching ratio for $`BK\mu ^+\mu ^{}`$ for $`C_{RL}`$ is shown. In this figure we have used three different values of $`C_{RL}`$ and have varied the phase ($`\varphi _{RL}`$) in the range $`0\varphi _{RL}2\pi `$. All the other Wilsons are taken to have their SM values. In a similar graph, given in fig. 6, we have varied $`C_{RR}`$.
In the present work we have studied the sensitivity of the CP violating asymmetry on the new weak phases appearing in the Wilson coefficients. We have also observed that the CP asymmetry in Region II is 4-5 times larger than that observed in Region I when we consider a weak phase $`\varphi _7`$. This can be understood in that in Region II contributions coming from other operators become comparable with the dipole operator $`O_7`$, where this operator is dominant in Region I. Having obtained the averaged $`A_{CP}`$ asymmetry we obtained a maximal value of approximately 3%. Note that an additional weak phase in $`C_{10}`$ will not give rise to any CP asymmetry, however, if non-standard<sup>5</sup><sup>5</sup>5by non-standard we mean operators which are not present within the SM electroweak operators are considered then the CP asymmetry in the region of high dilepton invariant mass can reach a value of up to 10%.
As stated earlier, we can also, in principle, have weak phases in scalar and pseudo-scalar operators. The presence of weak phases in these operators can also substantially effect the CP asymmetry. The popular extensions of the SM, such as SUSY and 2HDM, all predict the existence of such operators. However, the magnitude of these Wilson coefficients is predicted to be small when the lepton $`\mathrm{}=e`$ or $`\mu `$. In the presence of these operators one also gets a non-zero value for the FB asymmetry in $`BK\mathrm{}^+\mathrm{}^{}`$. The FB asymmetry could provide another measure of CP asymmetry which has not been considered in this work.
The observation of CP asymmetry in $`BK\mathrm{}^+\mathrm{}^{}`$ would not only tell us about the nature of weak phases but would also give us an insight in to the structure of the effective Hamiltonian. Therefore the measurement of the CP violating asymmetry would provide us with an useful insight into the mechanism of CP violation, which in turn would serve as a good test for physics beyond the SM.
## Acknowledgments
The work of SRC and NG was supported by the Department of Science & Technology (DST), India under grant no SP/S2/K-20/99. The work of ASC was supported by the Japan Society for the Promotion of Science (JSPS), under fellowship no P04764. |
warning/0506/quant-ph0506258.html | ar5iv | text | # Non-Markov dynamics and phonon decoherence of a double quantum dot charge qubit
## I Introduction
Solid state qubits are considered to be promising candidates for realizing building blocks of quantum information processors because they can be scaled up to large numbers. The double quantum dot (DQD) DQD1 ; DQD2 ; DQD3 ; DQD4 ; Brandes charge qubit is one of these qubits. Two low-energy charge states are used as the local states $`|0`$ and $`|1`$ in the qubit. The qubit can be controlled directly via external voltage sources. There are some effective schemes to prepare the initial states and read out the final states of the qubit prepare-and-measure . So it is considered that decoherence may be the central impediment for the qubit to be taken as the cell of quantum computer. Finding out the primary origin or the dominating mechanism of decoherence for the qubit is a basal task for overcoming the impediment. Experimentally, many attempts experiment1 ; experiment2 for detecting the decoherence time of this kind of qubit have been performed. The decoherence has also been investigated theoretically. In 2000, Fedichkin *et al.* Fedichkin-etal01 investigated the Born-Markov type electron-phonon decoherence at large times due to spontaneous phonon emission of the quantum dot charge qubits. Recently, Vorojtsov *et al.* Vorojtsov-etal studied the decoherence of the DQD charge qubit by using Born-Markov approximation. But, as it has been pointed out, the use of the Born-Markov approximation is inappropriate at large tunneling amplitudes. The method is expected to become increasingly unreliable at DQD with larger interdot tunneling amplitudes. Wu *et al.* Wu-etal investigated the decoherence in terms of a perturbation treatment based on a unitary transformation. The Born-Markov approximation has not been used in the method but it neglects some terms of the effective Hamiltonian with high excited states. This kind of processing introduces an new approximation which has not been estimated to the affects of the dynamics. Fedichkin *et al.* Fedichkin01 ; Fedichkin02 studied the error rate of DQD charge qubit with short-time approximation. This method is accurate enough in adequate short time. But the decoherence in a moderately long time is also interesting. Recently, Thorwart *et al.* Thorwartetal2 investigated the decoherence of the DQD charge qubit in a longer time with a numerically exact iterative quasiadiabatic propagator path integral (QUAPI) Makri . This method is proved valid in investigating the qubit decoherence Thorwartetal1 . In Ref. Thorwartetal2 , Thorwart *et al.* considered the coupling of longitudinal piezoelectic acoustic phonons with the investigated qubit and neglected the contribution of the deformation acoustic phonons to decoherence. These two kinds of phonons may constitute two kinds of different coupling baths in the environment of the qubit. We call the former the piezoelectric coupling phonon bath (PCPB) and the latter the deformation coupling phonon bath (DCPB). Comparing Thorwart’s result and the reported experimental value they found that the theory predicts the decoherence time of the DQD charge qubit is two orders of magnitudes smaller than the experimental one. Thus, Thorwart *et al.* conclude thatthe piezoelectric coupling phonon decoherence is a subordinate mechanism in decoherence of the DQD charge qubit. Recently, Wu *et al.* Wu-etal gave the spectral density functions of PCPB as well as DCPB. Then how about the DCPB to the decoherence of the DQD charge qubit? In other words, is the decoherence of the DQD charge qubit induced by DCPB also subordinate? In this paper we shall use an iterative tensor multiplication (ITM) Makri scheme derived from the QUAPI to study the decoherence times of the DQD charge qubit not only in PCPB but also in DCPB. In order to validate if our result is in accordance with Thorwart’s result we at first investigate the decoherence times of the qubit in PCPB. Then we shall investigate the decoherence times of the qubit in another bath, DCPB, which will show that the influence of the DCPB to the decoherence of the DQD charge qubit is also subordinate because it results in a shorter decoherence time than the experimental value of $`1`$ $`ns`$ experiment1 ; experiment2 .
## II Models
The DQD charge qubit consists of left and right dots connected through an interdot tunneling barrier. Due to Coulomb blockade, at most one excess electron is allowed to occupy the left and right dot, which defines two basis vectors $`|0`$ and $`|1.`$ The energy difference $`\epsilon `$ between these two states can be controlled by the source-drain voltage. Neglecting the higher order tunneling between leads and the dots, the effective Hamiltonian in the manipulation process reads Wu-etal ; Thorwartetal2
$$H_{eff}=\mathrm{}T_c\sigma _x+\mathrm{}\underset{q}{}\omega _qb_q^{}b_q+\mathrm{}\sigma _z\underset{q}{}\left(M_qb_q^{}+M_q^{}b_q\right).$$
(1)
Here, $`T_c`$ is the interdot tunneling, $`\sigma _x`$ and $`\sigma _z`$ are Pauli matrix, $`b_q^{}`$ $`(b_q)`$ are the creation (annihilation) operators of phonons, $`\mathrm{}\omega _q`$ is the energy of the phonons, and $`M_q=C_q/\sqrt{2m_q\omega _q\mathrm{}}`$ where $`C_q`$ are the classical coupling constants of qubit-phonon system. We call the collective coupling phonons to the qubit in the environment a phonon bath. In order to obtain the reduced density matrix of the qubit in the system, one should know the coupling coefficients $`M_q`$, but in fact we need not know the details of each $`M_q`$ because all characteristics of the bath pertaining to the dynamics of the observable system are captured in the spectral density function Weiss ; Leggett
$$J\left(\omega \right)=\underset{q}{}\left|M_q\right|^2\delta \left(\omega \omega _q\right).$$
(2)
It is pointed out that the spectral density of PCPB is
$$J^{pz}\left(\omega \right)=g_{pz}\omega \left(1\frac{\omega _d}{\omega }\mathrm{sin}\frac{\omega }{\omega _d}\right)e^{\frac{\omega ^2}{2\omega _l^2}}.$$
(3)
Here, $`\omega _d=s/d`$ and $`\omega _l=s/l`$, where $`d`$ denotes the center-to-center distance of two dots, $`l`$ the dot size, $`s`$ the sound velocity in the crystal, and
$$g_{pz}=\frac{M}{\pi ^2\varrho s^3}\left(\frac{6}{35}+\frac{1}{x}\frac{8}{35}\right).$$
Here, $`M`$ is the piezoconstant, $`\varrho `$ is the density of the crystal, and $`x`$ is the rate of transverse to the longitudinal of sound velocity in the crystal, (see for example Refs. Wu-etal ; Fedichkin01 ). As in Refs. Wu-etal ; Fedichkin01 in this paper we set the sound velocity in the GaAs crystal $`s5\times 10^3`$ $`m/s`$. With the parameters of GaAs in Ref. Mahan-etal , Wu *et al.* Wu-etal propose a value $`g_{pz}0.035`$ $`(ps)^2`$. The spectral density of DCPB is
$$J^{df}\left(\omega \right)=g_{df}\omega ^3\left(1\frac{\omega _d}{\omega }\mathrm{sin}\frac{\omega }{\omega _d}\right)e^{\frac{\omega ^2}{2\omega _l^2}},$$
(4)
where
$$g_{df}=\frac{\mathrm{\Xi }^2}{8\pi ^2\varrho s^5}.$$
Here, $`\mathrm{\Xi }`$ is the deformation potential. In the same paper, Wu *et al.* also propose a value $`g_{df}0.029`$ $`(ps)^2.`$ One can investigate the dynamics and then the decoherence of the open qubit with the help of the definite spectral density functions of the baths. Before investigations of decoherence of the DQD charge qubit we introduce an optimal numerical path integral method, the ITM method in the following section.
## III QUAPI and ITM
In the following, we firstly review the QUAPI and then the ITM Makri scheme. Suppose the initial state of the qubit-bath system has the form
$$R\left(0\right)=\rho \left(0\right)\rho _{bath}\left(0\right),$$
(5)
where $`\rho \left(0\right)`$ and $`\rho _{bath}\left(0\right)`$ are the initial states of the qubit and bath. The evolution of the reduced density operator of the open qubit
$$\stackrel{~}{\rho }(s^{\prime \prime },s^{};t)=\text{Tr}_{bath}s^{\prime \prime }\left|e^{iHt/\mathrm{}}\rho \left(0\right)\rho _{bath}\left(0\right)e^{iHt/\mathrm{}}\right|s^{},$$
(6)
is given by
$`\stackrel{~}{\rho }(s^{\prime \prime },s^{};t)`$
$`=`$ $`{\displaystyle \underset{s_0^+=\pm 1}{}}{\displaystyle \underset{s_1^+=\pm 1}{}}\mathrm{}{\displaystyle \underset{s_{N1}^+=\pm 1}{}}{\displaystyle \underset{s_0^{}=\pm 1}{}}{\displaystyle \underset{s_1^{}=\pm 1}{}}\mathrm{}{\displaystyle \underset{s_{N1}^{}=\pm 1}{}}`$
$`\times s^{\prime \prime }\left|e^{iH_0\mathrm{\Delta }t/\mathrm{}}\right|s_{N1}^+\mathrm{}s_1^+\left|e^{iH_0\mathrm{\Delta }t/\mathrm{}}\right|s_0^+`$
$`\times s_0^+\left|\rho \left(0\right)\right|s_0^{}`$
$`\times s_0^{}\left|e^{iH_0\mathrm{\Delta }t/\mathrm{}}\right|s_1^{}\mathrm{}s_{N1}^{}\left|e^{iH_0\mathrm{\Delta }t/\mathrm{}}\right|s^{}`$
$`\times I(s_0^+,s_1^+,\mathrm{},s_{N1}^+,s^{\prime \prime },s_0^{},s_1^{},\mathrm{},s_{N1}^{},s^{};\mathrm{\Delta }t),`$
where the influence functional is
$`I(s_0^+,s_1^+,\mathrm{},s_{N1}^+,s^{\prime \prime },s_0^{},s_1^{},\mathrm{},s_{N1}^{},s^{};\mathrm{\Delta }t)`$ (8)
$`=`$ $`\text{Tr}_{bath}[e^{iH_{env}\left(s^{\prime \prime }\right)\mathrm{\Delta }t/2\mathrm{}}e^{iH_{env}\left(s_{N1}^+\right)\mathrm{\Delta }t/2\mathrm{}}`$
$`\times \mathrm{}e^{iH_{env}\left(s_0^+\right)\mathrm{\Delta }t/2\mathrm{}}\rho _{bath}\left(0\right)e^{iH_{env}\left(s_0^{}\right)\mathrm{\Delta }t/2\mathrm{}}`$
$`\times \mathrm{}e^{iH_{env}\left(s_{N1}^{}\right)\mathrm{\Delta }t/2\mathrm{}}e^{iH_{env}(s)\mathrm{\Delta }t/2\mathrm{}}].`$
Here, $`H_0`$ is a reference Hamiltonian that in general depends on the coordinate and momentum of the system. In the qubit system, it usually depends on Pauli matrixes $`\sigma _x`$ and $`\sigma _z`$. The $`H_{env}`$ is defined as $`H_{env}=HH_0`$. In our system we set $`H_0=\mathrm{}T_c\sigma _x.`$ The discrete path integral representation of the qubit density matrix contains temporal nonlocal terms $`I(s_0^+,s_1^+,\mathrm{},s_{N1}^+,s^{\prime \prime },s_0^{},s_1^{},\mathrm{},s_{N1}^{},s^{};\mathrm{\Delta }t)`$ which denotes the process being non-Markovian. With the quasiadiabatic discretization of the path integral, the influence functional, Eq.(8) takes the form
$$I=\mathrm{exp}\left\{\frac{i}{\mathrm{}}_{k=0}^N_{k^{}=0}^k\left(s_k^+s_k^{}\right)\left(\eta _{kk^{}}s_k^{}^+\eta _{kk^{}}^{}s_k^{}^{}\right)\right\},$$
(9)
where $`s_N^+=s^{\prime \prime }`$ and $`s_N^{}=s^{}.`$ The coefficients $`\eta _{kk^{}}`$ can be obtained by substituting the discrete path into the Feynman-Vernon expression. Their expressions have been shown in Ref. Makri . Thus, the influence functional can be expressed with a product of terms corresponding to different $`\mathrm{\Delta }k`$ as
$`I`$ $`=`$ $`_{k=0}^NI_0\left(s_k^\pm \right)_{k=0}^{N1}I_1(s_k^\pm ,s_{k+1}^\pm )_{k=0}^{N\mathrm{\Delta }k}I_{\mathrm{\Delta }k}(s_k^\pm ,s_{k+\mathrm{\Delta }k}^\pm )`$ (10)
$`\mathrm{}_{k=0}^{N\mathrm{\Delta }k_{\mathrm{max}}}I_{\mathrm{\Delta }k_{\mathrm{max}}}(s_k^\pm ,s_{k+\mathrm{\Delta }k_{\mathrm{max}}}^\pm ).`$
Here, $`\mathrm{\Delta }k=kk^{},`$ where $`k^{}`$ and $`k`$ are points of discrete path integral expressions, (see Ref. Makri ) and
$`I_0\left(s_i^\pm \right)`$ $`=`$ $`\mathrm{exp}\left\{{\displaystyle \frac{1}{\mathrm{}}}\left(s_i^+s_i^{}\right)\left(\eta _{ii}s_i^+\eta _{ii}^{}s_i^{}\right)\right\},`$
$`I_{\mathrm{\Delta }k}(s_i^\pm ,s_{i+\mathrm{\Delta }k}^\pm )`$ $`=`$ $`\mathrm{exp}\{{\displaystyle \frac{1}{\mathrm{}}}(s_{i+\mathrm{\Delta }k}^+s_{i+\mathrm{\Delta }k}^{})`$
$`\times (\eta _{i+\mathrm{\Delta }k,i}s_i^+\eta _{i+\mathrm{\Delta }k,i}^{}s_i^{})\},\mathrm{\Delta }k1.`$
The length of the memory of the time can be estimated by the following bath response function
$$\alpha ^x\left(t\right)=\frac{1}{\pi }_0^{\mathrm{}}𝑑\omega J^x\left(\omega \right)\left[\mathrm{coth}\left(\frac{\beta \mathrm{}\omega }{2}\right)\mathrm{cos}\omega ti\mathrm{sin}\omega t\right].$$
(12)
Here, the superscript $`x`$ denotes the bath type, $`\beta =1/k_BT,`$ where $`k_B`$ is the Boltzmann constant, and $`T`$ is the temperature. It is shown that when the real and imaginary parts behave as the delta function $`\delta \left(t\right)`$ and its derivative $`\delta ^{}\left(t\right),`$ the dynamics of the reduced density matrix is Markovian. However, if the real and imaginary parts are broader than the delta function, the dynamics is non-Markovian. The broader the Re$`[\alpha ^x\left(t\right)]`$ and Im$`[\alpha ^x\left(t\right)]`$ are, the longer of the memory time will be. The broader the Re$`[\alpha ^x\left(t\right)]`$ and Im$`[\alpha ^x\left(t\right)]`$ are, the more serious the Markov approximation will distort the practical dynamics. In Fig.1 we plot the Re$`[\alpha ^{pz}\left(t\right)]`$ and Im$`[\alpha ^{pz}\left(t\right)]`$ of the PCPB and in Fig.2 we plot the Re$`[\alpha ^{df}\left(t\right)]`$ and Im$`[\alpha ^{df}\left(t\right)]`$ of the DCPB.
$`Fig.1,`$
$`Fig.2`$
We see that the memory times are about $`\tau _{mem}^{pz}=1\times 10^{11}`$ $`s`$ for PCPB and $`\tau _{mem}^{df}=2\times 10^{11}`$ $`s`$ for DCPB {where the points beyond $`\pm 1\times 10^{11}`$ $`s`$ have not been plotted for clearly distinguishing the Re$`[\alpha ^{df}\left(t\right)]`$ and Im$`[\alpha ^{df}\left(t\right)]`$ in the same figure}. Due to the nonlocality, it is impossible to calculate the reduced density matrix by Eq.(LABEL:eq8) in the matrix multiplication scheme. However, the short range nonlocality of the influence functional implies that the effects of the nonlocality should drop off rapidly as the “interaction distance” increases. In the ITM scheme the interaction can be taken into account at each iteration step. The reduced density matrix at time $`t=N\mathrm{\Delta }t`$ ($`N`$ even) is given as
$$\stackrel{~}{\rho }(s_N^\pm ,N\mathrm{\Delta }t)=A^{\left(1\right)}(s_N^\pm ;N\mathrm{\Delta }t)I_0\left(s_N^\pm \right),$$
where
$`A^{\left(1\right)}(s_{k+1}^\pm ;(k+1)\mathrm{\Delta }t)`$ $`=`$ $`{\displaystyle 𝑑s_k^\pm T^{\left(2\right)}(s_k^\pm ,s_{k+1}^\pm )}`$
$`\times A^{\left(1\right)}(s_k^\pm ;k\mathrm{\Delta }t).`$
Here,
$`T^{\left(2\mathrm{\Delta }k_{\mathrm{max}}\right)}(s_k^\pm ,s_{k+1}^\pm \mathrm{}s_{k+2\mathrm{\Delta }k_{\mathrm{max}}1}^\pm )`$
$`=`$ $`_{n=k}^{k+\mathrm{\Delta }k_{\mathrm{max}}1}K(s_k^\pm ,s_{k+1}^\pm )I_0\left(s_n^\pm \right)I_1(s_n^\pm ,s_{n+1}^\pm )`$
$`\times I_2(s_n^\pm ,s_{n+2}^\pm )\mathrm{}I_{\mathrm{\Delta }k_{\mathrm{max}}}(s_n^\pm ,s_{n+\mathrm{\Delta }k_{\mathrm{max}}}^\pm ),`$
and
$$A^{\left(\mathrm{\Delta }k_{\mathrm{max}}\right)}(s_0^\pm ,s_1^\pm ,\mathrm{},s_{\mathrm{\Delta }k_{\mathrm{max}}1}^\pm ;0)=s_0^+\left|\rho _s\left(0\right)\right|s_0^{},$$
where
$`K(s_k^\pm ,s_{k+1}^\pm )`$ $`=`$ $`s_{k+1}^+\left|\mathrm{exp}(iH_0\mathrm{\Delta }t/\mathrm{})\right|s_k^+`$
$`\times s_k^{}\left|\mathrm{exp}(iH_0\mathrm{\Delta }t/\mathrm{})\right|s_{k+1}^{}.`$
In the ITM scheme a short-time approximation instead of the Markov approximation is used. The approximation makes an error of the short-time propagator in order $`\left(\mathrm{\Delta }t\right)^3,`$ which is small enough as we set the time step $`\mathrm{\Delta }t`$ very small. It is shown that when the time step $`\mathrm{\Delta }t`$ is not larger than the characteristic time of the qubit system, which can be calculated with $`1/T_c,`$ the calculation is accurate enough Privman . In particular, the scheme does not discard the memory of the temporal evolution, which may be appropriate to solve the decoherence of qubit. In the following section we shall use the ITM scheme to study the decoherence times of the DQD charge qubit in PCPB and DCPB.
## IV Decoherence of DQD charge qubit
To measure effects of decoherence one can use the entropy, the first entropy, and many other measures, such as maximal deviation norm, etc. (see for example Ref. Privman ). However, essentially, the decoherence of a open quantum system is reflected through the decays of the off-diagonal coherent terms of its reduced density matrix. The decoherence is in general produced due to the interaction of the quantum system with other systems with a large number of degrees of freedom, for example the devices of the measurement or environment. The decoherence time denoted by $`\tau _2`$ measures the time of the initial coherent terms to their $`1/e`$ times, namely, $`\rho _i(n,m)\stackrel{\tau _2}{}\rho _f(n,m)=\rho _i(n,m)/e.`$ Here, $`nm,`$ and $`n,`$ $`m=0`$ or $`1`$ for qubits. In this paper, we investigate the decoherence times via directly describing the evolutions of the off-diagonal coherent terms, instead of using any measure of decoherence. In our following investigations, we suppose the temperature $`T=30`$ $`mK`$ and the cut-off frequency of the bath modes $`\omega _C=5`$ $`\left(ps\right)^1`$. We set the initial state of the qubit to $`\rho \left(0\right)=\frac{1}{2}\left(|0+|1\right)\left(0|+1|\right),`$ which is a pure state and it has the maximum coherent terms, and the initial state of the environment is $`\rho _{bath}\left(0\right)=_ke^{\beta M_k}/`$Tr$`{}_{k}{}^{}\left(e^{\beta M_k}\right).`$ In the calculations we set $`\omega _d=0.02`$ $`(ps)^1,`$ $`T_c=0.1\omega _l`$ according to Ref. Wu-etal , and two kinds of cases $`\omega _l=0.5`$ $`(ps)^1`$ and $`\omega _l=0.7`$ $`(ps)^1`$ are calculated.
*Decoherence time obtained from ITM scheme:* In the following, at first, we use the ITM scheme investigating the decoherence time of the DQD charge qubit. The evolutions of the coherent elements of the reduced density matrix of the DQD charge qubit in PCPB and DCPB are plotted in Figs. 3 and 4. Here, we simply choose $`\mathrm{\Delta }k_{\mathrm{max}}=1`$ and $`\mathrm{\Delta }t=1\times 10^{11}`$ $`s`$ for PCPB and $`\mathrm{\Delta }t=2\times 10^{11}`$ $`s`$ for DCPB in the ITM scheme. These choices of the time steps are feasible as we consider that it should be not smaller than the memory times of the baths, because the latter is about $`\tau _{mem}^{pz}1\times 10^{11}`$ $`s`$ for PCPB and $`\tau _{mem}^{df}2\times 10^{11}`$ $`s`$ for DCPB (see Figs. 1 and 2). It is also appropriate as we consider that the time steps should not be larger than the characteristic time of the qubit, because the characteristic time of the qubit is about $`2\times 10^{11}s`$.
$`Fig.3,`$
$`Fig.4`$
Helped with detailed numerical analyses, we can obtain that the decoherence times of the DQD charge qubit in PCPB are about $`\tau _2^{pz}97`$ $`ps`$ \[when $`\omega _l=0.7`$ $`(ps)^1`$\] and $`\tau _2^{pz}118`$ $`ps`$ \[when $`\omega _l=0.5`$ $`(ps)^1`$\]. Similarly, we can obtain that the decoherence times of this qubit in DCPB are about $`\tau _2^{df}=1.04`$ $`ps`$ \[when $`\omega _l=0.7`$ $`(ps)^1`$\] and $`\tau _2^{df}=3.5`$ $`ps`$ \[when $`\omega _l=0.5`$ $`(ps)^1`$\]. It is shown that the DCPB behaves more destructively than the PCPB does to the coherence of the DQD charge qubit. A further calculation shows that the decoherence time will increase with the decreasing of $`T_c.`$
*Decoherence time calculated on Bloch equations:* It is well known that the decoherence time can be calculated based on Bloch equations. In the following, we calculate the decoherence time of the DQD charge qubit in PCPB and DCPB with the Bloch equation method. In this method, the relaxation and dephasing times can be evaluated from the spin-bosonic model with Bloch equations Leggett ; Weiss . For our model, they are BBS
$$\tau _1^1=\tau _2^1=\frac{1}{2\mathrm{}}J\left(\omega _0\right)\mathrm{coth}\left(\beta \mathrm{}\omega _0/2\right),$$
where $`\omega _0=2T_c`$ is the natural frequency of the DQD charge qubit. By using the parameters of the DQD charge qubit and PCPB bath as above, we can calculate the decoherence times with this method as $`\tau _2^{pz}122.3`$ $`ps`$ \[when $`\omega _l=0.7`$ $`\left(ps\right)^1`$\] and $`\tau _2^{pz}192.2`$ $`ps`$ \[when $`\omega _l=0.5`$ $`\left(ps\right)^1`$\]. Similarly, we can obtain the decoherence times of the DQD charge qubit in the DCPB with this method as $`\tau _2^{df}3.18`$ $`ps`$ \[when $`\omega _l=0.7`$ $`\left(ps\right)^1`$\] and $`\tau _2^{df}12.6`$ $`ps`$ \[when $`\omega _l=0.5`$ $`\left(ps\right)^1`$\]. It is shown that the decoherence times obtained from the ITM scheme are shorter than those obtained on Bloch equations. We suggest that the differences are derived from the different choices of approximation schemes. The Bloch equations are in general derived from the Markov approximation which discards the memory of baths in the derivation of dynamical evolution. The decoherence of the qubit obtained on Bloch equations is similar to the “resonant decoherence” Openov obtained from the Fermi golden rule. It is not accurately equal to the actual decoherence except that the “nonresonant decoherence” very small.
*Decoherence time derived from the quality factor:* We like to compare our results obtained from the ITM scheme based on QUAPI with Thorwart’s results which are also obtained from QUAPI. Thorwart *et al.* Thorwartetal2 investigated the PCPB case and they obtained the quality factor instead of the decoherence time. By using a set of parameters of the DQD charge qubit and the PCPB they obtained the quality factor of the qubit as $`Q_{pz}=336,`$ which corresponds to decoherence time $`\tau _2^{pz}=Q_{pz}\pi /\omega _{pz}^{}115.9`$ $`ps,`$ where $`\omega ^{}=\omega _0+\mathrm{\Delta }\omega ,`$ and $`\mathrm{\Delta }\omega `$ is the bath-induced shift BBS in the natural frequency $`\omega _0=2T_c.`$ From Fig.1 of Ref. Wu-etal we see $`\mathrm{\Delta }\omega _{pz}1.75\omega _c`$ and $`\mathrm{\Delta }\omega _{df}1.65\omega _c.`$ Their used parameters $`[T=10`$ $`mk,`$ $`T_c0.07\left(ps\right)^1]`$ have a little difference from ours. But we have calculated that the difference does not result in much decoherence time departure. It is meant that our results is in accordance with Thorwart’s result. On the other hand, from our decoherence time $`\tau _2^{df}3.5`$ $`ps`$ of the qubit in DCPB we can obtain its quality factor $`Q_{df}=\tau _2^{df}\omega _{df}^{}/\pi 8.`$
## V Discussions and conclusions
In this paper we investigated the decoherence times of the DQD charge qubit in PCPB and DCPB with the ITM scheme based on QUAPI. The decoherence times are also calculated based on Bloch equations. The results derived from the two kinds of methods are compared to each other. It is shown that the latter are longer than the former. We think this results from the different choices of approximation schemes because the Markov approximation used in the latter method discards the memory of the baths. On the other hand, Hayashi *et al.* experiment1 ; experiment2 have detected that the decoherence time of the DQD charge qubit is about $`1`$ $`ns`$ as $`T_c0.07`$ $`\left(ps\right)^1.`$ The exact ITM theoretical decoherence times are two orders of magnitude and five orders of magnitude smaller than the experimental value even when we consider the DQD charge qubit in independent PCPB and DCPB. These can finally and without accident lead to the conclusion that the phonon decoherence is a subordinate mechanism in the DQD charge qubit. In general, besides the phonon couplings’ decoherence, the qubit can also result in decoherence from electromagnetic fluctuations (with Ohmic noise spectrum), cotunneling effect, background charge fluctuations (with $`1/f`$ noise spectrum), and so on. To find out the dominating mechanism of the DQD charge qubit decoherence and to the best of our abilities to suppress the central decoherence resources are important challenges in the quantum computation field.
###### Acknowledgement 1
The project was supported by National Natural Science Foundation of China (Grant No. 10347133) and Ningbo Youth Foundation (Grant No. 2004A620003).
## VI Figure captions
Fig.1: Real (line) and imaginary (short lines) parts of the response function of the piezoelectric coupling phonon bath (PCPB). Here, we set the temperature $`T=30`$ $`mK`$, and $`\omega _d=0.02`$ $`(ps)^1,`$ $`\omega _l=0.5`$ $`(ps)^1,`$ $`g^{pz}=0.035`$ $`(ps)^2.`$
Fig.2: Real (line) and imaginary (short lines) parts of the response function of the deformation coupling phonon bath (DCPB). Here, we set $`g^{df}=0.029`$ $`(ps)^2`$, and other parameters are same as those in Fig.1.
Fig.3: The evolutions of the off-diagonal elements of the reduced density matrix for the DQD charge qubit in PCPB when $`\omega _l=0.5`$ $`(ps)^1`$ (line) and $`\omega _l=0.7`$ $`(ps)^1`$ (short lines). Here, the cutoff frequency is $`\omega _c=5`$ $`(ps)^1`$, other parameters are same as those in Fig.1. The initial state of the qubit and environment are described in the text.
Fig.4: The evolutions of the off-diagonal elements of the reduced density matrix for the DQD charge qubit in DCPB when $`\omega _l=0.5`$ $`(ps)^1`$ (line) and $`\omega _l=0.7`$ $`(ps)^1`$ (short lines). Here, the parameters are same as those in Figs.2 and 3. |
warning/0506/hep-th0506243.html | ar5iv | text | # 1 The one-cut solution 𝒞,𝜌 with 𝑛=6 for some fixed value of 𝛼=𝐾/𝐿 and the distribution of fluctuation roots 𝜇_𝑘^±. The mode number of 𝜇⁺_𝑘 is 6+𝑘 for 𝑘>𝑘_c=4 for all other 𝜇⁺_𝑘,𝜇⁻_𝑘 it is 6-𝑘. Some of these mode numbers are indicated in the figure. As compared to the vacuum with 𝛼=0, the roots between 𝜇₄⁻ and 𝜇₄⁺, as well as their mode numbers, are distorted strongly by the cut.
## Acknowledgements.
We thank Volodya Kazakov, Charlotte Kristjansen, Matthias Staudacher and Arkady Tseytlin for discussions. N. B. would like to thank NORDITA for hospitality during work on this article. The work of N. B. is supported in part by the U.S. National Science Foundation Grant No. PHY02-43680. Any opinions, findings and conclusions or recommendations expressed in this material are those of the authors and do not necessarily reflect the views of the National Science Foundation. |
warning/0506/astro-ph0506024.html | ar5iv | text | # Massive Star Outflows
## 1 Introduction
The dynamics of outflow and infall associated with young stellar objects (YSOs) affect the turbulent support and dissipation of molecular clouds, the final mass of the central star, and ultimately, the conditions for planet formation. Outflows in the form of wide-angle winds and/or well-collimated jets are associated with YSOs of all luminosities. However, the conditions and timescales associated with massive star formation differ from their low-mass counterparts, e.g.: 1) Massive stars evolve to the Zero Age Main Sequence (ZAMS) more rapidly than low-mass stars: O spectral type stars have a Kelvin-Helmholtz timescale $`\stackrel{<}{}10^4`$ years while solar-type stars require more than $`10^7`$ years to reach the main sequence; 2) OB stars reach the ZAMS while still embedded and perhaps still accreting. Once on the ZAMS, UV photons generate a hyper-compact HII (HC HII) region and strong winds drastically affect the physical conditions, structure, and chemistry of their surrounding; 3) OB stars form in denser clusters than lower mass stars and a higher fraction of OB stars form in binary or multiple systems. Thus, there may be more dynamical interactions between (proto)stars (e.g. Bonnell, Bate, & Vine 2003).
Two scenarios have been proposed to explain the formation of massive O stars. The first scenario relies on non-spherical accretion which over comes radiation pressure even from the most luminous stars (e.g., Norberg & Maeder 2000, Behrend & Maeder 2001, Yorke & Sonnhalter 2002, McKee & Tan 2003). However, the most massive star that can form in the 2D models of Yorke & Sonnhalter is about 30 M – even when 120 M of gas is available in the initial cloud. A solution to this problem was presented by Krumholz et al. (2005). Using a 3D Monte Carlo radiative transfer code (Whitney et al. 2003) modified to account for time-dependent hydrodynamics Krumholz et al. found that the optically thin outflow channels carried away significant radiation, allowing the surrounding envelope to remain relatively cooler. This made it easier for the core to collapse to form very massive stars. The wider the outflow cavity, the more easily radiation can escape. Thus, the presence of a cavity created by a bipolar outflow may be critical for the formation of the most massive stars. Yet, the accretion-based formation of massive stars may differ from low-mass scenarios. In particular, recent simulations suggest that nearly 90% of the mass of massive stars is not due to accretion from an initial clump or envelope but from subsequent competitive accretion of material much farther from the stars (e.g., Bonnell, Vine, & Bate 2004).
The second scenario proposed for the formation of O stars involves the coalescence of lower mass (proto)stars at the centers of dense clusters (e.g., Bonnell, Bate, & Zinnecker 1998, Stahler, Palla & Ho 2000, Bonnell & Bate 2002). The mergers would likely destroy any accretion disks around the lower mass components and disrupt their outflows. The resulting massive star will likely have rotating circumstellar material but accretion is not a necessary criteria for formation. Simulations by Bonnell, Bate, & Vine (2003) suggest that, while true mergers of (proto)stars may be rare outside of binary systems, dynamical interactions between stars in dense star forming clusters are not: nearly 1/3 of all stars and most stars with M$`{}_{}{}^{}>3`$ M suffer disruptive interactions that can truncate circumstellar disks and bring the accretion and outflow process to an abrupt and possibly explosive halt (Bally & Zinnecker 2005). While the competitive accretion model predicts molecular outflows with properties that may be similar to low-mass outflows, it is unlikely that highly-collimated structures could exist in the coalescence scenario and the outflow energetics may be significantly different.
Previous reviews that focus mostly on outflows from low-mass YSOs include Lada (1995), Bachiller (1996), Bachiller & Tafalla (1999), and Reipurth & Bally (2001). Massive YSO outflows are reviewed by Garay & Lizano (1999), Churchwell (1999), Königl (1999), and Beuther & Shepherd (2005). Richer et al. (2000), Cabrit (2002), and Shepherd (2003) examine the similarities between massive and low-mass outflows. Here we examine new results associated with massive flows, discuss a possible evolutionary scenario, and examine evidence for morphological and dynamical differences between massive and low-mass flows.
## 2 Outflow Energetics
Young, low-luminosity YSOs ($`L_{bol}`$ few L) have mass outflow rates $`\dot{M}_f=10^8`$ to $`10^5`$ M yr<sup>-1</sup>, momentum rates $`\dot{P}_f=10^7`$ to $`10^4`$ M km s<sup>-1</sup>yr<sup>-1</sup>, and mechanical luminosity $`L_{mech}=10^3`$ to $`10^1`$ L. Young stars reach the main sequence at about the same time the outflow terminates ($`10^7`$ to $`10^8`$ years). The full opening angle, $`\theta `$, of flows from low-luminosity YSOs tends to be 10–30 close to the central source ($`<`$ 50 AU) but then re-collimates ($`\theta `$ few degrees) within 100 AU from the protostar (e.g. Ray et al. 1996; Dougados et al. 2000). Intermediate-mass YSOs cover a wide range of $`L_{bol}`$. For A-type stars $`L_{bol}`$ 10–30 L. Outflows from A-type YSOs are similar to those from low-mass YSOs although the energetics are roughly an order of magnitude higher. Well-collimated jets are common and the associated molecular flows tend to be collimated for young sources and less collimated for more evolved sources. For B-type stars $`L_{bol}`$ varies 2 orders of magnitude: from 30 L several $`\times 10^4`$ L. Early-B stars ($`L_{bol}10^4`$ L) generate UC HII regions and reach the ZAMS while still accreting and generating strong molecular outflows (e.g. Churchwell 1999, Garay & Lizano 1999, and references therein). Intermediate-mass YSOs have $`\dot{M}_f=10^5`$ to a few $`\times 10^3`$ M yr<sup>-1</sup>, $`\dot{P}_f=10^4`$ to $`10^2`$ M km s<sup>-1</sup>yr<sup>-1</sup>, and $`L_{mech}=10^1`$ to $`10^2`$ L. Outflows from early-B and late O stars can be well-collimated when the dynamical times scale is less than $`10^4`$ years although poorly collimated flows are more common in both young and old sources. O stars with $`L_{bol}>10^4`$ L generate powerful winds with $`\theta 90`$ within 50 AU of the star while accompanying molecular flows can have $`\theta >90`$. The flow momentum rate ($`>10^2`$ M km s<sup>-1</sup>yr<sup>-1</sup>) is more than an order of magnitude higher than what can be produced by stellar winds and $`L_{mech}`$ exceeds $`10^2`$ L (e.g. Churchwell 1999, Garay & Lizano 1999).
For both well-collimated and poorly collimated flows, independent studies have established correlations of the form: $`\dot{M}L_{bol}^{0.6}`$ where $`\dot{M}`$ is the bipolar molecular outflow rate and the ionized mass outflow rate in the wind from $`L_{bol}=0.3`$ to $`10^5`$ L (e.g. Levreault 1988, Cabrit & Bertout 1992, Shepherd & Churchwell 1996, Anglada 1996, Henning et al. 2000, Beuther et al. 2002a, Wu et al. 2004). There is also a strong correlation between bolometric luminosity and circumstellar mass from $`L_{bol}=0.1`$ to $`10^5`$ L (Saraceno et al. 1996, Chandler & Richer 2000). These correlations suggest that there is a strong link between accretion & outflow for a wide range of $`L_{bol}`$ \- even into the mid-O star range.
CO outflows from both low and high mass stars show a mass-velocity relation in the form of a power law $`dM(v)/dvv^\gamma `$ with $`\gamma `$ ranging from $`1`$ to $`8`$. The slope, $`\gamma `$, steepens with age and energy in the flow (e.g. Rodríguez et al. 1982, Lada & Fich 1996, Shepherd et al. 1998, Richer et al. 2000). A similar relation of H<sub>2</sub> flux-velocity also exists with $`\gamma `$ between $`1.8`$ and $`2.6`$ for low and high-mass outflows (Salas & Curz-González 2002).
For a few early B (proto)stars with outflows that have a well-defined jet, the jet appears to have adequate momentum to power the larger scale CO flow although this relation is not as well established as it is for lower luminosity sources. For example, IRAS 20126$`+`$4104 has a momentum rate in the SiO jet of $`2\times 10^1\left(\frac{2\times 10^9}{\mathrm{SiO}/\mathrm{H}_2}\right)`$ M km s<sup>-1</sup>yr<sup>-1</sup> while the CO momentum rate is $`6\times 10^3`$ M km s<sup>-1</sup>yr<sup>-1</sup> (Cesaroni et al. 1999, Shepherd et al. 2000). Although the calculated momentum rate in the SiO jet is adequate to power the CO flow, the uncertainties in the assumed SiO abundance makes this difficult to prove. Another example is IRAS 18151–1208 in which the H<sub>2</sub> jet appears to have adequate momentum to power the observed CO flow (Beuther et al. 2002a, Davis et al. 2004). For sources that have a weak jet component coupled with a wide-angle flow, it has not been determined what fraction of the momentum rate is supplied by the jet. Despite the detection of jets toward massive protostars, the molecular flows themselves tend to be less collimated than their low mass counterparts. Wu et al. (2004) find that the average collimation factor (major/minor radius) for outflows from sources with $`L_{bol}>10^3`$L is 2.05 compared with 2.81 for flows from lower luminosity sources. This is true even for sources in which the angular size of the flow is at least five times the resolution. Thus, the generally poorer collimation for massive flows is not due to low resolution of the observations.
There are an increasing number sources for which our observations are of comparable quality to those of low-mass YSOs and they provide a quantitative view of the range of the outflow/infall properties for these energetic sources. Wu et al. (2004) provide a statistical overview of the general properties of both massive and low-mass outflows. In the following sections, selected outflows from early B and O (proto)stars are discussed and compared.
## 3 Well-Collimated Outflows from Early B (Proto)Stars
Collimated, ionized jets can be generated by early-B protostars. The youngest sources ($`10^4`$ years or less) can be jet-dominated and can have either well-collimated or poorly collimated molecular flows. In a few sources, jets tend to have opening angles, $`\theta `$, between 25 and 30 but they do not re-collimate. Other sources appear to generate well-collimated jets ($`\theta `$ few degrees) that look like scaled up jets from low-luminosity protostars. Jet activity can continue as long as $`10^6`$ years although associated molecular flows have large opening angles and complex morphology. Below are some examples of sources with outflows that are at least partially driven by a jet.
One of the best examples of an early B ZAMS star with a jet is HH 80-81 which has a highly collimated ionized jet with a projected length $`5`$ pc and age $`10^6`$ years (Martí, Rodríguez, & Reipurth 1993; Heathcote et al. 1998). The truncated CO flow full opening angle is roughly 40 and does not re-collimate (Yamashita et al. 1989). The molecular flow momentum rate ($`\dot{P}_f=6\times 10^3`$ M km s<sup>-1</sup>yr<sup>-1</sup>) is an order of magnitude greater than $`\dot{P}_j`$ in the ionized jet. The CO flow position angle is misaligned with the jet by roughly 30. The jet itself has only a slight wobble of a few degrees and thus, jet precession is not likely to cause the wide opening in the molecular flow. The luminosity of the driving source is uncertain. The total luminosity of the system is $`2\times 10^4`$ L however there is at least one early B star near the jet center that could account for a third or more of the luminosity (Stecklum et al. 1997). Thus, HH 80-81 appears to be powered by a B3-B1 ZAMS star. A similar situation is seen toward Ceph A HW2 - an early B star with an ionized jet, complex HH objects and multiple molecular outflows (e.g. Sargent 1979, Hartigan et al. 1986, Torrelles et al. 1993, Rodríguez et al. 1994, and Garay et al. 1995).
While HH 80-81 looks like a scaled T-Tauri star with a jet that stays well-collimated far from the central source, a growing number of sources are being discovered which have jets with wider opening angles that do not appear to re-collimate. In particular, IRAS 20126$`+`$4104 (B0.5 spectral type, age $``$ few $`\times 10^4`$ years) has an ionized and molecular jet which are misaligned from the larger scale molecular outflow by more than 60 (Cesaroni et al. 1999, Cesaroni et al. 2005, Hofner et al. 1999, Shepherd et al. 2000). The large precession angle in the IRAS 20126 jet appears to be due to an interaction with a close binary companion. The full opening angle of the $``$0.1 pc SiO jet is roughly 40. This agrees with recent estimates based on water maser studies that suggest the jet full opening angle is $`34`$ roughly 200 AU from the central source (Moscadelli et al. 2005).
Another relatively wide-opening angle, ionized jet is IRAS 16547$``$4247 (Brooks et al. 2003, Garay et al. 2003, Rodríguez et al. 2005). The luminosity of the source suggests a B0-O8 spectral type. The molecular mass outflow rate is unknown but the radio luminosity of the ionized jet is consistent with being powered by a late O star. Based on the extent of the H<sub>2</sub> emission, the flow appears to be $`7,000`$ years old. Rodríguez et al. estimate the full opening angle of the jet to be $`25`$ out to about 10,000 AU ($`3^{\prime \prime }`$ at D = 2.9 kpc). Again, this is a wide-opening angle jet relative to the typical 1-3 opening angles often seen toward low-mass jets. Two other examples of early-B (proto)stars with well-collimated flows are: 1) IRAS 05358$`+`$3543 which is an early B star cluster (age $`4\times 10^4`$ years) that has multiple collimated molecular outflows (Beuther et al. 2002b, Sridharan et al. 2002); and 2) IRAS 18151$``$1208 which harbors two H<sub>2</sub> jets, one of which appears to be from an early B star (Davis et al. 2004).
## 4 Wide-Angle Outflows from Early B (Proto)Stars
Molecular outflows from early B (proto)stars are on average less collimated than low-mass flows (Wu et al. 2004). Low collimation factors can occur in molecular flows even when there is a well-defined ionized jet (see previous section). In at least some sources, both the ionized wind near the central source and the larger scale molecular flow are poorly collimated and there is no evidence for a jet. Several examples are given below.
The early-B YSO, G192.16–3.82, has a poorly collimated ionized wind ($`\theta 40`$) within 50 AU of the YSO that expands to $`\theta 90`$ 0.1 pc from the source (Shepherd et al. 1998, Shepherd & Kurtz 1999, Shepherd, Claussen, & Kurtz 2001). The 100 M molecular flow is a few $`\times 10^5`$ years old and forms the truncated base of the larger scale flow that extends more than 10 pc from end-to-end (Devine et al. 1999). The outflow is consistent with being produced by a wind blown bubble and there is no evidence for a collimated jet. A source that appears to be similar to G192.16 is AFGL 490. It has a wide angle wind and outflow from an early B star (e.g. Snell et al. 1984, Mitchell et al. 1995, Schreyer et al. 2005). The age of the CO flow is estimated to be $`2\times 10^4`$ years old, the outflow opening angle is roughly 50 and the collimation factor is about 1.0. The central star, detected in the near-IR (e.g. Davis et al. 1998) is surrounded by a 20,000 AU CS torus that appears to have an inner accretion disk that is less than 500 AU in diameter (Schreyer et al. 2002). There is no obvious sign of a collimated jet from the central object.
As a final example consider W75 N. At the center of the CO outflows is a cluster of UC HII regions. One, VLA 2, appears to power a wide-angle, red-shifted CO flow to the south-west (Davis et al. 1998, Shepherd 2003, Shepherd et al. 2004). Proper motions of water masers associated with the outflow from VLA 2 show that the opening angle of the flow appears to be nearly 180 within 100 AU of the star (Torrelles et al. 2003). The larger scale CO flow opening angle is roughly 50, consistent with a wind-blown bubble. There is no indication near the star that a well-collimated jet exists. The infrared line emission suggests that only slow, non-dissociative J-type shocks exist throughout the parsec-scale outflow. Fast, dissociative shocks, common in jet-driven outflows from low-mass stars, are absent in W75 N. Interestingly, VLA 1, just 2,000 AU from VLA 2, is a jet that drives a collimated CO flow. Torrelles et al. (2004) suggest that the increased collimation may be because VLA 1 is at a different evolutionary state than VLA 2.
## 5 Mid to Early O Star Outflows
To date, extremely collimated outflows have not been observed toward sources earlier than B0. It is possible that this is simply a selection effect because O stars form in dense clusters and reach the ZAMS in only a few $`\times 10^4`$ years. Thus, any collimated outflows may be confused by other flows. For a few sources that are relatively isolated and for which adequate high-resolution data are available, no inner accretion disk ($`50`$ AU diameter) has been conclusively detected. In a few sources, collimation close to the protostar appears to be due to pressure confinement from an equatorial torus of dense gas (e.g. K3-50A: Depree et al. 1994, Howard et al. 1997). This may also be the case for the Orion I outflow (e.g. McCaughrean & Mac Low 1997, Greenhill et al. 1998) however, there is now strong evidence that Source I and the BN object are moving away from each other and that they were within 225 AU of each other about 500 years ago (Rodríguez et al. 2005). Source I and BN may have been members of a multiple system or BN may have been a companion of $`\theta ^1`$ Ori C and made a close approach to Source I after it was ejected (Tan 2003). The poorly collimated, explosive-looking outflow seems to have occurred $`1,000`$ years ago. Could it have been caused by the event that ejected the BN object? If so, then the outflow is not actively powered by Source I and the geometry of the circumstellar material around Source I need not be linked to the large-scale flow. Such explosive events may not be uncommon as evidenced by the ’fingers’ of shocked gas seen in Spitzer images of the ionized outflow associated with G34.26$`+`$0.15 (Churchwell, personal communication).
As a final example, consider G5.89–0.39. The UC HII region appears to be powered by an O5 star. The star has a small excess at 3.5$`\mu `$m suggesting the presence of circumstellar material (Feldt et al. 2003). The O5 star is along the axis of two H<sub>2</sub> knots that appear to trace a collimated N-S flow (Puga et al., poster presented at this conference). There is also a N-S C<sup>34</sup>S(J=3-2) and the UC HII region is expanding in the N-S direction (Cesaroni et al. 1991, Acord et al. 1998). Although still circumstantial, the evidence is mounting that the O5 star in G5.89 produced the N-S outflow and thus is forming via accretion. Note, there is also an SiO flow (with a NE-SW orientation) that is powered by a star to the S-W of the O5 star (Sollins et al. 2004, Puga et al. in preparation) and a larger scale E-W CO flow for which the driving source has not been identified (Watson et al. 2002).
Additional evidence that O stars form via linked accretion and outflow comes from 7 mm continuum observations of very young O stars (van der Tak & Menten 2005). Models of the derived sizes, flux densities, and radio spectra of sources with luminosities up to $`10^5`$L suggest that at least one star has a dust disk and all are accreting material. Although not conclusive, this supports the accretion-based formation scenario for even mid-O stars.
## 6 Evidence for an Evolutionary Scenario
Current observations indicate that the outflow/infall mechanism is similar from T Tauri stars up to early B protostars. Although there is evidence that the energetics for at least some early-B stars may differ from their low-mass counterparts, the dynamics are still governed by the presence of linked accretion and outflow. A few young O stars show evidence for accretion as well. Although there are clearly explosive events, these may be due to close encounters that disrupt the accretion process. However, dynamical mergers of (proto)stars to create the most massive stars may still be a viable formation mechanism.
There are both well-collimated and poorly collimated molecular flows from massive stars. Well-collimated molecular flows tend to be in systems with ages less than a few times $`10^4`$ years old where the central object has not yet reached the main sequence. Hence the effects of increased irradiation on the disk and disk-wind due to the stellar radiation field are minimized. Observed jets often have opening angles between 25 and 30 with little evidence for recollimation of the jets on larger scales. This could be due to a change in the balance between magnetic and plasma pressure. Poorly collimated flows (opening angle greater than 50 that show no evidence for a more collimated component) are associated with more evolved sources that have detectable UC HII regions and the central star has reached the main sequence. Thus, the disk and outflow are subject to significantly more ionizing radiation.
To account for the differences seen in flow morphologies from early B to late O stars Beuther & Shepherd (2005) proposed two possible evolutionary sequences which could result in similar observable outflow signatures as shown in Figure 1. For both early B and O stars, the flows begin collimated and become less-collimated as the star reaches the ZAMS and generates significantly more Lyman continuum photons. This evolutionary sequence appears to qualitatively fit the observations however it must be tested against both theory and observations.
### 6.1 Impact of Luminosity on Outflows
Once a massive OB star reaches the main sequence, the increased radiation from the central star generates significant Lyman continuum photons and will likely ionize the outflowing gas even at large radii. The result may be an increase in the plasma pressure at the base of the flow which could overwhelm any collimating effects of a magnetic field (see, e.g., Shepherd 2003, Königl 1999). It may also affect the inherent collimation of the ionized wind from the stellar surface as suggested in the simulations by Yorke & Sonnhalter (2002). Recollimation of the outflowing gas from the disk is expected under ideal magneto-hydrodynamic (MHD) conditions where the recombination timescale is much greater than the dynamical timescale of the outflow. Thus, the ionized plasma becomes ’frozen-in’ to the magnetic field lines. Models of the ionization & density along beams of HH jets indicate that ideal MHD is valid for partially ionized jets from T Tauri stars (e.g. Bacciotti & Eislöffel 1999). However, the ideal MHD assumption breaks down as: 1) the plasma temperature and density increase; 2) the turbulence in the disk or wind increases; or 3) the toroidal component of the magnetic field, $`B_\varphi `$, decreases. As $`L_{bol}`$ increases plasma temperature and density will increase and the disk may be more turbulent. Thus, ideal MHD assumptions may break down for high-mass outflows and one may expect that energetics and re-collimation could be affected. To be more specific, conditions appropriate for luminous YSO must be included in simulations (e.g. Königl 1999). It is interesting to note that several early B protostars with strong jet components in the outflows have $`30`$ opening angle jets that do not appear to re-collimate. Perhaps this is due to a break-down in ideal MHD conditions along the outflow cone?
### 6.2 Impact of Mass-Loading on Outflows
Two conditions expected for luminous YSOs are increased disk surface heating (due to increased stellar radiation and shocks) and a higher level of disk turbulence when $`M_D/M_{}>0.3`$ where $`M_D`$ is the disk mass and $`M_{}`$ is the (proto)stellar mass. Increased disk surface heating causes increased mass loading (e.g. Ferriera 2002) and since winds carry angular momentum away, this means that a warmer, denser wind causes the disk to rotate slower. Pudritz (these proceedings) has evaluated the impact of mass loading on magnetized disk-winds under conditions appropriate for massive protostars that have not yet reached the ZAMS (e.g. no significant heating from the central protostar). Solutions show that a strong wide-angle wind develops as the mass-loading increases. Thus, the outflow dynamics and geometry appears to change as the outflow rate increases.
### 6.3 Impact of Disk Turbulence on Outflows
Disks around early-B YSOs are significantly more massive than those around A-type and T Tauri YSOs. Further, $`M_D/M_{}`$ is often greater than 0.3 which may cause the disk to be locally unstable and provide an additional means of angular momentum transport from the protostar (e.g. Laughlin & Bodenheimer 1994; Yorke, Bodenheimer, & Laughlin 1995; Shepherd, Claussen, & Kurtz 2001). There are a growing number of early-B (proto)stars that are actively powering molecular outflows and have circumstellar disks or molecular tori with masses ranging from 0.1 M to tens of solar masses. The increased $`M_D/M_{}`$ suggests that the detailed dynamics of infall and outflow in early-B YSOs may differ from their lower mass counterparts. Recently, Lodato & Rice (2005) have made simulations of circumstellar disks in which $`M_{disk}`$ = 0.1, 0.5, and 1.0 M to evaluate the behavior of spiral density waves. They find that the mass of the disk dramatically changes the characteristics of the spiral density waves that can transport angular momentum efficiently. As the disk mass increases, the waves become stronger and in some cases episodic. The associated accretion onto the central star increases during periods when spiral density waves are strong and carry away significant angular momentum. It is not clear how or even if such behavior will cause significant changes in the outflow characteristics but one might expect it would since accretion and outflow are so inherently linked. The simulations of Lodato & Rice are in a regime in which fragmentation of the disk is not allowed, further, only disk heating and cooling are considered (no stellar heating or heating due to the outflow or shocks above the disk plane). Additional improvements to the simulations are planned which will incorporate more realistic heating and cooling functions and disk heating from the star and outflow. Indeed, simulations of disks around low-mass stars which include realistic heating from the central star and outflow suggest that irradiation can quench fragmentation due to local gravitational instabilities because the disk temperature is raised above the parent cloud temperature (Matzner & Levin 2005).
The simulations of Lodato & Rice (2005) suggest an increase in $`\dot{M}_{acc}`$ if spiral density waves transport angular momentum to the outer parts of the disk - one would not expect a corresponding increase in the disk wind mass loss rate, $`\dot{M}_w`$, since there is not excess angular momentum in the inner disk where the wind is generated. Thus, significant angular momentum transport associated with spiral density waves in the disk may contribute to a decrease in $`\dot{M}_w/\dot{M}_{acc}`$ as well. Is there observational evidence for a decrease in $`\dot{M}_w/\dot{M}_{acc}`$? Richer et al. (2000) demonstrate that one can measure the Outflow force/Accretion Force vs $`L_{bol}`$ and that this quantity can be related to $`\dot{M}_w/\dot{M}_{acc}`$. Shepherd (2003) shows that there is marginal evidence for such a decrease but the errors in the data are large and any trend is at best, marginal. Further observations are necessary to determine if such a trend actually exists.
## 7 Summary
Mid- to early-B protostars and late O protostars have accretion disks and outflows that can be well-collimated. Ionized jets in very young sources (age $`\stackrel{<}{}`$$`10^4`$ years) can have opening angles $`30`$ with no recollimation detected or they can re-collimate similar to jets from low-luminosity protostars. Once an HII region forms, the associated ionized outflows have strong wide-angle winds and a collimated jet is often not detected. The non-detection of a jet in older sources is unlike low-mass outflows were there is evidence for a two-component wind (jet$`+`$wide-angle) even in older sources e.g. Arce & Goodman 2002). Mid to early O stars have poorly collimated outflows and explosive events suggest that protostellar or protostar/disk interactions are common. Current observations suggests that outflows from early OB protostars begin collimated and then become less collimated with time once the star reaches the main sequence. This evolution in outflow dynamics and morphology could be due to, e.g.: increased ionizing radiation from the star as it begins to burn hydrogen; mass-loading onto the disk wind; and/or disk instabilities which can be an efficient angular momentum transport mechanism in massive disks. |
warning/0506/math0506087.html | ar5iv | text | # Closures of Steinberg fibers in twisted wonderful compactifications
## 1. Introduction
Let $`G`$ be a simple linear algebraic group over an algebraically closed field $`k`$. Let $`B`$ be a Borel subgroup of $`G`$ and let $`TB`$ denote a maximal torus. Let $`W`$ denote the associated Weyl group and $`I`$ denote the associated set of simple roots. For a subset $`J`$ of $`I`$ we let $`W_J`$ denote the subgroup of $`W`$ generated by the simple roots in $`J`$.
The wonderful compactification $`X`$ of $`G`$ (see e.g. \[DP\], \[Str\]), is a smooth projective $`(G\times G)`$-variety containing $`G`$ as an open subset. The $`G\times G`$-orbits in $`X`$ are indexed by the subsets $`J`$ of $`I`$, and we fix certain base points $`h_J`$ for these orbits. Let $`\sigma `$ denote a diagram automorphism of $`G`$ and let $`X_\sigma `$ be the associated twisted wonderful compactification of $`G`$, i.e. as a variety $`X_\sigma `$ is just $`X`$ but the $`G\times G`$-action is twisted by $`\sigma `$ on the second coordinate. Let $`h_{J,\sigma }`$ denote the point in $`X_\sigma `$ identified with $`h_{\sigma (J)}`$ in $`X`$. Then the collection $`h_{\sigma (J)}`$, $`JI`$, are representatives for the $`G\times G`$-orbits in $`X_\sigma `$. A $`G`$-stable piece in $`X_\sigma `$ is then a locally closed and smooth subvariety in $`X`$ of the form $`Z_{J,\sigma }^w=\mathrm{diag}(G)(Bw,1)h_{J,\sigma }`$, where $`wW^{\sigma (J)}`$ is a minimal length coset representative of $`W/W_{\sigma (J)}`$ and $`\text{diag}(G)`$ denotes the diagonal in $`G\times G`$. We then have a decomposition $`X_\sigma =_{JI,wW^J}Z_J^w`$ (see \[L2, 12.3\] and \[H2, 1.12\]).
The $`G`$-stable pieces were first introduced by Lusztig to study the $`G`$-orbits and parabolic character sheaves. However, his original definition was based on some inductive method. The (equivalent) definition that we used above is due to the first author in \[H1\]. What we need in this paper is that the dimension of $`Z_{J,\sigma }^w`$ is equal to $`dim(G)l(w)|IJ|`$, where $`l(w)`$ is the length of $`w`$ and $`|IJ|`$ is the cardinality of the set $`IJ`$ (see \[L2, 8.20\]). More properties about the $`G`$-stable pieces can be found in \[L2\] and \[H2\]. The $`G`$-stable pieces were also used by Evens and Lu in \[EL\] to study the Poisson structure and symplectic leaves.
Consider $`G`$ as a $`G\times G`$-variety by left and right translation and define $`G_\sigma `$, similar to the definition of $`X_\sigma `$, by twisting the $`G`$-structure of $`G`$ on the second factor by $`\sigma `$. A $`\sigma `$-conjugacy class in $`G_\sigma `$ is then a $`\text{diag}(G)`$-orbit in $`G_\sigma `$. The set of elements in $`G_\sigma `$ whose semisimple part lies in a fixed $`\sigma `$-conjugacy class is then called a Steinberg fiber of $`G_\sigma `$. In this paper we study the closure of Steinberg fibers within $`X_\sigma `$.
In \[L2\], Lusztig gave an explicit description for the closure of the unipotent variety in the group compactification when $`G=PGL_2`$ or $`PGL_3`$. In \[Sp2\], Springer studied the closure of an arbitrary Steinberg fiber for any connected, simple algebraic group and obtained some partial results. Based on their results, the first author obtained an explicit description of the closure of Steinberg fibers in the non-twisted case. The result in \[H1\] was formulated using $`G`$-stable pieces and the proof was based on a case-by-case checking. The main purpose of this paper is to generalize the result of \[H1\] to the twisted case with a more conceptual (and easier) proof. More precisely, we prove,
###### Theorem.
Let $`F`$ be a Steinberg fiber of $`G_\sigma `$ and $`\overline{F}`$ be its closure in $`X_\sigma `$. Then
$$\overline{F}F=\underset{JI}{}\underset{\begin{array}{c}wW^{\sigma (J)}\\ \text{supp}_\sigma (w)=I\end{array}}{}Z_{J,\sigma }^w,$$
where $`\text{supp}_\sigma (w)`$ denotes the minimal $`\sigma `$-stable subset of $`I`$ such that $`w`$ is contained in $`W_{\text{supp}_\sigma (w)}`$.
As a consequence, the boundary of the closure is independent of the choice of the Steinberg fiber. Likewise it may be shown that the boundary of the closure of $`F`$ within any equivariant embedding of $`G`$ is independent of the choice of $`F`$ (see \[T\]). As a by-product, we will also give an explicit description of the “nilpotent cones” on $`X`$.
## 2. Wonderful compactifications and $`G`$-stable pieces
### 2.1
Let $`G`$ denote a simple linear algebraic group over an algebraically closed field $`k`$. We consider $`G`$ as a $`G\times G`$-variety by left and right translation. Let $`B`$ be a Borel subgroup of $`G`$, $`B^{}`$ be an opposite Borel subgroup and $`T=BB^{}`$. The unipotent radical of $`B`$ (resp. $`B^{}`$) will be denoted by $`U`$ (resp. $`U^{}`$). Let $`R`$ denote the set of roots defined by $`T`$ and let $`R^+`$ denote the set of positive roots defined by $`B`$. Let $`(\alpha _i)_{iI}`$ be the set of simple roots. For $`iI`$, we denote by $`\omega _i`$ and $`s_i`$ the fundamental weight and the simple reflection corresponding to $`\alpha _i`$.
We denote by $`W`$ the Weyl group associated to $`T`$. For any subset $`J`$ of $`I`$, let $`W_J`$ be the subgroup of $`W`$ generated by $`\{s_jjJ\}`$ and $`W^J`$ be the set of minimal length coset representatives of $`W/W_J`$.
For $`JI`$, let $`P_JB`$ be the standard parabolic subgroup defined by $`J`$ and let $`P_J^{}B^{}`$ be the parabolic subgroup opposite to $`P_J`$. Set $`L_J=P_JP_J^{}`$. Then $`L_J`$ is a Levi subgroup of $`P_J`$ and $`P_J^{}`$. The semisimple quotient of $`L_J`$ of adjoint type will be denoted by $`G_J`$. We denote by $`\pi _{P_J}`$ (resp. $`\pi _{P_J^{}}`$) the projection of $`P_J`$ (resp. $`P_J^{}`$) onto $`G_J`$.
### 2.2
Let $`X`$ denote the wonderful compactification of $`G`$ (\[DP\], \[Str\]). Then $`X`$ is an irreducible, smooth projective $`(G\times G)`$-variety with finitely many $`G\times G`$-orbits $`Z_J`$ indexed by the subsets $`J`$ of $`I`$. As a $`(G\times G)`$-variety the orbit $`Z_J`$ is uniquely isomorphic to the product $`(G\times G)\times _{P_J^{}\times P_J}G_J`$, where $`P_J^{}\times P_J`$ acts on $`G\times G`$ by $`(q,p)(g_1,g_2)=(g_1q^1,g_2p^1)`$ and on $`G_J`$ by $`(q,p)z=\pi _{P_J^{}}(q)z\pi _{P_J}(p)^1`$. Let $`h_J`$ be the image of $`(1,1,1)`$ in $`Z_J`$ under this isomorphism.
We denote by $`\mathrm{diag}(G)`$ the image of the diagonal embedding of $`G`$ in $`G\times G`$. For $`JI`$ and $`wW^J`$, set $`Z_J^w=\mathrm{diag}(G)(Bw,1)h_J`$. Then $`Z_J^w`$ is a locally closed subvariety of $`X`$ and (see \[L2, 12.3\] and \[H2, 1.12\])
$$X=\underset{\begin{array}{c}JI\\ wW^J\end{array}}{}Z_J^w.$$
We call $`Z_J^w`$ a $`G`$-stable piece.
## 3. Twisted actions
### 3.1
An automorphism $`\sigma `$ of $`G`$ which stabilizes the Borel subgroup $`B`$ and the maximal torus $`T`$ will induce a permutation of $`I`$. When the order of $`\sigma `$ as an automorphism of $`G`$ coincides with the order of the associated permutation of $`I`$, we say that $`\sigma `$ is a diagram automorphism. From now on $`\sigma `$ will denote a diagram automorphism of $`G`$. We also denote by $`\sigma `$ the corresponding bijection on $`I`$ and $`W`$.
Let $`G_\sigma `$ be the $`(G\times G)`$-variety which as a variety is isomorphic to $`G`$ and where the $`G\times G`$ action is twisted by the morphism $`G\times GG\times G`$, $`(g,h)(g,\sigma (h))`$ for $`g,hG`$. Then we define the wonderful compactification $`X_\sigma `$ of $`G_\sigma `$ to be the $`G\times G`$-variety which as a variety is isomorphic to the wonderful compactification $`X`$ of $`G`$ and where the $`G\times G`$ action is twisted in the same way as above. Notice that we may regard $`G_\sigma `$ as a connected component of the semidirect product $`G<\sigma >`$. In this case, $`X_\sigma `$ is the completion of $`G_\sigma `$ considered in \[L2, 12\].
The $`G\times G`$-orbits in $`X_\sigma `$ coincide with the associated orbits in $`X`$ and we let $`Z_{J,\sigma }`$ denote the orbit coinciding with $`Z_{\sigma (J)}`$. Accordingly we let $`h_{J,\sigma }`$ denote the point in $`Z_{J,\sigma }`$ identified with the base point $`h_{\sigma (J)}`$ of $`Z_{\sigma (J)}`$. For $`JI`$ and $`wW^{\sigma (J)}`$, set $`Z_{J,\sigma }^w=\mathrm{diag}(G)(Bw,1)h_{J,\sigma }`$. Then
$$X_\sigma =\underset{JI}{}\underset{wW^{\sigma (J)}}{}Z_{J,\sigma }^w.$$
We call $`(Z_{J,\sigma }^w)_{JI,wW^{\sigma (J)}}`$ the $`G`$-stable pieces of $`X_\sigma `$ (see \[L2, 12.3\] and \[H2, 1.12\]).
### 3.2
The orbits of $`\mathrm{diag}(G)`$ on $`G_\sigma `$ are called $`\sigma `$-conjugacy classes. Let $`G//_\sigma G`$ be the affine variety whose algebra is the subalgebra $`k[G]^{G,\sigma }`$ of functions in $`k[G]`$ invariant under $`\sigma `$-conjugacy. The inclusion $`k[G]^{G,\sigma }k[G]`$ induces a morphism $`\mathrm{St}:G_\sigma G//_\sigma G`$. If $`\sigma `$ is trivial, then $`\mathrm{St}`$ is just the Steinberg morphism of $`G`$. Thus for arbitrary $`\sigma `$, we call $`\mathrm{St}`$ the Steinberg morphism of $`G_\sigma `$ and the fibers the Steinberg fibers of $`G_\sigma `$.
An element $`gG_\sigma `$ is $`\sigma `$-conjugate to an element in $`B`$ \[Ste2, Lem.7.3\]. Write $`b=tu`$ where $`tT`$ and $`u`$ is an element of the unipotent radical $`U`$ of $`B`$. It is then easily seen that there exists an element $`t_1T`$, such that $`t_1t\sigma (t_1)^1T^\sigma `$. Hence, $`g`$ is $`\sigma `$-conjugate to some element in $`T^\sigma U`$, i.e. we may assume that $`tT^\sigma `$. Notice, that $`t`$ is contained in the closure of the $`\sigma `$-conjugacy class of $`tu`$ and thus, by geometric invariant theory, we find $`\mathrm{St}(tu)=\mathrm{St}(t)`$. Moreover, considering $`t\sigma `$ as an element of the semisimple group $`G<\sigma >`$ it follows that $`t\sigma `$ is quasi-semisimple in the sense of \[Ste2, Sect.9\], i.e. the automorphism of $`G`$ obtained by conjugation by $`t\sigma `$ will fix a Borel subgroup and a maximal torus thereof. As a consequence, the $`\sigma `$-conjugacy class of $`t`$ in $`G_\sigma `$ is closed \[Spa, II.1.15(f)\]. We conclude that any Steinberg fiber of $`G_\sigma `$ is of the form $`_{gG}g(tU)\sigma (g)^1`$ for some $`tT^\sigma `$. In particular, any Steinberg fiber is irreducible.
### 3.3
Let $`G_{\mathrm{sc}}`$ be the connected, simply connected, linear algebraic group associated to $`G`$, and let $`B_{\mathrm{sc}}`$ (resp. $`T_{\mathrm{sc}}`$) denote the Borel subgroup (resp. maximal torus) of $`G_{\mathrm{sc}}`$ associated to $`B`$ (resp. $`T`$). By \[Ste2, 9.16\] the automorphism $`\sigma `$ of $`G`$ may be lifted to an automorphism of $`G_{\mathrm{sc}}`$, which we also denote by $`\sigma `$. We then define the $`G_{\mathrm{sc}}\times G_{\mathrm{sc}}`$-variety $`G_{\mathrm{sc},\sigma }`$ similar to the definition of $`G_\sigma `$. We may also form the quotient $`G_{\mathrm{sc}}//_\sigma G_{\mathrm{sc}}`$ and define Steinberg fibers in $`G_{\mathrm{sc},\sigma }`$ similar to the considerations in 3.2 for $`G_\sigma `$.
The automorphism $`\sigma `$ of $`G_{\mathrm{sc}}`$ induces a natural action of $`\sigma `$ on the set $`\mathrm{\Lambda }`$ of $`T_{\mathrm{sc}}`$-characters, and we let $`\mathrm{\Lambda }_+^\sigma `$ denote the set of $`\sigma `$-invariant dominant weights. Let $`𝒞_1,𝒞_2,\mathrm{},𝒞_l`$ denote the $`\sigma `$-orbits in $`I`$, and set $`\omega _{𝒞_j}=_{i𝒞_j}\omega _i`$. Then the elements $`\omega _{𝒞_j}`$, $`j=1,\mathrm{},l`$, is a generating set of the semigroup $`\mathrm{\Lambda }_+^\sigma `$.
### 3.4
To any dominant weight $`\lambda =_{iI}a_i\omega _i`$ we associate the subset $`I(\lambda )=\{iIa_i0\}`$ of $`I`$. For $`wW`$, let $`\text{supp}(w)I`$ be the set of simple roots whose associated simple reflections occur in some (or equivalently, any) reduced decomposition of $`w`$ (see \[Bou, Prop.IV.1.7\]) and let $`\text{supp}_\sigma (w)=_{k0}\sigma ^k(\text{supp}(w))`$. We have the following characterization of $`\text{supp}(w)`$.
###### Lemma 3.5.
Let $`wW`$ and $`iI`$. Then $`w\omega _i\omega _i`$ if and only if $`isupp(w)`$. Hence for a dominant weight $`\lambda `$, $`w\lambda \lambda `$ if and only if $`I(\lambda )\text{supp}(w)\mathrm{}`$.
Proof. If $`i\text{supp}(w)`$, then $`w\omega _i=\omega _i`$. Now we fix a reduced expression $`w=s_{i_1}\mathrm{}s_{i_n}`$. Assume that $`i\mathrm{supp}(w)`$. We show that $`w\omega _i\omega _i`$ by induction on $`n`$. If $`i_ni`$, then we are done by induction in $`n`$. Hence, we may assume that $`i_n=i`$. But then $`w\alpha _i`$ is a negative root. Thus $`1=\omega _i,\alpha _i^{}=w\omega _i,(w\alpha _i)^{}`$ and, in particular, we cannot have $`w\omega _i=\omega _i`$. ∎
### 3.6
For any dominant weight $`\lambda \mathrm{\Lambda }_+`$ let $`\mathrm{H}(\lambda )`$ denote the dual Weyl module for $`G_{\mathrm{sc}}`$ with lowest weight $`\lambda `$. We then define $`{}_{}{}^{\sigma }\mathrm{H}(\lambda )`$ to be the $`G_{\mathrm{sc}}`$-module which as a vector space is $`\mathrm{H}(\lambda )`$ and with $`G_{\mathrm{sc}}`$-action twisted by the automorphism $`\sigma `$ of $`G_{\mathrm{sc}}`$. Notice that up to a nonzero constant there exists a unique $`G_{\mathrm{sc}}`$-isomorphism $`{}_{}{}^{\sigma }\mathrm{H}(\lambda )\mathrm{H}(\sigma (\lambda ))`$. In particular, when $`\lambda \mathrm{\Lambda }_+^\sigma `$ is $`\sigma `$-invariant there exists a $`G_{\mathrm{sc}}`$-equivariant isomorphism $`f_\lambda :\mathrm{H}(\lambda ){}_{}{}^{\sigma }\mathrm{H}(\lambda )`$. Fix $`f_\lambda `$ such that the its restriction to the lowest weight space $`k_\lambda `$ in $`\mathrm{H}(\lambda )`$ is the identity map (here we use the identification of $`{}_{}{}^{\sigma }\mathrm{H}(\lambda )`$ with $`\mathrm{H}(\lambda )`$ as vector spaces).
## 4. The “nilpotent cone” of $`X`$
### 4.1
For any dominant weight $`\lambda `$ there exists (see \[DS, 3.9\]) a $`G\times G`$-equivariant morphism
$$\rho _\lambda :X\left(\text{End}(\mathrm{H}(\lambda ))\right)$$
which extends the morphism $`G\left(\text{End}(\mathrm{H}(\lambda ))\right)`$ defined by $`gg[\mathrm{Id}_\lambda ]`$, where $`[\mathrm{Id}_\lambda ]`$ denotes the class representing the identity map on $`\mathrm{H}(\lambda )`$ and $`g`$ acts by the left action. By the definition of $`X_\sigma `$ we obtain a $`G\times G`$-equivariant morphism
$$X_\sigma \left(\text{Hom}_k(^\sigma \mathrm{H}(\lambda ),\mathrm{H}(\lambda ))\right).$$
When $`\lambda \mathrm{\Lambda }_+^\sigma `$ we may apply $`f_\lambda `$ to obtain an induced map
$$\rho _{\lambda ,\sigma }:X_\sigma \left(\text{End}(\mathrm{H}(\lambda ))\right)$$
which is $`G\times G`$-equivariant.
### 4.2
An element in $`\left(\text{End}(\mathrm{H}(\lambda ))\right)`$ is said to be nilpotent if it may be represented by a nilpotent endomorphism of $`\mathrm{H}(\lambda )`$. For $`\lambda \mathrm{\Lambda }_+^\sigma `$ we let
$$𝒩(\lambda )_\sigma =\{zX_\sigma \rho _{\lambda ,\sigma }(z)\text{ is nilpotent}\},$$
and call $`𝒩(\lambda )_\sigma `$ the nilpotent cone of $`X_\sigma `$ associated to the dominant weight $`\lambda `$. In 4.4, we will give an explicit description of $`𝒩(\lambda )_\sigma `$.
### 4.3
Define $`\mathrm{ht}`$ to be the height map on the root lattice, i.e., the linear map on the root lattice which maps all the simple roots to $`1`$.
Now fix $`\lambda \mathrm{\Lambda }_+`$. Choose a basis $`v_1,\mathrm{},v_m`$ for $`\mathrm{H}(\lambda )`$ consisting of $`T`$-eigenvectors with eigenvalues $`\lambda _1,\mathrm{},\lambda _m`$ satisfying $`\mathrm{ht}(\lambda _j+\lambda )\mathrm{ht}(\lambda _i+\lambda )`$ whenever $`ji`$. Then $`B`$ is upper triangular with respect to this basis.
Let $`A_J`$ be a representative of $`\rho _\lambda (h_J)`$ in $`\text{End}(\mathrm{H}(\lambda ))`$. Then when $`\lambda _j+\lambda `$ is a linear combination of the simple roots in $`J`$ we have that $`A_Jv_jk^\times v_j`$. If $`\lambda _j+\lambda `$ is not a linear combination of the simple roots in J then $`A_Jv_j=0`$. Assuming that $`\lambda `$ is $`\sigma `$-invariant we obtain, by the definitions in 4.1, a similar description for a representative $`A_{J,\sigma }`$ of $`\rho _{\lambda ,\sigma }(h_{J,\sigma })`$ : if $`\lambda _j+\lambda `$ is a linear combination of the simple roots in J then we have that $`A_{J,\sigma }v_jk^\times f_\lambda (v_j)`$; otherwise $`A_{J,\sigma }v_j=0`$. Notice that we regard $`f_\lambda (v_j)`$ as an element of $`\mathrm{H}(\lambda )`$ and, as such, $`f_\lambda (v_j)`$ is a $`T`$-eigenvector of weight $`\sigma (\lambda _j)`$.
We now obtain.
###### Proposition 4.4.
Let $`\lambda \mathrm{\Lambda }_+^\sigma `$, then
$$𝒩(\lambda )_\sigma =\underset{JI}{}\underset{\begin{array}{c}wW^{\sigma (J)}\\ I(\lambda )\text{supp}(w)\mathrm{}\end{array}}{}Z_{J,\sigma }^w.$$
Proof. Let $`wW^{\sigma (J)}`$. Assume that $`w\lambda \lambda `$. Note that if $`x`$ is a nonnegative linear combination of the simple roots in $`J`$ then $`\mathrm{ht}(w\sigma (x))\mathrm{ht}(x)`$. Hence,
$$\mathrm{ht}\left(w\sigma (\lambda +x)+\lambda \right)=\mathrm{ht}(w\sigma (x))+\mathrm{ht}(w\lambda +\lambda )>\mathrm{ht}(x).$$
Therefore, $`\rho _{\lambda ,\sigma }((w,1)h_{J,\sigma })`$ is represented by a strictly upper triangular matrix with respect to the chosen basis in 4.3 above. As a consequence for any $`bB`$, $`\rho _{\lambda ,\sigma }((bw,1)h_{J,\sigma })`$ is also represented by a strictly upper triangular matrix. So $`(Bw,1)h_{J,\sigma }𝒩(\lambda )_\sigma `$. Since $`𝒩(\lambda )_\sigma `$ is $`\mathrm{diag}(G)`$-stable it follows $`Z_{J,\sigma }^w=\mathrm{diag}(G)(Bw,1)h_{J,\sigma }𝒩(\lambda )_\sigma `$.
Now assume that $`w\lambda =\lambda `$. Let $`bB`$ and $`z=(bw,1)h_{J,\sigma }`$. Denote by $`A`$ a representative of $`\rho _{\lambda ,\sigma }(z)`$ in $`\text{End}(\mathrm{H}(\lambda ))`$. Let $`V`$ be the subspace of $`\mathrm{H}(\lambda )`$ spanned by $`v_1,\mathrm{},v_{m1}`$. Then $`Av_mk^\times v_m+V`$ and $`AVV`$. Hence, $`A^nv_m0`$ for all $`n`$. Thus $`z𝒩(\lambda )_\sigma `$. ∎
###### Corollary 4.5.
Let $`\lambda ,\mu \mathrm{\Lambda }_+^\sigma `$, then
$$𝒩(\lambda +\mu )_\sigma =𝒩(\lambda )_\sigma 𝒩(\mu )_\sigma .$$
Proof. This follows from the relation $`I(\lambda +\mu )=I(\lambda )I(\mu )`$. ∎
## 5. A compactification of $`G_{\mathrm{sc}}`$
### 5.1
Consider the morphism $`\psi _i:G_{\mathrm{sc}}\left(\text{End}(\mathrm{H}(\omega _i))k\right)`$ defined by $`g[(g\mathrm{Id}_{\mathrm{H}(\omega _i)},1)]`$, where $`\mathrm{Id}_{\mathrm{H}(\omega _i)}`$ denotes the identity map on $`\mathrm{H}(\omega _i)`$ and $`g`$ acts on $`\text{End}(\mathrm{H}(\omega _i))`$ by the left action. Let furthermore $`\pi :G_{\mathrm{sc}}X`$ denote the the natural $`G_{\mathrm{sc}}\times G_{\mathrm{sc}}`$-equivariant morphism and let $`X_{\mathrm{sc}}`$ denote the closure of the image of the product map
$$(\pi ,\underset{iI}{}\psi _i):G_{\mathrm{sc}}X\times \underset{iI}{}(\text{End}(\mathrm{H}(\omega _i)k),$$
Then $`X_{\mathrm{sc}}`$ is a projective $`G_{\mathrm{sc}}\times G_{\mathrm{sc}}`$-equivariant variety containing $`G_{\mathrm{sc}}`$ as an open subset. Unlike $`X`$ the variety $`X_{\mathrm{sc}}`$ need not be smooth and in general it is not even normal (e.g. in type $`A_3`$). By abuse of notation we use the notation $`\pi `$ and $`\psi _i`$, $`iI`$, for the natural extensions of the corresponding maps to $`X_{\mathrm{sc}}`$.
###### Lemma 5.2.
The projective morphism $`\pi :X_{\mathrm{sc}}X`$ defines a bijection between $`X_{\mathrm{sc}}G_{\mathrm{sc}}`$ and $`XG`$. In particular, $`\pi `$ is a finite morphism.
Proof. Notice that the $`G_{\mathrm{sc}}\times G_{\mathrm{sc}}`$-invariant homogeneous polynomial function on $`\text{End}(\mathrm{H}(\omega _i))k`$ defined by $`(f,a)\mathrm{det}(f)a^{\mathrm{dim}_k(\mathrm{H}(\omega _i))}`$ vanishes at $`(\mathrm{Id}_{\mathrm{H}(\omega _i)},1)`$. Thus
$$X_{\mathrm{sc}}X\times \underset{iI}{}\left(\right(\text{End}(\mathrm{H}(\omega _i)k)(0k)),$$
and we may consider the following commutative diagram
$$\text{},$$
where all the maps are the natural ones. Assume now that $`x`$ is an element of the boundary $`X_{\mathrm{sc}}G_{\mathrm{sc}}`$. As the dimensions of $`G_{\mathrm{sc}}`$ and $`X_{\mathrm{sc}}`$ coincide the $`(G_{\mathrm{sc}},1)`$-stabilizer of $`x`$ has strictly positive dimension. In particular, the images $`\psi _i(x)=[(f_i,a_i)]`$, $`iI`$, have the same property. Thus, the endomorphism $`f_i`$ is not invertible and thus $`a_i=0`$. This proves that
$$X_{\mathrm{sc}}G_{\mathrm{sc}}X\times \underset{iI}{}(\text{End}(\mathrm{H}(\omega _i)),$$
and hence $`\pi `$ maps $`X_{\mathrm{sc}}G_{\mathrm{sc}}`$ injectively to the boundary $`XG`$. As $`\pi `$ is dominant and projective, and thus surjective, this proves the first assertion. Finally $`\pi `$ is finite as it is projective and quasifinite. ∎
### 5.3
For a dominant weight $`\lambda \mathrm{\Lambda }_+`$ let $`\psi _\lambda :G_{\mathrm{sc}}\left(\text{End}(\mathrm{H}(\lambda ))k\right)`$ be the morphisms defined by $`\psi _\lambda (g)=[(g\mathrm{Id}_{\mathrm{H}(\lambda )},1)]`$, for $`gG_{\mathrm{sc}}`$. Then we let $`X_{\mathrm{sc}}^\lambda `$ denote the closure of the image of the map
$$G_{\mathrm{sc}}X_{\mathrm{sc}}\times \left(\text{End}(\mathrm{H}(\lambda ))k\right)$$
defined as the product of the inclusion $`G_{\mathrm{sc}}X_{\mathrm{sc}}`$ and $`\psi _\lambda `$. We claim
###### Lemma 5.4.
The canonical morphism $`\pi ^\lambda :X_{\mathrm{sc}}^\lambda X_{\mathrm{sc}}`$ is an isomorphism.
Proof. Let $`X_0`$ be the complement of the union of the closures $`\overline{Bs_iB^{}}`$, $`iI`$, within $`X`$, and let $`X_0^{}=\overline{T}X_0`$. Then, by \[B-K, Prop.6.2.3(i)\], the natural morphism
$$U\times U^{}\times X_0^{}X_0,$$
$$(g,h,x)(g,h)x,$$
is an isomorphism of varieties. Let $`X_{\mathrm{sc},0}=\pi ^1(X_0)`$ and $`X_{\mathrm{sc},0}^{}`$ denote the (scheme theoretic) inverse image $`\pi ^1(X_0^{})`$. As $`\pi `$ is $`G_{\mathrm{sc}}\times G_{\mathrm{sc}}`$-equivariant we obtain an induced isomorphism
$$U\times U^{}\times X_{\mathrm{sc},0}^{}X_{\mathrm{sc},0}.$$
In particular, $`X_{\mathrm{sc},0}^{}`$ is an irreducible closed subvariety of the open subvariety $`X_{\mathrm{sc},0}`$ containing $`T_{\mathrm{sc}}`$ as an open subset. Thus, $`X_{\mathrm{sc},0}^{}`$ is contained in the closure of $`T_{\mathrm{sc}}`$ within $`X_{\mathrm{sc}}`$. Let $`\pi _\lambda `$ denote the composition of $`\pi `$ and $`\pi ^\lambda `$. Defining $`X_{\mathrm{sc},0}^\lambda =\pi _\lambda ^1(X_0)`$ and $`(X_{\mathrm{sc},0}^\lambda )^{}=\pi _\lambda ^1(X_0^{})`$ we, similarly, obtain an isomorphism
$$U\times U^{}\times (X_{\mathrm{sc},0}^\lambda )^{}X_{\mathrm{sc},0}^\lambda .$$
Notice that $`X_{\mathrm{sc},0}^\lambda =\pi ^\lambda (X_{\mathrm{sc},0})`$. Moreover, the $`G\times G`$-translates of $`X_0`$ cover $`X`$ \[B-K, Thm.6.1.8\]. Thus, it suffices to show that the morphism $`(X_{\mathrm{sc},0}^\lambda )^{}X_{\mathrm{sc},0}^{}`$ induced by $`\pi ^\lambda `$ is an isomorphism. This will follow if $`\pi ^\lambda `$ induces an isomorphism between the closures of $`T_{\mathrm{sc}}`$ in $`X_{\mathrm{sc}}`$ and $`X_{\mathrm{sc}}^\lambda `$. Determining the latter closures of $`T_{\mathrm{sc}}`$ and checking that they are isomorphic is now an easy exercise. ∎
It follows that we may consider $`\psi _\lambda `$ as the extended morphism
$$\psi _\lambda :X_{\mathrm{sc}}\left(\text{End}(\mathrm{H}(\lambda ))k\right),$$
which we will do in the following. As in the proof of Lemma 5.2 we may prove that
$$\psi _\lambda (X_{\mathrm{sc}})\left(\left(\text{End}(\mathrm{H}(\lambda ))k\right)\left(0k\right)\right)$$
and that the induced map $`X_{\mathrm{sc}}\left(\text{End}(\mathrm{H}(\lambda ))\right)`$ is compatible with $`\pi :X_{\mathrm{sc}}X`$ and the map $`\rho _\lambda :X\left(\text{End}(\mathrm{H}(\lambda ))\right)`$.
### 5.5
The variety $`X_{\mathrm{sc}}`$ is a compactification of $`G_{\mathrm{sc}}`$ with the $`G_{\mathrm{sc}}\times G_{\mathrm{sc}}`$ action defined in the natural way. Let $`X_{\mathrm{sc},\sigma }`$ be the $`G_{\mathrm{sc}}\times G_{\mathrm{sc}}`$-variety which as a variety is isomorphic to $`X_{\mathrm{sc}}`$ and where the $`G_{\mathrm{sc}}\times G_{\mathrm{sc}}`$-action is twisted by the morphism $`G_{\mathrm{sc}}\times G_{\mathrm{sc}}G_{\mathrm{sc}}\times G_{\mathrm{sc}}`$ , $`(g,h)(g,\sigma (g))`$ for $`g,hG_{\mathrm{sc}}`$. Thus we may identify $`G_{\mathrm{sc},\sigma }`$ of 3.3 with an open subset of $`X_{\mathrm{sc},\sigma }`$. Notice that by Lemma 5.2 we may identify the boundaries of $`X_{\mathrm{sc},\sigma }`$ and $`X_\sigma `$ and we may therefore also regard $`Z_{J,\sigma }^w`$, for $`JI`$, as subsets of $`X_{\mathrm{sc},\sigma }`$.
## 6. Steinberg fibers and trace maps
### 6.1
Let $`\mathrm{Tr}_i`$ denote the trace function on $`\text{End}(\mathrm{H}(\omega _{𝒞_i}))`$. To each $`a_ik`$ we may associate a global section $`(\mathrm{Tr}_i,a_i)`$ of the line bundle $`𝒪_i(1):=𝒪_{\left(\text{End}(\mathrm{H}(\omega _{𝒞_i}))k\right)}(1)`$ on $`\left(\text{End}(\mathrm{H}(\omega _{𝒞_i}))k\right)`$. The pull back of $`(\mathrm{Tr}_i,a_i)`$ to $`X_{\mathrm{sc},\sigma }`$, by the morphism $`\psi _{\omega _{𝒞_i}}`$, is then a global section $`f_{i,a_i}^\sigma `$ of a line bundle on $`X_{\mathrm{sc},\sigma }`$. In the following, we will study the common zero set $`Z(a_1,\mathrm{},a_l)`$ of the sections $`f_{i,a_i}^\sigma `$, for varying $`a_ik`$. By choosing a trivialization of the pull back of $`𝒪_i(1)`$ to $`G_{\mathrm{sc},\sigma }`$ we may think of $`f_{i,a_i}^\sigma `$ as a function on $`G_{\mathrm{sc},\sigma }`$ and, by abuse of notation, we also denote this function by $`f_{i,a_i}^\sigma `$. We assume that the trivialization is chosen such that $`f_{i,a}^\sigma f_{i,0}^\sigma =a`$ as functions on $`G_{\mathrm{sc}}`$ (then the trivialization is actually unique). Then $`f_{i,a_i}^\sigma `$ is invariant under $`\sigma `$-conjugation by $`G_{\mathrm{sc}}`$ and thus $`f_{i,a_i}^\sigma `$ induces a morphism $`\overline{f}_{i,a_i}^\sigma :G_{\mathrm{sc}}//_\sigma G_{\mathrm{sc}}k`$. We then claim
###### Proposition 6.2.
The product morphism
$$(\overline{f}_{1,0}^\sigma ,\overline{f}_{2,0}^\sigma ,\mathrm{},\overline{f}_{l,0}^\sigma ):G_{\mathrm{sc}}//_\sigma G_{\mathrm{sc}}𝔸^l,$$
is an isomorphism.
Proof. Let $`f_i^\sigma `$ denote the restriction of $`f_{i,0}^\sigma `$ to $`T_{\mathrm{sc}}`$. An easy calculation shows that $`f_i^\sigma `$ equals
$$w_0𝒞_i+\underset{\begin{array}{c}\lambda \mathrm{\Lambda }^\sigma ,\lambda 𝒞_i\\ \mathrm{H}(𝒞_i)_\lambda 0\end{array}}{}q_{i,\lambda }\lambda .$$
Thus $`f_i^\sigma `$ is contained in the semigroup algebra $`k[\mathrm{\Lambda }^\sigma ]`$ generated by the $`\sigma `$-invariant weights of $`\mathrm{\Lambda }`$. Moreover, $`f_i^\sigma `$ is invariant under the action of the group $`W^\sigma `$ of $`\sigma `$-invariant elements of $`W`$. Hence, $`f_i^\sigma `$ is an element of the polynomial ring $`k[\mathrm{\Lambda }^\sigma ]^{W^\sigma }`$ in $`l`$ variables \[Sp3, proof of Cor.2\], generated (as an $`k`$-algebra) by the elements
$$\mathrm{sym}(𝒞_i):=\underset{wW^\sigma }{}w𝒞_i,i=1,\mathrm{},l.$$
As in \[Ste, proof of Lem.6.3\] we conclude that also $`f_i^\sigma `$, $`i=1,\mathrm{},l`$, generates $`k[\mathrm{\Lambda }^\sigma ]^{W^\sigma }`$ as an $`k`$-algebra. Now apply \[Sp3, Thm.1\]. ∎
###### Corollary 6.3.
The intersection of $`Z(a_1,\mathrm{},a_l)`$ with the boundary $`X_{\mathrm{sc},\sigma }G_{\mathrm{sc},\sigma }`$ of $`X_{\mathrm{sc},\sigma }`$ is independent of $`a_1,\mathrm{},a_l`$. Moreover, the intersection $`Z(a_1,\mathrm{},a_l)G_{\mathrm{sc},\sigma }`$ is a single Steinberg fiber.
Proof. As in the proof of Lemma 5.2 it follows that $`x`$ is an element of $`X_{\mathrm{sc},\sigma }G_{\mathrm{sc},\sigma }`$ exactly when the image $`\psi _{\omega _{𝒞_i}}(x)`$ is of the form $`[(f,0)]`$ for all $`i=1,\mathrm{},l`$. Thus, the section $`f_{i,a_i}^\sigma `$ coincides with $`f_{i,0}^\sigma `$ on the boundary of $`X_{\mathrm{sc},\sigma }`$. This proves the first statement. The latter statement follows by Proposition 6.2. ∎
## 7. Proofs of the main results
###### Lemma 7.1.
Let $`JI`$, $`wW^{\sigma (J)}`$ and $`bB`$. If $`f_{i,0}^\sigma ((bw,1)h_{J,\sigma })=0`$, then either (1) $`w\omega _{𝒞_i}\omega _{𝒞_i}`$ or (2) $`𝒞_iJ`$ and $`w\alpha _j=\alpha _j`$ for all $`j𝒞_i`$.
Proof. Assume that $`w\omega _{𝒞_i}=\omega _{𝒞_i}`$. Then the diagonal entry of a representative $`A`$ of $`\rho _{\omega _{𝒞_i},\sigma }((bw,1)h_{J,\sigma })`$ associated to the lowest weight space is nonzero. In particular, the relation $`f_{i,0}^\sigma ((bw,1)h_{J,\sigma })=0`$ cannot be satisfied unless there exists a weight $`x\omega _{𝒞_i}\omega _{𝒞_i}`$ of $`\mathrm{H}(\omega _{𝒞_i})`$ satisfying $`x=_{jJ}a_j\alpha _j`$, with $`a_j\{0\}`$, and $`w\sigma (x)=x`$.
Let $`KJ`$ denote the set of $`jJ`$ such that $`a_j0`$. As $`x\omega _{𝒞_i}`$ is a weight of $`\mathrm{H}(\omega _{𝒞_i})`$ we know that $`𝒞_iK`$ is nonempty. Now $`_{jK}a_jw\alpha _{\sigma (j)}=_{jK}a_j\alpha _j`$ and thus $`_{jK}a_j(\mathrm{ht}(w\alpha _{\sigma (j)})\mathrm{ht}(\alpha _j))=0`$. As $`wW^{\sigma (J)}`$ we conclude that $`\mathrm{ht}(w\alpha _{\sigma (j)})1`$ and consequently $`w\alpha _{\sigma (j)}`$ is a simple root for all $`jK`$. By the assumption $`w\omega _{𝒞_i}=\omega _{𝒞_i}`$ we know that $`w\alpha _{\sigma (j)}=\alpha _{\sigma (j)}`$ for each $`j𝒞_iK`$. In particular, when $`j𝒞_iK`$ then $`a_{\sigma (j)}=a_j`$. Hence, $`𝒞_iK`$ is invariant under $`\sigma `$ and as $`𝒞_i`$ is a single $`\sigma `$-orbit we have $`𝒞_iK=𝒞_i`$. This ends the proof. ∎
###### Lemma 7.2.
Let $`JI`$. Then
$$Z(a_1,\mathrm{},a_l)Z_{J,\sigma }=\underset{\begin{array}{c}wW^{\sigma (J)}\\ \text{supp}_\sigma (w)=I\end{array}}{}Z_{J,\sigma }^w.$$
Proof. By Corollary 6.3 it suffices to consider the case when all $`a_i`$ are zero. By Proposition 4.4,
$$\underset{JI}{}\underset{\begin{array}{c}wW^{\sigma (J)}\\ \text{supp}_\sigma (w)=I\end{array}}{}Z_{J,\sigma }^w=_i𝒩(\omega _{𝒞_i})_\sigma Z(0,\mathrm{},0).$$
This proves one inclusion. For $`zZ(0,\mathrm{},0)Z_{J,\sigma }`$, we have that $`z=(g,g)(bw,1)h_{J,\sigma }`$ for some $`gG,bB`$ and $`wW^{\sigma (J)}`$. Then $`f_{i,0}^\sigma ((bw,1)h_{J,\sigma })=0`$ for all $`i=1,\mathrm{},l`$. It suffices to prove that $`\text{supp}_\sigma (w)=I`$.
If $`w=1`$, then by Lemma 7.1, $`𝒞_iJ`$ for each $`\sigma `$-orbit $`𝒞_i`$. Thus $`I=J`$, which contradicts our assumption. Now assume that $`w1`$ and that $`\text{supp}_\sigma (w)I`$. As $`G`$ is assumed to be simple, there exist simple roots $`\alpha _i`$ and $`\alpha _j`$ with $`n=\alpha _j,\alpha _i^{}0`$ satisfying that $`i\text{supp}_\sigma (w)`$ and $`j\text{supp}_\sigma (w)`$. Let $`𝒞_i`$ and $`𝒞_j`$ denote the associated $`\sigma `$-orbits of $`\alpha _i`$ and $`\alpha _j`$. As $`\text{supp}_\sigma (w)`$ is $`\sigma `$-stable it follows that $`𝒞_i\text{supp}_\sigma (w)`$ and $`𝒞_jI\text{supp}_\sigma (w)`$.
Now there exists $`m`$, such that $`\sigma ^m(i)\text{supp}(w)`$ and thus $`w\omega _{\sigma ^m(i)}\omega _{\sigma ^m(i)}`$. Hence redefining, if necessary, $`\alpha _i`$ and $`\alpha _j`$, we may assume that $`w\omega _i\omega _i`$. Consider then the relation $`\alpha _j=2\omega _jn\omega _i\lambda `$ with $`\lambda `$ denoting a dominant weight. Then Lemma 7.1, applied to $`𝒞_j`$, implies that $`w\alpha _j=\alpha _j`$ and $`w\omega _j=\omega _j`$ and thus $`w(n\omega _i+\lambda )=n\omega _i+\lambda `$. As both $`\omega _i`$ and $`\lambda `$ are dominant we conclude that $`w\omega _i=\omega _i`$ which is a contradiction. ∎
Now we will prove the main theorem.
###### Theorem 7.3.
Let $`F`$ be a Steinberg fiber of $`G_\sigma `$ and $`\overline{F}`$ its closure in $`X_\sigma `$. Then
$$\overline{F}F=\underset{JI}{}\underset{\begin{array}{c}wW^{\sigma (J)}\\ \text{supp}_\sigma (w)=I\end{array}}{}Z_{J,\sigma }^w,$$
which also coincides with the set $`Z(a_1,\mathrm{},a_l)(X_{\mathrm{sc},\sigma }G_{\mathrm{sc},\sigma })`$ for all $`a_1,\mathrm{},a_l`$.
Proof. Let $`C`$ be an irreducible component of $`Z(a_1,\mathrm{},a_l)`$. Then by Krull’s principal ideal theorem, $`dim(C)dim(G_{\mathrm{sc}})l`$. By \[L2, 8.20\],
$$dim(Z_{J,\sigma }^w)=dim(G)l(w)|IJ|<dim(G_{\mathrm{sc}})l,$$
for $`JI`$ and $`wW^{\sigma (J)}`$ with $`\text{supp}_\sigma (w)=I`$. Thus by Lemma 7.2,
$$dim\left(C(X_{\mathrm{sc},\sigma }G_{\mathrm{sc},\sigma })\right)<dim(G_{\mathrm{sc}})ldim(C).$$
Hence $`CG_{\mathrm{sc},\sigma }`$ is dense in $`C`$. But, by Corollary 6.3, the intersection $`Z(a_1,\mathrm{},a_l)G_{\mathrm{sc},\sigma }`$ is a single Steinberg fiber $`F(a_1,\mathrm{},a_n)`$ which, as in 3.2, is irreducible. We conclude that $`C`$ is contained in the closure of $`F(a_1,\mathrm{},a_l)`$, and thus the closure of $`F(a_1,\mathrm{},a_l)`$ is $`Z(a_1,\mathrm{},a_l)`$. In particular, $`Z(a_1,\mathrm{},a_l)`$ is irreducible.
Let $`F`$ be a Steinberg fiber of $`G_\sigma `$. Then $`F=\pi (F(a_1,\mathrm{},a_l))`$ for some $`a_1,\mathrm{},a_lk`$. Hence $`\overline{F}=\pi (Z(a_1,\mathrm{},a_l))`$. The statement now follows from Lemma 7.2 and Lemma 5.2. ∎
###### Remark.
1. We call an element $`wW`$ a $`\sigma `$-twisted Coxeter element if $`l(w)=l`$ and $`\text{supp}_\sigma (w)=I`$. (The notation of twisted Coxeter elements was first introduced by Springer in \[Sp1\]. Our definition is slightly different from his). It follows easily from Theorem 7.3 that $`\overline{Z_{I\{i\},\sigma }^w}`$ are the irreducible components of $`\overline{F}F`$, where $`iI`$ and $`w`$ runs over all $`\sigma `$-twisted Coxeter elements that are contained in $`W^{I\{\sigma (i)\}}`$.
2. By the proof of Theorem 7.3 we may also deduce that the closure of a Steinberg fiber $`F`$ within $`X_{\mathrm{sc},\sigma }`$ coincides with $`Z(a_1,\mathrm{},a_l)`$ for certain uniquely determined $`a_1,\mathrm{},a_l`$ depending on $`F`$. This result may be considered as an extension of Corollary 2 in \[Sp3\] to the compactification $`X_{\mathrm{sc},\sigma }`$ of $`G_{\mathrm{sc},\sigma }`$. More precisely, notice that the statement of \[Sp3, Corollary 2\] is equivalent to saying that a Steinberg fiber $`F`$ of $`G_{\mathrm{sc},\sigma }`$ is the common zero set of the functions $`f_{i,a_i}^\sigma `$ for uniquely determined $`a_1,\mathrm{},a_l`$. Here we think of $`f_{i,a_i}^\sigma `$ as regular functions on $`G_{\mathrm{sc},\sigma }`$ as explained in 6.1. When generalizing to $`X_{\mathrm{sc},\sigma }`$ the only difference is that we have to regard $`\overline{f}_{i,a_i}^\sigma `$ as sections of certain line bundles on $`X_{\mathrm{sc},\sigma }`$.
Similar to \[H1, 4.6\], we have the following consequence.
###### Corollary 7.4.
Assume that $`G_\sigma `$ is defined and split over $`𝔽_q`$, then for any Steinberg fiber $`F`$ of $`G_\sigma `$, the number of $`𝔽_q`$-rational points of $`\overline{F}F`$ is
$$(\underset{wW}{}q^{l(w)})(\underset{\text{supp}_\sigma (w)=I}{}q^{l(w_0w)+L(w_0w)}),$$
where $`w_0`$ is the maximal element of $`W`$ and for $`wW`$, $`l(w)`$ is its length and $`L(w)`$ is the number of simple roots $`\alpha `$ satisfying $`w\alpha <0`$.
## Acknowledgements
We thank J.C.Jantzen, G.Lusztig and T.A.Springer for some useful discussions and comments. |
warning/0506/hep-ph0506251.html | ar5iv | text | # A Variable-Flavour-Number Scheme at NNLO
## Abstract
I present a formulation of a Variable Flavour Number Scheme for heavy quarks that is implemented up to NNLO in the strong coupling constant and may be used in NNLO global fits for parton distributions.
###### Keywords:
QCD, Structure Functions, Heavy Quarks
While up, down and strange quarks are treated as effectively massless partons, charm, bottom and top have to be regarded as heavy partons. There are two distinct regimes for these types of quarks. At low scales, $`Q^2m_H^2`$, they are only created in the final state and described using the Fixed Flavour Number Scheme (FFNS)
$$F_i(x,Q^2)=C_{i,k}^{FF}(Q^2/m_H^2)f_k^{n_f}(Q^2).$$
However, for $`Q^2m_H^2`$, the coefficient functions contain large $`\mathrm{ln}(Q^2/m_H^2)`$ terms, spoiling the perturbative expansion. In this regime it is more appropriate to treat the quarks like massless partons, and the large $`\mathrm{ln}(Q^2/m_H^2)`$ terms are summed via the DGLAP evolution equations. The simplest recipe involving this regime is the Zero Mass Variable Flavour Number Scheme (ZMVFNS). This ignores all $`𝒪(m_H^2/Q^2)`$ corrections, i.e.
$$F_i(x,Q^2)=C_{i,j}^{ZMVF}f_j^{n_f+1}(Q^2).$$
The partons in different flavour-number regions are related perturbatively,
$$f_k^{n_f+1}(Q^2)=A_{jk}(Q^2/m_H^2)f_k^{n_f}(Q^2),$$
where the perturbative matrix elements $`A_{jk}(Q^2/m_H^2)`$ containing $`\mathrm{ln}(Q^2/m_H^2)`$ terms guarantee the correct evolution for both descriptions. At LO, i.e. zeroth order in $`\alpha _S`$, the relationship between the two descriptions is trivial – $`q(g)_k^{n_f+1}(Q^2)q(g)_k^{n_f}(Q^2).`$ At NLO, i.e. first order in $`\alpha _S`$ ($`h^+(Q^2)=(h+\overline{h})(Q^2)`$),
$$h^+(Q^2)=\frac{\alpha _S}{4\pi }P_{qg}^0g^{n_f}(Q^2)\mathrm{ln}\left(\frac{Q^2}{m_H^2}\right),g^{n_f+1}(Q^2)=\left(1\frac{\alpha _S}{6\pi }\mathrm{ln}\left(\frac{Q^2}{m_H^2}\right)\right)g^{n_f}(Q^2),$$
i.e. the heavy flavour evolves from zero at $`Q^2=m_H^2`$ and the gluon loses corresponding momentum. It is natural to choose $`Q^2=m_H^2`$ as the transition point. At NNLO, i.e. second order in $`\alpha _S`$, there is much more complication
$$f_i^{n_f+1}(Q^2)=\left(\frac{\alpha _S}{4\pi }\right)^2\underset{ij}{}(A_{ij}^{2,0}+A_{ij}^{2,1}\mathrm{ln}(Q^2/m_H^2)+A_{ij}^{2,2}\mathrm{ln}^2(Q^2/m_H^2))f_j^{n_f}(Q^2),$$
where $`A_{ij}^{2,0}`$ is generally non-zero buza . There is no longer a smooth transition at this order, and in fact the heavy parton begins with a negative value at small $`x`$.
This leads to discontinuities in the partons and, without the correct treatment, also in the structure functions. ZMVFNS coefficient functions also lead to discontinuities at the transition point due to a sudden change in the flavour number in the coefficient functions. (This is already true at NLO, i.e. $`F_2^H(x,Q^2)=0Q^2<m_H^2,=\frac{\alpha _S}{4\pi }C_{2,g}g^{n_f+1}(Q^2)Q^2>m_H^2`$, but the effect is very small.) This is a large effect at NNLO and is also negative at smallish $`x`$ ($`x0.001`$). Hence, ZMVFNS is not really feasible at NNLO, leading to a huge discontinuity in $`F_2^c(x,Q^2)`$, which is significant in $`F_2^{Tot}(x,Q^2)`$, as shown in Fig. 1.
Hence we need a general Variable Flavour Number Scheme (VFNS) interpolating between the two well-defined limits of $`Q^2m_H^2`$ and $`Q^2m_H^2`$. The VFNS can be defined by demanding equivalence of the $`n_f`$ and $`n_f+1`$-flavour descriptions at all orders,
$`F_i(x,Q^2)`$ $`=`$ $`C_{i,k}^{FF}(Q^2/m_H^2)f_k^{n_f}(Q^2)=C_{i,j}^{VF}(Q^2/m_H^2)f_j^{n_f+1}(Q^2)`$
$``$ $`C_{i,j}^{VF}(Q^2/m_H^2)A_{jk}(Q^2/m_H^2)f_k^{n_f}(Q^2)`$
$`C_{i,k}^{FF}(Q^2/m_H^2)`$ $`=`$ $`C_{i,j}^{VF}(Q^2/m_H^2)A_{jk}(Q^2/m_H^2).`$
At $`𝒪(\alpha _S)`$ this gives
$$C_{2,g}^{FF,1}(Q^2/m_H^2)=C_{2,HH}^{VF,0}(Q^2/m_H^2)P_{qg}^0\mathrm{ln}(Q^2/m_H^2)+C_{2,g}^{VF,1}(Q^2/m_H^2).$$
The VFNS coefficient functions tend to the massless limits as $`Q^2/m_H^2\mathrm{}`$, as demonstrated to all orders in collins , and if we use the zeroth order cross-section for photon-heavy quark scattering we obtain the original ACOT scheme acot .
However, $`C_{2,HH}^{VF,0}(Q^2/m_H^2)`$ is only uniquely defined as $`Q^2/m_H^2\mathrm{}`$, i.e. one can swap $`𝒪(m_H^2/Q^2)`$ terms between $`C_{2,HH}^{VF,0}(Q^2/m_H^2)`$ and $`C_{2,g}^{VF,1}(Q^2/m_H^2)`$. Similar reasoning holds for $`C_{2,HH}^{VF,n}(Q^2/m_H^2)`$. The ACOT prescription violated the threshold $`W^2=Q^2(1x)/x>4M^2`$ since only one quark was needed in final state. The Thorne-Roberts variable flavour number scheme (TR-VFNS) trvfns recognized this ambiguity and removed it by imposing continuity of $`(dF_2/d\mathrm{ln}Q^2)`$ at the transition point. This guaranteed smoothness at $`Q^2=m_H^2`$, but was complicated and cumbersome when extended to higher orders.
There have been other alternatives, and most recently the ACOT($`\chi `$) prescription acotchi defines $`F_2^{H,0}(x,Q^2)=h^+(x/x_{max},Q^2)`$, where $`x_{max}=Q^2/(Q^2+4m_H^2)`$. The coefficient functions tend to the massless limit for $`Q^2/m_H^2\mathrm{}`$ but also respect the threshold requirement $`W^24m_H^2`$ for quark-antiquark production. Moreover it is very simple. For the VFNS to remain simple (and physical) at all orders I choose $`C_{2,HH}^{VF,n}(Q^2/m_H^2,z)=C_{2,HH}^{ZM,n}(z/x_{max}).`$<sup>1</sup><sup>1</sup>1It is also important to choose $`C_{L,HH}^{VF,n}(Q^2/m_H^2,z)C_{L,HH}^{ZM,n}(z/x_{max}).`$ Adopting this convention then at NNLO we have, for example,
$$C_{2,Hg}^{VF,2}\left(\frac{Q^2}{m_H^2}\right)=C_{2,Hg}^{FF,2}\left(\frac{Q^2}{m_H^2}\right)C_{2,HH}^{ZM,1}\left(\frac{z}{x_{max}}\right)A_{Hg}^1\left(\frac{Q^2}{m_H^2}\right)C_{2,HH}^{ZM,0}\left(\frac{z}{x_{max}}\right)A_{Hg}^2\left(\frac{Q^2}{m_H^2}\right).$$
Since $`A_{Hg}^2(1,z)0`$, $`C_{2,Hg}^2(Q^2/m_H^2,z)`$ is discontinuous at $`Q^2=m_H^2`$, and this compensates exactly for the discontinuity in the heavy flavour parton distribution.<sup>2</sup><sup>2</sup>2At NNLO there are also contributions due to heavy flavours in loops away from the photon vertex. These are included within the VFNS and lead to a discontinuity in the coefficient functions for light flavours cancelling that in the light quark distributions. Strictly, part of this contribution should be interpreted as light flavour structure functions, while part of it contributes to $`F_2^H(x,Q^2)`$ smith .
There is one more issue in defining the VFNS: the ordering for $`F_2^H(x,Q^2)`$, i.e.
$`n_f`$-flavour $`n_f+1`$-flavour
LO $`\frac{\alpha _S}{4\pi }C_{2,Hg}^{FF,1}g^{n_f}`$ $`C_{2,HH}^{VF,0}h^+`$
NLO $`\left(\frac{\alpha _S}{4\pi }\right)^2(C_{2,Hg}^{FF,2}g^{n_f}+C_{2,Hq}^{FF,2}\mathrm{\Sigma }^{n_f})`$ $`\frac{\alpha _S}{4\pi }(C_{2,HH}^{VF,1}h^++C_{2,Hg}^{FF,1}g^{n_{f+1}}).`$
Switching directly when going from $`n_f`$ to $`n_f+1`$ flavours leads to a discontinuity. We must decide how to deal with this. Up to now ACOT have used e.g. at NLO
$$\frac{\alpha _S}{4\pi }C_{2,Hg}^{FF,1}g^{n_f}\frac{\alpha _S}{4\pi }(C_{2,HH}^{VF,1}h^++C_{2,Hg}^{FF,1}g^{n_f+1})+C_{2,HH}^{VF,0}h^+,$$
i.e. the same order of $`\alpha _S`$ above and below, but LO below and NLO above. The Thorne-Roberts scheme proposed e.g. at LO
$$\frac{\alpha _S(Q^2)}{4\pi }C_{2,Hg}^{FF,1}\left(\frac{Q^2}{m_H^2}\right)g^{n_f}(Q^2)\frac{\alpha _S(m_H^2)}{4\pi }C_{2,Hg}^{FF,1}(1)g^{n_f}(m_H^2)+C_{2,HH}^{VF,0}\left(\frac{Q^2}{m_H^2}\right)h^+(Q^2)$$
i.e. the higher order $`\alpha _S`$ term is frozen when going upwards through $`Q^2=m_H^2`$. This difference in choice is extremely important at low $`Q^2`$.
Making this choice, in order to define the VFNS at NNLO we need the $`𝒪(\alpha _S^3)`$ heavy flavour coefficient functions for $`Q^2m_H^2`$. However, these are not yet calculated (making a NNLO FFNS problematic). We know the leading threshold logarithms thresh , and can derive the leading $`ln(1/x)`$ term from $`k_T`$-dependent impact factors asymp ,
$$C_{2,Hg}^{FF,3,lowx}(Q^2/m_H^2,z)=96\frac{\mathrm{ln}(1/z)}{z}f(Q^2/m_H^2),f(1)4,$$
and $`C_{2,Hq}^{FF,3,lowx}(Q^2/m_H^2,z)=4/9C_{2,Hg}^{FF,3,lowx}(Q^2/m_H^2,z).`$ By analogy with the known NNLO coefficient functions and splitting functions I hypothesize that
$$C_{2,Hg}^{FF,3,lowx}(Q^2/m_H^2,z)=\frac{96}{z}(\mathrm{ln}(1/z)4)(1z/x_{max})^{20}f(Q^2/m_H^2),$$
i.e. the leading $`\mathrm{ln}(1/z)`$ term is always accompanied by $`4`$, and the effect of the small $`z`$ term is damped as $`z1`$. Using the full (if slightly approximate) VFNS one can produce NNLO predictions for charm with discontinuous partons, but a continuous $`F^c(x,Q^2)`$. NNLO clearly improves the match to lowest $`Q^2`$ data zeus ; h1 , where NLO is generally too low, as seen in Fig. 2.
Hence, I have devised a full NNLO VFNS, with a small amount of necessary modelling. This seems to improve the fit to the lowest $`x`$ and $`Q^2`$ data greatly. It also guarantees continuity of the physical observables, such as structure functions, despite the discontinuity in NNLO parton distributions. It can now be used in a full NNLO global analysis. |
warning/0506/astro-ph0506748.html | ar5iv | text | # Observational constraints on low redshift evolution of dark energy: How consistent are different observations?
## I Introduction
Observational evidence for accelerated expansion in the universe has been growing in the last two decades crisis2 . Observations of high redshift supernovae nova\_data1 ; nova\_data3 provided an independent confirmation. Using these along with observations of cosmic microwave background radiation (CMB) boomerang ; wmap\_params and large scale structure 2df ; sdss , we can construct a “concordance” model for cosmology and study variations around it (e.g., see wmap\_params ; 2003Sci…299.1532B ; 2004PhRvD..69j3501T ; for an overview of our current understanding, see TPabhay ).
Observations indicate that dark energy should have an equation of state parameter $`wP/\rho <1/3`$ for the universe to undergo accelerated expansion. Indeed, observations show that dark energy is the dominant component of our universe. The cosmological constant is the simplest explanation for accelerated expansion ccprob ; review3 and it is known to be consistent with observations. In order to avoid theoretical problems related to cosmological constant ccprob , other scenarios have been investigated. In these models one can have $`w1`$ and in general $`w`$ varies with redshift. These models include quintessence quint1 , k-essence 2001PhRvD..63j3510A , tachyons tachyon1 ; 2003PhRvD..67f3504B , phantom fields 2002PhLB..545…23C , branes brane1 , etc. There are also some phenomenological models water , field theoretical and renormalisation group based models (see e.g. tp173 ), models that unify dark matter and dark energy unified\_dedm1 and many others like those based on horizon thermodynamics (e.g. see 2005astro.ph..5133S ). Even though these models have been proposed to overcome the fine tuning problem for cosmological constant, most of these require similar fine tuning of parameter(s) to be consistent with observations. Nevertheless, they raise the possibility of $`w(z)`$ evolving with time (or being different from $`1`$), which — in principle — can be tested by observations.
Given that $`w<1/3`$ for dark energy for the universe to undergo accelerated expansion, the energy density of this component changes at a much slower rate than that of matter and radiation. Indeed, $`w=1`$ for cosmological constant and in this case the energy density is a constant. Unless $`w`$ is a rapidly varying function of redshift and becomes $`w0`$ at ($`z1`$), the energy density of the dark energy component should be negligible at high redshifts ($`z1`$) as compared to that of non-relativistic matter. If dark energy evolves in a manner such that its energy density is comparable to, or greater than the matter density in the universe at high redshifts then the basic structure of the cosmological model needs to be modified. We do not consider such models here. We confine our attention to models with dark energy density being an insignificant component of the universe at $`z1`$ and choose observations which are sensitive to evolution of $`w(z)`$ at redshifts $`z1`$.
To put the present work in context, we recall that combining supernova observations with the WMAP data provides strong constraints on the variation of dark energy density 2005MNRAS.356L..11J . (A review of relevant observations for constraining dark energy models along with a summary of the previous work in this area is given in Section II.2.) Reproducing the location of acoustic features requires the angular diameter distance to the last scattering surface to be in the correct range. This analysis showed that while the data from SN observations allows for a large range in parameters of dark energy, combining with WMAP data limits this range significantly. However, in that work, we did not explore the cosmological parameter space widely and had fixed nearly all parameters other than those used to describe evolution of dark energy. In the present work, we allow many cosmological parameters to vary and include constraints from cluster abundance in addition to the supernova and WMAP constraints.
In addition to obtaining quantitative bounds on parameters in different contexts, we address the following key issues in this paper:
* Does allowing cosmological parameters to vary weaken the constraints on variation of dark energy?
* Conversely, how does the allowed range for different cosmological parameters change when we allow for a epoch dependent $`w(z)`$ ?
* Do the observational constraints agree with each other? In particular, what kind of cosmological models are preferred by SN and WMAP observations individually ?
The last point is important and requires elaboration. Different observational sets are combined together precisely because these observations are sensitive to different combinations of cosmological parameters and facilitate in breaking degeneracies between parameters. If we consider $`\mathrm{\Lambda }`$CDM models then the SN observations, for example, broadly depend on the combination ($`0.85\mathrm{\Omega }_{NR}0.53\mathrm{\Omega }_V`$) 2005A&A…429..807C while WMAP is sensitive to ($`\mathrm{\Omega }_{NR}+\mathrm{\Omega }_V)`$ wmap\_params , a feature which was originally highlighted in the literature as ‘cosmic complementarity’. Therefore, we cannot expect constraints from different observations to agree over the entire parameter space. At the same time, we do not expect models favored by one observation to be ruled out by another when such a divergence is not expected. This divergence may point to some shortcomings in the model, or to systematic errors in observations, or even to an incorrect choice of priors. If all observational sets are consistent then we should be able to derive similar constraints using subsets of observations, even though the final constraints may not be as tight as with the full set of observations.
In order to address the questions listed above in a systematic manner, we proceed in three steps. We choose a ‘base’ reference model with cold dark matter and cosmological constant, with neutrinos contributing a negligible amount to the energy density of the Universe. We assume that the Universe is flat and restrict ourselves to an unbroken power law for the primordial power spectrum of density fluctuations and we assume that the perturbations are adiabatic. Another assumption is that the perturbations in tensor mode are negligible and we take $`r=0`$ 2005PhRvD..71j3515S . We choose this to be our standard model as this can be described by a compact set of parameters.
Next, we generalize from $`\mathrm{\Lambda }`$CDM models ($`w=1`$) to study a wider class of dark energy models with a constant $`w`$ and address the issues listed above. In this case, we also study the effect of perturbations in dark energy. Finally, we generalize to models in which $`w`$ is allowed to vary with $`z`$ in a parameterized form. This approach allows us to delineate changes that come about from choosing a constant $`w1`$, from those allowed by a varying dark energy. We do not impose theoretically motivated constraints on models, e.g. we do not require $`w1`$ as the present work is focused on understanding the nature of models favored by observations.
The paper is organized as follows: In section II we discuss the background cosmological equations followed by a brief review on the various observation used to constrain dark energy equation of state and the observations we concentrate on. The Markov Chain Monte Carlo method is discussed in section III and detailed results are presented in section IV. We conclude with a discussion of the results and future prospects for constraining dark energy models in section V.
## II Theoretical background
### II.1 Cosmological equations
If we assume that each of the constituents of the homogeneous and isotropic universe can be considered to be an ideal fluid, and that the space is flat, the Friedman equations become:
$`\left({\displaystyle \frac{\dot{a}}{a}}\right)^2`$ $`=`$ $`{\displaystyle \frac{8\pi G}{3}}\rho `$ (1)
$`{\displaystyle \frac{\ddot{a}}{a}}`$ $`=`$ $`{\displaystyle \frac{4\pi G}{3}}(\rho +3P)`$ (2)
where $`P`$ is the pressure and $`\rho =\rho _{NR}+\rho _\gamma +\rho _{_{\mathrm{DE}}}`$ with the respective terms denoting energy densities for nonrelativistic matter, for radiation/relativistic matter and for dark energy. Pressure is zero for the non-relativistic component, whereas radiation and relativistic matter have $`P_\gamma =\rho _\gamma /3`$. If the cosmological constant is the source of acceleration then $`\rho _{_{\mathrm{DE}}}=`$ constant and $`P_{_{\mathrm{DE}}}=\rho _{_{\mathrm{DE}}}`$.
An obvious generalization is to consider models with a constant equation of state parameter $`wP/\rho =`$ constant. One can, in fact, further generalize to models with a varying equation of state parameter $`w(z)`$. Since a function is equivalent to an infinite set of numbers (defined e.g. by a Taylor-Laurent series coefficients), it is clearly not possible to constrain the form of an arbitrary function $`w(z)`$ using finite number of observations. One possible way of circumventing this issue is to parameterize the function $`w(z)`$ by a finite number of parameters and try to constrain these parameters by observations. There have been many attempts to describe varying dark energy with different parameterizations constraints\_2 ; 2004ApJ…617L…1B ; 2005MNRAS.356L..11J ; constraints\_10 ; holographic ; Hannestad:2004cb where the functional form of $`w(z)`$ is fixed and the variation is described with a small number of parameters. Observational constraints depend on the specific parameterization chosen, but it should be possible to glean some parameterization independent results from the analysis.
To model varying dark energy we use two parameterizations
$$w(z)=w_0+w^{}(z=0)\frac{z}{(1+z)^p};p=1,2$$
(3)
These are chosen so that, among other things, the high redshift behavior is completely different in these two parameterizations 2005MNRAS.356L..11J . If $`p=1`$ p1 , the asymptotic value $`w(\mathrm{})=w_0+w^{}(z=0)`$ and for $`p=2`$, $`w(\mathrm{})=w_0`$. For both $`p=1,2`$, the present value $`w(0)=w_0`$. Clearly, we must have $`w(z1)1/3`$ for the standard cosmological models with a hot big bang to be valid. This restriction is imposed over and above the priors used in our study.
Integrating $`d(\rho a^3)=w(z)\rho da^3`$, the energy density can be expressed as
$$\frac{\rho _{_{DE}}}{\rho _{_{DE_0}}}=\left(1+z\right)^{3\left(1+w_0+w_0^{}\right)}\mathrm{exp}\left[\frac{3w_0^{})z}{1+z}\right]$$
(4)
for $`p=1`$ (in Eq. 3) and
$$\frac{\rho _{_{DE}}}{\rho _{_{DE_0}}}=\left(1+z\right)^{3\left(1+w_0\right)}\mathrm{exp}\left[\frac{3w_0^{}}{2}\left(\frac{z}{1+z}\right)^2\right]$$
(5)
for $`p=2`$. The allowed range of parameters $`w_0`$ and $`w_0^{}w^{}(z=0)`$ is likely to be different for different $`p`$. However, the allowed variation at low redshifts in $`\rho _{_{DE}}`$ should be similar in both models as observations actually probe the variation of dark energy density. Indeed, in an earlier study 2005MNRAS.356L..11J where we had studied a restricted class of models, we found this to be the case. For example, $`\rho _{_{DE}}`$ can vary by at most a factor two up to $`z=2`$ when both the WMAP and SN data are taken into account 2005MNRAS.356L..11J . This reaffirms the expectation that the results are parameterization independent at some level.
### II.2 Observational constraints
In this subsection, we briefly review potential observational constraints on dark energy and we also summarize previous work in this area.
Constraints on dark energy models essentially arise as follows: To begin with, dark energy affects the rate of expansion of the universe and thus the luminosity distance and also the angular diameter distance. Constraints from observations of high redshift supernovae 1998ApJ…509…74G ; nova\_data1 ; dynamic\_de1 ; 2005A&A…429..807C ; 2004MNRAS.354..275A ; dynamic\_de5 and the location of peaks in the angular power spectrum of CMB anisotropies mainly use this feature 2004PhRvD..70j3523E . The signature of acoustic peaks in correlation function of galaxies also provides a similar geometric constraint 2005astro.ph..1171E . There is also an effect on weak lensing statistics through changes in distance-redshift relation 2002PhRvD..65f3001H ; 2003ApJ…583..566M ; 2003MNRAS.341..251W ; 2003PhRvL..91d1301A ; 2003PhRvL..91n1302J .
Second, the rate of expansion influences the growth of perturbations in the universe and this leads to another set of probes of dark energy 2001PhRvD..64h3501B . Abundance of rich clusters of galaxies, their evolution and the integrated Sachs-Wolfe (ISW) effect belong to this category of constraints, along with constraints from redshift space distortions 1999ApJ…517…40B ; cluster2 ; cluster4 ; 2004astro.ph..8252V ; 2004astro.ph..9207Y ; Calvao:2002zr . All these constraints are sensitive to different aspects of dark energy and a combination of all of these can put tight limits on models. The redshift space distortions are a local effect as these are sensitive to the rate of growth of density perturbations at a given epoch. The abundance of rich clusters of galaxies, the ISW effect and distances are integrated effects in that the effect of dark energy is averaged over a range of redshifts in some sense.
Observations of high redshift supernovae of type Ia provide the most unambiguous evidence for accelerated expansion 1998ApJ…509…74G ; nova\_data1 ; dynamic\_de1 ; 2005A&A…429..807C ; 2004MNRAS.354..275A . Assuming these sources to be standard candles, observations spanning a range of redshifts can be used to study the change in rate of expansion and this imposes direct constraints on the variation of dark energy density. Supernovae have been observed up to a redshift of $`z_{max}1.8`$ and hence can be used to constrain models of dark energy up to this redshift. Constraints from SN observations alone however, permit a large variation in the dark energy parameters 2004MNRAS.354..275A and in particular favor models with $`w<1`$ at the present epoch nova\_data3 ; dynamic\_de1 ; 2005A&A…429..807C .
Baryon oscillations in the matter-radiation fluid prior to decoupling provide a standard scale and the angle at which the acoustic peaks occur in the angular power spectrum of temperature anisotropies in the CMB fixes the distance to the surface of last scattering. This provides a useful constraint on models of dark energy 2004PhRvD..70j3523E as long as dark energy does not affect the dynamics of universe at the time of decoupling of matter and radiation. Unlike supernovae that are observed over a range of redshifts where dark energy is dominant, the surface of last scattering is at $`z1100`$. However, the exquisite quality of CMB anisotropy measurements makes this a very useful constraint and these observations offer a constraint that is different from SN observations 2005MNRAS.356L..11J . Indeed, as we shall see, WMAP and SN observations often favor models that are mutually unacceptable.
Recent detection of the baryon acoustic peak in galaxy correlation function using the luminous red galaxy sample of the SDSS survey has provided an additional handle to constrain cosmological parameters 2005astro.ph..1171E . The geometric constraint from these observations can, in principle, constrain models of dark energy. A measurement of the angular scale corresponding to the peak at different redshifts can indeed be a powerful constraint.
If we consider a given cosmological model that is consistent with observations of CMB anisotropy then the amplitude of fluctuations at the time of decoupling is fixed, and its linearly extrapolated value today can be computed using linear perturbation theory. The abundance of rich clusters of galaxies is related to the amplitude of perturbations in dark matter at a scale of about $`8`$h<sup>-1</sup>Mpc. If we study different models for dark energy while other parameters are not changed, the abundance of rich clusters constrains the net growth of structures between the epoch of decoupling and the present epoch lss\_de .
Redshift space distortions due to kinematics and the Alcock-Paczynski effect are also potential probes of dark energy 1979Natur.281..358A ; 2004astro.ph..9207Y . Ongoing surveys like the SDSS and future surveys will be able to distinguish between different dark energy models through these effects 2003NewAR..47..775W . However, this method does not provide useful constraints at present.
Dependence of the distance-redshift relation and a different rate of growth for perturbations, as well as changes in the matter power spectrum are also reflected in weak lensing statistics. Several studies have been carried out on the potential constraints that can be put on dark energy models from weak lensing observations and their degeneracy with other parameters. These studies indicate that future surveys will be able to put strong constraints on dark energy models 2003ApJ…583..566M ; 2003PhRvL..91d1301A .
Growth of perturbations also leaves a signature in the CMB anisotropy spectrum at large angular scales. The ISW effect leads to an enhancement in the angular power spectrum at these scales. The detailed form of this enhancement depends on the equation of state parameter $`w`$ and its variation. This effect can be detected by cross-correlation of temperature anisotropies with the foreground distribution of matter isw0 . It is difficult to distinguish the ISW effect from the effect of a small but non-zero optical depth $`\tau `$ due to re-ionisation by using only the temperature anisotropies in the CMB, cross-correlation with the matter distribution or polarization anisotropies in the CMB must be used.
Redshift surveys of galaxy clustering do not constrain properties of dark energy directly, however the shape of the power spectrum constrains the combination $`\mathrm{\Gamma }=\mathrm{\Omega }_{NR}h`$ 1992MNRAS.258P…1E . This provides an indirect constraint on dark energy through the well known degeneracy between $`\mathrm{\Omega }_{NR}`$ and $`w_0`$ (e.g. see 2005MNRAS.356L..11J ).
The large number of different observations that can be used to constrain dark energy models is encouraging. Indeed, many attempts have been made to use some of these observations to put constraints on models lss\_de ; 2004PhRvD..69j3501T ; 2005MNRAS.356L..11J ; constraints\_2 ; constraints\_3 .
### II.3 A choice of three observations
In this work, we concentrate on SN, WMAP and cluster abundance observations. We briefly explain the reason for this choice and the kind of constraints one can expect.
The left panel in Fig. 1 shows the degeneracy in $`\mathrm{\Omega }_{NR}`$ and $`w_0`$ (for models with constant $`w`$; i.e., with $`w_0^{}=0`$ in Eq. 3). The figure shows contours of constant luminosity distance $`H_0d_l(z)`$ at $`z=0.17`$ (red/solid curves) and $`z=1.17`$ (blue/dashed curves). The second panel displays the constant luminosity distance contours in the $`w^{}(z=0)w_0`$ plane for $`p=2`$ if $`\mathrm{\Omega }_{NR}=0.3`$. Given that SN observations constrain luminosity distance as a function of redshifts, these figures illustrate the shape of the allowed region that we are likely to get and also demonstrates the degeneracies between different parameters. The third panel shows how the redshift at which the expansion of the universe begins to accelerate depends on the parameters $`w_0`$ and $`w^{}(z=0)`$ for $`p=2`$. This epoch is constrained by SN observations and hence the allowed region in parameter space should lie between contours of this nature. Clearly, regions with a late onset of acceleration (upper right corner) as well as a very early onset of acceleration (lower right corner) will be ruled out by observations of supernovae.
The fourth panel of this figure shows the redshift at which matter and dark energy contribute equally in terms of the energy density of the universe. Structure formation constraints are likely to follow these contours as the rate of growth for density perturbations is significant only in the matter dominated era. Too little structure formation (upper right corner) as well as too much structure formation (lower left corner) are likely to restrict the allowed models along a diagonal (upper left to lower right) in this plane.
Lastly, the location of acoustic peaks in the angular power spectrum of temperature anisotropies in the CMB is the most significant constraint provided by CMB observations. This essentially constrains the distance to the surface of last scattering and hence a suitably defined (see eqn.(8)) mean value ($`w_{eff}`$) for $`w`$. The right panel shows contours of $`w_{eff}`$, which run almost diagonal in this plane. Thus a band of allowed models is the likely outcome of comparison with observations. The contours of $`w_{eff}`$ are the same as contours of equal distance to the surface of last scattering, or the $`l`$ corresponding to the first peak in the angular power spectrum of CMB temperature fluctuations.
SN data provides geometric constraints for dark energy evolution. These constraints are obtained by comparing the predicted luminosity distance to the SN with the observed one. The theoretical model and observations are compared for luminosity measured in magnitudes:
$$m_B(z)=+5log_{10}(D_L)$$
(6)
where $`=M5log_{10}(H_0)`$ and $`D_L=H_0d_L`$, $`M`$ being the absolute magnitude of the object and $`d_L`$ is the luminosity distance
$$d_L=(1+z)a(t_0)r(z);r(z)=c\frac{dz}{(1+z)H(z)}$$
(7)
where $`z`$ is the redshift. This depends on evolution of dark energy through $`H(z)`$. For our analysis we use the combined gold and silver SN data set in nova\_data3 . This data is a collection of supernova observations from nova\_data1 ; 1998ApJ…509…74G and many other sources with $`16`$ supernovae discovered with Hubble space telescope nova\_data3 . The parameter space for comparison of models with SN observations is small and we do a dense sampling of the parameter space.
CMB anisotropies constrain dark energy in two ways, through the distance to the last scattering surface and through the ISW effect. Given that the physics of recombination and evolution of perturbations does not change if $`w(z)`$ remains within some safe limits, any change in the location of peaks will be due to dark energy 2004PhRvD..70j3523E . The angular size $`\theta `$ of the Hubble radius at the time of decoupling can be written as:
$`\theta ^1`$ $`=`$ $`{\displaystyle \frac{H(z)/H_0}{\underset{0}{\overset{z}{}}𝑑y/\left(H(y)/H_0\right)}}`$ (8)
$``$ $`{\displaystyle \frac{\sqrt{\mathrm{\Omega }_{NR}\left(1+z\right)^3}}{\underset{0}{\overset{z}{}}𝑑y/\sqrt{\mathrm{\Omega }_{NR}\left(1+z\right)^3+\varrho ^{DE}(z)/\varrho _0^{DE}}}}`$
$``$ $`{\displaystyle \frac{\sqrt{\mathrm{\Omega }_{NR}\left(1+z\right)^3}}{\underset{0}{\overset{z}{}}𝑑y/\sqrt{\mathrm{\Omega }_{NR}\left(1+z\right)^3+\mathrm{\Omega }_{de}\left(1+z\right)^{3\left(1+w_{eff}\right)}}}}.`$
The second line is obtained as decoupling happens at a redshift where dark energy is not important, and if we ignore the contribution of radiation and relativistic matter; the last equation defines $`w_{eff}`$. Clearly, the value of the integral will be different if we change $`w_0`$, $`w^{}(z=0)`$ and there will also be some dependence on the parameterized form. Location of peaks in the angular power spectrum of the CMB provide a constraint, but this can only constrain $`w_{eff}`$ and not all of $`w_0`$, $`w^{}(z=0)`$ and $`p`$. Therefore if the present value $`w_0<w_{eff}`$ then it is essential that $`w^{}(z=0)>0`$, and similarly if $`w_0>w_{eff}`$ then $`w^{}(z=0)<0`$ is needed to ensure that the integrals match. Specifically, the combination of $`w_0`$, $`w^{}(z=0)`$ and $`p`$ should give us a $`w_{eff}`$ within the allowed range.
In our analysis, we use the angular power spectrum of the CMB temperature anisotropies cmbrev1 ; White and Cohn (2002) as observed by WMAP and these are compared to theoretical predictions using the likelihood program provided by the WMAP team wmap\_lik . We vary the amplitude of the spectrum till we get the best fit with WMAP observations. Note that this is different from the commonly used approach of normalizing the angular power spectrum at $`l=10`$. As we use the entire angular power spectrum for comparison with observations, the impact of ISW effect on the likelihood is relatively small.
It has been pointed out that constraints from structure formation restrict the allowed variation of dark energy in a significant manner lss\_de . We use observed abundance of rich clusters 1999ApJ…517…40B ; cluster2 ; cluster4 to apply this constraint. Since the mass of a typical rich cluster corresponds to the scale of $`8`$h<sup>-1</sup>Mpc, cluster abundance observations therefore constrain $`\sigma _8`$, the rms fluctuations in density contrast at $`8`$h<sup>-1</sup>Mpc. The number density of clusters depends strongly on $`\sigma _8`$ and $`\mathrm{\Omega }_{NR}`$. We use the $`\sigma _8`$ constraints given in 1999ApJ…517…40B from ROSAT deep cluster survey and are given by $`\sigma _8=(0.58\pm 0.1)\times \mathrm{\Omega }_{NR}^{0.47+0.16\mathrm{\Omega }_{NR}}`$ at $`99\%`$ confidence level. The cosmological model should predict $`\sigma _8`$ in the allowed range in order to be consistent with observations.
Recent detection of the baryon acoustic peak using luminous red galaxy sample of the SDSS survey has provided an additional handle to constrain the cosmological parameters 2005astro.ph..1171E . We also used distance scale $`D_V`$ at redshift $`z=0.35`$ introduced in the above reference to further constrain the cosmological models. Here $`D_V(0.35)^3D_M(z)^2cz/H(z)`$ and the observational constraint is $`D_V(0.35)=1370\pm 64Mpc`$ at the $`1\sigma `$ level. This fourth observation does not add significantly to other constraints listed here and we will not describe quantitative results from this constraint here.
Priors used in the present study are listed in Table 1. Apart from these limits on the models studied here, we also assumed that neutrinos are massless and the ratio $`r`$ of tensor to scalar mode is zero. These assumptions are consistent with the known upper bounds, and in any case these do not make any difference to the observations used here as constraints 2005PhRvD..71j3515S .
## III Markov chain monte Carlo method
We compute $`\chi ^2`$ using the routines provided by the WMAP team wmap\_lik . The CMBFAST package 1996ApJ…469..437S is used for computing the theoretical angular power spectrum for a given set of cosmological parameters. We have combined the likelihood program with the CMBFAST code and this required a few minor changes in the CMBFAST driver routine. Given the large number of parameters, the task of finding the minimum $`\chi ^2`$ and mapping its behavior in the entire range of values for parameters is computationally intensive.
We adapt the Metropolis algorithm metropolis (also known in the context of parameter estimation as the Markov Chain Monte Carlo (MCMC) method wmap\_lik ; mcmc2 ) for efficiently mapping regions with low values of $`\chi ^2`$. The algorithm used is as follows:
1. Start from a random point $`𝐫_i`$ in parameter space and compute $`C_l`$ and $`\chi ^2(𝐫_i)`$.
2. Consider a small random displacement $`𝐫_{i+1}=𝐫_i+d𝐫`$ and compute $`\chi ^2(𝐫_{i+1})`$.
3. If $`\chi ^2(𝐫_{i+1})\chi ^2(𝐫_i)`$ then $`ii+1`$. Go to the first step.
4. Else:
* Compute $`\mathrm{\Delta }\chi ^2=\chi ^2(𝐫_{i+1})\chi ^2(𝐫_i)`$ and $`\mathrm{exp}[\alpha \mathrm{\Delta }\chi ^2]`$.
* Compare this with a random number $`0\beta 1`$.
* If $`\beta \mathrm{exp}[\alpha \mathrm{\Delta }\chi ^2]`$ then $`ii+1`$. Go to the first step.
The size of the small displacement $`d𝐫`$ and the parameter $`\alpha `$ are chosen to optimally map the regions of low $`\chi ^2`$. We wish the chain to converge towards the minimum, starting from an arbitrary point, and we also want the Markov chain to map the region in parameter with low $`\chi ^2`$ exhaustively without getting bogged down near the minima. These two conflicting requirements are reconciled by choosing a small but non-zero value of $`\alpha `$. Maximum displacement allowed in one step is small compared with the range of parameters, but not small enough for the chain to get trapped in a small region around the minimum. The optimum values of maximum displacement and $`\alpha `$ are related to each other. We ran several chains with a varying number of points, a typical chain has about $`10^4`$ points. For each set, we have at least $`10^5`$ points (We have done calculations for five sets: cosmological constant, constant $`w`$ with and without perturbations in dark energy, time varying $`w(z)`$ for $`p=1`$ and $`p=2`$. Results presented here required an aggregate CPU time of nearly $`10000`$ hours on $`2.4`$ GHz Xeon CPUs). The convergence criteria for such chains is satisfied for all the sets, and for all the parameters in each set convergence .
We use the MCMC approach only for comparison of models with the CMB data. Observations of cluster abundance are compared with models from the Markov chain run for CMB, after the chain has been run. Comparison of models with observations of high redshift supernovae is done separately. This approach is more conducive to one of the questions that we wish to address, namely, are the observational constraints consistent with each other?
## IV Results
We present results in the form of likelihood functions for various parameters in sets of increasing complexity, starting with the standard $`\mathrm{\Lambda }`$CDM model. Before we proceed with a discussion of results in this form, we discuss a few specific models sampling a few interesting regions of the parameter space in order to develop an intuitive feel for different observational constraints. We call these fiducial or reference models. Along with the fiducial models, we also discuss the best fit models in each set. We find the best fit model for individual observations as well as for the combination of all the observations.
### IV.1 Fiducial Models
#### IV.1.1 The $`\mathrm{\Lambda }`$CDM model
The $`\mathrm{\Lambda }`$CDM model is our ’standard’ model and we begin our discussion with this class of models. Several studies have been carried out to constrain parameters for the $`\mathrm{\Lambda }`$CDM model wmap\_params ; 2004PhRvD..69j3501T . Our results for the $`\mathrm{\Lambda }`$CDM model bring out — among other things — the differences from previous work which arises due to a different method we use here for normalizing power spectra. (See section 2.3 for details.) Differences introduced by priors are also apparent. Our results are as follows:
* For $`\mathrm{\Lambda }`$CDM model, if we consider SN observations alone, we get a best fit at $`\mathrm{\Omega }_{NR}=0.28`$ with a $`\chi _S^2=233.1`$. (We will use subscript $`S`$ for $`\chi ^2`$ from SN analysis and $`W`$ for analysis with WMAP observations.) This model with $`\mathrm{\Omega }_{NR}=0.28`$ is allowed by WMAP observations and has a best fit $`\chi _W^2=974.3`$ for $`\mathrm{\Omega }_B=0.045`$, $`h=0.69`$, $`n=0.95`$ and $`\tau =0.008`$. SN observations do not fix these parameters so we varied the other parameters to get the best fit WMAP model for $`\mathrm{\Omega }_{NR}=0.28`$.
* The model which best fits the WMAP observations has $`\mathrm{\Omega }_B=0.05`$, $`\mathrm{\Omega }_{NR}=0.34`$, $`h=0.66`$, $`n=0.96`$ and $`\tau =0.002`$ with a $`\chi _W^2=972.5`$. The $`\chi ^2`$ value corresponding to SN fit is $`\chi _S^2=239.9`$. In the context of cosmological constant models alone, this model is away from the SN best fit by $`\mathrm{\Delta }\chi _S^2=6.8`$ and is allowed with probability $`0.009`$. In other words, the model most favored by WMAP observations is allowed by the SN observations only with less than one percent probability (We define probability $`𝒫`$ of a given model to be $`1𝒞/100`$, where $`𝒞\%`$ is the confidence limit at which the model is allowed. By using this definition we avoid dilution due to a large parameter space. While the statement about $`\chi ^2`$ is accurate and directly obtainable from the analysis, the conversion of confidence intervals to probabilities has well-known statistical caveats while dealing with multiparameter fits. This should be kept in mind while interpreting our statements about probability with which a model is allowed). In contrast, the model most favored by SN observations is allowed by WMAP observations with a probability $`𝒫=0.945`$.
We now restrict some of the parameters to values favored by other observations e.g. 2001ApJ…553…47F . We fix the baryon density parameter $`\mathrm{\Omega }_B=0.05`$, present day Hubble parameter $`h=0.7`$ and the spectral index $`n=1`$ for these models. Allowing $`\mathrm{\Omega }_{NR}`$ and $`\tau `$ to vary the best fit $`\mathrm{\Lambda }`$CDM model in this restricted class of models is with $`\mathrm{\Omega }_{NR}=0.31`$ and $`\tau =0.14`$ and $`\chi _W^2=974.8`$. This is fairly close to the best fit model found by the WMAP team using a large set of observations wmap\_params . This model is allowed by the rich cluster abundance observations, and also by SN observations ($`\chi _S^2=234.8`$, the corresponding probability being $`𝒫=0.2`$).
Thus convergence between the WMAP and SN observations happens only in a narrow window for flat $`\mathrm{\Lambda }`$CDM models, with the SN constraint being the tighter of the two. It is worth mentioning that in a wider class of models, (obtained by relaxing the prior $`\mathrm{\Omega }_{\mathrm{tot}}=1`$) SN data favors a closed universe with $`\mathrm{\Omega }_{tot}=1.44\pm 0.28`$ and — more importantly — allows the $`\mathrm{\Omega }_{tot}=1`$ models with $`𝒫=0.12`$ 2005A&A…429..807C .
#### IV.1.2 Models with a constant $`w`$
We now allow the dark energy equation of state parameter to have values different from $`w=1`$ but do not allow variation with time. We then find that:
* SN observations favor a model with $`w=1.99`$ and $`\mathrm{\Omega }_{NR}=0.47`$ and with $`\chi _S^2=227.5`$. This is the root cause of the phantom menace. This model is, however, ruled out at a very high probability by WMAP data (with a $`\mathrm{\Delta }\chi _W^2=13.6`$) and is allowed only with $`𝒫=0.022`$.
* WMAP observations favor higher values for $`w`$ (non-phantom models) and the best fit model has $`w=0.72`$. The other cosmological parameters corresponding to this best fit are $`\mathrm{\Omega }_B=0.06`$, $`\mathrm{\Omega }_{NR}=0.34`$, $`h=0.64`$, $`n=0.99`$ and $`\tau =0.19`$ with $`\chi _W^2=971.6`$. SN observations, on the other hand, rule out this model at a very high significance level ($`\mathrm{\Delta }\chi _S^253`$). This is the case when we do not take perturbations in dark energy into account for computing theoretical predictions.
* If we include dark energy perturbations<sup>1</sup><sup>1</sup>1The CMBFAST version ($`\mathrm{4.5.1}`$) being used here incorporates the effects of perturbations in dark energy for constant $`w`$ and non-phantom models with a variable $`w(z)`$. the model which fits best with WMAP observations is same as the best fit for the $`\mathrm{\Lambda }`$CDM models. This model is allowed by SN observations with $`\chi ^2=239.3`$ which is allowed with probability $`0.003`$. Clearly, the discrepancy between the WMAP and SN observations is reduced when we take perturbations in dark energy into account, but SN observations allow the WMAP best fit model only marginally even in this case.
* How well does the $`\mathrm{\Lambda }`$CDM fare when we allow a range of $`w`$ ? The cosmological constant is allowed with $`𝒫=0.063`$ by SN observations in the set of models with constant $`w`$. WMAP observations allow the $`\mathrm{\Lambda }`$CDM models with a very high probability, indeed the best fit model continues to be a $`\mathrm{\Lambda }`$CDM model if perturbations in dark energy are taken into account.
To illustrate these aspects, we consider two fiducial models with $`w=1.5`$ and $`w=0.9`$ allowing $`\mathrm{\Omega }_{NR}`$ and $`\tau `$ to vary. Other parameters are fixed as for the $`\mathrm{\Lambda }`$CDM model (see subsection 1 above). We do not take perturbations in dark energy into account here. The best fit values using WMAP observations for $`w=1.5`$ are $`\mathrm{\Omega }_{NR}=0.38`$ and $`\tau =0.04`$. Generically, models with lower $`w`$ require higher values of $`\mathrm{\Omega }_{NR}`$ which is clear from the form of the curves in Fig.1 (a). This model is allowed by WMAP observations with $`\chi _W^2=978.8`$ as well as by SN observations ($`\chi _S^2=229.9`$). But this model is not allowed by observations of cluster abundance. The value of $`\sigma _8=1.1`$ for this model is higher than the upper limit $`\sigma _8=1.02`$ allowed with $`99\%`$ confidence level for this value of $`\mathrm{\Omega }_{NR}`$.
With $`w=0.9`$, WMAP favors a model with $`\mathrm{\Omega }_{NR}=0.28`$ and $`\tau =0.18`$ ($`\chi _W^2=974.5`$). This model is allowed by cluster abundance observations. The model becomes marginal when SN observations are taken into account with $`\chi _S^2=237.9`$, as compared to a minimum of $`\chi _S^2=227.5`$ for models with constant $`w`$.
Once again, we find that WMAP observations and SN observations favor different regions in the parameter space. Generalizing to the class of models with a constant $`w`$ from $`w=1`$, we find that SN data has a distinct preference for $`w<1`$ and the $`\mathrm{\Lambda }`$CDM model is allowed only marginally. WMAP data also shows a mild preference for models with $`w1`$, though it continues to allow the $`\mathrm{\Lambda }`$CDM model. If perturbations in dark energy are taken into account then the $`\mathrm{\Lambda }`$CDM model continues to be the most favored model for WMAP observations.
#### IV.1.3 Models with varying $`w(z)`$
Next, we allow the dark energy equation of state parameter to vary. As it is not possible to take the effect of non-adiabatic perturbations into account for these models and as we do not wish to add another parameter, we work with models without any perturbations in dark energy perturb3 . It is clear that this will introduce a slight bias towards models with $`w>1`$ but will not change anything else. For the discussion of fiducial models, we choose the parameterization with $`p=2`$ in Eq.(3).
* Parameter values for the best fit model with SN data are $`w_0=1.95`$, $`w^{}(z=0)=4.5`$ and $`\mathrm{\Omega }_{NR}=0.498`$ with a $`\chi _S^2=227.4`$. As with the constant $`w`$ models, SN data favors large negative values of the equation of state parameter. SN observations do not favor models with varying $`w(z)`$ over models with $`w1`$, since the change in best fit $`\chi ^2`$ is only $`0.1`$ when the additional parameters are added. WMAP observations allow this model with $`𝒫=0.035`$ ($`\chi _W^2=987.6`$) if we have $`\mathrm{\Omega }_B=0.056`$, $`h=0.662`$, $`n=1.04`$ and $`\tau =0.002`$. However, this model is ruled out by cluster abundance observations as $`\sigma _8=1.25`$ for this model is higher than the allowed range at $`99\%`$ confidence level.
* The best fit model for WMAP data has $`w_0=1.48`$, $`w^{}(z=0)=3.86`$, $`\mathrm{\Omega }_B=0.05`$, $`\mathrm{\Omega }_{NR}=0.24`$, $`h=0.73`$, $`n=1.1`$ and $`\tau =0.35`$ with $`\chi _W^2=970.9`$. This model is allowed by SN observations with $`𝒫=0.33`$ ($`\mathrm{\Delta }\chi ^2=4.63`$).
The tension between the WMAP and SN observations is less serious for models with varying $`w`$ than for models with a constant $`w`$. Part of the reason is that with a larger number of parameters, the same $`\mathrm{\Delta }\chi _S^2`$ gives us a larger probability for a given model.
As in the previous cases, let us choose fiducial models by restricting some of the parameters. We consider a model with dark energy parameter values $`w_0=1.5`$ and $`w^{}(z=0)=1.0`$. The WMAP best fit with these parameters is with $`\mathrm{\Omega }_{NR}=0.37`$ and $`\tau =0.0005`$ with $`\chi _W^2=976.5`$ . The change in acceptance level of this model is due to our restricting the values of the spectral index and the Hubble parameter. This model has $`\chi _S^2=229.3`$ and is allowed by all the three observations used here.
If we move closer to the model most favored by SN observations, $`w_0=1.5`$ and $`w^{}(z=0)=5.0`$, then we find that WMAP data favors $`\mathrm{\Omega }_{NR}=0.42`$ and $`\tau =0.0`$. Even this model is outside the range allowed by WMAP observations at $`68\%`$ confidence limit and is allowed at $`𝒫=0.12`$. The model fares similarly with SN observations, i.e., it is allowed with $`𝒫=0.12`$.
We wrap up with a discussion of a model with $`w_0>1`$. We consider $`w_0=0.9`$ and $`w^{}(z=0)=3.0`$. WMAP observations allow this model with $`\chi _W^2=977.4`$ for $`\mathrm{\Omega }_{NR}=0.365`$ and $`\tau =0.0013`$. SN observations allow this model within the $`95\%`$ confidence limit with a probability $`𝒫=0.15`$. This model is also allowed by observations of cluster abundance.
Within the context of models with a variable $`w(z)`$, the $`\mathrm{\Lambda }`$CDM model is allowed by WMAP observations ($`\mathrm{\Delta }\chi _W^2=1.6`$) as well as by supernova observations ($`\mathrm{\Delta }\chi _S^2=5.73`$). SN observations clearly favor models other than the $`\mathrm{\Lambda }`$CDM model in context of $`\mathrm{\Omega }_{tot}=1`$, while no such preference is seen for WMAP observations. However, models favored by SN observations require $`\mathrm{\Omega }_{NR}`$ to be much larger than the values favored by observations of rich clusters crisis2 .
In summary, there is significant tension between the sets of observations we are studying and this tension does not reduce when the parameter space is enlarged. We see that best fit model for one set of observation is often ruled out by another set with a high level of significance. There is an overlap of region allowed with $`95\%`$ confidence limit in all cases and within $`68\%`$ confidence limit in some cases; but this does not take away the significance of the differences which the above analysis has thrown up. The results presented here are summarized in Table 2.
### IV.2 Results in Detail
In this section we outline the detailed results of our analysis. We marginalize the results in the multi-dimensional parameter space to derive likelihood function for each parameter of interest. The likelihood function is sensitive to the bin-size used for the given parameter and tends to be somewhat noisy. It is customary to smooth the likelihood function with a Gaussian filter in order to remove noise but the results are sensitive to the width of the filter used. Therefore we choose to plot the cumulative likelihood as this is insensitive to binning and smoothing is not required. We define the cumulative likelihood as follows:
$$I(x)=\frac{\underset{x_{min}}{\overset{x}{}}(y)𝑑y}{\underset{x_{min}}{\overset{x_{max}}{}}(y)𝑑y}$$
(9)
where $`x`$ is the parameter we are interested in and it has a range $`x_{min}xx_{max}`$, $`(x)`$ is the likelihood obtained by marginalizing over other parameters. The central value for the given variable is thus $`x_c`$, with $`I(x_c)=0.5`$. This can be, and in general it is different from the value of the parameter for maximum likelihood. $`I(x)`$ is like the cumulative probability function for the parameter $`x`$. Parameter values for which $`I(x)=0.025`$ and $`I(x)=0.975`$ define the range allowed at $`95\%`$ confidence limit, i.e. the probability that the variable lies within this range is $`0.95`$. We mark this limit by two vertical lines in all the likelihood plots.
### IV.3 The $`\mathrm{\Lambda }`$CDM model
Fig. 2 shows marginalized cumulative likelihood $`I(x)`$ for the $`\mathrm{\Lambda }`$CDM model, different frames correspond to the different parameters we have considered here. Curves are shown for $`\mathrm{\Omega }_{NR}`$, the shape parameter $`\mathrm{\Gamma }=\mathrm{\Omega }_{NR}h`$, Hubble parameter $`h=H_0/100kms^1Mpc^1`$, $`\mathrm{\Omega }_bh^2`$, spectral index $`n`$ and optical depth to the redshift of reionisation $`\tau `$. Red/dot-dashed curves show constraints from WMAP observations of CMB temperature anisotropies, dark-green/dotted curve shows constraints from WMAP observations and abundance of rich clusters, black/solid curve shows constraints from a combined analysis of WMAP observations, cluster abundance and high redshift supernovae. The blue/dashed curve shows the constraints from SN observations alone, this has been plotted only for $`\mathrm{\Omega }_{NR}`$ as SN data does not constrain other parameters directly. Combined analysis for other parameters does have an input from supernova observations as many models allowed by WMAP and cluster abundance observations are ruled out by SN observations. Vertical lines mark the $`95\%`$ confidence limit when all the observations are used together. CMB observations allow considerable range for each of these parameters whereas observations of high redshift supernovae provide tight constraints. The reason why SN observations provide a tight constraint is clear from the discussion in the previous section, namely, SN observations favor models with $`\mathrm{\Omega }_{tot}1`$. Within the context of models with $`\mathrm{\Omega }_{tot}=1`$, SN observations favor models with $`w1`$. The $`\mathrm{\Lambda }`$CDM model is only marginally allowed by SN observations within both the sets; flat $`\mathrm{\Lambda }`$CDM models are allowed with $`𝒫=0.12`$ 2005A&A…429..807C and within flat models the $`\mathrm{\Lambda }`$CDM model is allowed with $`𝒫=0.06`$.
The allowed range for all the parameters within $`68\%`$ confidence limit is given in Table 3 and range allowed in $`95\%`$ confidence limit is given in Table 4. Table 3 clearly illustrates the discrepancy between SN observations and WMAP observations. The values are comparable with those obtained in other analyses wmap\_params ; 2004PhRvD..69j3501T . The range of values for the shape parameter $`\mathrm{\Gamma }=\mathrm{\Omega }_{NR}h`$ favored by these observations is consistent with values obtained from galaxy surveys 2004PhRvD..69j3501T ; sdss ; 2005astro.ph..1174C . The allowed range of values for $`h`$ are in agreement with direct determination 2001ApJ…553…47F .
Constraints from abundance of rich clusters do not make a significant impact on the likelihood function of individual parameters even though these constraints reject a significant fraction of models allowed by WMAP observations. To illustrate this, we have plotted points in the Markov chain for $`\mathrm{\Lambda }`$CDM models that are allowed by WMAP observations in a few projections in the parameter space (Fig. 3). Also shown in the same plots are points allowed by abundance of rich clusters of galaxies. This clearly shows that the region in parameter space allowed by cluster abundance is distinctly smaller than that allowed by CMB observations. Well known degeneracies between parameters are also highlighted by this figure, e.g., there is a clear degeneracy between $`\mathrm{\Omega }_{NR}`$ and $`h`$. An important point to note is that if the preferred range for $`\sigma _8`$ were to be towards larger values than taken here cluster4 ; 2003ApJ…591..599S then the cluster abundance constraint will favor models with larger $`\mathrm{\Omega }_{NR}`$. There is also a related shift towards $`w<1`$.
### IV.4 Models with a constant $`w`$
We now consider models with a constant equation of state parameter $`w`$. Introduction of this additional parameter changes the relative effectiveness of different observations in constraining cosmological parameters. Observations of high redshift supernovae of type Ia constrain the key cosmological parameters much more strongly than CMB observations for $`\mathrm{\Lambda }`$CDM models. This is no longer the case once we introduce $`w`$ as a parameter. The main reason for this is the degeneracy between $`w`$ and $`\mathrm{\Omega }_{NR}`$.
In models with $`w1`$, it is necessary to take perturbations in the dark energy component into account 2004PhRvD..69h3503B ; perturb3 ; perturb4 . Here, we study models with and without perturbations in dark energy in order to illustrate the role played by these perturbations and to study how strongly these influence determination of cosmological parameters. Fig. 4 shows the likelihood $`I`$ for parameters $`\mathrm{\Omega }_{NR}`$, $`w`$, $`h`$, $`\mathrm{\Omega }_bh^2`$, $`n`$ and $`\tau `$ in models with and without perturbations in the dark energy component. Models with a larger $`\mathrm{\Omega }_{NR}`$ and smaller $`w`$ are better fits to SN observations, whereas CMB observations prefer models with smaller $`\mathrm{\Omega }_{NR}`$ and a larger equation of state parameter $`w`$. The combination of these observations and abundance of rich clusters constrains both the parameters to a fairly narrow range, much narrower than is allowed by SN observations alone. Models with $`w>1`$ fare badly with CMB observations when perturbations in the dark energy component are taken into account. For models with a constant $`w`$, observations allow higher values of $`\mathrm{\Omega }_bh^2`$ than for $`\mathrm{\Lambda }`$CDM models. The allowed range in $`\tau `$ is smaller then the case with no dark energy perturbations whereas the ranges for $`n`$ are similar in both.
We find that the range of $`w`$ allowed at the $`95\%`$ confidence limit is smaller when perturbations in dark energy are allowed. For other parameters, the allowed range is similar. In other words, if we ignore perturbations in dark energy we can still make a reasonable estimate of the range of parameters allowed by observations. This fact is of immense use when we work with models that have a varying equation of state parameter $`w`$. In order to take full effects of dark energy perturbations in these models it is essential to know full details of the model 2004PhRvD..69h3503B ; perturb3 . We cannot include the effect of non-adiabatic perturbations in a model independent study of dark energy models with a varying $`w`$. Given the fact that ignoring perturbations in dark energy does not lead to an incorrect estimate of the range of parameters (except for $`w`$) that is allowed, we can safely proceed with our analysis without taking perturbations in dark energy into account. As regards the equation of state parameter $`w`$, ignoring perturbations tends to allow $`w>1`$ models with a larger probability and we should keep this in mind while interpreting results.
### IV.5 Models with varying $`w(z)`$
We now proceed to the case of varying $`w`$, we use two parameterizations given in 2005MNRAS.356L..11J . The first of these, corresponding to $`p=1`$, is a Taylor series expansion for $`w`$ in scale factor and this is a very commonly used parameterization. The variation of the equation of state parameter is monotonic in this case and rapid increase of $`w`$ at low redshifts cannot be allowed as it will lead to $`w1/3`$ at high redshifts. The parameterization with $`p=2`$ avoids this problem to some extent as the value of $`w`$ at very high redshifts is the same as the present value, but there can be a large deviation from this at low redshifts with the deviation peaking at $`z=1`$. Rapid variation at low red-shifts has been reported 2004MNRAS.354..275A on the basis of SN observations nova\_data1 ; nova\_data3 , and even though these conclusions have been contested dynamic\_de6 it is useful to check if the larger set of observations support a rapid variation of $`w`$ at low redshifts.
We have already seen in the discussion of fiducial models, supernova observations do not distinguish between models with a constant $`w`$ and models with a variable equation of state parameter. Also, WMAP observations do not differentiate between the three classes of models being studied here in a statistically significant manner.
Fig. 5 shows the allowed ranges of parameters $`w_0`$ and $`w^{}(z=0)`$ for $`p=1`$ (left panels) and $`p=2`$ (right panels). In both the parameterizations, the preferred values for $`w_0`$ with supernova observations are for phantom models ($`w_01`$) and a tendency for a larger $`w`$ at intermediate redshift, i.e., $`w_0^{}0`$ though supernova observations do not provide a clear constraint on this parameter. This is perhaps related to the fact that there is a strong degeneracy in $`\mathrm{\Omega }_{NR}`$ and $`w_0`$. CMB observations and abundance of rich clusters of galaxies allow models around the $`\mathrm{\Lambda }`$CDM model, which is fairly close to the centre of the allowed region. A combination of these three observations rejects models with $`w_01`$ due to CMB constraints and $`w_0>1`$ due to SN constraints. Even though all the observations allow $`w_0^{}=0`$, the combination of these observations does not favor such models. This implies that the overlap of allowed regions for the three observations is stronger for models with $`w_0^{}0`$, it is clear from table 3 that there is no overlap between SN and WMAP at $`68\%`$ confidence limit for constant $`w`$. The $`\mathrm{\Lambda }`$CDM model, i.e., $`w_0=1`$ and $`w_0^{}=0`$ is a marginally allowed model for both the parameterizations.
To understand the nature of constraint from CMB observations 2003MNRAS.343..533D , we computed the likelihood of $`w_{eff}`$ (as defined in eqn.(8)) for models by comparing these with WMAP data. This is then compared with the likelihood for $`w`$ in models with a constant equation of state parameter. We have plotted this in Fig. 6 for $`p=1`$ as well as $`p=2`$, along with the curve for constant $`w`$ (with no perturbations in dark energy). All the three curves show very similar behavior and the $`95\%`$ confidence limit is identical for all three ($`1.5w_{eff}0.6`$). This also shows that the CMB observations primarily provide a constraint for $`w_{eff}`$. Given that adding perturbations reduces the likelihood for models with $`w>1`$, it is likely that detailed analysis of a model with perturbations in dark energy taken into account will limit the range for $`w_{eff}`$ in this region.
Lastly, we study the effect of varying dark energy on other parameters. The specific question we wish to address is, how the allowed ranges for these parameters change if we allow variation of dark energy. Fig. 7 shows likelihood for the parameters studied here. We have plotted the likelihood using all three observations for the $`\mathrm{\Lambda }`$CDM models, constant $`w`$ models as well as for $`p=1`$ and $`p=2`$. For most parameters, the effect of $`w1`$ and varying dark energy is to increase the range of allowed values. This increase in the allowed range is sometimes accompanied by a shift, e.g. for $`h`$ where varying dark energy models fit observations better with smaller values as compared to the $`\mathrm{\Lambda }`$CDM model as well as models with constant $`w`$. This shift is primarily due to models with $`w>1`$ and this point has been noted in other analyses as well wmap\_params . If this is the case then including perturbations in dark energy may well remove this shift.
Similarly, larger values of spectral index $`n`$ and optical depth to the epoch of reionisation $`\tau `$ fit observations better. As these parameters can be constrained using other observations, we may be able to restrict models with varying $`w`$ by constraining the values of these parameters. For example, polarization anisotropies in the CMB can be used to constrain $`\tau `$ cmbrev1 ; 2003ApJ…583…24K . We find that the presently available information from WMAP about polarization anisotropies does not lead to a significantly improved constraint on the parameter $`\tau `$.
### IV.6 Evolution of dark energy
We now summarize our results for the allowed variation of dark energy, once all three observational constraints have been taken into account. In the left panel of Fig. 8, we have plotted the cumulative likelihood for the equation of state parameter $`w`$ at redshift $`z=1`$. Here we have used models which lie within the range $`1.1<w_0<0.9`$. The upper and lower panels correspond to parameterizations with $`p=1`$ and $`p=2`$ respectively. The allowed range of variation in the equation of state by supernova observations is much larger than that allowed by WMAP results. This, again, is a reflection of the strong preference of supernova observations for $`w1`$ and of the large parameter space allowed by SN data. (Our result that SN data prefers $`w1`$ with large variation is consistent with previous published analysis e.g., in 2004ApJ…617L…1B .) The range of values in both the parameterizations are similar for this subset of models. In the middle panel we have shown the likelihood for variation in the equation of state parameter from the present to its value at redshift $`z=1`$. The allowed ranges of variation in dark energy equations of state are different for these two parameterizations. In fact, the constant dark energy equation of state is ruled out at $`95\%`$ confidence level for $`p=2`$ (that is, the probability of occurrence is less than $`0.05`$) when all the constraints are taken into account even though each observational constraint allows such models individually. Clearly, the models with constant $`w`$ allowed by each of these observations are ruled out by other observations. The $`\mathrm{\Lambda }`$CDM model is allowed for $`p=2`$ at $`77\%`$ C.L. (with probability of $`23\%`$) by SN observations and by $`0.9\%`$ C.L. ($`𝒫=0.991`$) by WMAP observations.
In the right panel, we have shown the ratio of dark energy density at $`z=1`$ and the present value. The variation allowed by SN observations is very large, whereas WMAP limits the variation to within a factor $`2.5`$ at $`95\%`$ confidence limits. This drives the joint analysis to restrict variation even further. That WMAP observations provide a much tighter constraint on the equation of state as compared to SN observations was earlier shown in 2005MNRAS.356L..11J .
In Fig. 9 we show the allowed range of variation of dark energy as function of redshift for $`w=constant`$ models, with and without perturbations at $`68\%`$, $`95\%`$ and $`99\%`$ confidence levels. The figure shows the disparity in allowed range by SN observations and WMAP observations at $`68\%`$ confidence level. Allowing perturbations in dark energy gives a similar range as compared to the case where dark energy perturbations are absent. In Fig. 10 we plot this range for varying dark energy models, top panel for models with $`p=1`$ and lower panel for $`p=2`$. As mentioned earlier (see also 2005MNRAS.356L..11J ), SN observations allow a much wider range in change of dark energy density with redshift. The variation allowed by WMAP is smaller in all cases except constant $`w`$. The combination of the three constraints allows very little variation, with maximum allowed variation in dark energy density being by a factor $`5`$ up to $`z=2`$ at $`68\%`$ confidence limit. The allowed variation in dark energy density is similar in both the cases, indicating that the constraints on this quantity are parameterization independent to a large extent.
Finally, we would like to make some comments regarding the fact that WMAP constrains the evolution of dark energy more effectively than SN. This arises essentially from the constraint on the angular diameter distance to the last scattering surface, or — equivalently — the effective equation of state parameter $`w_{eff}`$. (The Integrated Sachs-Wolfe effect and the contribution of other parameters turns out to be less important.) To illustrate this point, we have compared the constraints on dark energy density for $`p=2`$ with those implied by constraints on $`w_{eff}`$. The second figure in the lower panel in Fig. 10 shows the allowed range in evolution of dark energy density allowed by WMAP data alone for $`p=2`$. We have also plotted the allowed range for dark energy density as a function of redshift if $`1.6w_{eff}0.6`$, by thick black lines. This is derived by using all $`w_0`$ and $`w_0^{}`$ that lead to $`w_{eff}`$ in the range given above, and computing the highest and lowest dark energy density amongst this set of models at each redshift. We allow $`w_0`$ and $`w_0^{}`$ to vary in the range specified in the priors. The region allowed by the range in $`w_{eff}`$ and that directly obtained from all allowed models is similar, with the latter allowing larger variation for phantom models. This reiterates our claim that the main constraint from WMAP data on dark energy parameters is on the value of $`w_{eff}`$ at the last scattering surface. We believe that the larger range allowed at $`95\%`$ confidence limit is due mainly to the ISW effect.
## V Conclusions
In this paper we presented a detailed analysis of constraints on cosmological parameters from different observations. In particular we focussed on constraints on dark energy equation of state, its present value and the allowed range of variation in it.
It is demonstrated that the allowed range for the equation of state parameter $`w`$ is smaller if dark energy is allowed to cluster. Including perturbations mainly affects models with $`w>1`$.
We find that WMAP observations do not distinguish between the $`\mathrm{\Lambda }`$CDM model, models with a constant equations of state parameter $`w`$ and models with a variable $`w`$, the change in $`\chi _W^2`$ for best fit models is less than $`3`$ even as the number of parameters is increased by three. WMAP allows only a modest variation in energy density of dark energy, with maximum variation being less than a factor of three in $`99\%`$ confidence limit up to $`z=1`$. We infer that the main constraint from WMAP observations is for the derived quantity $`w_{eff}`$, essentially representing the distance to the last scattering surface.
SN observations favour models with $`w<1`$ and $`\mathrm{\Omega }_{nr}>0.4`$. A corollory is that if we restrict to models with $`w1`$ then the $`\mathrm{\Lambda }`$CDM model is the most favoured model. Without this restriction the $`\mathrm{\Lambda }`$CDM model is allowed only marginally by the combination of observations used here, this is driven mainly by SN observations.
Allowing variation in dark energy has an impact on other cosmological parameters as the allowed range for many of these parameters becomes larger. Conversely, better measurements of these parameters will allow us to constrain models of dark energy.
We find significant tension between different observations. Our key conclusions in this regard may be summarised as follows:
* SN observations favor models with large $`\mathrm{\Omega }_{NR}`$ and $`w1`$. Indeed, the best fit model is at the edge of our priors.
* Enlarging priors to $`0.1\mathrm{\Omega }_{NR}0.6`$, $`0.3w3.0`$ does not lead to a better fit model for SN observations, indicating that our default priors are sufficiently wide for joint estimation of parameters. This is because $`w1`$ is rejected by WMAP observations.
* WMAP observations favor models with $`w1`$ with a marginal preference for $`w>1`$. Including perturbations in dark energy removes this marginal preference as well.
* For constant $`w`$ models and models with variable $`w`$, the best fit model of each observation is ruled out by the other observations at a high significance level. As an example, the model that best fits the WMAP observations is completely ruled out by SN observations ($`\mathrm{\Delta }\chi _S^2=53`$). The problem is slightly less serious if perturbations in dark energy are taken into account ($`\mathrm{\Delta }\chi _S^2=12.5`$).
* There is overlap of allowed regions at $`95\%`$ (or better) by these observations, though there is little overlap of allowed regions at $`68\%`$ confidence limit (see table 3). It can, of course, be argued that situation is not alarming given that there is an overlap of allowed regions in parameter space at $`95\%`$. But we find this offset noteworthy.
* Using larger values for $`\sigma _8`$, as indicated by some recent analyses cluster4 ; 2003ApJ…591..599S favors models with a slightly larger $`\mathrm{\Omega }_{NR}`$ and slightly lower $`w_{eff}`$.
* Given that the preference of individual observations for different types of models is not understood, and the fact that the best fit model of one is ruled out by the other, it is necessary to use a combination of observations for reliable constraints on models of dark energy. Use of either one of the observations is likely to mislead.
* Our conclusions are not sensitive to priors used for parameters other than $`\mathrm{\Omega }_{NR}`$. Limiting priors for matter density to $`0.1\mathrm{\Omega }_{NR}0.3`$ enhances the overlap between SN and WMAP observations and removes the tension between SN and WMAP observations for constant $`w`$ models.
* If we repeat the analysis of models with variable $`w`$ with this restricted priors then we find that SN observations strongly favor a variable $`w`$ as compared to constant $`w`$. There is no significant tension between SN and WMAP models in this class of models even for the wider priors, but this is mainly due to a larger number of parameters.
* Using only the Gold data set for supernovae instead of the Gold $`+`$ Silver used here reduces the tension between the WMAP and supernova observations by a marginal amount.
Given the points noted here regarding tension between different observations, it is important that some effort is made to look for systematic effects in observations as well as in analysis of observations. We have tested our analysis for systematic effects by varying priors and our findings appear to be independent of the chosen priors, the only instance of change in results is mentioned above. Since the SN data set which is used by most people (including in this work) arises from different sources, one needs to be careful regarding hidden systematics (see e.g., the discussion in saul ). When larger, homogeneous SN datasets are available in future (like for example, from SNLS), it is likely that the tension between the SN observations and WMAP results disappear. If it does not, and the agreement continues to exist only at 3-sigma level, there is some cause for concern.
## Acknowledgements
JSB thanks Stefano Borgani and U. Seljak and TP thanks S.Perlmutter for useful comments. The numerical work in this paper was done using cluster computing facilities at the Harish-Chandra Research Institute (http://cluster.mri.ernet.in/). This research has made use of NASA’s Astrophysics Data System. |
warning/0506/astro-ph0506682.html | ar5iv | text | # The XMM-Newton Detection of Diffuse Inverse Compton X-rays from Lobes of the FR-II Radio Galaxy 3C 98
## 1 Introduction
Since the pioneering discoveries with ASCA and ROSAT (e.g. Feigelson et al. 1995; Kaneda et al. 1995; Tashiro et al. 1998), increasing numbers of diffuse hard X-ray emission have recently been detected from lobes in radio galaxies and quasars (e.g; Brunetti et al. 2001; Isobe et al. 2002; Hardcastle et al. 2002). These diffuse X-ray photons are thought to arise via inverse Compton (IC) scattering by electrons in the lobes, which also radiate synchrotron radio emission. In most of these sources, the seed photons of the IC scattering in the lobes are provided by the cosmic microwave background (CMB) radiation (Harris & Grindlay 1979), or infra-red (IR) photons from the active nucleus (Brunetti et al. 1999). By comparing the IC X-ray and synchrotron radio intensities, we can independently deduce energy densities of the electrons and magnetic fields in the lobes, $`u_\mathrm{e}`$ and $`u_\mathrm{m}`$ respectively, on condition that the seed-photon energy density is known. Therefore, these IC X-rays from radio lobes provide a valuable tool to directly measure the energetics in lobes, and infer those in jets, indirectly.
The observed intensities of the IC X-rays from the lobes frequently infer that $`u_\mathrm{e}`$ significantly dominate $`u_\mathrm{m}`$ (e.g. Isobe 2002). The possibilities of magnetic fields smaller than those estimated under the equipartition condition by some factors are also reported in the knots and hot spots of radio galaxies and quasars (e.g., Hardcastle et al., 2004; Kataoka & Stawaltz, 2005). However, the sample of the IC X-ray lobes still remains small, due mainly to the limited sensitivity of the X-ray instruments. Obviously, XMM-Newton has a great potential in this research subject owing to its large effective area and its moderately good angular resolution. Actually, a few results have already been reported based on the XMM-Newton data (Grandi et al. 2003; Croston et al. 2004).
We have conducted an XMM-Newton observation of a lobe-dominant radio galaxy 3C 98, and successfully detected the diffuse X-ray emission associated with its lobes. As reported in the present paper, we find a clear dominance of $`u_\mathrm{e}`$ over $`u_\mathrm{m}`$, in the lobes of this object. Throughout the present paper, we assume a Hubble constant of $`H_0=75`$ km sec<sup>-1</sup> Mpc<sup>-1</sup> and a deceleration parameter of $`q_0=0.5`$.
## 2 Observation
Located at a redshift of $`z=0.0306`$ (33.7 kpc/1 arcmin; Schmidt, 1965), 3C 98 is a radio galaxy with an elliptical host (Zirbel, 1996). It has a moderately high integrated radio flux of 10.3 Jy at 1.4 GHz (Laing, & Peacock, 1980), and shows a synchrotron spectral energy index of $`\alpha _\mathrm{R}=0.78`$ between 178 and 750 MHz (Laing, Riley, & Longair, 1983). Its VLA image (e.g., Leahy et al., 1997) reveals a classical FR II (Fanaroff, & Riley, 1974) double-lobe morphology, with a relatively flat surface brightness distribution and a total angular extent of $`5^{}\times 2^{}`$. It dose not lie in a rich cluster environment, and is hence relatively free from hot thermal X-ray emission. Moreover, a sign of extended X-ray emission on a scale of tens of arcseconds was reported based on an ROSAT observation (Hardcastle & Worral 1999). All these facts make this object very suitable to our search for the lobe IC X-rays.
Our 30 ksec XMM-Newton exposure on 3C 98 was performed on 2002 September 7 (hereafter, Exp1). The EPIC MOS and PN CCD cameras on-board XMM-Newton were both operated in the nominal full frame mode. Because some normal stars in the field of view are brighter than the optical magnitude of 10, we inevitably adopted the medium and thick optical blocking filters for MOS and PN, respectively. Although these filters significantly reduce the efficiency below 1 keV, they little affect our study on hard-spectrum IC emission. Due to an operational failure, however, only about 60% of the requested exposure was completed with MOS1, and no MOS2 data were taken in the observation. Therefore, another 10 ksec exposure was approved and conducted on 2003 February 5 (Exp2), about 5 month after Exp1. The log of these two exposures are summarized in Table 1.
We reduced the data using version 5.4.1 of the Science Analysis System (SAS). All of the obtained data were reprocessed, base on the latest Current Calibration Files (CCF), just after the data of Exp2 were delivered to us (2003 April). We selected only those events with PATTERN $``$ 12 (single, double triple and quadruple events) for MOS, and PATTERN $``$ 4 (single and double events.) for PN. In addition, we also imposed the criterion of #XMMEA\_EM and #XMMEA\_EP for MOS and PN respectively, in order to filter out artifact events. As is frequently reported, the MOS and the PN backgrounds were both found to be highly variable, and so-called background flares significantly contaminated the data in both exposures. Accordingly, we have selected the data obtained only when the 0.2 – 15 keV count rate was lower than 2.5 and 20 counts per sec for MOS and PN respectively, when accumulated within a source free region which is about 96% of the field of view. As a results, the remained good time interval became slightly short as summarized in Table 1. In addition to the above criteria; we adopt only those events with FLAG == 0, in the analysis below; this strictly rejects events next to the edges of the CCD chips, and next to bad pixels and columns.
## 3 Results
### 3.1 X-ray Image
Figure 1 shows the 0.2–12 keV EPIC image around 3C 98. The image is smoothed with a two dimensional Gaussian function of $`\sigma =5^{\prime \prime }`$, after all the MOS and PN data obtained in both exposures are summed up. The 4.86 GHz VLA image (Bridle; unpublished) is superposed with contours. The bright X-ray source located near the center of the image coincides with the optical host galaxy, and also with the radio nucleus within $`1^{\prime \prime }.5`$, which is reasonable for the astrometric accuracy of the EPIC. In addition, we notice faint diffuse X-ray emission associated with the radio lobes, especially the northern one.
In order to visualize more clearly the diffuse X-ray emission, we subtracted the point spread function (PSF) corresponding to the central point-like source at the host galaxy, based on the CCF data base. For this purpose, we used only the MOS data, because gaps of the PN CCD chips intersect the lobes. The in-orbit calibration of the MOS PSF is shown in detail, by Ghizzardi (2001), at least for an on-axis point source. The central point source is not so bright that the shape of the PSF is not distorted by significant CCD event pile-ups. We adequately took into account the energy dependence of the PSF, referring to the spectral information of the source, shown in §3.2. Before the subtraction, the PSF image is normalized to the observed one, by using the event counts within a circle of 15 arcsec radius, around the X-ray peak.
Figure 2 shows the MOS image of 3C 98, from which the PSF for the central source is subtracted, in 0.2 – 12 keV. This image is heavily smoothed with a Gaussian function of $`\sigma =20^{\prime \prime }`$. The diffuse X-ray emission is significantly detected, in the figure, with some X-ray point sources. Moreover, its association with the lobes of 3C 98 is found to be clear, and its overall angular extent almost coincides with that of the radio lobes.
We extracted linear X-ray brightness profiles parallel and orthogonal to the radio axis of 3C 98. in order to evaluate the significance of the diffuse X-ray emission. The integration strip for the parallel profile (P<sub>0</sub> in Figure 1) has an width of $`2^{}.5`$, to contain all the radio structure of 3C 98. The orthogonal profile utilizes a narrower strip, O<sub>0</sub>, of a $`1^{}.25`$ width, for minimizing contamination from the possible diffuse X-ray emission associated with the lobes. We estimated the X-ray background for the parallel profile from the neighboring regions P<sub>b1</sub> and P<sub>b2</sub>, and that for the orthogonal profile from O<sub>b1</sub> and O<sub>b2</sub>. All of these regions are shown in Figure 1 with dotted lines.
The X-ray profiles thus obtained are shown in Figure 3, together with the radio ones accumulated within the same sky regions. In the parallel profile, an excess over the PSF is clearly found with a significance of $`4\sigma `$, accumulated between $`0^{}`$ and $`3^{}`$, and $`2\sigma `$ between $`3^{}`$ and $`0^{}`$ from the X-ray peak. On the other hand, the X-ray profile perpendicular to the lobes exhibits no significant excess over the PSF ($`1\sigma `$, between $`3^{}`$ and $`3^{}`$) implying that the central source is point-like, within the MOS angular resolution. Therefore, we confirm that the EPIC image reveals the diffuse X-ray emission associated with the lobes of 3C 98.
### 3.2 X-ray spectra of the host galaxy
We derived EPIC spectra within a circle of $`0^{}.5`$ ($`16.9`$ kpc) radius centered on the nucleus of 3C 98 (denoted as N in Figure 1). The background spectra were derived from a neighboring source-free region with the same radius. We co-added the MOS1 and MOS2 data from Exp2. Figure 4 shows the background-subtracted EPIC MOS and PN spectra of the host galaxy obtained in the individual exposures. They exhibit at least two spectral components; hard and soft ones. The hard component seems heavily absorbed, and its flux varied clearly between the two exposures. The soft and hard components are naturally attributed to the emission from hot gaseous halo and the active nucleus, respectively, of the host galaxy.
We jointly fitted the MOS and PN spectra of each exposure with a two component model, consisting of a soft Raymond-Smith (RS) thermal plasma component and an intrinsically absorbed hard power-law (PL) component. Both components were subjected to the Galactic absorption, $`N_\mathrm{H}=1.17\times 10^{21}`$ cm<sup>-2</sup> (Hardcastle, & Worrall, 1999, and reference therein). We fixed the abundance in the RS component at 0.4 times the solar value, which is typical among nearby elliptical galaxies (Matsushita et al., 2000). The model has successfully reproduced the data, yielding the best-fit spectral parameters summarized in Table 2. We found no clear evidence of Fe K line emission around 6 keV. An apparent data excess below $`0.7`$ keV, seen in the PN spectrum of Exp1, is statistically insignificant.
Within the statistical uncertainty, the spectral parameters of the soft RS component have stayed constant between the two exposures. The best-fit temperature and the 0.5 – 10 keV luminosity becomes $`kT0.85`$ keV and $`L_{\mathrm{RS}}=(89)\times 10^{40}`$ erg s<sup>-1</sup>, respectively, in agreement with hot plasma emission from nearby elliptical galaxies (Matsushita et al., 2000). The PL component has a flat spectral slope with a photon index of $`\mathrm{\Gamma }_\mathrm{X}1.4`$, and a high intrinsic absorption of $`N_\mathrm{H}1\times 10^{23}`$ cm<sup>-2</sup>; these values did not change between the two exposures, whereas its absorption-corrected luminosity was nearly halved meantime. These results support the identification of the hard component with the active nucleus emission. The intrinsic X-ray luminosity (Table 2) lies in the midst of those of radio galaxies (Sambruna et al. 1999). The heavy absorption can be attributed to obscuration by outer regions of the accretion disk, in agreement with good symmetry between the two lobes.
### 3.3 X-ray spectra of the lobes
We accumulated X-ray spectra of the northern and southern lobes within the circular regions LN and LS shown in Figure 1, respectively. The radius of LN is set to $`1^{}.25`$ (42.2 kpc), and that of LS to $`1^{}.45`$ (48.9 kpc), to contain the whole lobe structure. To avoid contamination from the X-ray emission of the host galaxy, we rejected events in a circular region of $`1^{}`$ (33.7 kpc) radius centered on the nucleus, considering a typical size of a galactic halo. We summed the results over the two exposures, and co-added the MOS1 and MOS2 data.
A background estimation is of crucial importance, especially for diffuse and faint X-ray sources. The background spectrum should be accumulated within the same CCD chip as LN and/or LS. In the meantime, we had better select a region, which is in a similar situation to the signal integration area, with respect to the host galaxy. For PN, it is impossible to satisfy the second criterion, because of the gaps of the PN CCD chip. We, alternatively, select the circles of $`1^{}.45`$ radius (same as LS), BG1 and BG2 in Figure 1, as the background region for LN and LS, respectively. This is because a large fraction of LN and LS is within the same PN chip as BG1 and BG2. A dotted circle within BG1 is removed, since an X-ray point source is detected in the position, as shown in Figure 2. We also utilize the sum of the BG1 and BG2 spectra for MOS, after checking that the background around 3C 98 is spatially stable, within the statistical uncertainties.
The background-subtracted MOS and PN spectra of the northern and southern lobes of 3C 98 are shown in the left and right panels of Figure 5, respectively. Because instrumental fluorescent lines severely contaminated the 1.5 – 1.7 keV data, we here and hereafter exclude this energy range from our spectral study. The significance of the northern lobe signals is about $`7.0\sigma `$ and $`8.5\sigma `$ in the 0.6 – 3.5 keV MOS and 0.5 – 6 keV PN data, respectively, while that from the southern lobe signals is about $`2\sigma `$ and $`5.3\sigma `$ in the 0.6 – 3.5 keV MOS and 0.5 – 5 keV PN data. In the following, we analyzed the spectra in these energy range, although we rejected the MOS data of the southern lobe, because of their slightly low signal significance.
We fitted the background-subtracted lobe spectra with a single PL model modified by the Galactic absorption. This model has successfully reproduced the data as shown with histograms in Figure 5, yielding the best fit spectral parameters summarized in Table 3. The best-fit model to the PN spectrum of the southern lobe is consistent with the MOS spectrum, within statistical errors. Thus, the two lobes share almost the same set of parameters, though within rather large errors. In either lobe, the PL fit implies a 0.7 – 7 keV luminosity of $`9\times 10^{40}`$ erg s<sup>-1</sup>.
We also tried to describe the spectra with a thermal bremsstrahlung model (Bremss) modified by the Galactic absorption. As shown in Table 3, the fit turns out to be slightly worse, compared with the single PL model, especially for the northern lobe.
## 4 Discussion
In addition to emission from the host galaxy, including its nucleus, we have significantly detected diffuse faint X-ray emission from the lobes of 3C 98. The X-ray spectra of the northern and southern lobes are reproduced by a single PL model of $`\mathrm{\Gamma }_\mathrm{X}=2.2_{0.5}^{+0.6}`$ and $`1.7_{0.6}^{+0.7}`$, respectively, modified by Galactic absorption.
We are not able to exclude the thermal interpretation from the spectral fitting alone, and the best-fit Bremss model requires the thermal electron densities, $`n_{\mathrm{th}}2\times 10^3`$ cm<sup>-3</sup>, in the lobes of 3C 98. However, numbers of studies on radio polarization effects indicate typically $`n_{\mathrm{th}}10^3`$ cm<sup>-3</sup>, within lobes of FR II radio galaxies (e.g., Burch 1979, Spangler & Sakurai 1985). Therefore, we can conclude that thermal plasma in the lobes of 3C 98 themselves, has only a negligible contribution to the observed diffuse X-ray flux.
Recent Chandra and XMM-Newton observations frequently reveal that the lobes of radio galaxies produce the X-ray holes in the thermal X-ray emission associated with their host galaxies and/or clusters of galaxies around them (e.g., McNamara et al, 2000; Finogurnov, & Jones, 2001), and as a result the holes are usually surrounded by the enhanced X-ray shell of the displaced thermal plasma. The best-fit Bremss temperature which is higher than that of the thermal plasma in the host galaxy (see §3.2), may be consistent with the picture, in which the plasma around the lobes is heated through the bow shock (Kraft et al. 2003), or by kinetic work (Croston et al. 2003) caused by its expansion. However, the thermal plasma to be displaced by the lobes is not detected at the distance corresponding to the lobes, in the direction perpendicular to the lobe axis (see the right panel of Figure 3). The diffuse X-ray emission in 3C 98 clearly has different spatial distribution; i.e., they seem to fill all over the lobes. Moreover, the observed X-ray flux needs a thermal pressure ($`10^{11}`$ dyne cm<sup>-2</sup>) considerably higher than the non-thermal one of the lobes, and it should be difficult for the lobes to remove such a thermal plasma. Therefore, we strongly favor the non-thermal interpretation for the diffuse X-ray emission.
Figure 6 shows the radio and X-ray spectral energy distribution of the northern and southern lobes of 3C 98. Radio data referring to the total flux of 3C 98 are also shown. The contributions from the components other than the lobes (such as the nucleus and hot spots) are reported to be very small (at most 5 %; Hardcastle et al. 1998). At least from 20 MHz and 10 GHz, the total radio synchrotron spectrum is well described by a single PL model (except for only a few deviating points); the flux density at 1.4 GHz and the photon index becomes $`S_\mathrm{R}=11.1\pm 0.1`$ Jy and $`\mathrm{\Gamma }_\mathrm{R}=1.73\pm 0.01`$, respectively. This value of $`\mathrm{\Gamma }_\mathrm{R}`$ falls well within the range of two-point spectral index of both lobes between 1.4 GHz and 8.35 GHz, ($`1.61.75`$). $`\mathrm{\Gamma }_\mathrm{R}`$ is also consistent with the X-ray photon index of the lobes, within the statistical uncertainties. This agreement makes the IC X-ray interpretation fully self-consistent. Hence, we re-analyzed the X-ray spectra with the PL model of which the photon index is fixed at $`\mathrm{\Gamma }_\mathrm{R}`$ and obtained the parameters shown in Table 3.
The non-thermal X-ray spectrum may alternatively be interpreted as a high-energy extension of the synchrotron radio spectrum. However, as is clear from Figure 6, it is impossible to describe the radio and X-ray spectra simultaneously by synchrotron radiation from a single electron population with a PL energy distribution, even if we invoke a synchrotron cutoff in the electron spectrum. Moreover, a synchrotron X-ray photon would need the electron Lorentz factor to exceed $`\gamma _\mathrm{e}4.5\times 10^8E_{\mathrm{keV}}^{0.5}B_{\mu \mathrm{G}}^{0.5}`$ where $`E_{\mathrm{keV}}`$ is the energy of the synchrotron photon in keV, and $`B_{\mu \mathrm{G}}`$ is the magnetic field strength in $`\mu `$G; in that case, the synchrotron cooling time scale of those electrons which are responsible for the synchrotron X-ray photons would become unrealistically short, less than about $`5.5\times 10^4E_{\mathrm{keV}}^{0.5}B_{\mu \mathrm{G}}^{1.5}`$ yr. Therefore, we conclude that our working hypothesis, namely the IC scattering of some seed photons by the synchrotron electron, provides the most plausible explanation for the observed diffuse X-ray emission.
The soft seed photons of the IC scattering in the lobes are usually provided by either CMB photons or IR radiation from the nucleus (see §1). Base on the observed IR flux density, 85 mJy at 25 $`\mu `$m (Golombek et al., 1988), the IR luminosity of the nucleus of 3C 98 is estimated to be $`L_{\mathrm{IR}}2\times 10^{43}`$ erg s<sup>-1</sup>. Even though the orientation and the obscuration effects are taken into account (Heckman et al. 1994), $`L_{\mathrm{IR}}`$ dose not exceed $`10^{44}`$ erg s<sup>-1</sup>. This yields an IR photon energy density of $`u_{\mathrm{IR}}10^{14}(r/50\mathrm{kpc})^2`$ erg cm<sup>-3</sup>, where $`r`$ is the distance from the nucleus in kpc. At $`r8`$ kpc, $`u_{\mathrm{IR}}`$ becomes lower than the CMB energy density $`u_{\mathrm{CMB}}=4.6\times 10^{13}`$ erg cm<sup>-3</sup> at the redshift of 3C 98. We therefore conclude that $`u_{\mathrm{CMB}}`$ dominates over $`u_{\mathrm{IR}}`$ in the larger part of the lobes of 3C 98.
Table 4 summarizes relevant parameters to diagnose energetics in the lobes. We assumed the shape of each lobe to be a simple ellipsoid, with the major and minor axes evaluated from the radio VLA image (Figure 1), to calculate the volume $`V`$ of the lobes. The synchrotron radio flux densities of the individual lobes were derived from the PL fit to the total radio spectrum of Figure 6, and their relative contributions taken from Hardcastle et al. (1998). We calculated the corresponding IC X-ray flux density from the PL fit to the lobe X-ray spectra with the photon index fixed at $`\mathrm{\Gamma }_\mathrm{R}=1.73`$; this specifies the index of the electron number density spectrum as $`2\mathrm{\Gamma }_\mathrm{R}1=2.46`$. To evaluate $`u_\mathrm{e}`$, we assume an electron filling factor of unity, and integrated the electron spectrum between the Lorentz factor of $`\gamma _\mathrm{e}=10^3`$ and $`10^5`$, because the synchrotron or IC radiation from such electrons is directly observable; the lower limit corresponds to 1 keV IC X-ray photons, while the upper limit to synchrotron photons of $`10B_{\mu \mathrm{G}}`$ GHz.
Referring to Harris and Grindlay (1979), we have determined $`u_\mathrm{e}`$, $`u_\mathrm{m}`$, and the corresponding magnetic field $`B`$, as show in Table 4. The physical parameters are found to be quite similar between the two lobes. In both lobes of 3C 98, the electrons are inferred to highly dominate over the magnetic fields, parameterized as $`u_\mathrm{e}/u_\mathrm{m}=(4050)`$. This means that the magnetic field in the lobes is about 2.5 times lower than $`B_{\mathrm{me}}`$, where $`B_{\mathrm{me}}`$ is magnetic field strength which is estimated under the minimun energy condition without proton contribution (Miley 1980).
In order to examine the uncertainty of the result shown in Table 4, we evaluate possible systematic uncertainties in $`S_\mathrm{R}`$, $`S_\mathrm{X}`$, $`\mathrm{\Gamma }_\mathrm{R}`$ and $`V`$, all of which are fundamental observable parameters for the calculation of $`u_\mathrm{e}`$ and $`u_\mathrm{m}`$. Based on the following discussion, the uncertainties in $`u_\mathrm{e}`$ and $`u_\mathrm{m}`$are estimated to be at most 25 % and 35 %, respectively, which can reduce the electron dominance $`u_\mathrm{e}/u_\mathrm{m}`$, only a factor of two. Therefore, our conclusion basically holds.
We consider that the uncertainty in $`S_\mathrm{R}`$ dose not exceed 15 %, based on the radio spectrum in Figure 6. On the other hand, $`S_\mathrm{X}`$ is thought to have an error of about 20 %, mainly due to the integration regions for the signal and background spectrum. Obviously, these errors almost directly propagate to $`u_\mathrm{e}`$ and/or $`u_\mathrm{m}`$, since $`u_\mathrm{e}`$ and $`u_\mathrm{m}`$ is in proportion to $`S_\mathrm{X}`$, and to $`(S_\mathrm{R}S_\mathrm{X}^1)^{2/\mathrm{\Gamma }}`$, respectively.
If the major axis of the lobes (and hence, the jets) of 3C 98 is not precisely perpendicular to our line of sight, the adopted volume of the lobes become smaller than the actual one. This results in an overestimation of $`u_\mathrm{e}`$ with no impact on $`u_\mathrm{m}`$, because of inverse proportionality of $`u_\mathrm{e}`$ to $`V`$. Based on the theoretical prediction of jet to counter jet brightness ratio, taking a relativistic boosting into account (e.g, Giovannini, et al., 2001), and on the observed flux densities of the northern and the southern inner jets of 3C 98 (Hardcastle et al. 1998), we estimated the jet angle to our line of sight to be about 75 degree. The angle put only a negligible effect on $`V`$ and $`u_\mathrm{e}`$ of the lobes (at most 5 %).
The smaller spectral index can artificially enhance the electron dominance. Our employment of $`\mathrm{\Gamma }_\mathrm{R}`$ as the spectral index of the lobes should be safely justified, because the sum flux densities of the lobes highly dominates the total flux density of 3C 98 (about 95 %), and the spectral shape of the lobes are very similar to each other, and also to that of the total spectrum. Base on the radio spectrum of Figure 6 alone, we cannot fully reject the possibility of the index larger than $`\mathrm{\Gamma }_\mathrm{R}`$. However, even the upper limit on the two-point index of the lobes, 1.75, will reduce $`u_\mathrm{e}`$ only about 7 %, and simultaneously enlarge $`u_\mathrm{m}`$ about 15 %, compared with the values in Table 4.
Figure 7 is a compilation of the reported IC X-ray detections in terms of $`u_\mathrm{e}`$ and $`u_\mathrm{m}`$. This figure suggest that the electron dominance of $`u_\mathrm{e}/u_\mathrm{m}10`$ (corresponding to $`B0.5B_{\mathrm{me}}`$, slightly depending on the spectral index), is typically observed in the lobes of radio galaxies, from which the IC X-rays are detected. Recently, Croston et al. (2005) have reported an almost consistent result of $`B0.7B_{\mathrm{me}}`$ in lobes, based on the Chandra and XMM-Newton observations of 33 FR-II radio galaxies and quasars. The figure indicates that the present result on 3C 98 is very typical among these measurements.
The observed electron dominance will raise an important question; what does confine the non-thermal plasma within the lobes ? Although we regard that it is currently difficult to answer the question, the spatial distribution of $`u_\mathrm{e}`$ and $`u_\mathrm{m}`$ in the lobes may provide us an useful clue to the problem. In several lobes in which the electron dominance is already reported, almost uniform or center-filled distributions of IC X-rays are found, in spite of a rim-brightening synchrotron radio feature (Tashiro et al. 1998; Tashiro et al. 2001; Isobe et al. 2002; Comastri et al. 2003). A comparison of the X-ray and radio distributions indicates that the magnetic field is enhanced toward the edge of the lobes, to become even closer to the equipartition (Tashiro et al., 2001; Isobe et al., 2002), while the electrons are relatively uniformly distributed over the lobes. As shown in Figure 3, the lobes of 3C 98 have the tendency though with a little large statistical errors. However, the clear solution to the issue inevitably needs more detailed observations with higher statistics, and theoretical considerations.
We are gratefull to the anonymous referee for his constructive comments to improve the present paper. We thank Dr. I. Takahashi for his guidance of the XMM-Newton data analysis, especially the background estimation. This paper is based on observations obtained with XMM-Newton, an ESA science mission with instruments and contributions directly funded by ESA Member States and NASA. This research has made use of the NASA/IPAC Extra galactic Database (NED) (the Jet Propulsion Laboratory, California Institute of Technology, the National Aeronautics and Space Administration). The unpublished VLA image of 3C 98 (Bridle) was downloaded from “An Atlas of DRAGNs” (http://www.jb.man.ac.uk/atlas), edited by Leahy, Bridle, & Strom. |
warning/0506/quant-ph0506173.html | ar5iv | text | # Quantum Mechanics in Multiply-Connected Spaces
## 1 Introduction
We shall be concerned here with topological effects in quantum mechanics, and shall elaborate on some results first described in . The kind of statement on which we shall focus asserts that if the configuration space $`𝒬`$ is a multiply-connected<sup>6</sup><sup>6</sup>6Recall that a manifold $`𝒬`$ is *simply connected* if all closed curves in $`𝒬`$ are contractible. Otherwise, it is *multiply connected*. (Note that “multiply connected” is different from the notion “$`n`$-connected” for $`n0`$, which is sometimes used in the literature on algebraic topology \[37, p. 51\] and means that the first $`n`$ homotopy groups $`\pi _n(𝒬)`$ are all trivial.) Examples of simply-connected spaces are $`^d`$ for any $`d0`$, the spheres $`S^d`$ for $`d2`$, or the punctured spaces $`^d\{0\}`$ for $`d3`$; examples of multiply-connected spaces are the circle $`S^1`$, the torus $`S^1\times S^1`$, the punctured plane $`^2\{0\}`$, or generally $`^dU`$ where $`U`$ is a subspace of dimension $`d2`$, $`d2`$. Riemannian manifold then there exist *several* quantum theories in $`𝒬`$. More precisely, the dynamics is not completely determined by specifying $`𝒬`$ (whose metric we regard as incorporating the “masses of the particles”) together with the potential and the value space of the wave function; in addition, one can choose *topological factors*, which form a representation (or twisted representation) of the fundamental group $`\pi _1(𝒬)`$ of $`𝒬`$. In each of the theories, the Hamiltonian is locally equivalent to $`\frac{\mathrm{}^2}{2}\mathrm{\Delta }+V`$, though not globally. The investigation in this paper is continued with other methods in three follow-up papers .
Our interest lies in explaining why there is more than one quantum theory and how the several possibilities arise, and in classifying the possibilities. The formulation of quantum mechanics we use for this purpose is Bohmian mechanics , a quantum theory without observers; it describes a world in which particles have trajectories, guided by a wave function $`\psi _t`$; observers in this world would find that the results of their experiments obey the quantum formalism . We will give a brief review of Bohmian mechanics in Section 4. Most of our mathematical considerations and methods are equally valid, relevant, and useful in orthodox quantum mechanics, or any other version of quantum mechanics. Bohmian mechanics, however, provides a sharp mathematical justification of these considerations that is absent in the orthodox framework.
The motion of the configuration in a Bohmian $`N`$-particle system can be regarded as corresponding to a dynamical system in the configuration space $`𝒬=^{3N}`$, defined by a time-dependent vector field $`v^{\psi _t}`$ on $`𝒬`$ which in turn is defined, by the Bohmian law of motion, in terms of $`\psi _t`$. We are concerned here with the analogues of the Bohmian law of motion for the case that $`𝒬`$ is, instead of $`^{3N}`$, an arbitrary Riemannian manifold.<sup>7</sup><sup>7</sup>7Manifolds will throughout be assumed to be Hausdorff, paracompact, connected, and $`C^{\mathrm{}}`$. They need not be orientable. The main result is that, if $`𝒬`$ is multiply connected, there are several such analogues: several Bohmian dynamics, which we will describe in detail, corresponding to different choices of the topological factors.
To define a Bohmian theory in a manifold $`𝒬`$ ultimately amounts to defining trajectories in $`𝒬`$ and their probabilities. This leads to clear mathematical classification questions, while from the orthodox point of view the ground rules with respect to the issue of the existence of several quantum theories with different topological factors are less clear. We will review the differences between the two viewpoints in Section 2.
The topological factors consists of, in the simplest case, phase factors associated with non-contractible loops in $`𝒬`$, forming a character<sup>8</sup><sup>8</sup>8By a *character* of a group we refer to what is sometimes called a unitary multiplicative character, i.e., a one-dimensional unitary representation of the group. of the fundamental group $`\pi _1(𝒬)`$. All characters can be physically relevant; we emphasize this because it is easy to overlook the multitude of dynamics by focussing too much on just one, the simplest one, which we will define in Section 5: the *immediate generalization* of the Bohmian dynamics from $`^{3N}`$ to a Riemannian manifold, or, as we shall briefly call it, the *immediate Bohmian dynamics*.
Apart from the mathematical exercise, what do we gain from studying the possible Bohmian dynamics on manifolds?
* A new understanding of how topological factors in quantum mechanics can be regarded as arising.
* A presumably complete *classification* of the topological factors in quantum mechanics, including some, corresponding to what we call *twisted representations* of $`\pi _1(𝒬)`$, that have not, to our knowledge, been considered so far in the literature.
* An explanation of the fact that the wave function of a system of identical particles is either symmetric or anti-symmetric, a fact that (at least insofar as nonrelativistic quantum mechanics is concerned) is usually, instead of being derived, introduced as a *symmetrization postulate*. This application is discussed in detail in a sister paper to this one , and will only be touched upon briefly here.
Our main motivation for studying the question of Bohmian dynamics on manifolds was in fact the investigation of the symmetrization postulate for identical particles.
As we have already mentioned, one of the different Bohmian dynamics on a manifold $`𝒬`$ is special, as it is the immediate Bohmian dynamics on $`𝒬`$.<sup>9</sup><sup>9</sup>9That is, if one considers the value space of the wave function as given. In , in contrast, we obtain several bundles of spin spaces from the Riemannian geometry of the configuration space. We could take any of these bundles as the starting point for defining the dynamics, and then which one of the dynamics is immediate will depend on this choice. The other kinds of Bohmian dynamics come in a hierarchy of increasing complexity. There are three natural classes $`𝒞_1,𝒞_2,𝒞_3`$ of Bohmian dynamics, related according to
$$𝒞_0𝒞_1𝒞_2𝒞_3,$$
(1)
where $`𝒞_0`$ contains only the immediate Bohmian dynamics. The dynamics of class $`𝒞_1`$, defined in Section 6.4, involve topological phase factors forming a character of the fundamental group $`\pi _1(𝒬)`$. Those of class $`𝒞_2`$, defined in Section 8.1 and in a more general setting in Section 8.4, involve topological factors that are given by matrices, forming a unitary representation of $`\pi _1(𝒬)`$ or, in the case of a vector bundle, a twisted representation (see the end of Section 8.4 for the definition). Those of class $`𝒞_3`$ will not be discussed here but in ; they involve changes in connections and potentials that are not based on multiple connectivity. As we shall explain, the dynamics of bosons belongs to $`𝒞_0`$ while that of fermions belongs to $`𝒞_1`$. More precisely, fermions can be regarded as belonging either to $`𝒞_0`$, for a certain nontrivial vector bundle defined in (for which bosons are of class $`𝒞_1`$), or to $`𝒞_1`$, for the trivial bundle $`𝒬\times `$. In Section 7 we derive a dynamics of class $`𝒞_1`$ for the Aharonov–Bohm effect. We will define dynamics here in a non-rigorous way; a rigorous definition of the classes $`𝒞_0,𝒞_1,𝒞_2,𝒞_3`$ is given in .
It is not obvious what “all possible kinds of Bohmian dynamics” should mean. We will investigate one approach here, while others, as already mentioned, are studied in . The present approach is based on considering wave functions $`\psi `$ that are defined not on the configuration space $`𝒬`$ but on its universal covering space $`\widehat{𝒬}`$. We then study which kinds of periodicity conditions, relating the values on different levels of the covering fiber by a topological factor, will ensure that the Bohmian velocity vector field associated with $`\psi `$ is projectable from $`\widehat{𝒬}`$ to $`𝒬`$. This is carried out in Section 6 for scalar wave functions and in Section 8 for wave functions with values in a complex vector space (such as a spin-space) or a complex vector bundle. In the case of vector bundles, we derive a novel kind of topological factor, given by a twisted representation of $`\pi _1(𝒬)`$.
Let us mention the other approaches to defining the dynamics of classes $`𝒞_1,𝒞_2`$, and $`𝒞_3`$. Since wave functions can be regarded as sections of Hermitian bundles, i.e., complex vector bundles with a connection and parallel Hermitian inner products, one approach considers all Hermitian bundles that are locally (but not globally) isomorphic to a given one. Another approach expresses the dynamics in terms of the Hamiltonian and considers all Hamiltonians $`H`$ that are locally (though not necessarily globally) equivalent to $`\frac{\mathrm{}^2}{2}\mathrm{\Delta }+V`$. Another approach regards the value space of the wave function as a representation space (such as a spin-space) of a suitable group (such as the rotation group) and classifies the Hermitian bundles consisting of representation spaces. A last approach removes a surface $`\kappa `$ from configuration space $`𝒬`$, such that $`𝒬\kappa `$ is simply connected, and imposes a periodic boundary condition relating the wave function on both sides of the new, “virtual”, boundary $`\kappa `$ by a topological factor.
In some cases, some of the classes coincide: When $`𝒬`$ is simply connected, then $`𝒞_0=𝒞_1=𝒞_2`$. When the wave function is a scalar (as for spinless particles), then $`𝒞_1=𝒞_2=𝒞_3`$. For generic potentials, $`𝒞_1=𝒞_2=𝒞_3`$.
We encounter examples of multiply-connected configuration spaces in two ways: either, as in the Aharonov–Bohm effect, by ignoring an existing part of physical (or configuration) space, or, as for identical particles or multiply-connected cosmologies, from the very nature of the configuration space. In the former case the topological factors of the *effective* dynamics on the available configuration space depend on external fields, while in the latter case the topological factors of the *fundamental* dynamics should be compatible with any choice of external fields. This compatibility is a strong restriction, which allows only the dynamics of class $`𝒞_1`$. Hence, as we shall argue more fully later, the several fundamental quantum theories in $`𝒬`$ are those given by $`𝒞_1`$. This conclusion we call the *Character Quantization Principle*, since the dynamics of $`𝒞_1`$ are defined using the characters of the fundamental group of $`𝒬`$. It is formulated and discussed in Section 9. We conclude in Section 10.
The notion that multiply-connected spaces give rise to different quantum theories is not new. Here is a sampling of the literature. A covering space was used at least as early as 1950 by Bopp and Haag for the configuration space of the spinning top; it was more fully exploited by Dowker , and was used by Leinaas and Myrheim for the configuration space of identical particles. Vector potentials on multiply connected spaces were used by Aharonov and Bohm in . Path integrals on multiply connected spaces began largely with the work of Schulman and that of Laidlaw and DeWitt in ; see for details. There is also the current algebra approach of Goldin, Menikoff, and Sharp . Most of these works are dedicated to scalar wave functions. The study of arbitrary manifolds began with Laidlaw and DeWitt , which deals with path integration on the universal covering space for scalar wave functions. Nelson derives the topological phase factors for scalar wave functions from stochastic mechanics. Gamboa and Rivelles consider relativistic Hamiltonians using a path-integral approach. Ho and Morgan provide a study of quantum mechanics on $`^d\times S^1`$ for scalar wave functions.
## 2 Perspective on Orthodox Quantum Mechanics
When one considers a Bohmian dynamics and removes the Bohmian trajectories, there still remains the wave function $`\psi `$, and a number of nontrivial things can be said about it, such as, which space $`\psi `$ can be taken from and how its evolution is defined. As a consequence, much of the mathematical discussion in this paper would be equally valid, applicable and relevant for any other formulation of quantum mechanics. However, our analysis of the emergence of the further kinds of wave functions, whose main role in Bohmian mechanics is to define trajectories, would not work in the same way if one were to dispense with the trajectories.
One would meet, when trying to carry out the program of this paper in orthodox quantum mechanics, some difficulties that are absent in the Bohmian approach. This is mainly because of two traits of Bohmian mechanics: first, it is clear in the Bohmian framework at which point the specification of a theory is complete; and second, it is clear whether two variants of a theory are physically equivalent or not. Let us explain.
In the Bohmian framework, once the possible trajectories of all particles have been defined (together with the appropriate equivariant probability distribution, see Section 4) then the theory has been completely specified, and there is neither need nor room for further axioms. In orthodox quantum mechanics, in contrast, it is not obvious what it is that needs to be specified in order to have a variant of the theory. The Hilbert space $``$ and the Hamiltonian $`H`$? While they certainly must be specified, they are certainly not enough.
One could think, for example, of different possible position observables in the same Hilbert space, and these would lead to different predictions for position measurements. Thus, one should specify, it would seem, $``$, $`H`$, and the operator, or commuting set of operators, for the position observable. But would that be enough? Need we not also be told what operator represents the momentum observable? Need we not be told what operators represent *all* the observables? This should be contrasted with the fact that in Bohmian mechanics, once the dynamics of the particles is specified, also the outcomes of all experiments are specified.<sup>10</sup><sup>10</sup>10One could argue that for exactly the same reason, to specify the position observable in orthodox quantum mechanics would be sufficient, as it would fix the statistics of the outcomes of every measurement. This is true, and we think that this is a healthy attitude. However, it is also quite against the spirit of orthodox quantum mechanics which sets a high value on the “democracy” for all observables.
And what are, by the way, “all” observables? It seems clear that the list of all observables should begin with position, momentum, and energy, but where it should end is rather obscure. In addition, the notion of observable becomes somewhat problematic when the configuration space $`𝒬`$ is a manifold. The problem is not so much that the position observable can no longer be represented by a set of commuting position operators, as the manifold may not permit global coordinates (e.g., on the circle); one should conclude that the appropriate notion of position observable is then a PVM (projection-valued measure) on $`𝒬`$ acting on $``$, associating with every subset of $`𝒬`$ a projection in $``$. The more serious problem concerns the momentum observable: already on the half-line, the operator $`p=i\mathrm{}d/dq`$ does not have a self-adjoint extension. More generally, on a Riemannian manifold $`𝒬`$ the notion of the momentum observable becomes obscure, as it is based on a translation symmetry that may not exist in $`𝒬`$. Thus, a momentum observable may not exist. This brings us back to the point that it is not clear which observables need to be specified in order to specify an orthodox quantum theory.
The contrast between the clarity of Bohmian mechanics and the vagueness of orthodox quantum theory is perhaps even more striking when we consider the issue of the physical equivalence of theories. In this paper we shall always treat Bohmian theories, when they are mathematically different but lead to the same trajectories (and probabilities), as physically equivalent. For example, the dynamics we shall define using wave functions on the covering space with the trivial character is physically equivalent to the immediate Bohmian dynamics, using wave functions on the configuration space.
In orthodox quantum mechanics, when should we regard two variants of the theory as physically equivalent? The answer in the spirit of orthodox quantum mechanics is, when they predict the same statistics for outcomes for all experiments; that is, when they are empirically equivalent.<sup>11</sup><sup>11</sup>11This answer can be criticized on the grounds that there are known examples of theories that are empirically equivalent though physically inequivalent, such as Bohmian mechanics and stochastic mechanics , or the variants of Bohmian mechanics in which some of the particles do not possess actual positions while their coordinates get integrated over in the law of motion . This answer leads again to the question, what are “all” observables? In addition, it leads us to the possibly separate problem of identifying the observables of one theory with the observables of another. Within the Bohmian framework, based on a clear ontology and a correspondingly sharp specification of the relevant physical structures and their behavior, no such questions and problems can arise.
## 3 Perspective on Spontaneous Collapse Theories
Another approach besides Bohmian mechanics leading to quantum theories without observers is that of spontaneous wave function collapse ; the simplest and best known model of this kind is due to Ghirardi, Rimini, and Weber (GRW) . Its situation with respect to topological factors is very different from that of Bohmian mechanics. For example, the situation of identical particles in the GRW theory is different from that in Bohmian mechanics because the latter is (in a suitable sense) automatically compatible with bosons and fermions, whereas the equations of the GRW model require modification for identical particles as follows .
In the original GRW model (corresponding to $`N`$ distinguishable particles), collapses are associated with points in 3-space and labels $`i\{1,\mathrm{},N\}`$. Given the wave function $`\psi :^{3N}`$, a collapse with label $`i`$ and location $`𝒙^3`$ occurs with rate
$$r_i(𝒙|\psi )=\psi |\mathrm{\Lambda }_i(𝒙)\psi ,$$
(2)
where the collapse rate operator $`\mathrm{\Lambda }_i(𝒙)`$ is a multiplication operator defined by
$$\mathrm{\Lambda }_i(𝒙)\psi (𝒒_1,\mathrm{},𝒒_N)=\lambda \mathrm{exp}\left(\frac{(𝒙𝒒_i)^2}{2a^2}\right)\psi (𝒒_1,\mathrm{},𝒒_N).$$
(3)
The constants $`\lambda `$ and $`a`$ are parameters of the model. A collapse at time $`t`$ and location $`𝒙`$ with label $`i`$ changes the wave function according to
$$\psi _t\psi _{t+}=\frac{\mathrm{\Lambda }_i(𝒙)^{1/2}\psi _t}{\mathrm{\Lambda }_i(𝒙)^{1/2}\psi _t}.$$
(4)
In the version for identical particles, collapses are associated with locations $`𝒙`$ only, without labels. Letting $`\psi `$ be either a symmetric or an anti-symmetric function on $`^{3N}`$, a collapse occurs at location $`𝒙^3`$ with rate
$$r(𝒙|\psi )=\psi |\mathrm{\Lambda }(𝒙)\psi ,$$
(5)
where
$$\mathrm{\Lambda }(𝒙)=\underset{i=1}{\overset{N}{}}\mathrm{\Lambda }_i(𝒙),$$
(6)
and changes $`\psi `$ according to
$$\psi \frac{\mathrm{\Lambda }(𝒙)^{1/2}\psi }{\mathrm{\Lambda }(𝒙)^{1/2}\psi }.$$
(7)
Thus, the collapsed wave function is again symmetric respectively anti-symmetric.
The arguments used in the present paper for deriving the topological factors cannot be repeated in the context of the GRW theory because they rely on particle configurations, which do not exist in the GRW theory. As a consequence, indeed, configuration space does not play, in the GRW theory, the same central role as in Bohmian mechanics, but merely that of a convenient tool for representing the state vector as a function (similar to the role, in Bohmian mechanics, of momentum space or of the set of spin eigenvalues). Thus, it is hard to see how the GRW theory could provide any reasons for the existence of several possibilities, corresponding to different topological factors, in situations in which the configuration space of the Bohmian theory is multiply connected. Moreover, for the GRW theory, for which there are no particles to begin with, it is hard to see why the multiply-connected natural configuration space $`{}_{}{}^{N}_{}^{3}`$ for $`N`$ identical particles, see Section 5, should be considered at all. Nonetheless, topological factors can always be introduced into GRW theories, as in the example above, as we shall explain later in Section 6.5, Remark 6.
## 4 Bohmian Mechanics in $`^{3N}`$
Bohmian mechanics is a theory about particles with definite locations. The theory specifies the trajectories in physical space of these particles. The object which determines the trajectories is the wave function, familiar from quantum mechanics. More precisely, the state of the system in Bohmian mechanics is given by the pair $`(Q,\psi )`$; $`Q=(𝑸_1,\mathrm{},𝑸_N)^{3N}`$ is the configuration of the $`N`$ particles in our system and $`\psi `$ is a (standard quantum mechanical) wave function on the configuration space $`^{3N}`$, taking values in some *Hermitian vector space* $`W`$, i.e., a finite-dimensional complex vector space endowed with a positive-definite Hermitian (i.e., conjugate-symmetric and sesqui-linear) inner product $`(,)`$. The state of the system changes according to the guiding equation and Schrödinger’s equation:
$$\frac{d𝑸_k}{dt}=\frac{\mathrm{}}{m_k}\mathrm{Im}\frac{(\psi ,_k\psi )}{(\psi ,\psi )}(𝑸_1,\mathrm{},𝑸_N)=:v_k^\psi (Q),k=1,\mathrm{},N$$
(8)
$$i\mathrm{}\frac{\psi }{t}=\underset{k=1}{\overset{N}{}}\frac{\mathrm{}^2}{2m_k}\mathrm{\Delta }_k\psi +V\psi $$
(9)
where $`V`$ is the potential function with values given by Hermitian matrices (endomorphisms of $`W`$). We call $`(\varphi (q),\psi (q))`$, the inner product on the value space $`W`$, the *local inner product*, in distinction from the inner product $`\varphi ,\psi `$ on the Hilbert space of wave functions. For complex-valued wave functions, the potential is a real-valued function on configuration space and the local inner product is $`\overline{\varphi (q)}\psi (q)`$, where the bar denotes complex conjugation.
The empirical agreement between Bohmian mechanics and standard quantum mechanics is grounded in equivariance . In Bohmian mechanics, if the configuration is initially random and distributed according to $`|\psi _0|^2`$, then the evolution is such that the configuration at time $`t`$ will be distributed according to $`|\psi _t|^2`$. This property is called the equivariance of the $`|\psi |^2`$ distribution. It follows from comparing the transport equation
$$\frac{\rho _t}{t}=(\rho _tv^{\psi _t})$$
(10)
for the distribution $`\rho _t`$ of the configuration $`Q_t`$, where $`v^\psi =(v_1^\psi ,\mathrm{},v_N^\psi )`$, to the quantum continuity equation
$$\frac{|\psi _t|^2}{t}=(|\psi _t|^2v^{\psi _t}),$$
(11)
which is a consequence of Schrödinger’s equation (9). A rigorous proof of equivariance requires showing that almost all (with respect to the $`|\psi |^2`$ distribution) solutions of (8) exist for all times. This was done in . A more comprehensive introduction to Bohmian mechanics may be found in .
Spin is already incorporated in (8) and (9) if one chooses for $`W`$ a suitable spin space . By assumption, for one particle moving in $`^3`$, $`W`$ is a complex, irreducible representation space of $`SU(2)`$, the universal covering group<sup>12</sup><sup>12</sup>12The universal covering space of a Lie group is again a Lie group, the *universal covering group*. It should be distinguished from another group also called the *covering group*: the group $`Cov(\widehat{𝒬},𝒬)`$ of the covering (or deck) transformations of the universal covering space $`\widehat{𝒬}`$ of a manifold $`𝒬`$, which will play an important role later. of the rotation group $`SO(3)`$. If it is the spin-$`s`$ representation then $`W=^{2s+1}`$.
## 5 The Immediate Generalization to Riemannian Manifolds
We now consider, in the role of the configuration space, a Riemannian manifold $`𝒬`$ instead of $`^{3N}`$. The primary physical motivation is the study of identical particles, for which the natural configuration space is the set $`{}_{}{}^{N}_{}^{3}`$ of all $`N`$-element subsets of $`^3`$,
$${}_{}{}^{N}_{}^{3}:=\{S|S^3,|S|=N\},$$
(12)
which naturally carries the structure of a Riemannian manifold, in fact a multiply-connected one. This configuration space was first suggested in and ; for further discussion see .
But the generalization to manifolds is also very natural mathematically. In addition, there are further cases of physical relevance: One could consider, instead of $`^3`$, a curved physical space. And in cases like the Aharonov–Bohm effect, the phase shift that occurs can be attributed to the topology of the effectively available configuration space, a subset of the entire configuration space that can be viewed as a multiply-connected manifold.
### 5.1 Euclidean Vector Spaces
It is in fact easy to find a generalization of the Bohmian dynamics to a Riemannian manifold $`𝒬`$, which we call the immediate Bohmian dynamics on $`𝒬`$. One reason why it is so easy is that the law of motion for the point $`Q_t=Q(t)=(𝐐_1(t),\mathrm{},𝐐_N(t))`$ in the configuration space $`^{3N}`$ representing the positions of all particles at time $`t`$ is almost independent of the way in which $`^{3N}`$ is composed of $`N`$ copies of $`^3`$. In fact, (8) can be written as
$$\frac{dQ_t}{dt}=\mathrm{}𝔪^1\mathrm{Im}\frac{(\psi ,\psi )}{(\psi ,\psi )}(Q_t)$$
(13)
where $`𝔪`$ is the diagonal matrix with the masses as entries, each mass $`m_k`$ appearing 3 times. That is, as soon as $`𝔪`$ is given, the information about which directions in $`^{3N}`$ correspond to the single factors $`^3`$ becomes irrelevant for defining the dynamics of $`Q_t`$. Eq. (13) would as well define a dynamics on any Euclidean vector space $``$ of finite dimension, given a wave function $`\psi `$ on $``$ and a positive-definite symmetric endomorphism $`𝔪:`$.
It will be convenient to include the mass matrix $`𝔪`$ in the metric $`g_{ab}`$ of $``$,
$$g_{ab}=\underset{c=1}{\overset{dim}{}}g_{ac}^{}𝔪_b^c,$$
(14)
where $`g_{ab}^{}`$ is the metric of $``$ before the inclusion of masses, and indices $`a,b,c`$ run through the dimensions of $``$. In the standard example of $`^{3N}`$, this amounts to introducing the metric
$$g_{i\alpha ,j\beta }=m_i\delta _{ij}\delta _{\alpha \beta },$$
(15)
where $`i,j=1,\mathrm{},N`$ and $`\alpha ,\beta =1,2,3`$ (and the index $`i`$ occurring twice on the right is not summed over). With $``$ then defined using $`g`$ instead of $`g^{}`$, (13) becomes
$$\frac{dQ_t}{dt}=\mathrm{}\mathrm{Im}\frac{(\psi ,\psi )}{(\psi ,\psi )}(Q_t).$$
(16)
(Note that in order to turn the covector given by the differential of $`\psi `$ into a vector, one uses $`g^{ab}`$.)
Similarly, the Schrödinger equation (9) can then be written
$$i\mathrm{}\frac{\psi }{t}=\frac{\mathrm{}^2}{2}\mathrm{\Delta }\psi +V\psi ,$$
(17)
where $`\mathrm{\Delta }`$, the Laplacian on $``$, is to be understood as the metric trace of the second derivatives,
$$\mathrm{\Delta }=g^{ab}_a_b$$
(18)
(in abstract-index notation with sum convention). Thus, also the Schrödinger equation is well-defined on a Euclidean space $``$, or, in other words, it is independent of the product structure of $`^{3N}=(^3)^N`$.
### 5.2 Riemannian Manifolds
In order to transfer (16) and (17) to Riemannian manifolds, we need only replace $``$ by the tangent space $`T_{Q(t)}𝒬`$. In this subsection, the wave functions we consider are $`W`$-valued functions on $`𝒬`$, with $`W`$ a Hermitian vector space.
We begin by recalling the definitions of the gradient and the Laplacian on Riemannian manifolds. By the gradient $`f`$ of a function $`f:𝒬`$ we mean the tangent vector field on $`𝒬`$ metrically equivalent (by “raising the index”) to the 1-form $`df`$, the differential of $`f`$. For a function $`\psi :𝒬W`$, the differential $`d\psi `$ is a $`W`$-valued 1-form, and thus $`\psi (q)T_q𝒬W`$, where $`T_q𝒬`$ denotes the complexified tangent space at $`q`$, and the tensor product $``$ is, as always in the following, over the complex numbers. The Laplacian $`\mathrm{\Delta }f`$ of a function $`f`$ is defined to be the divergence of $`f`$, where the divergence of a vector field $`X`$ is defined by
$$\mathrm{div}X=D_aX^a$$
(19)
with $`D`$ the (standard) covariant derivative operator, corresponding to the Levi-Civita connection on the tangent bundle of $`𝒬`$ arising from the metric $`g`$. Since $`Dg=0`$, we can write
$$\mathrm{\Delta }f=g^{ab}D_aD_bf,$$
(20)
where the second $`D`$, the one which is applied first, actually does not make use of the Levi-Civita connection. In other words, the Laplacian is the metric trace of the second (covariant) derivative. Another equivalent definition is $`\mathrm{\Delta }f=ddf`$ where $`d`$ is the exterior derivative of differential forms and $``$ is the Hodge star operator (see, e.g., ).<sup>13</sup><sup>13</sup>13The Hodge operator $``$ depends on the orientation of $`𝒬`$ in such a way that a change of orientation changes the sign of the result. Thus, $``$ does not exist if $`𝒬`$ is not orientable. However, it exists locally for any chosen local orientation, and since the Laplacian contains two Hodge operators, it is not affected by the sign ambiguity. For $`W`$-valued functions $`\psi `$ the Laplacian $`\mathrm{\Delta }\psi `$ is defined correspondingly as the divergence of the “$`W`$-valued vector field” $`\psi `$, or equivalently by
$$\mathrm{\Delta }\psi =g^{ab}D_aD_b\psi $$
(21)
or by $`\mathrm{\Delta }\psi =dd\psi `$, using the obvious extension of the exterior derivative to $`W`$-valued differential forms.
The time evolution of the state $`(Q_t,\psi _t)`$ is simply given by the same formal equations as (16) and (17) with the appropriate interpretation of $``$ and $`\mathrm{\Delta }`$. We give the equations for future reference:
$`{\displaystyle \frac{dQ_t}{dt}}`$ $`=v^{\psi _t}(Q_t)`$ (22a)
$`i\mathrm{}{\displaystyle \frac{\psi _t}{t}}`$ $`=\frac{\mathrm{}^2}{2}\mathrm{\Delta }\psi _t+V\psi _t,`$ (22b)
where the Bohmian velocity vector field $`v^\psi `$ associated to the wave function $`\psi `$ is
$$v^\psi :=\mathrm{}\mathrm{Im}\frac{(\psi ,\psi )}{(\psi ,\psi )}.$$
(23)
Thus, given $`𝒬`$, $`W`$, and $`V`$, we have specified a Bohmian dynamics, the *immediate Bohmian dynamics*.<sup>14</sup><sup>14</sup>14The question arises whether these equations possess unique solutions, for all times or at least for short times. For some Riemannian manifolds $`𝒬`$ this may require the introduction of boundary conditions. Since the existence question is mathematically demanding and not our concern here, we make only a few remarks: For a discussion of the existence question of Bohmian trajectories in $`^{3N}`$, see . The existence of the evolution of the wave function amounts to defining the Hamiltonian $`H`$ as a self-adjoint operator, i.e., as a self-adjoint extension of $`H^0=\frac{\mathrm{}^2}{2}\mathrm{\Delta }+V`$, with $`H^0`$ defined on $`C_0^{\mathrm{}}(𝒬,W)`$, the space of smooth $`W`$-valued functions with compact support. As far as we know, it is not known in all cases whether a self-adjoint extension exists, and, when so, how many exist, and what the physical meaning of the different extensions is when there is more than one. When several extensions exist, they must perhaps be regarded as different possible Bohmian dynamics on $`𝒬`$, and thus as further possibilities, not captured in the classes $`𝒞_0,𝒞_1,𝒞_2`$ considered in this paper. We shall not pursue this idea further, and shall reason instead in terms of the “formal” dynamics. We introduce the notation $`𝒞_0(𝒬,W,V)`$ for the set containing just this one dynamics. We also write $`𝒞_0(𝒬,V)`$ for $`𝒞_0(𝒬,,V)`$. (A rigorous definition of what is meant here by a “dynamics,” avoiding the question of the existence of solutions, is given in . For now, we simply proceed as if we have the global existence of solutions and say (a bit vaguely) that a “dynamics” is defined by a set of wave functions, in this case $`C_0^{\mathrm{}}(𝒬,W)`$, and for every wave function $`\psi `$ a set $`𝒮^\psi `$ of trajectories in $`𝒬`$, in this case the solutions of (22a), together with a probability distribution $`\rho ^\psi `$ on $`𝒮^\psi `$, in this case given by
$$\rho ^\psi (dQ)=(\psi _t(Q_t),\psi _t(Q_t))dQ_t.$$
(24)
Due to equivariance, (24) is independent of $`t`$.)
### 5.3 An Example
An important case is that of several particles moving in a Riemannian manifold $`M`$, a possibly curved physical space. Then the configuration space for $`N`$ distinguished particles is $`𝒬:=M^N`$. Let the masses of the particles be $`m_i`$ and the metric of $`M`$ be $`g`$. Then the relevant metric on $`M^N`$, the analogue of (14) and (15) acting on the tangent space $`T_{(𝒒_1,\mathrm{},𝒒_N)}M^N=_{i=1}^NT_{𝒒_i}M`$, is
$$g^N(v_1\mathrm{}v_N,w_1\mathrm{}w_N):=\underset{i=1}{\overset{N}{}}m_ig(v_i,w_i).$$
Using $`g^N`$ allows us to write (23) and (22) instead of the equivalent equations
$$\frac{d𝑸_k}{dt}=\frac{\mathrm{}}{m_k}\mathrm{Im}\frac{(\psi ,_k\psi )}{(\psi ,\psi )}(𝑸_1,\mathrm{},𝑸_N),k=1,\mathrm{},N$$
(25)
$$i\mathrm{}\frac{\psi }{t}=\underset{k=1}{\overset{N}{}}\frac{\mathrm{}^2}{2m_k}\mathrm{\Delta }_k\psi +V\psi ,$$
(26)
where $`𝑸_k`$, the $`k^{th}`$ component of $`Q`$, lies in $`M`$, and $`_k`$ and $`\mathrm{\Delta }_k`$ are the gradient and the Laplacian with respect to $`g`$, acting on the $`k^{th}`$ factor of $`M^N`$. We take $`W=`$. Observe that (8) and (9) are special cases, corresponding to Euclidean space $`M=^3`$, of (25) and (26).
The configuration space of $`N`$ identical particles in $`M`$ is
$${}_{}{}^{N}M:=\{S|SM,|S|=N\},$$
(27)
which inherits a Riemannian metric from $`M`$, see .
### 5.4 Vector Bundles
Even more generally, we can consider a Bohmian dynamics for wave functions taking values in a complex vector bundle $`E`$ over the Riemannian manifold $`𝒬`$. That is, the value space then depends on the configuration, and wave functions become sections of the vector bundle.<sup>15</sup><sup>15</sup>15Recall that a *section* (also known as *cross-section*) of $`E`$ is a map $`\psi :𝒬E`$ such that $`\psi (q)E_q`$, i.e. it maps a point $`q`$ of $`𝒬`$ to an element of the vector fiber over $`q`$. For example, a vector field on a manifold $`M`$ is a section of the tangent bundle $`TM`$.
Such a case occurs for identical particles with spin $`s`$, where the bundle $`E`$ of spin spaces over the configuration space $`𝒬={}_{}{}^{N}_{}^{3}`$ defined in (12) consists of the $`(2s+1)^N`$-dimensional spaces
$$E_q=\underset{𝒒q}{}^{2s+1},q𝒬.$$
(28)
For a detailed discussion of this bundle, of why this is the right bundle, and of the notion of a tensor product over an arbitrary index set, see . Vector bundles also occur for particles with spin in a curved physical space. In addition to their physical relevance, bundles are a natural mathematical generalization of our previous setting involving wave functions defined on manifolds. Finally, the approaches we use in for suggesting natural classes of Bohmian dynamics are based on considerations concerning vector bundles (even for spinless particles).
We introduce some notation and terminology. $`C^{\mathrm{}}(E)`$ will denote the set of smooth sections while $`C_0^{\mathrm{}}(E)`$ will be the set of smooth sections with compact support.
###### Definition 1.
A *Hermitian vector bundle*, or *Hermitian bundle*, is a finite-dimensional complex vector bundle $`E`$ with a connection and a positive-definite, Hermitian local inner product $`(,)_q`$ on $`E_q`$ which is parallel.
Recall that a connection defines (and is defined by) a notion of parallel transport of vectors in $`E_q`$ along curves $`\beta `$ in $`𝒬`$ from $`q`$ to $`r`$, given by linear mappings $`P_\beta :E_qE_r`$. A section $`\psi `$ of $`E`$ is *parallel* if always $`P_\beta \psi (q)=\psi (r)`$. If $`\beta `$ is a loop, $`q=r`$, the mapping $`P_\beta `$ is called the *holonomy endomorphism* $`h_\beta `$ of $`E_q`$ associated with $`\beta `$. A connection also defines (and is defined by) a covariant derivative operator $`D`$, which allows us to form the derivative $`D\psi `$ of a section $`\psi `$ of $`E`$. A section $`\psi `$ is parallel if and only if $`D\psi =0`$. A bundle with connection is called *flat* if all holonomies of contractible loops are trivial, i.e., the identity endomorphism (this is the case if and only if the curvature of the connection vanishes everywhere).
Parallelity of the local inner product means that parallel transport preserves inner products; equivalently, $`D(\psi ,\varphi )=(D\psi ,\varphi )+(\psi ,D\varphi )`$ for all $`\psi ,\varphi C^{\mathrm{}}(E)`$. It follows in particular that holonomy endomorphisms are always unitary.
Our bundle, the one of which $`\psi `$ is a section, will always be a Hermitian bundle. Note that since a Hermitian bundle consists of a vector bundle and a connection, it can be nontrivial even if the vector bundle is trivial: namely, if the connection is nontrivial. The *trivial Hermitian bundle* $`𝒬\times W`$, in contrast, consists of the trivial vector bundle with the trivial connection, whose parallel transport $`P_\beta `$ is always the identity on $`W`$. The case of a $`W`$-valued function $`\psi :𝒬W`$ corresponds to the trivial Hermitian bundle $`𝒬\times W`$.
The global inner product on the Hilbert space of wave functions is the local inner product integrated against the Riemannian volume measure associated with the metric $`g`$,
$$\varphi ,\psi =_𝒬𝑑q(\varphi (q),\psi (q)).$$
The Hilbert space equipped with this inner product, denoted $`L^2(𝒬,E)`$, contains the square-integrable, measurable (not necessarily smooth) sections of $`E`$ modulo equality almost everywhere. In an obvious sense, $`C_0^{\mathrm{}}(E)L^2(𝒬,E)`$.
The covariant derivative $`D\psi `$ of a section $`\psi `$ is an “$`E`$-valued 1-form,” i.e., a section of $`T𝒬^{}E`$ (with $`T𝒬^{}`$ the cotangent bundle), while we write $`\psi `$ for the section of $`T𝒬E`$ metrically equivalent to $`D\psi `$. To define the covariant derivative of $`D\psi `$, one uses the connection on $`T𝒬^{}E`$ that arises in an obvious way from the Levi-Civita connection on $`T𝒬^{}`$ and the given connection on $`E`$, with the defining property $`D_{T𝒬^{}E}(\omega \psi )=(D_{T𝒬^{}}\omega )\psi +\omega (D_E\psi )`$ for every 1-form $`\omega `$ and every section $`\psi `$ of $`E`$. We take as the Laplacian $`\mathrm{\Delta }\psi `$ of $`\psi `$ the (Riemannian) metric trace of the second covariant derivative of $`\psi `$,
$$\mathrm{\Delta }\psi =g^{ab}D_aD_b\psi ,$$
(29)
where the second $`D`$, the one which is applied first, is the covariant derivative on $`E`$, and the first $`D`$ is the covariant derivative on $`T𝒬^{}E`$.<sup>16</sup><sup>16</sup>16While this is the natural definition of the Laplacian of a section of a Hermitian bundle, we note that for differential $`p`$-forms with $`p1`$ there are two inequivalent natural definitions of the Laplacian: one is $`\mathrm{\Delta }=(d^{}d+dd^{})`$ (sometimes called the de Rham Laplacian, with $`d^{}=(1)^{(dim𝒬)(p+1)+1}d`$ on $`p`$-forms \[20, p. 9\]), the other is (29) for $`E=\mathrm{\Lambda }^pT𝒬^{}`$ (sometimes called the Bochner Laplacian). They differ by a curvature term given by the Weitzenböck formula \[20, p. 11\]. Again, an equivalent definition is $`\mathrm{\Delta }\psi =dd\psi `$, using the obvious extension (based on the connection of $`E`$) of the exterior derivative to $`E`$-valued differential forms, i.e., sections of $`\mathrm{\Lambda }^pT𝒬^{}E`$.
The potential $`V`$ is now a self-adjoint section of the endomorphism bundle $`EE^{}`$ acting on the vector bundle’s fibers.
The equations defining the Bohmian dynamics are the same as before. Explicitly, we define $`v^\psi `$, the Bohmian velocity vector field associated with a wave function $`\psi `$, by
$$v^\psi :=\mathrm{}\mathrm{Im}\frac{(\psi ,\psi )}{(\psi ,\psi )}.$$
(30)
The time evolution of the state $`(Q_t,\psi _t)`$ is given by
$`{\displaystyle \frac{dQ_t}{dt}}`$ $`=v^{\psi _t}(Q_t)`$ (31a)
$`i\mathrm{}{\displaystyle \frac{\psi _t}{t}}`$ $`=\frac{\mathrm{}^2}{2}\mathrm{\Delta }\psi _t+V\psi _t`$ (31b)
The class $`𝒞_0(𝒬,E,V)`$ contains just this one dynamics, defined by (30) and (31). This agrees with the definition of $`𝒞_0(𝒬,W,V)`$ given in Section 5.2 in the sense that $`𝒞_0(𝒬,W,V)=𝒞_0(𝒬,E,V)`$ when $`E`$ is the trivial bundle $`𝒬\times W`$.
Equivariance of the distribution $`\rho =(\psi ,\psi )`$ is (on a formal level) obtained from the equations
$$\frac{\rho _t}{t}=(\rho _tv^{\psi _t})$$
(32)
for the distribution $`\rho _t`$ of the configuration $`Q_t`$ and
$$\frac{|\psi _t|^2}{t}=(|\psi _t|^2v^{\psi _t}),$$
(33)
which follow, since $`(,)`$ is parallel, from (30) and (31) just as (10) and (11) follow from (8) and (9).
## 6 Scalar Periodic Wave Functions on the Covering Space
We introduce now the Bohmian dynamics belonging to the class that we denote $`𝒞_1`$; in Section 8 we introduce the dynamics of class $`𝒞_2`$. To this end, we will consider wave functions on the universal covering space of $`𝒬`$. This idea is rather standard in the literature on quantum mechanics in multiply-connected spaces . However, the standard treatment lacks the precise justification that one can provide in Bohmian mechanics. Moreover, the complete classification of the possibilities that we give in Section 8 includes some, corresponding to what we call *holonomy-twisted representations* of $`\pi _1(𝒬)`$, that until recently had not been considered. The possibilities considered so far correspond to unitary representations of $`\pi _1(𝒬)`$ on the value space of the wave function. Each possibility has locally the same Hamiltonian $`\frac{\mathrm{}^2}{2}\mathrm{\Delta }+V`$, with the same potential $`V`$, and each possibility is equally well defined and equally reasonable. In this section all wave functions will be complex-valued; in Section 8 we consider wave functions with higher-dimensional value spaces.
### 6.1 The Circle, for Example
Let us start with the configuration space $`𝒬=S^1`$, the circle. This space is multiply connected since only those loops that surround the circle as many times clockwise as counterclockwise can be shrunk to a point. It is convenient to write the wave function $`\psi :S^1`$ as a function $`\widehat{\psi }(\theta )`$ of the angle coordinate, with $`\widehat{\psi }:`$ a $`2\pi `$-periodic function. From $`\widehat{\psi }`$ one obtains
$$\widehat{v}^{\widehat{\psi }}=\mathrm{}\mathrm{Im}\frac{\widehat{\psi }}{\widehat{\psi }}$$
(34)
as a $`2\pi `$-periodic function of $`\theta `$. The relevant observation is that for (34) to be $`2\pi `$-periodic, it is (sufficient but) not necessary that $`\widehat{\psi }`$ be $`2\pi `$-periodic. It would be sufficient as well to have a $`\widehat{\psi }`$ that is merely periodic up to a phase shift,
$$\widehat{\psi }(\theta +2\pi )=\gamma \widehat{\psi }(\theta ),$$
(35)
where $`\gamma `$ is a complex constant of modulus one, called a *topological phase factor*.
Another way of viewing this is to write $`\psi `$ in the polar form $`Re^{iS/\mathrm{}}`$, where $`R0`$ and the phase $`S`$ is real, and find that the Bohmian velocity (23) is given by $`v^\psi =S`$. If we view the phase $`S`$ as a function $`\widehat{S}(\theta )`$ of the angle coordinate, we see that $`\widehat{S}`$ will be $`2\pi `$-periodic if
$$\widehat{S}(\theta +2\pi )=\widehat{S}(\theta )+\beta $$
(36)
for some constant $`\beta `$. This corresponds to (35) with $`\gamma :=e^{i\beta /\mathrm{}}`$.
Since $`\widehat{v}^{\widehat{\psi }}`$ is $`2\pi `$-periodic, it makes sense to write the equation of motion
$$\frac{dQ_t}{dt}=\widehat{v}^{\widehat{\psi }}(\theta (Q_t))=\mathrm{}\mathrm{Im}\frac{\widehat{\psi }}{\widehat{\psi }}(\theta (Q_t))$$
(37)
where $`\theta (Q_t)`$ is any of the values of the angle coordinate that one can associate with $`Q_t`$. If we let $`\widehat{\psi }`$ evolve by the Schrödinger equation on the real line with a $`2\pi `$-periodic potential $`V`$,
$$i\mathrm{}\frac{\widehat{\psi }_t}{t}=\frac{\mathrm{}^2}{2}\mathrm{\Delta }\widehat{\psi }_t+V\widehat{\psi }_t,$$
(38)
then the periodicity condition (35) is preserved by the evolution, thanks to the linearity of the Schrödinger equation. Thus, for any fixed complex $`\gamma `$ of modulus one, (37), (38), and (35) together define a Bohmian dynamics, just as (23) and (22) do. This dynamics permits as many different wave functions as the one defined by (23) and (22), which corresponds to $`\gamma =1`$.
Since $`|\gamma |=1`$, so that $`|\widehat{\psi }|^2`$ is $`2\pi `$-periodic, this theory also has an equivariant probability distribution on the circle, with density $`\rho =|\widehat{\psi }|^2`$. This is the reason why we restrict the possibilities to $`\gamma `$ of modulus 1: otherwise we lose equivariance. (The trajectories may still exist globally even if $`|\gamma |1`$.)
We summarize the results of our reasoning.
###### Assertion 1.
For each potential $`V`$ and each complex number $`\gamma `$ of modulus one, there is a Bohmian dynamics on the circle, defined by (35), (37), and (38).
According to the notation that we will define later, these dynamics form the class $`𝒞_1(S^1,V)`$, a one-parameter class parametrized by $`\gamma `$. What are the physical factors that determine which $`\gamma `$ is to be used? It depends. The following subsection provides a concrete example.
### 6.2 Relation to the Aharonov–Bohm Effect
The additional possibilities associated with nontrivial phase shifts $`\gamma `$ occur in physics even in the case of the circle. We describe here a simplified version of the Aharonov–Bohm effect .
Consider a single particle confined to a loop in 3-space. Suppose there is a magnetic field $`𝑩`$ that vanishes at every point of the loop, but with field lines that pass through the interior of the loop. Thus, if $`D`$ is a 2-dimensional surface bounded by the loop, there may be a nonzero flux of the magnetic field across $`D`$,
$$\mathrm{\Phi }:=_D𝑩𝒏𝑑A,$$
(39)
where $`dA`$ is the area element and $`𝒏`$ is the unit normal on the surface $`D`$. (Note that by Maxwell’s equation $`𝑩=0`$ and the Ostrogradski–Gauss integral formula, the value of $`\mathrm{\Phi }`$ does not depend on the particular choice of the surface $`D`$.)
The appropriate quantum or Bohmian theory on $`S^1`$ corresponds, in the sense of Assertion 1, to the phase factor
$$\gamma =e^{ie\mathrm{\Phi }/\mathrm{}},$$
(40)
where $`e`$ in the exponent is the charge of the particle, provided that the orientation of the loop and the surface agree in the sense that the direction of increasing $`\theta `$ and the direction of $`𝒏`$ satisfy a right-hand-rule. A justification of this description will be given in Section 7.1.
From this example we conclude that the dynamics of class $`𝒞_1𝒞_0`$ *do actually occur* in a universe governed by Bohmian mechanics as the effective dynamics in a restricted configuration space.
### 6.3 Notation and Relevant Facts Concerning Covering Spaces
In the remainder of this section, we generalize the considerations of Section 6.1 to arbitrary $`𝒬`$. The relevant notation is summarized by the table:
| | Generic | Simplest example |
| --- | --- | --- |
| configuration space | $`𝒬`$ | $`S^1`$ |
| universal covering space | $`\widehat{𝒬}`$ | $``$ |
| points in $`𝒬`$, $`\widehat{𝒬}`$ | $`q`$, $`\widehat{q}`$ | $`e^{i\theta }`$, $`\theta `$ |
| projection map | $`\pi :\widehat{𝒬}𝒬`$ | $`e^i:S^1`$ |
| covering fiber over $`q`$, $`e^{i\theta }`$ | $`\pi ^1(q)`$ | $`\{\theta +2\pi k|k\}`$ |
| fundamental group of $`𝒬`$ | $`\pi _1(𝒬)`$ | $``$ |
| covering transformation | $`\sigma :\widehat{𝒬}\widehat{𝒬}`$ | $`\sigma _k:\theta \theta +2\pi k`$ |
| covering group | $`Cov(\widehat{𝒬},𝒬)`$ | $`\{\sigma _k|k\}`$ |
| character of fundamental group | $`\gamma _\sigma `$ | $`\gamma _k=\gamma ^k`$ |
| bundle, lifted bundle | $`E`$, $`\widehat{E}`$ | $`S^1\times `$, $`\times `$ |
Again, $`𝒬`$ is a Riemannian manifold with metric $`g`$, with universal covering space denoted by $`\widehat{𝒬}`$. Recall that the universal covering space is, by definition, a simply connected space, endowed with a covering map (a local diffeomorphism) $`\pi :\widehat{𝒬}𝒬`$, also called the projection. The *covering fiber* for $`q𝒬`$ is the set $`\pi ^1(q)`$ of points in $`\widehat{𝒬}`$ that project to $`q`$ under $`\pi `$. Every function or vector field on $`𝒬`$ can be lifted to a function, respectively vector field, on $`\widehat{𝒬}`$. The functions and vector fields on $`\widehat{𝒬}`$ arising in this way are called *projectable*. A function $`f:\widehat{𝒬}`$ is projectable if and only if $`f(\widehat{q})=f(\widehat{r})`$ whenever $`\pi (\widehat{q})=\pi (\widehat{r})`$. In that case it is the lift of $`\stackrel{~}{f}:𝒬`$ given by $`\stackrel{~}{f}(\pi (\widehat{q})):=f(\widehat{q})`$, called the projection of $`f`$. A vector field $`w`$ on $`\widehat{𝒬}`$ is projectable if and only if, whenever $`\pi (\widehat{q})=\pi (\widehat{r})`$, $`\pi ^{}w(\widehat{q})=\pi ^{}w(\widehat{r})`$ where $`\pi ^{}`$ is the (push-forward) action of $`\pi `$ on tangent vectors.
We shall always take $`\widehat{𝒬}`$ to be endowed with the lifted metric $`\widehat{g}`$, which makes $`\widehat{𝒬}`$ a Riemannian manifold as well and assures that $`\pi `$ is a local isometry. As a consequence, if $`\widehat{f}`$ is the lift of the function $`f`$, $`\mathrm{\Delta }_{\widehat{𝒬}}\widehat{f}=\widehat{\mathrm{\Delta }_𝒬f}`$.
A *covering transformation* is an isometry $`\sigma `$ mapping the covering space to itself which preserves the covering fibers, $`\pi \sigma =\pi `$. The group of such transformations is the *covering group* and is denoted by $`Cov(\widehat{𝒬},𝒬)`$. It acts freely and transitively on every covering fiber, i.e., for every $`\widehat{q}`$ and $`\widehat{r}`$ in the same fiber there is precisely one $`\sigma `$ such that $`\widehat{r}=\sigma \widehat{q}`$. As a consequence, projectability of the vector field $`w`$ on $`\widehat{𝒬}`$ is equivalent to the condition that $`w(\sigma \widehat{q})=\sigma ^{}w(\widehat{q})`$ for all $`\widehat{q}\widehat{𝒬}`$ and all $`\sigma Cov(\widehat{𝒬},𝒬)`$.
The *fundamental group at a point* $`q`$, denoted by $`\pi _1(𝒬,q)`$, is the set of equivalence classes of closed loops through $`q`$, where the equivalence relation is that of homotopy, i.e. smoothly deforming one curve into the other. The product in this group is concatenation; more precisely, $`\sigma \tau `$ corresponds to the loop obtained by first following $`\tau `$ and then following $`\sigma `$. (This is in contrast to the common definition of the product, with the opposite order. We do it this way as it seems more natural for parallel transport.) The fundamental groups at different points are isomorphic to each other as well as to the covering group, but the isomorphisms are not canonical. However, for every given $`\widehat{q}\widehat{𝒬}`$ there is a canonical isomorphism $`\phi _{\widehat{q}}:Cov(\widehat{𝒬},𝒬)\pi _1(𝒬,\pi (\widehat{q}))`$; for different choices of $`\widehat{q}`$ in the same fiber, the different $`\phi _{\widehat{q}}`$’s are conjugate: they are related by $`\phi _{\sigma \widehat{q}}(\tau )=\phi _{\widehat{q}}(\sigma ^1\tau \sigma )=\phi _{\widehat{q}}(\sigma )^1\phi _{\widehat{q}}(\tau )\phi _{\widehat{q}}(\sigma )`$. By *the fundamental group of* $`𝒬`$, written $`\pi _1(𝒬)`$, we shall mean any one of the fundamental groups $`\pi _1(𝒬,q)`$.
A *character* of a group $`G`$ is a unitary, 1-dimensional representation of $`G`$, i.e., a homomorphism $`GU(1)`$ where $`U(1)`$ is the multiplicative group of the complex numbers of modulus one. The characters of $`G`$ form a group denoted by $`G^{}`$.
### 6.4 Scalar Periodic Wave Functions
The motion of the configuration $`Q_t`$ in $`𝒬`$ is determined by a velocity vector field $`v_t`$ on $`𝒬`$, which may arise from a wave function $`\psi `$ not on $`𝒬`$ but instead on $`\widehat{𝒬}`$ in the following way.
Suppose we are given a map $`\gamma :Cov(\widehat{𝒬},𝒬)`$, and suppose that a wave function $`\psi :\widehat{𝒬}`$ satisfies the *periodicity condition associated with the topological factors $`\gamma `$*, i.e.,
$$\psi (\sigma \widehat{q})=\gamma _\sigma \psi (\widehat{q})$$
(41)
for every $`\widehat{q}\widehat{𝒬}`$ and $`\sigma Cov(\widehat{𝒬},𝒬)`$. (We no longer put the hat $`\widehat{}`$ on top of $`\psi `$ that served for emphasizing that $`\psi `$ lives on the covering space.) For (41) to be possible for a $`\psi `$ that does not identically vanish, $`\gamma `$ must be a representation of the covering group, as was first emphasized in . To see this, let $`\sigma _1`$, $`\sigma _2Cov(\widehat{𝒬},𝒬)`$. Then we have the following equalities
$$\gamma _{\sigma _1\sigma _2}\psi (\widehat{q})=\psi (\sigma _1\sigma _2\widehat{q})=\gamma _{\sigma _1}\psi (\sigma _2\widehat{q})=\gamma _{\sigma _1}\gamma _{\sigma _2}\psi (\widehat{q}).$$
(42)
We thus obtain the fundamental relation
$$\gamma _{\sigma _1\sigma _2}=\gamma _{\sigma _1}\gamma _{\sigma _2},$$
(43)
establishing (since $`\gamma _{\mathrm{Id}}=1`$) that $`\gamma `$ is a representation.
The 1-dimensional representations of the covering group are, via the canonical isomorphisms $`\phi _{\widehat{q}}:Cov(\widehat{𝒬},𝒬)\pi _1(𝒬,q),\widehat{q}\pi ^1(q)`$, in canonical correspondence with the 1-dimensional representations of any fundamental group $`\pi _1(𝒬,q)`$: The different isomorphisms $`\phi _{\widehat{q}},\widehat{q}\pi ^1(q)`$, will transform a representation of $`\pi _1(𝒬,q)`$ into representations of $`Cov(\widehat{𝒬},𝒬)`$ that are conjugate. But the 1-dimensional representations are homomorphisms to the *abelian* multiplicative group of $``$ and are thus invariant under conjugation.
From (41) it follows that $`\psi (\sigma \widehat{q})=\gamma _\sigma \sigma ^{}\psi (\widehat{q})`$, where $`\sigma ^{}`$ is the (push-forward) action of $`\sigma `$ on tangent vectors, using that $`\sigma `$ is an isometry. Thus, the velocity field $`\widehat{v}^\psi `$ on $`\widehat{𝒬}`$ associated with $`\psi `$ according to
$$\widehat{v}^\psi (\widehat{q}):=\mathrm{}\mathrm{Im}\frac{\psi }{\psi }(\widehat{q})$$
(44)
is projectable, i.e.,
$$\widehat{v}^\psi (\sigma \widehat{q})=\sigma ^{}\widehat{v}^\psi (\widehat{q}),$$
(45)
and therefore gives rise to a velocity field $`v^\psi `$ on $`𝒬`$,
$$v^\psi (q)=\pi ^{}\widehat{v}^\psi (\widehat{q})$$
(46)
where $`\widehat{q}`$ is an arbitrary element of $`\pi ^1(q)`$.
If we let $`\psi `$ evolve according to the Schrödinger equation on $`\widehat{𝒬}`$,
$$i\mathrm{}\frac{\psi }{t}(\widehat{q})=\frac{\mathrm{}^2}{2}\mathrm{\Delta }\psi (\widehat{q})+\widehat{V}(\widehat{q})\psi (\widehat{q})$$
(47)
with $`\widehat{V}`$ the lift of the potential $`V`$ on $`𝒬`$, then the periodicity condition (41) is preserved by the evolution, since, according to
$$i\mathrm{}\frac{\psi }{t}(\sigma \widehat{q})\stackrel{(\text{47})}{=}\frac{\mathrm{}^2}{2}\mathrm{\Delta }\psi (\sigma \widehat{q})+\widehat{V}(\sigma \widehat{q})\psi (\sigma \widehat{q})=\frac{\mathrm{}^2}{2}\mathrm{\Delta }\psi (\sigma \widehat{q})+\widehat{V}(\widehat{q})\psi (\sigma \widehat{q})$$
(48)
(note the different arguments in the potential), the functions $`\psi \sigma `$ and $`\gamma _\sigma \psi `$ satisfy the same evolution equation (47) with, by (41), the same initial condition, and thus coincide at all times.
Therefore we can let the Bohmian configuration $`Q_t`$ move according to $`v^{\psi _t}`$,
$$\frac{dQ_t}{dt}=v^{\psi _t}(Q_t)=\mathrm{}\pi ^{}\left(\mathrm{Im}\frac{\psi }{\psi }\right)(Q_t)=\mathrm{}\pi ^{}\left(\mathrm{Im}\frac{\psi }{\psi }|_{\widehat{q}\pi ^1(Q_t)}\right).$$
(49)
One can also view the motion in this way: Given $`Q_0`$, choose $`\widehat{Q}_0\pi ^1(Q_0)`$, let $`\widehat{Q}_t`$ move in $`\widehat{𝒬}`$ according to $`\widehat{v}^{\psi _t}`$, and set $`Q_t=\pi (\widehat{Q}_t)`$. Then the motion of $`Q_t`$ is independent of the choice of $`\widehat{Q}_0`$ in the fiber over $`Q_0`$, and obeys (49).
If, as we shall assume from now on, $`|\gamma _\sigma |=1`$ for all $`\sigma Cov(\widehat{𝒬},𝒬)`$, i.e., if $`\gamma `$ is a *unitary* representation (in $``$) or a *character*, then the motion (49) also has an equivariant probability distribution, namely
$$\rho (q)=|\psi (\widehat{q})|^2.$$
(50)
To see this, note that we have
$$|\psi (\sigma \widehat{q})|^2\stackrel{(\text{41})}{=}|\gamma _\sigma |^2|\psi (\widehat{q})|^2=|\psi (\widehat{q})|^2,$$
(51)
so that the function $`|\psi (\widehat{q})|^2`$ is projectable to a function on $`𝒬`$ which we call $`|\psi |^2(q)`$ in this paragraph. From (47) we have that
$$\frac{|\psi _t(\widehat{q})|^2}{t}=\left(|\psi _t(\widehat{q})|^2\widehat{v}^{\psi _t}(\widehat{q})\right)$$
and, by projection, that
$$\frac{|\psi _t|^2(q)}{t}=\left(|\psi _t|^2(q)v^{\psi _t}(q)\right),$$
which coincides with the transport equation for a probability density $`\rho `$ on $`𝒬`$,
$$\frac{\rho _t(q)}{t}=\left(\rho _t(q)v^{\psi _t}(q)\right).$$
Hence,
$$\rho _t(q)=|\psi _t|^2(q)$$
(52)
for all times if it is so initially; this is equivariance.
This also makes clear that the relevant wave functions are those with
$$_𝒬𝑑q|\psi (\widehat{q})|^2=1$$
(53)
where the choice of $`\widehat{q}\pi ^1(q)`$ is arbitrary by (51). The relevant Hilbert space, which we denote $`L^2(\widehat{𝒬},\gamma )`$, thus consists of the measurable functions $`\psi `$ on $`\widehat{𝒬}`$ (modulo changes on null sets) satisfying (41) with
$$_𝒬𝑑q|\psi (\widehat{q})|^2<\mathrm{}.$$
(54)
It is a Hilbert space with the scalar product
$$\varphi ,\psi =_𝒬𝑑q\overline{\varphi (\widehat{q})}\psi (\widehat{q}).$$
(55)
Note that the value of the integrand at $`q`$ is independent of the choice of $`\widehat{q}\pi ^1(q)`$ since, by (41) and the fact that $`|\gamma _\sigma |=1`$,
$$\overline{\varphi (\sigma \widehat{q})}\psi (\sigma \widehat{q})=\overline{\gamma _\sigma \varphi (\widehat{q})}\gamma _\sigma \psi (\widehat{q})=\overline{\varphi (\widehat{q})}\psi (\widehat{q}).$$
We summarize the results of our reasoning.
###### Assertion 2.
Given a Riemannian manifold $`𝒬`$ and a smooth function $`V:𝒬`$, there is a Bohmian dynamics in $`𝒬`$ with potential $`V`$ for each character $`\gamma `$ of the fundamental group $`\pi _1(𝒬)`$; it is defined by (41), (47), and (49), where the wave function $`\psi _t`$ lies in $`L^2(\widehat{𝒬},\gamma )`$ and has norm one.
We define $`𝒞_1(𝒬,V)`$ to be the class of Bohmian dynamics provided by Assertion 2. It contains as many elements as there are characters of $`\pi _1(𝒬)`$ because different characters $`\gamma ^{}\gamma `$ always define different dynamics; we give a proof of this fact in .<sup>17</sup><sup>17</sup>17But, essentially, this is already clear for the same reason as why one can, in the Aharonov–Bohm effect, read off the phase shift from a shift in the interference pattern: If one splits a wave packet, located at $`q`$, into two pieces and, say, lets them move along curves $`\beta _1`$ and $`\beta _2`$ from $`q`$ to $`r`$ that are such that the loop $`\beta _1^1\beta _2`$ is incontractible, then one obtains interference between the two packets at $`r`$ in a way that depends on the phase shift associated with the loop $`\beta _1^1\beta _2`$. Interestingly, different characters can define the same dynamics when we consider, instead of complex-valued wave functions, sections of Hermitian bundles. Here is an example of such a nontrivial Hermitian bundle $`E`$: consider $`𝒬={}_{}{}^{N}_{}^{3}`$, whose fundamental group is the permutation group $`S_N`$ with two characters, and the bundle $`E=FB`$ over $`{}_{}{}^{N}_{}^{3}`$, where $`F`$ is what we call the fermionic line bundle (the unique flat Hermitian line bundle over $`{}_{}{}^{N}_{}^{3}`$ whose holonomy representation is the alternating character) and $`B={}_{}{}^{N}_{}^{3}\times `$, the trivial line bundle; see section 7.2 of for a discussion.
### 6.5 Remarks
1. Since the law of motion (49) involves a derivative of $`\psi `$, the merely measurable functions in $`L^2(\widehat{𝒬},\gamma )`$ will of course not be adequate for defining trajectories. However, we will leave aside the question, from which dense subspace of $`L^2(\widehat{𝒬},\gamma )`$ should one choose $`\psi `$.
2. In our example $`𝒬=S^1`$, the possible dynamics are precisely those mentioned in Section 6.1. The covering group is isomorphic to $``$, and every homomorphism $`\gamma `$ is of the form $`\gamma _k=\gamma _1^k`$. Thus a character is determined by a complex number $`\gamma _1`$ of modulus one; the periodicity condition (41) reduces to (35).
3. For the trivial character $`\gamma _\sigma =1`$, we obtain the immediate dynamics, as defined by (23) and (22). Thus, $`𝒞_0(𝒬,V)𝒞_1(𝒬,V)`$.
4. Another example, or application of Assertion 2, is provided by identical particles without spin. The natural configuration space $`{}_{}{}^{N}_{}^{3}`$ for identical particles, defined in (12), has fundamental group $`S_N`$, the group of permutations of $`N`$ objects, which possesses two characters, the trivial character, $`\gamma _\sigma =1`$, and the alternating character, $`\gamma _\sigma =\mathrm{sgn}(\sigma )=1`$ or $`1`$ depending on whether $`\sigma S_N`$ is an even or an odd permutation. As explained in detail in , the Bohmian dynamics associated with the trivial character is that of bosons, while the one associated with the alternating character is that of fermions.
5. When $`|\gamma _\sigma |1`$ for some $`\sigma Cov(\widehat{𝒬},𝒬)`$, in which case the equivariant distribution (50) is not defined, one could think of obtaining instead an equivariant distribution by setting
$$\rho (q)=\underset{\widehat{q}\pi ^1(q)}{}|\psi (\widehat{q})|^2.$$
(56)
However, this ansatz does not work for providing an equivariant distribution in this case. Any $`\sigma `$ for which $`|\gamma _\sigma |1`$ must be an element of infinite order, since otherwise $`\gamma _\sigma `$ would have to be a root of unity. Thus $`\pi _1(𝒬)`$ is infinite, and so is the covering fiber $`\pi ^1(q)`$, which is in a canonical (given $`\widehat{q}\pi ^1(q)`$) correspondence with $`\pi _1(𝒬,q)`$, and the sum on the right hand side of (56) is divergent unless $`\psi `$ vanishes everywhere on this covering fiber. (To see this, note that either $`|\gamma _\sigma |>1`$ or $`|\gamma _{\sigma ^1}|>1`$ since $`\gamma `$ is a representation; without loss of generality we suppose $`|\gamma _\sigma |>1`$. If $`\psi (\widehat{q})0`$ for some $`\widehat{q}`$, then already the sum over just the fiber elements $`\sigma ^k\widehat{q}`$, $`k=1,2,3,\mathrm{}`$, is divergent, since by the periodicity condition (41), $`_k|\psi (\sigma ^k\widehat{q})|^2=_k|\gamma _{\sigma ^k}|^2|\psi (\widehat{q})|^2=|\psi (\widehat{q})|^2_k|\gamma _\sigma |^{2k}=\mathrm{}`$.)
6. As stated already in Section 3, topological factors can also be introduced into GRW theories, provided we start with a GRW theory of the following kind: Wave functions $`\psi `$ are functions on a Riemannian manifold $`𝒬`$, and collapses according to (7) occur with rate (5) with collapse rate operators $`\mathrm{\Lambda }(𝒙)`$ (where $`𝒙`$ is in for example $`^3`$) that are multiplication operators on configuration space:
$$\mathrm{\Lambda }(𝒙)\psi (q)=f_𝒙(q)\psi (q).$$
(57)
Then we may define a GRW theory with topological factor given by the character $`\gamma `$ of $`\pi _1(𝒬)`$ by using wave functions $`\psi `$ on the covering space $`\widehat{𝒬}`$ satisfying the periodicity condition (41) associated with $`\gamma `$, with collapse rate operators the lifted multiplication operators on $`\widehat{𝒬}`$:
$$\mathrm{\Lambda }(𝒙)\psi (\widehat{q})=f_𝒙(\pi (\widehat{q}))\psi (\widehat{q}).$$
(58)
Collapse then maps periodic to periodic wave functions, with the same topological factor. Note, however, that the theory would work as well with aperiodic wave functions.
## 7 The Aharonov–Bohm Effect
We now give a more detailed treatment of the Aharonov–Bohm effect in the framework of Bohmian mechanics and its relation to topological phase factors. In doing so, we repeat various standard considerations on this topic.
### 7.1 Derivation of the Topological Phase Factor
To justify the dynamics described in Section 6.2, we consider a less idealized description of the Aharonov–Bohm effect. Consider a particle moving in $`^3`$ that cannot enter the solid cylinder
$$𝒞=\left\{𝒒=(q_1,q_2,q_3)^3\right|q_1^2+q_2^21\}$$
(59)
because of, say, a potential $`V`$ that goes to $`+\mathrm{}`$ as $`𝒒`$ approaches the cylinder from outside. The effective configuration space $`𝒬=^3𝒞`$ has the same fundamental group $``$ as the circle since it is diffeomorphic to $`S^1\times ^+\times `$ (the fundamental group of a Cartesian product is the direct product of the fundamental groups, and the half plane $`^+\times `$ is simply connected). The magnetic field $`𝑩`$, which vanishes outside $`𝒞`$ but not inside, is included in the equations by means of a vector potential $`𝑨`$ with $`\times 𝑨=𝑩`$:
$`{\displaystyle \frac{dQ_t}{dt}}`$ $`=v^\psi (Q_t)=\mathrm{}\mathrm{Im}{\displaystyle \frac{(\psi ,(\frac{ie}{\mathrm{}}𝑨)\psi )}{(\psi ,\psi )}}(Q_t)`$ (60a)
$`i\mathrm{}{\displaystyle \frac{\psi _t}{t}}`$ $`=\frac{\mathrm{}^2}{2}(\frac{ie}{\mathrm{}}𝑨)^2\psi _t+V\psi _t`$ (60b)
with $`e`$ the charge of the particle. These equations are in fact best regarded as instances of the immediate dynamics (31) on a Hermitian bundle for a nontrivial connection on the trivial vector bundle $`𝒬\times `$, a point of view that we will discuss in Section 7.2.
For now we observe that the vector potential can be gauged away—not on $`𝒬`$ but on the covering space $`\widehat{𝒬}`$ (which is diffeomorphic to $`\times ^+\times `$). More precisely, the dynamics is unaffected by (i) lift to the covering space with $`\gamma =1`$, and (ii) change of gauge.
Concerning (i), note that $`v^\psi `$ is the projection of the vector field
$$\widehat{v}^{\widehat{\psi }}=\mathrm{}\mathrm{Im}\frac{(\widehat{\psi },(\frac{ie}{\mathrm{}}\widehat{𝑨})\widehat{\psi })}{(\widehat{\psi },\widehat{\psi })}$$
(61)
where $`\widehat{\psi }`$ is the lift of $`\psi `$ (and thus a periodic wave function with $`\gamma =1`$) and evolves according to the lift of (60b),
$$i\mathrm{}\frac{\widehat{\psi }_t}{t}=\frac{\mathrm{}^2}{2}(\frac{ie}{\mathrm{}}\widehat{𝑨})^2\widehat{\psi }_t+\widehat{V}\widehat{\psi }_t.$$
(62)
We will now write again $`\psi `$, rather than $`\widehat{\psi }`$, for the wave function on $`\widehat{𝒬}`$.
Concerning (ii), a change of gauge means the simultaneous replacement
$$\widehat{𝑨}(\widehat{q})\widehat{𝑨}^{}(\widehat{q})=\widehat{𝑨}(\widehat{q})+f(\widehat{q}),\psi (\widehat{q})\psi ^{}(\widehat{q})=e^{ief(\widehat{q})/\mathrm{}}\psi (\widehat{q})$$
(63)
for an arbitrary function $`f:\widehat{𝒬}`$. This does not change the dynamics. A vector field with vanishing curl, such as $`𝑨`$ on $`𝒬`$ or $`\widehat{𝑨}`$ on $`\widehat{𝒬}`$, is a gradient in every simply connected region; thus, while $`𝑨`$ is locally but not globally a gradient, $`\widehat{𝑨}`$ is globally a gradient,
$$\widehat{𝑨}=g.$$
(64)
The (not projectable) function $`g:\widehat{𝒬}`$ is given by
$$g(\widehat{q})=_{\widehat{q}_0}^{\widehat{q}}\widehat{𝑨}(\widehat{r})𝑑\widehat{r}+C$$
(65)
for arbitrary $`\widehat{q}_0\widehat{𝒬}`$, where the integration path is an arbitrary curve from $`\widehat{q}_0`$ to $`\widehat{q}`$ and $`C`$ is a constant depending on $`\widehat{q}_0`$. By setting $`f=g`$, we can change the gauge in such a way that $`\widehat{𝑨}^{}`$ vanishes.
However, the change of gauge affects the periodicity of the wave function $`\psi `$: instead of $`\psi (\sigma \widehat{q})=\psi (\widehat{q})`$ we have that
$$\psi ^{}(\sigma \widehat{q})=\gamma _\sigma \psi ^{}(\widehat{q})$$
(66)
with $`\gamma _\sigma =\gamma ^{k(\sigma )}`$, where $`\gamma =\mathrm{exp}(ie\mathrm{\Phi }/\mathrm{})`$ as in (40), with $`\mathrm{\Phi }`$ the magnetic flux given by (39), and $`k(\sigma )`$ is the number of full counterclockwise rotations that the covering transformation $`\sigma `$ induces on $`\widehat{𝒬}`$.
To see this, note first that $`\psi ^{}(\sigma \widehat{q})=\mathrm{exp}(ieg(\sigma \widehat{q})/\mathrm{})\psi (\sigma \widehat{q})=\mathrm{exp}(ieg(\sigma \widehat{q})/\mathrm{})\psi (\widehat{q})`$. Since, by (65),
$$g(\sigma \widehat{q})=_{\widehat{q}_0}^{\widehat{q}}\widehat{𝑨}(\widehat{q})𝑑\widehat{r}+_{\widehat{q}}^{\sigma \widehat{q}}\widehat{𝑨}(\widehat{r})𝑑\widehat{r}+C=_{\widehat{q}}^{\sigma \widehat{q}}\widehat{𝑨}(\widehat{r})𝑑\widehat{r}+g(\widehat{q})$$
for arbitrary integration paths with the indicated end points, we have that
$$\psi ^{}(\sigma \widehat{q})=\mathrm{exp}(i\frac{e}{\mathrm{}}_{\widehat{\alpha }}\widehat{𝑨}(\widehat{r})d\widehat{r})\psi ^{}(\widehat{q})=:\gamma _\sigma \psi ^{}(\widehat{q})$$
(67)
for arbitrary path $`\widehat{\alpha }`$ from $`\widehat{q}`$ to $`\sigma \widehat{q}`$. To evaluate the integral and show that it is independent of $`\widehat{q}`$, note that it agrees with the corresponding integral along the projected path $`\alpha =\pi (\widehat{\alpha })`$, a loop in $`𝒬`$:
$$_{\widehat{\alpha }}\widehat{𝑨}(\widehat{r})𝑑\widehat{r}=_\alpha 𝑨(r)𝑑r,$$
(68)
which depends only on the homotopy class of $`\alpha `$. Now consider for $`\alpha `$ a loop in $`𝒬`$ that surrounds the cylinder once counterclockwise. By the Stokes integral formula, the last integral then agrees with the integral of $`𝑩`$ over any surface $`D`$ in $`^3`$ bounded by $`\alpha `$,
$$_{\widehat{\alpha }}\widehat{𝑨}(\widehat{r})𝑑\widehat{r}=_D𝑩𝒏𝑑A=\mathrm{\Phi }.$$
(69)
This completes the proof.
The dynamics can thus be described without a vector potential by a wave function on $`\widehat{𝒬}`$ satisfying the periodicity condition (66). Ignoring the radial and $`q_3`$ coordinates, keeping only the circle $`S^1`$, yields the model of Section 6.2.
### 7.2 The Bundle View
We re-express our discussion in Section 7.1 of the Aharonov–Bohm effect in terms of Hermitian bundles.
The dynamics that fundamentally takes place in $`^3`$ is the immediate dynamics, given by (31), for the Hermitian bundle $`E^0`$ consisting of the trivial vector bundle $`^3\times `$, the local inner product $`(\varphi (q),\psi (q))_q=\overline{\varphi (q)}\psi (q)`$, and the nontrivial connection whose gradient operator is $`=_{\mathrm{trivial}}\frac{ie}{\mathrm{}}𝑨`$. (The inner product is parallel iff $`\mathrm{Im}𝑨=0`$.) The curvature of this connection is proportional to the magnetic field $`𝑩`$; therefore $`E^0`$ is curved but its restriction $`E=E^0|_𝒬`$ to the effective configuration space $`𝒬=^3𝒞`$ is flat.
Let, for the moment, $`E`$ be any flat Hermitian line bundle over any Riemannian manifold $`𝒬`$. Its lift $`\widehat{E}`$ is a trivial Hermitian bundle, like every flat bundle over a simply connected base manifold. We obtain the same dynamics (as from $`E`$) from periodic sections of $`\widehat{E}`$ with $`\gamma =1`$. However, the description of $`\widehat{E}`$ in which the periodicity condition has trivial topological factor $`\gamma =1`$, namely the description as the lift of $`E`$, is not the description in which the triviality of $`\widehat{E}`$ is manifest, namely the description relative to a trivialization, which may change the topological factor in the periodicity condition as follows.
A trivialization corresponds to a parallel choice of orthonormal basis in every fiber $`\widehat{E}_{\widehat{q}}`$ (i.e., of an identification of $`\widehat{E}_{\widehat{q}}`$ with $``$), and thus to a parallel section $`\varphi `$ of $`\widehat{E}`$ with $`(\varphi _{\widehat{q}},\varphi _{\widehat{q}})_{\widehat{q}}=1`$. Relative to this trivialization, a section $`\psi `$ of $`\widehat{E}`$ corresponds to a function $`\psi ^{}:\widehat{𝒬}`$ according to
$$\psi (\widehat{r})=\psi ^{}(\widehat{r})\varphi (\widehat{r}).$$
(70)
If $`\psi (\sigma \widehat{q})=\psi (\widehat{q})`$ corresponding to $`\gamma =1`$ then $`\psi ^{}`$ satisfies the periodicity condition
$$\psi ^{}(\sigma \widehat{q})=h_\alpha ^1\psi ^{}(\widehat{q}),$$
(71)
where $`\alpha `$ is any loop in $`𝒬`$ based at $`\pi (\widehat{q})`$ whose lift $`\widehat{\alpha }`$ starting at $`\widehat{q}`$ leads to $`\sigma \widehat{q}`$, and $`h_\alpha `$ is the associated holonomy (which in this case, with $`\mathrm{rank}E=1`$, is a complex number of modulus 1). This follows from (70) by parallel transport along $`\widehat{\alpha }`$, using that, by parallelity, $`\varphi (\sigma \widehat{q})=h_\alpha \varphi (\widehat{q})`$.
As a consequence, every dynamics from $`𝒞_0(𝒬,E,V)`$ for a flat Hermitian line bundle $`E`$ exists also in $`𝒞_1(𝒬,V)`$. In other words, we can avoid the use of a nontrivial flat Hermitian line bundle if we use a dynamics of class $`𝒞_1𝒞_0`$ with suitable topological phase factor.
Let us return to the concrete bundle defined in the beginning of this section. It remains to determine the holonomy. For a loop $`\alpha `$ in $`𝒬`$ that surrounds the cylinder once counterclockwise,
$$h_\alpha =\mathrm{exp}\left(\frac{ie}{\mathrm{}}_\alpha 𝑨(r)𝑑r\right).$$
(72)
Since the integral equals, according to our computation in Section 7.1, the magnetic flux $`\mathrm{\Phi }`$, the topological phase factor is given by $`\gamma =\mathrm{exp}(ie\mathrm{\Phi }/\mathrm{})`$.
## 8 Vector-Valued Periodic Wave Functions on the Covering Space
When the wave function is not a scalar but rather a mapping to a vector space $`W`$ of dimension greater than 1, such as for a particle with spin, the topological factors can be matrices, forming a unitary representation of $`\pi _1(𝒬)`$, as we shall derive presently. The more complicated case in which $`\psi _t`$ is a section of a vector bundle is discussed in Section 8.4. The possibility of topological factors given by representations more general than characters was first mentioned in , Notes to Section 23.3.
### 8.1 Vector Spaces
Suppose that the wave functions assume values in a Hermitian vector space $`W`$. Then in a periodicity condition analogous to (41),
$$\psi (\sigma \widehat{q})=\mathrm{\Gamma }_\sigma \psi (\widehat{q}),$$
(73)
we can allow the topological factor $`\mathrm{\Gamma }_\sigma `$ to be an endomorphism $`WW`$, rather than just a complex number. By the same argument as in the scalar case, using that $`\psi (\widehat{q})`$ can be any element of $`W`$, $`\mathrm{\Gamma }`$ must be a representation of $`Cov(\widehat{𝒬},𝒬)`$ on $`W`$.
It follows from (73) that $`\psi (\sigma \widehat{q})=(\sigma ^{}\mathrm{\Gamma }_\sigma )\psi (\widehat{q})`$, where $`\psi (\widehat{q})`$ is viewed as an element of $`T_{\widehat{q}}\widehat{𝒬}W`$. Assume now that, in addition, $`\mathrm{\Gamma }`$ is a *unitary* representation of $`Cov(\widehat{𝒬},𝒬)`$. Then the velocity field $`\widehat{v}^\psi `$ on $`\widehat{𝒬}`$ associated with $`\psi `$ according to
$$\widehat{v}^\psi (\widehat{q}):=\mathrm{}\mathrm{Im}\frac{(\psi ,\psi )}{(\psi ,\psi )}(\widehat{q})$$
(74)
is projectable, $`\widehat{v}^\psi (\sigma \widehat{q})=\sigma ^{}\widehat{v}^\psi (\widehat{q})`$, and gives rise to a velocity field $`v^\psi `$ on $`𝒬`$. (While we used unitarity in the scalar case of Section 6.4 only for obtaining the equivariant probability density, we use it here already for having a projectable velocity field. For this purpose, we could have allowed $`\mathrm{\Gamma }_\sigma `$ to be, rather than unitary, a complex multiple of a unitary endomorphism; unitarity would then be required in order to obtain an equivariant density.)
The potential $`V`$ can now assume values in the Hermitian endomorphisms of $`W`$; the space of endomorphisms can be written $`WW^{}`$, so that $`V`$ is a function $`𝒬WW^{}`$. We let $`\psi `$ evolve according to the Schrödinger equation on $`\widehat{𝒬}`$,
$$i\mathrm{}\frac{\psi }{t}(\widehat{q})=\frac{\mathrm{}^2}{2}\mathrm{\Delta }\psi (\widehat{q})+\widehat{V}(\widehat{q})\psi (\widehat{q}).$$
(75)
The periodicity condition (73) is preserved by the evolution (75) when and only when every $`\mathrm{\Gamma }_\sigma `$ commutes with every $`V(q)`$,
$$\mathrm{\Gamma }_\sigma V(q)=V(q)\mathrm{\Gamma }_\sigma $$
(76)
for all $`\sigma Cov(\widehat{𝒬},𝒬)`$ and all $`q𝒬`$. To see this, note that $`\psi \sigma `$ is a solution of (75) if $`\psi `$ is. Thus (73) is preserved if and only if $`\mathrm{\Gamma }_\sigma \psi `$ satisfies (75), which is the case precisely when multiplication by $`\mathrm{\Gamma }_\sigma `$ commutes with the Hamiltonian. Since it trivially commutes with the Laplacian, the relevant condition is that $`\mathrm{\Gamma }_\sigma `$ commute with the potential $`\widehat{V}(\widehat{q})`$ at every $`\widehat{q}\widehat{𝒬}`$, or, what amounts to the same, with $`V(q)`$ at every $`q𝒬`$.
Given (76), we can let the configuration $`Q_t`$ move according to $`v^{\psi _t}`$,
$$\frac{dQ_t}{dt}=v^{\psi _t}(Q_t)=\mathrm{}\pi ^{}\left(\mathrm{Im}\frac{(\psi ,\psi )}{(\psi ,\psi )}\right)(Q_t).$$
(77)
Since $`\mathrm{\Gamma }`$ is a *unitary* representation of $`Cov(\widehat{𝒬},𝒬)`$, the motion (77) has an equivariant probability distribution, namely
$$\rho (q)=(\psi (\widehat{q}),\psi (\widehat{q})).$$
(78)
The right hand side does not depend on the choice of $`\widehat{q}\pi ^1(q)`$ since, by (73), $`(\psi (\sigma \widehat{q}),\psi (\sigma \widehat{q}))=(\mathrm{\Gamma }_\sigma \psi (\widehat{q}),\mathrm{\Gamma }_\sigma \psi (\widehat{q}))=(\psi (\widehat{q}),\psi (\widehat{q}))`$. Equivariance can be established in the same way as in the scalar case.
We define the Hilbert space $`L^2(\widehat{𝒬},W,\mathrm{\Gamma })`$ to be the set of measurable functions $`\psi :\widehat{𝒬}W`$ (modulo changes on null sets) satisfying (73) with
$$_𝒬𝑑q(\psi (\widehat{q}),\psi (\widehat{q}))<\mathrm{},$$
(79)
endowed with the scalar product
$$\varphi ,\psi =_𝒬𝑑q(\varphi (\widehat{q}),\psi (\widehat{q})).$$
(80)
Again, the value of the integrand at $`q`$ is independent of the choice of $`\widehat{q}\pi ^1(q)`$.
We summarize the results of our reasoning.
###### Assertion 3.
Given a Riemannian manifold $`𝒬`$, a Hermitian vector space $`W`$, and a Hermitian function $`V:𝒬WW^{}`$, there is a Bohmian dynamics for each unitary representation $`\mathrm{\Gamma }`$ of $`Cov(\widehat{𝒬},𝒬)`$ on $`W`$ that commutes with all the endomorphisms $`V(q)`$; it is defined by (73), (75), and (77), where the wave function $`\psi _t`$ lies in $`L^2(𝒬,W,\mathrm{\Gamma })`$ and has norm 1.
We define $`𝒞_2(𝒬,W,V)`$ to be the class of Bohmian dynamics provided by Assertion 3.
The characters $`\gamma `$ of $`Cov(\widehat{𝒬},𝒬)`$ (which are in a canonical one-to-one correspondence with the characters of $`\pi _1(𝒬)`$) are contained in Assertion 3 as special cases of unitary representations $`\mathrm{\Gamma }`$ by setting
$$\mathrm{\Gamma }_\sigma =\gamma _\sigma \mathrm{Id}_W.$$
(81)
These are precisely those unitary representations $`\mathrm{\Gamma }`$ for which all $`\mathrm{\Gamma }_\sigma `$ are multiples of the identity. We define $`𝒞_1(𝒬,W,V)`$ to be the class of those Bohmian dynamics from $`𝒞_2(𝒬,W,V)`$ arising from a unitary representation $`\mathrm{\Gamma }`$ of the form (81), i.e., arising from a character. This class contains as many elements as there are characters of $`\pi _1(𝒬)`$, since different characters define different dynamics; we give a proof of this fact in \[13, Sect. 6.4\] (see also footnote 17). The definition of $`𝒞_1(𝒬,W,V)`$ agrees with that of $`𝒞_1(𝒬,V)`$ given in Section 6.4 in the sense that the latter is the special case $`W=`$, $`𝒞_1(𝒬,V)=𝒞_1(𝒬,,V)`$. Trivially, $`𝒞_0(𝒬,W,V)𝒞_1(𝒬,W,V)𝒞_2(𝒬,W,V)`$.
### 8.2 Remarks
1. The condition that $`\mathrm{\Gamma }`$ be a representation of $`Cov(\widehat{𝒬},𝒬)`$ that commutes with $`V`$ can alternatively be expressed by saying that $`\mathrm{\Gamma }`$ is a homomorphism $`Cov(\widehat{𝒬},𝒬)C(V)`$ where $`C(V)`$ denotes the *centralizer* of $`V`$, i.e., the subgroup of $`U(W)`$ (the unitary group of $`W`$) containing all elements that commute with each $`V(q)`$.
2. The dynamics defined by $`W`$, $`V`$, and $`\mathrm{\Gamma }`$ is the same as the one defined by $`W^{}`$, $`V^{}`$, and $`\mathrm{\Gamma }^{}`$ (another vector space, a potential on $`W^{}`$, and a representation on $`W^{}`$) if there is a unitary isomorphism $`U:WW^{}`$ such that
$$V^{}=UVU^1$$
(82)
and
$$\mathrm{\Gamma }^{}=U\mathrm{\Gamma }U^1.$$
(83)
To see this, define a mapping $`\psi \psi ^{}`$, from $`L^2(\widehat{𝒬},W,\mathrm{\Gamma })`$ to $`L^2(\widehat{𝒬},W^{},\mathrm{\Gamma }^{})`$, by $`\psi ^{}(\widehat{q}):=U\psi (\widehat{q})`$. Here we use that
$$\psi ^{}(\sigma \widehat{q})=U\psi (\sigma \widehat{q})=U\mathrm{\Gamma }_\sigma \psi (\widehat{q})=U\mathrm{\Gamma }_\sigma U^1\psi ^{}(\widehat{q})=\mathrm{\Gamma }_\sigma ^{}\psi ^{}(\widehat{q}).$$
Since $`(\frac{\mathrm{}^2}{2}\mathrm{\Delta }+\widehat{V}^{})\psi ^{}=U(\frac{\mathrm{}^2}{2}\mathrm{\Delta }+\widehat{V})\psi `$, $`U`$ intertwines the time evolutions on $`L^2(\widehat{𝒬},W,\mathrm{\Gamma })`$ and $`L^2(\widehat{𝒬},W^{},\mathrm{\Gamma }^{})`$ based on $`V`$ and $`V^{}`$, i.e., $`(\psi ^{})_t=(\psi _t)^{}`$. Since, moreover, at any fixed time $`\psi ^{}`$ and $`\psi `$ lead to the same probability distribution $`\rho `$ on $`𝒬`$ and to the same velocity fields $`\widehat{v}^\psi ^{}=\widehat{v}^\psi `$ and $`v^\psi ^{}=v^\psi `$, $`\psi ^{}`$ and $`\psi `$ lead to the same trajectories with the same probabilities. That is, the dynamics defined by $`W,V,\mathrm{\Gamma }`$ and the one defined by $`W^{},V^{},\mathrm{\Gamma }^{}`$ are the same.
3. As a consequence of the previous remark, we can use, in Assertion 3, representations of the *fundamental group* $`\pi _1(𝒬)`$ instead of representations of the *covering group* $`Cov(\widehat{𝒬},𝒬)`$.
With a unitary representation $`\stackrel{~}{\mathrm{\Gamma }}`$ of $`\pi _1(𝒬,q)`$ on the vector space $`W`$ (for any $`q`$) there are naturally associated several unitary representations $`\mathrm{\Gamma }(\widehat{q})`$ of $`Cov(\widehat{𝒬},𝒬)`$ on $`W`$, one for each $`\widehat{q}\pi ^1(q)`$, defined by $`\mathrm{\Gamma }(\widehat{q})=\stackrel{~}{\mathrm{\Gamma }}\phi _{\widehat{q}}`$, i.e., by $`\mathrm{\Gamma }_\tau (\widehat{q})=\stackrel{~}{\mathrm{\Gamma }}_{\phi _{\widehat{q}}(\tau )}`$, using the isomorphism $`\phi _{\widehat{q}}:Cov(\widehat{𝒬},𝒬)\pi _1(𝒬,q)`$ introduced in Section 6.3. However, these representations $`\mathrm{\Gamma }(\widehat{q})`$ lead to the same dynamics. To see this, consider two points $`\widehat{q},\widehat{r}\pi ^1(q)`$ with $`\widehat{r}=\sigma \widehat{q},\sigma Cov(\widehat{𝒬},𝒬)`$. Then $`\mathrm{\Gamma }_\tau (\widehat{r})=\stackrel{~}{\mathrm{\Gamma }}_{\phi _{\widehat{r}}(\tau )}=\stackrel{~}{\mathrm{\Gamma }}_{\phi _{\widehat{q}}(\sigma ^1\tau \sigma )}=\mathrm{\Gamma }_{\sigma ^1\tau \sigma }(\widehat{q})=\mathrm{\Gamma }_\sigma (\widehat{q})^1\mathrm{\Gamma }_\tau (\widehat{q})\mathrm{\Gamma }_\sigma (\widehat{q})=U\mathrm{\Gamma }_\tau (\widehat{q})U^1`$ with $`U=\mathrm{\Gamma }_\sigma (\widehat{q})^1`$ a unitary endomorphism of $`W`$. Since $`\stackrel{~}{\mathrm{\Gamma }}`$ commutes with $`V`$ so does $`U`$, and by virtue of Remark 8, $`W,V,\mathrm{\Gamma }(\widehat{q})`$ defines the same dynamics as does $`W,V,\mathrm{\Gamma }(\widehat{r})`$.
4. As a further consequence of Remark 8, corresponding to the case in which $`W^{}=W`$ and $`V^{}=V`$, if $`\mathrm{\Gamma }^{}=U\mathrm{\Gamma }U^1`$ for $`UC(V)`$ (so that $`UVU^1=V`$) then $`W`$, $`V`$, and $`\mathrm{\Gamma }^{}`$ define the same dynamics as $`W`$, $`V`$, and $`\mathrm{\Gamma }`$. Therefore, $`𝒞_2(𝒬,W,V)`$ contains at most as many elements as there are homomorphisms $`\stackrel{~}{\mathrm{\Gamma }}:\pi _1(𝒬)C(V)`$ modulo conjugation by elements $`U`$ of $`C(V)`$.
5. The characters—more precisely, the representations of the form (81)—commute with all endomorphisms of $`W`$, and are thus compatible with *every* potential. All other unitary representations $`\mathrm{\Gamma }`$ are compatible only with *some* potentials. (If $`\mathrm{\Gamma }_\sigma `$ is not a multiple of the identity, then there is a Hermitian endomorphism, which could occur as a $`V(q)`$ for some $`q`$, that does not commute with it.)
6. Moreover, characters are the only representations that commute with a potential $`V`$ when (and, if $`\pi _1(𝒬)`$ has a nontrivial character, only when) the algebra $`\mathrm{Alg}(V(𝒬))`$ generated by the $`V(q)`$ is the full endomorphism algebra $`\mathrm{End}(W)`$ of $`W`$, a condition that is satisfied for a generic potential. Thus, for a generic potential $`V`$, $`𝒞_2(𝒬,W,V)=𝒞_1(𝒬,W,V)`$.
7. Similarly, characters are the only representations that commute with several potentials $`V_1,\mathrm{},V_m`$ when the algebra $`\mathrm{Alg}\left(V_1(𝒬)\mathrm{}V_m(𝒬)\right)`$ generated by the $`V_1(q),\mathrm{},V_m(q)`$ is $`\mathrm{End}(W)`$.
8. Even for $`V`$ such that $`\mathrm{Alg}(V(𝒬))\mathrm{End}(W)`$, generically only characters are necessary. This is because for a generic such $`V`$ there will be a $`q𝒬`$ such that $`V(q)`$ is nondegenerate. Since $`\mathrm{\Gamma }_\sigma `$ must commute with $`V(q)`$, $`\mathrm{\Gamma }_\sigma `$ and $`V(q)`$ must be simultaneously diagonalizable, and if $`V(q)`$ is nondegenerate we have that the representation $`\mathrm{\Gamma }`$ is diagonal, with diagonal entries $`\gamma ^{(i)}`$ given by characters, in the basis $`|iW`$ of eigenvectors of $`V(q)`$. In other words, the representation is of the form
$$\mathrm{\Gamma }_\sigma =\underset{i}{}\gamma _\sigma ^{(i)}P_{W^{(i)}},$$
(84)
where $`P_{W^{(i)}}`$ is the projection onto the $`i`$-th eigenspace $`W^{(i)}=|i`$ of $`V(q)`$. Moreover, when the $`\gamma ^{(i)}`$’s are all different, we then have that every $`V(r)`$ is diagonal in the basis $`|i`$ and the corresponding Schrödinger dynamics, of class $`𝒞_2(𝒬,W,V)`$, can be decomposed into a direct sum of dynamics of class $`𝒞_1(𝒬,W^{(i)},V^{(i)})`$, given by characters, where $`V^{(i)}`$ is the action of $`V`$ on $`W^{(i)}`$. Thus, the set of dynamics corresponding to representations $`\mathrm{\Gamma }`$ of the form (84) could be denoted
$$\underset{i}{}𝒞_1(𝒬,W^{(i)},V^{(i)}).$$
(85)
When the $`\gamma ^{(i)}`$’s are not distinct, a similar decomposition holds, with the sum over $`i`$ replaced by the sum over the distinct characters $`\gamma `$ and with the $`W^{(i)}`$’s replaced by the spans of the $`W^{(i)}`$’s corresponding to the same $`\gamma `$.
9. In the situation described in the previous remark, the representation $`\mathrm{\Gamma }`$ on $`W`$ is reducible, as it clearly is when it is given by a character (unless $`dimW=1`$). In fact, by Schur’s lemma, $`\mathrm{\Gamma }`$ can be irreducible only when the potential $`V`$ is a scalar, i.e., of the form $`V(q)=\stackrel{~}{V}(q)\mathrm{Id}_W`$ with $`\stackrel{~}{V}(q)`$.
10. We have so far considered the possible Bohmian dynamics associated with a configuration space $`𝒬`$, a Hermitian vector space $`W`$ and a Hermitian function $`V:𝒬WW^{}`$, and have argued that we have one such dynamics for each representation $`\mathrm{\Gamma }`$ of $`Cov(\widehat{𝒬},𝒬)`$ that commutes with $`V`$. Let us now consider the class $`𝒞_2(𝒬,W,\mathrm{\Gamma })`$ of possible Bohmian dynamics associated with a configuration space $`𝒬`$, a Hermitian vector space $`W`$, and a representation $`\mathrm{\Gamma }`$ of $`Cov(\widehat{𝒬},𝒬)`$. There is of course one such dynamics for every choice of $`V`$ that commutes with $`\mathrm{\Gamma }`$ but there are more. In fact, in addition to these there is also a dynamics for every Hermitian function $`V^{}:\widehat{𝒬}WW^{}`$ satisfying
$$V^{}(\sigma \widehat{q})=\mathrm{\Gamma }_\sigma V^{}(\widehat{q})\mathrm{\Gamma }_\sigma ^1,$$
(86)
a dynamics involving a potential $`V^{}(\widehat{q})`$ on $`\widehat{𝒬}`$ that need not be the lift of any potential $`V`$ on $`𝒬`$. If $`\mathrm{\Gamma }`$ is given by a character then $`V^{}`$ must in fact be the lift of a potential $`V`$ on $`𝒬`$, $`V^{}(\widehat{q})=V(\pi (q))`$, and we obtain nothing new, but when $`\mathrm{\Gamma }`$ is not given by a character, many new possibilities occur.
### 8.3 Examples
Let us give an example of matrices as topological factors: a higher-dimensional version of the Aharonov–Bohm effect. We may replace the vector potential in the Aharonov–Bohm setting by a non-abelian gauge field (à la Yang–Mills) whose field strength (curvature) vanishes outside the cylinder $`𝒞`$ but not inside; the value space $`W`$ (now corresponding not to spin but to, say, quark color) has dimension greater than one, and the difference between two wave packets that have passed the cylinder $`𝒞`$ on different sides is in general, rather than a phase, a unitary endomorphism $`\mathrm{\Gamma }`$ of $`W`$.
A more practical version is provided by the Aharonov–Casher variant of the Aharonov–Bohm effect, according to which a neutral spin-1/2 particle that carries a magnetic moment $`\mu `$ acquires a nontrivial phase while encircling a charged wire $`𝒞`$. Start with the Dirac equation for a neutral particle with nonzero magnetic moment $`\mu `$ (such as a neutron),
$$i\mathrm{}\gamma ^\mu _\mu \psi =m\psi +\frac{1}{2}\mu F^{\mu \nu }\sigma _{\mu \nu }\psi ,$$
(87)
where $`\psi :^4^4`$, $`\gamma ^\mu `$ are the four Dirac matrices, $`F^{\mu \nu }`$ is the field tensor of the external electromagnetic field, and $`\sigma _{\mu \nu }=\gamma _\mu \gamma _\nu \gamma _\nu \gamma _\mu `$. The last term in (87) should be regarded as phenomenological. Consider now the nonrelativistic limit, in which the wave function assumes values in spin space $`W=^2`$, acted upon by the vector $`𝝈`$ of spin matrices. Suppose that the magnetic field is zero and the electric field $`𝑬`$ is generated by a charge distribution $`\varrho (𝒒)`$ inside $`𝒞`$ which is invariant under translations in the direction $`𝒆^3`$, $`𝒆^2=1`$ in which the wire is oriented. Then the charge per unit length $`\lambda `$ is given by the integral
$$\lambda =_D\varrho (𝒒)𝑑A$$
(88)
over the cross-section disk $`D`$ in any plane perpendicular to $`𝒆`$. The Hamiltonian is
$$H=\frac{\mathrm{}^2}{2m}\left(\frac{i\mu }{\mathrm{}}𝑬\times 𝝈\right)^2\frac{\mu ^2}{m}𝑬^2,$$
(89)
where $`\times `$ denotes the vector product in $`^3`$. This looks like a Hamiltonian $`\frac{\mathrm{}^2}{2}\mathrm{\Delta }+V`$ based on a nontrivial connection $`=_{\mathrm{trivial}}\frac{i\mu }{\mathrm{}}𝑬\times 𝝈`$ on the vector bundle $`^3\times ^2`$. The restriction of this connection, outside of $`𝒞`$, to any plane $`\mathrm{\Sigma }`$ orthogonal to the wire turns out to be flat<sup>18</sup><sup>18</sup>18The curvature is $`\mathrm{\Omega }=d_{\mathrm{trivial}}𝝎+𝝎𝝎`$ with $`𝝎=i\frac{\mu }{\mathrm{}}𝑬\times 𝝈`$. The 2-form $`\mathrm{\Omega }`$ is dual to the vector $`_{\mathrm{trivial}}\times 𝝎+𝝎\times 𝝎=i\frac{\mu }{\mathrm{}}(𝑬)𝝈i\frac{\mu }{\mathrm{}}(𝝈)𝑬2i(\frac{\mu }{\mathrm{}})^2(𝝈𝑬)𝑬.`$ Outside the wire, the first term vanishes and, noting that $`𝑬𝒆=0,`$ the other two terms have vanishing component in the direction of $`𝒆`$ and thus vanish when integrated over any region within an orthogonal plane. so that its restriction to the intersection $`𝒬`$ of $`^3𝒞`$ with the orthogonal plane can be replaced, as in the Aharonov–Bohm case, by the trivial connection if we introduce a periodicity condition on the wave function with the topological factor
$$\mathrm{\Gamma }_1=\mathrm{exp}\left(\frac{4\pi i\mu \lambda }{\mathrm{}}𝒆𝝈\right).$$
(90)
In this way we obtain a representation $`\mathrm{\Gamma }:\pi _1(𝒬)SU(2)`$ that is not given by a character. (For further discussion of the link between gauge connections and topological factors, see .)
Though the $`\mathrm{\Gamma }`$’s are matrices in the above examples, the representation is still abelian since $`\pi _1(𝒬)`$ is an abelian group. To obtain a non-abelian representation, let $`𝒬`$ be $`^3`$ minus two disjoint solid cylinders; its fundamental group is isomorphic to the non-abelian group $``$ where $``$ denotes the free product of groups; it is generated by two loops, $`\sigma _1`$ and $`\sigma _2`$, each surrounding one of the cylinders. Using again non-abelian gauge fields, one can arrange that the matrices $`\mathrm{\Gamma }_{\sigma _i}`$ corresponding to $`\sigma _i`$, $`i=1,2`$, do not commute with each other.
Another example, concerning a system of $`N`$ spin 1/2 fermions: $`𝒬={}_{}{}^{N}_{}^{3}`$ (whose fundamental group is the permutation group $`S_N`$), $`W=_{i=1}^N^2`$, and $`\mathrm{\Gamma }_\sigma =\mathrm{sgn}(\sigma )R_\sigma `$, where $`R_\sigma `$ is the natural action of permutations on the tensor product. The Pauli spin interaction is well defined on $`\widehat{𝒬}`$ (but not on $`𝒬`$) for this $`W`$ (unlike for the natural spin bundle (28)). It is given by
$$V^{}(\widehat{q})=\mu \underset{i}{}𝑩(𝒒_i)𝝈_i$$
(91)
with $`𝝈_i`$ the vector of spin matrices acting on the $`i`$-th component of the tensor product. $`V^{}`$ is not the lift of any potential $`V`$ on $`𝒬`$ (since there is no continuous section of the bundle, over $`𝒬={}_{}{}^{N}_{}^{3}`$, of maps $`q\{1,\mathrm{},N\}`$). This is an example of class $`𝒞_2(𝒬,W,\mathrm{\Gamma })`$, see Remark 16.
### 8.4 Vector Bundles
We now consider wave functions that are sections of vector bundles. The topological factors will be expressed as *periodicity sections*, i.e., parallel unitary sections of the endomorphism bundle indexed by the covering group and satisfying a certain composition law, or, equivalently, as *twisted representations* of $`\pi _1(𝒬)`$.
If $`E`$ is a vector bundle over $`𝒬`$, then the lift of $`E`$, denoted by $`\widehat{E}`$, is a vector bundle over $`\widehat{𝒬}`$; the fiber space at $`\widehat{q}`$ is defined to be the fiber space of $`E`$ at $`q`$, $`\widehat{E}_{\widehat{q}}:=E_q`$, where $`q=\pi (\widehat{q})`$. It is important to realize that with this construction, it makes sense to ask whether $`v\widehat{E}_{\widehat{q}}`$ is equal to $`w\widehat{E}_{\widehat{r}}`$ whenever $`\widehat{q}`$ and $`\widehat{r}`$ are elements of the same covering fiber. Equivalently, $`\widehat{E}`$ is the pull-back of $`E`$ through $`\pi :\widehat{𝒬}𝒬`$. As a particular example, the lift of the tangent bundle of $`𝒬`$ to $`\widehat{𝒬}`$ is canonically isomorphic to the tangent bundle of $`\widehat{𝒬}`$. Sections of $`E`$ or $`EE^{}`$ can be lifted to sections of $`\widehat{E}`$ respectively $`\widehat{E}\widehat{E}^{}`$.
If $`E`$ is a Hermitian vector bundle, then so is $`\widehat{E}`$. The wave function $`\psi `$ that we consider here is a section of $`\widehat{E}`$, so that the $`\widehat{q}`$-dependent Hermitian vector space $`\widehat{E}_{\widehat{q}}`$ replaces the fixed Hermitian vector space $`W`$ of the previous subsection. $`V`$ is a section of the bundle $`EE^{}`$, i.e., $`V(q)`$ is an element of $`E_qE_q^{}`$. To indicate that every $`V(q)`$ is a Hermitian endomorphism of $`E_q`$, we say that $`V`$ is a Hermitian section of $`EE^{}`$.
Since $`\psi (\sigma \widehat{q})`$ and $`\psi (\widehat{q})`$ lie in the same space $`E_q=\widehat{E}_{\widehat{q}}=\widehat{E}_{\sigma \widehat{q}}`$, a periodicity condition can be of the form
$$\psi (\sigma \widehat{q})=\mathrm{\Gamma }_\sigma (\widehat{q})\psi (\widehat{q})$$
(92)
for $`\sigma Cov(\widehat{𝒬},𝒬)`$, where $`\mathrm{\Gamma }_\sigma (\widehat{q})`$ is an endomorphism $`E_qE_q`$. By the same argument as in (42), the condition for (92) to be possible, if $`\psi (\widehat{q})`$ can be any element of $`\widehat{E}_{\widehat{q}}`$, is the composition law
$$\mathrm{\Gamma }_{\sigma _1\sigma _2}(\widehat{q})=\mathrm{\Gamma }_{\sigma _1}(\sigma _2\widehat{q})\mathrm{\Gamma }_{\sigma _2}(\widehat{q}).$$
(93)
Note that this law differs from the one $`\mathrm{\Gamma }(\widehat{q})`$ would satisfy if it were a representation, which reads $`\mathrm{\Gamma }_{\sigma _1\sigma _2}(\widehat{q})=\mathrm{\Gamma }_{\sigma _1}(\widehat{q})\mathrm{\Gamma }_{\sigma _2}(\widehat{q})`$, since in general $`\mathrm{\Gamma }(\sigma \widehat{q})`$ need not be the same as $`\mathrm{\Gamma }(\widehat{q})`$ .
For periodicity (92) to be preserved under the Schrödinger evolution,
$$i\mathrm{}\frac{\psi }{t}(\widehat{q})=\frac{\mathrm{}^2}{2}\mathrm{\Delta }\psi (\widehat{q})+\widehat{V}(\widehat{q})\psi (\widehat{q}),$$
(94)
we need that multiplication by $`\mathrm{\Gamma }_\sigma (\widehat{q})`$ commute with the Hamiltonian. Observe that
$$[H,\mathrm{\Gamma }_\sigma ]\psi (\widehat{q})=\frac{\mathrm{}^2}{2}(\mathrm{\Delta }\mathrm{\Gamma }_\sigma (\widehat{q}))\psi (\widehat{q})\mathrm{}^2(\mathrm{\Gamma }_\sigma (\widehat{q}))(\psi (\widehat{q}))+[\widehat{V}(\widehat{q}),\mathrm{\Gamma }_\sigma (\widehat{q})]\psi (\widehat{q}).$$
(95)
Since we can choose $`\psi `$ such that, for any one particular $`\widehat{q}`$, $`\psi (\widehat{q})=0`$ and $`\psi (\widehat{q})`$ is any element of $`T_{\widehat{q}}\widehat{𝒬}E_q`$ we like, we must have that
$$\mathrm{\Gamma }_\sigma (\widehat{q})=0$$
(96)
for all $`\sigma Cov(\widehat{𝒬},𝒬)`$ and all $`\widehat{q}\widehat{𝒬}`$, i.e., that $`\mathrm{\Gamma }_\sigma `$ is parallel. . Inserting this in (95), the first two terms on the right hand side vanish. Since we can choose for $`\psi (\widehat{q})`$ any element of $`E_q`$ we like, we must have that
$$[\widehat{V}(\widehat{q}),\mathrm{\Gamma }_\sigma (\widehat{q})]=0$$
(97)
for all $`\sigma Cov(\widehat{𝒬},𝒬)`$ and all $`\widehat{q}\widehat{𝒬}`$. Conversely, assuming (96) and (97), we obtain that $`\mathrm{\Gamma }_\sigma `$ commutes with $`H`$ for every $`\sigma Cov(\widehat{𝒬},𝒬)`$, so that the periodicity (92) is preserved.
From (92) and (96) it follows that $`\psi (\sigma \widehat{q})=(\sigma ^{}\mathrm{\Gamma }_\sigma (\widehat{q}))\psi (\widehat{q})`$. If every $`\mathrm{\Gamma }_\sigma (\widehat{q})`$ is *unitary*, as we assume from now on, the velocity field $`\widehat{v}^\psi `$ on $`\widehat{𝒬}`$ associated with $`\psi `$ according to
$$\widehat{v}^\psi (\widehat{q}):=\mathrm{}\mathrm{Im}\frac{(\psi ,\psi )}{(\psi ,\psi )}(\widehat{q})$$
(98)
is projectable, $`\widehat{v}^\psi (\sigma \widehat{q})=\sigma ^{}\widehat{v}^\psi (\widehat{q})`$, and gives rise to a velocity field $`v^\psi `$ on $`𝒬`$. We let the configuration move according to $`v^{\psi _t}`$,
$$\frac{dQ_t}{dt}=v^{\psi _t}(Q_t)=\mathrm{}\pi ^{}\left(\mathrm{Im}\frac{(\psi ,\psi )}{(\psi ,\psi )}\right)(Q_t).$$
(99)
###### Definition 2.
Let $`E`$ be a Hermitian bundle over the manifold $`𝒬`$. A *periodicity section* $`\mathrm{\Gamma }`$ over $`E`$ is a family indexed by $`Cov(\widehat{𝒬},𝒬)`$ of unitary parallel sections $`\mathrm{\Gamma }_\sigma `$ of $`\widehat{E}\widehat{E}^{}`$ satisfying the composition law (93).
Since $`\mathrm{\Gamma }_\sigma (\widehat{q})`$ is unitary, one sees as before that the probability distribution
$$\rho (q)=(\psi (\widehat{q}),\psi (\widehat{q}))$$
(100)
does not depend on the choice of $`\widehat{q}\pi ^1(q)`$ and is equivariant.
As usual, we define for any periodicity section $`\mathrm{\Gamma }`$ the Hilbert space $`L^2(\widehat{𝒬},\widehat{E},\mathrm{\Gamma })`$ to be the set of measurable sections $`\psi `$ of $`\widehat{E}`$ (modulo changes on null sets) satisfying (92) with
$$_𝒬𝑑q(\psi (\widehat{q}),\psi (\widehat{q}))<\mathrm{},$$
(101)
endowed with the scalar product
$$\varphi ,\psi =_𝒬𝑑q(\varphi (\widehat{q}),\psi (\widehat{q})).$$
(102)
As before, the value of the integrand at $`q`$ is independent of the choice of $`\widehat{q}\pi ^1(q)`$.
We summarize the results of our reasoning.
###### Assertion 4.
Given a Hermitian bundle $`E`$ over the Riemannian manifold $`𝒬`$ and a Hermitian section $`V`$ of $`EE^{}`$, there is a Bohmian dynamics for each periodicity section $`\mathrm{\Gamma }`$ commuting (pointwise) with $`\widehat{V}`$ (cf. (97)); it is defined by (92), (94), and (99), where the wave function $`\psi _t`$ lies in $`L^2(\widehat{𝒬},\widehat{E},\mathrm{\Gamma })`$ and has norm 1.
The situation of Section 8.1, where the wave function assumed values in a fixed Hermitian space $`W`$ instead of a bundle, corresponds to the trivial Hermitian bundle $`E=𝒬\times W`$ (i.e., with the trivial connection, for which parallel transport is the identity on $`W`$). Then, parallelity (96) implies that $`\mathrm{\Gamma }_\sigma (\widehat{r})=\mathrm{\Gamma }_\sigma (\widehat{q})`$ for any $`\widehat{r},\widehat{q}\widehat{𝒬}`$, or $`\mathrm{\Gamma }_\sigma (\widehat{q})=\mathrm{\Gamma }_\sigma `$, so that (93) becomes the usual composition law $`\mathrm{\Gamma }_{\sigma _1\sigma _2}=\mathrm{\Gamma }_{\sigma _1}\mathrm{\Gamma }_{\sigma _2}`$. As a consequence, $`\mathrm{\Gamma }`$ is a unitary representation of $`Cov(\widehat{𝒬},𝒬)`$ , and Assertion 3 is a special case of Assertion 4.
Every character $`\gamma `$ of $`Cov(\widehat{𝒬},𝒬)`$ (or of $`\pi _1(𝒬)`$) defines a periodicity section by setting
$$\mathrm{\Gamma }_\sigma (\widehat{q}):=\gamma _\sigma \mathrm{Id}_{\widehat{E}_{\widehat{q}}}.$$
(103)
It commutes with every potential $`V`$. Conversely, a periodicity section $`\mathrm{\Gamma }`$ that commutes with every potential must be such that every $`\mathrm{\Gamma }_\sigma (\widehat{q})`$ is a multiple of the identity, $`\mathrm{\Gamma }_\sigma (\widehat{q})=\gamma _\sigma (\widehat{q})\mathrm{Id}_{\widehat{E}_{\widehat{q}}}`$. By unitarity, $`|\gamma _\sigma |=1`$; by parallelity (96), $`\gamma _\sigma (\widehat{q})=\gamma _\sigma `$ must be constant; by the composition law (93), $`\gamma `$ must be a homomorphism, and thus a character.
We define $`𝒞_2(𝒬,E,V)`$ to be the class of Bohmian dynamics provided by Assertion 4. We define $`𝒞_1(𝒬,E,V)`$ to be the class of those Bohmian dynamics from $`𝒞_2(𝒬,E,V)`$ arising from characters: The class $`𝒞_1(𝒬,E,V)`$ contains at most<sup>19</sup><sup>19</sup>19For nontrivial Hermitian bundles, different characters can lead to the same dynamics; we gave an example in footnote 17. as many elements as there are characters of $`\pi _1(𝒬)`$. These definitions agree with the definitions of $`𝒞_1(𝒬,W,V)`$ and $`𝒞_2(𝒬,W,V)`$ given in Section 8.1 in the sense that $`𝒞_1(𝒬,W,V)=𝒞_1(𝒬,E,V)`$ and $`𝒞_2(𝒬,W,V)=𝒞_2(𝒬,E,V)`$ when $`E`$ is taken to be the trivial bundle $`𝒬\times W`$.
We briefly indicate how a periodicity section $`\mathrm{\Gamma }`$ corresponds to something like a representation of $`\pi _1(𝒬)`$, in fact to a (holonomy-) twisted representation of $`\pi _1(𝒬)`$. Fix a $`\widehat{q}\widehat{𝒬}`$. Then $`Cov(\widehat{𝒬},𝒬)`$ can be identified with $`\pi _1(𝒬)=\pi _1(𝒬,\pi (\widehat{q}))`$ via $`\phi _{\widehat{q}}`$. Since the sections $`\mathrm{\Gamma }_\sigma `$ of $`\widehat{E}\widehat{E}^{}`$ are parallel, $`\mathrm{\Gamma }_\sigma (\widehat{r})`$ is determined for every $`\widehat{r}`$ by $`\mathrm{\Gamma }_\sigma (\widehat{q})`$. (Note in particular that the parallel transport $`\mathrm{\Gamma }_\sigma (\tau \widehat{q})`$ of $`\mathrm{\Gamma }_\sigma (\widehat{q})`$ from $`\widehat{q}`$ to $`\tau \widehat{q},\tau Cov(\widehat{𝒬},𝒬)`$, may differ from $`\mathrm{\Gamma }_\sigma (\widehat{q})`$.) Thus, the periodicity section $`\mathrm{\Gamma }`$ is completely determined by the endomorphisms $`\mathrm{\Gamma }_\sigma :=\mathrm{\Gamma }_\sigma (\widehat{q})`$ of $`E_q`$, $`\sigma Cov(\widehat{𝒬},𝒬)`$, which satisfy the composition law
$$\mathrm{\Gamma }_{\sigma _1\sigma _2}=h_{\alpha _2}\mathrm{\Gamma }_{\sigma _1}h_{\alpha _2}^1\mathrm{\Gamma }_{\sigma _2},$$
(104)
where $`\alpha _2`$ is any loop in $`𝒬`$ based at $`\pi (\widehat{q})`$ whose lift starting at $`\widehat{q}`$ leads to $`\sigma _2\widehat{q}`$, and $`h_{\alpha _2}`$ is the associated holonomy endomorphism of $`E_q`$. Since (104) is not the composition law $`\mathrm{\Gamma }_{\sigma _1\sigma _2}=\mathrm{\Gamma }_{\sigma _1}\mathrm{\Gamma }_{\sigma _2}`$ of a representation, the $`\mathrm{\Gamma }_\sigma `$ form, not a representation of $`\pi _1(𝒬)`$, but what we call a *twisted representation*. See for further discussion.
### 8.5 Further Remarks
1. The dynamics defined by $`E`$, $`V`$, and $`\mathrm{\Gamma }`$ is the same as the one defined by $`E^{}`$, $`V^{}`$, and $`\mathrm{\Gamma }^{}`$ (another Hermitian bundle, a potential on $`E^{}`$, and a periodicity section over $`E^{}`$) if there is a unitary parallel section $`U`$ of $`\widehat{E}^{}\widehat{E}^{}`$ such that
$$\widehat{V}^{}(\widehat{q})=U(\widehat{q})\widehat{V}(\widehat{q})U(\widehat{q})^1$$
(105)
and
$$\mathrm{\Gamma }_\sigma ^{}(\widehat{q})=U(\sigma \widehat{q})\mathrm{\Gamma }_\sigma (\widehat{q})U(\widehat{q})^1.$$
(106)
To see this, define a mapping $`\psi \psi ^{}`$, from $`L^2(\widehat{𝒬},\widehat{E},\mathrm{\Gamma })`$ to $`L^2(\widehat{𝒬},\widehat{E}^{},\mathrm{\Gamma }^{})`$, by $`\psi ^{}(\widehat{q}):=U(\widehat{q})\psi (\widehat{q})`$. Here we use that
$$\psi ^{}(\sigma \widehat{q})=U(\sigma \widehat{q})\psi (\sigma \widehat{q})=U(\sigma \widehat{q})\mathrm{\Gamma }_\sigma (\widehat{q})\psi (\widehat{q})$$
$$=U(\sigma \widehat{q})\mathrm{\Gamma }_\sigma (\widehat{q})U(\widehat{q})^1\psi ^{}(\widehat{q})=\mathrm{\Gamma }_\sigma ^{}(\widehat{q})\psi ^{}(\widehat{q}).$$
Since, by the parallelity of $`U`$, $`(\frac{\mathrm{}^2}{2}\mathrm{\Delta }+\widehat{V}^{})\psi ^{}=U(\frac{\mathrm{}^2}{2}\mathrm{\Delta }+\widehat{V})\psi `$, $`U`$ intertwines the time evolutions on $`L^2(\widehat{𝒬},\widehat{E},\mathrm{\Gamma })`$ and $`L^2(\widehat{𝒬},\widehat{E}^{},\mathrm{\Gamma }^{})`$ based on $`V`$ and $`V^{}`$, i.e., $`(\psi ^{})_t=(\psi _t)^{}`$. Since, moreover, at any fixed time $`\psi ^{}`$ and $`\psi `$ lead to the same probability distribution $`\rho `$ on $`𝒬`$ (by the unitarity of $`U`$) and to the same velocity fields $`\widehat{v}^\psi ^{}=\widehat{v}^\psi `$ and $`v^\psi ^{}=v^\psi `$ (by the parallelity and unitarity of $`U`$), $`\psi ^{}`$ and $`\psi `$ lead to the same trajectories with the same probabilities. That is, the dynamics defined by $`E,V,\mathrm{\Gamma }`$ and the one defined by $`E^{},V^{},\mathrm{\Gamma }^{}`$ are the same.
2. As a consequence of the previous remark, we can use, in Assertion 4, twisted representations $`\stackrel{~}{\mathrm{\Gamma }}`$ of the fundamental group $`\pi _1(𝒬)`$, satisfying
$$\stackrel{~}{\mathrm{\Gamma }}_{\alpha _1\alpha _2}=h_{\alpha _2}\stackrel{~}{\mathrm{\Gamma }}_{\alpha _1}h_{\alpha _2}^1\stackrel{~}{\mathrm{\Gamma }}_{\alpha _2},$$
(107)
instead of periodicity sections (twisted representations of the covering group $`Cov(\widehat{𝒬},𝒬)`$).
3. As a further consequence of Remark 17, corresponding to the case in which $`E^{}=E`$ and $`V^{}=V`$, if $`\mathrm{\Gamma }_\sigma ^{}(\widehat{q})=U(\sigma \widehat{q})\mathrm{\Gamma }_\sigma (\widehat{q})U(\widehat{q})^1`$ for a unitary parallel section $`U`$ of $`\widehat{E}^{}\widehat{E}^{}`$ that commutes with $`\widehat{V}`$ then $`E`$, $`V`$, and $`\mathrm{\Gamma }^{}`$ define the same dynamics as $`E`$, $`V`$, and $`\mathrm{\Gamma }`$. Therefore, $`𝒞_2(𝒬,E,V)`$ contains at most as many elements as there are twisted representations $`\stackrel{~}{\mathrm{\Gamma }}`$ of $`\pi _1(𝒬)`$, i.e., periodicity sections $`\mathrm{\Gamma }`$ over $`E`$, that commute with $`\widehat{V}`$, modulo conjugation by such $`U`$’s.
4. As we have already seen, the characters—the periodicity sections of the form (103)—are compatible with every potential, and all other periodicity sections are compatible only with some potentials.
5. A potential $`V`$ does not commute with any periodicity section save the characters when (and, if $`\pi _1(𝒬)`$ has a nontrivial character, only when) for arbitrary $`q𝒬`$,
$$\mathrm{Alg}\left(V(𝒬)_q\mathrm{\Theta }_q\right)=\mathrm{End}(E_q),$$
(108)
where $`V(𝒬)_q=\{P_\beta ^1V(r)P_\beta :r𝒬,\beta \text{ a curve from }q\text{ to }r\}`$ with $`P_\beta :E_qE_r`$ denoting parallel transport, and $`\mathrm{\Theta }_q=\{h_\alpha :\alpha `$ a contractible loop based at q} with $`h_\alpha =P_\alpha `$, the holonomy of $`\alpha `$. This follows from the fact that a periodicity section, by parallelity, must commute with $`\mathrm{\Theta }_q`$. The condition (108) holds, for example, for the potential occurring in the Pauli equation for $`N`$ identical particles with spin,
$$V(q)=\mu \underset{𝒒q}{}𝑩(𝒒)𝝈_𝒒$$
(109)
on the spin bundle (28) over $`{}_{}{}^{N}_{}^{3}`$, with $`𝝈_𝒒`$ the vector of spin matrices acting on the spin space of the particle at $`𝒒`$, provided merely that the magnetic field $`𝑩`$ is not parallel. Thus, for a generic potential $`V`$, $`𝒞_2(𝒬,E,V)=𝒞_1(𝒬,E,V)`$.
6. A periodicity section $`\mathrm{\Gamma }`$ defining a Bohmian dynamics of class $`𝒞_2(𝒬,E,V)`$ can be irreducible only when the potential $`V`$ is a scalar. (When $`\mathrm{\Gamma }`$ is reducible, its decomposition may involve sub-bundles $`E^{(i)}`$ of $`\widehat{E}`$ that are not the lifts of any sub-bundles of $`E`$.)
7. Consider the class $`𝒞_2(𝒬,E,\mathrm{\Gamma })`$ of possible Bohmian dynamics associated with a Riemannian manifold $`𝒬`$, a Hermitian bundle $`E`$ over $`𝒬`$, and a periodicity section $`\mathrm{\Gamma }`$ over $`E`$. There is one such dynamics for every choice of Hermitian section $`V^{}`$ of $`\widehat{E}\widehat{E}^{}`$ satisfying
$$V^{}(\sigma \widehat{q})=\mathrm{\Gamma }_\sigma (\widehat{q})V^{}(\widehat{q})\mathrm{\Gamma }_\sigma (\widehat{q})^1,$$
(110)
a dynamics involving a potential $`V^{}(\widehat{q})`$ on $`\widehat{𝒬}`$ that need not be the lift of any potential $`V`$ on $`𝒬`$.
8. For a *generic* (curved) Hermitian bundle $`E`$, and any fixed $`V`$ we have that $`𝒞_2(𝒬,E,V)=𝒞_1(𝒬,E,V)`$; in other words, there are no more possibilities than the characters. This follows from the fact that, generically, for every unitary endomorphism $`U`$ of $`E_q`$ there is a contractible curve $`\alpha `$ based at $`q`$ whose holonomy is $`U`$. That is, $`\mathrm{\Theta }_q`$ is the full unitary group of $`E_q`$, and thus, by (108), all periodicity sections correspond to characters.
9. Consider a dynamics of class $`𝒞_2(𝒬,W,\mathrm{\Gamma })`$ or class $`𝒞_2(𝒬,E,\mathrm{\Gamma })`$, given by a potential $`V^{}`$ on $`\widehat{𝒬}`$ satisfying (86), respectively (110). We show in that there is a Hermitian bundle $`E^{}`$ over $`𝒬`$ that is locally isomorphic to $`𝒬\times W`$, respectively $`E`$, such that this dynamics, of class $`𝒞_2`$, coincides with the dynamics of class $`𝒞_0(𝒬,E^{},V)`$, i.e., the immediate dynamics for $`𝒬`$, $`E^{}`$, and a potential $`V`$ on $`𝒬`$. For example, the dynamics associated with $`𝒬={}_{}{}^{N}_{}^{3}`$, $`W=_{i=1}^N^{2s+1}`$, and $`V^{}`$ on $`\widehat{𝒬}`$ given by (91) coincides with the dynamics associated with $`{}_{}{}^{N}_{}^{3}`$, the natural spin bundle $`E^{}`$ (28) over $`{}_{}{}^{N}_{}^{3}`$, and the Pauli interaction $`V`$ on $`𝒬`$ given by (109). (By “dynamics” here we refer to the evolution of the configuration.)
### 8.6 Examples Involving Vector Bundles
We close this section with two examples of topological factors on vector bundles.
The most important example is provided by identical particles with spin. In fact, for this case, Assertion 4 entails the same conclusions we arrived at in Remark 4, the alternative between bosons and fermions, even for particles with spin. To understand how this comes about, consider the potential occurring in the Pauli equation (109) for $`N`$ identical particles with spin, on the spin bundle (28) over $`{}_{}{}^{N}_{}^{3}`$, and observe that the algebra generated by $`\{V(q)\}`$ arising from all possible choices of the magnetic field $`𝑩`$ is $`\mathrm{End}(E_q)`$. Thus the only holonomy-twisted representations that define a dynamics for all magnetic fields (or even for a single magnetic field provided only that it is not parallel, see Remark 21) are those given by a character.
Our last example involves a holonomy-twisted representation $`\mathrm{\Gamma }`$ that is not a representation in the ordinary sense. Consider $`N`$ fermions, each as in the examples at the beginning of Section 8.3, moving in $`M=^3_i𝒞_i`$, where $`𝒞_i`$ are one or more disjoint solid cylinders. More generally, consider $`N`$ fermions, each having 3-dimensional configuration space $`M`$ and value space $`W`$ (which may incorporate spin or “color” or both). Then the configuration space $`𝒬`$ for the $`N`$ fermions is the set $`{}_{}{}^{N}M`$ of all $`N`$-element subsets of $`M`$, with universal covering space $`\widehat{𝒬}=\widehat{{}_{}{}^{N}M}=\widehat{M}^N\mathrm{\Delta }`$ with $`\mathrm{\Delta }`$ the extended diagonal, the set of points in $`\widehat{M}^N`$ whose projection to $`M^N`$ lies in its coincidence set. Every diffeomorphism $`\sigma Cov(\widehat{{}_{}{}^{N}M},{}_{}{}^{N}M)`$ can be expressed as a product
$$\sigma =p\stackrel{~}{\sigma }$$
(111)
where $`pS_N`$ and $`\stackrel{~}{\sigma }=(\sigma ^{(1)},\mathrm{},\sigma ^{(N)})Cov(\widehat{M},M)^N`$ and these act on $`\widehat{q}=(\widehat{}𝒒_1,\mathrm{},\widehat{}𝒒_N)`$ $`\widehat{M}^N`$ as follows:
$$\stackrel{~}{\sigma }\widehat{q}=(\sigma ^{(1)}\widehat{}𝒒_1,\mathrm{},\sigma ^{(N)}\widehat{}𝒒_N)$$
(112)
and
$$p\widehat{q}=(\widehat{}𝒒_{p^1(1)},\mathrm{},\widehat{}𝒒_{p^1(N)}).$$
(113)
Thus
$$\sigma \widehat{q}=(\sigma ^{(p^1(1))}\widehat{}𝒒_{p^1(1)},\mathrm{},\sigma ^{(p^1(N))}\widehat{}𝒒_{p^1(N)}).$$
(114)
Moreover, the representation (111) of $`\sigma `$ is unique. Thus, since
$$\sigma _1\sigma _2=p_1\stackrel{~}{\sigma }_1p_2\stackrel{~}{\sigma }_2=(p_1p_2)(p_2^1\stackrel{~}{\sigma }_1p_2\stackrel{~}{\sigma }_2)$$
(115)
with $`p_2^1\stackrel{~}{\sigma }_1p_2=(\sigma _1^{(p_2(1))},\mathrm{},\sigma _1^{(p_2(N))})Cov(\widehat{M},M)^N`$, we find that $`Cov(\widehat{{}_{}{}^{N}M},{}_{}{}^{N}M)`$ is a semidirect product of $`S_N`$ and $`Cov(\widehat{M},M)^N`$, with product given by
$$\sigma _1\sigma _2=(p_1,\stackrel{~}{\sigma }_1)(p_2,\stackrel{~}{\sigma }_2)=(p_1p_2,p_2^1\stackrel{~}{\sigma }_1p_2\stackrel{~}{\sigma }_2).$$
(116)
Wave functions for the $`N`$ fermions are sections of the lift $`\widehat{E}`$ to $`\widehat{𝒬}`$ of the bundle $`E`$ over $`𝒬`$ with fiber
$$E_q=\underset{𝒒q}{}W$$
(117)
and (nontrivial) connection inherited from the trivial connection on $`M\times W`$. If the dynamics for $`N=1`$ involves wave functions on $`\widehat{M}`$ obeying (92) with topological factor $`\mathrm{\Gamma }_\sigma (\widehat{𝒒})=\mathrm{\Gamma }_\sigma `$ given by a unitary representation of $`\pi _1(M)`$ (i.e., independent of $`\widehat{𝒒}`$), then the $`N`$ fermion wave function obeys (92) with topological factor
$$\mathrm{\Gamma }_\sigma (\widehat{q})=\mathrm{sgn}(p)\underset{𝒒\pi (\widehat{q})}{}\mathrm{\Gamma }_{\sigma ^{(i_{\widehat{q}}(𝒒))}}\mathrm{sgn}(p)\mathrm{\Gamma }_{\stackrel{~}{\sigma }}(\widehat{q})$$
(118)
where for $`\widehat{q}=(\widehat{}𝒒_1,\mathrm{},\widehat{}𝒒_N),\pi (\widehat{q})=\{\pi _M(\widehat{}𝒒_1),\mathrm{},\pi _M(\widehat{}𝒒_N)\}`$ and $`i_{\widehat{q}}(\pi _M(\widehat{}𝒒_j))=j`$. Since
$$\mathrm{\Gamma }_{\stackrel{~}{\sigma }_1\stackrel{~}{\sigma }_2}(\widehat{q})=\mathrm{\Gamma }_{\stackrel{~}{\sigma }_1}(\widehat{q})\mathrm{\Gamma }_{\stackrel{~}{\sigma }_2}(\widehat{q})$$
(119)
we find, using (116) and (119), that
$`\mathrm{\Gamma }_{\sigma _1\sigma _2}(\widehat{q})`$ $`=\mathrm{sgn}(p_1p_2)\mathrm{\Gamma }_{p_2^1\stackrel{~}{\sigma }_1p_2\stackrel{~}{\sigma }_2}(\widehat{q})`$ (120a)
$`=\mathrm{sgn}(p_1)\mathrm{\Gamma }_{p_2^1\stackrel{~}{\sigma }_1p_2}(\widehat{q})\mathrm{sgn}(p_2)\mathrm{\Gamma }_{\stackrel{~}{\sigma }_2}(\widehat{q})`$ (120b)
$`=P_2\mathrm{\Gamma }_{\sigma _1}(\widehat{q})P_2^1\mathrm{\Gamma }_{\sigma _2}(\widehat{q}),`$ (120c)
which agrees with (104) since the holonomy on the bundle $`E`$ is given by permutations $`P`$ acting on the tensor product (117).
## 9 The Character Quantization Principle
We have seen that for a Riemannian manifold $`𝒬`$ that is multiply connected, there are additional possibilities for a Bohmian dynamics beyond the usual ones. These new possibilities correspond to (twisted) representations of $`\pi _1(𝒬)`$, the most important of which are given by the characters. In fact, unless the potential $`V`$ is very special, the characters are the only representations that define a possible dynamics involving that potential.
We summarize our discussion so far, invoking the special status of the characters, in the
Character Quantization Principle. Consider a quantum system whose configuration space is given by the Riemannian manifold $`𝒬`$ and whose value space for the wave function is given by the the Hermitian vector space $`W`$ \[or the Hermitian bundle $`E`$ over $`𝒬`$\]. Then for every character $`\gamma `$ of the fundamental group $`\pi _1(𝒬)`$, there is a family $`_\gamma =\{_\gamma (V)\}`$ of Bohmian dynamics, one for each potential, i.e., Hermitian function $`V:𝒬WW^{}`$ \[or Hermitian section $`V`$ of $`EE^{}`$\]. The dynamics $`_\gamma (V)`$ associated with the potential $`V`$ can be taken to be defined by
$`\psi (\sigma \widehat{q})`$ $`=\gamma _\sigma \psi (\widehat{q}),`$ (121a)
$`i\mathrm{}{\displaystyle \frac{\psi }{t}}(\widehat{q})`$ $`=\frac{\mathrm{}^2}{2}\mathrm{\Delta }\psi (\widehat{q})+\widehat{V}(\widehat{q})\psi (\widehat{q}),`$ (121b)
$`{\displaystyle \frac{dQ_t}{dt}}`$ $`=\mathrm{}\pi ^{}\left(\mathrm{Im}{\displaystyle \frac{(\psi ,\psi )}{(\psi ,\psi )}}\right)(Q_t)`$ (121c)
with $`\psi L^2(\widehat{𝒬},W,\gamma )`$ \[or $`\psi L^2(\widehat{𝒬},\widehat{E},\gamma )`$\].
Equations (121) are identical with (41), (94), and (99). Recall that the characters of $`\pi _1(𝒬)`$ are canonically identified with those of $`Cov(\widehat{𝒬},𝒬)`$.
The Character Quantization Principle corresponds to the symmetrization postulate for the case of $`N`$ identical particles in $`^3`$: in this case the natural configuration space $`𝒬={}_{}{}^{N}_{}^{3}`$ and the fundamental group $`\pi _1(𝒬)=S_N`$, the group of permutations of $`N`$ elements, which has two characters, the trivial character, corresponding to bosons, and the alternating character, corresponding to fermions.
We now wish to elaborate upon why the theories given by characters deserve special attention, and are, arguably, the only possibilities for theories that can be regarded as fundamental. There are at least four crucial considerations: (i) freedom, (ii) genericity, (iii) theoretical stability, and (iv) irreducibility.
It seems within our power to expose a physical system, for example a system of $`N`$ identical particles (which has the multiply-connected configuration space $`{}_{}{}^{N}_{}^{3}`$), to a wide variety of potentials. As we have noted already in Remarks (11), (12), (13), \[and (20)\], if we can arrange any potential we like, or if the potentials we can arrange are sufficient to generate together the algebra $`\mathrm{End}(W)`$ \[respectively $`\mathrm{End}(E_q)`$\], then the (twisted) representations defining the dynamics must be given by a character. For example, as we show in , if we can expose a system of $`N`$ identical particles to arbitrary magnetic fields $`𝑩`$ then the potentials (109) on the natural spin bundle (28) over $`{}_{}{}^{N}_{}^{3}`$, which occur in the Pauli equation, generate $`\mathrm{End}(E_q)`$.
The second consideration is based on the hypothesis that, to the extent that a Hamiltonian defining a fundamental physical theory can be regarded as a Schrödinger operator $`\frac{\mathrm{}^2}{2}\mathrm{\Delta }+V`$, the potential $`V`$ is rather generic, or at least not too special. But generically we have that $`\mathrm{Alg}(V(𝒬))=\mathrm{End}(W)`$ \[and that $`\mathrm{Alg}(V(𝒬)_q\mathrm{\Theta }_q)=\mathrm{End}(E_q)`$\]. It then follows, as we have pointed out in Remarks 12 and 21, that the dynamics belongs to $`𝒞_1`$. And even systems that we can describe to a very good degree of approximation by special (e.g., scalar) potentials \[and, if appropriate, special (e.g., flat) Hermitian bundles\] then still cannot have a dynamics from $`𝒞_2𝒞_1`$.
The third consideration concerns the stability of the theory under perturbations. The idea is that the theoretical description of a system (such as, again, $`N`$ identical particles) should not be so delicately contrived as to make sense only for a single potential $`V`$, but should also make sense for all potentials close to $`V`$. (One reason why one might require theoretical stability is the idea that our theoretical descriptions may be idealized, for example when we take physical space to be Euclidean $`^3`$ or a magnetic field to vanish, neglecting small perturbations.) This implies that the theory should be well defined for a generic potential, allowing only dynamics of class $`𝒞_1`$.
Finally, and perhaps most importantly, it seems reasonable to demand of a fundamental physical theory that it be suitably irreducible. But it follows from Remark (12) \[respectively (21)\] that a (twisted) representation can fail to be given by a character only when $`\mathrm{Alg}(V(𝒬))\mathrm{End}(W)`$ \[respectively when $`\mathrm{Alg}(V(𝒬)_q\mathrm{\Theta }_q)\mathrm{End}(E_q)`$\], and in this case the Schrödinger dynamics is decomposable into a direct sum of dynamics corresponding to subspaces of the value-space $`W`$ \[or to sub-bundles\]. One might wonder, in this case, why the full value-space \[or bundle\] was involved to begin with.
These considerations are of course related. For example, a generic potential clearly corresponds to an irreducible dynamics. Freedom relies on genericity in the following way. Since our one universe has in fact just one Hamiltonian and thus just one potential $`V=V_{\mathrm{univ}}`$, what must be meant when one speaks of exposing a system to various potentials $`V_{\mathrm{sys}}`$ is that
$$V_{\mathrm{sys}}(q_{\mathrm{sys}})=V_{\mathrm{univ}}(q_{\mathrm{sys}},Q_{\mathrm{env}}),$$
(122)
for all configurations $`q_{\mathrm{sys}}`$ of the system, where $`Q_{\mathrm{env}}`$ is the actual configuration of the environment of the system (i.e., the rest of the universe), which we can control to a certain extent. In words, $`V_{\mathrm{sys}}`$ is the *conditional potential* of a subsystem of the universe. Thus, the diversity of potentials that we can arrange for a system is inherited from the richness of the potential of the universe: if $`V_{\mathrm{univ}}`$ were scalar, we would be unable to arrange potentials $`V_{\mathrm{sys}}`$ other than scalars. Thus, the origin of the freedom of potentials must lie in genericity. On the other hand, freedom, since it requires the genericity hypothesis, lends support to it.
## 10 Conclusions
We have studied the possible quantum theories on a topologically nontrivial configuration space $`𝒬`$ from the point of view of Bohmian mechanics, which is fundamentally concerned with the motion of matter in physical space, represented by the evolution of a point in configuration space.
Our goal was to find, define, and classify all Bohmian dynamics in $`𝒬`$, where the wave functions may be sections of a Hermitian vector bundle $`E`$. What “all” Bohmian dynamics means is not obvious; we have followed one approach to what it can mean; other approaches are described in . The present approach uses wave functions $`\psi `$ that are defined on the universal covering space $`\widehat{𝒬}`$ of $`𝒬`$ and satisfy a periodicity condition ensuring that the Bohmian velocity vector field on $`\widehat{𝒬}`$ defined in terms of $`\psi `$ can be projected to $`𝒬`$. We have arrived in this way at two natural classes $`𝒞_1𝒞_2`$ of Bohmian dynamics beyond the immediate Bohmian dynamics. A dynamics from $`𝒞_1`$ is defined by a potential and some topological information encoded in a character (one-dimensional unitary representation) of the fundamental group of the configuration space, $`\pi _1(𝒬)`$. A dynamics from $`𝒞_2`$ is defined by a potential and a more general algebraic-geometrical object, a “periodicity section” $`\mathrm{\Gamma }`$.
The dynamics of $`𝒞_2𝒞_1`$ exist only for special potentials. Those of $`𝒞_1`$, however, are compatible with *every* potential, as one would desire for what could be considered a version of quantum mechanics in $`𝒬`$. We have thus arrived at the known fact that for every character of $`\pi _1(𝒬)`$ there is a version of quantum mechanics in $`𝒬`$; we have formulated this in terms of Bohmian mechanics as the “Character Quantization Principle.” A consequence, which will be discussed in detail in a sister paper , is the symmetrization postulate for identical particles. These different quantum theories emerge naturally when one contemplates the possibilities for defining a Bohmian dynamics in $`𝒬`$.
## Acknowledgments
We thank Kai-Uwe Bux (Cornell University), Frank Loose (Eberhard-Karls-Universität Tübingen, Germany) and Penny Smith (Lehigh University) for helpful discussions.
R.T. gratefully acknowledges support by the German National Science Foundation (DFG) through its Priority Program “Interacting Stochastic Systems of High Complexity”, by INFN, and by the European Commission through its 6th Framework Programme “Structuring the European Research Area” and the contract Nr. RITA-CT-2004-505493 for the provision of Transnational Access implemented as Specific Support Action. N.Z. gratefully acknowledges support by INFN. The work of S. Goldstein was supported in part by NSF Grant DMS-0504504.
Finally, we appreciate the hospitality that some of us have enjoyed, on more than one occasion, at the Mathematisches Institut of Ludwig-Maximilians-Universität München (Germany), the Dipartimento di Fisica of Università di Genova (Italy), the Institut des Hautes Études Scientifiques in Bures-sur-Yvette (France), and the Mathematics Department of Rutgers University (USA). |
warning/0506/hep-ph0506266.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The pseudoscalar-meson decay constants (e.g., $`f_D`$, $`f_{D_s}`$, $`f_B`$ and $`f_{B_s}`$) play an important role in extracting the CKM matrix elements (e.g., the leptonic decay width of $`D_s^+l^+\nu _l`$ is proportional to $`f_{D_s}^2|V_{cs}|^2`$), which are crucial for testing the flavor sector of the standard model via the unitarity of CKM matrix. Experimentally<sup>1</sup><sup>1</sup>1See, for example, Refs. , and other experimental results complied by PDG ., precise determination of $`f_{D_s}`$ and $`f_D`$ will soon result from the high-statistics program of CLEO-c, however, the determination of $`f_B`$ and $`f_{B_s}`$ remains beyond the reach of current experiments.
Theoretically, lattice QCD provides a solid framework to compute the masses and decay constants of pseudoscalar mesons (as well as other physical observables) nonperturbatively from the first principles of QCD. Thus reliable lattice QCD determinations of $`f_B`$ and $`f_{B_s}`$ are of fundamental importance, in view of their experimental determinations are still lacking. Obviously, the first step for lattice QCD is to check whether the lattice determinations of $`f_D`$ and $`f_{D_s}`$ will agree with those coming soon from the high-statistics charm program of CLEO-c. This motivates our present study.
In this paper, we compute quenched quark propagators for 30 quark masses in the range $`0.03m_qa0.80`$, in the framework of optimal domain-wall fermion proposed by Chiu -. Then we determine the inverse lattice spacing $`a^1=2.237(76)`$ GeV from the pion time-correlation function, with the experimental input of pion decay constant $`f_\pi =132`$ MeV. The strange quark bare mass $`m_sa=0.08`$ and the charm quark bare mass $`m_ca=0.80`$ are fixed such that the corresponding masses extracted from the vector meson correlation function agree with $`\varphi (1020)`$ and $`J/\psi (3097)`$ respectively. Then the masses and decay constants of any hadrons containing $`c,s`$, and $`u(d)`$ quarks<sup>2</sup><sup>2</sup>2In this paper, we work in the isospin limit $`m_u=m_d`$. are predictions of QCD from the first principles, with the understanding that chiral extrapolation to physical $`m_{u,d}m_s/25`$ (or equivalently $`m_\pi =135`$ MeV) is required for any observables containing $`u(d)`$ quarks.
For pseudoscalar and vector mesons, we measure their time correlation functions for the following three categories: (i) two quarks have the same mass; (ii) one quark mass is fixed at $`m_s`$; (iii) one quark mass is fixed at $`m_c`$. Note that for mesons which are composed of strange and/or charm quarks, their masses and decay constants can be measured directly without chiral extrapolation.
The outline of this paper is as follows. In section 2, we outline our formulation of exact chiral symmetry on the lattice, and our computation of quark propagators. In section 3, we determine the inverse lattice spacing spacing, the strange quark bare mass, and the charm quark bare mass. In section 4, we present our results of $`m_K`$ and $`f_K`$. In section 5, we present our results of $`m_D`$, $`m_{D_s}`$, $`f_D`$, and $`f_{D_s}`$. In section 6, we summarize our results and conclude with some remarks.
## 2 Lattice quarks with exact chiral symmetry
To implement exact chiral symmetry on the lattice , we consider the optimal domain-wall fermion proposed by Chiu -. The action of optimal domain-wall fermion can be written as
$`𝒜_F`$ $`=`$ $`{\displaystyle \underset{s,s^{}=0}{\overset{N_s+1}{}}}{\displaystyle \underset{x,x^{}}{}}\overline{\psi }(x,s)\{(\omega _sD_w(x,x^{})+\delta _{x,x^{}})\delta _{ss^{}}`$ (1)
$`+(\omega _sD_w(x,x^{})\delta _{x,x^{}})(P_+\delta _{s^{},s1}+P_{}\delta _{s^{},s+1})\}\psi _(x^{},s^{})`$
with boundary conditions
$`P_+\psi (x,1)`$ $`=`$ $`rm_qP_+\psi (x,N_s+1),`$
$`P_{}\psi (x,N_s+2)`$ $`=`$ $`rm_qP_{}\psi (x,0),r={\displaystyle \frac{1}{2m_0}},`$
where $`m_q`$ is the bare quark mass, and $`\{\omega _s,s=0,\mathrm{},N_s+1\}`$ are specified by the exact formula derived in Ref. . Here $`H_w=\gamma _5D_w`$, and $`D_w`$ is the standard Wilson Dirac operator plus a negative parameter $`m_0`$ ($`0<m_0<2`$). The quark fields are constructed from the boundary modes at $`s=0`$ and $`s=N_s+1`$ with $`\omega _0=\omega _{N_s+1}=0`$ :
$`q(x)`$ $`=`$ $`\sqrt{r}\left[P_{}\psi (x,0)+P_+\psi (x,N_s+1)\right],`$ (2)
$`\overline{q}(x)`$ $`=`$ $`\sqrt{r}\left[\overline{\psi }(x,0)P_++\overline{\psi }(x,N_s+1)P_{}\right].`$ (3)
After introducing pseudofermions with $`m_q=2m_0`$, the generating functional for $`n`$-point Green’s function of the quark fields can be derived as ,
$`Z[J,\overline{J}]={\displaystyle \frac{[dU]e^{𝒜_g}\text{det}[(D_c+m_q)(1+rD_c)^1]\mathrm{exp}\left\{\overline{J}(D_c+m_q)^1J\right\}}{[dU]e^{𝒜_g}\text{det}[(D_c+m_q)(1+rD_c)^1]}}`$ (4)
where $`𝒜_g`$ is the action of the gauge fields, $`\overline{J}`$ and $`J`$ are the Grassman sources of $`q`$ and $`\overline{q}`$ respectively, and
$`D_c`$ $`=`$ $`2m_0{\displaystyle \frac{1+\gamma _5S_{opt}}{1\gamma _5S_{opt}}},`$ (5)
$`S_{opt}`$ $`=`$ $`{\displaystyle \frac{1\underset{s=1}{\overset{N_s}{}}T_s}{1+_{s=1}^{N_s}T_s}},`$ (6)
$`T_s`$ $`=`$ $`{\displaystyle \frac{1\omega _sH_w}{1+\omega _sH_w}}.`$ (7)
Using the exact formula of $`\omega _s`$ , one immediately obtains
$`S_{opt}=\{\begin{array}{cc}H_wR_Z^{(n,n)}(H_w^2),\hfill & N_s=2n+1,\hfill \\ H_wR_Z^{(n1,n)}(H_w^2),\hfill & N_s=2n,\hfill \end{array}`$ (10)
where $`R_Z(H_w^2)`$ is the Zolotarev optimal rational polynomial for the inverse square root of $`H_w^2`$,
$`R_Z^{(n,n)}(H_w^2)`$ $`=`$ $`{\displaystyle \frac{d_0}{\lambda _{min}}}{\displaystyle \underset{l=1}{\overset{n}{}}}{\displaystyle \frac{1+h_w^2/c_{2l}}{1+h_w^2/c_{2l1}}}`$ (11)
$`=`$ $`{\displaystyle \frac{1}{\lambda _{min}}}(h_w^2+c_{2n}){\displaystyle \underset{l=1}{\overset{n}{}}}{\displaystyle \frac{b_l}{h_w^2+c_{2l1}}},h_w^2=H_w^2/\lambda _{min}^2`$
and
$`R_Z^{(n1,n)}(H_w^2)={\displaystyle \frac{d_0^{}}{\lambda _{min}}}{\displaystyle \frac{\underset{l=1}{\overset{n1}{}}(1+h_w^2/c_{2l}^{})}{_{l=1}^n(1+h_w^2/c_{2l1}^{})}}={\displaystyle \frac{1}{\lambda _{min}}}{\displaystyle \underset{l=1}{\overset{n}{}}}{\displaystyle \frac{b_l^{}}{h_w^2+c_{2l1}^{}}},`$ (12)
where the coefficients $`d_0`$, $`d_0^{}`$, $`c_l`$ and $`c_l^{}`$ are expressed in terms of elliptic functions with arguments depending on $`N_s`$ and $`\lambda _{max}^2/\lambda _{min}^2`$, and $`\lambda _{min}`$ ($`\lambda _{max}`$) is fixed to be the greatest lower bound (least upper bound) of the eigenvalues of $`|H_w|`$ for the set of gauge configurations under investigation.
From (4), the effective 4D lattice Dirac operator for the fermion determinant is
$`D(m_q)=(D_c+m_q)(1+rD_c)^1=m_q+(m_0m_q/2)\left[1+\gamma _5H_wR_Z(H_w^2)\right]`$ (13)
and the quark propagator in background gauge field is
$`q(x)\overline{q}(y)`$ $`=`$ $`{\displaystyle \frac{\delta ^2Z[J,\overline{J}]}{\delta \overline{J}(x)\delta J(y)}}|_{J=\overline{J}=0}`$ (14)
$`=`$ $`(D_c+m_q)_{x,y}^1=(1rm_q)^1[D_{x,y}^1(m_q)r\delta _{x,y}]`$
Note that $`D_c`$ is exactly chirally symmetric (i.e. $`D_c\gamma _5+\gamma _5D_c=0`$) in the limit $`N_s\mathrm{}`$, and its deviation from exact chiral symmetry due to finite $`N_s`$ is the minimal provided that the weights $`\{\omega _s\}`$ are fixed according to the formula derived in Ref. . Further, the bare quark mass $`m_q`$ (whether heavy or light) in the quark propagator $`(D_c+m_q)^1`$ is well-defined for any gauge configurations.
In practice, we have two ways to evaluate the quark propagator (14) in background gauge field:
(i) To solve the linear system of the 5D optimal DWF operator;
(ii) To solve $`D_{x,y}^1(m_q)`$ from the system
$`D(m_q)Y=\left[m_q+(m_0m_q/2)\left(1+\gamma _5H_wR_Z(H_w^2)\right)\right]Y=\text{1I},`$ (15)
with nested conjugate gradient , and then substitute the solution vector $`Y`$ into (14).
Since either (i) or (ii) yields exactly the same quark propagator, in principle, it does not matter which linear system one actually solves. However, in practice, one should choose the most efficient scheme for one’s computational system (hardware and software). For our present system (a Linux PC cluster of 100 nodes ), it has been geared to the scheme (ii), and it attains the maximum efficiency if the inner conjugate gradient loop of (15) is iterated with Neuberger’s 2-pass algorithm . So we use the scheme (ii) to compute the quark propagator, with the quark fields (2)-(3) defined by the boundary modes at $`s=0`$ and $`s=N_s+1`$. Note that Neuberger’s 2-pass algorithm not only provides very high precision of chiral symmetry with fixed amount of memory, but also is faster than the single pass algorithm for $`n>1225`$ (where $`n`$ is the order of the rational polynomial $`R^{(n1,n)}`$) for most computer platforms, as discussed by Chiu and Hsieh .
We generate 100 gauge configurations with single plaquette gauge action at $`\beta =6.1`$ on the $`20^3\times 40`$ lattice. Fixing $`m_0=1.3`$, we project out 16 low-lying eigenmodes of $`|H_w|`$ and perform the nested conjugate gradient in the complement of the vector space spanned by these eigenmodes. For $`N_s=128`$, the weights $`\{\omega _s\}`$ are fixed with $`\lambda _{min}=0.18`$ and $`\lambda _{max}=6.3`$, where $`\lambda _{min}\lambda (|H_w|)\lambda _{max}`$ for all gauge configurations. For each configuration, point to point quark propagators are computed for 30 bare quark masses in the range $`0.03m_qa0.8`$, with stopping criteria $`10^{11}`$ and $`2\times 10^{12}`$ for the outer and inner conjugate gradient loops respectively. Then the norm of the residual vector of each column of the quark propagator is less than $`2\times 10^{11}`$
$`(D_c+m_q)Y\text{1I}<2\times 10^{11},`$
and the chiral symmetry breaking due to finite $`N_s`$ is less than $`10^{14}`$,
$`\sigma =\left|{\displaystyle \frac{Y^{}S^2Y}{Y^{}Y}}1\right|<10^{14},`$
for every iteration of the nested conjugate gradient. Further details of our scheme have been described in Refs. .
In this paper, we measure the time-correlation functions for pseudoscalar ($`PS`$) and vector ($`V`$) mesons,
$`C_{PS}(t)`$ $`=`$ $`{\displaystyle \underset{\stackrel{}{x}}{}}\text{tr}\{\gamma _5(D_c+m_Q)_{x,0}^1\gamma _5(D_c+m_q)_{0,x}^1\}_U`$ (16)
$`C_V(t)`$ $`=`$ $`{\displaystyle \frac{1}{3}}{\displaystyle \underset{\mu =1}{\overset{3}{}}}{\displaystyle \underset{\stackrel{}{x}}{}}\text{tr}\{\gamma _\mu (D_c+m_Q)_{x,0}^1\gamma _\mu (D_c+m_q)_{0,x}^1\}_U`$ (17)
where the subscript $`U`$ denotes averaging over gauge configurations. Here $`C_{PS}(t)`$ and $`C_V(t)`$ are measured for the following three categories:
(i) Symmetric masses $`m_Q=m_q`$ ,
(ii) Asymmetric masses with fixed $`m_Q=m_s=0.08a^1`$ ,
(iii) Asymmetric masses with fixed $`m_Q=m_c=0.80a^1`$ ,
where $`m_q`$ is varied for 30 masses in the range $`0.03m_qa0.80`$.
## 3 Determination of $`a^1`$, $`m_s`$, and $`m_c`$
For symmetric masses $`m_Q=m_q`$, the pseudoscalar time-correlation function $`C_\pi (t)`$ (16) is measured, and is fitted to the usual formula
$`{\displaystyle \frac{Z}{2m_\pi a}}[e^{m_\pi at}+e^{m_\pi a(Tt)}]`$ (18)
to extract the pion mass $`m_\pi a`$ and the pion decay constant
$`f_\pi a=2m_qa{\displaystyle \frac{\sqrt{Z}}{m_\pi ^2a^2}}.`$ (19)
In Figs. 1 and 2, we plot the decay constant $`f_\pi a`$ and pion mass square $`(m_\pi a)^2`$ versus bare quark mass $`m_qa`$, respectively.
The data of $`f_\pi a`$ (see Fig. 1) is well fitted by the straight line
$`f_\pi a=0.059(2)+0.235(38)\times (m_qa).`$
Then taking $`f_\pi a`$ at $`m_qa=0`$ equal to $`0.132`$ GeV times the lattice spacing $`a`$, we can determine the lattice spacing $`a`$ and its inverse,
$`a^1`$ $`=`$ $`{\displaystyle \frac{0.132}{f_0}}\text{ GeV}=2.237(76)\text{ GeV},`$
$`a`$ $`=`$ $`0.088(3)\text{ fm}.`$ (20)
Thus the size of our lattice is about $`(1.8\text{ fm})^3\times 3.6\text{ fm}`$. Since the smallest pion mass is $`439\text{ MeV}`$, the lattice size is about $`(3.9)^3\times 7.8`$, in units of the Compton wavelength ($`0.45\text{ fm}`$) of the smallest pion mass.
The data of $`m_\pi ^2`$ (see Fig. 2) can be fitted by the form
$`m_\pi ^2a^2=A_1(m_qa)^{\frac{1}{1+\delta }}+B(m_qa)^2`$ (21)
in quenched chiral perturbation theory (q$`\chi `$PT). The fitted parameters are
$`\delta `$ $`=`$ $`0.187(21)`$ (22)
$`A_1`$ $`=`$ $`0.669(45)`$ (23)
$`B`$ $`=`$ $`2.666(357)`$ (24)
with $`\chi ^2`$/d.o.f.=0.54. Evidently, the coefficient of quenched chiral logarithm $`\delta =0.187(21)`$ is in good agreement with the theoretical estimate $`\delta 0.176`$ in q$`\chi `$PT.
The bare mass of strange quark is determined by extracting the mass of vector meson from the time-correlation function $`C_V(t)`$. At $`m_qa=0.08`$, $`m_Va=0.460(4)`$, which gives $`m_V=1029(10)`$ MeV, in good agreement with the mass of $`\varphi (1020)`$. Thus we take the strange quark bare mass to be $`m_sa=0.08`$. Similarly, at $`m_qa=0.80`$, $`m_Va=1.368(2)`$, which gives $`m_V=3060(5)`$ MeV, in good agreement with the mass of $`J/\mathrm{\Psi }(3097)`$. Thus, we fix the charm quark bare mass to be $`m_ca=0.80`$.
## 4 $`f_K`$ and $`m_K`$
Next we measure the time-correlation function of kaon $`C_K(t)`$ (16) with $`m_Q`$ fixed at $`m_s=0.08a^1`$, while $`m_q`$ is varied for 30 masses in the range $`0.03m_qa0.80`$. Then the data of $`C_K(t)`$ is fitted by the formula analogous to (18) to extract the kaon mass $`m_Ka`$ and the kaon decay constant $`f_Ka`$.
In Fig. 3, the kaon mass $`m_K`$ is plotted versus $`m_\pi `$, for 15 quark masses in the range $`0.03m_qa0.10`$. The data of $`m_Ka`$ can be fitted by
$`m_Ka=0.197(1)+0.255(4)(m_\pi a)+0.389(8)(m_\pi a)^2.`$
At the physical limit $`m_\pi =135`$ MeV, it gives $`m_K=478(16)`$ MeV, in good agreement with the experimental value of kaon mass ($`495`$ MeV).
In Fig. 4, $`f_Ka`$ is plotted versus bare quark mass $`m_qa`$. The data is well fitted by the straight line
$`f_Ka=0.068(0)+0.116(1)\times (m_qa)`$
At $`m_qa=0`$, it gives $`f_K=152(6)`$ MeV, in agreement with the value $`f_{K^+}=159.8\pm 1.4\pm 0.44`$ MeV complied by PDG .
## 5 $`f_D`$, $`f_{D_s}`$, $`m_D`$, and $`m_{D_s}`$
Now we turn to charmed pseudoscalar mesons. We measure the time-correlation function $`C_D(t)`$ (16) with $`m_Q`$ fixed at $`m_c=0.80a^1`$, while $`m_q`$ is varied for 30 different masses in the range $`0.03m_qa0.80`$. Then the data of $`C_D(t)`$ is fitted by the formula analogous to (18) to extract the mass $`m_Da`$ and decay constant $`f_Da`$.
In Fig. 5, $`m_Da`$ is plotted versus $`m_\pi a`$, for 15 quark masses in the range $`0.03m_qa0.10`$. The data of $`m_Da`$ can be fitted by
$`m_Da=0.816(0)+0.101(3)(m_\pi a)+0.298(6)(m_\pi a)^2`$
At $`m_\pi =135`$ MeV, it gives $`m_D=1842(15)`$ MeV, in good agreement with the mass of $`D`$ meson ($`1865`$ MeV). In Fig. 6, the decay constant $`f_Da`$ is plotted versus bare quark mass $`m_qa`$. The data is well fitted by the straight line
$`f_Da=0.105(1)+0.172(1)\times (m_qa)`$
At $`m_qa=0`$, it gives $`f_D=235(8)`$ MeV, which serves as a prediction of lattice QCD with exact chiral symmetry.
The pseudoscalar meson of $`c\overline{s}`$ or $`s\overline{c}`$ corresponds to $`m_Qa=m_ca=0.80`$ and $`m_qa=m_sa=0.08`$. Its mass and decay constant are extracted directly from the time-correlation function, which are plotted as the eleventh data point (counting from the smallest one) in Figs. 5 and 6 respectively. The results are $`m_{D_s}a=0.878(2)`$ and $`f_{D_s}a=0.119(2)`$. The mass gives $`m_{D_s}=1964(5)`$ MeV, in good agreement with the mass of $`D_s(1968)`$. The decay constant gives $`f_{D_s}=266(10)`$ MeV, which agrees with the value $`f_{D_s^+}=267\pm 33`$ MeV complied by PDG .
## 6 Summary and Concluding Remarks
In this paper, we have determined the masses and decay constants of pseudoscalar mesons $`K`$, $`D`$ and $`D_s`$, in quenched lattice QCD with exact chiral symmetry. Our results are:
$`m_K`$ $`=`$ $`478\pm 16\pm 20\text{ MeV},`$
$`m_D`$ $`=`$ $`1842\pm 15\pm 21\text{ MeV},`$
$`m_{D_s}`$ $`=`$ $`1964\pm 5\pm 10\text{ MeV},`$
$`f_K`$ $`=`$ $`152\pm 6\pm 10\text{ MeV},`$
$`f_D`$ $`=`$ $`235\pm 8\pm 14\text{ MeV},`$
$`f_{D_s}`$ $`=`$ $`266\pm 10\pm 18\text{ MeV},`$
where in each case, the first error is statistical, while the second is our crude estimate of combined systematic uncertainty. It is interesting to see whether the values of $`f_D`$ and $`f_{D_s}`$ coming soon from the high-statistics charm program of CLEO-c would agree with our values determined by lattice QCD with exact chiral symmetry. Further, we note that in a recent 3-flavor unquenched lattice QCD study with $`O(a^2)`$ improved staggered light quarks and $`O(a)`$-improved charm quark, their results of $`f_D`$ and $`f_{D_s}`$ agree with our values.
Obviously, our next task is to determinate $`f_B`$ and $`f_{B_s}`$ which are of fundamental importance, in view of their experimental determinations are still lacking. Since we will not use any approximations for the heavy b quark, our lattice spacing must be small enough such that $`m_ba<1`$. Even though this does not seem to be formidable for $`f_{B_s}`$, it is unclear whether we can determine $`f_B`$ reliably via chiral extrapolation. Presumably, $`f_B`$ would behave like a function linear in $`m_q`$ for a wide range of $`m_q`$, similar to $`f_D`$ (Fig 6) and $`f_K`$ (Fig. 4), then one should be able to obtain a reliable chiral extrapolation even for data points with $`m_q>m_s`$.
Acknowledgement
This work was supported in part by the National Science Council, Republic of China, under the Grant No. NSC93-2112-M002-016, and by National Center for High Performance Computation at Hsinchu. T.W.C. would like to thank Andreas Kronfeld for a timely remark at the International Conference on QCD and Hadron Physics, Beijing, June 16-20, 2005. |
warning/0506/astro-ph0506386.html | ar5iv | text | # Nuclear fusion in dense matter: Reaction rate and carbon burning
## I Introduction
We will study nuclear fusion rates of identical nuclei in dense stellar matter. This problem is of utmost importance for understanding the structure and evolution of stars of various types. Despite the efforts of many authors the theoretical reaction rates are still rather uncertain, especially at high densities. The uncertainties have two aspects. The first one is related to nuclear physics and is concerned with the proper treatment of nuclear interaction transitions (conveniently described in terms of the astrophysical factor $`S`$). The other issue is associated with aspects of plasma physics and concerns the proper description of Coulomb barrier penetration in a high density many-body system. We will analyze both aspects and illustrate the results taking the carbon fusion reaction as an example.
Considerable experimental effort has been spent on the study of low energy fusion reactions such as <sup>12</sup>C+<sup>12</sup>C, to investigate the impact on the nucleosynthesis, energy production and time scale of late stellar evolution. Nevertheless, it has been difficult to develop a global and reliable reaction formalism to extrapolate the energy dependence of the fusion cross section into the stellar energy range. The overall energy dependence of the cross section is determined by the Coulomb-barrier tunnel probability. One goal of the present work is to apply the São Paulo potential model to provide a general description of the stellar fusion processes. This model does not contain any free parameter and represents a powerful tool to predict average low energy cross sections for a wide range of fusion reactions, as long as the density distribution of the nuclei involved in the reaction can be determined. In this context we also seek to introduce a phenomenological formalism for a generalized reaction rate to describe all the regimes of nuclear burning in a one-component plasma ion system. In this paper we want to demonstrate the applicability of the method on the specific example of <sup>12</sup>C+<sup>12</sup>C, to evaluate the reliability and uncertainty range of the proposed formalism through the comparison with the available low energy data. In a subsequent publication we want to extend the model to multi-ion systems with the aim of simulating a broad range of heavy-ion nucleosynthesis scenarios from thermonuclear burning in hot stellar plasma, to pycnonuclear burning in high density crystalline stellar matter.
Carbon burning represents the third phase of stellar evolution for massive stars ($`M>\mathrm{\hspace{0.33em}8}M_{}`$); it follows helium burning that converts He-fuel to <sup>12</sup>C via the triple $`\alpha `$ process. Carbon burning represents the first stage during stellar evolution determined by heavy-ion fusion processes (e.g., Ref. Wallerstein ). The most important reaction during the carbon burning phase is the <sup>12</sup>C+<sup>12</sup>C fusion Barnes ; additional processes can be <sup>12</sup>C+<sup>16</sup>O and <sup>16</sup>O+<sup>16</sup>O, depending on the <sup>12</sup>C/<sup>16</sup>O abundance ratio which is determined by the <sup>12</sup>C($`\alpha ,\gamma `$)<sup>16</sup>O reaction rate WW93 ; WHR03 . The most important reaction branches are <sup>12</sup>C(<sup>12</sup>C,$`\alpha `$)<sup>20</sup>Ne ($`Q`$=4.617 MeV) and <sup>12</sup>C(<sup>12</sup>C,$`p`$)<sup>23</sup>Na ($`Q`$=2.241 MeV). Carbon burning in evolved massive stars takes place at typical densities $`\rho 10^5`$ g cm<sup>-3</sup> and temperatures $`T(68)\times 10^8`$ K.
Carbon burning is also crucial for type Ia supernovae. These supernova explosions are driven by carbon ignition in cores of accreting massive CO white dwarfs SNI-1 . The burning process proceeds from the carbon ignition region near the center of a white dwarf by detonation or deflagration through the entire white dwarf body. The ignition conditions and time scale are defined by the <sup>12</sup>C+<sup>12</sup>C reaction rate, typically, at $`T(1.57)\times 10^8`$ K and $`\rho (25)\times 10^9`$ g cm<sup>-3</sup> SNI ; bhw04 . Depending on the <sup>16</sup>O abundance, other fusion reactions may also contribute. For these high densities the reaction cross sections are affected by strong plasma screening, which reduces the repulsive Coulomb barrier between interacting <sup>12</sup>C or <sup>16</sup>O nuclei (e.g., Refs. Salp ; CLL02 ; also see Section III).
Explosive carbon burning in the crust of accreting neutron stars has recently been proposed as a possible trigger and energy source for superbursts CB01 ; SB02 ; cn04 . In this scenario small amounts of carbon (3%–10%), which have survived in the preceding rp-process phase during the thermonuclear runaway, ignite after the rp-process ashes are compressed by accretion to a density $`\rho 1.3\times 10^9`$ g cm<sup>-3</sup>. The ignition of a carbon flash requires an initial temperature $`T>\mathrm{\hspace{0.33em}10}^9`$ K, triggering a photodisintegration runaway of the rp-process ashes after a critical temperature $`T2\times 10^9`$ K is reached SBC03 . For these scenarios, carbon burning proceeds in the thermonuclear regime with strong plasma screening (see Section III for details).
At high densities and/or low temperatures the thermonuclear reaction rate formalism is insufficient since the fusion process is mainly driven by the high density conditions in stellar matter (Sections III and IV). This is particularly important for nuclear fusion in the deeper layers of the crust of an accreting neutron star hz . At sufficiently high $`\rho `$ and low $`T`$ nuclei form a crystalline lattice. Neighboring nuclei may penetrate the Coulomb barrier and fuse owing to zero-point vibrations in their lattice sites. In this pycnonuclear burning regime the reaction rate depends mainly on the density and is nearly independent of temperature (e.g. svh69 ; Schramm ). Pycnonuclear burning regimes may not be limited to carbon induced fusion reactions only, but may be driven by a broad range of fusion reactions between stable and neutron rich isotopes hz .
In the following Section II, we discuss the theory of fusion cross sections and calculate the astrophysical $`S`$-factor for carbon burning in the framework of a generalized parameter-free potential model. In Section III we study the Coulomb barrier problem for identical nuclei, and propose an expression for the reaction rate which describes all the regimes of nuclear burning in a dense one-component plasma of atomic nuclei. In Section IV we analyze, for illustration, the main features of <sup>12</sup>C burning from high temperature gaseous or liquid plasma to high density crystalline matter. We summarize and conclude in Section V.
## II Fusion Cross Section and Astrophysical $`S`$-factor
Nuclear reactions are possible after colliding nuclei tunnel through the Coulomb barrier. Recently, a parameter-free model for the real part of the nuclear interaction (São Paulo potential) based on nonlocal quantum effects was developed ref6 ; ref7 ; ref8 ; Toward . In previous work Fusion , this model was applied to the study of fusion processes using the barrier penetration (BP) formalism for about 2500 cross section data, corresponding to approximately 165 different systems. Within the nonlocal model, the bare interaction $`V_N(r,E)`$ is connected with the folding potential $`V_F(r)`$,
$$V_N(r,E)=V_F(r)e^{4\mathrm{v}^2/c^2},$$
(1)
where $`c`$ is the speed of light, $`E`$ is the particle collision energy (in the center-of-mass reference frame), v is the local relative velocity of the two nuclei 1 and 2,
$$\text{v}^2(r,E)=\frac{2}{\mu }\left[EV_C(r)V_N(r,E)\right],$$
(2)
$`V_C(r)`$ is the Coulomb potential, $`\mu =A_1A_2m_\mathrm{u}/(A_1+A_2)`$ is the reduced mass, and $`m_\mathrm{u}`$ is the atomic mass unit. The folding potential depends on the matter densities of the nuclei involved in the collision:
$$V_F(R)=\rho _1(𝒓_1)\rho _2(𝒓_2)V_0\delta (𝑹𝒓_1+𝒓_2)𝑑𝒓_1,$$
(3)
with $`V_0=456`$ MeV fm<sup>3</sup>. The use of the matter densities and delta function in Eq. (3) corresponds to the zero-range approach for the folding potential, which is equivalent Toward to the more usual procedure of using the M3Y effective nucleon-nucleon interaction with the nucleon densities of the nuclei. The advantage in adopting the São Paulo potential to describe the fusion cross section relies on the fact that no additional parameter is necessary once the density distribution of the participating nuclei has been determined. The model is therefore a good choice for a generalized treatment of low energy heavy-ion fusion reactions.
There are several ways to determine the nuclear density distribution Toward ; Afanasjev ; Afanasjev1 ; Gonzalez ; NL2 ; NL3 ; DD-ME1 ; DD-ME2 . Density Functional Theories (DFT) provide for example a successful description of many nuclear ground state properties, in particular, of charge distributions in the experimentally known region. Since these theories are universal in the sense that their parameter sets are carefully adjusted and valid all over the periodic table, one can expect that they also yield reliable predictions for nuclei far from stability. Non-relativistic density functionals, such as the Skyrme or Gogny functional, have been widely used in the literature. In recent years, relativistic density functionals have played an increasingly important role since they provide a fully consistent description of the spin-orbit splitting. This is of greatest importance for nuclei far from stability. The spin-orbit splitting determines the shell structure, the most basic ingredient in any microscopic theory of finite nuclei. In fact, the results obtained with relativistic functionals are in very good agreement with experimental data, throughout the periodic table, despite having a smaller number of adjustable parameters in comparison with the non-relativistic case. Best known is the relativistic Hartree-Bogoliubov theory Afanasjev ; Afanasjev1 ; Gonzalez , which includes pairing correlations with finite range pairing forces. It provides a unified description of mean-field and pairing correlations in nuclei.
These functionals contain a strong density dependence, either through non-linear coupling terms between the meson fields (e.g., in the Lagrangians with the parameter sets NL2 NL2 and NL3 NL3 ), or by using an explicit density dependence for the meson-nucleon vertices (e.g., in the parameter sets DD-ME1 DD-ME1 and DD-ME2 DD-ME2 ).
In the present paper, we consider only spherical nuclear shapes. Pairing correlations are in principle included, but they vanish for the <sup>12</sup>C nucleus. In Fig. 1 we compare the calculated densities with experimental data c12-dens . The RHB calculations are in good agreement with surface properties best described by the NL2, DD-ME1, and DD-ME2 effective interactions.
To apply the BP model for calculating fusion cross sections one needs the effective potential defined as a sum of the Coulomb, nuclear and centrifugal components:
$$V_{\mathrm{eff}}(r,E)=V_C(r)+V_N(r,E)+\frac{\mathrm{}(\mathrm{}+1)\mathrm{}^2}{2\mu r^2}.$$
(4)
Following the BP model one can associate the fusion cross section with the particle flux transmitted through the barrier,
$$\sigma _{ij}(E)=\frac{\pi }{k^2}\underset{\mathrm{}=0}{\overset{\mathrm{}_{cr}}{}}(2\mathrm{}+1)T_{\mathrm{}}.$$
(5)
It is important to point out that the sum in Eq. (5) is performed up to a maximum $`\mathrm{}`$ wave ($`\mathrm{}_{cr})`$, which corresponds to the greatest value of angular momentum that produces a pocket (and a barrier) in the corresponding effective potential, Eq. (4). For $`\mathrm{}`$-waves with effective barrier heights $`V_B\mathrm{}<E`$, the shape of the effective potential can be approximated by a parabola with curvature defined as
$$\mathrm{}\omega _{\mathrm{}}=\left|\frac{\mathrm{}^2}{\mu }\frac{d^2V_{\mathrm{eff}}}{dr^2}\right|_{R_B\mathrm{}}^{1/2},$$
(6)
where $`R_B\mathrm{}`$ is the barrier radius. In such cases, the transmission coefficients have been obtained through the Hill-Wheeler formula Hill :
$$T_{\mathrm{}}=\left\{1+\mathrm{exp}\left[\frac{2\pi (V_B\mathrm{}E)}{\mathrm{}\omega _{\mathrm{}}}\right]\right\}^1.$$
(7)
On the other hand, for $`\mathrm{}`$-waves with $`V_B\mathrm{}>E`$, instead of the Hill-Wheeler formula, we employ a more appropriate heuristic treatment based on a WKB approximation Schiff :
$$T_{\mathrm{}}=[1+\mathrm{exp}(S_{\mathrm{}})]^1,$$
(8)
$$S_{\mathrm{}}=_{r_1}^{r_2}\sqrt{\frac{8\mu }{\mathrm{}^2}[V_{\mathrm{eff}}(r,E)E]}𝑑r,$$
(9)
where $`r_1`$ and $`r_2`$ are the classical turning points. At low energies, the WKB method gives values for the transmission coefficients which are quite different from those of the Hill-Wheeler formula. In this case, we define the barrier curvature by connecting Eqs. (6) and (7):
$$\mathrm{}\omega _{\mathrm{}}=\frac{2\pi (V_B\mathrm{}E)}{S_{\mathrm{}}}.$$
(10)
The overall results provided by the BP model are in very good agreement with the fusion data for energies above the s-wave barrier height. For light systems $`(\mu 8m_\mathrm{u})`$ the model also shows very good agreement with fusion data at sub-barrier energies Fusion . Therefore, the use of the BP model in calculating the fusion cross section at energies of astrophysical interest for the <sup>12</sup>C+<sup>12</sup>C system is entirely justified.
Historically, reaction cross sections $`\sigma (E)`$ at very low energies, typical for astrophysical conditions, have been expressed in terms of the astrophysical $`S`$-factor (e.g., Ref. FCZII ),
$$S(E)=\sigma (E)Ee^{2\pi \eta },$$
(11)
where $`\eta =(Z_1Z_2e^2/\mathrm{})\sqrt{\mu /(2E)}`$ is the usual Gamow parameter.
Considerable efforts have been made over the last decades to measure the <sup>12</sup>C+<sup>12</sup>C fusion cross section at very low energies Patterson ; Mazarakis ; High ; Eli ; Kettner ; Becker . The experimentally determined $`S`$-factors are shown in Fig. 2. For reaction rate calculations the experimental $`S`$-factor needs to be extrapolated towards the stellar energy range, the Gamow window, which depends sensitively on the temperature and density conditions of the stellar environment. The typical range of energy $`E`$ for thermonuclear carbon burning, in the center-of-mass reference, varies from 1 to 4 MeV. For pycnonuclear carbon burning in the neutron star crusts the energies can be as low as 10 keV. Large discrepancies between the different experimental results at low energies complicate a reliable extrapolation of $`S(E)`$ towards such low $`E`$. In addition, the $`S`$-factor shows pronounced resonant structures, presumably resulting from quasimolecular doorway states. Theoretical calculations of $`S(E)`$ using the effective interactions NL2, NL3, DD-ME1, DD-ME2 agree reasonably well, within a factor $``$ 3.5 in the limit $`E`$ $``$ 0 (Fig. 2). Furthermore, the resonant behavior of the data cannot be described with the BP calculations because the effects of nuclear structure were neglected. However, an average description of the data (neglecting resonant oscillations) for the sub-barrier region ($`E`$$`<`$$``$ 6.0 MeV) is reproduced satisfactorily. Such a description of an average $`S`$-factor is quite sufficient since the reaction rate formalism relies on the average $`S`$-factor behavior over the entire Gamow range.
In this context it is important to emphasize that the main purpose of this paper is not to investigate the oscillations in the <sup>12</sup>C+<sup>12</sup>C fusion excitation function. In order to reproduce the resonances we could, for example, use the concept of internal and barrier waves based on a semiclassical description Oh87 or adopt the R-matrix formalism (e.g., Michaud and Vogt Mi72 ). However, neither theoretical approach would allow us to extrapolate with confidence the fusion cross section to the energy region of astrophysical interest.
In order to calculate the carbon burning rate, we use the values of $`S(E)`$ obtained on the basis of the well established NL2 effective interaction. As one can see from Fig. 1, the <sup>12</sup>C density distribution obtained using the parameter set NL2 can describe satisfactorily the surface properties, which is the most important region for the fusion process at low energies. The values of $`S(E)`$ calculated at $`E19.8`$ MeV can be fitted by an analytic expression
$$S(E)=5.15\times 10^{16}\mathrm{exp}\left\{0.428E\frac{3E^{0.308}}{1+\mathrm{e}^{0.613(8E)}}\right\}\mathrm{MeV}\mathrm{barn},$$
(12)
where the center-of-mass energy $`E`$ is expressed in MeV. The formal maximum fit error, 16%, occurs at $`E=5.8`$ MeV. However, let us bear in mind that the values of $`S(E)`$ provided by the NL2 model and given by Eq. (12) are actually uncertain within a factor of $`3.5`$.
With the aim of investigating the validity of our assumption for the real part of the nuclear interaction, we performed an optical model (OM) analysis of the <sup>12</sup>C+<sup>12</sup>C elastic scattering data at energies around and slightly above the Coulomb barrier Treu80 . We defined the imaginary part of the optical potential, which accounts for the nuclear absorption process, as
$$W(r,E)=N_iV_N(r,E),$$
(13)
where $`V_N(r,E)`$ is described by Eq. (1) and $`N_i`$ = 0.78 was determined by adjusting thirty elastic scattering angular distributions corresponding to seven different heavy-ion systems and measured in a very wide energy range Al03 . Figure 3 illustrates a comparison between our OM analysis and five elastic scattering angular distribution data of the <sup>12</sup>C+<sup>12</sup>C system. As one can note, it is possible to obtain a reasonable description of the data by adopting the São Paulo potential to account for the real part of the nuclear interaction, combined with a simple model to describe the imaginary part of the optical potential. This means that both elastic scattering and fusion processes can be described by the same real part of the nuclear interaction, which has been well accounted by the São Paulo potential, Eq. (1). As discussed in Ref. Al03 , details on the absorption part of the interaction are not very important for describing the elastic scattering data, which allows us to get reasonable estimates for the <sup>12</sup>C+<sup>12</sup>C system.
Further experiments at lower energies are necessary to confirm the validity of the predicted <sup>12</sup>C+<sup>12</sup>C $`S`$-factor and its impact on the reaction rate. However, the $`S`$-factor is not the only uncertainty for a reliable description of the <sup>12</sup>C+<sup>12</sup>C fusion process in stellar matter. The reaction rates are also uncertain because of the problems in calculating the probability of Coulomb barrier penetration in a dense many-body environment. We shall discuss these problems in Section III and show that the associated uncertainties are higher than the current uncertainties in the values of $`S(E)`$.
## III Nuclear fusion rate in dense matter
### III.1 Physical conditions and reaction regimes
In the following we will turn our attention to the plasma physics aspects of nuclear burning in dense matter. We will focus on the formalism of fusion reactions between identical nuclei $`(A,Z)+(A,Z)`$ in the wide domain of temperatures $`T`$ and densities $`\rho `$, characteristic for the range of stellar environments outlined above.
As an example, we consider the <sup>12</sup>C+<sup>12</sup>C reaction in stellar matter at conditions displayed in the $`\rho T`$ phase diagram in Fig. 4. Under these conditions, carbon is fully ionized (either by electron pressure and/or by high temperature) and immersed in an almost uniform electron background. The electrons are typically strongly degenerate; their degeneracy temperature $`T_\mathrm{F}`$ is shown in the figure.
The state of ions (nuclei) is determined by the Coulomb coupling parameter $`\mathrm{\Gamma }=Z^2e^2/(aT)`$, where $`a=[3/(4\pi n_i)]^{1/3}`$ is the ion-sphere radius and $`n_i`$ is the number density of ions; the Boltzmann constant is set $`k_\mathrm{B}1`$. If $`\mathrm{\Gamma }<\mathrm{\hspace{0.33em}1}`$ (which happens at $`T>T_l=Z^2e^2/a`$, see Fig. 4), the ions constitute a Boltzmann gas, while at higher $`\mathrm{\Gamma }`$ they constitute a strongly coupled Coulomb liquid. The gas transforms smoothly into the liquid, without any phase transition. At small $`T`$ (large $`\mathrm{\Gamma }`$) the liquid can solidify. In the density range displayed in Fig. 4, the solidification occurs at $`T=T_\mathrm{m}=Z^2e^2/a\mathrm{\Gamma }_\mathrm{m}`$, where $`\mathrm{\Gamma }_\mathrm{m}=175`$ (e.g., De Witt et al. dewittetal03 ). The important measure of quantum effects in ion motion is provided by the ion plasma frequency $`\omega _p=\sqrt{4\pi Z^2e^2n_i/m}`$ or the associated ion plasma temperature $`T_p=\mathrm{}\omega _p`$ ($`m`$ being the ion mass). As a rule, the quantum effects are strongly pronounced at $`T`$ below $`T_p`$.
Figure 4 shows that the ion system can have very different properties, depending on $`T`$ and $`\rho `$. As a result, there are five qualitatively different regimes of nuclear burning in dense matter (Salpeter and Van Horn svh69 ). These are (1) the classical thermonuclear regime; (2) the thermonuclear regime with strong plasma screening; (3) the thermo-pycnonuclear regime; (4) the thermally enhanced pycnonuclear regime; and (5) the zero-temperature pycnonuclear regime. The regimes differ mainly in the character of the Coulomb barrier penetration of reacting nuclei. The penetration can be greatly complicated by Coulomb fields of ions which surround the reacting nuclei. These fields are fluctuating and random (e.g., Alastuey and Jancovici aj78 ).
A strict solution of the barrier penetration problem should imply the calculation of the tunneling probability in a random potential, with subsequent averaging over an ensemble of random potentials. This program has not been fully realized so far. The exact theory should take into account a range of effects which can be subdivided (somewhat conventionally) into classical and quantum ones. The classical effects are associated with classical motion of plasma ions and with related structure of Coulomb plasma fields (including spatial and temporal variability of these fields). The quantum effects manifest themselves in ion motion (e.g., zero-point ion vibrations), quantum “widths” of ion trajectories during Coulomb barrier penetration, and quantum statistics of reacting nuclei. The effects of quantum statistics are usually small due to the obvious reason that quantum tunneling lengths are typically much larger than nuclear radii. The smallness of these effects has been confirmed by Ogata ogata97 in path-integral Monte Carlo (PIMC) simulations.
The reaction rates in the classical thermonuclear regime are well known (e.g., Fowler, Caughlan, and Zimmerman FCZII ); they have been tested very successfully by the theory and observations of the evolution of normal stars. This theory will be only shortly reviewed in the following section. The reaction rates in other regimes have been calculated by a number of authors in different approximations. In the following we summarize the main results published after the seminal paper by Salpeter and Van Horn svh69 (see that paper for references to earlier works). Let us stress that the reaction rate is a rapidly varying function of plasma parameters. In the most important density-temperature domain it varies over tens orders of magnitude (Section IV). In this situation, a very precise calculation of the reaction rate is very difficult but not required for many applications.
### III.2 Classical thermonuclear reaction rate
The classical thermonuclear regime takes place at sufficiently high $`T`$ and low $`\rho `$ so that the ions constitute a Boltzmann gas ($`TT_l`$, Fig. 4). The tunnel probability (penetrability) through the Coulomb barrier depends on the energy of the interacting ions; the main contribution to the reaction rate comes from ion collisions with energies approximately equal to the Gamow peak energy $`E_{\mathrm{pk}}`$ (that is much higher than $`T`$). This regime is typical for all nuclear burning stages in “normal” stars (from the main sequence to pre-supernovae).
The thermonuclear reaction rate is expressed by
$$R_{\mathrm{th}}=\frac{n_i^2}{2}\mathrm{\hspace{0.17em}4}\sqrt{\frac{2E_{\mathrm{pk}}}{3\mu }}\frac{S(E_{\mathrm{pk}})}{T}\mathrm{exp}(\tau ),$$
(14)
where $`E_{\mathrm{pk}}=T\tau /3`$ is the Gamow peak energy and
$$\tau =\left(\frac{27\pi ^2\mu Z_1^2Z_2^2e^4}{2T\mathrm{}^2}\right)^{1/3}=\left(\frac{27\pi ^2mZ^4e^4}{4T\mathrm{}^2}\right)^{1/3}$$
(15)
is the parameter which characterizes the penetrability $`\mathrm{exp}(\tau `$). The parameter $`\tau `$ can be rewritten as
$$\tau =3(\pi /2)^{2/3}(E_\mathrm{a}/T)^{1/3},E_\mathrm{a}mZ^4e^4/\mathrm{}^2.$$
(16)
Now the reaction rate can be presented as
$$R_{\mathrm{th}}=\frac{n_i^2}{2}S(E_{\mathrm{pk}})\frac{\mathrm{}}{mZ^2e^2}P_{\mathrm{th}}F_{\mathrm{th}},$$
(17)
where $`\mathrm{}/(mZ^2e^2)`$ is a convenient dimensional factor, $`F_{\mathrm{th}}`$ is the exponential function, and $`P_{\mathrm{th}}`$ is the pre-exponent:
$$F_{\mathrm{th}}=\mathrm{exp}(\tau ),P_{\mathrm{th}}=\frac{8\pi ^{1/3}}{\sqrt{3}\mathrm{\hspace{0.17em}2}^{1/3}}\left(\frac{E_\mathrm{a}}{T}\right)^{2/3}.$$
(18)
The classical thermonuclear reaction rate decreases exponentially with decreasing $`T`$.
### III.3 Thermonuclear burning with strong plasma screening
The thermonuclear regime with strong plasma screening operates in a colder and denser plasma ($`T_p<T<T_l`$), where ions constitute a strongly coupled classical Coulomb system (liquid or solid). The majority of ions in such a system are confined in deep Coulomb potential wells ($`Z^2e^2/a>T`$). The main contribution into the reaction rate comes from a small amount of higher-energy, unbound ions with $`EE_{\mathrm{pk}}T`$ from the tail of the Boltzmann distribution. The plasma screening effects are produced by surrounding plasma ions and simplify close approaches of the reacting nuclei, required for a successful Coulomb tunneling. This enhances the reaction rate with respect to that given by Eqs. (17) and (18).
The enhancement has been studied by a number of authors, beginning with Salpeter Salp ; it can reach many orders of magnitude. Calculations show that the equations given in Section III.2 remain valid in this regime, but the penetrability function $`F_{\mathrm{th}}`$ has to be corrected for the screening effects:
$$F_{\mathrm{th}}=F_{\mathrm{sc}}\mathrm{exp}(\tau )F_{\mathrm{sc}}=\mathrm{exp}(h),$$
(19)
where $`F_{\mathrm{sc}}`$ is the enhancement factor, and $`h`$ is a function of plasma parameters.
Plasma screening effects are usually modeled by introducing a mean-force plasma potential $`H(r)`$. In this approximation, the reacting nuclei move in a potential $`W(r)=Z^2e^2/rH(r)`$. The mean-force plasma potential $`H(r)`$ is static and spherically symmetric. It cannot take into account dynamical variations of plasma microfields and their instantaneous spatial structures in the course of an individual tunneling event. In the mean-force approximation, the function $`h`$ consists of two parts, $`h=h_0+h_1`$, where the leading term $`h_0=H(0)/T`$ ($`|h_1|`$) is calculated assuming a constant plasma potential $`H(r)=H(0)`$ during the quantum tunneling, while $`h_1`$ is a correction due to a weak variation of $`H(r)`$ along the tunneling path. Note that according to simple estimates (e.g., Ref. ys89 ) typical tunneling lengths of reacting ions in the thermonuclear regime (where $`T>T_p`$) are considerably smaller than the ion sphere radius $`a`$, and typical tunneling times are much smaller than the plasma oscillation period $`\omega _p^1`$. This justifies the assumption of almost constant and static plasma potential during a tunneling event.
The mean-force plasma potential $`H(r)`$ for a classical strongly coupled system of ions (liquid or solid) can be determined using classical Monte Carlo (MC) sampling (e.g., DeWitt et al. dgc73 ). MC sampling gives the static radial-pair distribution function of ions $`g(r)=\mathrm{exp}(W(r)/T)`$ which enables one to find $`H(r)`$. In this way one can accurately determine $`g(r)`$ and $`H(r)`$ at not too small $`r`$ (typically, at $`r>a`$), because of poor MC statistics of close ion separations. The potential $`H(r)`$ at small $`r`$, required for a tunneling problem, is obtained by extrapolating MC values of $`H(r)`$ to $`r0`$; the extrapolation procedure is a delicate subject and may be ambiguous (as discussed, e.g., by Rosenfeld rosenfeld96 ).
It is only $`H(0)`$ which is required for finding $`h_0`$. For a classical ion system, $`H(0)`$ can be determined by $`H(0)=\mathrm{\Delta }`$, where $`\mathrm{\Delta }`$ is a difference of Coulomb free energies (for a given system and for a system with two nuclei merging into one compound nucleus; e.g., DeWitt et al. dgc73 ). In this approximation, the enhancement factor of the nuclear reaction becomes a thermodynamic quantity and acquires a Boltzmann form, $`\mathrm{exp}(h_0)=\mathrm{exp}(\mathrm{\Delta }/T)`$, showing that plasma screening increases the probability of close separations (and subsequent quantum tunneling); $`h_0`$ becomes the function of one argument $`\mathrm{\Gamma }`$. Assuming a linear mixing rule in a multi-component strongly coupled ion system, Jancovici jancovici77 obtained $`h_0=2f_0(\mathrm{\Gamma })f_0(2^{5/3}\mathrm{\Gamma })`$, where $`f_0(\mathrm{\Gamma })`$ is a Coulomb free energy of one ion in a one-component plasma of ions (in units of $`T`$). In a Coulomb liquid at $`\mathrm{\Gamma }>\mathrm{\hspace{0.33em}1}`$ the linear mixing rule is highly accurate (DeWitt and Slattery ds03 ); the function $`f_0(\mathrm{\Gamma })`$ is now determined from MC sampling with very high accuracy (e.g., Ref. pc00 ; ds03 ). In this way the function $`h_0(\mathrm{\Gamma })`$ has been calculated in many papers (e.g., Refs. jancovici77 ; ys89 ; rosenfeld96 ; ds99 ), and the results are in very good agreement. Let us present the analytical approximation of $`h_0(\mathrm{\Gamma })`$ which follows from the recent MC results of DeWitt and Slattery ds99 for a Coulomb liquid at $`1\mathrm{\Gamma }170`$:
$$h_0=1.0563\mathrm{\Gamma }+1.0208\mathrm{\Gamma }^{0.3231}0.2748\mathrm{ln}\mathrm{\Gamma }1.0843.$$
(20)
However, this accurate expression is inconvenient for further use, and we propose another fit
$$h_0=C_{\mathrm{sc}}\mathrm{\Gamma }^{3/2}/[(C_{\mathrm{sc}}/\sqrt{3})^4+\mathrm{\Gamma }^2]^{1/4},$$
(21)
where $`C_{\mathrm{sc}}=1.0754`$. It approximates $`\mathrm{e}^{h_0}`$ with the maximum error of $``$40% at $`\mathrm{\Gamma }=170`$, quite sufficient for our purpose. There may be still some uncertainty of the reaction rate associated with the choice of $`C_{\mathrm{sc}}`$ but it seems to be not higher than the uncertainty in the $`S`$-factor (Section II). Our fit function in Eq. (21) is chosen in such a way to reproduce also the well known expression $`h_0\sqrt{3}\mathrm{\Gamma }^{3/2}`$ derived by Salpeter Salp for the classical thermonuclear regime ($`\mathrm{\Gamma }1`$), where $`h_01`$ and the plasma screening is weak. In the Coulomb liquid, at $`\mathrm{\Gamma }>\mathrm{\hspace{0.33em}1}`$, we have actually the linear function $`h_0=C_{\mathrm{sc}}\mathrm{\Gamma }`$. Such a function was obtained by Salpeter Salp using a simple model of ion spheres (with slightly lower coefficient, $`C_{\mathrm{sc}}^{\mathrm{Salp}}=1.057`$).
Some authors calculated $`h_0`$ and the associated enhancement factor $`\mathrm{e}^{h_0}`$ by extrapolating MC $`H(r)`$ to $`r0`$ (as discussed above). In particular, Ogata et al. oii91 ; oiv93 employed this formalism to study the enhancement of nuclear reactions in one-component and two-component strongly coupled ion liquids. The enhancement factor $`\mathrm{e}^{h_0}`$ for a one-component ion liquid, calculated in these papers (e.g., Eq. (6) in Ref. oii91 ), is systematically higher than the factor given by Eq. (20) or (21). The difference reaches a factor of approximately 40 for $`\mathrm{\Gamma }170`$. Because the enhancement factor itself becomes as high as $`\mathrm{e}^{h_0}10^{74}`$ at $`\mathrm{\Gamma }170`$, such a difference is insignificant for many applications. As shown by Rosenfeld rosenfeld96 , the difference comes from the problems of extrapolation of $`H(r)`$ to $`r0`$ in Refs. oii91 ; oiv93 . The function $`h_0`$ was also calculated by Ogata ogata97 using direct PIMC method. His result (his Eq. (19)) is in much better agreement with Eq. (20). The maximum difference of $`\mathrm{e}^{h_0}`$ reaches only a factor of approximately 6 at $`\mathrm{\Gamma }=170`$. Recently new PIMC calculations have been performed by Pollock and Militzer pm04 but the authors have not calculated directly $`h_0(\mathrm{\Gamma })`$.
Let us emphasize that the enhancement factor $`\mathrm{e}^{h_0}`$, derived in a constant mean-force plasma potential $`H(0)`$, is invariant with respect to the order of the mean-force averaging and the tunneling probability calculation. One can consider a real (random) plasma potential, constant over a tunneling path in an individual tunneling event. Calculating the tunneling probability and averaging over an ensemble of realizations of plasma potentials, one comes (e.g., Ref. ys89 ) to the same expression for $`h_0`$ as given by the mean-force potential.
In addition to $`\mathrm{e}^{h_0}`$, the enhancement factor $`F_{\mathrm{sc}}`$ in Eq. (19) contains a smaller factor $`\mathrm{e}^{h_1}`$, associated with variations of the plasma potential along the tunneling path. Numerous calculations of $`h_1`$ have commonly employed the mean-force potential $`H(r)`$. The results are sensitive to the behavior of $`H(r)`$ at small $`r`$ (where this behavior is not very certain). For example, Jancovici jancovici77 got $`h_1=(5/32)\mathrm{\Gamma }(3\mathrm{\Gamma }/\tau )^2`$. Note that for the thermonuclear burning ($`T>T_p`$, Section III.1), the ratio $`3\mathrm{\Gamma }/\tau r_\mathrm{t}/a(T/T_p)^{2/3}`$ can be regarded as a small parameter ($`r_\mathrm{t}`$ being the tunneling length). It is possible that the mean-force approximation is too crude for calculating $`h_1`$. For that reason, we will not specify $`h_1`$ in this section. Our final expression for the reaction rate will include $`h_1`$, but phenomenologically, when we combine reaction rates in all regimes (Section III.7).
### III.4 Zero-temperature pycnonuclear fusion
The zero-temperature pycnonuclear regime operates in a cold and dense matter ($`T`$ well below $`T_p`$) in a strongly coupled quantum system of nuclei. In this regime the Coulomb barrier is penetrated owing to zero-point vibrations of neighboring nuclei which occupy their ground states in a strongly coupled system. One usually considers pycnonuclear reactions in a crystalline lattice of nuclei but they are also possible in a quantum liquid. The main contribution into the reaction rate comes from pairs of nuclei which are most closely spaced. The reaction rate is temperature-independent but increases exponentially with increasing density as we will discuss in the following.
Pycnonuclear reaction rates between identical nuclei in crystalline lattice have been calculated by many authors using different approximations. In analogy with Eq. (17), the resulting reaction rates can be written as
$$R_{\mathrm{pyc}}=\frac{n_i^2}{2}S(E_{\mathrm{pk}})\frac{\mathrm{}}{mZ^2e^2}P_{\mathrm{pyc}}F_{\mathrm{pyc}},$$
(22)
where $`F_{\mathrm{pyc}}`$ and $`P_{\mathrm{pyc}}`$ depend on the density and have the form
$$F_{\mathrm{pyc}}=\mathrm{exp}\left(C_{\mathrm{exp}}/\sqrt{\lambda }\right),P_{\mathrm{pyc}}=8C_{\mathrm{pyc}}\mathrm{\hspace{0.17em}11.515}/\lambda ^{C_{\mathrm{pl}}}.$$
(23)
The dimensionless parameters $`C_{\mathrm{exp}}`$, $`C_{\mathrm{pl}}`$ and $`C_{\mathrm{pyc}}`$ are model dependent (see below). The dimensionless parameter $`\lambda `$ is expressed in terms of the mass fraction $`X_i`$ contained in atomic nuclei (in a one-component ion plasma under study) and the mass density $`\rho `$ of the medium
$$\lambda =\frac{\mathrm{}^2}{mZ^2e^2}\left(\frac{n_i}{2}\right)^{1/3}=\frac{1}{AZ^2}\left(\frac{1}{A}\frac{\rho X_i}{1.3574\times 10^{11}\mathrm{g}\mathrm{cm}^3}\right)^{1/3}.$$
(24)
For densities $`\rho `$ lower than the neutron drip density ($`4\times 10^{11}`$ g cm<sup>-3</sup>; e.g., Ref. st83 ), one can set $`X_i=1`$, while for higher $`\rho `$ one has $`X_i<1`$ because of the presence of free (dripped) neutrons.
The reaction rate can be expressed numerically as
$$R_{\mathrm{pyc}}=\rho X_iAZ^4S(E_{\mathrm{pk}})C_{\mathrm{pyc}}\mathrm{\hspace{0.17em}10}^{46}\lambda ^{3C_{\mathrm{pl}}}\mathrm{exp}\left(C_{\mathrm{exp}}/\sqrt{\lambda }\right)\mathrm{s}^1\mathrm{cm}^3,$$
(25)
where $`\rho `$ is in g cm<sup>-3</sup> and $`S(E_{\mathrm{pk}})`$ is in MeV barn. The typical energy of the interacting nuclei is $`E_{\mathrm{pk}}\mathrm{}\omega _p`$.
Table 1 lists the values of $`C_{\mathrm{exp}}`$, $`C_{\mathrm{pl}}`$ and $`C_{\mathrm{pyc}}`$ reported in the literature for two models (1 and 2) of Coulomb barrier penetration by Salpeter and Van Horn svh69 , for six models (3–8) by Schramm and Koonin Schramm , and for one model (9) by Ogata, Iyetomi and Ichimaru oii91 . The corresponding carbon burning rates are plotted as a function of density in Fig. 5. In this figure (as well as in Figs. 4 and 6) we use the astrophysical factors given by the fit expression (12). Actually, the $`S`$-factors are uncertain within one order of magnitude (Section II) but we ignore these uncertainties (because they seem to be much lower than those associated with the Coulomb barrier problem).
All the authors cited above have treated quantum tunneling by fixing the center-of-mass of reacting nuclei in its equilibrium position. All models, except for models 5–8, focus on nuclear reactions in the body-centered cubic (bcc) lattice of atomic nuclei. This lattice is thought to be preferable over other lattices, particularly, over the face-centered cubic (fcc) lattice. The main reason is that the bcc lattice is more tightly bound in the approximation of a rigid electron background. However, the difference in binding energies of bcc and fcc lattices is small (see, e.g., Ref. Schramm ), and a finite polarizability of the electron background complicates the problem baiko . Therefore, one cannot exclude that the lattice type is fcc.
Salpeter and Van Horn svh69 calculated the quantum tunneling probability of interacting nuclei in a bcc Coulomb lattice using the three-dimensional WKB approximation (most adequate for the given problem). The authors employed two models, static and relaxed lattice (models 1 and 2 in Table 1), to account for the lattice response to the motion of tunneling nuclei. The static lattice model assumes that surrounding nuclei remain in their original lattice sites during the tunneling process. The relaxed lattice model assumes that the surrounding nuclei are promptly rearranged into new equilibrium positions in response to the motion of the reacting nuclei. Simple estimates show that the actual tunneling is dynamical (neither static not relaxed). Thus, the static-lattice and relaxed-lattice models impose constraints on the actual reaction rate. In Ref. svh69 the screening potential for the relaxed-lattice model was calculated approximately; the energy difference between the initial and fused states was evaluated by subtracting the energies of the corresponding Wigner-Seitz (WS) spheres. The relaxed lattice simplifies Coulomb tunneling and increases the reaction rate with respect to the static lattice (cf. curves 1 and 2 in Fig. 5).
Schramm and Koonin Schramm applied this treatment to the bcc and fcc, static and relaxed lattices in the same WKB approximation. They calculated the screening potential for the relaxed-lattice model with improved accuracy (model 3 of Table 1 for the bcc lattice and model 7 for fcc). For comparison, they also used the screening potential for the relaxed-lattice obtained in the WS approximation (as in Ref. svh69 ). Unfortunately, they calculated the tunneling probability neglecting the correction $`\mathrm{e}^K`$ for to the “curvature of trajectories” of reacting ions. This is the main reason for the formal disagreement between the results of Salpeter and Van Horn svh69 and Schramm and Koonin Schramm for the static-lattice and relaxed-lattice-WS models (1 and 2) of bcc crystals. The inclusion of the curvature correction should reduce the constant $`C_{\mathrm{pyc}}`$ in Eq. (25) and the reaction rates calculated in Ref. Schramm . Fortunately, this correction can be extracted by comparing Eq. (38) of Ref. svh69 with Eq. (31) of Ref. Schramm . In this way we get $`\mathrm{e}^K=0.067`$ for the static bcc lattice, and $`\mathrm{e}^K=0.050`$ for the relaxed-WS bcc lattice. After introducing this correction into the coefficients $`C_{\mathrm{pyc}}`$, obtained formally from the results of Schramm and Koonin, these coefficients become identical to those given by Salpeter and Van Horn. Thus, Schramm and Koonin actually exactly reproduce models 1 and 2 of Salpeter and Van Horn. The curvature correction for the models 3–8 of Schramm and Koonin have not been determined. We expect it to be $`\mathrm{e}^K0.050`$ for model 3 and $`\mathrm{e}^K=0.067`$ for model 4 (bcc crystals), and we introduced such corrections in Table 1. We introduced, somewhat arbitrarily, the correction $`\mathrm{e}^K=0.05`$ in all fcc lattice models 5–8.
The two versions of the screening potential for the relaxed lattice (WS and more accurate) almost coincide. Accordingly, models 2 and 3 yield almost the same reaction rates for the bcc lattice, while models 6 and 7 yield nearly identical rates for the fcc lattice (Fig. 5).
Schramm and Koonin Schramm also took into account the dynamical effect of motion of the surrounding ions in response to the motion of tunneling nuclei in the relaxed lattice (models 4 and 8). This effect was described by introducing the effective mass of the reacting nuclei. The effective mass appears to be noticeably higher than the real nucleus mass, reducing the tunneling probability. It turns out that the reduction almost exactly compensates the increase of the tunneling probability due to the lattice relaxation neglecting the effective mass effects svh69 . Accordingly, model 4 of Schramm and Koonin Schramm gives almost the same reaction rate as model 1 (for bcc); and model 8 gives almost the same rate as model 5 (for fcc). This means that the two limiting approximations, the static-lattice and relaxed-lattice, yield very similar reaction rates. It is natural to expect that the actual reaction rate (to be calculated for the dynamically responding lattice) would be the same, and the problem of dynamical tunneling is thus solved Schramm . Notice that this conclusion is made using the curvature corrections $`\mathrm{e}^K`$ adopted above (whereas the accurate curvature correction for the effective mass model has not been calculated).
Zero-temperature pycnonuclear reactions in bcc crystals were also studied by Ogata, Iyetomi, and Ichimaru oii91 and Ichimaru, Ogata, and Van Horn iovh92 using MC lattice screening potentials. These authors considered one-component and two-component ion systems. Model 9 of Table 1 represents their results for one-component bcc crystals. In order to calculate the tunneling probability, the authors used the mean-force plasma screening potential $`H(r)`$, obtained from MC sampling in a classical bcc crystal at $`r>a`$ and extrapolated to $`r0`$ (see Section III.3). This potential is static and spherically symmetric. It cannot take into account the dynamics of lattice response and the anisotropic character of the real screening potential in a lattice. Furthermore, the barrier penetration was calculated by solving numerically the effective radial Schrödinger equation. This procedure is more approximate than the direct WKB approach of Salpeter and Van Horn svh69 and Schramm and Koonin Schramm (particularly, it neglects the curvature corrections). Numerically, Ichimaru, Ogata and Iyetomi oii91 give a reaction rate which is close to the relaxed lattice model (model 2) of Salpeter and Van Horn svh69 . The main reason for the coincidence of these rates is that the screening potential in the radial equation of Ichimaru, Ogata and Iyetomi is close to the relaxed-lattice screening potential of Salpeter and Van Horn at ion separations $`r1.5a`$, most important for pycnonuclear tunneling problem (see Fig. 2 in Ref. oii91 ).
In spite of the differences in theoretical models 1–9, they result in similar reaction rates (Fig. 5). According to the above discussion, models 1 and 4 seem to be the most reliable among all available models for the bcc lattice, while models 5 and 8 seem to be the most reliable for fcc. These reaction rates may be modified, for instance, by taking into account the quantum effect of the spreading of WKB trajectories or by a more careful treatment of the center-of-mass motion of reacting nuclei. Such effects will possibly reduce the reaction rate (as discussed in Ref. pm04 with regard to the spreading of WKB trajectories). This could have been studied by direct PIMC simulations (e.g., Refs. ogata97 ; pm04 ). PIMC is also a good tool to confirm the conclusions on dynamical effects of lattice response. However, PIMC is time consuming and requires very powerful computers. It is not clear whether today’s computer capabilities are sufficient to obtain accurate PIMC pycnonuclear reaction rates.
We suggest to calculate the reaction rates from Eq. (25) taking into account that the constants $`C_{\mathrm{exp}}`$, $`C_{\mathrm{pl}}`$ and $`C_{\mathrm{pyc}}`$ are not known very precisely. In particular, we propose two “limiting” purely phenomenological sets of these constants labeled as models 10 and 11 in Table 1. These limiting parameters define the maximum and minimum reaction rates which enclose all model reaction rates 1–4 and 9 (proposed in the literature for the bcc lattice in a density range where the pycnonuclear carbon burning is important). They also enclose the most reliable models 5 and 8 for the fcc lattice.
The crucial parameter for modeling pycnonuclear fusion is the exponent $`F_{\mathrm{pyc}}=\mathrm{exp}(C_{\mathrm{pyc}}/\sqrt{\lambda })`$ in Eq. (23) that characterizes the probability of Coulomb tunneling. It is easy to show that the exponent argument behaves as $`C_{\mathrm{pyc}}/\sqrt{\lambda }=\alpha (r_{12}/r_{\mathrm{qm}})^2\rho ^{1/6}`$, where $`r_{12}`$ is the equilibrium distance between the interacting nuclei in their lattice sites, $`r_{\mathrm{qm}}`$ is the rms displacement of the nucleus due to zero-point vibrations in its lattice site, and $`\alpha 1`$ is a numerical factor which depends on a model of Coulomb tunneling. The usual condition is $`r_{\mathrm{qm}}r_{12}`$ (and the tunneling length $`r_{\mathrm{qm}}`$). The exponent argument $`C_{\mathrm{pyc}}/\sqrt{\lambda }`$ is typically large but decreases with growing $`\rho `$, making the Coulomb barrier more transparent. The tunneling is actually possible for closest neighbors (smallest $`r_{12}`$); the tunneling of more distant nuclei (higher $`r_{12}`$) is exponentially suppressed. Elastic lattice properties specify $`r_{\mathrm{qm}}`$ and $`\alpha `$, and are, thus, most important for the reaction rate. The presence of different ion species, lattice impurities and imperfections may drastically affect the rate svh69 .
### III.5 Thermally enhanced pycnonuclear regime
The thermally enhanced pycnonuclear burning occurs with increasing $`T`$; it operates svh69 in a relatively narrow temperature interval $`0.5T_p/\mathrm{ln}(1/\sqrt{\lambda })<T<\mathrm{\hspace{0.33em}0.5}T_p`$. In this interval the majority of the nuclei occupy their ground states in a strongly coupled quantum Coulomb system, but the main contribution to the reaction rate comes from a tiny fraction of nuclei which occupy excited bound energy states. The increase of the excitation energy increases the penetrability of the Coulomb barrier, and makes the excited states more efficient than the ground state.
The thermally enhanced pycnonuclear regime has been studied less accurately than the zero-temperature pycnonuclear regime. Salpeter and Van Horn svh69 calculated the thermally enhanced pycnonuclear reaction rate for models 1 and 2 of a bcc lattice in the WKB approximation. The spectrum of excited quantum states was determined for a relative motion of interacting nuclei in an anisotropic harmonic oscillator field; the summation over discrete quantum states in the expression for the reaction rate was replaced by the integration. According to their Eq. (45), the enhancement of the reaction rate is approximately described by
$$\frac{R_{\mathrm{pyc}}(T)}{R_{\mathrm{pyc}}(0)}1=\frac{\mathrm{\Omega }}{\lambda ^{1/2}}\mathrm{exp}\left(\mathrm{\Lambda }\frac{T_p}{T}+\frac{\mathrm{\Omega }_1}{\sqrt{\lambda }}\mathrm{e}^{\mathrm{\Lambda }_1T_p/T}\right),$$
(26)
where $`\mathrm{\Omega }`$, $`\mathrm{\Omega }_1`$, $`\mathrm{\Lambda }`$, and $`\mathrm{\Lambda }_1`$ are model-dependent dimensionless constants. The exponent $`\mathrm{exp}(\mathrm{\Lambda }T_p/T)`$ reflects the Boltzmann probability to occupy excited quantum states while the double exponent $`\mathrm{exp}\{(\mathrm{\Omega }_1/\sqrt{\lambda })\mathrm{e}^{\mathrm{\Lambda }_1T_p/T}\}`$ describes the enhancement itself. In this case the characteristic energy of the reacting nuclei is
$$E_{\mathrm{pk}}C_1\mathrm{}\omega _p+C_2\frac{Z^2e^2}{a}\mathrm{exp}\left(\mathrm{\Lambda }_1\frac{T_p}{T}\right),$$
(27)
where $`C_1`$ and $`C_2`$ are new dimensionless constants ($`1`$). When $`T`$ increases from $`T=0`$ to $`T0.5T_p`$, the characteristic energy $`E_{\mathrm{pk}}`$ increases from the ground state level, $`E_{\mathrm{pk}}\mathrm{}\omega _p`$, to the top of the Coulomb potential well, $`E_{\mathrm{pk}}Z^2e^2/a`$.
The thermally enhanced pycnonuclear reaction rate was studied also by Kitamura and Ichimaru ki95 adopting the formalism of Ogata, Iyetomi and Ichimaru oii91 (Sections III.3 and III.4). The relative motion of interacting nuclei was described by a model radial Schrödinger equation which employed the angle-averaged static MC plasma screening potential. The excited energy states were determined from the solution of this equation. Such an approach seems to be oversimplified. It gives the temperature dependence of the reaction rate (Eqs. (14) and (15) in Ref. ki95 ) which, functionally, differs from the temperature dependence, Eq. (27), predicted by Salpeter and Van Horn. Nevertheless, numerically, both temperature dependencies at $`T<\mathrm{\hspace{0.33em}0.5}T_p`$ are in a reasonable qualitative agreement.
We expect that the reaction rate in the thermally enhanced pycnonuclear regime will be further elaborated in the future.
### III.6 The intermediate thermo-pycnonuclear regime
The intermediate thermo-pycnonuclear regime is realized at temperatures $`TT_p`$ (roughly, at $`T_p/2<T<T_p`$) which separate the domains of quantum and classical ion systems. The main contribution to the reaction rate stems then from nuclei which are either slightly bound, or slightly unbound, with respect to their potential wells. The calculation of the reaction rate in this regime is complicated. We will describe this rate by a phenomenological expression presented in the following section.
### III.7 Single analytical approximation in all regimes
Let us propose a phenomenological expression for the reaction rate which combines all the five burning regimes:
$`R`$ $`=`$ $`R_{\mathrm{pyc}}(T=0)+\mathrm{\Delta }R(T),\mathrm{\Delta }R(T)={\displaystyle \frac{n_i^2}{2}}S(E_{\mathrm{pk}}){\displaystyle \frac{\mathrm{}}{mZ^2e^2}}PF,`$
$`F`$ $`=`$ $`\mathrm{exp}\left(\stackrel{~}{\tau }+C_{\mathrm{sc}}\stackrel{~}{\mathrm{\Gamma }}\phi \mathrm{e}^{\mathrm{\Lambda }T_p/T}\mathrm{\Lambda }{\displaystyle \frac{T_p}{T}}\right),P={\displaystyle \frac{8\pi ^{1/3}}{\sqrt{3}\mathrm{\hspace{0.33em}2}^{1/3}}}\left({\displaystyle \frac{E_\mathrm{a}}{\stackrel{~}{T}}}\right)^\gamma .`$ (28)
In this case, $`\phi =\sqrt{\mathrm{\Gamma }}/[(C_{\mathrm{sc}}^4/9)+\mathrm{\Gamma }^2]^{1/4}`$; $`R_{\mathrm{pyc}}(T=0)`$ is the zero-temperature pycnonuclear reaction rate (Section III.4); $`\mathrm{\Delta }R(T)`$ is the temperature-dependent part (with a product of an exponential function $`F`$ and a pre-exponent $`P`$). The quantities $`\stackrel{~}{\tau }`$ and $`\stackrel{~}{\mathrm{\Gamma }}`$ are similar to the familiar quantities $`\tau `$ and $`\mathrm{\Gamma }`$, but contain a “renormalized” temperature $`\stackrel{~}{T}`$:
$$\stackrel{~}{\tau }=3\left(\frac{\pi }{2}\right)^{2/3}\left(\frac{E_\mathrm{a}}{\stackrel{~}{T}}\right)^{1/3},\stackrel{~}{\mathrm{\Gamma }}=\frac{Z^2e^2}{a\stackrel{~}{T}},\stackrel{~}{T}=\sqrt{T^2+C_T^2T_p^2},$$
(29)
where $`C_T`$ is a dimensionless renormalization parameter ($`1`$). For high temperatures $`TT_p`$ we have $`\stackrel{~}{\tau }\tau `$, $`\stackrel{~}{\mathrm{\Gamma }}\mathrm{\Gamma }`$, and $`\stackrel{~}{T}T`$. In this case the temperature dependent term tends to $`\mathrm{\Delta }R(T)R_{\mathrm{th}}(T)R_{\mathrm{pyc}}`$, and Eq. (28) reproduces the thermonuclear reaction rate (Sections III.2 and III.3). At low temperatures $`T<T_p`$ the quantities $`\stackrel{~}{\tau }`$, $`\stackrel{~}{\mathrm{\Gamma }}`$ and $`\stackrel{~}{T}`$, roughly speaking, contain “the quantum temperature” $`T_p`$ rather than the real temperature $`T`$ in the original quantities $`\tau `$, $`\mathrm{\Gamma }`$ and $`T`$. In the limit of $`T0`$ we obtain $`\stackrel{~}{\mathrm{\Gamma }}=1/[\sqrt{\lambda }(72\pi )^{1/6}C_T]`$ and $`\stackrel{~}{\tau }=\left(3\sqrt{\pi /\lambda }\right)/\left(2^{7/6}C_T^{1/3}\right)`$.
At this point, let us require that at $`TT_p`$ the factor $`\mathrm{exp}(\stackrel{~}{\tau })`$ in the exponential function $`F`$, Eq. (28), reduces to $`\mathrm{exp}(C_{\mathrm{exp}}/\sqrt{\lambda })`$ in the exponential function $`F_{\mathrm{pyc}}`$, Eq. (23). This would allow us to obey Eq. (26) by satisfying the equality
$$3\sqrt{\pi }/(2^{7/6}C_T^{1/3})=C_{\mathrm{exp}}.$$
(30)
Taking $`C_{\mathrm{exp}}`$ we can determine $`C_T`$ (see Table 1). The double exponent factor in $`F`$, Eq. (28), will correspond to the double exponent factor in Eq. (26). Strictly speaking, Eq. (26) contains two different constants $`\mathrm{\Lambda }`$ and $`\mathrm{\Lambda }_1`$. However, taking into account the uncertainties of $`R`$ in the thermally enhanced pycnonuclear regime (Section III.5), we replace two constants by one.
Finally, the quantity $`\gamma `$ in Eq. (28) should be taken in such a way as to reproduce the correct limit $`\gamma _1=2/3`$ at $`TT_p`$ (Sect. III.2) and $`\gamma _2=(2/3)(C_{\mathrm{pl}}+0.5)`$ at $`TT_p`$ (see Eq. (26)). The natural interpolation expression for $`\gamma `$ would be
$$\gamma =(T^2\gamma _1+T_p^2\gamma _2)/(T^2+T_p^2).$$
(31)
In addition, we need the reaction energy $`E_{\mathrm{pk}}`$ to evaluate the astrophysical factor $`S(E_{\mathrm{pk}})`$. Since the $`S`$-factor is a slowly varying function of energy, it is reasonable to approximate $`E_{\mathrm{pk}}`$ by the expression
$$E_{\mathrm{pk}}=\mathrm{}\omega _p+\left(\frac{Z^2e^2}{a}+\frac{T\tau }{3}\right)\mathrm{exp}\left(\frac{\mathrm{\Lambda }T_p}{T}\right)$$
(32)
which combines the expressions in the thermonuclear and pycnonuclear regimes. To avoid the introduction of many fit parameters (unnecessary at the present state of investigation), we set $`C_1=C_2=1`$ in Eq. (27).
Thus, we propose to adopt the analytic expression (28) using the following parameters:
(1) $`C_{\mathrm{sc}}=1.0754`$ for the case of strong plasma screening in the thermonuclear regime (Section III.3);
(2) $`C_{\mathrm{exp}}`$, $`C_{\mathrm{pyc}}`$, and $`C_{\mathrm{pl}}`$ for conditions of zero-temperature pycnonuclear burning (see Table 1 and Section III.4);
(3) The quantum-temperature constant $`C_T`$ (Section III.4), which is important at $`TT_p`$ and expressed through $`C_{\mathrm{exp}}`$ via Eq. (30); the corresponding values of $`C_T`$ are listed in Table 1;
(4) The last constant $`\mathrm{\Lambda }`$, that is important at $`TT_p`$, is still free.
We have checked (Fig. 6) that taking the optional model 1 from Table 1 and the value $`\mathrm{\Lambda }=0.5`$ results in a good agreement with the carbon burning rate calculated at $`\rho 10^910^{10}`$ g cm<sup>-3</sup> and $`T<\mathrm{\hspace{0.33em}0.5}T_p`$ from Eq. (45) of Salpeter and Van Horn svh69 (for the thermally enhanced pycnonuclear regime). (Notice that model 2 requires slightly lower $`\mathrm{\Lambda }0.45`$.) Taking $`\mathrm{\Lambda }=0.35`$ leads to a noticeably higher reaction rate at $`T<\mathrm{\hspace{0.33em}0.5}T_p`$, while taking $`\mathrm{\Lambda }=0.65`$ leads to a noticeably lower rate.
Accordingly, for any model of zero-temperature pycnonuclear burning from Table 1 we suggest to adopt $`\mathrm{\Lambda }=0.5`$ as optional, $`\mathrm{\Lambda }=0.35`$ to maximize and $`\mathrm{\Lambda }=0.65`$ to minimize the reaction rate. In particular, model 1 with $`\mathrm{\Lambda }=0.5`$ seems to be the “most optional”; our limiting model 10 from Table 1 with $`\mathrm{\Lambda }=0.35`$ is expected to give the upper theoretical limit for the reaction rate, while the other limiting model 11 with $`\mathrm{\Lambda }=0.65`$ is expected to give the lower theoretical limit. We also need the astrophysical factor $`S(E)`$, which was described in Section II for the carbon burning. We could easily introduce additional constants to tune our phenomenological model when precise calculations of reaction rates appear in the future.
For illustration, Fig. 6 shows the temperature dependence of the carbon burning rate at $`\rho =5\times 10^9`$ g cm<sup>-3</sup>. The solid curve is the most optimal model (based on both – zero-temperature and thermally enhanced – pycnonuclear burning models of Salpeter and Van Horn svh69 for the bcc static lattice). The double hatched region shows assumed uncertainties of this model associated with variations of $`\mathrm{\Lambda }`$ from 0.35 to 0.65 (as if we accept the zero-temperature model but question the less elaborated model of thermal enhancement). The singly hatched region indicates overall uncertainties (limited by the models of the maximum and minimum reaction rates). The lower long-dashed line is obtained assuming classical thermonuclear burning without any screening (Section III.2). The upper long-dashed line is calculated using the formalism of thermonuclear burning with screening (Section III.3). The screening enhancement of the reaction rate becomes stronger with the decrease of $`T`$. The formalism for describing this enhancement is expected to be valid at $`T>T_p`$, but we intentionally extend the upper long-dashed curve to $`T=0.5T_p`$, where the formalism breaks down and the curve diverges from the expected (solid) curve. The short-dash curve is calculated from the equations of Salpeter and Van Horn svh69 derived in the thermally enhanced pycnonuclear regime (model 1) and valid at $`T<\mathrm{\hspace{0.33em}0.5}T_p`$. We intentionally extend the curve to higher $`T`$, where the formalism of thermally enhanced pycnonuclear burning becomes invalid and the curve diverges from the expected curve. Our phenomenological solid curve provides a natural interpolation at $`TT_p`$ between the short-dashed curve and the upper long-dashed curve.
More complicated expressions for the reaction rate $`R`$ in wide ranges of $`\rho `$ and $`T`$ were proposed by Kitamura kitamura00 who combined the results of recent calculations of $`R`$ in the different regimes. His expressions are mainly based on the results of Refs. oii91 ; iovh92 ; oiv93 ; ki95 ; ogata97 which are not free of approximations (as discussed in Sections III.3, III.4, and III.5). In contrast to our formula, Kitamura took into account the effects of electron screening (finite polarizability of the electron gas) and considered the case of equal and non-equal reacting nuclei. However, the electron screening effects are relatively weak; their strict inclusion in the pycnonuclear regime represents a complicated problem. We do not include them but, instead, take into account theoretical uncertainties of the reaction rates without electron screening. We have checked that the results by Kitamura kitamura00 for carbon burning in the most important $`T\rho `$ domain lie well within these uncertainties.
Our formula gives a smooth behavior of the reaction rate as a function of temperature and density, without any jump at the melting temperature $`T=T_\mathrm{m}`$. We do not expect any strong jump of such kind since the liquid-solid phase transition in dense stellar matter is tiny. A careful analysis shows the absence of noticeable jumps of transport coefficients baikoetal98 and the neutrino emissivity owing to electron-nucleus bremsstrahlung. A direct example is given by the theory of nuclear burning. Ichimaru and Kitamura ik99 predicted a noticeable jump of the reaction rate at $`T=T_\mathrm{m}`$, while a more careful analysis of Kitamura kitamura00 considerably reduced this jump.
## IV Carbon burning and ignition in dense matter
In this section we will analyze the rate of the <sup>12</sup>C+<sup>12</sup>C reaction as a function of $`T`$ and $`\rho `$ and investigate the conditions for carbon burning in dense stellar matter.
Because the probability for Coulomb tunneling depends exponentially on plasma parameters, changes in density $`\rho `$ and temperature $`T`$ have dramatic effects on the burning rate $`R`$. In thermonuclear regimes (Sections III.2 and III.3) the <sup>12</sup>C+<sup>12</sup>C rate is more sensitive to changes in temperature $`T`$ than in density $`\rho `$. On the contrary, in pycnonuclear regimes (Sections III.4 and III.5) the rate depends significantly on the density $`\rho `$. For instance, if $`T`$ decreases from $`3\times 10^9`$ K to $`3\times 10^8`$ K at $`\rho =5\times 10^9`$ g cm<sup>-3</sup> (Fig. 6; thermonuclear burning with strong screening), the reaction rate drops by $`20`$ orders of magnitude. Neglecting the enhancement due to plasma screening, the rate will drop by ten more orders of magnitude. An increase in density $`\rho `$ from $`10^8`$ g cm<sup>-3</sup> to $`10^{11}`$ g cm<sup>-3</sup> at $`T<\mathrm{\hspace{0.33em}3}\times 10^7`$ K (in the zero-temperature pycnonuclear regime) results in a rate increase of $`100`$ orders of magnitude (Fig. 5). Note that no carbon can survive in a degenerate matter at $`\rho >3.90\times 10^{10}`$ g cm<sup>-3</sup> because of the double electron capture $`{}_{6}{}^{12}\mathrm{C}{}_{5}{}^{12}\mathrm{B}{}_{4}{}^{12}\mathrm{Be}`$ (e.g., Shapiro and Teukolsky st83 ). The electron capture has a well defined density threshold, $`3.9\times 10^{10}`$ g cm<sup>-3</sup>, and proceeds quickly after the threshold is exceeded. We will ignore this process in the present section.
The strong dependence of the rate $`R`$ on density $`\rho `$ and temperature $`T`$ leads to huge variations of the characteristic time scale $`\tau _{\mathrm{burn}}=n_i/R`$ for carbon burning. Figure 4 shows two solid lines in the $`\rho T`$ plane, along which $`\tau _{\mathrm{burn}}=1`$ s and $`10^{10}`$ years (nearly the Universe age), respectively. They are calculated using the most optional carbon burning model (model 1 from Table 1, $`\mathrm{\Lambda }=0.5`$). The lines are almost horizontal in the thermonuclear burning regime ($`R`$ is a slowly varying function of $`\rho `$) and almost vertical in the pycnonuclear regime ($`R`$ is a slowly varying function of $`T`$). The bending part of the lines corresponds to the thermally enhanced pycnonuclear and intermediate thermo-pycno nuclear regimes. At $`T`$ and $`\rho `$ above the upper line the burning time is even shorter than 1 s; at these conditions no carbon will survive in the dense matter of astrophysical objects. For conditions below the lower solid line, $`\tau _{\mathrm{burn}}`$ is longer than $`10^{10}`$ years, and carbon burning can be disregarded for most applications.
Thus, the studies of carbon burning can be focused on the narrow strip in the $`\rho T`$ plane between the lines of $`\tau _{\mathrm{burn}}=1`$ s and $`\tau _{\mathrm{burn}}=10^{10}`$ yr. Hatched regions show theoretical uncertainties of each line (limited by the maximum and minimum reaction rate models, Section III.7). The uncertainties are relatively small in the thermonuclear regime where all models give nearly the same reaction rate. The uncertainties are higher in other burning regimes.
Having a carbon burning model we can plot the carbon ignition curve. This curve is a necessary ingredient for modeling nuclear explosions of massive white dwarfs (producing supernova Ia events, so important for cosmology; see, e.g., Refs. SNI ; filippenko04 ) and for modeling carbon explosions of matter in accreting neutron stars (viable models of superbursts observed recently from some accreting neutron stars; e.g., Refs. CB01 ; SB02 ; cn04 ).
The ignition curve is commonly determined as the line in the $`\rho T`$ plane (Fig. 4) where the nuclear energy generation rate is equal to the local neutrino emissivity of dense matter (the neutrino emission carries the generated energy out of the star). At higher $`\rho `$ and $`T`$ (above the curve) the nuclear energy generation rate exceeds the neutrino losses and carbon ignites. In Fig. 4 we present the carbon ignition (solid) curve, calculated using the most optional model of carbon burning, together with its uncertainties (limited by the minimum and maximum rate models). The neutrino energy losses are assumed to be produced by plasmon decay and by electron-nucleus bremsstrahlung. The neutrino emissivity due to plasmon decay is obtained from extended tables calculated by M. E. Gusakov (unpublished); they are in good agreement with the results by Itoh et al. itohetal92 . The neutrino bremsstrahlung emissivity is calculated using the formalism of Kaminker et al. kaminkeretal99 , which takes into account electron band structure effects in crystalline matter.
For $`\rho <\mathrm{\hspace{0.33em}10}^9`$ g cm<sup>-3</sup> theoretical uncertainties of the ignition curve are seen to be small. They become important at $`\rho >\mathrm{\hspace{0.33em}10}^9`$ g cm<sup>-3</sup> and $`T(13)\times 10^8`$ K in the intermediate thermo-pycnonuclear burning regime and the thermally enhanced burning regime. This $`\rho T`$ range is appropriate for central regions of massive and warm white dwarfs which may produce type Ia supernova explosions. Lower $`T`$ are also interesting for these studies (e.g., Baraffe et al. bhw04 ).
If we formally continue the ignition curve to lower $`T`$, it will bend and shift to lower densities, where the nuclear burning time scale $`\tau _{\mathrm{burn}}`$ is exceptionally slow exceeding the age of the Universe. The bend is associated with a very weak neutrino emission at $`T<\mathrm{\hspace{0.33em}10}^8`$ K. These parts of the ignition curve are oversimplified because at low $`T`$ the energy outflow produced by thermal conduction becomes more efficient than the outflow due to the neutrino emission. These parts are shown by the long-dash line (and their uncertainties are indicated by thin dash-and-dot lines). Unfortunately, the conduction energy outflow is non-local and “non-universal”. It depends on specific conditions of the burning environment (a white dwarf core or a neutron star crust) and the associated thermal conductivity (provided mainly by strongly degenerate electrons). In this case the ignition becomes especially complicated. A very crude estimate shows that the ignition curve, governed by the thermal conduction, is nearly vertical and close to the $`\tau _{\mathrm{burn}}=10^{10}`$ yr curve in the range of $`T`$ from $`10^8`$ K to $`10^6`$ K in Fig. 4. At $`T<\mathrm{\hspace{0.33em}10}^6`$ K the curve is strongly affected by the thermal conductivity model. In a cold ideal carbon crystal, umklapp processes of electron-phonon scattering are frozen out (e.g., Ref. ry82 ). Under these conditions the electron conduction is determined by inefficient normal electron-phonon scattering, leading to high conductivity values. This shifts the ignition to higher $`\rho `$. On the other hand, carbon matter may contain randomly located ions of other elements (charged impurities) which can keep the electron Coulomb scattering rather efficient and maintain a low electron thermal conductivity. In this case the ignition curve at $`T<\mathrm{\hspace{0.33em}10}^6`$ K remains nearly vertical.
## V Conclusion
The goal of this paper was to develop a phenomenological formalism for calculating fusion reaction rates between identical nuclei. This formalism should be applicable for a broad range of thermonuclear and pycnonuclear burning scenarios. It involves a generalized treatment for calculation of the fusion probability at low energies and the development of a single simple phenomenological expression for the fusion rate valid in a wide range of temperatures and densities.
We have introduced a generalized model approach for calculating the $`S`$-factor of heavy-ion fusion reactions relevant for stellar nucleosynthesis processes. We have demonstrated the applicability and reliability of the approach by calculating the astrophysical factor $`S(E)`$ for the carbon fusion reaction <sup>12</sup>C+<sup>12</sup>C (Section II) and by comparing the theoretical results with experimental data.
Furthermore, we have analyzed (Section III) previous calculations of the fusion rate for identical nuclei in stellar matter, with emphasis on the complicated problem of Coulomb barrier penetration in a dense-plasma environment. Combining the results of previous studies, we have proposed a single simple phenomenological expression for the fusion rate, valid in all five fusion regimes (that can be realized in the different $`\rho T`$ regions). Our formula contains adjustable parameters whose variations reflect theoretical uncertainties of the reaction rates.
For illustration, we have considered (Section IV) the efficiency of carbon burning in dense matter and the conditions for carbon ignition in white dwarf cores and neutron star crusts. We show that carbon burning is actually important in a sufficiently narrow $`\rho T`$ strip which is mainly determined by the temperature $`T(415)\times 10^8`$ K as long as $`\rho <\mathrm{\hspace{0.33em}3}\times 10^9`$ g cm<sup>-3</sup>, and by the density $`\rho (350)\times 10^9`$ g cm<sup>-3</sup> as long as $`T<\mathrm{\hspace{0.33em}10}^8`$ K. On the basis of these results we suggest that the current knowledge of nuclear fusion is sufficient to understand the main features of carbon burning in stellar matter, especially at $`\rho <\mathrm{\hspace{0.33em}3}\times 10^9`$ g cm<sup>-3</sup>.
We have focused on the simplest case of heavy-ion burning in a one-component Coulomb system; particularly, in a perfect crystal. There is no doubt that dense matter of white dwarfs and neutron stars are more complicated and require a more complex approach taking into account mixtures of different heavy nuclei and imperfections in dense matter. The complexity ranges from essentially two-component plasma conditions anticipated in the carbon-oxygen cores of white dwarfs to the multi-component isotope distribution in the ashes of accreting neutron stars woos04 .
In a forthcoming paper we will expand the presented analysis to the case of the fusion rates between different isotopes. We will employ this formalism for calculating the $`S`$-factors for a broad range of heavy-ion fusion reactions. We will include the results in a pycno-thermonuclear reaction network and simulate the nucleosynthesis in high density stellar matter.
###### Acknowledgements.
We are grateful to H. DeWitt for critical remarks and to M. Gusakov for providing the tables of neutrino emissivities due to plasmon decay. This work was partially supported by The Joint Institute for Nuclear Astrophysics (JINA) NSF PHY 0216783, Fundação de Amparo à Pesquisa do Estado de São Paulo (FAPESP), CONACYT (México), DoE grant DE-F05-96ER-40983 and BMBF (Germany), under the project 06 MT 193, RFBR (grants 03-07-90200 and 05-02-16245) and RLSSP (project 1115.2003.2). |
warning/0506/gr-qc0506057.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Rotating solutions of cosmological Einstein gravity in $`D`$ dimensions, $`R_{\mu \nu }=(D1)\mathrm{\Lambda }g_{\mu \nu }`$, have been constructed recently , extending earlier $`\mathrm{\Lambda }=0`$ solutions of , themselves generalizations of the well-known $`D`$=4 metrics of and , and of in $`D`$=5. These geometries provide a useful application of our recent generalized “conserved charge” definitions, which are also extensions – of the original ADM , and AD charges – to cover wider classes of actions : We will compute the energy and angular momenta of these new solutions, as well as of the $`D=3`$ BTZ metric as calculated within topologically massive gravity.
Gravity theories have been historically endowed with a variety of seemingly different charge definitions, with different degrees of applicability and coordinate invariance. This topic has also seen much very recent activity, for example . A summary and comparison of some of them is given in which also includes a computation of the charges for Kerr-AdS black holes, using thermodynamic arguments; see also . Our results will agree with those, but we emphasize that in a general context, certain coincidences between charge definitions are suspect: For example, the frequently invoked “Komar” charges, are in general not applicable, being highly gauge-dependent .
## 2 Mass and Angular Momenta of Kerr-AdS
Let us briefly recapitulate the formulations of . The field equations of any metric model coupled to a (necessarily covariantly conserved) matter source $`\tau _{\mu \nu }`$ are
$`{\displaystyle \frac{\delta I}{\delta g_{\mu \nu }}}\mathrm{\Phi }_{\mu \nu }(g,R,R,..)=\kappa \tau _{\mu \nu },`$ (1)
where $`\mathrm{\Phi }_{\mu \nu }`$ is an identically conserved tensor that can depend on curvatures and their derivatives. Decompose the metric into the sum of a background “vacuum”, $`\overline{g}_{\mu \nu }`$ (which solves (1) for $`\tau _{\mu \nu }=0`$), plus a deviation $`h_{\mu \nu }`$, not necessarily small, that vanishes sufficiently rapidly far from the matter source: $`g_{\mu \nu }=\overline{g}_{\mu \nu }+h_{\mu \nu }`$. The field equations can be divided into a part linear in $`h_{\mu \nu }`$ plus a non-linear remainder, which (with $`\tau _{\mu \nu }`$) constitutes the total source $`T_{\mu \nu }`$. If the background $`\overline{g}_{\mu \nu }`$ admits Killing vectors $`\overline{\xi }_\mu `$, obeying $`\overline{}_\mu \overline{\xi }_\nu +\overline{}_\nu \overline{\xi }_\mu =0`$, then, up to normalization factors (which we shall fix later), the conserved Killing charges are
$`Q^\mu (\overline{\xi })={\displaystyle _{}}d^{D1}x\sqrt{\overline{g}}T^{\mu \nu }\overline{\xi }_\nu ={\displaystyle _\mathrm{\Sigma }}𝑑S_i^{\mu i}.`$ (2)
Here $`\mathrm{\Sigma }`$ is a $`D2`$ dimensional space-like asymptotic hypersurface of the space $``$ and $`^{\mu i}`$ is an anti-symmetric tensor, whose explicit form is model-dependent. For Einstein’s theory with a cosmological constant,
$`Q^\mu ={\displaystyle \frac{1}{4\mathrm{\Omega }_{D2}G_D}}{\displaystyle _\mathrm{\Sigma }}`$ $`dS_i`$ $`\{\overline{\xi }_\nu \overline{}^\mu h^{i\nu }\overline{\xi }_\nu \overline{}^ih^{\mu \nu }+\overline{\xi }^\mu \overline{}^ih\overline{\xi }^i\overline{}^\mu h`$ (3)
$`+h^{\mu \nu }\overline{}^i\overline{\xi }_\nu h^{i\nu }\overline{}^\mu \overline{\xi }_\nu +\overline{\xi }^i\overline{}_\nu h^{\mu \nu }\overline{\xi }^\mu \overline{}_\nu h^{i\nu }+h\overline{}^\mu \overline{\xi }^i\},`$
where $`i`$ takes values in $`1,2,\mathrm{}D2`$ and the charge is normalized as shown, by dividing with the $`D`$-dimensional Newton’s constant and the solid angle. These charges are background gauge invariant under the diffeomorphisms $`\delta _\zeta h_{\mu \nu }=\overline{}_\mu \zeta _\nu +\overline{}_\nu \zeta _\mu `$: $`\delta _\zeta Q^\mu =0`$.
Let us now calculate the conserved charges of the metrics for $`D>3`$. \[We shall treat the special $`D=3`$ case at the end\]. They have the Kerr-Schild form
$`ds^2=d\overline{s}^2+{\displaystyle \frac{2M}{U}}(k_\mu dx^\mu )^2,`$ (4)
in terms of the de Sitter metric
$`d\overline{s}^2`$ $`=`$ $`W(1\mathrm{\Lambda }r^2)dt^2+Fdr^2+{\displaystyle \underset{i=1}{\overset{N+ϵ}{}}}{\displaystyle \frac{r^2+a_i^2}{1+\mathrm{\Lambda }a_i^2}}d\mu _i^2+{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{r^2+a_i^2}{1+\mathrm{\Lambda }a_i^2}}\mu _i^2d\varphi _i^2`$ (5)
$`+{\displaystyle \frac{\mathrm{\Lambda }}{W(1\mathrm{\Lambda }r^2)}}\left({\displaystyle \underset{i=1}{\overset{N+ϵ}{}}}{\displaystyle \frac{(r^2+a_i^2)\mu _id\mu _i}{1+\mathrm{\Lambda }a_i^2}}\right)^2.`$
Here $`ϵ=0/1`$ for odd/even, dimensions and $`D=2N+1+ϵ`$. The null 1-form reads
$`k_\mu dx^\mu `$ $`=`$ $`Fdr+Wdt{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{a_i\mu _i^2}{1+\mathrm{\Lambda }a_i^2}}d\varphi _i,`$ (6)
with
$`U`$ $``$ $`r^ϵ{\displaystyle \underset{i=1}{\overset{N+ϵ}{}}}{\displaystyle \frac{\mu _i^2}{r^2+a_i^2}}{\displaystyle \underset{j=1}{\overset{N}{}}}(r^2+a_j^2),W{\displaystyle \underset{i=1}{\overset{N+ϵ}{}}}{\displaystyle \frac{\mu _i^2}{1+\mathrm{\Lambda }a_i^2}},F{\displaystyle \frac{1}{1\mathrm{\Lambda }r^2}}{\displaystyle \underset{i=1}{\overset{N+ϵ}{}}}{\displaystyle \frac{r^2\mu _i^2}{r^2+a_i^2}}.`$ (7)
To find the energy and angular momenta corresponding to (4), we must compute the charges $`Q^0`$ for the corresponding Killing vectors: for the energy we shall take $`\overline{\xi }^\mu =(1,\stackrel{}{0})`$ and each angular momentum has the appropriate unit entry $`(0,\mathrm{}1_i\mathrm{}0)`$. Then
$`Q^0`$ $`=`$ $`{\displaystyle \frac{1}{4\mathrm{\Omega }_{D2}G_D}}{\displaystyle _\mathrm{\Sigma }}𝑑S_r\left\{g_{00}\overline{}^0h^{r0}+g_{00}\overline{}^rh^{00}+h^{0\nu }\overline{}^r\overline{\xi }_\nu h^{r\nu }\overline{}^0\overline{\xi }_\nu +\overline{}_\nu h^{r\nu }\right\}.`$ (8)
Using the energy Killing vector, we obtain<sup>1</sup><sup>1</sup>1We are assuming that the background spacetime is AdS rather than dS, whose cosmological horizon causes complications. Some of these issues were addressed in . For details of acceptable asymptotic falloff to (A)dS in various dimensions, we refer to .
$`E_D`$ $`=`$ $`{\displaystyle \frac{1}{4\mathrm{\Omega }_{D2}G_D}}{\displaystyle _\mathrm{\Sigma }}dS_r\{g_{00}g^{rr}_rh^{00}+{\displaystyle \frac{1}{2}}h^{00}g^{rr}_rg_{00}{\displaystyle \frac{m}{U}}g^{00}_rg_{00}+2m_rU^1`$ (9)
$`+{\displaystyle \frac{2m}{U}}g^{rr}_rg_{rr}{\displaystyle \frac{m}{U}}g^{rr}k^ik^j_rg_{ij}+{\displaystyle \frac{m}{U}}g^{ij}_rg_{ij}\}.`$
To compute $`E_D`$, one needs the large $`r`$ behavior of the integrand $`I`$ of (9); since
$`g_{00}W\mathrm{\Lambda }r^2,F{\displaystyle \frac{1}{\mathrm{\Lambda }r^2}},Ur^{D3},k^\varphi `$ $`{\displaystyle \frac{a_\varphi }{r^2}},`$ (10)
then
$`I={\displaystyle \frac{2m}{r^{D2}}}[(D1)W1].`$ (11)
For completeness, let us also note how the determinant is calculated,
$`\text{det}g=W(1\mathrm{\Lambda }r^2)F{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{(r^2+a_i^2)\mu _i^2}{1+\mathrm{\Lambda }a_i^2}}\text{det}M.`$ (12)
Here $`M`$ is the matrix representing the coefficients of the form $`d\mu _id\mu _j`$ in the metric, which can be expressed as (no repeated index summation),
$`M_{ij}=A_i\delta _{ij}+B_iB_j+C_iC_j`$ (13)
where
$`A_i`$ $`=`$ $`{\displaystyle \frac{(r^2+a_i^2)}{1+\mathrm{\Lambda }a_i^2}},B_i=\sqrt{{\displaystyle \frac{(r^2+a_{N+ϵ}^2)}{1+\mathrm{\Lambda }a_{N+ϵ}^2}}}{\displaystyle \frac{\mu _i}{\mu _n}}`$
$`C_i`$ $`=`$ $`\sqrt{{\displaystyle \frac{\mathrm{\Lambda }}{W(1\mathrm{\Lambda }r^2)}}}({\displaystyle \frac{(r^2+a_i^2)}{1+\mathrm{\Lambda }a_i^2}}{\displaystyle \frac{(r^2+a_{N+ϵ}^2)}{1+\mathrm{\Lambda }a_{N+ϵ}^2}})\mu _i.`$ (14)
Then we have
$`\text{det}M={\displaystyle \underset{i=1}{\overset{N+ϵ1}{}}}A_i{\displaystyle \underset{i=1}{\overset{N+ϵ1}{}}}\left\{{\displaystyle \frac{B_i^2}{A_i}}+{\displaystyle \frac{C_i^2}{A_i}}+{\displaystyle \underset{ji}{\overset{N+ϵ1}{}}}{\displaystyle \frac{B_i^2C_i^2}{A_iA_j}}{\displaystyle \underset{ji}{\overset{N+ϵ1}{}}}{\displaystyle \frac{B_iB_jC_jC_i}{A_iA_j}}\right\}.`$ (15)
Inserting (14) in the above equation, one gets
$`\text{det}M={\displaystyle \frac{1}{W\mu _{N+ϵ}^2}}{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{1}{1+\mathrm{\Lambda }a_i^2}}.`$ (16)
Using equations (16,12,11) the energy of the $`D`$ dimensional rotating black hole becomes
$`E_D={\displaystyle \frac{m}{\mathrm{\Xi }}}{\displaystyle \underset{i=1}{\overset{\frac{D1ϵ}{2}}{}}}\left\{{\displaystyle \frac{1}{\mathrm{\Xi }_i}}(1ϵ)({\displaystyle \frac{1}{2}})\right\}.`$ (17)
where
$`\mathrm{\Xi }{\displaystyle \underset{i=1}{\overset{\frac{D1ϵ}{2}}{}}}(1+\mathrm{\Lambda }a_i^2),\mathrm{\Xi }_i1+\mathrm{\Lambda }a_i^2.`$ (18)
This expression reduces to the standard limits $`a_i0`$ and $`\mathrm{\Lambda }0`$, and agrees (up to a constant factor) with those of .
The computation of angular momenta follows along similar lines. Consider a given, say that $`i^{\text{th}}`$ (which we call the $`\varphi `$) component, i.e., the Killing vector $`\xi _{(i)}^\mu =(0,\mathrm{},0,1_i,0,..)`$. Then the corresponding Killing charge becomes
$`Q^0`$ $`=`$ $`{\displaystyle \frac{1}{4\mathrm{\Omega }_{D2}G_D}}{\displaystyle _\mathrm{\Sigma }}𝑑S_r\left\{g_{\varphi \varphi }\overline{}^0h^{r\varphi }g_{\varphi \varphi }\overline{}^rh^{0\varphi }+h^{0\nu }\overline{}^r\overline{\xi }_\nu h^{r\nu }\overline{}^0\overline{\xi }_\nu \right\}`$ (19)
$`=`$ $`{\displaystyle \frac{1}{4\mathrm{\Omega }_{D2}G_D}}{\displaystyle _\mathrm{\Sigma }}𝑑S_r\left\{g_{\varphi \varphi }g^{rr}g^{00}_rh_0^\varphi \right\}.`$
Once again the integrand can be calculated to be
$`I={\displaystyle \frac{(D1)2ma_i\mu _i^2}{r^{D2}(1+\mathrm{\Lambda }a_i^2)}}.`$ (20)
Putting the pieces together, the angular momentum is
$`J_i={\displaystyle \frac{ma_i}{\mathrm{\Xi }\mathrm{\Xi }_i}}.`$ (21)
This expression again agrees with . Note that, unlike in the energy expression, $`ϵ`$ does not appear here since even dimensional spaces have as many independent 2-planes as the odd dimensional spaces with one lower dimension.<sup>2</sup><sup>2</sup>2For even dimensions, there is a nice relation between the energy and the angular momentum $`E=_i\frac{J_i}{a_i}`$.
Having computed the desired conserved charges (17,21) for Kerr-AdS spacetimes in $`D>3`$, let us briefly turn our attention to the $`D=3`$ BTZ black hole . This solution has long been studied but we recompute the charges with our method for the sake of completeness. The BTZ black hole differs from its higher dimensional counterparts in one very important aspect: for it, AdS is not the correct-vacuum-background . The full metric is
$`ds^2=(M\mathrm{\Lambda }r^2)dt^2+{\displaystyle \frac{dr^2}{M+\mathrm{\Lambda }r^2+\frac{a^2}{4r^2}}}adtd\varphi +r^2d\varphi ^2,`$ (22)
The background metric corresponds to $`M=0`$ and AdS corresponds to $`M=1`$. Only AdS with $`J=0`$ is allowed for $`M<0`$: the others have naked singularities. So we consider $`M>0`$ and compute the charges following our calculations above (about the $`M=0`$ background.) We get the usual answers
$`E=M,J=a.`$ (23)
BTZ black holes also solve the more general topologically massive gravity equations, where the Einstein term is augmented by the Cotton tensor ,
$`G_{\mu \nu }+\mathrm{\Lambda }g_{\mu \nu }+{\displaystyle \frac{1}{\mu }}C_{\mu \nu }=\kappa \tau _{\mu \nu }.`$ (24)
Conserved charges in this model were obtained in , in terms of those of the Einstein model $`Q_E^\mu `$,
$`Q^\mu (\overline{\xi })`$ $`=`$ $`Q_E^\mu (\overline{\xi })+{\displaystyle \frac{1}{2\mu }}{\displaystyle 𝑑S_i\left\{ϵ^{\mu i\beta }𝒢_{\nu \beta }^L\overline{\xi }^\nu +ϵ_\beta ^{\nu i}𝒢_L^{\mu \beta }\overline{\xi }_\nu +ϵ^{\mu \nu \beta }𝒢_\beta ^{Li}\overline{\xi }_\nu \right\}}`$
$`+{\displaystyle \frac{1}{2\mu }}Q_E^\mu (ϵ\overline{}\overline{\xi }),`$
where $`Q_E^\mu (ϵ\overline{}\overline{\xi })`$ is the Einstein form but $`\overline{\xi }`$ is replaced with its curl. Once the contributions of the Cotton parts are computed the mass and the angular momentum of the BTZ black hole reads:
$`E=M{\displaystyle \frac{\mathrm{\Lambda }a}{\mu }},J=a{\displaystyle \frac{M}{\mu }},`$ (26)
a shift in values that may be compared with those for gravitational anyons , (linearized) solutions of TMG but not of pure D=3 Einstein.
## 3 Mass and Angular Momenta in Higher Curvature Models
We turn now to a slightly more formal exercise, which is to indicate the stability of our generic charge definition framework as it applies to a wider range of models, specifically higher derivative gravities. While Kerr-like solutions to $`R+R^2`$ gravity models have yet to be discovered, it is not unlikely that they would approach the Einstein ones asymptotically. In that case, we could compute their conserved charges-defined as integrals at infinity, using the definitions for generic quadratic models . Let us stick to the quadratic models of the form<sup>3</sup><sup>3</sup>3Note that we changed normalization of the cosmological constant compared to the previous section.
$`I={\displaystyle d^Dx\sqrt{g}\left\{\frac{R}{2\kappa }+2\mathrm{\Lambda }_0+\alpha R^2+\beta R_{\mu \nu }^2+\gamma (R_{\mu \nu \rho \sigma }^24R_{\mu \nu }^2+R^2)\right\}}.`$ (27)
This model allows constant curvature spacetimes with an effective cosmological constant given as
$`\mathrm{\Lambda }={\displaystyle \frac{1}{4f(\alpha ,\beta ,\gamma _)\kappa }}\left\{1\pm \sqrt{1+8\kappa f(\alpha ,\beta ,\gamma )\mathrm{\Lambda }_0}\right\}\text{for}\text{ }f(\alpha ,\beta ,\gamma )0,`$ (28)
where
$`f(\alpha ,\beta ,\gamma )={\displaystyle \frac{(D4)}{(D2)^2}}(D\alpha +\beta )+{\displaystyle \frac{\gamma (D4)(D3)}{(D2)(D1)}}.`$ (29)
When the bare cosmological constant vanishes ($`\mathrm{\Lambda }_0=0`$ ), (A)dS spaces are still allowed and one has the $`+`$ sign branch in (28). Conserved charges in this model, which we quote below, were defined in
$`Q^\mu (\overline{\xi })`$ $`=`$ $`\left\{{\displaystyle \frac{1}{\kappa }}+{\displaystyle \frac{4\mathrm{\Lambda }D\alpha }{D2}}+{\displaystyle \frac{4\mathrm{\Lambda }\beta }{D1}}+{\displaystyle \frac{4\mathrm{\Lambda }\gamma (D4)(D3)}{(D2)(D1)}}\right\}{\displaystyle d^{D1}x\sqrt{\overline{g}}\overline{\xi _\nu }𝒢_L^{\mu \nu }}`$ (30)
$`+(2\alpha +\beta ){\displaystyle 𝑑S_i\sqrt{g}\left\{\overline{\xi }^\mu \overline{}^iR_L+R_L\overline{}^\mu \overline{\xi }^i\overline{\xi }^i\overline{}^\mu R_L\right\}}`$
$`+\beta {\displaystyle 𝑑S_i\sqrt{g}\left\{\overline{\xi }_\nu \overline{}^i𝒢_L^{\mu \nu }\overline{\xi }_\nu \overline{}^\mu 𝒢_L^{i\nu }𝒢_L^{\mu \nu }\overline{}^i\overline{\xi }_\nu +𝒢_L^{i\nu }\overline{}^\mu \overline{\xi }_\nu \right\}}.`$
where $`𝒢_L^{\mu \nu }`$ and $`R_L`$ are the linear parts of the Einstein tensor and the scalar curvature, respectively. The second and the third line vanish for Einstein spaces. The first line, on the other hand is just a factor times the cosmological Einstein theory’s charges (3). Therefore for asymptotic Kerr-AdS solutions, their conserved charges are given by the first term in (30), under the condition (28). Let us specifically consider the popular Einstein–Gauss–Bonnet theory, $`\alpha =\beta =0`$. Also, implementing the condition (28) (with the + sign) we have,
$`Q^\mu =\sqrt{1+8\kappa f(\gamma ,0,0)\mathrm{\Lambda }_0}{\displaystyle \frac{1}{\kappa }}{\displaystyle d^{D1}x\sqrt{\overline{g}}\overline{\xi _\nu }𝒢_L^{\mu \nu }}.`$ (31)
Although the energy seems to have the wrong sign, this is a red herring: As shown in , for the non-rotating case the exact metric reads
$`ds^2=g_{00}dt^2+g_{rr}dr^2+r^2d\mathrm{\Omega }_{D2}`$ (32)
$`g_{00}=g_{rr}^1=1+{\displaystyle \frac{r^2}{4\kappa \gamma (D3)(D4)}}\left\{1\pm \left\{1+32\gamma \kappa (D3)(D4){\displaystyle \frac{m}{(D2)r^{D1}}}\right\}^{\frac{1}{2}}\right\},`$ (33)
whose asymptotic forms branch into Schwarzschild and Schwarschild-de-Sitter respectively,
$`g_{00}=1{\displaystyle \frac{r_0}{r^{D3}}},g_{00}=1+{\displaystyle \frac{4m}{(D2)}}r^{D3}+{\displaystyle \frac{r^2}{\gamma (D3)(D4)}}.`$ (34)
We see that SdS branch comes with the “wrong” sign compared to the usual Schwarzschild one. Therefore, the minus sign in the energy becomes positive once $`𝒢_L^{\mu \nu }`$ is explicitly computed. We conclude that the conserved charges in the Einstein–Gauss–Bonnet theory for such asymptotic solutions would be simply proportional to those of (17, 21) cosmological gravity:
$`E_{\text{GB}}=\sqrt{1+8\kappa f(\gamma ,0,0)\mathrm{\Lambda }_0}E_DJ_i(\text{GB})=\sqrt{1+8\kappa f(\gamma ,0,0)\mathrm{\Lambda }_0}J_i,`$ (35)
It is important to note that if the coefficient $`\sqrt{1+8\kappa f(\gamma ,0,0)\mathrm{\Lambda }_0}`$ does not vanish, then one can simply rescale the Killing charges to get the Einstein charges (17, 21).
## 4 Conclusions
Using the charge definitions via background Killing charges of we have computed the mass and angular momenta of the rotating Kerr-AdS black holes for $`D`$ dimensions for cosmological Einstein gravity. As a test of stability, we checked that the corresponding charge definitions for higher order would lead to the same values for asymptotically similar geometries up to the indicated constant rescaling.
## 5 Acknowledgments
The work of S.D. is supported by NSF grant PHY 04-01667; that of B.T. by the “Young Investigator Fellowship” of Turkish Academy of Sciences (TUBA) and by a TUBITAK Kariyer Grant. B.T. thanks Özgür Sarıoğlu for a fruitful discussion on TMG. |
warning/0506/hep-ph0506169.html | ar5iv | text | # 1 Abstracts of Lectures
## 1 Abstracts of Lectures
### 1.1 Lecture I
In this lecture I shall begin by tracing the physical origins of chiral symmetry. By asking for the conditions for the conservation of the Axial Vector Current one arrives at the Goldberger Triemann Condition and the need for a triplet of massless scalars to be identified with Pions. The concept of Partially Conserved Axial Current or PCAC is shown to emerge naturally for the real world pions which are not massless. I shall work out the full group structure of chiral symmetry. The Linear Sigma Model is shown to be the simplest realisation.
### 1.2 Lecture II
In this lecture I shall show how one needs Spontaneous breakdown of chiral symmetry to make the Sigma-model to be compatible with nucleon mass. The phenomenon of spontaneous breaking of continuous symmetries will be elaborated with the example of $`O(4)`$ group. In this lecture I shall also discuss the notion of Nonlinear Realisation of chiral symmetry.
### 1.3 Lecture III
In this lecture I shall discuss the dramatic effects brought about by fluctuations in the Goldstone Boson fields. I shall show that in the large-N limit of $`O(N)`$-models in four dimensions the fluctuations completely destroy the spontaneously broken phase.
If time permits, I shall briefly discuss the implications for nuclear forces and I shall conclude my talks with a Renormalisation Group Overview and with a proposal for how to do Chiral Perturbation Theory better.
## 2 Lecture I: Why Chiral Symmetry?
Let us begin by first reviewing Noether’s Theorem which is central to any discussion of Symmetries and Conservation Laws. Consider the action functional
$$S=d^dx(\psi ,\psi ,..)$$
(1)
where $`\psi `$ stands collectively for all the relevant fields. Let the action be invariant under the global transformations
$$\delta \psi =i\mathrm{\Lambda }\psi $$
(2)
Now let us consider the variation of the action under eqn(2) but now for a space-time dependent $`\mathrm{\Lambda }(x)`$. The action is not invariant now but the variation of the action must depend only on the derivatives of $`\mathrm{\Lambda }(x)`$. For a generic $`\mathrm{\Lambda }(x)`$ this variation will be proportional to $`_\mu \mathrm{\Lambda }(x)`$:
$$\delta S=_\mu \mathrm{\Lambda }(x)j^\mu (x)$$
(3)
Now we consider the variational principle determining the equations of motion of the field theory. Here one considers variations $`\stackrel{~}{\delta }`$ of fields that vanish on some initial and final configurations that leaves the extremises the action. In these situations
$$\stackrel{~}{\delta }S=\frac{\stackrel{~}{\delta }S}{\stackrel{~}{\delta }\psi }\stackrel{~}{\delta }\psi $$
(4)
We are looking for a specific field configuration, usually called $`\psi _{cl}`$ for which $`\stackrel{~}{\delta }S=0`$:
$$\frac{\stackrel{~}{\delta }S}{\stackrel{~}{\delta }\psi }|_{\psi =\psi _{cl}}=0$$
(5)
Returning to the context of Noethers theorem the variation of the action eqn(3) can be rewritten as
$$\delta S=\mathrm{\Lambda }(x)_\mu j^\mu $$
(6)
where we have tacitly assumed that fields off sufficiently fast at ’infinity’. This is equivalent to demanding vanishing fields on a boundary and we can use eqn(4) along with eqn(2) for space-time dependent $`\mathrm{\Lambda }(x)`$ to get
$$\delta S=\frac{\delta S}{\delta \psi }\mathrm{\Lambda }(x)\psi $$
(7)
finally arriving at the Noether Theorem
$$_\mu j^\mu =\frac{\delta S}{\delta \psi }\psi $$
(8)
In particular, if the fields $`\psi `$ satisfy the equations of motion, the currents $`j_\mu (x)`$ are conserved:
$$_\mu j^\mu =0$$
(9)
Let us apply this to the case of Dirac action of the form
$$S_{gen}=d^dx\overline{\psi }(x)(i\gamma ^\mu _\mu M)\psi (x)$$
(10)
Let us consider the transformation law
$$\delta \psi =i\lambda \psi $$
(11)
when $`\lambda `$ is a real constant the action is obviously invariant. When $`\lambda `$ is position dependent one gets
$$\delta S=\overline{\psi }(x)\gamma ^\mu \psi (x)_\mu \lambda (x)$$
(12)
We identify the conserved current to be
$$j^\mu (x)=\overline{\psi }(x)\gamma ^\mu \psi (x)$$
(13)
Exercise
Show that this current is conserved only when $`\psi `$ satisfies the equation of motion.
Now let us turn to the case where $`\psi `$ includes protons $`\psi _p`$ and neutrons $`\psi _n`$. The electromagnetic current of the protons
$$j_{em}^\mu =\overline{\psi }_p(x)\gamma ^\mu \psi _p(x)$$
(14)
is conserved
$`i_\mu j_{em}^\mu `$ $`=`$ $`\overline{\psi }_p(x)\gamma ^\mu _\mu \psi _p(x)+(_\mu \overline{\psi }_p(x))\gamma ^\mu \psi _p`$ (15)
$`=`$ $`\overline{\psi }_pM_p\psi _p+(M_p)\overline{\psi }_p\psi _p=0`$
In contrast the ’current’ $`\overline{\psi }_p\gamma ^\mu \psi _n`$ is not conserved but is instead equal to
$$i_\mu (\overline{\psi }_p\gamma ^\mu \psi _n)=(M_pM_n)\overline{\psi }_p\psi _n$$
(16)
However, this current is also conserved if $`M_p=M_n`$ which obviously is a situation of greater symmetry. In fact the vector current $`\stackrel{}{J}^\mu =\overline{\psi }\gamma ^\mu \stackrel{}{\tau }\psi `$(now $`\psi `$ is an isospinor-spinor) is conserved then
$$_\mu \stackrel{}{J}^\mu =0$$
(17)
This enhanced symmetry is called Isospin invariance.
### 2.1 Axial Vector Currents
Now we can ask whether the axial vector current $`\stackrel{}{J}_5^{(p)\mu }=\overline{\psi }_p\gamma ^\mu \gamma _5\psi _p`$ can likewise be made to be conserved with special choice of parameters. It is easily checked that
$$i_\mu \overline{\psi }_i\gamma ^\mu \gamma _5\psi _j=(M_i+M_j)\overline{\psi }_i\gamma _5\psi _j$$
(18)
Since the mass matrix $``$ whose eigenvalues are $`M_i`$(i=1,2 corresponds to protons and neutrons) is positive, it follows that no choice of $``$ will lead to a conserved axial current, unless all the masses are zero! This latter possibility is not as unphysical as it may sound!!
One can ask whether at all axial vector current could ever be conserved. Considering the symmetry with which axial and vector currents occur in the Weak Interaction Hamiltonian it would indeed be very unusual if there was a such a disparity between them.
Let us probe this further by concentrating on the matrix elements of the axial vector current between one proton and one neutron states. The direct contribution corresponding to the fig(1) is given by
$$p|A_\mu ^{(+)}|n_{dir}=G_A\overline{u}(k_p)\gamma _\mu \gamma _5u(k_n)$$
(19)
Here the superscript $`(+)`$ on the axial current refers to the component that transfers one unit of charge i.e it connects ingoing neutron states to outgoing proton states. As discussed above this part of the axial vector is not conserved.
Suppose there existed a pseudoscalar positively charged particle $`\varphi ^{(+)}`$ which interacts with nucleons and whose axial current has a contribution of the form
$$A_\mu ^{(+)}=f_\varphi _\mu \varphi ^{(+)}+\mathrm{}$$
(20)
Diagrammatically this can be represented by fig(2).
This gives rise to an additional contribution to the matrix element which is diagrammatically shown in fig(3) given by
$$p|A_\mu ^{(+)}|n_{med}=ig_{NN\pi }\frac{if_\varphi k_\mu }{k^2m_\varphi ^2}\overline{u}(k_p)\gamma _5u(k_n)$$
(21)
Thus the total contribution to the single-nucleon matrix elements of the axial current is given by
$$p|A_\mu ^{(+)}|n_{tot}=G_A\overline{u}(k_p)\gamma ^\mu \gamma _5u(k_n)\frac{f_\varphi g_{NN\pi }k^\mu }{(k^2m_\varphi ^2)}\overline{u}(k_p)\gamma _5u(k_n)$$
(22)
It is then easy to verify that
$$p|_\mu A^{(+)\mu }|n=[2G_AM_Nf_\varphi g_{NN\pi }\frac{k^2}{(k^2m_\varphi ^2)}]\overline{u}(k_p)\gamma _5u(k_n)$$
(23)
This is zero if $`m_\varphi ^2=0`$ and
$$2G_AM_N=f_\varphi g_{NN\varphi }$$
(24)
In nature there is indeed the pion satisfying the properties conjectured here for $`\varphi `$. Thus the principle of a conserved axial vector current could have actually predicted the existence of pions!
The new symmetry associated with the conservation of axial vector currents is called Chiral Symmetry.
In eqn(24), called the Goldberger-Trieman Relation, $`f_\pi `$ can be determined from pion decay and all the other quantities are also determined from direct observations. That the observed quantities satisfy this relation to very good accuracy is a strong evidence in favour of the ideas expounded here.
But in reality $`m_\pi ^20`$ though it is very small compared to $`M_N^2`$. Keeping this in mind, one has
$$p|_\mu A^{(+)\mu }|n=f_\pi g_{NN\pi }m_\pi ^2\frac{1}{k^2m_\pi ^2}\overline{u}(k_p)\gamma _5u(k_n)$$
(25)
In the above-mentioned considerations the coupling between pions and nucleons has been described by
$$_{NN\pi }=g_{NN\pi }\overline{\psi }_p\gamma _5\pi ^{}\psi _n$$
(26)
Completing the argument for all the three components of the axial vector current means invoking a triplet of pions forming an iso-vector and the full pion-nucleon coupling takes the form
$$_{NN\pi }=g_{NN\pi }\overline{\psi }_p\stackrel{}{\tau }\stackrel{}{\pi }\gamma _5\psi _n$$
(27)
The pions can at this stage be described by
$$_\pi =\frac{1}{2}[(_\mu \stackrel{}{\pi })^2m_\pi ^2\pi ^2]$$
(28)
The generalisation of eqn(25) to the triplet of pions is
$$N,p|_\mu \stackrel{}{A}^\mu |N,q=f_\pi g_{NN\pi }m_\pi ^2\frac{1}{k^2m_\pi ^2}\overline{U}(p)\gamma _5\stackrel{}{\tau }U(q)$$
(29)
The Klein-Gordon equation for pions resulting from eqns(28,27) is
$$(_\mu ^\mu m_\pi ^2)\stackrel{}{\pi }=g_{NN\pi }\overline{N}\gamma _5\stackrel{}{\tau }N$$
(30)
One sees that atleast as far as the single nucleon matrix elements are concerned eqn(29) can be reproduced by
$$_\mu \stackrel{}{A}^\mu =f_\pi m_\pi ^2\stackrel{}{\pi }$$
(31)
This equation called Partially Conserved Axial Current is valid more generally than what is indicated here. Let us also record the equation of motion for the nucleon fields arising out of
$$_N=\overline{N}(i\gamma ^\mu _\mu M)N$$
(32)
and eqn(27)
$$(i\gamma ^\mu _\mu M)N=g_{NN\pi }\gamma _5\stackrel{}{\tau }\stackrel{}{\pi }N$$
(33)
The candidate expression for the (conserved) axial-vector vector current is
$$\stackrel{}{A}^\mu =\overline{N}\gamma ^\mu \gamma _5\frac{\stackrel{}{\tau }}{2}N+f_\pi ^\mu \stackrel{}{\pi }$$
(34)
The form of the pionic contribution to the axial current is strongly supported from Pion weak decays(in fact $`f_\pi `$ was introduced there)
## 3 Lecture II: More on axial vector currents
If one starts from the pion-nucleon Lagrangean introduced so far
$$_{\pi N}=\frac{1}{2}[(_\mu \stackrel{}{\pi })^2+m_\pi ^2\stackrel{}{\pi }^2]+\overline{N}(i\gamma ^\mu _\mu M)Nig_{NN\pi }\overline{N}\gamma _5\stackrel{}{\tau }\stackrel{}{\pi }N$$
(35)
If one makes the natural generalisation of the electromagnetic gauge transformation to
$$\delta _V\psi =i\stackrel{}{\mathrm{\Lambda }}\frac{\stackrel{}{\tau }}{2}\psi $$
(36)
one can see that the pion-nucleon lagrangean is invariant if the pions transform according to
$$\delta _V\stackrel{}{\pi }=\stackrel{}{\mathrm{\Lambda }}\times \stackrel{}{\pi }$$
(37)
when $`\mathrm{\Lambda }`$ is constant. If however $`\mathrm{\Lambda }`$ is position dependent one finds the variation of the Lagrangean to be
$$_\mu \stackrel{}{\mathrm{\Lambda }}\{_\mu \stackrel{}{\pi }\times \stackrel{}{\pi }+\overline{N}\gamma ^\mu \frac{\stackrel{}{\tau }}{2}N\}$$
leading to the isovector vector current
$$\stackrel{}{V}_\mu =\overline{N}\gamma _\mu \frac{\stackrel{}{\tau }}{2}N_\mu \stackrel{}{\pi }\times \stackrel{}{\pi }$$
(38)
The equations of motion resulting from the above lagrangean are
$`(i\gamma ^\mu _\mu M)N`$ $`=`$ $`ig_{NN\pi }\gamma _5\stackrel{}{\tau }\stackrel{}{\pi }N`$
$`\overline{N}(i\gamma ^\mu _\mu M)`$ $`=`$ $`ig_{NN\pi }\overline{N}\gamma _5\stackrel{}{\tau }\stackrel{}{\pi }`$
$`(_\mu ^\mu m_\pi ^2)\stackrel{}{\pi }`$ $`=`$ $`ig_{NN\pi }\overline{N}\gamma _5\stackrel{}{\tau }N`$ (39)
Using these it is easy to show that
$$^\mu \stackrel{}{V}_\mu =0$$
(40)
Now let us see how far we can mimic this procedure for the axial vector currents. We can introduce the axial transformations
$$\delta _A\psi =i\stackrel{}{\omega }\frac{\stackrel{}{\tau }}{2}\gamma _5\psi ;\delta _A\stackrel{}{\pi }=\mathrm{?}$$
(41)
It follows that
$$\delta _A\overline{\psi }=\overline{\psi }\gamma _5i\stackrel{}{\omega }\frac{\stackrel{}{\tau }}{2}$$
(42)
The variation of the lagrangean omitting the purely pionic parts is
$$\delta _A^\stackrel{}{\pi }=iM\overline{N}\gamma _5\stackrel{}{\tau }\stackrel{}{\omega }Nig_{NN\pi }\overline{N}\gamma _5\stackrel{}{\tau }\delta _A\stackrel{}{\pi }N+g_{NN\pi }\overline{N}\stackrel{}{\pi }\stackrel{}{\omega }N\overline{N}\gamma ^\mu \gamma _5\frac{\stackrel{}{\tau }}{2}_\mu \stackrel{}{\omega }$$
(43)
One sees that by virtue of the Goldberger-trieman relation, the first two terms would cancel each other if we choose
$$\delta _A\stackrel{}{\pi }=f_\pi \stackrel{}{\omega }$$
(44)
Under this the variation of the pionic lagrangean is
$$\delta _A_\stackrel{}{\pi }=f_\pi m_\pi ^2\stackrel{}{\omega }\stackrel{}{\pi }f_\pi _\mu \stackrel{}{\pi }^\mu \stackrel{}{\omega }$$
(45)
Thus the Noether prescription would point to an axial current of the form
$$\stackrel{}{A}_\mu =\overline{N}\gamma ^\mu \gamma _5\frac{\stackrel{}{\tau }}{2}N+f_\pi _\mu \stackrel{}{\pi }$$
(46)
which is not exactly conserved because of terms in the variation still proportional to $`\stackrel{}{\omega }`$. While the term proportional to $`m_\pi ^2`$ is what we would have anticipated on the basis of PCAC there is an additional term in $`\delta _A`$ equal to $`g_{NN\pi }\overline{N}\stackrel{}{\pi }\stackrel{}{\omega }N`$ which is clearly undesirable.
To get a clue as to what is happening note
$$\delta _V(i\overline{N}\gamma _5\stackrel{}{\tau }N)=\stackrel{}{\mathrm{\Lambda }}\times (i\overline{N}\gamma _5\stackrel{}{\tau }N)$$
(47)
Thus under vector transformations $`(i\overline{N}\gamma _5\stackrel{}{\tau }N)`$ transforms exactly as $`\stackrel{}{\pi }`$. Let us see how this transforms under the axial transformations:
$$\delta _A(i\overline{N}\gamma _5\stackrel{}{\tau }N)=\overline{N}N\stackrel{}{\omega }$$
(48)
But this does not look like $`\delta _A\stackrel{}{\pi }=f_\pi \stackrel{}{\omega }`$ that we proposed earlier!
But note that
$$\delta _A\overline{N}N=\stackrel{}{\omega }(i\overline{N}\gamma _5\stackrel{}{\tau }N)$$
(49)
This suggests postulating a new field $`\sigma `$ that, along with $`\stackrel{}{\pi }`$,forms an irreducible representation of the Chiral Symmetry Group! The corresponding transformation laws are
$`\delta _V\stackrel{}{\pi }`$ $`=`$ $`\stackrel{}{\mathrm{\Lambda }}\times \stackrel{}{\pi };\delta _A\stackrel{}{\pi }=\sigma \stackrel{}{\omega }`$
$`\delta _V\sigma `$ $`=`$ $`0;\delta _A\sigma =\stackrel{}{\omega }\stackrel{}{\pi }`$
$`\delta _VN`$ $`=`$ $`i\stackrel{}{\mathrm{\Lambda }}{\displaystyle \frac{\stackrel{}{\tau }}{2}}N;\delta _AN=i\stackrel{}{\omega }{\displaystyle \frac{\stackrel{}{\tau }}{2}}\gamma _5N`$ (50)
### 3.1 Group Properties
The way to compute the composition law for various transformations is to look at
$$\delta _{[1,2]}=[\delta _1,\delta _2]$$
(51)
Introducing the total variation of any field $`\psi `$ as
$$\delta \psi =\delta _V\psi +\delta _A\psi $$
(52)
and the generators of vector transformations $`\stackrel{}{T}`$ and of axial transformations $`\stackrel{}{X}`$ we can write
$$\delta \psi =[\stackrel{}{\mathrm{\Lambda }}\stackrel{}{T}+\stackrel{}{\omega }\stackrel{}{X},\psi ]$$
(53)
Then it is easily shown that
$`\delta _{12}\sigma `$ $`=`$ $`\stackrel{}{\omega }_{12}\stackrel{}{\pi }=\stackrel{}{\pi }[\stackrel{}{\omega }_2\times \stackrel{}{\mathrm{\Lambda }}_1\stackrel{}{\omega }_1\times \stackrel{}{\mathrm{\Lambda }}_2]`$
$`\delta _{12}\stackrel{}{\pi }`$ $`=`$ $`(\stackrel{}{\omega }_1\times \stackrel{}{\omega }_2)\times \stackrel{}{\pi }(\stackrel{}{\mathrm{\Lambda }}_1\times \stackrel{}{\mathrm{\Lambda }}_2)\times \stackrel{}{\pi }(\stackrel{}{\mathrm{\Lambda }}_1\times \stackrel{}{\omega }_2\stackrel{}{\mathrm{\Lambda }}_2\times \stackrel{}{\omega }_1)\sigma `$ (54)
Leading to the Lie algebra
$$[T_i,T_j]=iϵ_{ijk}T_k;[T_i,X_j]=iϵ_{ijk}X_k;[X_i,X_j]=iϵ_{ijk}T_k$$
(55)
### 3.2 Invariants
It is easy to verify that the following are invariants:
$`I_1`$ $`=`$ $`\sigma ^2+\stackrel{}{\pi }^2`$
$`I_2`$ $`=`$ $`(_\mu \sigma )^2+(_\mu \stackrel{}{\pi })^2`$
$`I_3`$ $`=`$ $`\overline{N}i\gamma ^\mu _\mu N`$
$`I_4`$ $`=`$ $`\overline{N}N\sigma +i\overline{N}\gamma _5\stackrel{}{\tau }N\stackrel{}{\pi }`$ (56)
With these invariants we can build the following invariant lagrangean which is bound to yield conserved axial and vector currents.
$$_{inv}=\overline{N}i\gamma ^\mu _\mu Ng_{NN\pi }\overline{N}(\sigma +i\gamma _5\stackrel{}{\tau }\stackrel{}{\pi })N+\frac{1}{2}[(_\mu \sigma )^2+(_\mu \stackrel{}{\pi })^2]\frac{\mu ^2}{2}(\sigma ^2+\stackrel{}{\pi }^2)+\frac{\lambda }{4!}(\sigma ^2+\stackrel{}{\pi }^2)^2$$
(57)
Let us illustrate Noether procedure by explicitly working out the conserved axial and vector currents coming only from the $`(\sigma ,\stackrel{}{\pi })`$ part of the invariant lagrangean. Remember we should work out the variation of the lagrangean under position dependent $`(\stackrel{}{\mathrm{\Lambda }}(x),\stackrel{}{\omega }(x)`$:
$`\delta _A\{{\displaystyle \frac{1}{2}}[(_\mu \sigma )^2+(_\mu \stackrel{}{\pi })^2]\}`$ $`=`$ $`_\mu \stackrel{}{\omega }(\sigma _\mu \stackrel{}{\pi }\stackrel{}{\pi }_\mu \sigma )`$
$`\delta _V\{{\displaystyle \frac{1}{2}}[(_\mu \sigma )^2+(_\mu \stackrel{}{\pi })^2]\}`$ $`=`$ $`_\mu \stackrel{}{\mathrm{\Lambda }}(\stackrel{}{\pi }\times _\mu \stackrel{}{\pi })`$ (58)
giving
$`\stackrel{}{A}_\mu `$ $`=`$ $`\sigma _\mu \stackrel{}{\pi }\stackrel{}{\pi }_\mu \sigma `$
$`\stackrel{}{V}_\mu `$ $`=`$ $`\stackrel{}{\pi }\times _\mu \stackrel{}{\pi }`$ (59)
Adding the nucleon contribution, the full currents are
$`\stackrel{}{A}_\mu `$ $`=`$ $`\overline{N}\gamma _\mu \gamma _5{\displaystyle \frac{\stackrel{}{\tau }}{2}}N+\sigma _\mu \stackrel{}{\pi }\stackrel{}{\pi }_\mu \sigma `$
$`\stackrel{}{V}_\mu `$ $`=`$ $`\overline{N}\gamma _\mu {\displaystyle \frac{\stackrel{}{\tau }}{2}}N+\stackrel{}{\pi }\times _\mu \stackrel{}{\pi }`$ (60)
The full set of equations of motion are
$`_\mu ^\mu \stackrel{}{\pi }+\mu ^2\stackrel{}{\pi }+{\displaystyle \frac{\lambda }{6}}(\sigma ^2+\stackrel{}{\pi }^2)\stackrel{}{\pi }+ig_{NN\pi }\overline{N}\gamma _5\stackrel{}{\tau }N`$ $`=`$ $`0`$
$`_\mu ^\mu \sigma +\mu ^2\sigma +{\displaystyle \frac{\lambda }{6}}(\sigma ^2+\stackrel{}{\pi }^2)\sigma +g_{NN\pi }\overline{N}N`$ $`=`$ $`0`$
$`(i\gamma ^\mu _\mu ig_{NN\pi }\gamma _5\stackrel{}{\tau }\stackrel{}{\pi }g_{NN\pi }\sigma )N`$ $`=`$ $`0`$ (61)
It is easy to verify that the currents in eqn(3.2) are exactly conserved when eqn(3.2) are used.
### 3.3 Spontaneous breaking of chiral symmetry
While the fully invariant lagrangean eqn(57) gives an axial current that is exactly conserved we note that eqn(57) has the following severe shortcomings:
i) There is no nucleon mass term in the lagrangean. As it stands it describes massless nucleons. However, there is an additional term $`g_{NN\pi }\overline{N}N`$.
ii) The masses of the pions and sigma are equal. Phenomenologically no isoscalar with such a property exists.
iii)The axial current is conserved even with nonvanishing pion mass!
iv) The axial current does not have a $`f_\pi _\mu \stackrel{}{\pi }`$ term. However, there is a $`\sigma _\mu \stackrel{}{\pi }`$ term.
v) The pion transformation in eqn(3) is $`\delta _A\stackrel{}{\pi }=\sigma \stackrel{}{\omega }`$ while what we had sought in eqn(44) was $`\delta _A\stackrel{}{\pi }=f_\pi \stackrel{}{\omega }`$.
It is remarkable that the difficulties i),iv) and v) can be simultaneously resolved if we let $`\sigma =f_\pi `$! On the other hand that would break chiral symmetry explicitly. What is really needed to fix the problem is something like $`\sigma =f_\pi +somefields`$.
Let us consider the equations of motion only in the $`(\sigma ,\stackrel{}{\pi })`$ sector:
$`_\mu ^\mu \stackrel{}{\pi }+\mu ^2\stackrel{}{\pi }+{\displaystyle \frac{\lambda }{6}}(\sigma ^2+\stackrel{}{\pi }^2)\stackrel{}{\pi }`$ $`=`$ $`0`$
$`_\mu ^\mu \sigma +\mu ^2\sigma +{\displaystyle \frac{\lambda }{6}}(\sigma ^2+\stackrel{}{\pi }^2)\sigma `$ $`=`$ $`0`$ (62)
These are covariant under chiral transformations in the sense that the chiral variation of one of the equations is proportional to the other. When $`\mu ^2>0`$ $`\mu `$ can be interpreted as the mass of the chiral multiplet. Let us now now look at translationally invariant or vacuum solutions to these equations. They satisfy
$$(\sigma _c,\stackrel{}{\pi }_c)(\mu ^2+\frac{\lambda }{6}(\sigma _c^2+\stackrel{}{\pi }_c^2))=0$$
(63)
Vacuum stability requires $`\lambda >0`$. Now if $`\mu ^2>0`$ the only way to satisfy eqn(63) is
$$\sigma _c=0\stackrel{}{\pi }_c=0$$
(64)
This is a chirally symmetric solution in the sense that the chiral variations of this solution vanish on this background.
On the other hand if $`\mu ^2<0`$ there are non-trivial solutions to eqn(63) Symmetric Solution
$$\sigma _c=0\stackrel{}{\pi }_c=0$$
(65)
Asymmetric Solution
$$\sigma _c^2+\stackrel{}{\pi }_c^2=\frac{6|\mu ^2|}{\lambda }$$
(66)
This condition is in itself an invariant but the physical vacuum picks one of the infinitely many solutions to eqn(66). We can use chiral transformations to bring any solution to the form $`\sigma _c=\sqrt{\frac{6|\mu ^2|}{\lambda }}=f_\pi ,\stackrel{}{\pi }_c=0`$.
The potential corresponding to the invariant lagrangean in eqn(57) can be rearranged as
$$V(\sigma ,\stackrel{}{\pi })=\frac{\lambda }{4!}(\sigma ^2+\stackrel{}{\pi }^2f_\pi ^2)^2\frac{\lambda f_\pi ^4}{4!}$$
(67)
Now it is clearly seen that the symmetric solution corresponds to $`V_{sym}=0`$ while the asymmetric solution corresponds to $`V_{asym}=\frac{\lambda f_\pi ^4}{4!}`$. Hence the asymmetric solution lies lower and in this case, the true vacuum.
To get the field content of this choice write
$$\sigma =f_\pi +\stackrel{~}{\sigma }\stackrel{}{\pi }=\stackrel{}{\stackrel{~}{\pi }}$$
(68)
The physical fields defined by vanishing vacuum values are now $`(\stackrel{~}{\sigma },\stackrel{~}{\stackrel{}{\pi }})`$ fields.
Let us see how this scenario resolves all the five difficulties mentioned earlier.
i)Nucleon Mass The interaction term now becomes $`g_{NN\pi }f_\pi \overline{N}N+g_{NN\pi }\overline{N}\stackrel{~}{\sigma }N`$. Not only is a nucleon mass term generated, it is so done as to satisfy the Goldberger-Trieman relation!
iv) The $`\sigma _\mu \stackrel{}{\pi }`$ term in the axial current now becomes $`f_\pi _\mu \stackrel{}{\pi }+\mathrm{}`$ thus resolving this problem.
v) It is easy to see that $`\delta _V`$ rules in terms of the physical fieldshas the same form as before. The axial transformations now become
$$\delta _A\stackrel{~}{\sigma }=\delta _A\sigma =\stackrel{}{\omega }\stackrel{}{\pi };\delta _A\stackrel{}{\pi }=f_\pi \stackrel{}{\omega }\stackrel{~}{\sigma }\stackrel{}{\omega }$$
(69)
Thus we see that the difficulty with the pion transformation law is also resolved.
Now it will be shown that, rather remarkably, the difficulties ii) and iii) also get resolved. To see this let us recast eqn(3.3) in terms of the physical fields
$`_\mu ^\mu \stackrel{}{\pi }+\mu ^2\stackrel{}{\pi }+{\displaystyle \frac{\lambda }{6}}(f_\pi ^2+2f_\pi \stackrel{~}{\sigma }+\stackrel{~}{\sigma }^2+\stackrel{}{\pi }^2)\stackrel{}{\pi }`$ $`=`$ $`0`$
$`_\mu ^\mu \stackrel{~}{\sigma }+\mu ^2f_\pi +\mu ^2\stackrel{~}{\sigma }+{\displaystyle \frac{\lambda }{6}}(f_\pi ^2+2f_\pi \stackrel{~}{\sigma }+\stackrel{~}{\sigma }^2+\stackrel{}{\pi }^2)(f_\pi +\stackrel{~}{\sigma })`$ $`=`$ $`0`$ (70)
These equations are covariant wrt the new chiral transformations. On using the relationship between $`\mu ,\lambda `$ and $`f_\pi `$ we can recast these as
$`_\mu ^\mu \stackrel{}{\pi }+{\displaystyle \frac{\lambda f_\pi }{3}}\stackrel{~}{\sigma }\stackrel{}{\pi }+{\displaystyle \frac{\lambda }{6}}(\stackrel{~}{\sigma }^2+\stackrel{}{\pi }^2)`$ $`=`$ $`0`$
$`_\mu ^\mu \stackrel{~}{\sigma }+{\displaystyle \frac{\lambda f_\pi ^2}{3}}\stackrel{~}{\sigma }+\mathrm{}`$ $`=`$ $`0`$ (71)
The pion mass term has disappeared and the sigma mass term has turned from negative mass squared to positive mass squared as it should be for a physical particle! Thus all the difficulties have been resolved. But now a new physical particle with a mass $`\sqrt{\frac{\lambda f_\pi ^2}{3}}`$ is predicted. Does such a particle exist or not?
While the equations of motion are covariant under the new chiral transformation laws, the vacuum solution is not invariant. This is called Spontaneous breakdown of a symmetry. General theorems exist to show that whenever a continuous symmetry is broken spontaneously, massless particlesappear in the spectrum of the theory. This is called the Goldstone Theorem and the resultant massless particles are called Goldstone particles. Thus pions are the Goldstone bosons of spontaneously broken chiral symmetry. That they are pseudoscalars is an additional consistency requirement.
One of the consequences of spontaneous symmetry breaking of chiral symmetry is that additional interactions are produced. For example the $`\frac{\lambda }{4!}(\sigma ^2+\stackrel{}{\pi }^2)^2`$ term becomes
$`_{int}^{ssb}`$ $`=`$ $`{\displaystyle \frac{\lambda }{4!}}(f_\pi ^2+2f_\pi \stackrel{~}{\sigma }+\stackrel{~}{\sigma }^2+\stackrel{}{\pi }^2)^2`$ (72)
$`=`$ $`massterms+{\displaystyle \frac{\lambda f_\pi }{3!}}\stackrel{~}{\sigma }(\stackrel{~}{\sigma }^2+\stackrel{}{\pi }^2)+{\displaystyle \frac{\lambda }{4!}}(\stackrel{~}{\sigma }^2+\stackrel{}{\pi }^2)^2`$
The additional interaction
$$\frac{\lambda f_\pi }{3!}\stackrel{~}{\sigma }(\stackrel{~}{\sigma }^2+\stackrel{}{\pi }^2)$$
has many important consequences including the phenomenon of chiral cancellations discussed later. Another important consequence is that the sigma particle is unstable with a large width. This is believed to be one of the reasons for the difficulty in detecting this particle.
### 3.4 Finite pion mass
In real life the pion is not massless but has a small mass and is therefore referred to as a pseudo-Goldstone boson. This is not an apt name. So how does one incorporate this into the picture so far developed? One just adds an explicit symmetry breaking term $`\frac{m_\pi ^2}{2}\stackrel{}{\pi }^2`$ to the lagrangean. Actually we would like to add a term that has definite transformation properties. So we could add $`f_\pi m_\pi ^2\sigma `$ instead. This leads to
$$_\mu \stackrel{}{A}^\mu =f_\pi m_\pi ^2\stackrel{}{\pi }$$
(73)
which is nothing but the earlier stated PCAC. In the light of modern developments like QCD, the symmetry breaking is taken to be $`m_q\overline{\psi }_q\psi _q`$ where $`q`$ refers to the quarks. Consistency between the two pictures then requires
$$f_\pi ^2m_\pi ^2=m_q<\overline{\psi }_q\psi _q>$$
(74)
which is called the Gellmann-Renner-Oakes relation.
## 4 Lecture III: Non Linear Realisation of <br>Chiral Symmetry
Let us recapitualate the linear realisation of Chiral Symmetry in terms of the $`(\sigma ,\stackrel{}{\pi })`$ fields:
$`\delta _V\sigma =0`$ $`\delta _V\stackrel{}{\pi }=\stackrel{}{\mathrm{\Lambda }}\times \stackrel{}{\pi }`$
$`\delta _A\sigma =\stackrel{}{\omega }\stackrel{}{\pi }`$ $`\delta _A\stackrel{}{\pi }=\sigma \stackrel{}{\omega }`$ (75)
As we saw already $`\sigma ^2+\stackrel{}{\pi }^2`$ is an invariant. So it is possible to obtain a reduced representation by imposing the invariant condition
$$\sigma ^2+\stackrel{}{\pi }^2=f_\pi ^2$$
(76)
The $`\sigma `$ field is no longer an independent degree of freedom and hence it can be completely eliminated in terms of the $`\stackrel{}{\pi }`$ fields i.e $`\sigma =\sqrt{f_\pi ^2\pi ^2}`$. This induces a non-linear chiral transformation of the pion fields
$$\delta _A\stackrel{}{\pi }=\sqrt{f_\pi ^2\pi ^2}\stackrel{}{\omega }$$
(77)
In terms of the generators $`X_i`$ eqn(77) can be recast as
$$[X_i,\pi _j]=i\sqrt{f_\pi ^2\pi ^2}\delta _{ij}$$
(78)
The isospin transformations are unchanged and remain linear. Weinberg showed that the most general non-linear transformations are of the type
$$[X_i,\pi _j]=i\delta _{ij}f(\pi ^2)+\pi _i\pi _jg(\pi ^2)$$
(79)
The functions $`f,g`$ are not independent but constrained by Jacobi Identity
$$[X_k,[X_i,\pi _j]]+[\pi _j,[X_k,X_i]]+[X_i,[\pi _j,X_k]]=0$$
(80)
On using
$$[X_i,X_j]=iϵ_{ijk}T_k;[T_i,\pi _j]=iϵ_{ijk}\pi _k$$
(81)
it is easy to show that
$$fg2ff^{}2f^{}g\pi ^2=1$$
(82)
It is easy to see that $`f`$ can not be chosen arbitrarily. For example $`f`$ can never vanish. On the other hand $`g`$ can be chosen arbitrarily and given $`g`$ we can determine $`f`$ accordingly. If we take $`g=0`$ then $`f=\sqrt{c^2\pi ^2}`$. This corresponds to eqn(78) with $`c=f_\pi `$.
The invariant lagrangean
$$_{\pi ,\sigma }=\frac{1}{2}[(_\mu \sigma )^2+(_\mu \stackrel{}{\pi })^2+m^2(\sigma ^2+\stackrel{}{\pi }^2)]$$
(83)
becomes, upon elimination of $`\sigma `$
$$_\pi =\frac{1}{2}[(_\mu \stackrel{}{\pi })^2+\frac{\stackrel{}{\pi }_\mu \stackrel{}{\pi })^2}{1\pi ^2}]$$
(84)
where we have set $`f_\pi =1`$. This can, following Weinberg, be written in the form
$$_\pi =\frac{1}{2}𝒟_\mu \stackrel{}{\pi })^2$$
(85)
where
$$𝒟_\mu \stackrel{}{\pi }=_\mu \stackrel{}{\pi }+\frac{1}{2}\frac{1}{\sigma +\sigma ^2}\stackrel{}{\pi }\stackrel{}{\pi }_\mu \stackrel{}{\pi }$$
(86)
is called the chiral covariant derivative and is like the covariant derivatives of general covariant theories.
Going back to the axial current and the chiral invariant pion-nucleon Lagrangean, these become
$`\stackrel{}{A}_\mu `$ $`=`$ $`\sigma _\mu \stackrel{}{\pi }\stackrel{}{\pi }_\mu \sigma `$ (87)
$``$ $`_\mu \stackrel{}{\pi }+(\stackrel{}{\pi }\stackrel{}{\pi }_\mu \stackrel{}{\pi }\pi ^2_\mu \stackrel{}{\pi })+\mathrm{}.`$
and
$`_{\pi N}`$ $`=`$ $`\overline{N}i\gamma ^\mu _\mu N+{\displaystyle \frac{1}{2}}(𝒟_\mu \stackrel{}{\pi })^2g_{NN\pi }\overline{N}(\sigma +i\gamma _5\stackrel{}{\tau }\stackrel{}{\pi })N`$ (88)
$`=`$ $`\overline{N}(i\gamma ^\mu _\mu M)N+{\displaystyle \frac{1}{2}}(𝒟_\mu \stackrel{}{\pi })^2g_{NN\pi }\overline{N}({\displaystyle \frac{\pi ^2}{2}}+i\gamma _5\stackrel{}{\tau }\stackrel{}{\pi })N`$
Another way to understand the non-linear realisation is to imagine studying the pion-nucleon system at low energies. In the spontaneously broken phase $`M_\sigma 57M_\pi `$. Thus when all the momenta are small, the $`\sigma `$ can considered heavy and integrated out leaving behind an effective chiral invariant description of pions and nucleons.
This brings us to the very important notion of the infrared behaviour of Goldstone Bosons. First let us consider the low energy scattering of pions. The linear theory predicts a constant amplitude which is diagrammatically represented in fig(4):
$$A(\pi \pi \pi \pi )i\lambda $$
(89)
But the spontaneously broken phase of this theory introduces additional processes like $`\sigma \pi \pi `$ leading to additional contributions shown in fig(5) which actually cancel this contribution to yield a scattering amplitude that vanishes as $`k^2`$!.
Same happens in the nonlinear model also. From eqn() the effective interaction is given by
$$_{int}=\frac{1}{2}(\stackrel{}{\pi }_\mu \stackrel{}{\pi })^2$$
(90)
This means that the pion-pion scattering amplitude, shown in fig(6), vanishes $`k^2`$!
This feature, known as chiral cancellations is not just restricted to pion-pion scattering. In fact it is generic to all physical processes involving Goldstone Bosons. We illustrate this now for pion-nucleon scattering. In the linear theory this approaches a constant in the infrared limit. The relevant amplitude is shown in fig(7).
However in the spontaneously broken phase additional contributions as shown in fig(8) exactly cancel this contribution!
The same cancellation takes place in the non-linear model also.
### 4.1 No spontaneous symmetry breaking in the d=4 large N model
The model we considered can also be thought of as $`O(4)`$ as $`(\sigma ,\stackrel{}{\pi })`$ can be arranged into a four-component vector multiplet. In this section we show that spontaneous symmetry breaking does not happen in d=4 when we consider $`O(N)`$ models in the large-N limit.
Let us start with the definition of the generating functional of connected Green’s functions $`W[\stackrel{}{H}]`$ of a field theory given by
$$e^{iW[\stackrel{}{H}]}=𝒟\stackrel{}{\mathrm{\Phi }}e^{i(S[\stackrel{}{\mathrm{\Phi }}]+\stackrel{}{H}\stackrel{}{\mathrm{\Phi }})}$$
(91)
where
$$\stackrel{}{H}\stackrel{}{\mathrm{\Phi }}=d^4x\stackrel{}{H}(x)\stackrel{}{\mathrm{\Phi }}(x)$$
(92)
The generating functional of 1PI vertices, $`\gamma [\stackrel{}{\mathrm{\Phi }}]`$, is obtained from $`W[\stackrel{}{H}]`$ by a Legendre transformation
$$\mathrm{\Gamma }[\stackrel{}{\mathrm{\Phi }}]=W[\stackrel{}{H}]\stackrel{}{H}\stackrel{}{\mathrm{\Phi }}$$
(93)
where
$$\mathrm{\Phi }_i(x)=\frac{\delta W[\stackrel{}{H}]}{\delta H_i(x)}$$
(94)
For the linear $`O(N)`$-model we have
$$e^{iW[\stackrel{}{H}]}=𝒟\stackrel{}{\mathrm{\Phi }}e^{i{\scriptscriptstyle d^4x[{\scriptscriptstyle \frac{1}{2}}(_\mu \stackrel{}{\mathrm{\Phi }}(x))^2{\scriptscriptstyle \frac{U}{4}}(\stackrel{}{\mathrm{\Phi }}^2C^2)^2+\stackrel{}{H}(x)\stackrel{}{\mathrm{\Phi }}(x)]}}$$
(95)
It helps to introduce an auxilliary field $`\chi (x)`$ to make the exponent in eqn(95) bilinear in $`\stackrel{}{\mathrm{\Phi }}`$:
$$e^{iW[\stackrel{}{H}]}=𝒟\chi 𝒟\stackrel{}{\mathrm{\Phi }}e^{i{\scriptscriptstyle d^4x[{\scriptscriptstyle \frac{1}{2}}(_\mu \stackrel{}{\mathrm{\Phi }}(x))^2{\scriptscriptstyle \frac{\chi (x)}{2}}(\stackrel{}{\mathrm{\Phi }}^2C^2)+{\scriptscriptstyle \frac{1}{4U}}\chi (x)^2+\stackrel{}{H}(x)\stackrel{}{\mathrm{\Phi }}(x)]}}$$
(96)
As $`\stackrel{}{\mathrm{\Phi }}`$ only appears quadratically it can easily be integrated out to get
$$e^{iW[\stackrel{}{H}]}=𝒟\chi e^{i{\scriptscriptstyle d^4x[{\scriptscriptstyle \frac{\chi ^2(x)}{4U}}+{\scriptscriptstyle \frac{C^2}{2}}\chi (x)]}\frac{N}{2}Trln\frac{i}{2}\stackrel{}{H}^1\stackrel{}{H}}$$
(97)
where
$$(x,y)=(^2\chi (x))\delta (xy)$$
(98)
On performing the rescalings
$$NU=u,C^2=Nc^2,H_i=\sqrt{N}h_i,W[\stackrel{}{H}]=Nw[\stackrel{}{h}]$$
(99)
we get
$$e^{iNw[\stackrel{}{h}]}=𝒟\chi e^{iN[{\scriptscriptstyle d^4x({\scriptscriptstyle \frac{\chi ^2}{4u}}+{\scriptscriptstyle \frac{c^2}{2}}\chi )}+\frac{i}{2}Trln\frac{1}{2}\stackrel{}{h}^1\stackrel{}{h}]}$$
(100)
Following the standard large-N techniques, the leading order result is
$$w_{\mathrm{}}[\stackrel{}{h}]=d^4x(\frac{\chi ^2}{4u}+\frac{c^2}{2}\chi )+\frac{i}{2}Trln\frac{1}{2}\stackrel{}{h}^1\stackrel{}{h}$$
(101)
The classical field $`\varphi _i`$ where $`\stackrel{}{\mathrm{\Phi }}=\sqrt{N}\stackrel{}{\varphi }`$ is given by
$$\varphi _i(x)=^1h_i(x)$$
(102)
Finally $`\mathrm{\Gamma }_{\mathrm{}}[\stackrel{}{\varphi },\chi ]`$ is given by
$$\mathrm{\Gamma }_{\mathrm{}}[\stackrel{}{\varphi },\chi ]=d^4x[\frac{\chi ^2(x)}{4u}+\frac{c^2}{2}\chi (x)+\frac{1}{2}(_\mu \stackrel{}{\varphi }(x))^2\frac{1}{2}\chi (x)\stackrel{}{\varphi }(x)^2]+\frac{i}{2}Trln$$
(103)
The resulting effective potential is
$$V_{eff}[\stackrel{}{\varphi },\chi ]=\frac{\chi ^2}{4u}\frac{c^2}{2}\chi +\frac{1}{2}\chi \stackrel{}{\varphi }^2+\frac{\chi ^2}{64\pi ^2}(ln\frac{\chi }{M^2}\frac{3}{2})$$
(104)
The vacua are determined by the stationarity conditions
$$\frac{V_{eff}}{\chi }=0\frac{V_{eff}}{\varphi _a}=0$$
(105)
That is
$`\chi \varphi _a`$ $`=`$ $`0`$
$`{\displaystyle \frac{\chi }{2u}}+{\displaystyle \frac{1}{2}}\varphi ^2{\displaystyle \frac{c^2}{2}}+{\displaystyle \frac{1}{32\pi ^2}}\chi (ln{\displaystyle \frac{\chi }{M^2}}1)`$ $`=`$ $`0`$ (106)
The spontaneously broken state requires $`\chi =0`$ and hence $`V_{eff}(asy)=0`$. From the second of eqn(4.1) one sees that this is possible only if $`c^2>0`$.
On the other hand, the symmetric vacuum condition $`\varphi ^2=0`$ and eqn(4.1) first of all says that $`\chi _c0`$ and in fact gives
$$\frac{\chi _c}{2u}\frac{c^2}{2}+\frac{1}{32\pi ^2}\chi _c(ln\frac{\chi _c}{M^2}1)=0$$
(107)
Using this the effective potential at the symmetric solution can be rewritten as
$`V_{eff}(sym)`$ $`=`$ $`{\displaystyle \frac{\chi _c}{2}}({\displaystyle \frac{\chi _c}{2u}}{\displaystyle \frac{c^2}{2}}+{\displaystyle \frac{chi_c}{32\pi ^2}}(ln{\displaystyle \frac{\chi _c}{M^2}}1)){\displaystyle \frac{c^2}{4}}\chi _c{\displaystyle \frac{\chi _c^2}{128\pi ^2}}`$ (108)
$`=`$ $`{\displaystyle \frac{c^2}{4}}\chi _c{\displaystyle \frac{\chi _c^2}{128\pi ^2}}`$
now eqn(107) says that $`\chi _c>0`$ and combining this with the positivity of $`c^2`$ one sees that $`V_{eff}(sym)<0`$ i.e the symmetric solution lies below the spontaneously broken solution! |
warning/0506/hep-th0506200.html | ar5iv | text | # Non-perturbative effects in the BMN limit of 𝒩=4 supersymmetric Yang–Mills
## 1 Introduction
According to there is a very interesting limit of the AdS/CFT correspondence that relates a special sector of the $`𝒩`$=4 supersymmetric Yang–Mills (SYM) theory to type IIB string theory in a maximally supersymmetric plane-wave background. A notable feature of this proposal is that it provides the first example of a gauge/gravity duality which can be studied in a quantitative way beyond the supergravity approximation. This is possible because string theory in the relevant background, which is obtained as a Penrose limit of AdS$`{}_{5}{}^{}\times S^5`$ , can be quantised in the light cone gauge . Moreover there exists a regime in which both the string and the gauge theory are weakly coupled, although there are subtleties associated with the different way in which the limits leading to this regime are taken on the two sides. This has allowed very precise comparisons between perturbative corrections in the two theories .
The duality relates the string mass spectrum to the spectrum of scaling dimensions of gauge theory operators in the so called BMN sector of $`𝒩`$=4 SYM. This consists of gauge invariant operators of large conformal dimension, $`\mathrm{\Delta }`$, and large charge, $`J`$, with respect to a U(1) subgroup of the SU(4) R-symmetry group. The duality involves the double limit $`\mathrm{\Delta }\mathrm{}`$, $`J\mathrm{}`$. The combination $`\mathrm{\Delta }J`$, which is kept finite, is related to the string theory hamiltonian,
$$\mathrm{\Delta }J=\frac{1}{\mu }H^{(2)},$$
(1.1)
where $`\mu `$ is the background value of the $`\mathrm{R}\mathrm{R}`$ five-form and is related to the mass parameter, $`m`$, which appears in the light cone string action by $`m=\mu p_{}\alpha ^{}`$ , where $`p_{}`$ is the light cone momentum.
The correspondence between the spectra of the two theories is thus the statement that the eigenvalues of the operators on the two sides of the equality (1.1) coincide. A quantitative comparison is possible if one considers the large $`N`$ limit in the gauge theory focusing on operators in the BMN sector. As a result of combining the large $`N`$ limit with the limit of large $`\mathrm{\Delta }`$ and $`J`$, new effective parameters arise , which are related to the ordinary ’t Hooft parameters, $`\lambda `$ and $`1/N`$, by a rescaling,
$$\lambda ^{}=\frac{g_{_{\mathrm{YM}}}^2N}{J^2},g_2=\frac{J^2}{N}.$$
(1.2)
The correspondence relates these effective gauge theory couplings to string theory parameters in the plane-wave background,
$$m^2=(\mu p_{}\alpha ^{})^2=\frac{1}{\lambda ^{}},4\pi g_\mathrm{s}m^2=g_2.$$
(1.3)
The double scaling limit, $`N\mathrm{}`$, $`J\mathrm{}`$ with $`J^2/N`$ fixed, connects the weak coupling regime of the gauge theory to string theory at small $`g_s`$ and large $`m`$.
In this limit the leading perturbative corrections to the scaling dimensions of BMN operators have been successfully compared to the leading quantum corrections to the masses of the dual plane-wave string states , see also the reviews for further references. In the present paper we will study one-instanton effects in the $`𝒩`$=4 Yang–Mills theory. These will be compared with $`D`$-instanton induced corrections to the plane-wave string mass spectrum that were computed in in order to check the validity of the BMN conjecture in non-perturbative sectors. In the original formulation of the AdS/CFT correspondence very good agreement was found between the effects of instantons in the $`𝒩`$=4 Yang–Mills theory and of $`D`$-instantons in type IIB string theory in AdS$`{}_{5}{}^{}\times S^5`$ . It is therefore of interest to see if similar agreement can be established in the BMN limit and whether the results of can be reproduced from the study of instanton contributions to the anomalous dimensions of BMN operators.
The possibility of testing the correspondence at the non-perturbative level is especially relevant since several aspects of the perturbative tests of the duality are only partially understood. A precise holographic formulation of the duality connecting the dynamics of the two theories beyond the identification of the spectra is still lacking and even the explicit tests at the level of the spectrum are not comprehensive. A limited class of states/operators has been studied and agreement has been explicitly verified only at leading order in $`g_2`$ (the planar limit). This is the limit in which the string is free and the $`\lambda ^{}`$ perturbation series on the gauge side reproduces the free string spectrum. The first non-planar contributions, of order $`\lambda ^{}g_2^2`$ or one-loop on the string side, have also been compared successfully (although in this case simplifying assumptions were made about the one-loop string theory calculation, which have since been questioned ). The systematics of the perturbative expansion beyond these leading order contributions has not been studied and the fact that the double scaling parameters (1.2) that arise at low orders indeed represent the correct expansion parameters at all orders remains a conjecture.
Results obtained in different but related limits of the AdS/CFT duality, both in string theory and on the gauge theory side , suggest the possibility that BMN scaling, i.e. the order by order reorganisation of the perturbative expansion into a double series in $`\lambda ^{}`$ and $`g_2`$, might break down at higher orders. In the strict BMN sector the scaling (1.2) has been verified to three loops in perturbation theory , but a deviation was observed in a related matrix model calculation at four loops.
In this paper we will show that instanton contributions to the conformal dimensions of BMN operators display BMN scaling. Two-point functions of BMN operators computed in the semi-classical approximation will be shown to be in striking agreement with the $`D`$-instanton induced two-point amplitudes computed in . The agreement includes not only the dependence on the parameters $`\lambda ^{}`$ and $`g_2`$, but also the dependence on the mode numbers characterising the states. In particular the agreement with in the mode number dependence is highly non-trivial and requires dramatic cancellations. These results combined with the three loop perturbative result provide substantial evidence indicating that BMN scaling should persist at all orders.
Instanton contributions to the anomalous dimensions of BMN operators are extracted from two-point correlation functions computed in the semi-classical approximation. Conformal invariance determines the form of two-point functions of primary operators, $`𝒪`$ and $`\overline{𝒪}`$, to be
$$𝒪(x_1)\overline{𝒪}(x_2)=\frac{c}{(x_1x_2)^{2\mathrm{\Delta }}},$$
(1.4)
where $`\mathrm{\Delta }`$ is the scaling dimension. In general in the quantum theory $`\mathrm{\Delta }`$ acquires an anomalous term, $`\mathrm{\Delta }(g_{_{\mathrm{YM}}})=\mathrm{\Delta }_0+\gamma (g_{_{\mathrm{YM}}})`$. At weak coupling the anomalous dimension $`\gamma (g_{_{\mathrm{YM}}})`$ is small and substituting in (1.4) gives
$$𝒪(x_1)\overline{𝒪}(x_2)=\frac{c\mathrm{\Lambda }^{2\gamma (g_{_{\mathrm{YM}}})}}{(x_1x_2)^{2\mathrm{\Delta }_0}}\left(1\gamma (g_{_{\mathrm{YM}}})\mathrm{log}\left[\mathrm{\Lambda }^2(x_1x_2)^2\right]+\mathrm{}\right),$$
(1.5)
where $`\mathrm{\Lambda }`$ is an arbitrary renormalisation scale. As a function of the coupling constant the anomalous dimension admits an expansion consisting of a perturbative series plus non-perturbative corrections. The generic two-point function at weak coupling takes the form
$`𝒪(x_1)\overline{𝒪}(x_2)`$ $`=`$ $`{\displaystyle \frac{c(g_{_{\mathrm{YM}}})}{(x_1x_2)^{2\mathrm{\Delta }_0}}}(1g_{_{\mathrm{YM}}}^2\gamma ^{(1)}\mathrm{log}\left[\mathrm{\Lambda }^2(x_1x_2)^2\right]`$ (1.6)
$`+\mathrm{}\mathrm{e}^{2\pi i\tau }\gamma ^{(\mathrm{inst})}\mathrm{log}\left[\mathrm{\Lambda }^2(x_1x_2)^2\right]+\mathrm{}).`$
Therefore perturbative and instanton contributions to the anomalous dimension are extracted from the coefficients of the logarithmically divergent terms in a two-point function. When there is more than one operator with the same quantum numbers operator mixing occurs. In this case the resulting set of two-point functions determines a matrix of anomalous dimensions and the eigenvalues of this matrix are the physical anomalous dimensions. The issue of operator mixing was first discussed in the context of the BMN limit in .
The procedure for calculating the instanton-induced contribution to the anomalous dimensions in semi-classical approximation is as follows. The gauge-invariant operators in the BMN sector are defined by colour traces involving a large number of elementary scalar fields together with a finite number of bosonic or fermionic ‘impurities’. In the semi-classical approximation correlation functions of such operators are computed by replacing the fields by the solution to the corresponding field equations in the presence of an instanton, expressed in terms of the fermionic and bosonic moduli, and integrating the resulting profiles over these moduli. These moduli encode the broken superconformal symmetries together with the (super)symmetries associated with the orientation of a SU($`2`$) instanton within SU($`N`$). For large $`N`$ integration over these moduli is carried out by a saddle point procedure (as in ).
The general structure of the anomalous dimensions of gauge invariant operators in the $`𝒩`$=4 Yang–Mills theory with SU($`N`$) gauge group is an expansion of the form
$$\gamma (g_{_{\mathrm{YM}}},\theta ,N)=\underset{n=1}{\overset{\mathrm{}}{}}\gamma _n^{\mathrm{pert}}(N)g_{_{\mathrm{YM}}}^{2n}+\underset{K>0}{}\underset{m=0}{\overset{\mathrm{}}{}}[\gamma _m^{(K)}(N)g_{_{\mathrm{YM}}}^{2m}\mathrm{e}^{2\pi i\tau K}+\mathrm{c}.\mathrm{c}.],$$
(1.7)
where $`\tau =\frac{\theta }{2\pi }+i\frac{4\pi }{g_{_{\mathrm{YM}}}^2}`$. The double series in the second term in (1.7) contains the contributions of multi-instanton sectors as well as the perturbative fluctuations in each such sector. One reason for studying instanton effects even though they are exponentially suppressed in the small coupling limit, is that they determine the dependence on the $`\theta `$-angle in $`𝒩`$=4 SYM. They therefore play an essential rôle in implementing $`S`$-duality which is a symmetry of the theory, just as $`D`$-instantons are crucial for the $`S`$-duality in type IIB string theory.
If the BMN sector of the gauge theory scales appropriately (1.7) becomes a series in the scaled couplings $`\lambda ^{}`$ and $`g_2`$. In particular, we will show that the leading one-instanton contribution to the two-point functions of a class of four impurity BMN operators scales as it should in the BMN limit and has the form $`(1/n_1n_2)^2g_2^{7/2}\mathrm{exp}\left(8\pi ^2/g_2\lambda ^{}+i\theta \right)`$, where the integers $`n_1`$ and $`n_2`$ correspond to the mode numbers of the dual string state. This result is in striking agreement with the corresponding $`D`$-instanton induced mass matrix on the string side found in .
For certain other classes of operators the leading one-instanton contribution vanishes and the first non-zero correction is of higher order in $`\lambda ^{}`$ (or a lower power of $`m`$ in the string calculation). In such cases the calculation requires knowledge of a non-leading term in the scalar solution – a term involving six fermionic moduli (whereas the leading term is quadratic in fermionic moduli). We have not evaluated the precise form of this contribution and so have not determined the precise form of the matrix elements in these cases. However, there is strong evidence that these also match the string calculations. For example, for two impurity operators, with some mild assumptions about the manner in which the fermion moduli are distributed in the profile of the operators, we will find a contribution to the two-point function of the form $`\lambda _{}^{}{}_{}{}^{2}g_2^{7/2}\mathrm{exp}\left(8\pi ^2/g_2\lambda ^{}+i\theta \right)`$, in accord with expectations from the string side. Later we will comment on the systematics of the expansion in the one-instanton sector and on how the higher order corrections can give rise to a double series in $`\lambda ^{}`$ and $`g_2`$.
This paper is organised as follows. In section 2 we review some general aspects of the $`𝒩`$=4 Yang–Mills theory and the BMN limit. The general method for evaluating instanton-induced contributions to two-point functions of BMN operators in terms of zero modes and the integration over super-moduli is described in section 3. The manner in which the profiles of the fields depend on these moduli is presented in section 4. In section 5 we consider some specific examples of two-point functions of BMN operators and derive expressions for the anomalous dimensions that arise after integration over all the moduli. We first consider the case of two-impurity operators (which presents the technical difficulty alluded to above) and then four-impurity operators. We conclude with a discussion in section 6, which includes a comparison with the string results in . Some technical details of the calculations are presented in the appendices.
## 2 Fields and operators in the BMN limit
The purpose of this section is mainly to present the notation used in the paper and to define the dictionary to be used for the comparison with string theory in the plane-wave background. We will only consider a small set of BMN operators with scalar impurities which are dual to the string states studied in . A more detailed discussion of the various types of operators relevant for the comparison with string theory in the plane-wave background can be found in the review papers .
### 2.1 Fields in $`𝒩`$=4 SYM
The $`𝒩`$=4 multiplet comprises six real scalars, $`\widehat{\phi }^i`$, $`i=1,\mathrm{},6`$, four Weyl fermions, $`\lambda _\alpha ^A`$, $`A=1,\mathrm{},4`$, and a vector, $`A_\mu `$, with field strength $`F_{\mu \nu }`$, all transforming in the adjoint representation of the gauge group. These are the building blocks used to construct gauge invariant composite operators which are classified according to the irreducible representations of the superconformal group, SU(2,2$`|`$4). The latter are identified by the quantum numbers $`(\mathrm{\Delta },j_1,j_2;a,b,c)`$ of the maximal bosonic subgroup SO(2,4)$`\times `$SO(6), where $`\mathrm{\Delta }`$ is the scaling dimension, $`j_1`$ and $`j_2`$ the Lorentz spins and $`[a,b,c]`$ the SU(4)$``$SO(6) Dynkin labels.
Under the SU(4) R-symmetry group the scalars transform in the $`\mathrm{𝟔}`$, the fermions in the $`\mathrm{𝟒}`$ (and their conjugates in the $`\overline{\mathrm{𝟒}}`$) and the gauge field is a singlet. It is often convenient to label the scalars by an antisymmetric pair of indices in the $`\mathrm{𝟒}`$, $`\phi ^{[AB]}`$, subject to the reality condition
$$\overline{\phi }_{AB}\left(\phi ^{AB}\right)^{}=\frac{1}{2}\epsilon _{ABCD}\phi ^{CD}.$$
(2.1)
The two parametrisations of the $`𝒩`$=4 scalars, $`\widehat{\phi }^i`$ and $`\phi ^{AB}`$, are related by
$$\widehat{\phi }^i=\frac{1}{\sqrt{2}}\mathrm{\Sigma }_{AB}^i\phi ^{AB},\phi ^{AB}=\frac{1}{\sqrt{8}}\overline{\mathrm{\Sigma }}_i^{AB}\widehat{\phi }^i,$$
(2.2)
where $`\mathrm{\Sigma }_{AB}^i`$ ($`\overline{\mathrm{\Sigma }}_i^{AB}`$) are Clebsch–Gordan coefficients projecting the product of two $`\mathrm{𝟒}`$’s ($`\overline{\mathrm{𝟒}}`$’s) onto the $`\mathrm{𝟔}`$. They are defined in appendix A. The representation of the scalars in terms of the $`\phi ^{AB}`$ fields is the most convenient for instanton calculations since, as we shall see, it makes manifest which fermion zero modes in a correlation function can be soaked up by each scalar field.
In the limit relevant for the comparison with string theory in the plane-wave background the symmetry group is a contraction of the original group and the operators are classified according to representations of the bosonic subgroup SO(4)$`\times `$SO(4)$`\times `$U(1)$`\times `$U(1). We shall denote by $`𝒟`$ the dilation operator and by $`𝒥`$ the U(1) generator selected by the Penrose limit in the dual AdS background. The SO(4)$`\times `$SO(4)$`\times `$U(1)$`\times `$U(1) quantum numbers are $`(s_1,s_2;s_1^{},s_2^{};\mathrm{\Delta },J)`$, where $`\mathrm{\Delta }`$ and $`J`$ refer to $`𝒟`$ and $`𝒥`$ and the spins $`s_i`$ and $`s_i^{}`$ refer to the two SO(4) factors. These can be considered to be respectively subgroups of the original SO(6) and SO(2,4) groups. This identification is not completely correct. The generators of the two SO(4)’s corresponding to the isometries of the dual string background, $`\stackrel{~}{G}_i`$, are related to the generators of the Euclidean Lorentz group and to those of an SO(4) subgroup of the R-symmetry group, $`G_i`$, $`i=1,2`$, by a similarity transformation, $`\stackrel{~}{G}_i=TG_iT^1`$. This distinction, however, will not be relevant for our analysis.
Since a precise formulation of the gauge theory dual to the plane wave string theory is not known, the rules for the decomposition of the $`𝒩`$=4 fields according to representations of SO(4)$`\times `$SO(4)$`\times `$U(1)$`\times `$U(1) are determined by the quantum numbers of the dual string excitations. The gauge invariant operators corresponding to states in the string spectrum will be discussed in the next subsection. String excitations created by bosonic and fermionic oscillators are associated respectively with the insertion of bosonic and fermionic elementary fields (“impurities”) in composite operators.
Bosonic excitations in the plane wave string theory originate from the vector of SO(8) which decomposes under SO(4)$`\times `$SO(4) as
$$\mathrm{𝟖}_\mathrm{v}=[(\frac{1}{2},\frac{1}{2});(0,0)][(0,0);(\frac{1}{2},\frac{1}{2})],$$
(2.3)
i.e. they are vectors of one SO(4) and singlets of the second or vice versa. Correspondingly in the $`𝒩`$=4 theory the six real scalars are reorganised into a complex field, $`Z`$, and its conjugate, $`\overline{Z}`$, which are singlets of SO(4)$`\times `$SO(4) and have $`\mathrm{\Delta }=1`$ and $`J=\pm 1`$ respectively, and four real fields which transform in the $`\mathrm{𝟒}=(\frac{1}{2},\frac{1}{2})`$ of the first SO(4) and are singlets with respect to the second and have $`J=0`$ and $`\mathrm{\Delta }=1`$. The insertion of the four real scalars in a composite operator corresponds to the insertion of bosonic creation operators with an index in one of the two SO(4) factors in the dual string state. States created by bosonic oscillators which are vectors of the second SO(4) correspond to operators involving insertions of $`D_\mu Z`$. The fields $`D_\mu Z`$ are in the $`\mathrm{𝟒}=(\frac{1}{2},\frac{1}{2})`$ of the second SO(4) and have $`J=1`$ and $`\mathrm{\Delta }=2`$. Explicitly the scalar fields are
$$\begin{array}{cc}Z=\varphi ^1=2\phi ^{14},\hfill & \overline{Z}=\varphi _1^{}=2\phi ^{23},\text{}\hfill \\ \phi ^1=\widehat{\phi }^2=\frac{1}{\sqrt{2}}\left(\phi ^{13}+\phi ^{24}\right),\hfill & \phi ^2=\widehat{\phi }^3=\frac{1}{\sqrt{2}}\left(\phi ^{12}+\phi ^{34}\right),\text{}\hfill \\ \phi ^3=\widehat{\phi }^5=\frac{i}{\sqrt{2}}\left(\phi ^{13}\phi ^{24}\right),\hfill & \phi ^4=\widehat{\phi }^6=\frac{i}{\sqrt{2}}\left(\phi ^{12}\phi ^{34}\right),\text{}\hfill \end{array}$$
(2.4)
Here, for convenience of notation, we have introduced the scalars $`\phi ^i`$, $`i=1,\mathrm{},4`$, related to four of the $`\widehat{\phi }^i`$’s by a relabelling. This should not cause any confusion since in the following we shall only work with (2.4) and we shall not need the SO(6) fields (A.5).
Unlike the scalar fields the fermions transform non trivially with respect to both SO(4)’s. The four Weyl fermions of the $`𝒩`$=4 SYM theory, $`\lambda _\alpha ^A`$, transform in the $`\mathrm{𝟒}`$ of SU(4)$``$SO(6) and their conjugates, $`\overline{\lambda }_A^{\dot{\alpha }}`$, transform in the $`\overline{\mathrm{𝟒}}`$. Their decomposition with respect to SO(4)$`\times `$SO(4)$`\times `$U(1) is dictated by the SO(4)$`\times `$SO(4) decomposition of the SO(8) fermions of the light-cone string. The type IIB fermions transform in the $`\mathrm{𝟖}_\mathrm{s}`$ of SO(8), which under SO(4)$`\times `$SO(4) decomposes as
$$\mathrm{𝟖}_\mathrm{s}=[(0,\frac{1}{2});(0,\frac{1}{2})][(\frac{1}{2},0);(\frac{1}{2},0)].$$
(2.5)
In terms of the $`\mathrm{𝟖}_\mathrm{s}`$ fermions $`S^a`$ and $`\stackrel{~}{S}^a`$ the decomposition is achieved via a projector , $`\frac{1}{2}(1\pm \mathrm{\Pi })`$,
$`S^a(S^{})_\alpha ^a(S^+)_{\dot{a}}^{\dot{\alpha }}`$
$`\stackrel{~}{S}^a(\stackrel{~}{S}^{})_\alpha ^a(\stackrel{~}{S}^+)_{\dot{a}}^{\dot{\alpha }}.`$ (2.6)
The Yang–Mills fermions, $`\lambda _\alpha ^A`$, have $`\mathrm{\Delta }=\frac{3}{2}`$ and U(1) charge $`\frac{1}{2}`$ for $`A=1,4`$ and $`\frac{1}{2}`$ for $`A=2,3`$. Similarly their conjugates, $`\overline{\lambda }_A^{\dot{\alpha }}`$, have $`\mathrm{\Delta }=\frac{3}{2}`$ and U(1) charge $`\frac{1}{2}`$ for $`A=2,3`$ and $`\frac{1}{2}`$ for $`A=1,4`$. To match the quantum numbers of the string oscillators we choose the following decomposition
$$\lambda _\alpha ^A\psi _{+\frac{1}{2};\alpha }^a\overline{\psi }_{\frac{1}{2};\alpha a}^+,a=1,4,$$
(2.7)
where the fermions $`\overline{\psi }_{\alpha a}^+`$ are defined as
$$\overline{\psi }_{\frac{1}{2};\alpha a}^+=\left(M^+\lambda \right)_{\frac{1}{2};\alpha a},$$
(2.8)
where the matrix $`M^+`$ is related to the matrix $`\mathrm{\Pi }`$ used in the plane-wave string theory to project the SO(8) fermions onto chiral SO(4)$`\times `$SO(4) spinors. The spinors $`\psi _{+\frac{1}{2};\alpha }^a`$ and $`\overline{\psi }_{\frac{1}{2};\alpha a}^+`$ transform under SO(4)$`\times `$SO(4) in the $`(\mathrm{𝟐}_L;\mathrm{𝟐}_L)=[(\frac{1}{2},0);(\frac{1}{2},0)]`$ and have $`(J=\frac{1}{2},\mathrm{\Delta }=\frac{3}{2})`$ and $`(J=\frac{1}{2},\mathrm{\Delta }=\frac{3}{2})`$ respectively.
Similarly
$$\overline{\lambda }_A^{\dot{\alpha }}\psi _{+\frac{1}{2};\dot{a}}^{+\dot{\alpha }}\overline{\psi }_{\frac{1}{2}}^{\dot{\alpha }\dot{a}},\dot{a}=2,3$$
(2.9)
where
$$\overline{\psi }_{\frac{1}{2}}^{\dot{\alpha }\dot{a}}=\left(M^{}\overline{\lambda }\right)_{\frac{1}{2}}^{\dot{\alpha }\dot{a}}$$
(2.10)
and $`M^{}`$ is also related to the projector used to define the SO(4)$`\times `$SO(4) spinors in the dual plane-wave string theory. The fermions $`\psi _{+\frac{1}{2};\dot{a}}^{+\dot{\alpha }}`$ and $`\overline{\psi }_{\frac{1}{2}}^{\dot{\alpha }\dot{a}}`$ transform in the $`(\mathrm{𝟐}_R;\mathrm{𝟐}_R)=[(0,\frac{1}{2});(0,\frac{1}{2})]`$ representation and have respectively $`(J=\frac{1}{2},\mathrm{\Delta }=\frac{3}{2})`$ and $`(J=\frac{1}{2},\mathrm{\Delta }=\frac{3}{2})`$. Some aspects of instanton contributions to BMN operators involving fermionic impurities will be discussed in .
The field strength, $`F_{\mu \nu }`$, is a singlet with respect to the first SO(4) and decomposes into $`F_{\mu \nu }^\pm `$ transforming in the $`\mathrm{𝟑}^{}=(1,0)`$ and $`\mathrm{𝟑}^+=(0,1)`$ with respect to the second. $`F_{\mu \nu }^\pm `$ both have $`J=0`$ and $`\mathrm{\Delta }=2`$.
In the string amplitudes in the plane-wave background $`P_+`$ and $`P_{}`$ are conserved, so the operators in the gauge theory are conveniently classified according to the dual quantities, i.e. $`\mathrm{\Delta }J`$ and $`\mathrm{\Delta }+J`$ respectively ($`\mathrm{\Delta }+J`$ is actually infinite in the limit; it is proportional to $`P_{}`$, but the proportionality constant diverges). Table 1 summarises the $`\mathrm{\Delta }`$, $`J`$ and SO(4)$`\times `$SO(4) quantum numbers for the $`𝒩`$=4 elementary fields. The notation SO(4)<sub>R</sub> and SO(4)<sub>C</sub> has been introduced to denote the SO(4) groups descending from the original SO(6) (R-symmetry) and SO(2,4) (conformal) groups respectively.
### 2.2 BMN operators
The composite operators dual to states in the spectrum of string theory in the plane wave background are also classified in terms of the same quantum numbers. In particular, $`\mathrm{\Delta }J`$, which is dual to the light-cone hamiltonian, measures the number of “impurities” and will be used to classify the operators
At finite $`J`$ and $`\mathrm{\Delta }`$ the selection rules of the $`𝒩`$=4 theory, implied by the superconformal symmetry, apply. So only two-point functions of (primary) operators of the same dimension can be non-zero. More precisely the two-point functions that are relevant for the calculation of anomalous dimensions are
$$\overline{𝒪}_{\overline{𝐫}_i,\mathrm{\Delta }_i}^i(x)𝒪_{𝐫_j,\mathrm{\Delta }_j}^j(y),$$
(2.11)
where the subscripts denote the SU(4) representation and the scaling dimension. The SU(2,2$`|`$4) symmetry imposes $`\mathrm{\Delta }_i=\mathrm{\Delta }_j`$ and $`𝐫_i=𝐫_j`$ so that both $`\mathrm{\Delta }`$ and $`J`$ are conserved in two-point functions. In the case of the U(1) charge $`J`$ this means that the two operators in a non-zero two-point function must have equal and opposite charges.
Gauge invariant composite operators which correspond to physical string states are of the form
$`𝒪_{J;n_1\mathrm{}n_k}^{i_1\mathrm{}i_k}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{J^{\mathrm{\Delta }J+1}\left(\frac{g_{_{\mathrm{YM}}}^2N}{8\pi ^2}\right)^{J+k}}}}\begin{array}{c}{\displaystyle \underset{p_1,\mathrm{},p_k=0}{\overset{J}{}}}\\ p_1p_2\mathrm{}p_k\end{array}\mathrm{e}^{2\pi i(p_1n_1+p_2n_2+\mathrm{}+p_kn_k)/J}`$
$`\times \mathrm{Tr}\left(Z^{p_1}X_\mathrm{}_1Z^{p_2p_1}X_\mathrm{}_2\mathrm{}Z^{p_kp_{k1}}X_\mathrm{}_kZ^{Jp_k}\right)`$
$`=`$ $`{\displaystyle \frac{1}{\sqrt{J^{\mathrm{\Delta }J1}\left(\frac{g_{_{\mathrm{YM}}}^2N}{8\pi ^2}\right)^{J+k}}}}\begin{array}{c}{\displaystyle \underset{q_2,\mathrm{},q_k=0}{\overset{J}{}}}\\ q_2+\mathrm{}+q_kJ\end{array}\mathrm{e}^{2\pi i[(n_2+\mathrm{}+n_k)q_2+(n_3+\mathrm{}+n_k)q_3+\mathrm{}+n_kq_k]/J}`$ (2.18)
$`\times \mathrm{Tr}\left(Z^{J(q_2+\mathrm{}+q_k)}X_\mathrm{}_1Z^{q_2}X_\mathrm{}_2\mathrm{}Z^{q_k}X_\mathrm{}_k\right),`$
where $`\mathrm{\Delta }J`$ denotes the total number of impurities. Here the cyclicity of the trace has been used and, after the change of variables, $`p_1q_1`$, $`p_iq_iq_{i1}`$ ($`i=2,\mathrm{},k`$), the sum over $`q_1`$ has been performed resulting in the condition
$$n_1=(n_2+\mathrm{}+n_k).$$
(2.19)
In (2.18) the $`X_{\mathrm{}}`$’s denote generic impurities, i.e. any of the elementary fields discussed in the previous subsection. String states dual to these operators are created acting on the vacuum with bosonic and fermionic oscillators. The integers $`n_1,\mathrm{},n_k`$ in (2.18) are identified with the mode numbers in the dual string state and the relation (2.19) corresponds to the level matching condition obeyed by the physical string states.
String creation operators are in one to one correspondence with $`\mathrm{\Delta }J=1`$ impurities, see table 1. Bosonic oscillators $`\alpha _n^i`$ and $`\alpha _n^\mu `$ in the $`(\mathrm{𝟒};\mathrm{𝟏})`$ and $`(\mathrm{𝟏};\mathrm{𝟒})`$ of SO(4)$`\times `$SO(4) correspond to the insertion of $`\phi ^i`$ and $`D^\mu Z`$ impurities respectively <sup>1</sup><sup>1</sup>1Here we are using a different notation with respect to , where the oscillators $`\alpha _n^\mu `$ were denoted by $`\alpha _n^i^{}`$.. Fermionic oscillators, $`S_n^+`$ and $`S_n^{}`$, in the $`(\mathrm{𝟐}_L;\mathrm{𝟐}_L)`$ and $`(\mathrm{𝟐}_R;\mathrm{𝟐}_R)`$ correspond to the insertion of $`\psi _{+\frac{1}{2},\alpha }^a`$ and $`\psi _{+\frac{1}{2},\dot{a}}^{+\dot{\alpha }}`$ impurities respectively. In string theory for each type of excitation one must consider left- and right-moving oscillators. These correspond to the insertion of the same field, but with the associated $`n_i`$ in the phase factor in (2.18) being respectively positive or negative.
The normalisation of operators of the form (2.18) involving only $`\mathrm{\Delta }J=1`$ impurities is such that their tree level two-point functions are of order 1 in the BMN limit, $`N\mathrm{}`$, $`J\mathrm{}`$ with $`J^2/N`$ finite. Operators involving $`\mathrm{\Delta }J=2`$ impurities have vanishing two-point functions at tree level because for equal total $`\mathrm{\Delta }J`$ they are normalised by the same prefactor but their definition involves fewer sums. Therefore these operators do not correspond to independent degrees of freedom in the BMN limit. In some cases, however, the insertion of $`\mathrm{\Delta }J=2`$ impurities is necessary in order to construct combinations which are well behaved in the double limit $`N\mathrm{}`$, $`J\mathrm{}`$ at higher orders in perturbation theory. For instance it is necessary to consider terms in which pairs of $`\phi ^i`$ impurities are replaced by a $`\overline{Z}`$ insertion in order to cancel divergences which arise at the level of the leading non planar perturbative corrections .
Operators with $`\mathrm{\Delta }J=0,1`$ are protected and so their two-point functions do not receive instanton contributions. At the level of two and more impurities the situation is more interesting. The spectrum is significantly richer and more importantly there appear unprotected operators. In the following we shall only discuss a small selection of gauge invariant composite operators with scalar impurities which are dual to the string states analysed in . A complete analysis would require computing the two-point functions involving all the operators in each sector and then diagonalising the resulting matrix of anomalous dimensions. We shall not carry out this program in this paper, but we shall concentrate on a few specific cases which illustrate the striking agreement with the corrections to the string mass spectrum calculated in . The generic operator with $`k`$ scalar impurities is of the form
$`𝒪_{J;n_1\mathrm{}n_k}^{i_1\mathrm{}i_k}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{J^{k1}\left(\frac{g_{_{\mathrm{YM}}}^2N}{8\pi ^2}\right)^{J+k}}}}\begin{array}{c}{\displaystyle \underset{p_1,\mathrm{},p_{k1}=0}{\overset{J}{}}}\\ p_1+\mathrm{}+p_{k1}J\end{array}\mathrm{e}^{2\pi i[(n_1+\mathrm{}+n_{k1})p_1+(n_2+\mathrm{}+n_{k1})p_2+\mathrm{}+n_{k1}p_{k1}]/J}`$ (2.23)
$`\times \mathrm{Tr}\left(Z^{J(p_1+\mathrm{}+p_{k1})}\phi ^{i_1}Z^{p_1}\phi ^{i_2}\mathrm{}Z^{p_{k1}}\phi ^{i_k}\right),`$
In our discussion we shall only consider operators with an even number of impurities. This is because operators with odd $`\mathrm{\Delta }J`$, which receive perturbative corrections , are expected not to receive contributions in the one-instanton sector. This is a prediction following from the calculation of string amplitudes in . In the plane-wave string theory the absence of instanton contributions to two-point amplitudes of states with an odd number of non-zero mode excitations is a straightforward consequence of the structure of the $`D`$-instanton boundary state. In the $`𝒩`$=4 theory, however, the corresponding statement is far from obvious.
#### 2.2.1 Two impurity operators
No field in the $`𝒩`$=4 multiplet has negative $`\mathrm{\Delta }J`$ hence the two impurity operators are obtained with the insertion of either two $`\mathrm{\Delta }J=1`$ fields or of a single field with $`\mathrm{\Delta }J=2`$. Because of the normalisation (2.18) only operators with two $`\mathrm{\Delta }J=1`$ insertions are relevant in the BMN limit.
Even restricting the attention to SO(4)<sub>C</sub> singlets, already at the two impurity level there is a rather rich spectrum of operators, which becomes even richer when multi-trace operators with the same quantum numbers are taken into account. In the SO(4)<sub>R</sub> singlet sector one can construct gauge invariant operators in the representations $`(0,0)\mathrm{𝟏}`$, $`(1,0)\mathrm{𝟑}^+`$, $`(0,1)\mathrm{𝟑}^{}`$ and $`(1,1)\mathrm{𝟗}`$ of SO(4)<sub>R</sub>. The singlet can be realised with the insertion of two scalars, two gauge fields (through covariant derivatives) or two fermions of the same chirality. Operators in the $`\mathrm{𝟑}^\pm `$ include combinations of two scalar or two fermionic impurities. The $`\mathrm{𝟗}`$ can only be obtained with the insertion of two scalar impurities.
The operators with two $`\phi ^i`$ insertions are <sup>2</sup><sup>2</sup>2Here and in the following we use square brackets to denote antisymmetrisation, curly brackets to denote symmetrisation and subtraction of the trace and parentheses to indicate symmetrisation without subtraction of the trace part.
$`𝒪_{\mathrm{𝟏};J;n}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{J\left(\frac{g_{_{\mathrm{YM}}}^2N}{8\pi ^2}\right)^{J+2}}}}\left[{\displaystyle \underset{p=0}{\overset{J}{}}}\mathrm{e}^{2\pi ipn/J}\mathrm{Tr}\left(Z^{Jp}\phi ^iZ^p\phi ^i\right)\mathrm{Tr}\left(Z^{J+1}\overline{Z}\right)\right]`$ (2.24)
$`𝒪_{\mathrm{𝟑}^\pm ;J;n}^{[ij]}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{J\left(\frac{g_{_{\mathrm{YM}}}^2N}{8\pi ^2}\right)^{J+2}}}}{\displaystyle \underset{p=0}{\overset{J}{}}}\mathrm{e}^{2\pi ipn/J}\Gamma _\pm ^{ijkl}\mathrm{Tr}\left(Z^{Jp}\phi ^{[k}Z^p\phi ^{l]}\right)`$ (2.25)
$`𝒪_{\mathrm{𝟗};J;n}^{\{ij\}}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{J\left(\frac{g_{_{\mathrm{YM}}}^2N}{8\pi ^2}\right)^{J+2}}}}{\displaystyle \underset{p=0}{\overset{J}{}}}\mathrm{e}^{2\pi ipn/J}\mathrm{Tr}\left(Z^{Jp}\phi ^{\{i}Z^p\phi ^{j\}}\right)`$
$``$ $`{\displaystyle \frac{1}{\sqrt{J\left(\frac{g_{_{\mathrm{YM}}}^2N}{8\pi ^2}\right)^{J+2}}}}{\displaystyle \underset{p=0}{\overset{J}{}}}\mathrm{e}^{2\pi ipn/J}\left[\mathrm{Tr}\left(Z^{Jp}\phi ^{(i}Z^p\phi ^{j)}\right){\displaystyle \frac{\delta ^{ij}}{2}}\mathrm{Tr}\left(Z^{Jp}\phi ^kZ^p\phi ^k\right)\right],`$
where the projectors onto the $`\mathrm{𝟑}^+`$ and $`\mathrm{𝟑}^{}`$ are defined as $`\Gamma _\pm ^{ijkl}=\frac{1}{4}\left(\delta ^{ik}\delta ^{jl}\delta ^{il}\delta ^{jk}\pm \epsilon ^{ijkl}\right)`$. The singlet operator (2.24) provides an example of what mentioned earlier about $`\overline{Z}`$ insertions. In order to define a well behaved BMN operator it is necessary to consider the combination in (2.24). The second term is needed to cancel a divergent contribution to the two-point function of $`𝒪_{\mathrm{𝟏};J;n}`$ arising at the leading non planar level .
The operators (2.24)-(2.2.1) are dual to string states in the plane-wave background of the form
$`\alpha _n^i\stackrel{~}{\alpha }_n^i|0_h`$ (2.27)
$`\Gamma _\pm ^{ijkl}\alpha _n^k\stackrel{~}{\alpha }_n^l|0_h`$ (2.28)
$`\alpha _n^{\{i}\stackrel{~}{\alpha }_n^{j\}}|0_h,`$ (2.29)
where $`|0_h`$ denotes the BMN ground state and the indices $`i,j,\mathrm{}`$ are taken to be in one of the two SO(4) factors (to be identified with SO(4)<sub>R</sub>). The integer $`n`$ in (2.24)-(2.2.1) corresponds to the level of the dual string excitation. The $`n=0`$ operators are protected and correspond to supergravity states.
As already remarked, in order to compute the instanton induced anomalous dimensions of the various operators in each sector it is in principle necessary to diagonalise the appropriate matrix. In the following we shall not consider the problem of mixing between single- and multi-trace operators, since it is a subleading effect in the large $`N`$ limit. In general, however, at the instanton level mixing occurs at leading order among all the single trace operators in each sector . In the case of the two impurity operators it has been shown that all the operators in different sectors have the same anomalous dimension as a consequence of superconformal invariance. Therefore in the following we shall only analyse the single operator in the $`\mathrm{𝟗}`$ for which there is no mixing to resolve. Superconformal symmetry guarantees that the results apply to operators in others sectors as well.
In general the problem of resolving the operator mixing and computing anomalous dimensions can be drastically simplified using the constraints imposed by superconformal invariance, in particular, the fact that all the operators in a multiplet have the same anomalous dimension as the superconformal primary operator.
#### 2.2.2 Four impurity operators
The number of independent BMN operators grows very rapidly with the number of impurities and at the four impurity level the spectrum is already extremely rich. Bosonic operators in the SO(4)<sub>C</sub> singlet sector exist in the following representations of SO(4)<sub>R</sub>
$`(0,0)=\mathrm{𝟏},(0,1)=\mathrm{𝟑}^+,(1,0)=\mathrm{𝟑}^{},(0,2)=\mathrm{𝟓}^+,(2,0)=\mathrm{𝟓}^{},`$
$`(1,1)=\mathrm{𝟗},(1,2)=\mathrm{𝟏𝟓}^+,(2,1)=\mathrm{𝟏𝟓}^{},(2,2)=\mathrm{𝟐𝟓}.`$ (2.30)
Operators relevant in the BMN limit involve four $`\mathrm{\Delta }J=1`$ impurities. The combinations which contribute to SO(4)<sub>C</sub> scalars are listed (up to permutations of the four fields) in table 2 <sup>3</sup><sup>3</sup>3$`t_{\mu _1\mu _2\mu _3\mu _4}`$ is a projector onto the singlet, i.e. $`\delta _{\mu _1\mu _2}\delta _{\mu _3\mu _4}`$, $`\delta _{\mu _1\mu _3}\delta _{\mu _2\mu _4}`$ or $`\epsilon _{\mu _1\mu _2\mu _3\mu _4}`$..
The SO(4)<sub>C</sub> singlet sector contains the largest number of operators, involving all the combinations $`(i)`$-$`(xi)`$ in table 2. Operators in the $`\mathrm{𝟑}^+`$ can contain $`(i),(ii),(iv)`$-$`(vi),(viii)`$ and $`(ix)`$ and those in the $`\mathrm{𝟑}^{}`$ $`(i),(ii),(iv),(vii),(viii)`$ and $`(x)`$. The $`\mathrm{𝟗}`$ involves $`(i),(ii),(iv),(v)`$ and $`(viii)`$. Operators in the $`\mathrm{𝟓}^+`$ and $`\mathrm{𝟓}^{}`$ can be obtained from $`(i),(iv)`$ and $`(ix)`$ and from $`(i),(v)`$ and $`(x)`$ respectively, those in the $`\mathrm{𝟏𝟓}^+`$ and $`\mathrm{𝟏𝟓}^{}`$ from $`(i)`$ and $`(iv)`$ and from $`(i)`$ and $`(v)`$ respectively. In the $`\mathrm{𝟐𝟓}`$ there is only one operator corresponding to the combination $`(i)`$ with indices fully symmetrised.
In the following we shall concentrate on a few specific two-point functions corresponding to the amplitudes computed in . This will be sufficient to show how instanton contributions to gauge theory correlation functions precisely reproduce various features observed in string theory amplitudes. The operators we study in detail are those containing four scalar impurities. These are of the form
$`𝒪_{𝐫;J;n_1,n_2,n_3}`$ $`=`$ $`{\displaystyle \frac{t_{ijkl}^𝐫}{\sqrt{J^3\left(\frac{g_{_{\mathrm{YM}}}^2N}{8\pi ^2}\right)^{J+4}}}}\begin{array}{c}{\displaystyle \underset{p_1,p_2,p_3=0}{\overset{J}{}}}\\ p_1+p_2+p_3J\end{array}\mathrm{e}^{2\pi i[(n_1+n_2+n_3)p_1+(n_2+n_3)p_2+n_3p_3]/J}`$ (2.34)
$`\times \mathrm{Tr}\left(Z^{J(p_1+p_2+p_3)}\phi ^iZ^{p_1}\phi ^jZ^{p_2}\phi ^kZ^{p_3}\phi ^l\right),`$
where $`t_{ijkl}^𝐫`$ is a projector onto the representation $`𝐫`$ of SO(4)<sub>R</sub>. In particular in the singlet sector there are three independent operators in this class, corresponding to the three tensors
$$t_{ijkl}^{(1)}=\epsilon _{ijkl},t_{ijkl}^{(2)}=\delta _{ij}\delta _{kl},t_{ijkl}^{(3)}=\delta _{ik}\delta _{jl}.$$
(2.35)
In section 5.2 we study instanton contributions to two-point functions of operators of the type (2.34). We discuss in detail the case of the singlet corresponding the choice of the $`t_{ijkl}^{(1)}`$ projector. We will show that the dependence on both the coupling constants, $`\lambda ^{}`$ and $`g_2`$, and the mode numbers, $`n_i`$, is in exact agreement with the results of . We also find that for all the operators with four scalar impurities in sectors other than the singlet instanton contributions to the matrix of anomalous dimensions are suppressed by powers of $`\lambda ^{}`$. This result is also in agreement with the string prediction.
## 3 Instanton contributions to two-point functions
In this section we recall some general aspects of the calculation of instanton contributions to correlation functions, in particular two-point functions, in $`𝒩`$=4 SYM.
In the semi-classical limit correlation functions are evaluated using a saddle point approximation around the classical instanton configuration. In this limit the computation of expectation values reduces to an integration over the finite dimensional instanton moduli space as parametrised in the ADHM construction . Before focusing on operators of large dimension, $`\mathrm{\Delta }`$, and R-charge, $`J`$, in the following sections, we briefly discuss the general formalism for extracting instanton contributions to the anomalous dimensions of gauge invariant local operators. Comprehensive reviews of instanton calculus in supersymmetric gauge theories can be found in and instanton contributions to anomalous dimensions of scalar operators in $`𝒩`$=4 SYM were studied in detail in .
Contributions to the (matrix of) anomalous dimensions are extracted from the logarithmically divergent terms in two-point functions. In the semi-classical approximation in the background of an instanton the two-point function of a generic local operator, $`𝒪(x)`$, and its conjugate takes the form
$$\overline{𝒪}(x_1)𝒪(x_2)_{\mathrm{inst}}=d\mu _{\mathrm{inst}}(𝓂_\mathrm{b},𝓂_\mathrm{f})\mathrm{e}^{S_{\mathrm{inst}}}\widehat{\overline{𝒪}}(x_1;𝓂_\mathrm{b},𝓂_\mathrm{f})\widehat{𝒪}(x_2;𝓂_\mathrm{b},𝓂_\mathrm{f}),$$
(3.1)
where we have denoted the bosonic and fermionic collective coordinates by $`𝓂_\mathrm{b}`$ and $`𝓂_\mathrm{f}`$ respectively. In (3.1) $`\mathrm{d}\mu _{\mathrm{inst}}(𝓂_\mathrm{b},𝓂_\mathrm{f})`$ is the integration measure on the instanton moduli space, $`S_{\mathrm{inst}}`$ is the classical action evaluated on the instanton solution and $`\widehat{𝒪}`$ and $`\widehat{\overline{𝒪}}`$ denote the classical expressions for the operators $`𝒪`$ and $`\overline{𝒪}`$ computed in the instanton background.
A one-instanton configuration in SU($`N`$) Yang–Mills theory is characterised by 4$`N`$ bosonic collective coordinates. With a particular choice of parametrisation these bosonic moduli can be identified with the size, $`\rho `$, and position, $`x_0`$, of the instanton as well as its global gauge orientation. The latter can be described by three angles identifying the iso-orientation of a SU(2) instanton and 4$`N`$ additional constrained variables, $`w_{u\dot{\alpha }}`$ and $`\overline{w}^{\dot{\alpha }u}`$ (where $`u=1,\mathrm{},N`$ is a colour index), in the coset SU($`N`$)/(SU($`N2`$)$`\times `$U(1)) describing the embedding of the SU(2) configuration into SU($`N`$). In the one-instanton sector in the $`𝒩`$=4 theory there are additionally 8$`N`$ fermionic collective coordinates corresponding to zero modes of the Dirac operator in the background of an instanton. They comprise the 16 moduli associated with Poincaré and special supersymmetries broken by the instanton and denoted respectively by $`\eta _\alpha ^A`$ and $`\overline{\xi }^{\dot{\alpha }A}`$ (where $`A`$ is an index in the fundamental of the SU(4) R-symmetry group) and 8$`N`$ additional parameters, $`\nu _u^A`$ and $`\overline{\nu }^{Au}`$, which can be considered as the fermionic superpartners of the gauge orientation parameters. The sixteen superconformal moduli are exact, i.e. they enter the expectation values (3.1) only through the classical profiles of the operators. The other fermion modes, $`\nu _u^A`$ and $`\overline{\nu }^{Au}`$, appear explicitly in the integration measure via the classical action, $`S_{\mathrm{inst}}`$. This distinction plays a crucial rôle in the calculation of correlation functions. The $`\nu _u^A`$ and $`\overline{\nu }^{Au}`$ modes satisfy the fermionic ADHM constraints
$$\overline{w}^{\dot{\alpha }u}\nu _u^A=0,\overline{\nu }^{Au}w_{u\dot{\alpha }}=0,$$
(3.2)
which effectively reduce their number to $`8(N2)`$.
In the one-instanton sector the gauge-invariant measure on the instanton moduli space takes the form
$`{\displaystyle d\mu _{\mathrm{phys}}\mathrm{e}^{S_{\mathrm{inst}}}}`$ (3.3)
$`={\displaystyle \frac{\pi ^{4N}g_{_{\mathrm{YM}}}^{4N}\mathrm{e}^{2\pi i\tau }}{(N1)!(N2)!}}{\displaystyle d\rho \mathrm{d}^4x_0\underset{A=1}{\overset{4}{}}\mathrm{d}^2\eta ^A\mathrm{d}^2\overline{\xi }^A\mathrm{d}^{N2}\nu ^A\mathrm{d}^{N2}\overline{\nu }^A\rho ^{4N13}\mathrm{e}^{S_{4F}}},`$
where the instanton action is
$$S_{\mathrm{inst}}=2\pi i\tau +S_{4F}=2\pi i\tau +\frac{\pi ^2}{2g_{_{\mathrm{YM}}}^2\rho ^2}\epsilon _{ABCD}^{AB}^{CD}$$
(3.4)
with
$$\tau =\frac{4\pi i}{g_{_{\mathrm{YM}}}^2}+\frac{\theta }{2\pi },^{AB}=\frac{1}{2\sqrt{2}}(\overline{\nu }^{Au}\nu _u^B\overline{\nu }^{Bu}\nu _u^A).$$
(3.5)
In (3.3) we have omitted an overall ($`N`$-independent) numerical constant that will be reinstated in the final expression.
The two-point function (3.1) thus becomes
$`\overline{𝒪}(x_1)𝒪(x_2)={\displaystyle \frac{\pi ^{4N}g_{_{\mathrm{YM}}}^{4N}\mathrm{e}^{2\pi i\tau }}{(N1)!(N2)!}}{\displaystyle d\rho \mathrm{d}^4x_0\underset{A=1}{\overset{4}{}}\mathrm{d}^2\eta ^A\mathrm{d}^2\overline{\xi }^A\mathrm{d}^{N2}\nu ^A\mathrm{d}^{N2}\overline{\nu }^A\rho ^{4N13}}`$
$`\mathrm{e}^{\frac{\pi ^2}{16g_{_{\mathrm{YM}}}^2\rho ^2}\epsilon _{ABCD}(\overline{\nu }^{[A}\nu ^{B]})(\overline{\nu }^{[C}\nu ^{D]})}\widehat{\overline{𝒪}}(x_1;x_0,\rho ,\eta ,\overline{\xi },\nu ,\overline{\nu })\widehat{𝒪}(x_2;x_0,\rho ,\eta ,\overline{\xi },\nu ,\overline{\nu }).`$ (3.6)
Following the integration over the non-exact fermion modes can be reduced to a gaussian form introducing auxiliary bosonic coordinates, $`\chi ^i`$, $`i=1,\mathrm{},6`$, to rewrite the gauge invariant measure as
$`{\displaystyle \frac{\pi ^{4N}g_{_{\mathrm{YM}}}^{4N}\mathrm{e}^{2\pi i\tau }}{(N1)!(N2)!}}{\displaystyle d\rho \mathrm{d}^4x_0\mathrm{d}^6\chi \underset{A=1}{\overset{4}{}}\mathrm{d}^2\eta ^A\mathrm{d}^2\overline{\xi }^A\mathrm{d}^{N2}\nu ^A\mathrm{d}^{N2}\overline{\nu }^A}`$
$`\rho ^{4N7}\mathrm{exp}\left[2\rho ^2\chi ^i\chi ^i+{\displaystyle \frac{4\pi i}{g_{_{\mathrm{YM}}}}}\chi _{AB}^{AB}\right],`$ (3.7)
where $`\chi _{AB}=\frac{1}{\sqrt{8}}\mathrm{\Sigma }_{AB}^i\chi ^i`$ and the symbols $`\mathrm{\Sigma }_{AB}^i`$ are defined in (A.1).
The fermion modes $`\nu _u^A`$ and $`\overline{\nu }^{Bu}`$ enter explicitly in the classical profiles of the operators in the instanton background as well as in the measure through the instanton action. It is thus convenient to construct a generating function as in , which allows to deal easily with the otherwise complicated combinatorics associated with the integration over $`\nu _u^A`$ and $`\overline{\nu }^{Au}`$. We introduce sources, $`\overline{\vartheta }_A^u`$ and $`\vartheta _{Au}`$, coupled to $`\nu _u^A`$ and $`\overline{\nu }^{Au}`$ and define
$`Z[\vartheta ,\overline{\vartheta }]`$ $`=`$ $`{\displaystyle \frac{\pi ^{4N}g_{_{\mathrm{YM}}}^{4N}\mathrm{e}^{2\pi i\tau }}{(N1)!(N2)!}}{\displaystyle d\rho \mathrm{d}^4x_0\mathrm{d}^6\chi \underset{A=1}{\overset{4}{}}\mathrm{d}^2\eta ^A\mathrm{d}^2\overline{\xi }^A\mathrm{d}^{N2}\overline{\nu }^A\mathrm{d}^{N2}\nu ^A}`$ (3.8)
$`\rho ^{4N7}\mathrm{exp}\left[2\rho ^2\chi ^i\chi ^i+{\displaystyle \frac{\sqrt{8}\pi i}{g_{_{\mathrm{YM}}}}}\overline{\nu }^{Au}\chi _{AB}\nu _u^B+\overline{\vartheta }_A^u\nu _u^A+\vartheta _{Au}\overline{\nu }^{Au}\right].`$
Performing the gaussian integrals over $`\overline{\nu }`$ and $`\nu `$ and introducing polar coordinates,
$$\chi ^i(r,\mathrm{\Omega }),\underset{i=1}{\overset{6}{}}(\chi ^i)^2=r^2,$$
(3.9)
we find
$`Z[\vartheta ,\overline{\vartheta }]`$ $`=`$ $`{\displaystyle \frac{2^{29}\pi ^{13}g_{_{\mathrm{YM}}}^8\mathrm{e}^{2\pi i\tau }}{(N1)!(N2)!}}{\displaystyle d\rho \mathrm{d}^4x_0\mathrm{d}^5\mathrm{\Omega }\underset{A=1}{\overset{4}{}}\mathrm{d}^2\eta ^A\mathrm{d}^2\overline{\xi }^A\rho ^{4N7}}`$ (3.10)
$`{\displaystyle _0^{\mathrm{}}}drr^{4N3}\mathrm{e}^{2\rho ^2r^2}𝒵(\vartheta ,\overline{\vartheta };\mathrm{\Omega },r),`$
where all the numerical coefficients have been reinstated. In (3.10) we have introduced the density
$$𝒵(\vartheta ,\overline{\vartheta };\mathrm{\Omega },r)=\mathrm{exp}\left[\frac{ig_{_{\mathrm{YM}}}}{\pi r}\overline{\vartheta }_A^u\mathrm{\Omega }^{AB}\vartheta _{Bu}\right],$$
(3.11)
where the symplectic form $`\mathrm{\Omega }^{AB}`$ is given by
$$\mathrm{\Omega }^{AB}=\overline{\mathrm{\Sigma }}_i^{AB}\mathrm{\Omega }^i,\underset{i=1}{\overset{6}{}}\left(\mathrm{\Omega }^i\right)^2=1.$$
(3.12)
Gauge invariant operators depend on the $`\nu _u^A`$ and $`\overline{\nu }^{Au}`$ variables via colour singlet bilinears. These arise in symmetric or anti-symmetric combinations transforming respectively in the $`\mathrm{𝟏𝟎}`$ and $`\mathrm{𝟔}`$ dimensional representations of the SU(4) R-symmetry
$`(\overline{\nu }^A\nu ^B)_{\mathrm{𝟏𝟎}}`$ $``$ $`\overline{\nu }_{\text{}}^{u(A}\nu _u^{B)}=(\overline{\nu }^{Au}\nu _u^B+\overline{\nu }^{Bu}\nu _u^A),`$ (3.13)
$`(\overline{\nu }^A\nu ^B)_\mathrm{𝟔}`$ $``$ $`\overline{\nu }_{\text{}}^{u[A}\nu _u^{B]}=(\overline{\nu }^{Au}\nu _u^B\overline{\nu }^{Bu}\nu _u^A).`$ (3.14)
Using the generating function defined in (3.10) the $`\overline{\nu }^A\nu ^B`$ bilinears in the operators $`𝒪`$ and $`\overline{𝒪}`$ in (3.6) can be rewritten in terms of derivatives of $`𝒵(\vartheta ,\overline{\vartheta };\mathrm{\Omega },r)`$ with respect to the sources, $`\vartheta _A`$ and $`\overline{\vartheta }_B`$. The result for a two-point function in which the operator insertions contain a total of $`p`$ $`(\overline{\nu }\nu )_\mathrm{𝟔}`$ and $`q`$ $`(\overline{\nu }\nu )_{\mathrm{𝟏𝟎}}`$ bilinears is of the form
$`\overline{𝒪}(x_1)𝒪(x_2)={\displaystyle \frac{g_{_{\mathrm{YM}}}^8\mathrm{e}^{2\pi i\tau }}{(N1)!(N2)!}}{\displaystyle d\rho \mathrm{d}^4x_0\mathrm{d}^5\mathrm{\Omega }\underset{A=1}{\overset{4}{}}\mathrm{d}^2\eta ^A\mathrm{d}^2\overline{\xi }^A\rho ^{4N7}}`$ (3.15)
$`{\displaystyle drr^{4N3}\mathrm{e}^{2\rho ^2r^2}\frac{\delta ^{2p+2q}𝒵[\vartheta ,\overline{\vartheta };\mathrm{\Omega },r]}{\delta \vartheta _{u_1[A_1}\delta \overline{\vartheta }_{B_1]}^{u_1}\delta \vartheta _{v_1(C_1}\delta \overline{\vartheta }_{D_1)}^{v_1}\mathrm{}}}|_{\vartheta =\overline{\vartheta }=0}\stackrel{~}{\overline{𝒪}}(x_1;x_0,\rho ,\eta ,\overline{\xi })\stackrel{~}{𝒪}(x_2;x_0,\rho ,\eta ,\overline{\xi }),`$
where $`\stackrel{~}{𝒪}`$ and $`\stackrel{~}{\overline{𝒪}}`$ contain the dependence on the exact moduli, $`\eta ^A`$ and $`\overline{\xi }^A`$, and on the bosonic collective coordinates after extracting the $`\overline{\nu }^A\nu ^B`$ bilinears. Computing the $`r`$ integral gives
$`\overline{𝒪}(x_1)𝒪(x_2)`$ $``$ $`\alpha (p,q;N)g_{_{\mathrm{YM}}}^{8+p+q}\mathrm{e}^{2\pi i\tau }{\displaystyle d\rho \mathrm{d}^4x_0\mathrm{d}^5\mathrm{\Omega }\underset{A=1}{\overset{4}{}}\mathrm{d}^2\eta ^A\mathrm{d}^2\overline{\xi }^A\rho ^{p+q5}}`$ (3.16)
$`f(\mathrm{\Omega })\stackrel{~}{\overline{𝒪}}(x_1;\rho ,x_0;\eta ,\overline{\xi })\stackrel{~}{𝒪}(x_2;\rho ,x_0;\eta ,\overline{\xi }),`$
where $`f(\mathrm{\Omega })`$ contains the dependence on the $`\mathrm{\Omega }^{AB}`$ variables obtained from the derivatives of $`𝒵(\vartheta ,\overline{\vartheta };\mathrm{\Omega },r)`$. The coefficient $`\alpha (p,q;N)`$ contains the $`N`$ dependence and in the large $`N`$ limit we find
$`\alpha (p,q;N)`$ $`=`$ $`{\displaystyle \frac{2^{2N+\frac{1}{2}(p+q)}\pi ^{(p+q)}\mathrm{\Gamma }\left(2N1\frac{1}{2}(p+q)\right)}{(N1)!(N2)!}}\left(N^{p+\frac{q}{2}}+O(N^{p+\frac{q}{2}1})\right)`$ (3.17)
$`=`$ $`{\displaystyle \frac{N^{\frac{1}{2}(p+1)}}{4\pi ^{p+q+\frac{1}{2}}}}\left(1+O(1/N)\right).`$
From (3.16) and (3.17) it follows that the insertion of a $`(\overline{\nu }\nu )_{\mathrm{𝟏𝟎}}`$ or $`(\overline{\nu }\nu )_\mathrm{𝟔}`$ bilinear in a correlation function produces a factor of $`g_{_{\mathrm{YM}}}`$ or $`g_{_{\mathrm{YM}}}\sqrt{N}`$ respectively.
In computing the moduli space integrations in expressions for two-point functions of the type (3.16) it will prove convenient to calculate first the fermionic integrals over $`\eta ^A`$ and $`\overline{\xi }^A`$ and the angular integration over the five-sphere. These give rise to selection rules that determine which operators receive instanton contributions to their scaling dimensions. In particular since the superconformal modes are exact a correlation function can only receive instanton contribution if the operator expressions contain exactly sixteen such modes in the combination
$$\underset{A=1}{\overset{4}{}}\left(\eta ^{\alpha A}\eta _\alpha ^A\right)\left(\overline{\xi }_{\dot{\alpha }}^A\overline{\xi }^{\dot{\alpha }A}\right).$$
(3.18)
The integration over the five-sphere parametrised by the angular variables $`\mathrm{\Omega }^{AB}`$ factorises and gives rise to further selection rules. It gives a non-vanishing result only if the SU(4) indices carried by the $`\mathrm{\Omega }`$’s in the two operators can be combined to form a SU(4) singlet. The SU(4) indices are originally carried by the fermion modes which are all in the $`\mathrm{𝟒}`$, so the only possible singlet combinations correspond to products of $`\epsilon ^{ABCD}`$ tensors. The generic five-sphere integral is of the form
$$\mathrm{d}^5\mathrm{\Omega }\mathrm{\Omega }^{A_1B_1}\mathrm{}\mathrm{\Omega }^{A_{2n}B_{2n}}=c(n)\left(\epsilon ^{A_1B_1A_2B_2}\mathrm{}\epsilon ^{A_{2n1}B_{2n1}A_{2n}B_{2n}}+\mathrm{permutations}\right),$$
(3.19)
where the normalisation constant $`c(n)`$ is
$$c(n)=\frac{\pi ^{5/2}\mathrm{\Gamma }\left(n+\frac{1}{2}\right)}{2\mathrm{\Gamma }\left(n+4\right)}.$$
(3.20)
Equations (3.18) and (3.19) imply that a two-point function can receive a non-zero contribution only if the combined profiles of the two operators contain fermion modes of the four flavours with the same multiplicity.
The bosonic integrations over the position and size of the instanton are left as a last step. In the case of two-point functions these integrals are logarithmically divergent, signalling a contribution to the matrix of anomalous dimensions.
## 4 Fermion zero modes
In order to evaluate instanton induced correlation functions we need to integrate the classical profiles of the relevant composite operators over the instanton moduli space. We are interested in the dependence on the collective coordinates and of particular relevance will be the way the fermionic modes enter into the expressions for the various fields. The zero-mode dependence in the elementary fields of the $`𝒩`$=4 SYM multiplet was reviewed in detail in . Here we briefly summarise the features which will be relevant for the analysis of two-point functions of BMN operators.
The field equations of the $`𝒩`$=4 SYM theory admit a solution in which the gauge potential corresponds to a standard instanton of SU($`N`$) pure Yang–Mills theory and all the other fields vanish,
$$A_\mu =A_\mu ^I,\phi ^{AB}=\lambda _\alpha ^A=\overline{\lambda }_A^{\dot{\alpha }}=0.$$
(4.1)
However, the Dirac operator has zero modes in the background of this solution, i.e. the equation $`\overline{/D}_{\dot{\alpha }\alpha }\lambda ^{\alpha A}=0`$ has non-trivial solutions when the covariant derivative is evaluated in the background of an instanton. The general solution to the Dirac equation is linear in the instanton fermion zero modes. This non-trivial solution gives rise to a non-zero solution for the scalar fields when plugged into the corresponding equation, $`D^2\phi ^{AB}=\sqrt{2}[\lambda ^A,\lambda ^B]`$. The latter admits a solution for the scalar which is bilinear in the fermion modes. Proceeding with this iterative solution of the field equations one generates a complete supermultiplet and further iterations give rise to additional terms with more fermion modes in each field. The general solution obtained through this procedure is schematically of the form
$`A_\mu =\begin{array}{c}{\displaystyle \underset{n=0}{}}\\ 4n8N\end{array}A_\mu ^{(4n)},\phi ^{AB}=\begin{array}{c}{\displaystyle \underset{n=0}{}}\\ 4n+28N\end{array}\phi ^{(2+4n)AB}`$ (4.6)
$`\lambda _\alpha ^A=\begin{array}{c}{\displaystyle \underset{n=0}{}}\\ 4n+18N\end{array}\lambda _\alpha ^{(1+4n)A},\overline{\lambda }_{\dot{\alpha }A}=\begin{array}{c}{\displaystyle \underset{n=0}{}}\\ 4n+38N\end{array}\overline{\lambda }_{\dot{\alpha }A}^{(3+4n)},`$ (4.11)
where the notation $`\mathrm{\Phi }^{(n)}`$ is used to denote a term in the solution for the field $`\mathrm{\Phi }`$ containing $`n`$ fermion zero modes. It is also understood that in (4.11) the number of superconformal modes in each field does not exceed 16 and the remaining modes are of $`\nu _u^A`$ and $`\overline{\nu }^{Au}`$ type.
In computing the expressions for gauge invariant composite operators we shall make use of the ADHM description in which the elementary fields are written as $`[N+2]\times [N+2]`$ matrices. In particular, the leading order term in the solution for the scalars $`\phi ^{AB}`$ is given explicitly in appendix A. The solution of the iterative equations becomes very involved after a few steps. However the flavour structure of the combination of fermion zero modes in each term can be determined without actually solving the equations and is sufficient to identify which operators can get an instanton correction to their scaling dimension.
All the fermion zero modes, both the superconformal ones, $`\eta _\alpha ^A`$ and $`\overline{\xi }^{\dot{\alpha }A}`$, and the modes of type $`\nu _u^A`$ and $`\overline{\nu }^{Au}`$, transform in the $`\mathrm{𝟒}`$ of SU(4). We shall denote a generic fermion mode by $`𝓂_\mathrm{f}^A`$. The starting point for the construction of the instanton supermultiplet is the classical instanton, $`A_\mu ^{(0)}`$, which has no fermions. The first term in $`\lambda _\alpha ^A`$ is linear in the fermion modes
$$\lambda _\alpha ^{(1)A}𝓂_\mathrm{f}^A.$$
(4.12)
For the term $`\phi ^{(2)AB}`$ in the scalar solution one finds
$$\phi ^{(2)AB}𝓂_\mathrm{f}^{[A}𝓂_\mathrm{f}^{B]},$$
(4.13)
i.e. the two fermion modes are antisymmetrised in order to form a combination in the $`\mathrm{𝟔}`$. The $`\overline{\mathrm{𝟒}}`$ spinor $`\overline{\lambda }_A^{(3)\dot{\alpha }}`$ contains three fermion modes in the combination
$$\overline{\lambda }_A^{(3)\dot{\alpha }}\epsilon _{ABCD}𝓂_\mathrm{f}^B𝓂_\mathrm{f}^C𝓂_\mathrm{f}^D,$$
(4.14)
so that the component $`\lambda _A^{(3)}`$ has one mode of each flavour apart from $`A`$. Proceeding in the multiplet we find the quartic term in the solution for the vector, $`A_\mu ^{(4)}`$, which contains one fermion mode of each flavour in a singlet combination
$$A_\mu ^{(4)}\epsilon _{ABCD}𝓂_\mathrm{f}^A𝓂_\mathrm{f}^B𝓂_\mathrm{f}^C𝓂_\mathrm{f}^D.$$
(4.15)
The following term is $`\lambda _\alpha ^{(5)A}`$, which has flavour structure
$$\lambda _\alpha ^{(5)A}\epsilon _{BCDE}𝓂_\mathrm{f}^A𝓂_\mathrm{f}^B𝓂_\mathrm{f}^C𝓂_\mathrm{f}^D𝓂_\mathrm{f}^E,$$
(4.16)
i.e. it involves a mode of flavour $`A`$ plus one of each flavour. Then we find $`\phi ^{(6)AB}`$ that contains an antisymmetric combination of a mode of flavour $`A`$ and one of flavour $`B`$ plus one mode of each flavour antisymmetrised in a singlet,
$$\phi ^{(6)AB}\epsilon _{CDEF}𝓂_\mathrm{f}^{[A}𝓂_\mathrm{f}^{B]}𝓂_\mathrm{f}^C𝓂_\mathrm{f}^D𝓂_\mathrm{f}^E𝓂_\mathrm{f}^F.$$
(4.17)
The previous expressions are symbolic and the products of $`𝓂_\mathrm{f}`$’s in (4.13)-(4.17) correspond to different combinations of the modes $`\eta _\alpha ^A`$, $`\overline{\xi }^{\dot{\alpha }A}`$, $`\nu _u^A`$ and $`\overline{\nu }^{Au}`$ in the various entries of the ADHM matrix for each field. The structure of the terms with more fermion modes in the multiplet can be determined analogously. The iterative solution of the field equation to construct the first few terms in the multiplet was carried out explicitly in .
From the above equations we can deduce the form of the component fields in the decomposition relevant for the BMN limit. For the scalars in (2.4) we have
$`Z^{(2)}𝓂_\mathrm{f}^{[1}𝓂_\mathrm{f}^{4]},\overline{Z}^{(2)}𝓂_\mathrm{f}^{[2}𝓂_\mathrm{f}^{3]}`$
$`\phi ^{(2)\mathrm{\hspace{0.17em}1},3}𝓂_\mathrm{f}^{[1}𝓂_\mathrm{f}^{3]}+𝓂_\mathrm{f}^{[2}𝓂_\mathrm{f}^{4]},\phi ^{(2)\mathrm{\hspace{0.17em}2},4}𝓂_\mathrm{f}^{[1}𝓂_\mathrm{f}^{2]}+𝓂_\mathrm{f}^{[3}𝓂_\mathrm{f}^{4]}.`$ (4.18)
For the fermions in (2.7) and (2.9) we have respectively
$$\psi ^{(1)a}𝓂_\mathrm{f}^a,\overline{\psi }_a^{+(1)}\left(M^+𝓂_\mathrm{f}\right)_a,a=1,4\dot{a}=2,3$$
(4.19)
and
$$\psi _{\dot{a}}^{+(3)}\epsilon _{\dot{r}BCD}𝓂_\mathrm{f}^B𝓂_\mathrm{f}^C𝓂_\mathrm{f}^D,\overline{\psi }^{(3)\dot{a}}\left(M^{}\epsilon \right)^{\dot{a}}{}_{BCD}{}^{}𝓂_{\mathrm{f}}^{B}𝓂_\mathrm{f}^C𝓂_\mathrm{f}^D,a=1,4,\dot{a}=2,3,$$
(4.20)
so that
$`\psi ^{(1)\mathrm{\hspace{0.17em}1}}𝓂_\mathrm{f}^1,\psi ^{(1)\mathrm{\hspace{0.17em}4}}𝓂_\mathrm{f}^4,\overline{\psi }^{+(1)\mathrm{\hspace{0.17em}1}}𝓂_\mathrm{f}^2,\overline{\psi }^{+(1)\mathrm{\hspace{0.17em}4}}𝓂_\mathrm{f}^3,`$ (4.21)
$`\psi ^{+(3)\mathrm{\hspace{0.17em}2}}𝓂_\mathrm{f}^1𝓂_\mathrm{f}^3𝓂_\mathrm{f}^4,\psi ^{+(3)\mathrm{\hspace{0.17em}3}}𝓂_\mathrm{f}^1𝓂_\mathrm{f}^2𝓂_\mathrm{f}^4,\overline{\psi }^{(3)\mathrm{\hspace{0.17em}2}}𝓂_\mathrm{f}^2𝓂_\mathrm{f}^3𝓂_\mathrm{f}^4,\overline{\psi }^{(3)\mathrm{\hspace{0.17em}3}}𝓂_\mathrm{f}^1𝓂_\mathrm{f}^2𝓂_\mathrm{f}^3.`$ (4.22)
The terms of higher order are easily deduced from the previous general discussion. Notice that we can assign U(1) charge $`+\frac{1}{2}`$ to the fermion modes $`𝓂_\mathrm{f}^1`$ and $`𝓂_\mathrm{f}^4`$ and charge $`\frac{1}{2}`$ to the modes $`𝓂_\mathrm{f}^2`$ and $`𝓂_\mathrm{f}^3`$.
The dependence on the superconformal modes, $`\eta _\alpha ^A`$ and $`\overline{\xi }^{\dot{\alpha }A}`$, can be obtained using supersymmetry without solving the field equations. These modes are associated with superconformal symmetries broken in the instanton background. The corresponding terms in the $`𝒩`$=4 supermultiplet can thus be generated acting with the broken Poincaré and special supersymmetries, $`Q_{\alpha A}`$ and $`\overline{S}_{\dot{\alpha }A}`$, on the classical instanton solution for the gauge potential. In the case of SU(2) gauge group there are no additional fermion modes and the complete solution can be obtained in this way. In general, however, the dependence on the $`\nu _u^A`$ and $`\overline{\nu }^{Au}`$ modes can be determined only by solving the equations of motion.
It is useful to discuss the derivation of the dependence on the superconformal modes using supersymmetry since it allows us to clarify how different combinations of fermion modes appear in various operators.
Substituting $`A_\mu ^{(0)}A_\mu ^I`$ in the supersymmetry transformation of $`\lambda _\alpha ^A`$ gives $`\lambda _\alpha ^{(1)A}`$, which is linear in $`\eta _\alpha ^A`$ and $`\overline{\xi }^{\dot{\alpha }A}`$ and solves the corresponding field equation. Replacing $`\lambda _\alpha ^A`$ by $`\lambda ^{(1)A}`$ in the variation of $`\phi ^{AB}`$ generates the solution $`\phi ^{(2)AB}`$ for the scalar. The iteration of this procedure gives rise to $`\overline{\lambda }_{\dot{\alpha }A}^{(3)}`$, then to the correction $`A_\mu ^{(4)}`$ to the gauge field and so on.
In the examples studied in the superconformal modes always appear in the expressions of gauge-invariant composite operators in the combination
$$\zeta _\alpha ^A(x)=\frac{1}{\sqrt{\rho }}\left[\rho \eta _\alpha ^A(xx_0)_\mu \sigma _{\alpha \dot{\alpha }}^\mu \overline{\xi }^{\dot{\alpha }A}\right].$$
(4.23)
In general, however, the dependence on the fermion superconformal modes, even in gauge-invariant operators, is not only through this combination and instead the moduli $`\eta _\alpha ^A`$ and $`\overline{\xi }^{\dot{\alpha }A}`$ appear explicitly. This can be understood analysing the form of the Poincaré and special supersymmetry variations of the fields. Under a combination of the broken supersymmetries, $`\eta ^{\alpha A}Q_{\alpha A}+\overline{\xi }_{\dot{\alpha }}^A\overline{S}_A^{\dot{\alpha }}`$, we have
$`\delta A_\mu =(\eta ^A+\sigma x\overline{\xi }^A)\sigma _\mu \overline{\lambda }_A\text{}`$ (4.24)
$`\delta \lambda ^A=F_{\mu \nu }\sigma ^{\mu \nu }(\eta ^A+\sigma x\overline{\xi }^A)+[\phi ^{AB},\overline{\phi }_{BC}](\eta ^C+\sigma x\overline{\xi }^C)\text{}`$ (4.25)
$`\delta \phi ^{AB}=\lambda ^A(\eta ^B+\sigma x\overline{\xi }^B)(AB)`$ (4.26)
$`\delta \overline{\lambda }_A=/D\overline{\phi }_{AB}(\eta ^B+\sigma x\overline{\xi }^B)+\overline{\phi }_{AB}\overline{\xi }^B,\text{}`$ (4.27)
which shows that, whereas the variations of $`A_\mu `$, $`\lambda _\alpha ^A`$ and $`\phi ^{AB}`$ involve the combination $`\zeta _\alpha ^A`$, the superconformal variation of $`\overline{\lambda }_A^{\dot{\alpha }}`$ contains an extra term. Therefore the profiles of operators involving $`\overline{\lambda }_A^{\dot{\alpha }}`$ in general depend separately on $`\eta _\alpha ^A`$ and $`\overline{\xi }_{\dot{\alpha }}^A`$. Since further application of the broken supersymmetries generates new terms in the solution for the elementary fields, it follows that not only operators containing $`\overline{\lambda }_{\dot{\alpha }A}`$, but also those in which any elementary field contain a non-minimal number of fermion modes (e.g. $`A_\mu ^{(4)}`$, $`\lambda _\alpha ^{(5)A}`$, $`\phi ^{(6)AB}`$) will depend on $`\eta _\alpha ^A`$ and $`\overline{\xi }_{\dot{\alpha }}^A`$ not only via $`\zeta _\alpha ^A`$. This observation will play an important rôle in the case of two impurity BMN operators. As will be shown in the next section, a naive counting of zero-modes including only terms with the minimal number of fermion modes in each field would lead to conclude that these operators have vanishing two-point functions in the instanton background. We will, however, argue that the inclusion of the term $`\phi ^{(6)AB}`$ in the solution is needed in order to compute the leading non-zero instanton contributions to these two-point functions.
## 5 Two-point functions of BMN operators
In this section we analyse instanton contributions to two-point functions of the BMN operators described in section 2.2. Using the results of the previous sections we shall determine which operators have non-zero two-point functions in the instanton background and the dependence of the instanton induced anomalous dimensions on the parameters, $`g_{_{\mathrm{YM}}}`$, $`N`$ and $`J`$ as well as the integers corresponding to the mode numbers in the dual string states. Zero and one impurity operators are protected, therefore their two-point functions are not renormalised and receive no instanton contribution. We shall therefore discuss two and four impurity operators. Operators with an odd number of impurities are expected not to receive instanton contributions. The analysis of two impurity operators in the next subsection will be rather qualitative because the leading non-zero contribution to their two-point functions involves the six-fermion term in the scalar solution which is not known explicitly. The four impurity case which is fully under control will be discussed in the following subsection.
### 5.1 Two impurity operators
At the two impurity level we focus on the bosonic SO(4)<sub>C</sub> singlet operator (2.2.1) in the $`\mathrm{𝟗}`$ of SO(4)<sub>R</sub>. Since this sector contains only one operator there is no problem of mixing and the anomalous dimension of the operator $`𝒪_{J,\mathrm{𝟗};n}^{\{ij\}}`$ can be read directly from the coefficient of the two-point function $`𝒪_{J,\mathrm{𝟗};n}^{\{ij\}}(x_1)\overline{𝒪}_{J,\mathrm{𝟗};m}^{\{kl\}}(x_2)`$.
As usual it is convenient to compute this two-point function for a particular choice of components, rather than work in a manifestly SO(4)<sub>R</sub> covariant way. Therefore we consider <sup>4</sup><sup>4</sup>4The subscripts indicating the SO(4)<sub>R</sub> representation and the U(1) charge will be omitted except in situations where this may cause confusion.
$$G_\mathrm{𝟗}(x_1,x_2)=𝒪_n^{\{13\}}(x_1)\overline{𝒪}_m^{\{13\}}(x_2)_{\mathrm{inst}},$$
(5.1)
so that there is no trace to subtract.
The component $`𝒪_n^{\{13\}}`$ is
$$𝒪_n^{\{13\}}=\frac{i}{\sqrt{J\left(\frac{g_{_{\mathrm{YM}}}^2N}{8\pi ^2}\right)^{J+2}}}\underset{p=0}{\overset{J}{}}\mathrm{e}^{2\pi ipn/J}\left[\mathrm{Tr}\left(Z^{Jp}\phi ^{13}Z^p\phi ^{13}\right)\mathrm{Tr}\left(Z^{Jp}\phi ^{24}Z^p\phi ^{24}\right)\right]$$
(5.2)
and the conjugate operator is
$$\overline{𝒪}_n^{\{13\}}=\frac{i}{\sqrt{J\left(\frac{g_{_{\mathrm{YM}}}^2N}{8\pi ^2}\right)^{J+2}}}\underset{p=0}{\overset{J}{}}\mathrm{e}^{2\pi ipn/J}\left[\mathrm{Tr}\left(\overline{Z}^{Jp}\phi ^{13}\overline{Z}^p\phi ^{13}\right)\mathrm{Tr}\left(\overline{Z}^{Jp}\phi ^{24}\overline{Z}^p\phi ^{24}\right)\right].$$
(5.3)
The semi-classical approximation in the one-instanton sector for (5.1) gives
$`G_\mathrm{𝟗}(x_1,x_2)={\displaystyle \frac{\pi ^{4N}g_{_{\mathrm{YM}}}^{4N}\mathrm{e}^{2\pi i\tau }}{(N1)!(N2)!}}{\displaystyle d\rho \mathrm{d}^4x_0\underset{A=1}{\overset{4}{}}\mathrm{d}^2\eta ^A\mathrm{d}^2\overline{\xi }^A\mathrm{d}^{N2}\nu ^A\mathrm{d}^{N2}\overline{\nu }^A\rho ^{4N13}}`$ (5.4)
$`\mathrm{e}^{\frac{\pi ^2}{16g_{_{\mathrm{YM}}}^2\rho ^2}\epsilon _{ABCD}(\overline{\nu }^{[A}\nu ^{B]})(\overline{\nu }^{[C}\nu ^{D]})}\widehat{𝒪}_{J;n}^{\{13\}}(x_1;x_0,\rho ,\eta ,\overline{\xi },\nu ,\overline{\nu })\widehat{\overline{𝒪}}_{J;m}^{\{13\}}(x_2;x_0,\rho ,\eta ,\overline{\xi },\nu ,\overline{\nu }).`$
In order to have a non-zero contribution to this two-point function the classical profiles of the two operators must contain, when combined, the sixteen fermion modes corresponding to the broken supersymmetries.
It is easy to verify that substituting for each scalar field in (5.1) the leading order solution, $`\phi ^{(2)}`$, does not allow to soak up all the superconformal modes. Using for each scalar field the bilinear term (4.18) in the solution we find that $`𝒪^{\{13\}}`$ contains the following combinations of fermion modes
$$\left(𝓂_\mathrm{f}^{[1}𝓂_\mathrm{f}^{4]}\right)^J\left(𝓂_\mathrm{f}^{[1}𝓂_\mathrm{f}^{3]}𝓂_\mathrm{f}^{[1}𝓂_\mathrm{f}^{3]}+𝓂_\mathrm{f}^{[2}𝓂_\mathrm{f}^{4]}𝓂_\mathrm{f}^{[2}𝓂_\mathrm{f}^{4]}\right).$$
(5.5)
Similarly, the conjugate operator $`\widehat{\overline{𝒪}}^{\{13\}}`$ contains
$$\left(𝓂_\mathrm{f}^{[2}𝓂_\mathrm{f}^{3]}\right)^J\left(𝓂_\mathrm{f}^{[1}𝓂_\mathrm{f}^{3]}𝓂_\mathrm{f}^{[1}𝓂_\mathrm{f}^{3]}+𝓂_\mathrm{f}^{[2}𝓂_\mathrm{f}^{4]}𝓂_\mathrm{f}^{[2}𝓂_\mathrm{f}^{4]}\right).$$
(5.6)
The argument given at the end of section 4 shows that in these traces, involving only $`Z^{(2)}`$ and $`\phi ^{(2)AB}`$, the superconformal modes always appear in the combination $`\zeta ^A`$ of (4.23). This can be easily verified using the explicit expression for the scalar ADHM matrices given in (A.8). Then, because of the condition $`\left(\zeta ^A(x)\right)^3=0`$ satisfied by the Weyl spinors $`\zeta ^A`$, in order to soak up the sixteen superconformal modes in the two-point function (5.4) each of the operators should contain two factors of $`\zeta ^A`$ for each flavour. In other words the sixteen superconformal modes should appear in the two-point function in the form
$$\underset{A=1}{\overset{4}{}}\left[\zeta ^A(x_1)\right]^2\left[\zeta ^A(x_2)\right]^2.$$
(5.7)
Examining the combinations (5.5) and (5.6) it is clear that this is not possible. $`\widehat{𝒪}^{\{13\}}`$ cannot soak up the required superconformal modes of flavour 2 and 3, since it does not contain two factors of both $`\zeta ^2`$ and $`\zeta ^3`$, while $`\widehat{\overline{𝒪}}^{\{13\}}`$ cannot soak up all the superconformal modes of flavour 1 and 4, since it does not contain two factors of both $`\zeta ^1`$ and $`\zeta ^4`$. This simple analysis of the flavour structure of the superconformal modes in the classical profiles of the operators shows that the two-point function (5.1) vanishes at leading order in the instanton background. This argument does not rely on the way the remaining $`J2`$ $`(\overline{\nu }^A\nu ^B)`$ bilinears are distributed in the two operators. According to the discussion in section 3 the leading contribution in the large $`N`$ limit would come from terms in which all the $`(\overline{\nu }^A\nu ^B)`$ bilinears are antisymmetrised. However, since the above argument is based only on the analysis of the superconformal modes the conclusion that the leading $`g_{_{\mathrm{YM}}}`$ contribution to the two-point function (5.1) vanishes is valid at all orders in $`1/N`$.
In order to saturate the integrations over the superconformal modes in (5.4) we need to use for some of the scalar fields the solution containing six fermionic modes, $`\phi ^{(6)AB}`$. Inspecting the combinations (5.5) and (5.6) found at leading order and recalling (4.17) it is easy to verify that it is sufficient to consider one $`\phi ^{(6)AB}`$ (or $`Z^{(6)}`$ and $`\overline{Z}^{(6)}`$ respectively) insertion in each operator. These are the leading order contributions, the insertion of more six-fermion scalars leads to contributions of higher order in $`g_{_{\mathrm{YM}}}`$ since in this case more $`\overline{\nu }^A\nu ^B`$ bilinears appear.
Recalling the form of $`\phi ^{(6)AB}`$ given in (4.17) we find that the combinations of fermionic modes in $`\mathrm{Tr}\left(Z^{Jp}\phi ^{AB}Z^p\phi ^{CD}\right)`$ and $`\mathrm{Tr}\left(\overline{Z}^{Jp}\phi ^{AB}\overline{Z}^p\phi ^{CD}\right)`$ are respectively
$$\epsilon _{A^{}B^{}C^{}D^{}}𝓂_\mathrm{f}^A^{}𝓂_\mathrm{f}^B^{}𝓂_\mathrm{f}^C^{}𝓂_\mathrm{f}^D^{}\left(𝓂_\mathrm{f}^{[1}𝓂_\mathrm{f}^{4]}\right)^J\left(𝓂_\mathrm{f}^{[A}𝓂_\mathrm{f}^{B]}𝓂_\mathrm{f}^{[C}𝓂_\mathrm{f}^{D]}\right)$$
(5.8)
and
$$\epsilon _{A^{}B^{}C^{}D^{}}𝓂_\mathrm{f}^A^{}𝓂_\mathrm{f}^B^{}𝓂_\mathrm{f}^C^{}𝓂_\mathrm{f}^D^{}\left(𝓂_\mathrm{f}^{[2}𝓂_\mathrm{f}^{3]}\right)^J\left(𝓂_\mathrm{f}^{[A}𝓂_\mathrm{f}^{B]}𝓂_\mathrm{f}^{[C}𝓂_\mathrm{f}^{D]}\right).$$
(5.9)
Notice that here the superconformal modes do not necessarily enter via (4.23), since we are using the term with six fermions in the solution for one of the fields in each operator. More precisely the structure of the supersymmetry transformations (4.24)-(4.27) shows that both traces contain one single $`\overline{\xi }`$ mode which is not part of a $`\zeta `$. This is crucial in order to get a non-zero result from the moduli space integration, because it allows, for two flavours, to distribute three fermionic superconformal modes in one operator and one in the other.
The resulting non-zero contribution to the two-point function is
$`G_\mathrm{𝟗}(x_1,x_2){\displaystyle \frac{g_{_{\mathrm{YM}}}^4\mathrm{e}^{2\pi i\tau }}{JN^{3/2}}}{\displaystyle \mathrm{d}^4x_0d\rho \rho ^{2J5}f(x_1,x_2;x_0,\rho )\mathrm{d}^5\mathrm{\Omega }\left(\mathrm{\Omega }^{14}\right)^{J1}\left(\mathrm{\Omega }^{23}\right)^{J1}\mathrm{\Omega }^{13}\mathrm{\Omega }^{24}}`$
$`\times {\displaystyle }{\displaystyle \underset{A=1}{\overset{4}{}}}\mathrm{d}^2\eta ^A\mathrm{d}^2\overline{\xi }^A\{[\left(\zeta ^1\right)^2\zeta ^2\left(\zeta ^3\right)^2\left(\zeta ^4\right)^2\overline{\xi }^1](x_1)[\zeta ^1\left(\zeta ^2\right)^2\left(\zeta ^3\right)^2\left(\zeta ^4\right)^2\overline{\xi }^2](x_2)`$
$`+[\left(\zeta ^1\right)^2\left(\zeta ^2\right)^2\zeta ^3\left(\zeta ^4\right)^2\overline{\xi }^4](x_1)[\left(\zeta ^1\right)^2\left(\zeta ^2\right)^2\left(\zeta ^3\right)^2\zeta ^4\overline{\xi }^3](x_2)\},`$ (5.10)
where the $`(\overline{\nu }^A\nu ^B)`$ bilinears have been rewritten in terms of $`\mathrm{\Omega }^{AB}`$’s as described in section 3. The overall powers of $`g_{_{\mathrm{YM}}}`$ and $`N`$ come from the normalisation of the operators, the moduli space integration measure and the $`(\overline{\nu }\nu )_\mathrm{𝟔}`$ bilinears, see equations (3.16) and (3.17). The fermion superconformal modes are saturated and the corresponding integration is non-zero. In (5.10) the dependence on the bosonic moduli has been collected in the function $`f(x_1,x_2;x_0,\rho )`$, which can only be computed knowing the explicit form of the solution $`\phi ^{(6)AB}`$ which we have not determined. The exact form of the solution is also needed in order to compute the overall coefficient and, in particular, the dependence on $`J`$. More details of the derivation of (5.10) as well as of the evaluation of the moduli space integrals are given in appendix B.
The final result for the two-point function is of the form
$$G_\mathrm{𝟗}(x_1,x_2)\frac{g_{_{\mathrm{YM}}}^4J^3\mathrm{e}^{2\pi i\tau }}{N^{3/2}}\frac{1}{(x_1x_2)^{2(J+2)}}I,$$
(5.11)
where $`I`$ is a logarithmically divergent integral, to be regulated e.g. by dimensional regularisation of the $`x_0`$ integral. The logarithmic divergence is due to the bosonic integrations over $`x_0`$ and $`\rho `$, as can be verified by dimensional analysis. The presence of this divergence signals an instanton contribution to the anomalous dimension of the operator $`𝒪_\mathrm{𝟗}^{\{ij\}}`$.
As already observed there is only one operator in the representation $`\mathrm{𝟗}`$ of SO(4)<sub>R</sub> and thus there is no mixing to resolve and the present analysis directly determines the instanton correction to the scaling dimension. We thus find that the instanton induced anomalous dimension of $`𝒪_\mathrm{𝟗}^{\{13\}}`$ behaves as
$$\gamma _\mathrm{𝟗}^{\mathrm{inst}}\frac{g_{_{\mathrm{YM}}}^4J^3}{N^{3/2}}\mathrm{e}^{\frac{8\pi ^2}{g_{_{\mathrm{YM}}}^2}+i\theta }\left(g_2\right)^{7/2}\left(\lambda ^{}\right)^2\mathrm{e}^{\frac{8\pi ^2}{g_2\lambda ^{}}+i\theta }.$$
(5.12)
This is in agreement with the non-perturbative correction to the mass of the dual string state computed in . In particular, the anomalous dimension (5.12) is independent of the parameter $`n`$ corresponding to the mode number of the plane-wave string state. Apart from the exponential factor characteristic of instanton effects, (5.12) contains an additional factor of $`\left(\lambda ^{}\right)^2`$. This is due to the inclusion of six-fermion scalars which give rise to additional $`(\overline{\nu }\nu )_\mathrm{𝟔}`$ bilinears, each of which brings one more power of $`g_{_{\mathrm{YM}}}`$. As will be shown in the next subsection in the case of four impurity SO(4)<sub>R</sub> singlets it is sufficient to consider the bilinear solution for all the scalars and as a consequence we shall find a leading contribution of order $`(g_2)^{7/2}\mathrm{e}^{8\pi ^2/g_2\lambda ^{}}`$.
Contributions in which some of the $`\overline{\nu }\nu `$ bilinears are in the $`\mathrm{𝟏𝟎}`$ of SU(4) give rise to subleading corrections which are suppressed by powers of $`1/N`$.
Another class of contributions to (5.1) which are suppressed in the large $`N`$ limit are those in which pairs of scalars are contracted. In these terms the analysis of the superconformal modes is unaltered and in order to soak them up it is again necessary to use the solution $`\phi ^{(6)AB}`$ for two of the scalars. Two of the scalars which were previously replaced by $`\overline{\nu }\nu `$ bilinears are now contracted and do not contain any fermion modes. Hence the integration over the moduli space produces one less power of $`g_{_{\mathrm{YM}}}^2N`$. However, with the normalisation we are using the propagator is proportional to $`g_{_{\mathrm{YM}}}^2`$, so that in conclusion the contribution of these terms is down by $`1/N`$ with respect to (5.11) because there is no power of $`N`$ associated with the contraction.
A careful analysis of both types of $`1/N`$ corrections shows that they give a contribution to the anomalous dimension of the operator $`𝒪_\mathrm{𝟗}`$ of order $`(g_2)^{9/2}(\lambda ^{})^2`$. These are the leading terms in a power series in $`g_2`$. In general the corrections to the semi-classical approximation in the BMN limit can be reorganised into a double series in $`g_2`$ and $`\lambda ^{}`$.
Operators in different sectors can be studied along the same lines. However, superconformal invariance implies that all the two impurity operators have the same anomalous dimension and thus the above result can be extended to two impurity operators in all the other sectors with no further calculations required.
Arguments similar to those discussed here, showing the vanishing of the leading one-instanton contribution to the two-point function $`G_\mathrm{𝟗}(x_1,x_2)`$, have been used to prove various non-renormalisation properties in . In view of the results we found for $`𝒪_\mathrm{𝟗}`$, one can expect that some of the non-renormalisation results of these papers may not be extended to higher orders in the coupling.
### 5.2 Four impurity operators
The calculation of two-point functions of four impurity operators is more involved than the corresponding calculation in the two impurity case from the point of view of the combinatorial analysis. However, at the four impurity level, in the case of SO(4)<sub>R</sub> singlets, the leading instanton contributions do not involve the six fermion solution for the scalar fields. A non-zero result is obtained using only the bilinear solution, which is known explicitly and given in (A.8), in computing the classical profiles of the operators. Therefore we can analyse in a quantitative way the semi-classical contributions to the two-point functions. The fact that non-zero correlation functions of singlet operators are obtained using the minimal number of fermion modes for each field also implies that in this case a contribution to the matrix of anomalous dimensions arises at leading order in the instanton background. As we shall see these operators have instanton induced anomalous dimension of order $`(g_2)^{7/2}\mathrm{e}^{2\pi i\tau }`$. Another difference with respect to the two impurity case studied in the previous section is that two-point functions of four impurity operators depend explicitly on the integers dual to the string mode numbers. We shall discuss in detail a SO(4)$`{}_{R}{}^{}\times `$SO(4)<sub>C</sub> singlet with four scalar impurities and show that the behaviour of its two-point functions is in remarkable agreement with the corresponding string calculation of . Other singlet operators can be analysed in a similar fashion. Operators in other sectors will be shown to receive contribution only at higher order in $`\lambda ^{}`$. This result follows simply from the analysis of fermion zero modes and is also in agreement with the string theory prediction.
#### 5.2.1 $`\epsilon `$-singlet operator
In this subsection we present the calculation of the one-instanton contribution to the two-point function of one particular SO(4)<sub>R</sub> singlet. More details are provided in appendix B. We focus on the four scalar impurity operator in which the SO(4)<sub>R</sub> indices are contracted via an $`\epsilon `$-tensor,
$`𝒪_{\mathrm{𝟏};J;n_1,n_2,n_3}`$ $`=`$ $`{\displaystyle \frac{\epsilon _{ijkl}}{\sqrt{J^3\left(\frac{g_{_{\mathrm{YM}}}^2N}{8\pi ^2}\right)^{J+4}}}}\begin{array}{c}{\displaystyle \underset{q,r,s=0}{\overset{J}{}}}\\ q+r+sJ\end{array}\mathrm{e}^{2\pi i[(n_1+n_2+n_3)q+(n_2+n_3)r+n_3s]/J}`$ (5.16)
$`\times \mathrm{Tr}\left(Z^{J(q+r+s)}\phi ^iZ^q\phi ^jZ^r\phi ^kZ^s\phi ^l\right).`$
The string state in the plane-wave background which is naturally identified as being dual to this operator is of the form
$$\epsilon _{ijkl}\alpha _{n_1}^i\alpha _{n_2}^j\stackrel{~}{\alpha }_{n_3}^k\stackrel{~}{\alpha }_{(n_1+n_2n_3)}^l|0_h,$$
(5.17)
where $`|0_h`$ is the BMN ground state and the contraction runs over values of the indices in one of the two SO(4) factors. $`D`$-instanton contributions to the renormalisation of the mass of this state were computed in . We shall return to the comparison with the string results at the end of this section. Notice, however, that the state (5.17) is antisymmetric under the exchange of the two left-moving or right-moving modes. The operator (5.16) on the other hand has no definite symmetry under permutations of the parameters $`n_1`$, $`n_2`$ and $`n_3`$. Therefore in order to construct a gauge theory operator that can be identified with (5.17) it will be necessary to explicitly antisymmetrise (5.16). This point will prove crucial when comparing instanton corrections to the scaling dimension of $`𝒪_\mathrm{𝟏}`$ to $`D`$-instanton induced corrections to the mass of the string state.
We are interested in the two-point function
$`G_\mathrm{𝟏}(x_1,x_2;n_1,n_2,n_3;m_1,m_2,m_3)=𝒪_{\mathrm{𝟏};n_1,n_2,n_3}(x_1)\overline{𝒪}_{\mathrm{𝟏};m_1,m_2,m_3}(x_2)_{\mathrm{inst}}`$
$`={\displaystyle \frac{\pi ^{4N}g_{_{\mathrm{YM}}}^{4N}\mathrm{e}^{2\pi i\tau }}{(N1)!(N2)!}}{\displaystyle d\rho \mathrm{d}^4x_0\underset{A=1}{\overset{4}{}}\mathrm{d}^2\eta ^A\mathrm{d}^2\overline{\xi }^A\mathrm{d}^{N2}\nu ^A\mathrm{d}^{N2}\overline{\nu }^A\rho ^{4N13}}`$ (5.18)
$`\times \mathrm{e}^{\frac{\pi ^2}{16g_{_{\mathrm{YM}}}^2\rho ^2}\epsilon _{ABCD}(\overline{\nu }^{[A}\nu ^{B]})(\overline{\nu }^{[C}\nu ^{D]})}\widehat{𝒪}_{\mathrm{𝟏};n_1,n_2,n_3}(x_1;x_0,\rho ,\eta ,\overline{\xi },\nu ,\overline{\nu })\widehat{\overline{𝒪}}_{\mathrm{𝟏};m_1,m_2,m_3}(x_2;x_0,\rho ,\eta ,\overline{\xi },\nu ,\overline{\nu }).`$
As usual the semi-classical approximation requires the calculation of the classical profiles of $`𝒪_\mathrm{𝟏}`$ and $`\overline{𝒪}_\mathrm{𝟏}`$ in the instanton background.
Summing over the SO(4) indices in (5.16) and using the relations (2.4) we find that the operator $`𝒪_\mathrm{𝟏}`$ contains the independent traces
$`+\mathrm{Tr}\left(Z^p\phi ^{12}Z^q\phi ^{13}Z^r\phi ^{24}Z^s\phi ^{34}\right)\mathrm{Tr}\left(Z^p\phi ^{12}Z^q\phi ^{34}Z^r\phi ^{24}Z^s\phi ^{13}\right)`$
$`+\mathrm{Tr}\left(Z^p\phi ^{12}Z^q\phi ^{24}Z^r\phi ^{34}Z^s\phi ^{13}\right)+\mathrm{Tr}\left(Z^p\phi ^{12}Z^q\phi ^{34}Z^r\phi ^{13}Z^s\phi ^{24}\right)`$
$`\mathrm{Tr}\left(Z^p\phi ^{12}Z^q\phi ^{13}Z^r\phi ^{34}Z^s\phi ^{24}\right)\mathrm{Tr}\left(Z^p\phi ^{12}Z^q\phi ^{24}Z^r\phi ^{13}Z^s\phi ^{34}\right),`$ (5.19)
where $`p=J(q+r+s)`$, plus three other groups of six traces obtained by cyclic permutations of the indices on the impurities in (5.19). The conjugate operator, $`\overline{𝒪}_\mathrm{𝟏}`$, contains the same terms, but with the $`Z`$’s replaced by $`\overline{Z}`$’s.
It is straightforward to verify that these traces, when evaluated in the instanton background, contain the correct combination of fermions required to soak up the superconformal modes in a two-point function and that this can be achieved using only the bilinear solution for all the scalars. In this case all the $`\eta ^A`$ and $`\overline{\xi }^A`$ modes in the gauge invariant traces are combined into $`\zeta ^A`$’s. In order to give rise to a non-zero two-point function in the one instanton sector both operators should then contain the combination $`_{A=1}^4\left(\zeta ^A\right)^2`$. To achieve this in each trace in (5.19) the four impurities must provide two superconformal modes of flavours 2 and 3, whereas the modes of flavour 1 and 4 can be taken from the impurities or from the $`Z`$’s. Similarly in the case of $`\overline{𝒪}_\mathrm{𝟏}`$ the superconformal modes of flavour 1 and 4 come necessarily from the impurities and those of flavour 2 and 3 can be provided by the impurities or by the $`\overline{Z}`$’s. As in the two impurity case studied in the previous section the leading contribution is obtained taking all the remaining modes in $`(\overline{\nu }\nu )_\mathrm{𝟔}`$ bilinears. In all the traces appearing in both $`𝒪_\mathrm{𝟏}`$ and $`\overline{𝒪}_\mathrm{𝟏}`$ the impurities contain two fermion modes of each flavour. The combination of fermion modes entering into all the terms in $`𝒪_\mathrm{𝟏}`$ is
$$\left(𝓂_\mathrm{f}^1\right)^{J+2}\left(𝓂_\mathrm{f}^2\right)^2\left(𝓂_\mathrm{f}^3\right)^2\left(𝓂_\mathrm{f}^4\right)^{J+2},$$
(5.20)
whereas all the terms in the expansion of $`\overline{𝒪}_\mathrm{𝟏}`$ contain
$$\left(𝓂_\mathrm{f}^1\right)^2\left(𝓂_\mathrm{f}^2\right)^{J+2}\left(𝓂_\mathrm{f}^3\right)^{J+2}\left(𝓂_\mathrm{f}^4\right)^2.$$
(5.21)
The leading contribution to the two-point function $`G_\mathrm{𝟏}(x_1,x_2)`$ in the semi-classical approximation arises from terms in the profiles of the operators containing the following combinations of fermion modes
$`𝒪_\mathrm{𝟏}`$ $``$ $`\left(\zeta ^1\right)^2\left(\zeta ^2\right)^2\left(\zeta ^3\right)^2\left(\zeta ^4\right)^2\left(\overline{\nu }^{[1}\nu ^{4]}\right)^J`$
$`\overline{𝒪}_\mathrm{𝟏}`$ $``$ $`\left(\zeta ^1\right)^2\left(\zeta ^2\right)^2\left(\zeta ^3\right)^2\left(\zeta ^4\right)^2\left(\overline{\nu }^{[2}\nu ^{3]}\right)^J.`$ (5.22)
As previously observed these combinations can be obtained in many different ways corresponding to the choice of which field, $`Z`$ or $`\phi `$, provides each of the $`\zeta `$’s of flavour 1 and 4 in $`𝒪`$. An equal number of different contributions arises from the ways of distributing the $`\zeta `$’s of flavour 2 and 3 among the impurities or the $`\overline{Z}`$’s in $`\overline{𝒪}`$. In order to simplify the discussion of the associated combinatorics it is convenient to introduce the following notation. We denote by $`\stackrel{ˇ}{\phi }^{AB}`$ a scalar solution in which only the $`\nu `$ and $`\overline{\nu }`$ modes are kept and all the superconformal modes are set to zero; scalars containing only bilinears in the superconformal modes are indicated by $`\stackrel{~}{\phi }^{AB}`$; the symbol $`\widehat{\phi }^{AB}`$ is used for scalar profiles in which only mixed terms, $`\zeta \nu `$ or $`\zeta \overline{\nu }`$, are included,
$`\stackrel{ˇ}{\phi }^{AB}`$ $``$ $`\phi ^{AB}(x,x_0,\rho ;\nu ,\overline{\nu };\eta =\overline{\xi }=0)`$ (5.23)
$`\stackrel{~}{\phi }^{AB}`$ $``$ $`\phi ^{AB}(x,x_0,\rho ;\eta ,\overline{\xi };\nu =\overline{\nu }=0)`$ (5.24)
$`\widehat{\phi }^{AB}`$ $``$ $`\phi ^{AB}(x,x_0,\rho ;\eta \nu ,\eta \overline{\nu },\overline{\xi }\nu ,\overline{\xi }\overline{\nu };\zeta \zeta =\overline{\nu }\nu =0).`$ (5.25)
The same notation is also used for $`Z\phi ^{14}`$ and $`\overline{Z}\phi ^{23}`$.
We are only interested in contributions to the two-point function $`G_\mathrm{𝟏}(x_1,x_2)`$ which survive in the BMN limit, $`N\mathrm{}`$, $`J\mathrm{}`$, with $`J^2/N`$ fixed. The leading large $`N`$ contributions are those in which the combinations (5.22) are selected, i.e. the superconformal modes are soaked up and all the remaining fields are replaced by $`(\overline{\nu }\nu )_\mathrm{𝟔}`$ bilinears. Within this class of terms the dominant ones in the large $`J`$ limit are those in which as many superconformal modes as possible are extracted from the $`Z`$’s and $`\overline{Z}`$’s, because there is roughly a multiplicity factor of $`J`$ associated with the choice of each $`Z`$ or $`\overline{Z}`$ providing one such mode. We first discuss these leading terms and we will then show that these are the only non-vanishing contributions in the BMN limit.
As shown by the previous preliminary analysis, in the operator $`𝒪_\mathrm{𝟏}`$ the modes $`\zeta ^2`$ and $`\zeta ^3`$ necessarily come from the impurities and thus the leading large $`J`$ terms arise from traces in which we take the two $`\zeta ^1`$ and the two $`\zeta ^4`$ modes from four distinct $`Z`$’s. Similar considerations apply to the $`\overline{𝒪}_\mathrm{𝟏}`$ operator with the rôle of the flavours $`(1,4)`$ and $`(2,3)`$ exchanged. Using the notation introduced in (5.23)-(5.25) this means that we consider traces in which all four impurities are $`\widehat{\phi }^{AB}`$ matrices and we choose four $`Z`$’s to be $`\widehat{Z}`$ matrices, with all the others being $`\stackrel{ˇ}{Z}`$’s. There is a total of 35 different traces of this type for each of the $`6\times 4`$ terms in the operator (5.16) and a similar counting applies to its conjugate. The 35 traces correspond to the inequivalent ways of choosing the four $`\widehat{Z}`$’s from the four groups of $`Z`$’s in (5.16). For the generic term in the operator, $`\mathrm{Tr}\left(Z^p\phi ^{A_1B_1}Z^q\phi ^{A_2B_2}Z^r\phi ^{A_3B_3}Z^s\phi ^{A_4B_4}\right)`$, with $`p=J(q+r+s)`$, we need to consider
$`\mathrm{Tr}\left(\stackrel{ˇ}{Z}^{p_1}\widehat{Z}\stackrel{ˇ}{Z}^{p_2}\widehat{Z}\stackrel{ˇ}{Z}^{p_3}\widehat{Z}\stackrel{ˇ}{Z}^{p_4}\widehat{Z}\stackrel{ˇ}{Z}^{p_5}\widehat{\phi }^{A_1B_1}\stackrel{ˇ}{Z}^q\widehat{\phi }^{A_2B_2}\stackrel{ˇ}{Z}^r\widehat{\phi }^{A_3B_3}\stackrel{ˇ}{Z}^s\widehat{\phi }^{A_4B_4}\right)`$
$`\mathrm{Tr}\left(\stackrel{ˇ}{Z}^{p_1}\widehat{Z}\stackrel{ˇ}{Z}^{p_2}\widehat{Z}\stackrel{ˇ}{Z}^{p_3}\widehat{Z}\stackrel{ˇ}{Z}^{p_4}\widehat{\phi }^{A_1B_1}\stackrel{ˇ}{Z}^{q_1}\widehat{Z}\stackrel{ˇ}{Z}^{q_2}q\widehat{\phi }^{A_2B_2}\stackrel{ˇ}{Z}^r\widehat{\phi }^{A_3B_3}\stackrel{ˇ}{Z}^s\widehat{\phi }^{A_4B_4}\right)`$
$`\mathrm{}`$
$`\mathrm{Tr}\left(\stackrel{ˇ}{Z}^p\widehat{\phi }^{A_1B_1}\stackrel{ˇ}{Z}^q\widehat{\phi }^{A_2B_2}\stackrel{ˇ}{Z}^r\widehat{\phi }^{A_3B_3}\stackrel{ˇ}{Z}^{s_1}\widehat{Z}\stackrel{ˇ}{Z}^{s_2}\widehat{Z}\stackrel{ˇ}{Z}^{s_3}\widehat{Z}\stackrel{ˇ}{Z}^{s_4}\widehat{Z}\stackrel{ˇ}{Z}^{s_5}\widehat{\phi }^{A_4B_4}\right),`$ (5.26)
where in the first trace $`_ip_i=p4=J(q+r+s+4)`$, in the second $`_ip_i=p3`$ and $`_iq_i=q1`$ and so on until the last sum where $`_is_i=s4`$. The ellipsis in (5.26) refers to other combinations in which the four $`\widehat{Z}`$’s are gradually moved to the right. All these traces can be evaluated using the ADHM solution for the scalars given in (A.8) and selecting for each factor the matrix elements containing the appropriate fermion bilinears. The calculation is rather involved. As explained in appendix B it can be carried out most efficiently defining a more general trace from which the 35 distinct traces (5.26) can be obtained for different choices of indices.
In order to compute the relevant part of the profile of $`𝒪_\mathrm{𝟏}`$ we need to sum the contributions of the traces (5.26) corresponding to the 6$`\times `$4 choices of indices, $`(A_i,B_i)`$, $`i=1,\mathrm{},4`$, on the impurities. A key feature of all these traces is that they do not depend on the way the $`Z`$’s are grouped, i.e. they do not depend on the exponents, $`(p_1,\mathrm{},p_5,q,r,s)`$, $`(p_1,\mathrm{},p_4,q_1,q_2,r,s)`$ etc. in (5.26), but only on the ordering of the four $`\widehat{Z}`$’s with respect to the four impurities, $`\widehat{\phi }^{A_iB_i}`$, $`i=1,\mathrm{},4`$. This is a consequence of the structure of the ADHM matrices and the restrictions imposed by the ADHM constraints. Keeping only the terms with two $`\zeta `$’s of each flavour all the traces in $`𝒪_\mathrm{𝟏}`$ produce expressions which after simple Fierz rearrangements can be brought to the form
$$\frac{\rho ^8}{[(x_1x_0)^2+\rho ^2]^{J+8}}\left(\overline{\nu }^{[1}\nu ^{4]}\right)^J[\left(\zeta ^1\right)^2\left(\zeta ^2\right)^2\left(\zeta ^3\right)^2\left(\zeta ^4\right)^2](x_1).$$
(5.27)
Similarly all the contributions from the traces in $`\overline{𝒪}_\mathrm{𝟏}`$ containing the required eight superconformal modes can be put in the form
$$\frac{\rho ^8}{[(x_2x_0)^2+\rho ^2]^{J+8}}\left(\overline{\nu }^{[2}\nu ^{3]}\right)^J[\left(\zeta ^1\right)^2\left(\zeta ^2\right)^2\left(\zeta ^3\right)^2\left(\zeta ^4\right)^2](x_2).$$
(5.28)
Each set of indices $`(A_i,B_i)`$, $`i=1,\mathrm{},4`$ on the impurities in each of the 35 traces leads to a contribution of the form (5.27)-(5.28) with a different numerical coefficient.
The fact that the result of all the traces can be reduced to the above expressions implies that when substituting into the definition of the operator (5.16) and its conjugate a common factor (5.27) or, respectively, (5.28) can be taken out of the traces. Associated with each of the 35 types of traces there are, however, multiplicity factors which make the sums in the definition of the operator non-trivial. For instance in the last trace in (5.26) there are $`s`$ choices for the first $`\widehat{Z}`$ among the $`Z`$’s, $`(s1)`$ choices for the second $`\widehat{Z}`$, $`(s2)`$ for the third and $`(s3)`$ for the fourth. After substituting into the definition (5.16) and factoring out the moduli dependence in the form (5.27), the contribution of the last trace in (5.26) involves the sums
$$\begin{array}{c}\underset{q,r,s=0}{\overset{J}{}}\\ q+r+sJ\end{array}\mathrm{e}^{2\pi i[(n_1+n_2+n_3)q+(n_2+n_3)r+n_3s]/J}s(s1)(s2)(s3),$$
(5.29)
with a numerical coefficient resulting from the contributions of the 6$`\times `$4 permutations of indices of the impurities. Repeating the same analysis for all the traces means combining a huge number of terms which makes the calculation extremely laborious. Completely analogous steps go into the calculation of the profile of the conjugate operator.
After lengthy algebraic manipulations and the use of the formalism described in section 3, the semi-classical result for the two-point function (5.18) takes the form
$`G_\mathrm{𝟏}(x_1,x_2)`$ $`=`$ $`{\displaystyle \frac{\mathrm{e}^{2\pi i\tau }}{J^3N^{7/2}}}{\displaystyle \frac{\mathrm{d}^4x_0\mathrm{d}\rho }{\rho ^5}\frac{\rho ^{J+8}}{[(x_1x_0)^2+\rho ^2]^{J+8}}\frac{\rho ^{J+8}}{[(x_2x_0)^2+\rho ^2]^{J+8}}}`$ (5.30)
$`\times {\displaystyle }{\displaystyle \underset{A=1}{\overset{4}{}}}\mathrm{d}^2\eta ^A\mathrm{d}^2\overline{\xi }^A{\displaystyle \underset{B=1}{\overset{4}{}}}\left[\left(\zeta ^B\right)^2(x_1)\right]\left[\left(\zeta ^B\right)^2(x_2)\right]`$
$`\times {\displaystyle }\mathrm{d}^5\mathrm{\Omega }\left(\mathrm{\Omega }^{14}\right)^J\left(\mathrm{\Omega }^{23}\right)^JK(n_1,n_2,n_3;J)K(m_1,m_2,m_3;J),`$
where following the discussion in section 3 the $`(\overline{\nu }\nu )_\mathrm{𝟔}`$ bilinears have been expressed in terms of the angular variables $`\mathrm{\Omega }^{AB}`$. In (5.30) overall numerical coefficients have been omitted. The $`J`$ and $`N`$ dependence in the prefactor in (5.30) is obtained combining the normalisation of the operators, the contribution of the measure on the instanton moduli space and the factors of $`g_{_{\mathrm{YM}}}\sqrt{N}`$ associated with $`(\overline{\nu }\nu )_\mathrm{𝟔}`$ bilinears. The origin of the various factors which determine the dependence on the parameters $`g_{_{\mathrm{YM}}}`$, $`N`$ and $`J`$ will be summarised shortly. The expression (5.30) contains integrations over the bosonic moduli, $`x_0`$ and $`\rho `$, the sixteen superconformal fermion modes and the five-sphere coordinates $`\mathrm{\Omega }^{AB}`$. The dependence on the integers $`n_i`$, $`m_i`$, $`i=1,2,3`$, dual to the mode numbers of the corresponding string state is contained in the functions $`K(n_1,n_2,n_3;J)`$ and $`K(m_1,m_2,m_3;J)`$. These are given by the sum of 35 terms,
$$K(n_1,n_2,n_3;J)=\underset{a=1}{\overset{35}{}}c_a𝒮_a(n_1,n_2,n_3;J),$$
(5.31)
where the symbols $`𝒮_a`$ indicate 35 different sums over the indices $`q,r,s`$ in which the summands are given by the phase factor $`\mathrm{exp}\{2\pi i[(n_1+n_2+n_3)q+(n_2+n_3)r+n_3s]/J\}`$ times the multiplicity factors associated with the different distributions of $`\widehat{Z}`$’s in each case. The numerical coefficients $`c_a`$ are obtained combining the contributions of the different permutations of indices on the impurities for each of the 35 terms. See appendix B.2 for more details.
In the large $`J`$ limit the leading order contribution to the sums $`𝒮_a`$ can be obtained using a continuum approximation by setting $`x=q/J`$, $`y=r/J`$, $`z=s/J`$, so that $`x,y,z[0,1]`$. For instance the sum (5.29) is approximated as
$`\begin{array}{c}{\displaystyle \underset{q,r,s=0}{\overset{J}{}}}\\ q+r+sJ\end{array}\mathrm{e}^{2\pi i[(n_1+n_2+n_3)q+(n_2+n_3)r+n_3s]/J}s(s1)(s2)(s3)`$ (5.34)
$`J^7{\displaystyle _0^1}dx{\displaystyle _0^{1x}}dy{\displaystyle _0^{1xy}}dzz^4\mathrm{e}^{2\pi i[(n_1+n_2+n_3)x+(n_2+n_3)y+n_4z]},`$ (5.35)
which shows that it behaves as $`J^7`$ for large $`J`$. We shall denote by $`\kappa (n_1,n_2,n_3)`$ the function of the mode numbers arising from these sums/integrals after extracting a factor of $`J^7`$,
$$K(n_1,n_2,n_3;J)=J^7\kappa (n_1,n_2,n_3).$$
(5.36)
The two-point function is then
$`G_\mathrm{𝟏}(x_1,x_2)`$ $`=`$ $`{\displaystyle \frac{J^{11}\mathrm{e}^{2\pi i\tau }}{N^{7/2}}}{\displaystyle \frac{\mathrm{d}^4x_0\mathrm{d}\rho }{\rho ^5}\frac{\rho ^{J+8}}{[(x_1x_0)^2+\rho ^2]^{J+8}}\frac{\rho ^{J+8}}{[(x_2x_0)^2+\rho ^2]^{J+8}}}`$ (5.37)
$`\times {\displaystyle }{\displaystyle \underset{A=1}{\overset{4}{}}}\mathrm{d}^2\eta ^A\mathrm{d}^2\overline{\xi }^A{\displaystyle \underset{B=1}{\overset{4}{}}}\left[\left(\zeta ^B\right)^2(x_1)\right]\left[\left(\zeta ^B\right)^2(x_2)\right]`$
$`\times {\displaystyle }\mathrm{d}^5\mathrm{\Omega }\left(\mathrm{\Omega }^{14}\right)^J\left(\mathrm{\Omega }^{23}\right)^J\kappa (n_1,n_2,n_3;J)\kappa (m_1,m_2,m_3;J).`$
Unlike the case of two impurity operators discussed in the previous subsection, here the dependence on the instanton moduli is known explicitly and we can compute the associated integrations. More details are given in appendix B.2. The integration over the five-sphere in (5.37) is a special case of the general integral (3.19) and gives
$$I_{S^5}=\mathrm{d}^5\mathrm{\Omega }\left(\mathrm{\Omega }^{14}\right)^J\left(\mathrm{\Omega }^{23}\right)^J=\frac{\pi ^3}{(J+1)(J+2)}.$$
(5.38)
The integration over the superconformal modes is also straightforward. It does not depend on $`N`$ or $`J`$. For each flavour the result is
$$I_\zeta =\mathrm{d}^2\eta \mathrm{d}^2\overline{\xi }\left[\left(\zeta \right)^2(x_1)\right]\left[\left(\zeta \right)^2(x_2)\right]=(x_1x_2)^2,$$
(5.39)
so that the fermionic integrals contribute a factor of $`(x_1x_2)^8`$. The integration over the bosonic part of the moduli space must be treated carefully since it is logarithmically divergent as expected in the presence of a contribution to the matrix of anomalous dimensions. The integrals need to be regulated for instance by dimensional regularisation of the $`x_0`$ integral and can then be computed using standard techniques, e.g. introducing Feynman parameters. The result is
$`I_\mathrm{b}`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}^4x_0\mathrm{d}\rho }{\rho ^5}\frac{\rho ^{J+8}}{[(x_1x_0)^2+\rho ^2]^{J+8}}\frac{\rho ^{J+8}}{[(x_2x_0)^2+\rho ^2]^{J+8}}}`$ (5.40)
$`=`$ $`\text{}{\displaystyle \frac{1}{ϵ}}{\displaystyle \frac{\mathrm{\Gamma }(J+6)\mathrm{\Gamma }(J+8+ϵ)}{[\mathrm{\Gamma }(J+8)]^2}}\pi ^{2ϵ}{\displaystyle \frac{1}{(x_{12}^2)^{J+8+ϵ}}},ϵ0.`$
The $`1/ϵ`$ pole is the manifestation of a logarithmic divergence in dimensional regularisation. The contribution (5.40) behaves as $`1/J^2`$ in the large $`J`$ limit.
Putting together all the contributions the dependence on the parameters, $`g_{_{\mathrm{YM}}}`$, $`N`$ and $`J`$, in the correlation function can be summarised as follows
$`\underset{\mathrm{normalised}\mathrm{op}.\mathrm{profile}}{\underset{}{\left({\displaystyle \frac{1}{\sqrt{J^3(g_{_{\mathrm{YM}}}^2N)^{J+4}}}}\right)^2}}\underset{\nu ,\overline{\nu }\mathrm{bilinears}}{\underset{}{\left(g_{_{\mathrm{YM}}}\sqrt{N}\right)^{2J}}}\underset{\mathrm{measure}}{\underset{}{\mathrm{e}^{2\pi i\tau }g_{_{\mathrm{YM}}}^8\sqrt{N}}}\underset{S^5\mathrm{integral}}{\underset{}{{\displaystyle \frac{1}{J^2}}}}\underset{x_0,\rho \mathrm{integrals}}{\underset{}{{\displaystyle \frac{1}{J^2}}}}\underset{𝒮_a\mathrm{sums}}{\underset{}{(J^7)^2}}`$
$`{\displaystyle \frac{J^7}{N^{7/2}}}\mathrm{e}^{2\pi i\tau }=(g_2)^{7/2}\mathrm{e}^{\frac{8\pi ^2}{g_2\lambda ^{}}+i\theta }.`$ (5.41)
The final result for the two-point function is thus, up to a numerical coefficient,
$$G_\mathrm{𝟏}(x_1,x_2)=(g_2)^{7/2}\mathrm{e}^{\frac{8\pi ^2}{g_2\lambda ^{}}+i\theta }\kappa (n_1,n_2,n_3)\kappa (m_1,m_2,m_3)\frac{1}{(x_{12}^2)^{J+4}}\mathrm{log}\left(\mathrm{\Lambda }^2x_{12}^2\right).$$
(5.42)
where the scale $`\mathrm{\Lambda }`$ appears as a consequence of the $`1/ϵ`$ divergence. It has no observable effect. The physical information contained in the two-point function is in the contribution to the matrix of anomalous dimensions which is read from the coefficient in (5.42) and does not depend on $`\mathrm{\Lambda }`$.
The result is expressed in terms of the double scaling parameters $`\lambda ^{}`$ and $`g_2`$. Note that, unlike the two-point functions of two impurity operators (5.42) is independent of $`\lambda ^{}`$ apart from the dependence in the exponential instanton weight.
The non-perturbative mass correction computed in for the state (5.17), in terms of the same gauge theory parameters, is of the form
$$\delta m(g_2)^{7/2}\mathrm{e}^{\frac{8\pi ^2}{g_2\lambda ^{}}+i\theta }\frac{1}{(n_1n_2)^2}.$$
(5.43)
so that the $`\lambda ^{}`$ and $`g_2`$ dependence is in agreement with the gauge theory calculation. The mode number dependence in (5.43) is remarkably simple and a special feature of the string result, which is a direct consequence of the structure of the $`D`$-instanton boundary state, is that it is non-vanishing only if the mode numbers in both the incoming and the outgoing states are pairwise equal. The only states which couple to the $`D`$-instanton are of the form
$$\epsilon _{ijkl}\alpha _{n_1}^i\alpha _{n_2}^j\stackrel{~}{\alpha }_{n_1}^k\stackrel{~}{\alpha }_{n_2}^l|0_h.$$
(5.44)
On the other hand in the gauge theory result (5.42) obtained for the operator (5.16) the mode number dependence is contained in $`\kappa (n_1,n_2,n_3)`$ and $`\kappa (m_1,m_2,m_3)`$, which are extremely complicated rational functions of their arguments. In particular the condition that the integers $`n_i`$ be equal in pairs does not seem to be required.
However, as observed after (5.17) in order to correctly match the properties of the dual string state, the operator (5.16) must be explicitly antisymmetrised under the exchange of pairs of mode numbers. This antisymmetrisation induces dramatic simplifications. Working with the correctly antisymmetrised operators the result for the two-point function is
$$G_1(x_1,x_2)=\frac{3^2(g_2)^{7/2}\mathrm{e}^{\frac{8\pi ^2}{g_2\lambda ^{}}+i\theta }}{2^{41}\pi ^{9/2}}\frac{1}{(n_1n_2)(m_1m_2)}\frac{1}{(x_{12}^2)^{J+4}}\mathrm{log}\left(\mathrm{\Lambda }^2x_{12}^2\right)$$
(5.45)
if the mode numbers in each of the two operators are equal in pairs and vanishes otherwise. In (5.45) we have reinstated all the numerical coefficients coming from the profiles of the operators, the combinatorics previous described and the moduli space measure. This result is in perfect agreement with the string result (5.43) of . It is worth stressing that the simplification found after the antisymmetrisation is extraordinary given the complexity of the function $`\kappa (n_1,n_2,n_3)`$. Moreover the condition of pairwise equal mode numbers which is also imposed in this way is far from obvious and highly non-trivial from the point of view of the gauge theory calculation.
As we have seen, in the two-point function computed in the semi-classical approximation the mode number dependence factorises. A consequence of this is that in $`G_\mathrm{𝟏}(x_1,x_2)`$ the two independent mode numbers in $`𝒪_\mathrm{𝟏}`$, $`n_1`$ and $`n_2`$, do not have to equal those in $`\overline{𝒪}_\mathrm{𝟏}`$, $`m_1`$ and $`m_2`$. This appears to contradict energy conservation in the dual string amplitude, which requires the mode numbers of the outgoing state to match one to one those of the incoming state. However, the fact that the condition $`m_i=n_i`$, $`i=1,2`$ does not arise is an effect of the semi-classical approximation. This is valid in the $`\lambda ^{}0`$ limit which corresponds to the $`m\mathrm{}`$ limit in the plane-wave string theory (where $`m`$ is the mass parameter entering the string action). In this strict limit energy conservation in a two-point string amplitude only requires that the number of oscillators in the incoming and outgoing states be equal, with no constraint on the associated mode numbers. Therefore (5.45) is indeed in agreement with the string theory result. On the other hand the instanton corrections discussed here should be considered as subleading corrections on top of the perturbative effects. The condition $`m_i=n_i`$ on the operators in a two-point function is already imposed at the perturbative level and should therefore be assumed when computing instanton contributions in the semi-classical approximation.
The calculation presented here is not sufficient to determine the actual instanton induced anomalous dimension of the operator $`𝒪_\mathrm{𝟏}`$. This requires the diagonalisation of the matrix of anomalous dimensions of which we have not computed all the entries. Other entries are determined by the corresponding two-point functions whose calculation follows the same steps described here and results in expressions similar to (5.45). From this we can conclude that the behaviour of the leading instanton contribution to the anomalous dimensions of singlet operators is
$$\gamma _\mathrm{𝟏}^{\mathrm{inst}}(g_2)^{7/2}\mathrm{e}^{\frac{8\pi ^2}{g_2\lambda ^{}}+i\theta }\frac{1}{(n_1n_2)^2}.$$
(5.46)
As a further test of the result we find that the two point function vanishes in the limit of zero mode numbers. The function $`\kappa (n_1,n_2,n_3)`$ is identically zero when $`n_1=n_2=n_3=0`$. This is the expected behaviour because in this limit the operator is expected to become protected and no corrections to its free theory two-point functions should arise. The string theory counterpart of this result is the decoupling of the supergravity modes (dual to the protected operators with $`\{n_i=0\}`$), which was also verified in .
In the previous analysis we have considered only a class of contributions in which in each operator as many superconformal modes as possible were taken from the $`Z`$’s. It is easy to verify that these are the only relevant terms at leading order in the BMN limit. All the other types of traces are suppressed and vanish in the $`J\mathrm{}`$ limit. As an example consider a contribution to the profile of $`𝒪_\mathrm{𝟏}`$ in which the $`\zeta ^2`$ and $`\zeta ^3`$ modes as well as one of either the $`\zeta ^1`$ or $`\zeta ^4`$ modes are taken from the impurities. Instead of the last trace in (5.26) we would then consider traces of the type
$$\mathrm{Tr}\left(\stackrel{ˇ}{Z}^p\stackrel{~}{\phi }^{A_1B_1}\stackrel{ˇ}{Z}^q\widehat{\phi }^{A_2B_2}\stackrel{ˇ}{Z}^r\widehat{\phi }^{A_3B_3}\stackrel{ˇ}{Z}^{s_1}\widehat{Z}\stackrel{ˇ}{Z}^{s_2}\widehat{Z}\stackrel{ˇ}{Z}^{s_3}\widehat{Z}\stackrel{ˇ}{Z}^{s_4}\widehat{\phi }^{A_4B_4}\right),$$
(5.47)
where the first impurity contains two superconformal modes and thus only three $`\widehat{Z}`$’s are needed. An analysis similar to that carried out for the traces (5.26) can be repeated in this case and one finds that associated with such a trace there is sum of the form
$$\begin{array}{c}\underset{q,r,s=0}{\overset{J}{}}\\ q+r+sJ\end{array}\mathrm{e}^{2\pi i[(n_1+n_2+n_3)q+(n_2+n_3)r+n_3s]/J}\frac{1}{3!}s(s1)(s2),$$
(5.48)
which behaves as $`J^6`$ in the large $`J`$ limit. The remaining $`J`$, $`N`$ and $`g_{_{\mathrm{YM}}}`$ dependence in the correlation function is unmodified and thus the combined behaviour of this contribution can be read from (5.41) replacing the last factor on the first line by $`(J^6)^2`$ leading to
$$\mathrm{e}^{2\pi i\tau }\frac{J^5}{N^{7/2}}\frac{\mathrm{e}^{2\pi i\tau }(g_2)^{7/2}}{J^2},$$
which vanishes in the BMN limit. Similar arguments can be repeated for all the contributions other than those leading to (5.45), which is therefore the complete leading instanton contribution to this singlet two-point function in the BMN limit. We shall briefly comment on corrections to this result of higher order in $`\lambda ^{}`$ and $`g_2`$ in the discussion section.
#### 5.2.2 Other four impurity singlets
There are two other independent four impurity singlet operators involving four scalar impurities. They correspond to the two inequivalent ways of contracting the SO(4)<sub>R</sub> indices with Kronecker delta’s,
$`𝒪_{\mathrm{𝟏};J;n_1,n_2,n_3}^{(d_1)}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{J^3\left(\frac{g_{_{\mathrm{YM}}}^2N}{8\pi ^2}\right)^{J+4}}}}\begin{array}{c}{\displaystyle \underset{q,r,s=0}{\overset{J}{}}}\\ q+r+sJ\end{array}\mathrm{e}^{2\pi i[(n_1+n_2+n_3)q+(n_2+n_3)r+n_3s]/J}`$ (5.52)
$`\times \mathrm{Tr}\left(Z^{J(q+r+s)}\phi ^iZ^q\phi ^iZ^r\phi ^jZ^s\phi ^j\right),`$
$`𝒪_{\mathrm{𝟏};J;n_1,n_2,n_3}^{(d_2)}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{J^3\left(\frac{g_{_{\mathrm{YM}}}^2N}{8\pi ^2}\right)^{J+4}}}}\begin{array}{c}{\displaystyle \underset{q,r,s=0}{\overset{J}{}}}\\ q+r+sJ\end{array}\mathrm{e}^{2\pi i[(n_1+n_2+n_3)q+(n_2+n_3)r+n_3s]/J}`$ (5.56)
$`\times \mathrm{Tr}\left(Z^{J(q+r+s)}\phi ^iZ^q\phi ^jZ^r\phi ^iZ^s\phi ^j\right).`$
The calculation of instanton contributions to the two-point functions of these operators proceeds in complete analogy with the discussion in the previous subsection. These operators are expected to receive contributions of the same type as the $`\epsilon `$-singlet (5.16) and to mix with the latter much in the same way as was found in the string theory analysis of .
Other singlet operators can be constructed using all the other combinations of impurities in table 2. In all these cases the analysis of fermion zero modes shows that a non-zero contribution to the corresponding two-point functions can arise at the same leading order as (5.46). Operators containing $`D_\mu Z`$ insertions correspond to string states involving bosonic oscillators which are vectors of the second SO(4). Operators containing the $`\psi _a^a`$ and $`\psi _{\dot{a}}^{+\dot{\alpha }}`$ fermions are dual to states created by the $`S^\pm `$ oscillators. The calculation of two-point functions of all these operators is similar to that described in the previous section with the additional technical complication that in the presence of covariant derivatives the solution $`A_\mu ^{(4)}`$ for the gauge potential is needed and for operators containing fermions the solution $`\lambda _\alpha ^{(5)A}`$ is needed.
As observed in the case of the operator (5.16), two-point functions in the semi-classical approximation factorise, with the two operators being related only by the five-sphere integration. Because of this property mixing is expected among all the operators which receive instanton contributions.
In the $`𝒩`$=4 theory it is in principle possible to construct a large number of other operators which potentially mix with those considered here, being SO(4)$`{}_{R}{}^{}\times `$SO(4)<sub>C</sub> singlets with $`\mathrm{\Delta }J=4`$. These involve $`\mathrm{\Delta }J=2`$ impurities and thus do not correspond to new independent states having vanishing two-point functions in free theory. However, it is known that the inclusion of such operators is needed in perturbation theory to properly resolve the mixing beyond the zeroth order approximation in the $`g_2`$ expansion. Since instanton effects are exponentially suppressed in $`g_2`$ one should in principle expect these operators to be relevant at leading order in the instanton background. This is, however, not the case. The combinatorial analysis involved in computing the classical profiles of the operators shows that those containing $`\mathrm{\Delta }J=2`$ impurities are suppressed in the large $`J`$ limit.
#### 5.2.3 Operators in other sectors
As observed in section 2.2.2 the spectrum of four impurity BMN operators is rather rich. Instanton contributions to the anomalous dimensions of operators in other sectors can be studied with the same methods used for the singlets. $`D`$-instanton induced amplitudes for string states in the plane wave background dual to non-singlet operators are suppressed with respect to those in the singlet sector. Hence string theory predicts that the leading instanton contributions to the anomalous dimensions of non-singlet four impurity operators should be suppressed with respect to (5.46). More precisely the string prediction is that the leading non-zero contributions should arise at order $`\mathrm{e}^{2\pi i\tau }(g_2)^{7/2}(\lambda ^{})^2`$. We shall not discuss in detail the calculation of two-point functions needed to verify this prediction, but we present here an argument indicating that the gauge theory result is indeed in agreement with string theory. We focus on an operator with four scalar impurities which is a singlet of SO(4)<sub>C</sub> and belongs to the $`\mathrm{𝟑}^+\mathrm{𝟑}^{}`$ of SO(4)<sub>R</sub>,
$`𝒪_{\mathrm{𝟑}^+\mathrm{𝟑}^{};J;n_1,n_2,n_3}^{[ij]}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{J^3\left(\frac{g_{_{\mathrm{YM}}}^2N}{8\pi ^2}\right)^{J+4}}}}\begin{array}{c}{\displaystyle \underset{q,r,s=0}{\overset{J}{}}}\\ q+r+sJ\end{array}\mathrm{e}^{2\pi i[(n_1+n_2+n_3)q+(n_2+n_3)r+n_3s]/J}`$ (5.60)
$`\times \mathrm{Tr}\left(Z^{J(q+r+s)}\phi ^kZ^q\phi ^kZ^r\phi ^{[i}Z^s\phi ^{j]}\right).`$
The study of other non-singlet operators is completely analogous. Considering for concreteness the component $`i=1`$, $`j=2`$ in (5.60), we find that using for all the scalars the bilinear solution the combinations of fermion modes contained in the classical profiles of the operator and its conjugate are respectively
$`\left(𝓂_\mathrm{f}^1\right)^{J+3}\left(𝓂_\mathrm{f}^2\right)^2\left(𝓂_\mathrm{f}^3\right)^2\left(𝓂_\mathrm{f}^4\right)^{J+1}`$ (5.61)
$`+`$ $`\left(𝓂_\mathrm{f}^1\right)^{J+2}\left(𝓂_\mathrm{f}^2\right)^3\left(𝓂_\mathrm{f}^3\right)\left(𝓂_\mathrm{f}^4\right)^{J+2}`$ (5.62)
$`+`$ $`\left(𝓂_\mathrm{f}^1\right)^{J+2}\left(𝓂_\mathrm{f}^2\right)\left(𝓂_\mathrm{f}^3\right)^3\left(𝓂_\mathrm{f}^4\right)^{J+2}`$ (5.63)
$`+`$ $`\left(𝓂_\mathrm{f}^1\right)^{J+1}\left(𝓂_\mathrm{f}^2\right)^2\left(𝓂_\mathrm{f}^3\right)^2\left(𝓂_\mathrm{f}^4\right)^{J+3}`$ (5.64)
and
$`\left(𝓂_\mathrm{f}^1\right)^3\left(𝓂_\mathrm{f}^2\right)^{J+2}\left(𝓂_\mathrm{f}^3\right)^{J+2}\left(𝓂_\mathrm{f}^4\right)`$ (5.65)
$`+`$ $`\left(𝓂_\mathrm{f}^1\right)^2\left(𝓂_\mathrm{f}^2\right)^{J+3}\left(𝓂_\mathrm{f}^3\right)^{J+1}\left(𝓂_\mathrm{f}^4\right)^2`$ (5.66)
$`+`$ $`\left(𝓂_\mathrm{f}^1\right)^2\left(𝓂_\mathrm{f}^2\right)^{J+1}\left(𝓂_\mathrm{f}^3\right)^{J+3}\left(𝓂_\mathrm{f}^4\right)^2`$ (5.67)
$`+`$ $`\left(𝓂_\mathrm{f}^1\right)\left(𝓂_\mathrm{f}^2\right)^{J+2}\left(𝓂_\mathrm{f}^3\right)^{J+2}\left(𝓂_\mathrm{f}^4\right)^3.`$ (5.68)
This shows that the integrations over the superconformal modes in the two-point function can be saturated, selecting terms containing (5.61) or (5.64) in $`𝒪^{[12]}`$ and terms containing (5.66) and (5.67) in $`\overline{𝒪}^{[12]}`$. However, with these choices the resulting five-sphere integrals vanish. For instance combining (5.61) and (5.66) leads to the following moduli space integrals
$`{\displaystyle \mathrm{d}^8\eta \mathrm{d}^8\overline{\xi }[\left(\zeta ^1\right)^2\left(\zeta ^2\right)^2\left(\zeta ^3\right)^2\left(\zeta ^4\right)^2](x_1)[\left(\zeta ^1\right)^2\left(\zeta ^2\right)^2\left(\zeta ^3\right)^2\left(\zeta ^4\right)^2](x_2)}`$ (5.69)
$`\times `$ $`{\displaystyle \mathrm{d}^{4(N2)}\nu \mathrm{d}^{4(N2)}\overline{\nu }\left(\overline{\nu }^{[1}\nu ^{4]}\right)^{J1}\left(\overline{\nu }^{(1}\nu ^{1)}\right)\left(\overline{\nu }^{[2}\nu ^{3]}\right)^{J1}\left(\overline{\nu }^{(2}\nu ^{2)}\right)}.`$
The integration over the five-sphere arising from the second line of (5.69) vanishes because the multiplicity of the flavours 1 and 2 exceeds that of the flavours 3 and 4.
In order to soak up the superconformal modes while avoiding the obstruction from the five sphere integral it is necessary to include a six-fermion term in each operator. In this way the combinations of modes in the two operators are the same as in (5.61)-(5.68) with the addition of one mode of each flavour. The same arguments given in section 5.1 in connection with two impurity operators can be repeated here and for instance combining (5.61) and (5.68) we get moduli space integrations of the type
$`{\displaystyle \mathrm{d}^8\eta \mathrm{d}^8\overline{\xi }[\left(\zeta ^1\right)^2\zeta ^2\left(\zeta ^3\right)^2\left(\zeta ^4\right)^2\overline{\xi }^1](x_1)[\zeta ^1\left(\zeta ^2\right)^2\left(\zeta ^3\right)^2\left(\zeta ^4\right)^2\overline{\xi }^2](x_2)}`$ (5.70)
$`\times `$ $`{\displaystyle \mathrm{d}^{4(N2)}\nu \mathrm{d}^{4(N2)}\overline{\nu }\left(\overline{\nu }^{[1}\nu ^{4]}\right)^{J+1}\left(\overline{\nu }^{[2}\nu ^{3]}\right)^{J+1}\left(\overline{\nu }^{[1}\nu ^{2]}\right)\left(\overline{\nu }^{[3}\nu ^{4]}\right)}.`$
Just as in the two impurity case the resulting non-vanishing contribution to the two-point function is suppressed by a factor of $`(\lambda ^{})^2`$ due to the additional $`(\overline{\nu }\nu )_\mathrm{𝟔}`$ bilinears in (5.70). In conclusion the analysis of fermion zero modes confirms that the leading non-zero instanton contribution to the anomalous dimensions of four impurity operators in the $`\mathrm{𝟑}^+`$ and $`\mathrm{𝟑}^{}`$ representations behaves as
$$\gamma _{\mathrm{𝟑}^+\mathrm{𝟑}^{}}^{\mathrm{inst}}\mathrm{e}^{\frac{8\pi ^2}{g_2\lambda ^{}}+i\theta }(g_2)^{7/2}(\lambda ^{})^2,$$
(5.71)
in agreement with the string prediction of .
Other sectors can be analysed in a similar fashion and we find that all non-singlet operators receive leading non-zero contributions of the same order as (5.71).
## 6 Discussion and conclusions
This paper has considered one-instanton contributions to two-point correlation functions of gauge invariant operators in the BMN sector of the $`𝒩`$=4 supersymmetric Yang–Mills theory. These determine the leading instanton contributions to the anomalous dimensions of operators dual to physical string states in the maximally supersymmetric plane-wave background obtained as Penrose limit of AdS$`{}_{5}{}^{}\times S^5`$. The basic message is that we find striking agreement between these instanton effects in the gauge theory and those of the plane-wave string theory calculated in .
We focused on operators with two and four scalar impurities. The four impurity case, although more involved, is fully under control, whereas the two impurity case presents subtleties due to the fact the leading semi-classical approximation vanishes and the first non-zero contribution arises at higher order. We have explicitly computed a two-point function of four impurity operators which are SO(4)$`{}_{R}{}^{}\times `$SO(4)<sub>C</sub> singlets. Our analysis shows that instanton induced contributions to the anomalous dimensions of operators in this sector behave as $`1/(n_1n_2)^2\mathrm{exp}\left(8\pi ^2/g_2\lambda ^{}+i\theta \right)g_2^{7/2}`$, where $`\lambda ^{}`$ and $`g_2`$ are the effective coupling constant and genus counting parameter in the BMN limit and $`n_1`$ and $`n_2`$ correspond to the mode numbers of the dual string state. The result is in perfect agreement with the $`D`$-instanton correction to the mass matrix elements of the corresponding states in the plane-wave string theory which was computed in . Even without directly matching the numerical values of the anomalous dimensions and the string mass renormalisation, the agreement with the string calculation appears highly non-trivial. The correct dependence on the parameters $`\lambda ^{}`$ and $`g_2`$ is obtained by combining contributions arising from the integrations over the instanton moduli space and various combinatorial factors. Even more impressively, the mode number dependence found in is reproduced after spectacular cancellations.
In the case of two impurity operators the leading instanton correction vanishes. The first non-zero contribution is awkward to calculate completely, but with mild assumptions we showed that it has the form $`\mathrm{exp}\left(8\pi ^2/g_2\lambda ^{}+i\theta \right)g_2^{7/2}\lambda _{}^{}{}_{}{}^{2}`$ and does not depend on the single mode number characterising the dual string state. This behaviour is also in agreement with the results of , although the subtleties presented by the gauge theory calculation did not arise on the string side. Four impurity operators in sectors other than the singlet have the same $`\lambda ^{}`$ and $`g_2`$ dependence as two impurity operators, again in agreement with the string prediction of .
Our results provide a significant new test of the duality proposed in . The fact that non-perturbative contributions obey BMN scaling, i.e. can be re-expressed in terms of the effective parameters $`\lambda ^{}`$ and $`g_2`$, strongly supports the conjecture that this property should hold at all orders. This can be further tested by analysing subleading effects in the instanton background. A class of higher order contributions can easily be obtained from the calculations presented in this paper, relaxing the requirement that all the fermion modes of type $`\nu `$ and $`\overline{\nu }`$ be combined in $`(\overline{\nu }\nu )_\mathrm{𝟔}`$ bilinears. As already observed, replacing a $`(\overline{\nu }\nu )_\mathrm{𝟔}`$ bilinear with a $`(\overline{\nu }\nu )_{\mathrm{𝟏𝟎}}`$ leads to a suppression by a factor of $`1/\sqrt{N}`$, see (3.17). The profile of BMN operators of the type that we considered contains $`J`$ $`(\overline{\nu }\nu )`$’s after the superconformal modes have been soaked up. Hence each factor of $`1/\sqrt{N}`$ coming from the replacement of a $`(\overline{\nu }\nu )_\mathrm{𝟔}`$ by a $`(\overline{\nu }\nu )_{\mathrm{𝟏𝟎}}`$ is associated with a factor of $`J`$ corresponding to the number of choices for the $`(\overline{\nu }\nu )_{\mathrm{𝟏𝟎}}`$ bilinear, resulting in a suppression by $`g_2^{1/2}`$. Moreover in order to get a non-zero result from the five sphere integration an even number of $`(\overline{\nu }\nu )_{\mathrm{𝟏𝟎}}`$ is required. Therefore contributions of this type with an increasing number of $`(\overline{\nu }\nu )_{\mathrm{𝟏𝟎}}`$ insertions give rise to subleading corrections which form a series in integer powers of $`g_2`$. More complicated terms with the same behaviour correspond to contributions in which pairs of fields are contracted between the two-operators. We can also identify a class of subleading corrections suppressed by powers of $`\lambda ^{}`$. These are generated by including in the profile of the operators a number of fermion modes greater than the minimal number required by the moduli space integration. For instance in the case of the four impurity singlets that we studied in section 5.2.1 this is achieved by including one six-fermion scalar in one of the two operators. In this case the calculation is analogous to that of section 5.1 for the two impurity case and we can argue that the resulting contribution should be suppressed by a factor of $`\lambda ^{}`$. Including more six-fermion terms gives rise to higher powers of $`\lambda ^{}`$. Although these arguments are rather qualitative they indicate how the perturbative series of corrections to the semi-classical one-instanton contributions can be reorganised into a double series in $`\lambda ^{}`$ and $`g_2`$.
Some rather striking general properties of the $`D`$-instanton induced corrections to the string mass spectrum of can be immediately deduced from the structure of the $`D`$-instanton boundary state in the plane-wave background, whereas the corresponding effects in the gauge theory are far from obvious. In fact, the string theory results suggest a number of extensions and generalisations of the gauge theory results. For example, one generic feature of the string calculation is that only states with an even number of non-zero mode insertions receive $`D`$-instanton corrections. Zero mode oscillators can appear in odd numbers with the condition that they be contracted into a SO(4)$`\times `$SO(4) scalar between the incoming and outgoing states. The simplest example in which these properties can be verified involves five impurity operators and the calculation of the necessary two-point function in the gauge theory is extremely complicated <sup>5</sup><sup>5</sup>5Three impurity operators present technical difficulties similar to those encountered in the two impurity case.. In general, contributions to operators with a larger (even) number of impurities are expected to be non-zero at leading order in the instanton background. However, the complexity of the combinatorics involved in such calculations grows rapidly with the number of impurities.
Another peculiarity observed in the string theory calculation is that the $`D`$-instanton contribution to the masses of certain states with a large number of fermionic non-zero mode excitations involves large powers of the mass parameter $`m`$. When expressed in terms of gauge theory parameters this corresponds to large inverse powers of $`\lambda ^{}`$. As observed in the behaviour of these mass corrections is not pathological in the $`\lambda ^{}0`$ limit, because the inverse powers of $`\lambda ^{}`$ are accompanied by the instanton factor $`\mathrm{exp}\left(8\pi ^2/g_2\lambda ^{}\right)`$. From the point of view of the gauge theory this result is particularly intriguing not only because of the unusual coupling constant dependence that the anomalous dimensions of the dual operators should display, but also because there are no other known examples of operators in $`𝒩`$=4 SYM whose anomalous dimension receives instanton but not perturbative corrections. We will study this particular class of BMN operators in a future publication .
In the original formulation of the AdS/CFT duality, relating $`𝒩`$=4 SYM to type IIB string theory in AdS$`{}_{5}{}^{}\times S^5`$, the effects of multi-instantons in the large-$`N`$ limit of the gauge theory and of multi $`D`$-instantons in string theory were shown to be in remarkable agreement . Clearly it would be of interest to generalise the present work from the one-instanton sector to the multi-instanton sector. However, such a generalisation is technically very challenging both on the string and on the gauge side.
Instanton effects have been studied in a number of different supersymmetric field theories in the context of the AdS/CFT correspondence, and agreement has been found between string and gauge theory. The example of the $`𝒩`$=2 Sp($`N`$) superconformal field theory studied in is particularly interesting in connection with our work because in this case the analogue of the BMN limit has been studied in . In this case the duality involves a theory of open and closed strings in a plane-wave background and the dual gauge theory has a rich spectrum of gauge-invariant operators and possesses a Higgs branch . It would be interesting to study instanton effects in the BMN sector of this theory.
In a conformal field theory, the problem of computing the scaling dimensions of gauge invariant operators can be reformulated as an eigenvalue problem for the dilation operator of the theory. At the perturbative level this observation leads to a very efficient approach to the calculation of anomalous dimensions in $`𝒩`$=4 SYM . Some comments about the possibility of extending this approach to non-perturbative sectors were made in , but there has been no further progress in this direction. A remarkable consequence of recasting the problem of computing anomalous dimensions as an eigenvalue problem for the dilation operator is the emergence of connections with integrable systems. In the planar limit the dilation operator can be related to the hamiltonian of an integrable spin chain, leading to the possibility of applying techniques such as the Bethe ansatz to the computation of anomalous dimensions , see for a review and references. The integrability structure observed in the $`𝒩`$=4 Yang–Mills theory appears, however, to be spoiled by the inclusion of non-planar contributions. Therefore instanton effects, which are exponentially suppressed in the large-$`N`$ limit, are unlikely to be relevant in connection with integrability.
Instantons play a special rôle in the $`𝒩`$=4 theory in connection with the SL(2,$``$) $`S`$-duality symmetry, which transforms the complex coupling so that $`\tau \frac{a\tau +b}{c\tau +d}`$ (with $`a,b,c,d`$ satisfying $`adbc=1`$) thereby mixing perturbative and non-perturbative effects. In the conformal phase invariance of the theory requires that the full spectrum of scaling dimensions be invariant. The anomalous dimensions are therefore naturally expressed as functions $`\gamma (\tau ,\overline{\tau })`$. Similarly $`D`$-instantons are instrumental in the implementation of $`S`$-duality in type IIB string theory. Their rôle is well understood at the level of the effective action for the supergravity states, but little is known at the level of the massive string excitations. As in the SYM case, invariance of the theory requires that the complete spectrum be invariant. In general SL(2,$``$) transformations relate operators of small and large dimension, just as in string theory they relate fundamental strings to $`D`$-strings, which have large masses of order $`1/g_s`$, in the limit of weak string coupling, $`g_s1`$. It would be interesting to understand how $`S`$-duality is realised in type IIB string theory in the plane-wave background. A corresponding symmetry should exist in the BMN sector of $`𝒩`$=4 SYM and the instanton effects which we studied in this paper should be important in its implementation.
###### Acknowledgments.
AS acknowledges financial support from PPARC and Gonville and Caius college, Cambridge. We also wish to acknowledge support from the European Union Marie Curie Superstrings Network MRTN-CT-2004-512194.
## Appendix A Useful formulae
This appendix contains some definitions and formulae used in the paper. The Clebsch–Gordan coefficients $`\mathrm{\Sigma }_{AB}^i`$ ($`\overline{\mathrm{\Sigma }}_i^{AB}`$) projecting the product of two $`\mathrm{𝟒}`$’s ($`\overline{\mathrm{𝟒}}`$’s) onto the $`\mathrm{𝟔}`$ are defined as
$`\mathrm{\Sigma }_{AB}^i=(\mathrm{\Sigma }_{AB}^a,\mathrm{\Sigma }_{AB}^{a+3})=(\eta _{AB}^a,i\overline{\eta }_{AB}^a)`$
$`\overline{\mathrm{\Sigma }}_i^{AB}=(\overline{\mathrm{\Sigma }}_{AB}^a,\overline{\mathrm{\Sigma }}_{AB}^{a+3})=(\eta _a^{AB},i\overline{\eta }_a^{AB}),\text{}`$ (A.1)
where $`a=1,2,3`$ and the ’t Hooft symbols $`\eta _{AB}^a`$ and $`\overline{\eta }_{AB}^a`$ are
$`\eta _{AB}^a=\overline{\eta }_{AB}^a=\epsilon _{aAB},A,B=1,2,3,`$
$`\eta _{A4}^a=\overline{\eta }_{4A}^a=\delta _A^a,`$
$`\eta _{AB}^a=\eta _{BA}^a,\overline{\eta }_{AB}^a=\overline{\eta }_{BA}^a.`$ (A.2)
In some situations the $`𝒩`$=1 formulation proves very useful. The $`𝒩`$=1 decomposition of the $`𝒩`$=4 supermultiplet consists of three chiral multiplets and one vector multiplet and under this decomposition only a SU(3)$`\times `$U(1) subgroup of the SU(4) R-symmetry group is manifest. The six scalars are combined into three complex fields, $`\varphi ^I`$, $`I=1,2,3`$, according to
$`\varphi ^I`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(\widehat{\phi }^I+i\widehat{\phi }^{I+3}\right)`$
$`\varphi _I^{}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(\widehat{\phi }^Ii\widehat{\phi }^{I+3}\right).`$ (A.3)
The complex scalars $`\varphi ^I`$ and $`\varphi _I^{}`$ transform respectively in the $`\mathrm{𝟑}_1`$ and $`\overline{\mathrm{𝟑}}_1`$ of SU(3)$`\times `$U(1) (where the subscript indicates the U(1) charge). The fermions in the chiral multiplets are
$$\psi _\alpha ^I=\lambda _\alpha ^I,\overline{\psi }_I^{\dot{\alpha }}=\overline{\lambda }_I^{\dot{\alpha }},I=1,2,3,$$
(A.4)
transforming in the $`\mathrm{𝟑}_{3/2}`$ and $`\overline{\mathrm{𝟑}}_{3/2}`$. The fourth fermion and the vector form the $`𝒩`$=1 vector multiplet, $`\{\lambda _\alpha =\lambda _\alpha ^4,A_\mu \}`$, and are SU(3)$`\times `$U(1) singlets.
Using these definitions we find the following relations among the scalars in the different formulations
$$\begin{array}{ccc}\widehat{\phi }^1=\frac{\sqrt{2}}{2}\left(\phi ^{14}+\phi ^{23}\right),\hfill & \widehat{\phi }^2=\frac{\sqrt{2}}{2}\left(\phi ^{13}+\phi ^{24}\right),\hfill & \widehat{\phi }^3=\frac{\sqrt{2}}{2}\left(\phi ^{12}+\phi ^{34}\right),\hfill \\ \widehat{\phi }^4=\frac{i\sqrt{2}}{2}\left(\phi ^{14}+\phi ^{23}\right),\hfill & \widehat{\phi }^5=\frac{i\sqrt{2}}{2}\left(\phi ^{13}\phi ^{24}\right),\hfill & \widehat{\phi }^6=\frac{i\sqrt{2}}{2}\left(\phi ^{12}\phi ^{34}\right)\text{}\hfill \end{array}$$
(A.5)
and
$$\begin{array}{ccc}\varphi ^1=2\phi ^{14},\hfill & \varphi ^2=2\phi ^{24},\hfill & \varphi ^3=2\phi ^{34},\hfill \\ \varphi _1^{}=2\phi ^{23},\hfill & \varphi _2^{}=2\phi ^{13},\hfill & \varphi _3^{}=2\phi ^{12}.\text{}\hfill \end{array}$$
(A.6)
In the ADHM formalism the expressions for the $`𝒩`$=4 elementary fields in the background of an instanton are conveniently given as $`[N+2]\times [N+2]`$ matrices. In particular, the two-fermion solution for the scalar field $`\phi ^{AB}`$ in the one instanton sector can be written in the block-form
$$(\widehat{\phi }^{AB})_{u\beta ;}{}_{}{}^{v\gamma }=\left(\begin{array}{cc}\left(\widehat{\phi }^{(1)AB}\right)_{u;}^v& \left(\widehat{\phi }^{(2)AB}\right)_{u;}^\gamma \\ \left(\widehat{\phi }^{(3)AB}\right)_{\beta ;}^v& \left(\widehat{\phi }^{(4)AB}\right)_{\beta ;}^\gamma \end{array}\right),$$
(A.7)
where $`u,v=1,2,\mathrm{},N`$, $`\alpha ,\beta =1,2`$ and
$`\left(\widehat{\phi }^{(1)AB}\right)_{u;}^v`$ $`=`$ $`{\displaystyle \frac{1}{4(y^2+\rho ^2)^2}}\{y^2[16(\overline{\xi }^{\dot{\alpha }B}\overline{\xi }_{\dot{\beta }}^A\overline{\xi }^{\dot{\alpha }A}\overline{\xi }_{\dot{\beta }}^B)w_{u;\dot{\alpha }}\overline{w}^{\dot{\beta };v}`$
$`+`$ $`4w_{u;\dot{\alpha }}(\overline{\xi }^{\dot{\alpha }B}\overline{\nu }^{Av}\overline{\xi }^{\dot{\alpha }A}\overline{\nu }^{Bv})]`$
$`+`$ $`(y^2+\rho ^2)\left[4(\overline{\xi }_{\dot{\alpha }}^B\nu _u^A\overline{\xi }_{\dot{\alpha }}^A\nu _u^B)\overline{w}^{\dot{\alpha };v}+(\nu _u^B\overline{\nu }^{Av}\nu _u^A\overline{\nu }^{Bv})\right]`$
$`+`$ $`y^{\dot{\alpha }\delta }[16(\eta _\delta ^B\overline{\xi }_{\dot{\beta }}^A\eta _\delta ^A\overline{\xi }_{\dot{\beta }}^B)w_{u;\dot{\alpha }}\overline{w}^{\dot{\beta };v}4w_{u;\dot{\alpha }}(\eta _\delta ^B\overline{\nu }^{Av}\eta _\delta ^A\overline{\nu }^{Bv})]\}`$
$`\left(\widehat{\phi }^{(2)AB}\right)_{u;}^\gamma `$ $`=`$ $`{\displaystyle \frac{1}{4(y^2+\rho ^2)^2}}\{16y^2w_{u;\dot{\alpha }}(\overline{\xi }^{\dot{\alpha }B}\eta ^{\gamma A}\overline{\xi }^{\dot{\alpha }A}\eta ^{\gamma B})+4(y^2+\rho ^2)(\nu _u^B\eta ^{\gamma A}\nu _u^A\eta ^{\gamma B})\text{}`$
$``$ $`w_{u;\dot{\alpha }}[16y^{\dot{\alpha }\delta }(\eta _\delta ^B\eta ^{\gamma A}\eta _\delta ^A\eta ^{\gamma B})+{\displaystyle \frac{1}{2}}{\displaystyle \frac{y^2+\rho ^2}{\rho ^2}}y^{\dot{\alpha }\gamma }(\overline{\nu }^{Au}\nu _u^B\overline{\nu }^{Br}\nu _r^A)]\}`$
$`\left(\widehat{\phi }^{(3)AB}\right)_{\beta ;}^v`$ $`=`$ $`{\displaystyle \frac{1}{4(y^2+\rho ^2)^2}}\{\rho ^2[16y_{\beta \dot{\alpha }}(\overline{\xi }^{\dot{\alpha }B}\overline{\xi }_{\dot{\beta }}^A\overline{\xi }^{\dot{\alpha }A}\overline{\xi }_{\dot{\beta }}^B)\overline{w}^{\dot{\beta };v}4y_{\beta \dot{\alpha }}(\overline{\xi }^{\dot{\alpha }B}\overline{\nu }^{Av}\overline{\xi }^{\dot{\alpha }A}\overline{\nu }^{Bv})`$
$``$ $`16(\eta _\beta ^B\overline{\xi }_{\dot{\alpha }}^A\eta _\beta ^A\overline{\xi }_{\dot{\alpha }}^B)\overline{w}^{\dot{\alpha };v}+4(\eta _\beta ^B\overline{\nu }^{Av}\eta _\beta ^A\overline{\nu }^{Bv})]\}`$
$`\left(\widehat{\phi }^{(4)AB}\right)_{\beta ;}^\gamma `$ $`=`$ $`{\displaystyle \frac{\rho ^2}{4(y^2+\rho ^2)^2}}[16y_{\beta \dot{\alpha }}(\overline{\xi }^{\dot{\alpha }B}\eta ^{\gamma A}\overline{\xi }^{\dot{\alpha }A}\eta ^{\gamma B})+16(\eta _\beta ^B\eta ^{\gamma A}\eta _\beta ^B\eta ^{\gamma A})`$ (A.8)
$`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{y^2+\rho ^2}{\rho ^2}}\delta _\beta ^\gamma (\overline{\nu }^{Ar}\nu _r^B\overline{\nu }^{Br}\nu _r^A)].`$
## Appendix B Instanton induced two-point functions of BMN operators
In this appendix we present some details of the calculations of one-instanton contributions to the two-point functions of BMN operators discussed in section 5.
### B.1 Two-impurity operator in the $`\mathrm{𝟗}`$ of SO(4)<sub>R</sub>
As shown in section 5.1 the leading semi-classical contribution to the two-point functions of two impurity operators in the $`\mathrm{𝟗}`$ of SO(4)<sub>R</sub> vanishes because the superconformal modes cannot be soaked up. A non-zero result is obtained including for one scalar field in each operator the six fermion solution.
In the case of the component considered in section 5.1 the terms in the two-point function which contain the correct combination of fermion modes to give a non vanishing contribution are
$`G_\mathrm{𝟗}(x_1,x_2)`$ $`=`$ $`{\displaystyle \frac{1}{J\left(\frac{g_{_{\mathrm{YM}}}^2N}{8\pi ^2}\right)^{J+2}}}{\displaystyle \underset{p,q=0}{\overset{J}{}}}\mathrm{cos}\left({\displaystyle \frac{2\pi ipn}{J}}\right)\mathrm{cos}\left({\displaystyle \frac{2\pi iqm}{J}}\right)`$ (B.1)
$`\times \{\text{}\mathrm{Tr}\left[\left(Z^{Jp}\phi ^{13}Z^p\phi ^{13}\right)(x_1)\right]\mathrm{Tr}\left[\left(\overline{Z}^{Jq}\phi ^{24}\overline{Z}^q\phi ^{24}\right)(x_2)\right]`$
$`\text{}+\mathrm{Tr}\left[\left(Z^{Jp}\phi ^{24}Z^p\phi ^{24}\right)(x_1)\right]\mathrm{Tr}\left[\left(\overline{Z}^{Jq}\phi ^{13}\overline{Z}^q\phi ^{13}\right)(x_2)\right]\}.`$
The other terms vanish in the instanton background either because they do not contain all the required superconformal modes or because of the integration over the five-sphere. For instance if one considers $`\mathrm{Tr}\left[\left(Z^{Jp}\phi ^{13}Z^p\phi ^{13}\right)(x_1)\right]\mathrm{Tr}\left[\left(\overline{Z}^{Jq}\phi ^{13}\overline{Z}^q\phi ^{13}\right)(x_2)\right]`$ it is easy to verify that the superconformal modes can be soaked up, but the resulting five-sphere integral vanishes because among the remaining fermion modes different flavours appear with different multiplicities.
Let us consider the terms in (B.1) where in each trace one scalar is understood to be replaced with the six-fermion solution. In the first expectation value the two traces contain respectively the following combinations of fermion modes
$`\left(𝓂_\mathrm{f}^1\right)^{J+3}\left(𝓂_\mathrm{f}^2\right)^1\left(𝓂_\mathrm{f}^3\right)^3\left(𝓂_\mathrm{f}^4\right)^{J+1}`$
$`\left(𝓂_\mathrm{f}^1\right)^1\left(𝓂_\mathrm{f}^2\right)^{J+3}\left(𝓂_\mathrm{f}^3\right)^{J+1}\left(𝓂_\mathrm{f}^4\right)^3.`$ (B.2)
Using the fact that each trace contains one $`\overline{\xi }`$ mode not part of a $`\zeta `$ we can soak up the superconformal modes selecting the following combinations of fermion modes in the two traces
$`\mathrm{Tr}\left(Z^{Jp}\phi ^{13}Z^p\phi ^{13}\right)`$ $``$ $`\left(\zeta ^1\right)^2\zeta ^2\left(\zeta ^3\right)^2\left(\zeta ^4\right)^2\overline{\xi }^1\left(\overline{\nu }^{[1}\nu ^{4]}\right)^{J1}\left(\overline{\nu }^{[1}\nu ^{3]}\right)`$
$`\mathrm{Tr}\left(\overline{Z}^{Jp}\phi ^{24}\overline{Z}^p\phi ^{24}\right)`$ $``$ $`\zeta ^1\left(\zeta ^2\right)^2\left(\zeta ^3\right)^2\left(\zeta ^4\right)^2\overline{\xi }^2\left(\overline{\nu }^{[2}\nu ^{3]}\right)^{J1}\left(\overline{\nu }^{[2}\nu ^{4]}\right).`$ (B.3)
Similarly the two traces in the second term in (B.1) contain respectively
$$\left(𝓂_\mathrm{f}^1\right)^{J+1}\left(𝓂_\mathrm{f}^2\right)^3\left(𝓂_\mathrm{f}^3\right)^1\left(𝓂_\mathrm{f}^4\right)^{J+3}$$
(B.4)
and
$$\left(𝓂_\mathrm{f}^1\right)^3\left(𝓂_\mathrm{f}^2\right)^{J+1}\left(𝓂_\mathrm{f}^3\right)^{J+3}\left(𝓂_\mathrm{f}^4\right)^1$$
(B.5)
and we need to consider
$`\mathrm{Tr}\left(Z^{Jp}\phi ^{24}Z^p\phi ^{24}\right)`$ $``$ $`\left(\zeta ^1\right)^2\left(\zeta ^2\right)^2\zeta ^3\left(\zeta ^4\right)^2\overline{\xi }^4\left(\overline{\nu }^{[1}\nu ^{4]}\right)^{J1}\left(\overline{\nu }^{[2}\nu ^{4]}\right)`$
$`\mathrm{Tr}\left(\overline{Z}^{Jp}\phi ^{13}\overline{Z}^p\phi ^{13}\right)`$ $``$ $`\left(\zeta ^1\right)^2\left(\zeta ^2\right)^2\left(\zeta ^3\right)^2\zeta ^4\overline{\xi }^3\left(\overline{\nu }^{[2}\nu ^{3]}\right)^{J1}\left(\overline{\nu }^{[1}\nu ^{3]}\right).`$ (B.6)
These expressions contain the correct combinations of fermion superconformal modes such that the corresponding integration is non-zero. The two terms in (B.1) give rise to
$`{\displaystyle }{\displaystyle \underset{A=1}{\overset{4}{}}}\mathrm{d}^2\eta ^A\mathrm{d}^2\overline{\xi }^A\{[\left(\zeta ^1\right)^2\zeta ^2\left(\zeta ^3\right)^2\left(\zeta ^4\right)^2\overline{\xi }^1](x_1)[\zeta ^1\left(\zeta ^2\right)^2\left(\zeta ^3\right)^2\left(\zeta ^4\right)^2\overline{\xi }^2](x_2)`$
$`+[\left(\zeta ^1\right)^2\left(\zeta ^2\right)^2\zeta ^3\left(\zeta ^4\right)^2\overline{\xi }^4](x_1)[\left(\zeta ^1\right)^2\left(\zeta ^2\right)^2\left(\zeta ^3\right)^2\zeta ^4\overline{\xi }^3](x_2)\}`$
$`(x_1x_0)(x_2x_0)(x_1x_2)^4.`$ (B.7)
After re-expressing the $`(\overline{\nu }^A\nu ^B)`$ bilinears in terms of $`\mathrm{\Omega }^{AB}`$’s as described in section 3 both (B.3) and (B.6) lead to the same five-sphere integral,
$$I_{S^5}=\mathrm{d}^5\mathrm{\Omega }\left(\mathrm{\Omega }^{14}\right)^{J1}\left(\mathrm{\Omega }^{23}\right)^{J1}\mathrm{\Omega }^{13}\mathrm{\Omega }^{24}.$$
(B.8)
This can be evaluated rewriting it as
$$I_{S^5}=_{_{i=1}^6\mathrm{\Omega }_i^2=1}\mathrm{d}^6\mathrm{\Omega }\left(\mathrm{\Sigma }_i^{14}\mathrm{\Omega }^i\right)^{J1}\left(\mathrm{\Sigma }_j^{23}\mathrm{\Omega }^j\right)^{J1}\left(\mathrm{\Sigma }_k^{13}\mathrm{\Omega }^k\right)\left(\mathrm{\Sigma }_l^{24}\mathrm{\Omega }^l\right)$$
(B.9)
where the symbols $`\mathrm{\Sigma }_i^{AB}`$ are defined in (A.1). Defining $`\mathrm{\Omega }\mathrm{\Sigma }_i^{14}\mathrm{\Omega }^i=(\mathrm{\Omega }^1+i\mathrm{\Omega }^4)`$, $`\overline{\mathrm{\Omega }}\mathrm{\Sigma }_i^{23}\mathrm{\Omega }^i=(\mathrm{\Omega }^1i\mathrm{\Omega }^4)`$, $`\stackrel{~}{\mathrm{\Omega }}\mathrm{\Sigma }_i^{13}\mathrm{\Omega }^i=(\mathrm{\Omega }^2+i\mathrm{\Omega }^5)`$ and $`\overline{\stackrel{~}{\mathrm{\Omega }}}\mathrm{\Sigma }_i^{24}\mathrm{\Omega }^i=(\mathrm{\Omega }^2i\mathrm{\Omega }^5)`$ the integral reduces to
$`I_{S^5}`$ $`=`$ $`{\displaystyle \mathrm{d}^6\mathrm{\Omega }\delta \left(\underset{i=1}{\overset{6}{}}\mathrm{\Omega }_i^21\right)\left(\mathrm{\Omega }\overline{\mathrm{\Omega }}\right)^{J1}\left(\stackrel{~}{\mathrm{\Omega }}\overline{\stackrel{~}{\mathrm{\Omega }}}\right)}`$ (B.10)
$`=`$ $`{\displaystyle d\mathrm{\Omega }d\overline{\mathrm{\Omega }}d\stackrel{~}{\mathrm{\Omega }}d\overline{\stackrel{~}{\mathrm{\Omega }}}\mathrm{d}^2\mathrm{\Omega }^I\delta (\mathrm{\Omega }^I\mathrm{\Omega }^I+\mathrm{\Omega }\overline{\mathrm{\Omega }}+\stackrel{~}{\mathrm{\Omega }}\overline{\stackrel{~}{\mathrm{\Omega }}}1)\left(\mathrm{\Omega }\overline{\mathrm{\Omega }}\right)^{J1}\left(\stackrel{~}{\mathrm{\Omega }}\overline{\stackrel{~}{\mathrm{\Omega }}}\right)},`$
where $`\mathrm{\Omega }^I=(\mathrm{\Omega }^3,\mathrm{\Omega }^6)`$. Introducing polar coordinates for the $`\mathrm{\Omega }^I`$ directions
$`I_{S^5}`$ $`=`$ $`2\pi {\displaystyle drrd\mathrm{\Omega }d\overline{\mathrm{\Omega }}d\stackrel{~}{\mathrm{\Omega }}d\overline{\stackrel{~}{\mathrm{\Omega }}}\delta (r^2+\mathrm{\Omega }\overline{\mathrm{\Omega }}+\stackrel{~}{\mathrm{\Omega }}\overline{\stackrel{~}{\mathrm{\Omega }}}1)\left(\mathrm{\Omega }\overline{\mathrm{\Omega }}\right)^{J1}\left(\stackrel{~}{\mathrm{\Omega }}\overline{\stackrel{~}{\mathrm{\Omega }}}\right)}`$ (B.11)
$`=`$ $`\pi {\displaystyle _{\mathrm{\Omega }\overline{\mathrm{\Omega }}+\stackrel{~}{\mathrm{\Omega }}\overline{\stackrel{~}{\mathrm{\Omega }}}1}}d\mathrm{\Omega }d\overline{\mathrm{\Omega }}d\stackrel{~}{\mathrm{\Omega }}d\overline{\stackrel{~}{\mathrm{\Omega }}}\left(\mathrm{\Omega }\overline{\mathrm{\Omega }}\right)^{J1}\left(\stackrel{~}{\mathrm{\Omega }}\overline{\stackrel{~}{\mathrm{\Omega }}}\right).`$
The remaining integrals are straightforward
$`I_{S^5}`$ $`=`$ $`2\pi ^2{\displaystyle _{\mathrm{\Omega }\overline{\mathrm{\Omega }}1}}d\mathrm{\Omega }d\overline{\mathrm{\Omega }}\left(\mathrm{\Omega }\overline{\mathrm{\Omega }}\right)^{J1}{\displaystyle _0^{\sqrt{1\mathrm{\Omega }\overline{\mathrm{\Omega }}}}}dzz^3`$ (B.12)
$`=`$ $`{\displaystyle \frac{\pi ^2}{2}}{\displaystyle _{\mathrm{\Omega }\overline{\mathrm{\Omega }}1}}d\mathrm{\Omega }d\overline{\mathrm{\Omega }}\left(1\mathrm{\Omega }\overline{\mathrm{\Omega }}\right)^2\left(\mathrm{\Omega }\overline{\mathrm{\Omega }}\right)^{J1}={\displaystyle \frac{\pi ^3}{2}}{\displaystyle _0^1}dzz\left(1z^2\right)^2z^{2J}`$
$`=`$ $`{\displaystyle \frac{\pi ^3}{J(J+1)(J+2)}}.`$
As observed in section 5.1 the exact dependence on the bosonic moduli in the two-point function $`G_\mathrm{𝟗}(x_1,x_2)`$ cannot be determined without knowing the six-fermion solution. Dimensional analysis indicates that the final bosonic integrations over position and size of the instanton are logarithmically divergent as expected in the presence of an instanton contribution to the anomalous dimension of the operator $`𝒪^{\{ij\}}`$.
### B.2 $`\epsilon `$-singlet four impurity operator
In order to compute the profiles of the operator $`𝒪_\mathrm{𝟏}`$ in (5.16) and its conjugate, which are needed in the calculation of the two-point function (5.18), we have to evaluate the traces (5.26). In the instanton background these are rewritten as traces of $`[N+2]`$-dimensional ADHM matrices. To compute these traces more efficiently it is convenient to define the $`[N+2]\times [N+2]`$ matrix
$$\left[U_{k_1,k_2}^{C_1,D_1;C_2,D_2}(\zeta ,\nu ,\overline{\nu })\right]_{u,\alpha }{}_{}{}^{v,\beta }=\left[\left(\stackrel{ˇ}{\phi }^{14}\right)^{k_1}\widehat{\phi }^{C_1D_1}\left(\stackrel{ˇ}{\phi }^{14}\right)^{k_2}\widehat{\phi }^{C_2D_2}\right]_{u,\alpha }{}_{}{}^{v,\beta }.$$
(B.13)
where the notation used is that introduced in (5.23)-(5.25). This has the standard block-form of ADHM matrices and the range of the indices here is the same as in (A.8) for the elementary scalar fields.
In terms of the matrix $`U_{k_i,k_j}^{C_i,D_i;C_j,D_j}(\zeta ,\nu ,\overline{\nu })`$ all the 35 traces we are interested in can be written as
$$\mathrm{Tr}\left[U_{k_1,k_2}^{C_1,D_1;C_2,D_2}(\zeta ,\nu ,\overline{\nu })U_{k_3,k_4}^{C_3,D_3;C_4,D_4}(\zeta ,\nu ,\overline{\nu })U_{k_5,k_6}^{C_5,D_5;C_6,D_6}(\zeta ,\nu ,\overline{\nu })U_{k_7,k_8}^{C_7,D_7;C_8,D_8}(\zeta ,\nu ,\overline{\nu })\right],$$
(B.14)
for appropriate choices of the indices $`C_i`$, $`D_i`$ and the exponents $`k_i`$, $`i=1,\mathrm{},8`$. For example the three traces written explicitly in (5.26) become
$$\mathrm{Tr}\left[U_{p_1,p_2}^{1,4;1,4}(\zeta ,\nu ,\overline{\nu })U_{p_3,p_4}^{1,4;1,4}(\zeta ,\nu ,\overline{\nu })U_{p_5,q}^{A_1,B_1;A_2,B_2}(\zeta ,\nu ,\overline{\nu })U_{r,s}^{A_3,B_3;A_4,B_4}(\zeta ,\nu ,\overline{\nu })\right],$$
$$\mathrm{Tr}\left[U_{p_1,p_2}^{1,4;1,4}(\zeta ,\nu ,\overline{\nu })U_{p_3,p_4}^{1,4;A_1,B_1}(\zeta ,\nu ,\overline{\nu })U_{q_1,q_2}^{1,4;A_2,B_2}(\zeta ,\nu ,\overline{\nu })U_{r,s}^{A_3,B_3;A_4,B_4}(\zeta ,\nu ,\overline{\nu })\right]$$
and
$$\mathrm{Tr}\left[U_{p,q}^{A_1,B_1;A_2,B_2}(\zeta ,\nu ,\overline{\nu })U_{r,s_1}^{A_3,B_3;1,4}(\zeta ,\nu ,\overline{\nu })U_{s_2,s_3}^{1,4;1,4}(\zeta ,\nu ,\overline{\nu })U_{s_4,s_5}^{1,4;A_4,B_4}(\zeta ,\nu ,\overline{\nu })\right].$$
The generic trace (B.14) is thus the only one that needs to be evaluated. It can be computed using the building blocks (A.8) and the result is
$`\mathrm{Tr}\left[U_{k_1,k_2}^{C_1,D_1;C_2,D_2}(\zeta ,\nu ,\overline{\nu })U_{k_3,k_4}^{C_3,D_3;C_4,D_4}(\zeta ,\nu ,\overline{\nu })U_{k_5,k_6}^{C_5,D_5;C_6,D_6}(\zeta ,\nu ,\overline{\nu })U_{k_7,k_8}^{C_7,D_7;C_8,D_8}(\zeta ,\nu ,\overline{\nu })\right]`$
$`={\displaystyle \frac{1}{2^{3J8}}}{\displaystyle \frac{\rho ^8}{[(xx_0)^2+\rho ^2]^{J+8}}}\left(\overline{\nu }^{[1}\nu ^{4]}\right)^{J8}\{\left[\left(\zeta ^{D_1}\zeta ^{D_2}\right)\left(\zeta ^{D_3}\zeta ^{D_4}\right)\left(\zeta ^{D_5}\zeta ^{D_6}\right)\left(\zeta ^{D_7}\zeta ^{D_8}\right)\right]\text{}`$
$`\left[\left(\overline{\nu }^{[C_8}\nu ^{4]}\right)\left(\overline{\nu }^{[1}\nu ^{C_1]}\right)+\left(\overline{\nu }^{[C_8}\nu ^{1]}\right)\left(\overline{\nu }^{[4}\nu ^{C_1]}\right)\right]\left[\left(\overline{\nu }^{[C_2}\nu ^{4]}\right)\left(\overline{\nu }^{[1}\nu ^{C_3]}\right)+\left(\overline{\nu }^{[C_2}\nu ^{1]}\right)\left(\overline{\nu }^{[4}\nu ^{C_3]}\right)\right]`$
$`\left[\left(\overline{\nu }^{[C_4}\nu ^{4]}\right)\left(\overline{\nu }^{[1}\nu ^{C_5]}\right)+\left(\overline{\nu }^{[C_4}\nu ^{1]}\right)\left(\overline{\nu }^{[4}\nu ^{C_5]}\right)\right]\left[\left(\overline{\nu }^{[C_6}\nu ^{4]}\right)\left(\overline{\nu }^{[1}\nu ^{C_7]}\right)+\left(\overline{\nu }^{[C_6}\nu ^{1]}\right)\left(\overline{\nu }^{[4}\nu ^{C_7]}\right)\right]`$
$`+\left[\left(\zeta ^{D_2}\zeta ^{D_3}\right)\left(\zeta ^{D_4}\zeta ^{D_5}\right)\left(\zeta ^{D_6}\zeta ^{D_7}\right)\left(\zeta ^{D_8}\zeta ^{D_1}\right)\right]\left[\left(\overline{\nu }^{[C_1}\nu ^{4]}\right)\left(\overline{\nu }^{[1}\nu ^{C_2]}\right)+\left(\overline{\nu }^{[C_1}\nu ^{1]}\right)\left(\overline{\nu }^{[4}\nu ^{C_2]}\right)\right]`$
$`\left[\left(\overline{\nu }^{[C_3}\nu ^{4]}\right)\left(\overline{\nu }^{[1}\nu ^{C_4]}\right)+\left(\overline{\nu }^{[C_3}\nu ^{1]}\right)\left(\overline{\nu }^{[4}\nu ^{C_4]}\right)\right]\left[\left(\overline{\nu }^{[C_5}\nu ^{4]}\right)\left(\overline{\nu }^{[1}\nu ^{C_6]}\right)+\left(\overline{\nu }^{[C_5}\nu ^{1]}\right)\left(\overline{\nu }^{[4}\nu ^{C_6]}\right)\right]`$
$`[\left(\overline{\nu }^{[C_7}\nu ^{4]}\right)\left(\overline{\nu }^{[1}\nu ^{C_8]}\right)+\left(\overline{\nu }^{[C_7}\nu ^{1]}\right)\left(\overline{\nu }^{[4}\nu ^{C_8]}\right)]+\mathrm{permutations}\text{}\},`$ (B.15)
where the permutations not indicated explicitly correspond to antisymmetrisation in all the $`(C_i,D_i)`$ pairs.
The key feature of (B.15) is that it depends on the set of indices $`(C_i,D_i)`$, but not on the exponents $`k_i`$. The traces (5.26) require four pairs of $`(C_i,D_i)`$ indices to be $`(1,4)`$ while the remaining four pairs correspond to the $`(A_k,B_k)`$ indices carried by the impurities. The fact that (B.15) does not depend on the $`k_i`$’s means that the traces (5.26) do not depend on the exponents on the $`Z`$’s, but only on the relative positions of the $`\widehat{Z}`$’s with respect to the impurities. Therefore when substituting into the definition (5.16) of the operator the traces (B.15) can be taken out of the sums over the indices $`q`$, $`r`$, $`s`$. After substituting the values of the indices corresponding to the various terms in the expansion (5.19) and some simple Fierz rearrangements the profile of the operator $`𝒪_\mathrm{𝟏}`$ takes the form of a common factor containing the dependence on the bosonic and fermionic moduli, multiplying the combination $`K(n_1,n_2,n_3;J)`$ of 35 sums which contain the dependence on the mode numbers $`n_1`$, $`n_2`$ and $`n_3`$, see (5.31). To illustrate more concretely how this works let us describe explicitly one particular term. We consider the first trace in (5.19) and compute the contribution of the last type in (5.26) for this trace. We have to evaluate
$`{\displaystyle \frac{1}{\sqrt{J^3\left(\frac{g_{_{\mathrm{YM}}}^2N}{8\pi ^2}\right)^{J+4}}}}\begin{array}{c}{\displaystyle \underset{q,r,s_1,\mathrm{},s_5=0}{\overset{J}{}}}\\ q+r+s_1+\mathrm{}+s_5J4\\ s_1+\mathrm{}+s_5=s4\end{array}\mathrm{e}^{2\pi i[(n_1+n_2+n_3)q+(n_2+n_3)r+n_3s]/J}`$ (B.19)
$`\times \mathrm{Tr}\left(\stackrel{ˇ}{Z}^{J(q+r+s)}\widehat{\phi }^{12}\stackrel{ˇ}{Z}^q\widehat{\phi }^{13}\stackrel{ˇ}{Z}^r\widehat{\phi }^{24}\stackrel{ˇ}{Z}^{s_1}\widehat{Z}\stackrel{ˇ}{Z}^{s_2}\widehat{Z}\stackrel{ˇ}{Z}^{s_3}\widehat{Z}\stackrel{ˇ}{Z}^{s_4}\widehat{Z}\stackrel{ˇ}{Z}^{s_5}\widehat{\phi }^{34}\right)`$
$`={\displaystyle \frac{1}{\sqrt{J^3\left(\frac{g_{_{\mathrm{YM}}}^2N}{8\pi ^2}\right)^{J+4}}}}\begin{array}{c}{\displaystyle \underset{q,r,s_1,\mathrm{},s_5=0}{\overset{J}{}}}\\ q+r+s_1+\mathrm{}+s_5J4\\ s_1+\mathrm{}+s_5=s4\end{array}\mathrm{e}^{2\pi i[(n_1+n_2+n_3)q+(n_2+n_3)r+n_3s]/J}`$ (B.23)
$`\times \mathrm{Tr}\left(U_{p,q}^{1,2;1,3}U_{r,s_1}^{2,4;1,4}U_{s_2,s_3}^{1,4;1,4}U_{s_4,s_5}^{1,4;3,4}\right).`$ (B.24)
Using (B.15) with the particular choice of indices in the trace in (B.24) we get (up to a numerical constant)
$`\mathrm{Tr}\left(U_{p,q}^{1,2;1,3}U_{r,s_1}^{2,4;1,4}U_{s_2,s_3}^{1,4;1,4}U_{s_4,s_5}^{1,4;3,4}\right)={\displaystyle \frac{1}{2^{3J+8}}}{\displaystyle \frac{\rho ^8}{[(x_1x_0)^2+\rho ^2]^{J+8}}}\left(\overline{\nu }^{[1}\nu ^{4]}\right)^J`$
$`\times \left[(\zeta ^1\zeta ^1)(\zeta ^2\zeta ^3)(\zeta ^2\zeta ^3)\left(\zeta ^4\zeta ^4\right)(\zeta ^1\zeta ^4)(\zeta ^4\zeta ^4)(\zeta ^2\zeta ^3)\left(\zeta ^2\zeta ^3\right)\right]\text{}`$
$`={\displaystyle \frac{3}{2}}{\displaystyle \frac{1}{2^{3J+8}}}{\displaystyle \frac{\rho ^8}{[(x_1x_0)^2+\rho ^2]^{J+8}}}\left(\overline{\nu }^{[1}\nu ^{4]}\right)^J\left[\left(\zeta ^1\right)^2\left(\zeta ^2\right)^2\left(\zeta ^3\right)^2\left(\zeta ^4\right)^2\right],\text{}`$ (B.25)
where the last line has been obtained using simple Fierz rearrangements on the $`\zeta `$’s. As anticipated the trace is independent of the exponents, $`q`$, $`r`$, $`s_i`$. Equation (B.24) then becomes
$`{\displaystyle \frac{3}{2}}{\displaystyle \frac{1}{2^{3J+8}}}{\displaystyle \frac{1}{\sqrt{J^3\left(\frac{g_{_{\mathrm{YM}}}^2N}{8\pi ^2}\right)^{J+4}}}}{\displaystyle \frac{\rho ^8}{[(x_1x_0)^2+\rho ^2]^{J+8}}}\left(\overline{\nu }^{[1}\nu ^{4]}\right)^J\left[\left(\zeta ^1\right)^2\left(\zeta ^2\right)^2\left(\zeta ^3\right)^2\left(\zeta ^4\right)^2\right]`$
$`\times \begin{array}{c}{\displaystyle \underset{q,r,s=0}{\overset{J}{}}}\\ q+r+sJ\end{array}\mathrm{e}^{2\pi i[(n_1+n_2+n_3)q+(n_2+n_3)r+n_3s]/J}{\displaystyle \frac{1}{4!}}s(s1)(s2)(s3),`$ (B.28)
where we have used the fact that there is no dependence on the single exponents, $`s_1`$,…,$`s_5`$, so that the result is independent of the way the four $`\widehat{Z}`$ are distributed among the last $`s`$ $`Z`$’s. This leads to the factor $`\frac{1}{4!}s(s1)(s2)(s3)`$ which is a multiplicity coefficient associated with the number of ways of picking four identical $`\widehat{Z}`$’s out of $`s`$ $`Z`$’s. Equation (B.28) illustrates the factorisation of the result into two terms, the first line containing the dependence on the instanton moduli and the second line containing the dependence on the mode numbers.
The function $`K(n_1,n_2,n_3;J)`$ takes the form
$$K(n_1,n_2,n_3;J)=\underset{a=1}{\overset{35}{}}c_a𝒮_a(n_1,n_2,n_3;J),$$
(B.29)
where each of the $`𝒮_a(n_1,n_2,n_3;J)`$ is a sum similar to the second line of (B.28) with different summand corresponding to the different multiplicity factors associated with the distributions of $`\widehat{Z}`$’s in the traces (5.26). Table 3 summarises the contributions to (B.29).
Using the coefficients given in table 3, and noting that the phase factor factorises, (B.29) can be written as
$`K(n_1,n_2,n_3)`$ $`=`$ $`\begin{array}{c}{\displaystyle \frac{3}{4}}{\displaystyle \underset{p,q,r,s=0}{\overset{J}{}}}\\ q+r+s+p=J\end{array}\mathrm{e}^{2\pi i[(n_1+n_2+n_3)q+(n_2+n_3)r+n_3s]/J}(pq+rs)(p+q+r+s)^3`$ (B.32)
$`=`$ $`\begin{array}{c}{\displaystyle \frac{3J^3}{4}}{\displaystyle \underset{p,q,r=0}{\overset{J}{}}}\\ p+q+rJ\end{array}\mathrm{e}^{2\pi i[(n_1+n_2+n_3)p+(n_2+n_3)q+n_3r]/J}(2p+2rJ).`$ (B.35)
The sums in (B.35) can be approximated with integrals in the $`J\mathrm{}`$ limit, which can then be evaluated differentiating a generating function. The relevant generating function is given by
$`g(a_1,a_2,a_3,a_4)`$ $`=`$ $`{\displaystyle _0^1}dx{\displaystyle _0^{1x}}dy{\displaystyle _0^{1xy}}dz\mathrm{e}^{2\pi i[a_1x+a_2y+a_3z+a_4(1xyz)]}`$ (B.36)
$`=`$ $`{\displaystyle \frac{i}{8\pi ^3}}{\displaystyle \underset{i=1}{\overset{4}{}}}{\displaystyle \frac{\mathrm{e}^{2\pi ia_i}}{_{j=1,ji}^4(a_ia_j)}}.`$
Thus the above sum requires evaluating
$$f(a_1,a_2,a_3,a_4)=\frac{3i}{8\pi }\left(2\frac{}{a_1}+2\frac{}{a_3}2\pi i\right)g(a_1,a_2,a_3,a_4).$$
(B.37)
As discussed in section 5.2.1 in order to get the correct mode number dependence we need to antisymmetrise (B.37) with respect to the exchange of pairs of mode numbers. Therefore we need to compute
$`\underset{n_3n_1}{lim}[f(n_1+n_2+n_3,n_2+n_3,n_3,0)f(n_1+n_2+n_3,n_1+n_3,n_3,0)`$
$`f(n_3,n_1n_3,n_1n_2n_3,0)+f(n_3,n_2n_3,n_1n_2n_3,0)]`$
$`={\displaystyle \frac{3}{8\pi ^2}}{\displaystyle \frac{1}{(n_1n_2)}},`$ (B.38)
where only in the case where we impose pairwise equality do we get a non-zero result. In conclusion the mode number dependence in the profile of the operator $`𝒪_\mathrm{𝟏}`$ is
$$K(n_1,n_2,n_3;J)=\frac{3}{8\pi ^2}\frac{J^7}{(n_1n_2)},$$
(B.39)
where the factor of $`J^7`$ is the combination of the $`J^3`$ in (B.35) and a $`J^4`$ arising from the conversion of the sums into integrals in the continuum limit.
The two-point function $`G_\mathrm{𝟏}(x_1,x_2)`$ thus becomes
$`G_\mathrm{𝟏}(x_1,x_2)={\displaystyle \frac{J^{11}\mathrm{e}^{2\pi i\tau }}{N^{7/2}}}{\displaystyle \frac{1}{(n_1n_2)(m_1m_2)}}{\displaystyle \frac{\mathrm{d}^4x_0\mathrm{d}\rho }{\rho ^5}\frac{\rho ^{J+8}}{[(x_1x_0)^2+\rho ^2]^{J+8}}\frac{\rho ^{J+8}}{[(x_2x_0)^2+\rho ^2]^{J+8}}}`$
$`\times {\displaystyle }{\displaystyle \underset{A=1}{\overset{4}{}}}\mathrm{d}^2\eta ^A\mathrm{d}^2\overline{\xi }^A{\displaystyle \underset{B=1}{\overset{4}{}}}\left[\left(\zeta ^B\right)^2(x_1)\right]\left[\left(\zeta ^B\right)^2(x_2)\right]{\displaystyle }\mathrm{d}^5\mathrm{\Omega }\left(\mathrm{\Omega }^{14}\right)^J\left(\mathrm{\Omega }^{23}\right)^J.`$ (B.40)
The integrations in the second line of (B.40) are straightforward. The integrals over the fermion superconformal modes give $`(x_1x_2)^8`$, see (5.39). The five sphere integral is similar to that encountered in the two impurity case and can be calculated in a similar fashion. Proceeding as in (B.8)-(B.12) we get
$$I_{S^5}=\mathrm{d}^5\mathrm{\Omega }\left(\mathrm{\Omega }^{14}\right)^J\left(\mathrm{\Omega }^{23}\right)^J=d\mathrm{\Omega }d\overline{\mathrm{\Omega }}\mathrm{d}^4\mathrm{\Omega }^I\delta (\mathrm{\Omega }^I\mathrm{\Omega }^I+\mathrm{\Omega }\overline{\mathrm{\Omega }}1)\left(\mathrm{\Omega }\overline{\mathrm{\Omega }}\right)^J,$$
(B.41)
where $`\mathrm{\Omega }=(\mathrm{\Omega }^1+i\mathrm{\Omega }^4)`$, $`\overline{\mathrm{\Omega }}=(\mathrm{\Omega }^1i\mathrm{\Omega }^4)`$ and $`\mathrm{\Omega }^I=(\mathrm{\Omega }^2,\mathrm{\Omega }^3,\mathrm{\Omega }^5,\mathrm{\Omega }^6)`$, so that introducing spherical coordinates
$`I_{S^5}`$ $`=`$ $`2\pi ^2{\displaystyle drrd\mathrm{\Omega }d\overline{\mathrm{\Omega }}\delta (r^2+\mathrm{\Omega }\overline{\mathrm{\Omega }}1)\left(\mathrm{\Omega }\overline{\mathrm{\Omega }}\right)^J}`$ (B.42)
$`=`$ $`2\pi ^2{\displaystyle _{\mathrm{\Omega }\overline{\mathrm{\Omega }}1}}\mathrm{d}\mathrm{\Omega }\mathrm{d}\overline{\mathrm{\Omega }}((1\mathrm{\Omega }\overline{\mathrm{\Omega }})\left(\mathrm{\Omega }\overline{\mathrm{\Omega }}\right)^J={\displaystyle \frac{\pi ^3}{(J+1)(J+2)}}.`$
The integration over the bosonic part of the moduli space must be treated carefully since it is logarithmically divergent as expected in the presence of a contribution to the matrix of anomalous dimensions. The integrals need to be regulated for instance by dimensional regularisation of the $`x_0`$ integral. Introducing Feynman parameters we get
$`I_\mathrm{b}`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}^4x_0\mathrm{d}\rho }{\rho ^5}\frac{\rho ^{J+8}}{[(x_1x_0)^2+\rho ^2]^{J+8}}\frac{\rho ^{J+8}}{[(x_2x_0)^2+\rho ^2]^{J+8}}}`$ (B.43)
$`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(2J+16)}{[\mathrm{\Gamma }(J+8)]^2}}{\displaystyle _0^1}d\alpha _1d\alpha _2\delta (1\alpha _1\alpha _2)\alpha _1^{J+7}\alpha _2^{J+7}`$
$`\times `$ $`{\displaystyle \mathrm{d}^4x_0d\rho \frac{\rho ^{2J+11}}{[(x_0\alpha _1x_1\alpha _2x_2)^2+\rho ^2+\alpha _1\alpha _2x_{12}^2]^{2J+16}}}.`$
After dimensional regularisation,
$`I_\mathrm{b}I_\mathrm{b}^{(ϵ)}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(J+6)\mathrm{\Gamma }(J+8+ϵ)}{[\mathrm{\Gamma }(J+8)]^2}}\pi ^{2ϵ}{\displaystyle \frac{1}{(x_{12}^2)^{J+8+ϵ}}}{\displaystyle _0^1}d\alpha {\displaystyle \frac{1}{[\alpha (1\alpha )]^{1+ϵ}}}`$ (B.44)
$`=`$ $`{\displaystyle \frac{1}{ϵ}}{\displaystyle \frac{\mathrm{\Gamma }(J+6)\mathrm{\Gamma }(J+8+ϵ)}{[\mathrm{\Gamma }(J+8)]^2}}\pi ^{2ϵ}{\displaystyle \frac{1}{(x_{12}^2)^{J+8+ϵ}}}.`$
The $`1/ϵ`$ pole corresponds to a logarithmic divergence in dimensional regularisation.
Substituting into (B.40) we finally get
$$G_1(x_1,x_2)(g_2)^{7/2}\mathrm{e}^{\frac{8\pi ^2}{g_2\lambda ^{}}+i\theta }\frac{1}{(n_1n_2)(m_1m_2)}\frac{1}{(x_{12}^2)^{J+4}}\mathrm{log}\left(\mathrm{\Lambda }^2x_{12}^2\right),$$
(B.45)
where the exact numerical coefficient was given in (5.45). |
warning/0506/cond-mat0506103.html | ar5iv | text | # Soliton solution of continuum magnetization-equation in conducting ferromagnet with a spin-polarized current
## Abstract
Exact soliton solutions of a modified Landau-Lifshitz equation for the magnetization of conducting ferromagnet in the presence of a spin-polarized current are obtained by means of inverse scattering transformation. From the analytical solution effects of spin-current on the frequency, wave number, and dispersion law of spin wave are investigated. The one-soliton solution indicates obviously current-driven precession and periodic shape-variation as well. The inelastic collision of solitons by which we mean the shape change before and after collision appears due to the spin current. We, moreover, show that complete inelastic collisions can be achieved by adjusting spectrum and current parameters. This may lead to a potential technique for shape control of spin wave.
Considerable attention has been paid to the dynamics of magnetization associated with spin-polarized current in layered materials and in Mn oxides recently. Both theoretical and experimental investigations mainly concentrated on large magnetoresistance are of fundamental importance in the understanding of magnetism and applied interest in the fabrication of magnetic devices. Spin-transfer from spin-polarized current to magnetization of conducting ferromagnetic films is one of intriguing features which is theoretically proposed in Ref. Slonczewski ; Berger and subsequently verifiedKatine ; Bazaliy in experiment. Many studiesSun ; Waintal ; Stile ; Zhang ; Murakami on this phenomenon have been followed since the spin-transfer mechanism was first explained conceptually. However, the dynamics of magnetization in the presence of spin current has not been well understood. In Ref. Bazaliy a continuum equation for the magnetization of conducting ferromagnet in the presence of spin-polarized current is derived and is seen to be a modified Landau-Lifshitz equation with an additional topological term. The spatial dependence of the magnetization is evaluated Bazaliy and several solutions of one-dimension are also discussed. We in this paper study the soliton solution of the modified Landau-Lifshitz equation given in Ref. Bazaliy from which the current induced precession of the magnetization and soliton-soliton collisions are investigated. The effect of spin-current on the dynamics of magnetization is demonstrated explicitly. We obtain exact soliton solutions by means of inverse scattering transformation in one-dimensional geometry. An intriguing feature is that inelastic collisions generally appear due to the spin-polarized current and the complete inelastic collisions which may lead to a interesting technique of soliton filter and switch can be achieved in special cases by adjusting of the parameters of spin-polarized current.
Following Ref. Bazaliy we consider a current propagating through conducting ferromagnet and assume that the conducting electrons interact only with the local magnetization $`𝐌`$. A continuum equation for the magnetization is obtained in the local magnetization frame as a Landau-Lifschitz typeBazaliy ,
$$\frac{}{t}𝐌=𝐌\times \stackrel{~}{J}𝐌_{zz}\gamma 𝐌_z,$$
(1)
where $`𝐌(z,t)=(M^x(z,t)`$, $`M^y(z,t)`$, $`M^z(z,t))`$ is the local magnetization and $`\stackrel{~}{J}=\frac{g\mu _BM}{\mathrm{}}J`$ with $`J`$ being the exchange interacting constant between the local magnets. The parameter $`\gamma `$ describes effect of the current which is the same as defined in Ref. Bazaliy . The length of magnetization vector is set to unit , $`𝐌^2(z,t)=1`$, for the sake of simplicity. We begin with the inverse scattering transformation which is a useful method to solve the nonlinear equation (1). By means of Ablowitz-Kaup-Newell-Segur method one can construct Lax equations for the Eq. (1) as
$`{\displaystyle \frac{}{z}}\mathrm{\Psi }(z,t,\lambda )`$ $`=U_1\left(\lambda \right)\mathrm{\Psi }(z,t,\lambda ),`$
$`{\displaystyle \frac{}{t}}\mathrm{\Psi }(z,t,\lambda )`$ $`=U_2\left(\lambda \right)\mathrm{\Psi }(z,t,\lambda ),`$ (2)
where $`\lambda `$ is a spectrum parameter, $`\mathrm{\Psi }(z,t,\lambda )`$ is eigenfunction corresponding to the spectrum $`\lambda `$. The operators $`U_1\left(\lambda \right)`$ and $`U_2\left(\lambda \right)`$ are given in the following form
$`U_1`$ $`=i\stackrel{~}{J}\left(\lambda \lambda _0\right)\left(M\sigma \right),`$
$`U_2`$ $`=i2\stackrel{~}{J}^3\left(\lambda ^2\lambda _0^2\right)\left(M\sigma \right)`$
$`\stackrel{~}{J}^2\left(\lambda \lambda _0\right)\left(M\sigma \right)\left(M_z\sigma \right),`$ (3)
where the real parameter $`\lambda _0=\gamma (4\stackrel{~}{J}^2)^1`$ indicates the effect of the current and $`\sigma `$ is Pauli matrix. Thus Eq. (1) can be recovered from the compatibility condition $`\frac{}{t}U_1\frac{}{z}U_2+[U_1,U_2]=0`$. We consider the following natural boundary condition at initial time($`t=0`$), $`𝐌\left(z\right)(M^x,M^y,M^z)(0,0,1)\text{ as}\left|z\right|\mathrm{}`$. We then have the asymptotic form of Eq. (2) at $`\left|z\right|\mathrm{}`$,
$$_zE(z,\lambda )=L_0(\lambda )E(z,\lambda ),$$
(4)
where
$$E(z,\lambda )=e^{i\stackrel{~}{J}\left(\lambda \lambda _0\right)z\sigma _3},L_0(\lambda )=i\stackrel{~}{J}\left(\lambda \lambda _0\right)\sigma _3,$$
(5)
Based on the Lax equations (2), we can derive the exact solution of $`N`$-soliton trains by employing of the inverse scattering transformationzdli ; Ablowitz . As a special case that $`N=1`$ the exact one-soliton solution is given as follows:
$`M^x`$ $`={\displaystyle \frac{1}{\mathrm{\Delta }_1}}[2(\alpha _1\lambda _0)\beta _1\mathrm{sin}(\mathrm{\Phi }_1\varphi _1)\mathrm{cosh}\mathrm{\Theta }_1`$
$`2\beta _1^2\mathrm{cos}(\mathrm{\Phi }_1\varphi _1)\mathrm{sinh}\mathrm{\Theta }_1],`$
$`M^y`$ $`={\displaystyle \frac{1}{\mathrm{\Delta }_1}}[2(\alpha _1\lambda _0)\beta _1\mathrm{cos}(\mathrm{\Phi }_1\varphi _1)\mathrm{cosh}\mathrm{\Theta }_1`$
$`2\beta _1^2\mathrm{sin}(\mathrm{\Phi }_1\varphi _1)\mathrm{sinh}\mathrm{\Theta }_1],`$
$`M^z`$ $`=1{\displaystyle \frac{2\beta _1^2}{\mathrm{\Delta }_1}},`$ (6)
where
$`\mathrm{\Delta }_1`$ $`=|\lambda _1\lambda _0|^2\mathrm{cosh}^2\mathrm{\Theta }_1,`$
$`\mathrm{\Theta }_1`$ $`=2\stackrel{~}{J}\beta _1(zV_{1,M}t)z_1,`$
$`\mathrm{\Phi }_1`$ $`=2\stackrel{~}{J}\left(\alpha _1\lambda _0\right)z4\stackrel{~}{J}^3\left[\alpha _1^2\beta _1^2\lambda _0^2\right]t\varphi _1.`$
Parameter $`V_{1,M}=4\stackrel{~}{J}^2\alpha _1`$ denotes the velocity of envelope, $`z_1=\mathrm{ln}[(2\stackrel{~}{J}\beta _1)^1c_1]`$ is the center position, and $`\varphi _1=\mathrm{arg}\left[\stackrel{~}{J}\left(\lambda _1\lambda _0\right)\right]`$ is the initial phase of the spin wave. The parameter $`\lambda _1=\alpha _1+i\beta _1`$ denotes eigenvalue with $`\alpha _1`$, $`\beta _1`$ being the real and imaginary parts respectively, and $`c_1`$ is a real constant of integration. The solution Eq. (6) describes a current-driven precession of magnetization with periodic shape variation. The center of solitary wave moves with velocity $`V_{1,M}`$, while the wave amplitude and width vary periodically with time. We see that the spin-polarized current imparts a torque to the magnetization due to local exchange interaction between electron-spin and the magnetic moment. This observation is in accord with the prediction in Ref.Slonczewski ; Berger ; Myers ; Bazaliy ; Zhang . As a consequence of reaction the current flow is strongly affected by the orientation of the magnetic moments. Thus a higher electrical resistance in magnetic layer may occur. The spin-polarized current can be used to adjust the precession of magnetic moment and the wave shape as well. We then provide in principle a mechanism of current-control of the spin wave.
To see closely the physical significance of one-soliton solution it is helpful to show the parameter-dependence of Euler angles of the magnetization vector which in a spherical coordinate is written as $`𝐌(z,t)(\mathrm{sin}\theta \mathrm{cos}\phi ,\mathrm{sin}\theta \mathrm{sin}\phi ,\mathrm{cos}\theta )`$. From the Eq. (6) we find
$$\mathrm{cos}\theta =1\frac{A_M}{\mathrm{cosh}^2\left[ϝ_1^1(zV_{1,M}t)z_1\right]},$$
(7)
$$\phi =\frac{\pi }{2}\varphi _1+k_1z\mathrm{\Omega }_1t+\mathrm{arctan}\left(\frac{\beta _1}{\alpha _1\lambda _0}\mathrm{tanh}\mathrm{\Theta }_1\right),$$
(8)
where $`A_M=\frac{2\beta _1^2}{|\lambda _1\lambda _0|^2}`$, $`ϝ_1=\frac{1}{2\stackrel{~}{J}\beta _1}`$ are amplitude and width of the soliton respectively. The wave number is $`k_1=k_0k_S`$ with $`k_0=2\stackrel{~}{J}\alpha _1`$ denoting the wave number in the absence of the current while $`k_S=2\stackrel{~}{J}\lambda _0`$ is the wave number shift induced by the spin-polarized current. The frequency of magnetization precession is seen to be $`\mathrm{\Omega }_1=\mathrm{\Omega }_0\mathrm{\Omega }_S`$ with $`\mathrm{\Omega }_0=4\stackrel{~}{J}^3\left(\alpha _1^2\beta _1^2\right)`$ being the frequency in the absence of current and $`\mathrm{\Omega }_S=4\stackrel{~}{J}^3\lambda _0^2`$ the frequency shift induced by spin-polarized current. We see that the effect of current reduces both the wave number and frequency. For a large enough current such that $`\lambda _0^2>\alpha _1^2\beta _1^2`$ an instability occurs Bazaliy . We can rewrite the frequency, i.e. the energy spectrum as
$$\mathrm{\Omega }_1=\stackrel{~}{J}k_1\left(k_1+2k_S\right)4\stackrel{~}{J}^3\beta _1^2.$$
(9)
We then see that in the absence of current the minimum of the energy spectrum i.e. $`\mathrm{\Omega }_{0,\mathrm{min}}=0`$ is located at $`k_{0,\mathrm{min}}=\sqrt{4\stackrel{~}{J}^2\beta _1^2}`$ while the current shifts the position of minimum by an amount $`\delta =\sqrt{k_S^2+4\stackrel{~}{J}^2\beta _1^2}(k_S+\sqrt{4\stackrel{~}{J}^2\beta _1^2})`$. In the absence of spin-polarized current, i.e. $`\lambda _0=0`$, the solution Eqs. (7) (8) reduces to the soliton solution in an isotropic spin chain Takhtajan .
In the limit case that the amplitude $`A_M`$ approaches zero, namely $`\beta _10`$, the soliton width $`ϝ_1`$ diverges, the envelope velocity $`V_{1,M}`$ attains its maximum value $`2\sqrt{\stackrel{~}{J}\mathrm{\Omega }_0}`$, and the solution shown in Eq. (7) and Eq. (8) takes the form such that
$$M^z1,\text{ }\phi \frac{\pi }{2}\varphi _1+k_1z\mathrm{\Omega }_1t$$
indicating a small linear solution of magnon. In this case the quadratic dispersion law is seen to be $`\mathrm{\Omega }_1=\stackrel{~}{J}\left(k_0^2k_S^2\right)=\stackrel{~}{J}k_1\left(k_1+2k_S\right)`$. We also notice that the phase velocity of the precession is $`\frac{\mathrm{\Omega }_1}{k_1}=\frac{V_{1,M}}{2}+\frac{V_S}{2}`$ which possesses a correction value $`\frac{V_S}{2}`$ which is the half of envelope velocity. Whereas the group velocity of precession $`\frac{d\mathrm{\Omega }_1}{dk_1}=V_{1,M}`$ coincides with the envelope velocity.
We now consider another special case of the general $`N`$-soliton trains, i.e., the two-soliton solution which is seen to be
$`M^x`$ $`=\mathrm{Re}\left[i2\mathrm{\Gamma }_2(1i\mathrm{\Gamma }_1)\right],`$
$`M^y`$ $`=\mathrm{Im}\left[i2\mathrm{\Gamma }_2(1i\mathrm{\Gamma }_1)\right],`$
$`M^z`$ $`=\left|1i\mathrm{\Gamma }_1\right|^2\left|\mathrm{\Gamma }_2\right|^2,`$ (10)
where
$`\mathrm{\Gamma }_1`$ $`={\displaystyle \frac{1}{W}}\left[\left(g_1g_3\right)g_6+\left(g_2g_4\right)g_5\right],`$
$`\mathrm{\Gamma }_2`$ $`={\displaystyle \frac{1}{W}}\left[(\overline{g_1}\overline{g_3})g_8(\overline{g_2}\overline{g_4})g_7\right],`$
with
$`g_1`$ $`=1+\left|q_1\right|^2+\chi _1\overline{\chi }_2q_1\overline{q}_2,`$
$`g_2`$ $`=1+\left|q_2\right|^2+\overline{\chi }_1\chi _2\overline{q}_1q_2,`$
$`g_3`$ $`=\overline{\chi }_1\left|q_1\right|^2+\chi _1q_1\overline{q}_2,`$
$`g_4`$ $`=\overline{\chi }_2\left|q_2\right|^2+\chi _2\overline{q}_1q_2,`$
$`g_5`$ $`=\xi _1\left|q_1\right|^2\chi _1\xi _2q_1\overline{q}_2,`$
$`g_6`$ $`=\xi _2\left|q_2\right|^2\chi _2\xi _1\overline{q}_1q_2,`$
$`g_7`$ $`=\xi _1\overline{q}_1,`$
$`g_8`$ $`=\xi _2\overline{q}_2,`$
$`\chi _1`$ $`={\displaystyle \frac{2\beta _1\left(\lambda _1\lambda _0\right)}{i\left(\lambda _1\overline{\lambda }_2\right)\left|\lambda _1\lambda _0\right|}},`$
$`\chi _2`$ $`={\displaystyle \frac{2\beta _2\left(\lambda _2\lambda _0\right)}{i\left(\lambda _2\overline{\lambda }_1\right)\left|\lambda _2\lambda _0\right|}},`$
$`W`$ $`=g_1g_2g_3g_4,`$
$`q_j`$ $`=e^{\mathrm{\Theta }_j+i\mathrm{\Phi }_j},`$
$`\xi _j`$ $`=2\beta _j\left|\lambda _j\lambda _0\right|^1,`$
and
$`\mathrm{\Theta }_j`$ $`=2\stackrel{~}{J}\beta _j(zV_{j,M}t)z_j,`$
$`\mathrm{\Phi }_j`$ $`=k_jz\mathrm{\Omega }_jt\varphi _j,`$ (11)
where $`V_{j,M}=4\stackrel{~}{J}^2\alpha _j`$ denotes the velocity of envelope , $`z_j=\mathrm{ln}[(2\stackrel{~}{J}\beta _j)^1c_j]`$ the center position, $`\varphi _j=\mathrm{arg}\left[\stackrel{~}{J}\left(\lambda _j\lambda _0\right)\right]`$ the initial phase, $`k_j=2\stackrel{~}{J}\left(\alpha _j\lambda _0\right)`$ the wave number, and $`\mathrm{\Omega }_j=4\stackrel{~}{J}^3\left[\alpha _j^2\beta _j^2\lambda _0^2\right]`$ is frequency. The parameter $`\lambda _j=\alpha _j+i\beta _j`$ is eigenvalue parameter, and $`c_j`$ is real constant of integration, $`j=1,2`$. The solutions (10) describe in general a inelastic scattering process of two solitary waves with different center velocities and different shape-variation frequencies. Before collision, the two solitons move towards each other, one with velocity $`V_1`$ and shape variation frequency $`\mathrm{\Omega }_1`$, the other with $`V_2`$ and $`\mathrm{\Omega }_2`$ respectively. The interaction potential between two solitons is a complicated function of current-dependent parameter $`\lambda _0`$ and eigenvalue $`\lambda _j`$. For the case that $`\alpha _j=\beta _j`$, the shape-variation frequencies $`\mathrm{\Omega }_j(j=1,2)`$ of two-soliton depend only on the parameters of spin-polarized current seen from Eq. (11). In the case of $`\lambda _0=0`$, the solutions (10) reduce to that of the usual two-soliton solution with two center velocities while without shape change Takhtajan . A interesting process in the absence of spin-polarized current is that the collision can result in the interchange of amplitude $`A_j`$ and phase $`\mathrm{\Phi }_j(j=1,2)`$ like exactly in the case of elastic collision of two particles Takhtajan .
It is interesting to show the inelastic collision graphically. The head on collision is explained in Fig. 1 and Fig. 2 for suppressed amplitudes of $`M_1`$ and $`M_2`$ respectively after collision. This result shows that we may adjust the incoming spin current and the spectral parameters to control the shape of soliton of the magnetization. The dissipationless quantum spin current at room temperature reported in Ref. Murakami may be used to realize experimentally the soliton-control in future. Our theoretical observations predict the magnetic random-access memories in which the memory elements are controlled by local exchange-effect forces induced by spin-polaried current rather than by long-range magnetic fields.
In terms of inverse scattering transformation the exact soliton solutions for the magnetization in conducting ferromagnet in the presence of a spin-polarized current are obtained. Our solutions predict two intrinsic features of the effect of spin-polarized current on the magnetization: (1) Spin-polarized current induces the precession and shape variation of the solitary waves of magnetization. (2) The inelastic collision of solitons. The effect of spin-polarized current on the magnetization is similar to that of the periodically time-varying external magnetic fields reported earlier zdli and is in agreement with the observations in Refs. Slonczewski ; Berger ; Myers ; Bazaliy ; Zhang .
This work was supported by National Natural Science Foundation of China under Grant Nos. 10075032, 10174095, 90103024 and provincial overseas scholar foundation of Shanxi.
Figure caption
Fig. 1 Inelastic head on collision expressed by Eq. (10) when $`M_1`$ suppressed, where $`\lambda _1=0.4i0.3`$, $`\lambda _2=0.5+i0.45`$, $`\stackrel{~}{J}=0.9`$, $`c_1=0.55`$, $`c_2=3.8`$, $`\lambda _0=0.2`$.
Fig. 2 Inelastic head on collision expressed by Eq. (10) when $`M_2`$ suppressed, where $`\lambda _1=0.45+i0.3`$, $`\lambda _2=0.52i0.47`$, $`\stackrel{~}{J}=0.9`$, $`c_1=0.55`$, $`c_2=3.8`$, $`\lambda _0=0.2`$. |
warning/0506/cond-mat0506191.html | ar5iv | text | # Superconductivity of the Ternary Boride Li2Pd3B Probed by 11B NMR
Phys. Rev. B 71, 220505(R) (2005)
## Abstract
We report a <sup>11</sup>B NMR measurement on the recently discovered superconductor Li<sub>2</sub>Pd<sub>3</sub>B. The nuclear spin lattice relaxation rate $`1/T_1`$ shows a well-defined coherence peak just below $`T_c`$ ($`H`$=1.46 T)=5.7 K, and the spin susceptibility measured by the Knight shift also decreases below $`T_c`$. These results indicate that the superconductivity is of conventional nature, with an isotropic gap. Our results also suggest that the $`p`$-electrons of boron and the $`d`$-electrons of palladium that hybridize with boron $`p`$-electrons are primarily responsible for the superconductivity.
Since the discovery of high temperature superconductivity in copper oxides Bednorz , compounds containing transition metal elements have become targets for searching new strongly-correlated superconductors. In fact, superconductivity was discovered in the cobalt oxide Na<sub>0.3</sub>CoO$`{}_{2}{}^{}`$1.3H<sub>2</sub>O (Ref. Takada ), which was found to be of unconventional nature with strong electron correlations in the normal state Fujimoto . Meanwhile, the discovery of superconductivity at 40 K in MgB<sub>2</sub> Akimitsu has generated recurred interest on the physical properties of borides. Very recently, Togano et al found that the ternary metallic compounds containing boron and palladium, Li<sub>2</sub>Pd<sub>3</sub>B, is superconducting at $`T_c`$ 7 K (Ref.Togano ). This compound has a cubic structure with the space group of P4<sub>3</sub>32, containing a distorted Pd<sub>6</sub>B octahedra, which is structurally similar to the superconductor MgCNi<sub>3</sub> ($`T_c`$=8 K) Cava and in some sense also similar to the high-$`T_c`$ copper oxides where the key structure is the oxygen-containing octahedra. Although the physical properties of this new compound are unexplored, it has been proposed that the correlations of Pd $`d`$-electrons may be dominant in the electronic properties and may also be responsible for the superconductivity. Sardar Band calculation has shown that Pd $`d`$-electrons contribute significantly to the density of states at the Fermi level. Chandra
In this Communication, we report the first <sup>11</sup>B nuclear magnetic resonance (NMR) measurement in the superconducting and the normal states. We find that the superconductivity is of Bardeen-Cooper-Schrieffer (BCS) type with an isotropic energy gap. In the normal state, the temperature ($`T`$) dependencce of the nuclear spin-lattice relaxation rate $`1/T_1`$ and the Knight shift due to spin susceptibility, $`K_s`$ satisfies the so-called Korringa relation, $`T_1TK_s^2`$=const, with no signature of electron correlations. Our results suggest that the superconductivity in Li<sub>2</sub>Pd<sub>3</sub>B is phonon mediated.
The samples were prepared by the arc-melting method with starting materials of Li (99.9% purity), Pd (99.95%) and B (99.5%). The two-step arc melting process Togano was used. The starting melt composition of Li:Pd:B=2.1:3:1 was adopted in the light of high vapor pressure of Li. The x-ray diffraction chart shows that the sample is single phase. For NMR measurements, the sample was crushed into powder. Figure 1 shows the ac susceptibility measured using the NMR coil. The $`T_c`$ is about 7.1 K at zero magnetic field.
A standard phase-coherent pulsed NMR spectrometer was used to collect data. The NMR spectra were obtained by fast Fourier transform (FFT) of the spin echo taken at a constant magnetic field of 1.4629 T. The nuclear spin-lattice relaxation rate, $`1/T_1`$, was measured by using a single saturation pulse and fitting the recovery of the nuclear magnetization after the saturation pulse.
Figure 2 shows the spectrum at $`T`$=10 K. The full width at the half maximum (FWHM) of the NMR line is less than 5 kHz. The very sharp transition ensures the high quality of the sample. The temperature dependence of the FWHM is shown in Figure 3. It can be seen that the spectrum broadens below $`T_c(H`$=1.46 T)=5.7 K, at which temperature both $`1/T_1`$ and the Knight shift shows anomaly (see below).
Figure 4 shows the temperature dependence of $`1/T_1`$. In the normal state above $`T_c`$, $`1/T_1`$ varies in proportion to $`T`$, as seen in conventional metals. Just below $`T_c`$ ($`H`$=1.46 T)=5.7 K, however, $`1/T_1`$ increases as $`T`$ is reduced, showing a coherence peak, then decreases exponentially upon further lowing $`T`$. This is the characteristic of $`s`$-wave, isotropic superconductivity. The $`1/T_{1s}`$ in the superconducting state is expressed as
$`{\displaystyle \frac{T_1(T=T_c)}{T_{1s}}}`$ $`=`$
$`{\displaystyle \frac{2}{k_BT_c}}{\displaystyle (N_s(E)^2+M_s(E)^2)f(E)(1f(E))𝑑E}`$ (1)
where $`N_s(E)=N_0E/(E^2\mathrm{\Delta }^2)^{1/2}`$ is the superconducting density of states (DOS) with $`\mathrm{\Delta }`$ being the BCS gap, $`M_s(E)=N_0\mathrm{\Delta }/(E^2\mathrm{\Delta }^2)^{1/2}`$ is the anomalous DOS due to the coherence factor Mac , $`N_0`$ is the DOS in the normal state and $`f(E)`$ is the Fermi function. Both $`M_s(E)`$ and $`N_s(E)`$ diverge at $`E=\mathrm{\Delta }`$, so theoretically $`1/T_1`$ diverges just below $`T_c`$. In real materials, however, broadening of the energy level and some anisotropy in the energy gap will remove the singularity at $`E=\mathrm{\Delta }`$ so that $`1/T_1`$ only shows a peak just below $`T_c`$. Following Hebel Hebel , we convolute $`M_s(E)`$ and $`N_s(E)`$ with a broadening function $`B(E)`$ which is approximated with a rectangular function centered at $`E`$ with a height of $`1/2\delta `$. The solid curve below $`T_c`$ shown in Figure 4 is a calculation with 2$`\mathrm{\Delta }(0)=2.2k_BT_c`$ and $`r\mathrm{\Delta }(0)/\delta `$=2. It fits the experimental data reasonably well. The parameter 2$`\mathrm{\Delta }(0)`$ is substantially smaller than the BCS value of 3.5$`k_BT_c(H)`$. This could be due to the effect of the magnetic field, which usually reduces the gap size Mac .
Figure 5 shows the Knight shift, $`K`$, as a function of temperature. The shift is small and $`T`$-independent within the experimental error in the normal state, but changes abruptly at $`T_c`$ where $`1/T_1`$ shows the coherence peak. Note that the change is positive, namely, $`K`$ increases below $`T_c`$. The shift may be expressed as,
$`K=K_{orb}+K_s`$ (2)
$`K_s=A_{hf}\chi _s`$ (3)
where $`K_{orb}`$ is the contribution due to orbital (Van Vleck) susceptibility which is $`T`$-independent, $`A_{hf}`$ is the hyperfine coupling constant and $`\chi _s`$ is the spin susceptibility. $`K_s`$ may be decomposed into
$`K_s=K_{cp}+K_0`$ (4)
where $`K_{cp}`$ is the shift due to core polarization interaction between the boron $`p`$ electrons and the nuclear spins, which is negative, and $`K_0`$ is due to other interactions including the usual Fermi contact interaction. The result that $`K_s`$ increases below $`T_c`$ indicates that $`K_{cp}`$ is dominant over $`K_0`$. Therefore, the DOS that decreases due to the onset of superconductivity are due to boron $`p`$ electrons, and/or $`d`$ electrons of palladium which hybridize with boron $`p`$ electrons. If we assume that the spin susceptibility vanishes completely, due to singlet pairing, at $`T`$=1.6 K which is well below $`T_c`$, then $`K_{orb}`$ is about 0.085$`\pm `$0.003% and $`K_s`$=-0.012$`\pm `$0.003%. The curve below $`T_c`$ depicts the calculated Knight shift in the superconducting state
$`K_s^{sc}\chi _s^{sc}{\displaystyle N_s(E)(f(E)/E)𝑑E}`$ (5)
using the same gap parameter as in Figure 4.
Finally, when the electrons responsible for the spin lattice relaxation and the Knight shift are $`s`$-electrons, $`T_1TK_s^2`$=$`\frac{\mathrm{}}{4\pi k_B}(\frac{\gamma _e}{\gamma _N})^2`$=2.55$`\times `$10<sup>-6</sup> Sec$``$K (Korringa constant), whre $`\gamma _e`$ and $`\gamma _N`$ are the gyromagnetic ratio of electron and <sup>11</sup>B nuclear spins, respectively. The observed value of $`T_1TK_s^2`$ in the present case is 3.21$`\times `$10<sup>-7</sup> Sec$``$K, which is much smaller than the Korringa constant. This result also supports that the dominant electrons that participate in the relaxation and the Knight shift are $`p`$ electrons, in which case the excess relaxation such as the orbital relaxation is important so that the observed $`T_1TK_s^2`$ deviates from $`\frac{\mathrm{}}{4\pi k_B}(\frac{\gamma _e}{\gamma _N})^2`$.
In summary, through the measurements of <sup>11</sup>B NMR, we find that the spin-lattice relaxation rate $`1/T_1`$ in the newly discovered Pd-containing boride Li<sub>2</sub>Pd<sub>3</sub>B exhibits a well-defined coherence peak just below $`T_c`$ and decreases exponentially with further decreasing temperature. The spin susceptibility as measured by the Knight shift $`K_s`$ also decreases below $`T_c`$, suggesting a singlet spin pairing. In the normal state, the Korringa relation, namely, $`T_1TK_s^2`$=const is obeyed, indicating the lack of electron correlations. Our results indicate that this compound is a superconductor with an isotropic gap and that the superconductivity occurs in the absence of electron correlations.
We thank T. Kato and K. Matano for assistance in some of the measurements. This work was supported in part by a Grant-in-Aid for Scientific Research from MEXT. |
warning/0506/math0506327.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Elliptic curves over finite fields have been largely studied in the literature (for example, see and the references cited there). For an elliptic curve, we have an addition law which endows to the rational points (those defined over the finite field) a group structure. This group is an abelian group of rank 1 or 2. The type of the group is $`(n_1,n_2)`$, i.e. $`E(𝔽_q)/n_1\times /n_2`$ where $`n_2n_1`$ and furthermore $`n_2q1`$.
This structure gives us some information about the torsion points of an elliptic curve, but it is not easy to find them. This problem arises since the computation of the cardinal of an elliptic curve is, in general, very hard. Nevertheless, we have an easy method to test if an elliptic curve has 2-torsion points. If the elliptic curve has equation $`y^2=x^3+Ax+B`$, the abscissas of the points of order 2, are given by the solutions of te cubic $`x^3+Ax+B`$ over the finite field. A study of elliptic curves with 2-torsion points can be found in using elliptic curves in Legendre form.
For an integer $`n2`$, the n-torsion points of an elliptic curve have as abscissas the roots of the classical division polynomials. These can be constructed recurrently:
$$\begin{array}{cc}P_1=1,P_0=0,P_1=1,P_2(x,y)=2y\hfill & \\ P_3=3x^4+6ax^2+12bxa^2\hfill & \\ P_{2k}=\frac{P_k}{2y}(P_{k+2}P_{k1}^2P_{k2}P_{k+1}^2),P_{2k+1}=P_{k+2}P_k^3P_{k+1}^3P_{k1}\hfill & \end{array}$$
But, with these polynomials it is not possible to characterize all elliptic curves with a point of order $`n`$. In , Mestre gives equations for $`\mathrm{\_}0(N)`$, $`N=2,3,5,7,13`$ which parametrizes the pairs $`(E,C)`$ of an elliptic curve $`E`$ and a subgroup $`C`$ of order $`N`$. He gives equations of the modular curve via the j- invariant of the curve. Our purpose is to present a family of the elliptic curves with a rational subgroup of order 3, but in their Weiertrass normal form. The interest of these curves comes from the study of volcanoes of 3- isogenies (see ).
The structure of the paper is the following. In section 2 we present the elliptic curves with, at least, a point of order 3. As it is known (), an elliptic curve with a 3-torsion point over a finite field, admits a model of the form $`y^2+a_1x+a_3y=x^3`$. Nevertheless, we can consider an elliptic curve isomorphic to the former one that depends only on one parameter (after fixing a non cube). We distinguish between the cases when the 3-torsion subgroup is cyclic or not. After this classification, we determine the number of isomorphism classes of elliptic curves with this property.
Section LABEL:nine is dedicated to study when a rational point $`P`$ of order 3 on an elliptic curve (among those mentioned in the previous section) has a point $`Q`$ which *trisects* it (i.e. $`3Q=P`$). This follows easily from and .
Finally, in section 3 we determine a family of representatives of all isomorphic classes of elliptic curves with a rational subgroup of order 3 whose points are not rational. For this purpose, we compare the 3 division polinomial of a curve with that of its twisted curve, to conclude that both factorize in the same way. Thus, the curves in section 2 are in correspondence with the desired curves in this section.
## 2 Elliptic curves with cardinal a multiple of 3
Let $`E`$ be an elliptic curve over a finite field $`𝔽_q`$. It’s well known that $`E`$ admits a model in Weierstrass normal form ():
$$E:y^2+a_1xy+a_3y=x^3+a_2x^2+a_4x+a_6$$
The rational points of order 3 can be obtained from the 3-division polynomial. The roots of this polynomial correspond to the abscissas of the 3-order points <sup>2</sup><sup>2</sup>2Although these abscissas are defined over $`𝔽_q`$, it is possible that the point is not. So, they are points over $`𝔽_{q^2}`$. Let $`E`$ be an elliptic curve with, at least, a point of order 3, we can suppose, without loss of generality, that this point is $`(0,0)`$. In this situation, $`E`$ admits an equation (translating the 3-order point to the origin) of the type
$$y^2+a_1x+a_3y=x^3$$
(1)
The discriminant of an elliptic curve of equation (1) is $`\mathrm{}=(a_1^327a_3)a_3^3`$ and the points $`(0,0)`$ and $`(0,a_3)`$ are of order $`3`$.
Equation (1) can be simplified as follow:
###### Lemma 1
Let $`E`$ be an elliptic curve over $`𝔽_q`$ with equation $`y^2+a_1x+a_3y=x^3`$ and let $`b_0𝔽_q`$ be a fixed non cube. If $`b`$ is a cube, then $`E`$ is isomorphic to a curve of type $`y^2+axy+y=x^3`$. Otherwise, then $`E`$ is isomorphic to $`y^2+axy+b_0y=x^3`$ or $`y^2+axy+b_0^2y=x^3`$.
Moreover, if $`q1mod3`$, all elliptic curves with a point of order 3 admit an equation of the form $`y^2+axy+y=x^3`$.
Proof:
Let us suppose that $`b`$ is a cube. Let $`E`$ be the elliptic curve with equation $`y^2+axy+\alpha ^3y=x^3`$. Replacing $`(x,y)`$ by $`(\alpha ^2x^{},\alpha ^3y^{})`$, it gives an equation of the form $`y^2+(a/\alpha )xy+y=x^3`$.
If $`b`$ is not a cube, let $`\rho `$ be a 3-order primitive root of the unity and let $`\chi `$ be a cubic character over $`𝔽_q`$. It is clear that $`\chi (b)=\rho `$ or $`\rho ^2`$. We can suppose that $`\chi (b_0)=\rho `$ ( otherwise, we take its square). If $`\chi (b)=\rho `$, replacing $`(x,y)`$ by $`(u^2x^{},u^3y^{})`$, $`u=\sqrt[3]{b/b_0}`$; the elliptic curve has equation $`y^2+axy+b_0y=x^3`$. If $`\chi (b)=\rho ^2`$, replace $`(x,y)`$ by $`(u^2x^{},u^3y^{})`$, $`u=\sqrt[3]{b/b_0^2}`$ to get the desired equation.
Finally, over $`𝔽_q`$, $`q1mod3`$, all elements are cubes, therefore, all elliptic curves with a 3-order point, admit an equation of the type $`y^2+axy+y=x^3`$.
### 2.1 Elliptic curves with cyclic 3-torsion subgroup
The 3-division polynomial of an elliptic curve with equation $`y^2+axy+y=x^3`$ is $`P_3(x)=x(3x^3+a^2x^2+3ax+3)`$. If $`q1mod3`$, this polynomial has two roots over $`𝔽_q`$. Nevertheless, only the root 0 corresponds to an abscissa of a 3-order point. The other root corresponds to a 3-order point defined over a quadratic extension of $`𝔽_q`$. This means that there exists a 3-order rational subgroup with no rational points. So, we have the following result:
###### Theorem 2
Let $`𝔽_q`$, $`q1mod3`$ a finite field, then the only elliptic curves with 3-order rational points are isomorphic to a curve with equation $`y^2+axy+y=x^3`$, $`a3`$. Moreover, the 3-torsion group is, in this case, cyclic <sup>3</sup><sup>3</sup>3For all these elliptic curves, the 3-order points are $`(0,0)`$ and $`(0,0)=(0,1)`$.
Taking into account these equations for the elliptic curves, we can count the number of isomorphism classes of elliptic curves over $`𝔽_q`$, $`q1mod3`$, with a point of order three.
###### Proposition 3
Let $`𝔽_q`$, with $`q1mod3`$ a finite field,, there exist $`q1`$ isomorphic classes of elliptic curves with cardinal a multiple of 3. A parametrization of these classes is given by
$$\{y^2+axy+y=x^3|a3\}$$
(2)
Proof:
By theorem 2, let $`E:y^2+axy+y=x^3`$ be an elliptic curve over $`𝔽_q`$, $`q1mod3`$, with a 3-torsion point. The Weierstrass normal form of this curve is $`y^2=x^3+Ax+B`$ where
$$A_a=\frac{a(24+a^3)}{48}B_a=\frac{1}{4}+\frac{a^6}{864}\frac{a^3}{24}.$$
(3)
Over $`q1mod3`$, every element is a cube, and for each one, one only cubic root exists. Therefore, we can only consider the curves $`y^2+\sqrt[3]{a}xy+y=x^3`$ (this equation makes the calculations easier). The j-invariant of this curve is $`16\frac{a(24+a)^3}{a27}`$.
Let $`y^2+\sqrt[3]{b}xy+y=x^3`$ be another elliptic curve of this type. A necessary condition for the elliptic curve to be isomorphic to the previous one is that both possess the same j-invariant. So $`b`$ will be a root of
$$X^3+(a72)X^2+(a^272a+1728)X27\frac{(24+a^3)}{to27}.$$
The only root of this polynomial over the finite field considered is
$$b=\frac{72a}{3}\frac{(36+a)\sqrt[3]{a}}{3\sqrt[3]{27+a}}+\frac{2}{3}\sqrt[3]{27+a}\sqrt[3]{a^2}.$$
(4)
Also, if they are isomorphic, there exists $`u𝔽_{q}^{}{}_{}{}^{}`$ such that $`A_a=u^4A_b`$, $`B_a=u^6B_b`$. Replacing (4) in this second equation, one has that
$$u^6=3\frac{(\sqrt[3]{27+a})^2}{(\sqrt[3]{27+a}\sqrt[3]{a})^2}.$$
But, in this situation, $`3`$ is not a square and thus, isomorphic elliptic curves to $`y^2+\sqrt[3]{a}xy+y=x^3`$ do not exist.
For finite fields $`𝔽_q`$ with $`q1mod3`$, we have to consider all the different types of elliptic curves shown in lemma 1:
$$\{E_a^i:y^2+axy+b_0^iy=x^3;i=0,1,2\}.$$
(5)
where $`\chi (b_0)=\rho `$ ($`\chi `$ is a cubic character over $`𝔽_q`$ and $`\rho `$ a cubic root of the unity).
From the 3-division polynomial, it is easy to see that these curves will have cyclic 3-torsion subgroup if and only if $`\chi (27b_0^ia^3)1`$.
By Lemma 1, if $`E_{a_1}^iE_{a_2}^j`$ then $`i=j`$. We also have:
###### Lemma 4
Let $`E_{a_1}^i,E_{a_2}^i`$ be two elliptic curves over $`𝔽_q`$, $`q1mod3`$ such that $`\chi (27b_0^ia_1^3),\chi (27b_0^ia_2^3)1`$ (both curves possess cyclic 3-torsion subgroup). $`E_{a_1}^iE_{a_2}^i`$ if and only if:
$$a_2\{a_1,\rho a_1,(\rho +1)a_1\}$$
Proof:
It follows from the j-invariants of both curves. These are the only possibilities over $`𝔽_q`$, $`q1mod3`$ and it is easy to prove that these elliptic curves are isomorphic.
As $`\chi (27b_0^ia^3)1`$, determining the number of possible values for $`a`$, is equivalent to determine how many values of $`a`$ exist with $`a^3b_0^i`$ been a cube. The following result will be essential.
###### Lemma 5
Let $`𝔽_q`$ be a finite field, with $`q1mod3`$ and let $`A1mod3`$ be such that $`4q=A^2+27B^2`$ for some $`B`$ (with these conditions $`A`$ is unique).
1. The number of solutions of $`x^3+y^3=1`$ is exactly $`q2+A`$.
2. The number of pairs $`(x,y)`$ such that $`x^3+y^3`$ is not a cube t is $`\frac{(q1)(2q4A)}{3}`$.
Proof:
A proof of part i) can be found in .
For each $`a`$, the number of solutions of $`x^3+y^3=a^3`$ is also $`q2+A`$, except for $`a=0`$, in that case there are $`3q2`$. Therefore, since there exist $`(p1)/3`$ cubes different from 0, there will be $`q^2(\frac{(q1)}{3}(q2+A)+3q2)=\frac{(q1)(2q4A)}{3}`$ pairs $`(x,y)`$ such that $`x^3+y^3`$ is not a cube.
Using this result we can compute the number of isomorphism classes of elliptic curves over $`𝔽_q`$, $`q1mod3`$ with cyclic 3-torsion.
###### Proposition 6
Let $`𝔽_q`$, $`q1mod3`$ a finite field. There are $`(2q+4)/3`$ isomorphic classes of elliptic curves with cyclic 3-torsion subgroup. Let $`b_0𝔽_q`$ be a non cube. A family of representatives is given by:
$$\{y^2+m_axy+b_0^iy=x^3;i=0,1,2,\chi (a^327b_0^i)1\}$$
(6)
with $`m_a=min\{a,\rho a,(\rho +1)a\}`$.
Proof:
First, we compute the number of isomorphism classes of elliptic curves of the form $`y^2+axy+y=x^3`$. This curve has a cyclic 3-torsion if and only if $`a^327`$ is not a cube. By lemma 4, the curves $`y^2+\rho axy+y=x^3`$ and $`y^2+\rho ^2axy+y=x^3`$ are isomorphic to the previous one. Therefore, the number of isomorphism classes of these curves is given by the number of cubes $`a^3`$, such that $`\chi (a^327)=\chi (a^31)1`$.
By the previous lemma, there exist $`q2+A`$ pairs $`(x,y)`$ such that $`x^3+y^3=1`$. To study the number of elements $`x`$ such that both $`x`$ and $`x1`$ are cubes, it is necessary to keep in mind the following facts. The pairs $`(0,1),(0,\rho )`$ and $`(0,\rho ^2)`$ produce the same cube $`0`$. The same fact is true for $`(1,0),(\rho ,0)`$ and $`(\rho ^2,1)`$ and finally for $`(\rho ^ix,\rho ^jy)`$, $`0i,j2`$.
Therefore, there exist $`\frac{q+10+A}{9}`$ cubes $`x^3`$, such that $`x^31`$ is a cube. In consequence, there are $`(2q4+A)/9`$ cubes $`x^3`$ such that $`x^31`$ is not.
In a similar way, it can be shown that there exist $`(4q+16+A)/9`$ isomorphism classes of elliptic curves of the type $`y^2+a^3xy+b_0^iy=x^3;`$ with $`i=1,2,`$ and $`\chi (a^3b_0^i)1`$ and the result follows after adding this number to the previous one.
### 2.2 Elliptic curves with non cyclic 3-torsion subgroup
As shown in the previous section, a necessary condition for an elliptic curve to have non cyclic 3-torsion subgroup is that it be defined over $`𝔽_q`$ with $`q1mod3`$. If $`E`$ is defined by the equation $`y^2+axy+by=x^3`$, it is also necessary that $`\chi (a^327b)1`$. We can change the equation to distinguish between those with non cyclic subgroup as follows:
###### Proposition 7
The family of elliptic curves
$$E_a:y^2+(3a1)xy+a(\rho 1)(a\frac{\rho +1}{3})y=x^3$$
such that $`a(a(\rho +1)/3)(a\rho /3)0`$ corresponds to all the elliptic curves defined over $`𝔽_q`$ with non cyclic 3-torsion subgroup.
Proof:
Let $`y^2+axy+by=x^3`$ be an elliptic curve over $`𝔽_q`$ with non cyclic 3-torsion subgroup. It is clear that (by construction) one of the points of order 3 is the origin. Let $`(x_0,y_0)`$ be another point of order 3 which does not belong to $`<(0,0)>`$ (that is $`x_00`$). Let $`y=\lambda x+\mu `$ be the tangent line to the curve in $`(x_0,y_0)`$ ($`\lambda 0`$ because this point is not the origin).
The change of variables $`(x,y)(\lambda ^2x,\lambda ^3y)`$ transforms the equation of the elliptic curve in to another of the same type. The point of order 3 is now $`(x_0/\lambda ^2,y_0/\lambda ^3)`$ and the tangent line at this point has slope $`1`$.
The result follows easily by considering the conditions in the cubic such that both $`y=0`$ and $`y=x+u`$ have a triple intersection point with the cubic (this can be found in ).
Let us see now when two elliptic curves $`E_a`$ and $`E_b`$ are isomorphic. It is clear that $`a\{0,\rho /3,(\rho +1)/3\}`$, otherwise it would be a singular cubic. Using the expressions for the j-invariants and the equations in Weierstrass normal form of these curves we can prove that:
###### Lemma 8
Let $`𝔽_q`$ be a finite field with $`q1mod3`$ and let $`E_a`$ and $`E_b`$ be two elliptic curves over $`𝔽_q`$. The equations for $`E_a`$ and $`E_b`$ are
$`y^2+(3a1)xy+a(\rho 1)(a{\displaystyle \frac{\rho +1}{3}})y=x^3`$
$`y^2+(3b1)xy+b(\rho 1)(b{\displaystyle \frac{\rho +1}{3}})y=x^3`$
respectively (such that $`a,b\{0,\rho /3,(\rho +1)/3\}`$). $`E_aE_b`$ if and only if
$$\begin{array}{cc}bG_a=\{a,\frac{a(1+\rho )}{(3a\rho )},\frac{a\rho }{(3a1\rho )},\frac{1}{9a},\frac{(1+\rho )(3a1\rho )}{3},\frac{\rho (3a\rho )}{9a},\frac{\rho (3a1\rho )}{3(3a\rho )},\hfill & \\ & \\ \frac{\rho }{3(3a1\rho )},\frac{(1+\rho )(3a\rho )}{3(3a1\rho )},\frac{(1+\rho )(3a1\rho )}{9a},\frac{(1+\rho )}{3(3a\rho )},\frac{\rho (3a\rho )}{3}\}.\hfill & \end{array}$$
$`G_a`$ has a group structure and it is isomorphic to $`/2\times /6`$. $`G_a`$ acts over $`𝔽_q\{0,\rho /3,(\rho +1)/3\}`$. Then, the number of isomorphism classes of elliptic curves with non cyclic coincides whit the number of orbits under this action. To compute this number we use the well-known Burnside formula
$$\mathrm{}\text{ of orbits}=\frac{1}{|G_a|}\underset{gG_a}{}|N^g|$$
where $`N^g`$ is the set of elements fixed by $`g`$. We have:
| $`g`$ | $`N^g`$ | $`g`$ | $`N^g`$ |
| --- | --- | --- | --- |
| $`g`$ | $`𝔽_q\{0,\rho /3,(\rho +1)/3\}`$ | $`\frac{\rho (3a1\rho )}{3(3a\rho )}`$ | $`\frac{\rho +\sqrt{\rho }}{3},\frac{\rho \sqrt{\rho }}{3}`$ |
| $`\frac{a(1+\rho )}{(3a\rho )}`$ | $`\frac{1+2\rho }{3}`$ | $`\frac{\rho }{3(3a1\rho )}`$ | $`\frac{1}{3}`$ |
| $`\frac{a\rho }{(3a1\rho )}`$ | $`\frac{1+2\rho }{3}`$ | $`\frac{(1+\rho )(3a\rho )}{3(3a1\rho )}`$ | $`\frac{\rho +1+\rho \sqrt{1}}{3},\frac{\rho +1\rho \sqrt{1}}{3}`$ |
| $`\frac{1}{9a}`$ | $`\frac{\sqrt{1}}{3},\frac{\sqrt{1}}{3}`$ | $`\frac{(1+\rho )(3a1\rho )}{9a}`$ | $`\frac{1}{3}`$ |
| $`\frac{(1+\rho )(3a1\rho )}{3}`$ | $`\frac{1}{3}`$ | $`\frac{(1+\rho )}{3(3a\rho )}`$ | $`\frac{1}{3}`$ |
| $`\frac{\rho (3a\rho )}{9a}`$ | $`\frac{1}{3}`$ | $`\frac{\rho (3a\rho )}{3}`$ | $`\frac{1+2\rho }{9}`$ |
Considering whether these elements belong to $`𝔽_q\{0,\frac{\rho }{3},\frac{(\rho +1)}{3}\}`$ and the Burnside formula, we have the following result.
###### Proposition 9
Let $`𝔽_q`$ be a finite field with $`q1mod3`$. Over $`𝔽_q`$ there exist $`\frac{q+12(qmod12)}{12}`$ isomorphic classes of elliptic curves with non cyclic 3-torsion subgroup.
## 3 Elliptic curves with rational subgroups of order 3 containing no rational points
In this section we are interested in elliptic curves that have no rational points of order three, but they have a rational subgroup of order 3. These curves appear when we consider rational isogenies of degree three. This is because the points of this rational subgroup are not invariant under the action of the Frobenius endomorphism, whereas the subgroup is.
Let $`E:y^2=x^3+Ax+B`$ be an elliptic curve defined over a finite field $`𝔽_q`$ and let $`GE(𝔽_q)`$ be a subgroup of order $`3`$ whose points are not rational (except the point of the infinity). Let $`\sigma :𝔽_q𝔽_q`$, be the Frobenius endomorphism: $`\sigma ((x,y))=(x^q,y^q)`$. If $`PG`$, then $`\sigma (P)=P`$or $`P`$. If it coincides with $`P`$, it is clear that the point $`P`$ is rational, which is absurd. Therefore, $`\sigma (P)=P`$ and hence
$$(x^q,y^q)=(x,y).$$
From this, we deduce that the abscissa of $`P`$ should be rational and in consequence the ordinate should be defined over $`𝔽_{q^2}`$. Then $`\mathrm{}(E(𝔽_{q^2}))`$ is a multiple of 3.
By Weil theorem, if $`\mathrm{}(E(𝔽_q))=m(𝔽_q)=q+1t`$, then $`\mathrm{}(E(𝔽_{q^2}))=m(𝔽_{q^2})=q^2+1t^2+2q`$. If $`q1mod3`$, $`m(𝔽_{q^2})`$ is a multiple of 3 if and only if $`m(𝔽_q)`$ is . Therefore, all the elliptic curves with rational subgroups of order 3 over $`𝔽_q`$, $`q1mod3`$, have been already studied in section 2.
Otherwise, if $`q1mod3`$, $`m(𝔽_{q^2})`$ is a multiple of 3 if and only if $`t1,2mod3`$. In the first case, the curve defined over $`𝔽_q`$ has cardinal multiple of 3 (already studied in section 2). In the second case, $`m(𝔽_q)`$ is not a multiple of 3. In consequence, we must only study elliptic curves with points of order 3 defined over $`𝔽_{q^2}`$ but not over $`𝔽_q`$. Moreover, this implies that there exist a one-to-one correspondence between curves with rational points of order 3 and those with 3-order subgroups with no rational points (ones are the twisted of the others).
The 3-division polynomial (the one whose roots are the abscissas of the points of order three) of $`E`$ is
$$P_{A,B}=x^4+2Ax^2+4Bx\frac{A^2}{3}.$$
The following theorem gives us some information about the factorization of $`P_{A,B}`$ over $`𝔽_q`$, $`q1mod3`$.
###### Theorem 10 ()
Let $`p`$ be an odd prime and let $`f(x)`$ be a monic polynomial of degree $`n`$ and discriminant $`D`$ without multiple roots. Let $`f(x)=f_1(x)f_2(x)\mathrm{}f_r(x)modp`$ be the factorization of $`f`$ over $`𝔽_p`$. Then $`nrmod2`$ if and only if $`(D/p)=1`$.
The discriminant of $`P_{A,B}`$ is $`2^8(4A^3+27B^2)^2/27`$. Over $`𝔽_q`$, $`q1mod3`$, $`1/3`$ is a square. So, one has that $`P_{A,B}`$ factors into two polynomials or it has all its roots in $`𝔽_q`$. If it factors into two polynomials, the possibilities for its degree are $`(2,2)`$ or $`(3,1)`$. In the latter case, $`P_{A,B}`$ has a root in $`𝔽_q`$.
Using the results of , these three types can be distinguished in the following way. If $`\chi (2(27B^2+4A^3))1`$, then $`P_{A,B}=(xx_0)f_2(x)`$ with $`f_2(x)`$ an irreducible polynomial of degree 3. Otherwise, we define $`y_0=2\sqrt[3]{2(27B^2+4A^3)}/3`$. If $`y_0+16A/3`$ and $`\rho y_0+16A/3`$ are both squares, then $`P_{A,B}`$ factors completely over $`𝔽_q`$. Otherwise, $`P_{A,B}`$ factors in two irreducible quadratic polynomials.
With these results, we have that the type ot factorization of $`P_{A,B}`$ coincides with the that of $`P_{t^2A,t^3B}`$, where $`t`$ is a non quadratic residue. That corresponds to the 3-division polynomial of the twisted curve.
As we have shown previously, elliptic curves with rational points of order three are the twisted of those with subgroups of order 3 with no rational points. Thus, if $`y^2=x^3+Ax+B`$ is an equation for an elliptic curve with 3-torsion points, then $`y^2=x^3+t^2Ax+t^3B`$ ($`t`$ a non quadratic residue), is a curve with a rational subgroup of order 3 with no rational points. As we have studied in section 2 the curves of the first type, determine those of the second type from their Weierstrass form.
###### Proposition 11
Let $`𝔽_q`$, $`q1mod3`$ be a finite field and $`b_0𝔽_q`$, such that $`\chi (b_0)=\rho `$. Let $`t`$ be a non quadratic residue. There exist $`(2q+4)/3`$ isomorphism classes of elliptic curves over $`𝔽_q`$ with only one rational subgroup of order three, whose points are defined over $`𝔽_{q^2}𝔽_q`$. A family of representatives is given for $`y^2=x^3+t^2Ax+t^3B`$ with
$$A_a=\frac{a^3(a^924b_0^i)}{48}B_a=\frac{b_0^{2i}}{4}+\frac{a^{18}}{864}\frac{a^9b_0^i}{24}.$$
$$i=0,1,2\chi (a^327b_0^i)1$$
Proof:
Elliptic curves of equation (6) $`\{y^2+a^3xy+b_0^iy=x^3;i=0,1,2,\chi (a^327b_0^i)1\}`$ correspond to those curves with a non cyclic 3-torsion subgroup whose points are defined over $`𝔽_q`$. Therefore, their twisted curves will have a rational 3-torsion subgroup of points defined in a quadratic extension.
The Weierstrass normal form of elliptic curves with equation (6) is
$$y^2=x^3+\left(\frac{a^3(a^924b_0^i)}{48}\right)x+\frac{b_0^{2i}}{4}+\frac{a^{18}}{864}\frac{a^9b_0^i}{24}.$$
In consequence, the equation of the twisted curve coincides with the desired one.
Let $`y^2=x^3+t^2Ax+t^3B`$ and $`y^2=x^3+t^2A^{}x+t^3B^{}`$ be two of these curves. They will be isomorphic over $`𝔽_q`$, if and only if there exists $`u𝔽_{q}^{}{}_{}{}^{}`$ such that $`u^4=A/A^{},u^6=B/B^{}`$. Or equivalently, if and only if the curves $`y^2+a^3xy+b_0^iy=x^3,y^2+a^3xy+b_0^iy=x^3`$ are isomorphic. The number of isomorphism classes follows from proposition 6.
Two curves of this type will be isomorphic if and only if
$$a_2\{a_1,\rho a_1,(\rho +1)a_1\}.$$
We can prove, in the same way, the following result:
###### Proposition 12
Let $`𝔽_q`$ be a finite field with $`q1mod3`$. Over $`𝔽_q`$ there exist $`\frac{q+12(qmod12)}{12}`$ isomorphism classes of elliptic curves with 4 rational subgroups of order 3 composed by points defined over $`𝔽_{q^2}`$. A collection of representatives is given by $`y^2=x^3+t^2Ax+t^3B`$ with $`t`$ a non quadratic residue and
$$\begin{array}{c}A=\frac{(9a12\rho )(3a12\rho )(3a1)(3a+1)}{144}\hfill \\ B=\frac{(1+9a^2)(9a^26a6\rho to1)(9a^26\rho to1)}{864}\hfill \end{array}$$
Two curves of this type are isomorphic (the same as it happened with those with 8 rational points of order 3) if and only if
$$\begin{array}{cc}bG_a=\{a,\frac{a(1+\rho )}{(3a\rho )},\frac{a\rho }{(3a1\rho )},\frac{1}{9a},\frac{(1+\rho )(3a1\rho )}{3},\frac{\rho (3a\rho )}{9a},\frac{\rho (3a1\rho )}{3(3a\rho )},\hfill & \\ & \\ \frac{\rho }{3(3a1\rho )},\frac{(1+\rho )(3a\rho )}{3(3a1\rho )},\frac{(1+\rho )(3a1\rho )}{9a},\frac{(1+\rho )}{3(3a\rho )},\frac{\rho (3a\rho )}{3}\}.\hfill & \end{array}$$ |
warning/0506/physics0506170.html | ar5iv | text | # The Fundamental Diagram of Pedestrian Movement Revisited
## 1 Introduction
Pedestrian dynamics has a multitude of practical applications, like the evaluation of escape routes or the optimization of pedestrian facilities, along with some more theoretical questions . Empirical studies of pedestrian streams can be traced back to the year 1937 . To this day a central problem is the relation between density and flow or velocity. This dependency is termed the fundamental diagram and has been the subject of many investigations from the very beginning . This relation quantifies the capacity of pedestrian facilities and thus allows e.g. the rating of escape routes. Furthermore, the fundamental diagram is used for the evaluation of models for pedestrian movement , and is a primary test whether the model is suitable for the description of pedestrians streams .
The velocity-density relation differs for various facilities like stairs, ramps, bottlenecks or halls. Moreover one has to distinguish between uni- or bi-directional streams. The simplest system in this enumeration is the uni-directional movement of pedestrians in a plane without bottlenecks. In this context the fundamental diagram of Weidmann is frequently cited. It is a part of a review work and the author summarized 25 different investigations for the determination of the fundamental diagram. Apart from the fact, that with growing density the velocity decreases, the relation shows a non-trivial form. Weidmann notes that different authors choose different approaches to fit their data, indicating that the dependency is not completely analyzed. A multitude of possible effects can be considered which may influence the dependency. For instance we refer to passing maneuvers, internal friction, self-organization phenomena or ordering phenomena like the ‘zipper’ effect . A reduction of the degrees of freedom helps to restrict possible effects and allows an improved insight to the problem. To exclude in a natural way the influences of passing maneuvers, internal friction or ordering phenomena, we choose a one-dimensional system. At this we restrict the investigation on the velocity-density relation for the normal and not for the pushy or panic movement of pedestrians. Similar experiments with the focus on the time gap distribution and results for the clogging in pushy crowds can be found in .
Furthermore, we test a procedure for automatic recording of pedestrian flow characteristics. This method uses stereo video processing and allows the determination of trajectories of individual persons. We present preliminary results on measurement range and accuracy of this method.
## 2 Movement in a plane
In the field of pedestrian dynamics there are two different ways to quantify the capacity of pedestrian facilities. Either the relation between density $`\rho `$ and velocity $`v`$ or the relation between density and flow $`\mathrm{\Phi }=\rho v`$ is represented. Following Weidmann we choose the presentation of density and velocity.
Figure 1 shows the empirical velocity-density relation for pedestrian movement in a plane according to Weidmann<sup>1</sup><sup>1</sup>1Note that Weidmann’s combination does not distinguish between uni- (e.g. ) and bi-directional (e.g. ) movement. The maximal density in this diagram is $`5.4/m^2`$, though some authors have found higher densities. The slope varies for different density-domains indicating diverse effects which reduce the velocity. We discuss possible causes for this slope-variation by means of the Level of Service concept (LOS) .
At low densities there is a small and increasing decline of the velocity. The velocity is mostly determined by the individual free velocity of the pedestrians. Passing maneuvers for keeping the desired velocity are possible. The decrease is caused by the passing maneuvers.
The velocity decrease is nearly linear with growing density. Due to the reduction of the available space, passing maneuvers of slower pedestrians are hardly feasible and the possibility to choose the desired velocity is restricted. At least at uni-directional streams the available space is large enough to avoid contacts with other pedestrians. Thus the internal friction can only be of negligible influence.
The linear decrease of the velocity ends and the curvature changes. For growing density the velocity remains nearly constant. Contacts with other pedestrians are hardly avoidable. While the internal friction increases compared with domain II, the reduction of the velocity diminishes.
The velocity declines rapidly. The available space is strongly restricted. Internal friction may be a determining factor. There has to be a maximal density because the pedestrians behave like hard bodies.
As pointed out, there are several hints to effects which can influence the reduction of the velocity at different densities. But it is not clear which ‘microscopic’ properties of pedestrian movement lead to the nearly linear decrease of the velocity and to a slower decrease at high densities. Particularly the possible influence of the internal friction is contradictory to the slope of the velocity-density relation. We have to clarify the influences of collective behavior like marching in lock-step or self-organization phenomena like the ’zipper’ effect, which can be observed at bottlenecks .
## 3 Single-file movement
### 3.1 Description of the experiment
Our target is the measurement of the relation between density and velocity for the single-file movement of pedestrians. To facilitate this with a limited amount of test persons also for high densities and without boundary effects, we choose a experimental set-up similar to the set-up in .
The corridor, see Figure 2, is build up with chairs and ropes. The width of the passageway in the measurement section is $`0.8m`$. Thus passing is prevented and the single-file movement is enforced. The width of the corridor is not important as long it does not impede the free movement of the arms and on the other hand does not allow the formation of two - possible interleaving - lanes. The circular guiding of the passageway gives periodic boundary conditions. To reduce the effects of the curves on the measurement, we broaden the corridor in the curve and choose the position of the measured section in the center of the straight part of the passageway. The length of the measured section is $`l_m=2m`$ and the whole corridor is $`l_p=17.3m`$. The experiment is located in the auditorium Rotunde at the Central Institute for Applied Mathematics (ZAM) of the Research Centre Jülich. The group of test persons is composed of students of Technomathematics and staff of ZAM. To get a normal and not a pushy movement the test persons are instructed not to hurry and omit passing. This results in a rather relaxed conduct, i. e. the resulting free velocities are rather low. Most real life situations will show higher free velocities and more pushing. To enable measurements at different densities we execute six cycles with $`N=1,15,20,25,30,34`$ numbers of test persons in the passageway. For the cycle with $`N=1`$ every person passes alone through the corridor. For the other cycles we first distribute the persons uniformly in the corridor. After the instruction to get going, every person passes the passageway two to three times. After that we open the passageway and let the test persons get out. We observed that the test persons walked straight through the measurement section and follow each other up to a few centimeter of lateral variation. This does not mean that they walked exactly on the same path but lateral fluctuations are highly correlated. This observation is reinforced by the documentation of marching in lock-steps, which will be discussed later.
### 3.2 Measurement set-up
For the measurement of the flow characteristics we use both a manual and an automatic procedure. The manual procedure is based on standard video recordings with a DV camera (PAL format, 25 fps) of the measured section. These recordings are prepared on a computer to show time, frame number and the limits of the measured section. After the preparation the recordings were analyzed frame-wise, thus the accuracy of the extracted time is $`0.04s`$. To minimize the errors for extracting time data, the collected time is transferred by cut and paste between the video editing system and the spread sheet. Figure 3 shows one frame of the cycle with $`N=20`$ persons. For every test person $`i`$ we collect the entrance time (of the ear) in the measurement section $`t_i^{in}`$ and the exit time $`t_i^{out}`$. To ease the assignment of times the test persons carry numbers.
Additionally we test an automated procedure, based on a commercial system of Point Grey Research . The system is composed of the Bumblebee stereo vision camera and the software packages Digiclops, Triclops and Censys3d. The software uses stereo recordings for detection and tracking of peoples. The resulting trajectories allow the analysis of pedestrian movement in space and time. For this measurement we use a bumblebee BW-BB-20. The analysis device is an IBM ThinkPad T30 with an Intel Pentium 4 Mobile CPU (512 MB RAM) and the operating system WIN XP SP1. The data transfer from the camera to the analysis device is realized via FireWire and a PCMCIA FireWire card. Following the recommendation of the manufactures we decide to process the data directly without storing the pictures on a hard disk. In our measurement setup we are bounded on a transfer rate to the analysis device of about 2 times 20 fps with a resolution of 320 x 240 pixels.
### 3.3 Data analysis
The manual analysis uses the entrance and exit times $`t_i^{in}`$ and $`t_i^{out}`$. These two times allow the calculation of the individual velocities $`v_i^{man}=l_m/(t_i^{out}t_i^{in})`$ and the momentary number $`n(t)`$ of persons at time $`t`$ in the measured section. The concept of a momentary ‘density’ in the measurement region is problematic because of the small (1-5) number of persons involved. $`\stackrel{~}{\rho }^{man}(t)=n(t)/l_m`$ jumps between discrete values. For a better definition we choose $`\rho ^{man}(t)=_{i=1}^N\mathrm{\Theta }_i(t)/l_m`$, where $`\mathrm{\Theta }_i(t)`$ gives the ‘fraction’ to which the space between person $`i`$ and person $`i+1`$ is inside.
$$\mathrm{\Theta }_i(t)=\{\begin{array}{cc}\hfill \frac{tt_i^{in}}{t_{i+1}^{in}t_i^{in}}:& t[t_i^{in},t_{i+1}^{in}]\hfill \\ \hfill 1:& t[t_{i+1}^{in},t_i^{out}]\hfill \\ \hfill \frac{t_{i+1}^{out}t}{t_{i+1}^{out}t_i^{out}}:& t[t_i^{out},t_{i+1}^{out}]\hfill \\ \hfill 0:& \text{otherwise}\hfill \end{array}$$
(1)
In Figure 3 the ‘fractions’ are $`\mathrm{\Theta }_6(t)0.6`$, $`\mathrm{\Theta }_7(t)=1`$, $`\mathrm{\Theta }_8(t)0.8`$, $`\mathrm{\Theta }_9(t)=0`$, resulting in $`\rho (t)1.2m^1`$. This is 3 persons per $`2.5m`$, which is about the distance between person number 6 and number 9. The time average of $`\stackrel{~}{\rho }^{man}`$ and $`\rho ^{man}`$ is almost the same.
Figure 4 shows the development of the density and velocity of individual persons in time. At the beginning the pedestrians do not react simultaneously and they have to tune their movement. After the tuning phase one sees slow moderate size fluctuations of the density and the velocity. It becomes obvious that these ‘microscopic’ fluctuations are correlated. After the opening of the passageway the density declines and as a consequence the velocity grows. To consider this correlation in the analysis we regard a crossing of an individual pedestrian $`i`$ with velocity $`v_i^{man}`$ as one statistical event, which is associated to the density $`\rho _i^{man}`$. While $`\rho _i^{man}`$ is the mean value of the density during the time-slice $`[t_i^{in},t_i^{out}]`$. For this analysis we exclude the part of the data where the influence of the tuning phase and the rearrangement of the boundary conditions is explicit.
Figure 5 shows the distribution of the events $`(v_i^{man},\rho _i^{man})`$ of the cycles with $`N=15,20,25,30`$ and $`34`$. The cycles with $`N=15,20`$ and $`25`$ result in clearly separated areas of densities and velocities. These areas blend for the cycles with $`N=30`$ and $`34`$.
The automated procedure uses the trajectories for the analysis. The accuracy of the trajectory measurement depends on many factors, like the size of the covered area, the capability of the analysis device and the settings of parameters in software Censys3d for people detection and tracking. In test-runs we optimize these parameters for normal and not too crowded conditions. But during the experiment we realized that the system is too sensitive and with growing density the mismeasurements increase. For example we observe the loss of trajectories or persons with two trajectories, see figure 6. Due to the mismeasurements the determination of densities or other microscopic values like the distance between the pedestrians is not reliable. The resulting trajectories are thus only appropriate for the calculation of the mean value of the velocity. For this purpose we determine for every crossing pedestrian $`i`$ the first $`(x_i^f,t_i^f)`$ and the last point $`(x_i^l,t_i^l)`$ of the trajectory in the measured section and calculate the velocity through $`v_i^{aut}=(x_i^lx_i^f)/(t_i^lt_i^f)`$.
In the next step we test if the mismeasurements lead to systematic errors in the determination of the velocities and cross-check the data from the manual analysis. For this we exclude the data where the influence of the starting phase and opening of the passageway are apparent. To reduce the influence of mismeasurements we take into account only these trajectories, which last for at least one second.
For the test we determined the following values. Based on the individual velocities $`v_i^{aut}`$ and $`v_i^{man}`$ we calculate the mean value $`v^{man}`$ and $`v^{aut}`$ over all individual velocities for the different cycles. Furthermore we determine the mean value of the density $`\rho ^{man}`$ during one cycle based on the manual analysis and compare them with the densities calculated through $`\rho =N/l_p`$, where $`l_p`$ is the length of the passageway. The values are summarized in the following table.
Aside from local fluctuations, possible reasons for the deviations of the densities $`\rho ^{man}`$ and $`\rho `$ are the influence of the curve and the broadening of the passageway in the curve. The mean velocities calculated automatically agree roughly with the velocities extracted from the manual procedure. We get a very good agreement for medium densities $`N=15,20`$. Notable deviations arise at high velocities and high densities. The mismeasurements caused by loosing and picking up of trajectories increase the fluctuations ($`\sigma `$). The deviation of the mean value for the cycle with $`N=1`$ from the literature value $`v_{free}=1.34m/s`$ according to Weidmann, can be explained by the instruction to the test persons not to hurry.
## 4 Results
In the following we focus on the distance between the pedestrians. For the single-file movement the distance to the next close-by pedestrian can be regarded as the required length $`d`$ of one pedestrian to move with velocity $`v`$. Considering that in a one-dimensional system the harmonic average of this quantity is the inverse of the density, $`d=1/\rho `$, one can investigate the relation between required length and velocity by means of the velocity-density relation for the single-file movement.
Figure 7 shows the dependency between required length and velocity. We tested several approaches for the function $`d=d(v)`$ and found that a linear relationship with $`d=0.36+1.06v`$ gives the best fit to the data. According to the step length is a linear function of the velocity<sup>2</sup><sup>2</sup>2Lower average velocities arise from a lower step frequency. only for $`v0.5m/s`$. Thus it is surprising, that the linearity for entire distance holds even and persists for velocities smaller than $`0.5m/s`$. Possible explanations will be discussed later.
To compare the relation between velocity and density of the single-file movement $`(1d)`$ with the movement in a plane $`(2d)`$, we have to transform the line-density to a area-density.
$$\rho _{1d2d}=\rho _{1d}/c(v)\text{ }c(v)=\alpha +\beta v$$
(2)
The correction term $`c(v)`$ is introduced to take into account that the lateral required width is as well as the longitudinal required length $`d`$ a function of the velocity. For a first approximation we assume a linear dependency. According to the mean value of the width of the human body is $`0.46m`$ and thus we set $`\alpha =0.46m`$. The parameter $`\beta `$ considers that with increasing velocity the lateral required width grows.
The comparison of the relation between velocity and density for the single-file movement with the movement in a plane according to Weidmann shows a surprising conformity, see Figure 8. So far the value of $`\beta `$ is unknown, but the figure shows, that for a large domain of reasonable values we have an obvious accordance between the fundamental diagrams. For the lower bound of $`\beta =0.05s`$ one gets a lateral required width at $`v=1.3m/s`$ of $`c=0.5m`$, which is to definitely too small. For the upper bound we choose $`\beta =0.3s`$ which results in $`c(1.3m/s)=0.9m`$ which is certainly too high. But it becomes apparent, that the qualitative shape of the velocity-density relations and the magnitude of the velocities agree beetween the upper and the lower bound of $`\beta `$. This conformance indicates that two-dimensional specific properties, like internal friction and other lateral interferences, have no strong influence on the fundamental diagram at least at the density domains considered.
Instead, the visual analysis of the video recordings suggests that the following ‘microscopic’ properties of pedestrian movement determine the relation between velocity and density. At intermediate densities and velocities the step length is reduced with increasing density. The distance to the pedestrian in front is related to the step length as well as to the safety margin to avoid contacts with the pedestrian in front. Both, step length and safety margin are connected with the velocity. At high densities and small velocities we observed that small groups pass into marching in lock-step, see Figure 9. Furthermore the utilization of the available place is optimized. This is achieved by some persons setting their feet far right and left of the line of movement, giving some overlap in the space occupied with the pedestrian in front. While at intermediate densities and relative high velocities the pedestrians are concentrated on their movement, this concentration is reduced at smaller velocities and leads to a delayed reaction on the movement of the pedestrian in front. The marching in lock-steps and the optimized utilization of the available space, which compensate the slower step frequency, are possible explanation that the linearity between the required lenght and the velocity holds even, see Figure 7.
## 5 Summary and outlook
In the investigation presented we determine the fundamental diagram for the single-file movement of pedestrians under laboratory conditions. The data are appropriate to test, if microscopic models are able to reproduce the empirical relation between velocity and density in the simplest system.
The test of the automated procedure shows that this method is in principle capable to measure characteristics of pedestrian movement. The mean values of the velocities for different densities acquired automatically are comparable with those of the manual data analysis. To facilitate the measurement of microscopic characteristics we plan to increase the resolution of the stereo recordings and the transfer rate to the analysis system.
The comparison of the velocity-density relation for the single-file movement with the literature-data for the movement in a plane shows a surprising agreement for $`1m^2<\rho <5m^2`$. The conformance indicates that the internal friction and other lateral interferences, which are excluded in the single-file movement, have no influence on the relation at the density domains considered. The visual analysis of the video recording give hints to possible effects, like the self-organization through marching in lock-step, the optimized utilization of the available space at low velocities and the velocity dependence of step-length and safety margin. The investigation of the dependency between the required length of one pedestrian and velocity indicates a linear relation. For the domain $`0.1m/s<v<1.0m/s`$ we obtain $`d=0.36+1.06v`$. The investigation of the interplay between the self-organization effects and required length and thus a detailed quantification of these effects is necessary.
For real life situation including bottlenecks and changes in direction, and for higher densities, however, two-dimensional effects can certainly not be disregarded. Moreover, for most laboratory experiments, the emotional situation of the participants is relaxed, so some effects of real life will not be shown, as already observed in . Further studies will be required, and some of these will be vastly easier after the improvement of the automated procedure.
Acknowledgments
We thank Patrick Hartzsch for extracting the times from the video, Dr. Wolfgang Meyer for the advice with statistics, Oliver Passon for careful reading and Thomas Lippert for discussions. |
warning/0506/hep-ph0506257.html | ar5iv | text | # Neutrino Physics
## 1 The Fermi Theory
The history of the weak interactions dates back to the discovery of radioactivity by Becquerel in 1896 . In particular, $`\beta `$ decay, in which a nucleus emits an electron and increases its charge, apparently violated the conservation of energy (as well as momentum and, as we now understand, angular momentum). In 1931 Pauli postulated that a massless, chargeless, essentially non-interacting particle that he named the “neutron” (later renamed the neutrino by Fermi) was also emitted in the process and carried off the missing energy. Pauli’s hypothesis was verified around 1953 when the electron-type neutrino (actually the anti-neutrino $`\overline{\nu }_e`$) produced in a reactor was observed directly by its rescattering by Reines and Cowen. The second (muon-type) neutrino, $`\nu _\mu `$, associated with the $`\mu `$ in its interactions, was detected by its rescattering to produce a muon in 1962 by Lederman, Schwartz, and Steinberger at Brookhaven. The third charged lepton, the $`\tau `$, was discovered at SLAC in 1975. There was ample indirect evidence from the weak interactions of the $`\tau `$ that an associated neutrino, $`\nu _\tau `$, must exist, but it was not observed directly until 2000 at Fermilab.
In 1934, Enrico Fermi developed a theory of $`\beta `$ decay. The Fermi theory is loosely like QED, but of zero range (non-renormalizable) and non-diagonal (charged currents).
The Hamiltonian (see Figure 1) is
$$H\frac{G_F}{\sqrt{2}}J_\mu ^{}J^\mu ,$$
(1)
where $`J_\mu `$ is the charged current,
$$\begin{array}{ccccc}J_\mu ^{}\hfill & \hfill & \overline{p}\gamma _\mu n& +\hfill & \overline{\nu }_e\gamma _\mu e^{}\\ & & [np& ,\hfill & e^{}\nu _e]\end{array},$$
and $`G_F`$ is the Fermi constant, with the modern value $`G_F=1.16637(1)\times 10^5`$ GeV<sup>-2</sup>. The Fermi theory was later rewritten in terms of quarks, and modified to include: $`\mu ^{}`$ and $`\tau ^{}`$ decays; strangeness changing transitions (Cabibbo); the heavy quarks $`c`$, $`b`$, and $`t`$; quark mixing and CP violation (the Cabibbo-Kobayashi-Maskawa matrix); vector ($`V`$) and axial vector ($`A`$) currents (i.e, $`V\gamma _\mu `$ is replaced by $`VA\gamma _\mu (1\gamma ^5)`$), which implies parity violation (Lee-Yang; Wu; Feynman-Gell-Mann); and $`\nu `$ mass and mixing.
The Fermi theory gives an excellent description (at tree-level) of such processes as:
* Nuclear/neutron $`\beta `$decay: $`npe^{}\overline{\nu }_e`$.
* $`\mu ^{}`$, $`\tau ^{}`$ decays: $`\mu ^{}e^{}\overline{\nu }_e\nu _\mu `$; $`\tau ^{}\mu ^{}\overline{\nu }_\mu \nu _\tau ,\nu _\tau \pi ^{},\mathrm{}`$.
* $`\pi `$, $`K`$ decays: $`\pi ^+\mu ^+\nu _\mu ,\pi ^0e^+\nu _e`$; $`K^+\mu ^+\nu _\mu ,\pi ^0e^+\nu _e,\pi ^+\pi ^0,\mathrm{}`$.
* Hyperon decays: $`\mathrm{\Lambda }p\pi ^{}`$$`\mathrm{\Sigma }^{}n\pi ^{}`$, $`\mathrm{\Sigma }^+\mathrm{\Lambda }e^+\nu _e,\mathrm{}`$.
* Heavy quarks decays: $`bc\mu ^{}\overline{\nu }_\mu ,c\pi ^{};tb\mu ^+\nu _\mu ,\mathrm{}`$.
* Neutrino scattering: $`\nu _\mu e^{}\mu ^{}\nu _e`$, $`\underset{\text{elastic}}{\underset{}{\nu _\mu n\mu ^{}p}}`$, $`\underset{\text{deep inelastic}}{\underset{}{\nu _\mu n\mu ^{}X}},\mathrm{}`$.
However, the Fermi theory violates unitarity at high energy. Consider the process $`\nu _ee^{}\nu _ee^{}`$ (see Figure 2), which is described by the effective Lagrangian,
$$=\frac{G_F}{\sqrt{2}}J_\mu ^{}J^\mu .$$
(2)
The amplitude has only the $`L=0`$ partial wave, and the cross section, $`\sigma `$, is proportional to $`\frac{G_F^2s}{\pi }`$, where $`\sqrt{s}=E_{cm}`$ is the total energy in the center of mass reference frame. However, for a pure S-wave process unitarity requires $`\sigma <\frac{16\pi }{s}`$. Thus, for energies,
$$\frac{1}{2}E_{cm}\begin{array}{c}>\hfill \\ \hfill \end{array}\sqrt{\frac{\pi }{G_F}}500\text{ GeV},$$
the theoretical $`\sigma `$ would violate unitarity. The non-unitarity of Born-approximation amplitudes is usually restored by higher order terms. However, the Fermi theory involves divergent integrals in second order, such as,
$$d^4k\frac{\text{/}k+m_e}{k^2m_e^2}\frac{\text{/}k}{k^2},$$
which corresponds to the process in Figure 3. Moreover, it is non-renormalizable due to the dimension of the coupling $`[G_F]=(\text{mass})^2`$.
## 2 Intermediate Vector Bosons (IVB) and $`𝑺𝑼\mathbf{(}\mathrm{𝟐}\mathbf{)}_𝑾\mathbf{\times }𝑼\mathbf{(}\mathrm{𝟏}\mathbf{)}_𝒀`$
The intermediate vector boson $`W^\pm `$ was postulated by Yukawa in 1935 in the same paper as the meson theory, and reintroduced by Schwinger in 1957. It is a massive charged particle which mediates the weak interaction (in analogy with the photon in QED) as in Figure 4.
Assuming that the IVB is massive, in the low energy regime, $`M_W|Q|`$, we have $`\frac{G_F}{\sqrt{2}}\frac{g^2}{8M_W^2}`$, where (in modern normalization) the coupling of the $`W^\pm `$ to the current is $`\frac{g}{2\sqrt{2}}`$. The amplitudes for processes like $`\nu _ee^{}\nu _ee^{}`$ are now better behaved at high energy because they are no longer pure S-wave. However, unitarity violation occurs for $`E_{cm}/2\begin{array}{c}>\hfill \\ \hfill \end{array}500`$ GeV for processes involving external $`W^\pm `$, such as $`e^+e^{}W^+W^{}`$, because of the growth of the longitudinal polarization vector $`ϵ_\mu \frac{k_\mu }{M_W}`$. Even though the coupling $`g`$ is dimensionless, the theory is still nonrenormalizable.
The divergent parts of the diagrams can be canceled if one also introduces a massive neutral boson $`W^0`$, with appropriate $`W^0W^+W^{}`$ and $`e^+e^{}W^0`$ vertices. Requiring such cancellations in all processes leads to the relation $`[J^{},J^+]=J^0`$ for the charged and neutral currents, i.e., the couplings are those of an $`SU(2)`$ gauge model. However, the model is not realistic because it does not incorporate QED. In 1961 Glashow developed a gauge model with the symmetry $`SU(2)\times U(1)`$, in which there was the $`W^\pm `$ and a massive neutral $`Z`$ boson, as well the photon, $`\gamma `$, but there was no mechanism for generating the masses of the $`W^\pm `$ and $`Z`$ bosons and the fermions (bare masses would destroy renormalizability). In 1967 Weinberg and Salam in 1968 independently applied the idea of Higgs on how massless gauge particles can acquire mass through spontaneous symmetry breaking to the $`SU(2)_W\times U(1)_Y`$ model<sup>4</sup><sup>4</sup>4The subscripts $`W`$ and $`Y`$ refer to weak and hypercharge, respectively.. In 1971 ’t Hooft and Veltman proved the renormalizability of spontaneously broken gauge theories.
The $`SU(2)_W\times U(1)_Y`$ model worked well for leptons and predicted the existence of weak neutral current (WNC) transitions such as $`\nu _e\nu _e`$ and $`e^{}e^{}`$ mediated by the new $`Z`$ boson. However, the extension to hadrons (quarks) was problematic: the observed Cabibbo-type mixing between the $`s`$ and $`d`$ quarks in the charged current implied flavor changing neutral current (FCNC) transitions $`ds`$ mediated by the $`Z`$ because of the different transformations of the $`d`$ (doublet) and $`s`$ (singlet) under $`SU(2)_W`$. These were excluded by the non-observation of certain rare $`K`$ decays and because of the size of the $`K^0\overline{K}^0`$ mixing (for which there was also a too large contribution from a second order charged current diagram). In 1970, Glashow, Iliopoulos and Maiani proposed a new quark (the charm, $`c`$, quark) as the $`SU(2)`$ partner of the $`s`$, avoiding the FCNC and $`K^0\overline{K}^0`$ problems. The $`c`$ quark was discovered in 1974 through the observation of the $`J/\mathrm{\Psi }`$ meson. The (strangeness-conserving) WNC was discovered in 1973 by the Gargamelle collaboration at CERN and by HPW at Fermilab , and was subsequently studied in great detail experimentally , as were the weak interactions of heavy quarks, especially the $`b`$. The $`W`$ and $`Z`$ were observed directly with the predicted masses in 1983 by the UA1 and UA2 experiments at CERN, and the $`SU(2)_W\times U(1)_Y`$ theory was probed at the radiative correction level in high precision experiments at LEP (CERN) and SLC (SLAC) from 1989 until $`2000`$ . The existence of non-zero neutrino mass and mixing was firmly established by the Super-Kamiokande collaboration in 1998 , and intensively studied subsequently.
## 3 Aspects of the $`𝑺𝑼\mathbf{(}\mathrm{𝟐}\mathbf{)}_𝑾\mathbf{\times }𝑼\mathbf{(}\mathrm{𝟏}\mathbf{)}_𝒀`$ Theory
It is convenient to define the left (right) chiral projections of a fermion field $`\mathrm{\Psi }`$ by $`\mathrm{\Psi }_{L(R)}\frac{1}{2}(1\gamma ^5)\mathrm{\Psi }`$. $`\mathrm{\Psi }_{L(R)}`$ coincide with states of negative (positive) helicity for a massless fermion, while chirality and helicity differ in amplitudes by terms of $`O(m/E)`$ for relativistic fermions of mass $`m`$ and energy $`E`$. In the electroweak Standard Model (SM) the left ($`L`$) and right ($`R`$) chiral fermions are assigned to transform respectively as doublets and singlets of $`SU(2)_W`$ in order to incorporate the observed $`VA`$ structure of the weak charged current. The left-chiral quark and lepton doublets are
$$\begin{array}{ccccc}& & & & \\ q_{mL}^0=\left(\begin{array}{c}u_m^0\\ d_m^0\end{array}\right)_L\hfill & \text{ }& \text{ and }& \text{ }& l_{mL}^0=\left(\begin{array}{c}\nu _m^0\\ e_m^0\end{array}\right)_L\hfill \end{array},$$
(3)
where the $`0`$ superscript refers to weak eigenstates, i.e., fields associated to weak transitions, $`m=1,\mathrm{}F`$ is the family index, and $`F`$ is the number of families. The right-chiral singlets are $`u_{mR}^0`$, $`d_{mR}^0`$, $`e_{mR}^0`$, ($`\nu _{mR}^0`$). The weak eigenstates will in general be related to the mass eigenstate fields by unitary transformations. The quark color indices $`\alpha =r,g,b`$ have been suppressed (e.g., $`u_{m\alpha R}^0`$) to simplify the notation. The fermionic part of the weak Lagrangian is
$$_F=\underset{m=1}{\overset{F}{}}\left(\overline{q}_{mL}^0iD\text{/}q_{mL}^0+\overline{l}_{mL}^0iD\text{/}l_{mL}^0+\overline{u}_{mR}^0iD\text{/}u_{mR}^0+\overline{d}_{mR}^0iD\text{/}d_{mR}^0+\overline{e}_{mR}^0iD\text{/}e_{mR}^0\right).$$
(4)
The assignment of the chiral $`L`$ and $`R`$ fields to different representations leads to parity violation, which is maximal for $`SU(2)_W`$, and also implies that bare mass terms are forbidden by the gauge symmetry. The gauge covariant derivatives are given by
$$\begin{array}{ccc}D_\mu q_{mL}^0\hfill & =\hfill & \left(_\mu +\frac{ig}{2}\tau ^iW_\mu ^i+\frac{ig^{}}{6}B_\mu \right)q_{mL}^0,\hfill \\ D_\mu l_{mL}^0\hfill & =\hfill & \left(_\mu +\frac{ig}{2}\tau ^iW_\mu ^i\frac{ig^{}}{2}B_\mu \right)l_{mL}^0,\hfill \\ D_\mu u_{mR}^0\hfill & =\hfill & \left(_\mu +i\frac{2}{3}g^{}B_\mu \right)u_{mR}^0,\hfill \\ D_\mu d_{mR}^0\hfill & =\hfill & \left(_\mu i\frac{g^{}}{3}B_\mu \right)d_{mR}^0,\hfill \\ D_\mu e_{mR}^0\hfill & =\hfill & \left(_\mu ig^{}B_\mu \right)e_{mR}^0,\hfill \end{array}$$
(5)
where $`g`$ and $`g^{}`$ are respectively the $`SU(2)_W`$ and $`U(1)_Y`$ gauge couplings, while $`W^i,i=\mathrm{1..3},`$ and $`B`$ are the (massless) gauge bosons. The weak hypercharge ($`U(1)_Y`$) assignments are determined by $`Y=QT^3`$. The right-handed fields do not couple to the $`SU(2)_W`$ terms. One can easily extend (4) and (5) by the addition of SM-singlet right-chiral neutrino fields $`\nu _{mR}^0`$, with $`D_\mu \nu _{mR}^0=_\mu \nu _{mR}^0`$.
The Higgs mechanism may be used to give masses to the gauge bosons and chiral fermions. In particular, let us introduce a complex doublet of scalar fields,
$$\varphi =\left(\begin{array}{c}\varphi ^+\\ \varphi ^0\end{array}\right),$$
and assume that the potential for $`\varphi `$ is such that the neutral component acquires a vacuum expectation value (VEV)<sup>5</sup><sup>5</sup>5For a single Higgs doublet one can always perform a gauge rotation so that only the neutral component acquires an expectation value.. Then (in unitary gauge)
$$\varphi \frac{1}{\sqrt{2}}\left(\begin{array}{c}0\\ \nu +H\end{array}\right),$$
where $`\nu =\sqrt{2}\varphi ^o=246`$ GeV, and $`H`$ is the physical Higgs scalar. Three of the original four $`SU(2)_W\times U(1)_Y`$ gauge bosons will become massive and one, the photon, remains massless.
If we rewrite the Lagrangian in the new vacuum, the scalar kinetic energy terms take the form
$$\begin{array}{ccc}\left(D_\mu \varphi \right)^{}D^\mu \varphi \hfill & & \frac{1}{2}(\begin{array}{cc}0& \nu \end{array})\left[\frac{g}{2}\tau ^iW_\mu ^i+\frac{g^{}}{2}B_\mu \right]^2\left(\begin{array}{c}0\\ \nu \end{array}\right)\hfill \\ & & M_W^2W^{+\mu }W_\mu ^{}+\frac{M_Z^2}{2}Z^\mu Z_\mu ,\hfill \end{array}$$
(6)
where the gauge interaction and kinetic energy terms of the physical $`H`$ particle have been omitted. The mass eigenstate gauge bosons are
$$\begin{array}{ccc}W^\pm \hfill & =& \frac{1}{\sqrt{2}}\left(W^1iW^2\right),\hfill \\ Z\hfill & =& \mathrm{sin}\theta _WB+\mathrm{cos}\theta _WW^3,\hfill \\ A\hfill & =& \mathrm{cos}\theta _WB+\mathrm{sin}\theta _WW^3,\hfill \end{array}$$
(7)
where the weak angle is defined by $`\mathrm{tan}\theta _Wg^{}/g`$. The masses of the gauge bosons are predicted to be
$$\begin{array}{ccccc}M_W=\frac{g\nu }{2},\hfill & & M_Z=\sqrt{g^2+g^2}\frac{\nu }{2},& & \hfill M_A=0.\end{array}$$
$`W^\pm `$ are the IVB of the charged current, $`A=\gamma `$ corresponds to the photon, and the $`Z`$ is a massive neutral boson predicted by the theory, which mediates the new neutral current interaction. The Goldstone scalars transform into the longitudinal components of the $`W^\pm `$ and $`Z`$.
A typical weak process is shown in the Figure 5. The propagator of the $`W`$ boson goes like $`\frac{1}{Q^2M_W^2}`$. For $`|Q^2|M_W^2`$, $`\frac{1}{Q^2M_W^2}\frac{1}{M_W^2}`$, which leads to the effective zero-range four-Fermi interaction in Eq. (2), with $`\frac{G_F}{\sqrt{2}}\frac{g^2}{8M_W^2}`$ as in the IVB theory. The Fermi constant is determined from the lifetime of the muon, $`\tau _\mu ^1\frac{G_F^2m_\mu ^5}{192\pi ^3}`$, yielding the electroweak scale
$$\nu =2M_W/g(\sqrt{2}G_F)^{1/2}246\text{ GeV}.$$
(8)
Comparing the $`A`$ coupling to QED, one finds $`g=\frac{e}{\mathrm{sin}\theta _W}`$, where $`e`$ is the electric charge of the positron.
Similarly, the Yukawa couplings of the Higgs doublet to the fermions introduce fermion mass matrices when $`\varphi `$ is replaced by its vacuum expectation value.
### 3.1 The Weak Charged Current
The major quantitative tests of the SM involve gauge interactions and the properties of the gauge bosons. The charged current interactions of the Fermi theory are incorporated into the Standard Model and made renormalizable. The $`W`$-fermion interaction (Figure 6) is given by
$$=\frac{g}{2\sqrt{2}}\left(J_W^\mu W_\mu ^{}+J_W^\mu W_\mu ^+\right),$$
(9)
where the weak charge-raising current is
$$\begin{array}{ccc}J_W^\mu \hfill & =& \underset{m=1}{\overset{F}{}}\left[\overline{\nu }_m^0\gamma ^\mu \left(1\gamma ^5\right)e_m^0+\overline{u}_m^0\gamma ^\mu \left(1\gamma ^5\right)d_m^0\right]\hfill \\ & & \\ & =& \left(\begin{array}{ccc}\overline{\nu }_e& \overline{\nu }_\mu & \overline{\nu }_\tau \end{array}\right)\gamma ^\mu \left(1\gamma ^5\right)\left(\begin{array}{c}e^{}\\ \mu ^{}\\ \tau ^{}\end{array}\right)+\left(\begin{array}{ccc}\overline{u}& \overline{c}& \overline{t}\end{array}\right)\gamma ^\mu \left(1\gamma ^5\right)V_{CKM}\left(\begin{array}{c}d\\ s\\ b\end{array}\right),\hfill \end{array}$$
(10)
and where we have specialized to $`F=3`$ families in the second line and introduced the fields $`e,\mu ,\tau ;u,c,t;d,s,b`$ of definite mass (mass eigenstates). We are ignoring neutrino masses for now and simply define $`\nu _e`$, $`\nu _\mu `$ and $`\nu _\tau `$ as the weak partners
of the $`e`$, $`\mu `$ and $`\tau `$. The pure $`VA`$ form ensures maximal P and C violation, while CP is conserved except for phases in $`V_{CKM}`$. We define the vectors,
$$u_L\left(\begin{array}{c}u_L\\ c_L\\ t_L\end{array}\right),u_L^0\left(\begin{array}{c}u_L^0\\ c_L^0\\ t_L^0\end{array}\right),$$
(11)
and the unitary transformations from the weak basis to the mass basis,
$$u_L^0=A_L^uu_L,d_L^0=A_L^dd_L,e_L^0=A_L^ee_L,$$
(12)
with $`A_L^uA_L^u=A_L^uA_L^u=I`$, etc. In terms of these the unitary quark mixing matrix, due to the mismatch of the unitary transformations in the $`u`$ and $`d`$ sectors, is
$$V_{CKM}=A_L^uA_L^d=\left(\begin{array}{ccc}V_{ud}& V_{us}& V_{ub}\\ V_{cd}& V_{cs}& V_{cb}\\ V_{td}& V_{ts}& V_{tb}\end{array}\right),$$
(13)
called the Cabibbo-Kobayashi-Maskawa (CKM) matrix . After removing the unobservable $`q_L`$ phases, the $`V_{CKM}`$ matrix involves three angles and one CP-violating phase. As described previously, the low energy limit of the charged current interaction reproduces the Fermi theory (with the improvement that radiative corrections can be calculated because of the renormalizability of the theory), which successfully describes weak decays and neutrino scattering processes . In particular, extensive studies, especially in $`B`$ meson decays, have been done to test the unitarity of $`V_{CKM}`$ as a probe of new physics and to test the origin of CP violation . However, the CP-violating phase in $`V_{CKM}`$ is not enough to explain baryogenesis (the cosmological asymmetry between matter and antimatter), and therefore we need additional sources of CP violation<sup>6</sup><sup>6</sup>6Possibilities include electroweak baryogenesis in the supersymmetric extension of the SM, especially if there are additional scalar singlets, or leptogenesis in theories with a heavy Majorana neutrino..
The mixing of the third quark family with the first two is small, and it is sometimes a good zeroth-order approximation to ignore it and consider an effective $`F=2`$ theory. Then $`V_{CKM}`$ is replaced by the Cabibbo matrix
$$V=\left(\begin{array}{cc}\mathrm{cos}\theta _c& \mathrm{sin}\theta _c\\ \mathrm{sin}\theta _c& \mathrm{cos}\theta _c\end{array}\right),$$
(14)
which depends on only one parameter, the Cabibbo angle, $`\mathrm{sin}\theta _c0.22`$. In this case there is no CP violation because all the phases can be reabsorbed in the fields.
In addition to reproducing the Fermi theory, the charged current propagator of the SM has been probed at HERA in $`e^\pm p\underset{e}{\overset{()}{\nu }}X`$ at high energy. The theory has also been tested at the loop level in processes such as $`K^0\overline{K}^0`$ and $`B\overline{B}`$ mixing, and in the calculation of (finite) radiative corrections, which are necessary for the consistency of such processes as $`\beta `$ and $`\mu `$ decay.
### 3.2 The Weak Neutral Current
The weak neutral current was predicted by the $`SU(2)_W\times U(1)_Y`$ model and the Lagrangian for it is
$$\begin{array}{ccc}\hfill & =& \frac{\sqrt{g^2+g^2}}{2}J_Z^\mu \left(\mathrm{sin}\theta _WB_\mu +\mathrm{cos}\theta _WW_\mu ^3\right)\hfill \\ & =& \frac{g}{2\mathrm{cos}\theta _W}J_Z^\mu Z_\mu ,\hfill \end{array}$$
(15)
where we have used $`\mathrm{cos}\theta _W=\frac{g}{\sqrt{g^2+g^2}}`$. The neutral current is given by
$$\begin{array}{ccc}J_Z^\mu \hfill & =& \underset{m}{}\left[\overline{u}_{mL}^0\gamma ^\mu u_{mL}^0\overline{d}_{mL}^0\gamma ^\mu d_{mL}^0+\overline{\nu }_{mL}^0\gamma ^\mu \nu _{mL}^0\overline{e}_{mL}^0\gamma ^\mu e_{mL}^0\right]2\mathrm{sin}^2\theta _WJ_Q^\mu \hfill \\ & =& \underset{m}{}\left[\overline{u}_{mL}\gamma ^\mu u_{mL}\overline{d}_{mL}\gamma ^\mu d_{mL}+\overline{\nu }_{mL}\gamma ^\mu \nu _{mL}\overline{e}_{mL}\gamma ^\mu e_{mL}\right]2\mathrm{sin}^2\theta _WJ_Q^\mu ,\hfill \end{array}$$
(16)
where $`J_Q^\mu =\mathrm{\Sigma }_iq_i\overline{\psi }_i\gamma ^\mu \psi _i`$ is the electromagnetic current. $`J_Z^\mu `$ is flavor-diagonal in the Standard Model like the electromagnetic current, i.e., the weak neutral current has the same form in the weak and the mass bases, because the fields which mix have the same assignments under the $`SU(2)_W\times U(1)_Y`$ gauge group. This was the reason for introducing the GIM mechanism, discussed in section 2, into the theory. The neutral current has two contributions; one is purely $`VA`$, while the other is proportional to the electromagnetic current and is purely vector. Therefore parity is violated but not maximally.
The neutral current interaction of fermions in the low momentum limit is given by
$$_{eff}^{NC}=\frac{G_F}{\sqrt{2}}J_Z^\mu J_{Z\mu }.$$
(17)
The coefficients for the charged and neutral current interactions are the same because
$$\frac{G_F}{\sqrt{2}}=\frac{g^2}{8M_W^2}=\frac{g^2+g^2}{8M_Z^2}.$$
(18)
Since its discovery, the weak neutral current has been tested and extensively studied in many processes , These include pure weak processes such as $`\nu e\nu e`$, $`\nu N\nu N`$, and $`\nu N\nu X`$; weak-electromagnetic interference, including parity violating asymmetries in processes like polarized $`e^{}DeX`$, atomic parity violation, and $`e^+e^{}`$ annihilation above and below the $`Z`$-pole; and high precision (e.g., 0.1%) $`Z`$-pole reactions at CERN and SLAC.
## 4 Neutrino Preliminaries
Neutrinos are a unique probe of many aspects of physics and astrophysics on scales ranging from $`10^{33}`$ to $`10^{+28}cm`$ . Decays and scattering processes involving neutrinos have been powerful tests of many aspects of particle physics, and have been important in establishing the Fermi theory and the SM (at $``$ 1%), searching for new physics at the TeV scale, as a probe of the structure of hadrons, and as a test of QCD. Small neutrino masses are sensitive to new physics at scales ranging from a TeV up to grand unification and superstring scales. Similarly, neutrinos are important for the physics of the Sun, stars, core-collapse supernovae, the origins of cosmic rays, the large scale structure of the universe, big bang nucleosynthesis, and possibly baryogenesis.
### 4.1 Weyl Spinors
A Weyl two-component spinor is the minimal fermionic degree of freedom. A left-chiral Weyl spinor $`\mathrm{\Psi }_L`$ satisfies
$$P_L\mathrm{\Psi }_L=\mathrm{\Psi }_L,P_R\mathrm{\Psi }_L=0,$$
(19)
where $`P_{L(R)}=\frac{1\gamma ^5}{2}`$.A $`\mathrm{\Psi }_L`$ annihilates $`L`$ particles or creates $`R`$ antiparticles. Similarly, a right-chiral spinor $`\mathrm{\Psi }_R`$ satisfies $`P_R\mathrm{\Psi }_R=\mathrm{\Psi }_R,P_L\mathrm{\Psi }_R=0`$. A Weyl spinor can exist by itself (e.g., the $`\nu _L`$ in the SM), or can be considered a projection of a 4-component Dirac spinor $`\mathrm{\Psi }=\mathrm{\Psi }_L+\mathrm{\Psi }_R`$ (e.g., $`e^{}=e_L^{}+e_R^{}`$). A $`\mathrm{\Psi }_L`$ field is related by CP or CPT to a right-chiral antiparticle spinor<sup>7</sup><sup>7</sup>7Which is called the particle and which the antiparticle is a matter of convenience. $`\mathrm{\Psi }_R^c`$ ($`\mathrm{\Psi }_L\stackrel{\mathrm{CP},\mathrm{CPT}}{}\mathrm{\Psi }_R^c`$), where<sup>8</sup><sup>8</sup>8Thus, in our notation the subscript $`L`$ ($`R`$) always refers to a left (right) chiral spinor, independent of whether it is for a particle or antiparticle. $`\mathrm{\Psi }_R^c`$ annihilates $`R`$ antiparticles or creates $`L`$ particles and $`P_R\mathrm{\Psi }_R^c=\mathrm{\Psi }_R^c`$. Under CP, for example,
$$\mathrm{\Psi }_L(\stackrel{}{x},t)\gamma ^0\mathrm{\Psi }_R^c(\stackrel{}{x},t).$$
(20)
$`\mathrm{\Psi }_R^c`$ is essentially the adjoint of $`\mathrm{\Psi }_L`$:
$$\begin{array}{c}\mathrm{\Psi }_R^c=C\overline{\mathrm{\Psi }_L}^T=C\gamma _0^T(\mathrm{\Psi }_L^{})^T,\hfill \\ \mathrm{\Psi }_L=C\overline{\mathrm{\Psi }_R^c}^T,\hfill \end{array}$$
(21)
where $`T`$ represents the transpose and $`C`$ is the charge conjugation operator, defined by $`C\gamma _\mu C^1=\gamma _\mu ^T`$. In the Pauli-Dirac representation $`C=i\gamma ^2\gamma ^0`$. The free Weyl field is
$$\mathrm{\Psi }_L(x)=\frac{d^3\stackrel{}{p}}{\sqrt{(2\pi )^32E_p}}\left[a_L(\stackrel{}{p})u_L(\stackrel{}{p})e^{ipx}+b_R^{}(\stackrel{}{p})v_R(\stackrel{}{p})e^{ipx}\right]\text{ (no sum over spin).}$$
When two Weyl spinors are present, we can use either ($`\mathrm{\Psi }_L`$ and $`\mathrm{\Psi }_R`$) or ($`\mathrm{\Psi }_L`$ and $`\mathrm{\Psi }_L^c`$) to define the theory. In the Standard Model, for example, baryon ($`B`$) and lepton ($`L`$) numbers are conserved perturbatively, and it is convenient to work with, e.g., $`u_L`$ and $`u_R`$ and similarly for the other fermionic fields. Their CPT partners are $`u_R^cu_L`$ and $`u_L^cu_R`$. In some extensions of the SM, such as supersymmetry or grand unification, it is more convenient to work with $`L`$ spinors $`\mathrm{\Psi }_L`$ and $`\mathrm{\Psi }_L^c`$.
### 4.2 Active and Sterile Neutrinos
Active (a.k.a. ordinary or doublet) neutrinos are left-handed neutrinos which transform as $`SU(2)`$ doublets with a charged lepton partner, and which therefore have normal weak interactions. The $`L`$ doublets and their right-handed partners are
$$\begin{array}{ccc}\left(\begin{array}{c}\nu _e\\ e^{}\end{array}\right)_L\hfill & \stackrel{\mathrm{CPT}}{}& \left(\begin{array}{c}e^+\\ \nu _e^c\end{array}\right)_R\hfill \end{array}.$$
(22)
Sterile (a.k.a. singlet or “right-handed”) neutrinos, which are present in most extensions of the SM, are $`SU(2)`$ singlets. They do not interact except by mixing, Yukawa interactions, or beyond the SM (BSM) interactions. It is convenient to denote the right-chiral spinor as $`N_R`$ and its conjugate as $`N_L^c`$.
$$N_R\stackrel{\mathrm{CPT}}{}N_L^c.$$
(23)
## 5 Neutrino Mass Models
Most extensions of the Standard Model lead to nonzero neutrino mass at some level. There are many models of neutrino mass. Here, only a brief survey of the principal classes is given. For more detail see .
### 5.1 Fermion Masses
Mass terms convert a spinor of one chirality into one of the opposite chirality,
$$m(\overline{\mathrm{\Psi }}_{2L}\mathrm{\Psi }_{1R}+\overline{\mathrm{\Psi }}_{1R}\mathrm{\Psi }_{2L}).$$
(24)
An interpretation is that a massless fermion has the same helicity (chirality) in all frames of reference, while that of a massive particle depends on the reference frame and therefore can be flipped.
### 5.2 Dirac Mass
A Dirac mass (Figure 8) connects two distinct Weyl neutrinos<sup>9</sup><sup>9</sup>9These are usually active and sterile. However, there are variant forms involving two distinct active (Zeldovich-Konopinski-Mahmoud) or two distinct sterile neutrinos. $`\nu _L`$ and $`N_R`$. It is given by
$$_D=m_D(\overline{\nu }_LN_R+\overline{N}_R\nu _L)=m_D\overline{\nu }\nu .$$
(25)
Thus, $`\mathrm{\Psi }_{1R}\mathrm{\Psi }_{2R}^c`$ and we have two distinct Weyl spinors. These may be combined to form a Dirac field $`\nu \nu _L+N_R`$ with four components: $`\nu _L`$, $`\nu _R^c`$, $`N_R`$ and $`N_L^c`$. This can be generalized to three or more families.
This mass term allows a conserved lepton number $`L`$, which implies no mixing between $`\nu _L`$ and $`N_L^c`$, or between $`\nu _R^c`$ and $`N_R`$. However, it violates weak isospin $`I_W`$ by $`\mathrm{\Delta }I_W=\frac{1}{2}`$. A Dirac mass can be generated by the Higgs mechanism, as in Figure 9.
The Dirac neutrino masses $`m_D=h_\nu \varphi ^0=h_\nu \frac{\nu }{\sqrt{2}}`$ are analogous to the quark and charged lepton masses. The upper bound on the neutrino mass, $`m_\nu \begin{array}{c}<\hfill \\ \hfill \end{array}1`$ eV, implies a Yukawa coupling of the neutrino to the Higgs, $`h_\nu <10^{11}`$, which is extremely small in comparison with the coupling for the top quark, $`h_t=O(1)`$, or for the electron $`h_e^{}10^5`$.
### 5.3 Majorana Mass
A Majorana mass term describes a transition between a left-handed neutrino and its CPT-conjugate right-handed antineutrino, as shown in Figure 10. It can be viewed as the annihilation or creation of two neutrinos, and therefore violates lepton number by two units, $`\mathrm{\Delta }L=2`$.
A Majorana mass term has the form
$$=\frac{1}{2}m_T(\overline{\nu }_L\nu _R^c+\overline{\nu }_R^c\nu _L)=\frac{m_T}{2}(\overline{\nu }_LC\overline{\nu _L}^T+h.c)=\frac{1}{2}m_T\overline{\nu }\nu ,$$
(26)
where $`\nu =\nu _L+\nu _R^c`$ is a self-conjugate two-component (Majorana) field satisfying $`\nu =\nu ^c=C\overline{\nu }^T`$. A Majorana mass for an active neutrino violates weak isospin by one unit, $`\mathrm{\Delta }I_W=1`$, and can be generated either by the VEV of a Higgs triplet (see Figure 11)
or by a higher-dimensional operator involving two Higgs doublets (which could be generated, for example, in a seesaw model). A Majorana $`\nu `$ is its own antiparticle and can mediate neutrinoless double beta decay ($`\beta \beta _{0\nu }`$, in which two neutrons turn into two protons and two electrons, violating lepton number by two units, as shown in Figure 12.
A sterile neutrino can also have a Majorana mass term of the form,
$$=\frac{m_S}{2}[\overline{N}_L^cN_R+\overline{N}_RN_L^c].$$
(27)
In this case, weak isospin is conserved, $`\mathrm{\Delta }I_W=0`$, and thus $`m_S`$ can be generated by the VEV of a Higgs singlet<sup>10</sup><sup>10</sup>10$`m_S`$ could also be generated in principle by a bare mass, but this is usually forbidden by additional symmetries in extensions of the SM..
### 5.4 Mixed models
If active and sterile neutrinos are both present, one can have both Dirac and Majorana mass terms simultaneously. For one family, the Lagrangian has the form
$$=\frac{1}{2}\left(\begin{array}{cc}\overline{\nu }_L^0& \overline{N}_L^{0c}\end{array}\right)\left(\begin{array}{cc}m_T& m_D\\ m_D& m_S\end{array}\right)\left(\begin{array}{c}\nu _R^{0c}\\ N_R^0\end{array}\right)+h.c.,$$
(28)
where $`0`$ refers to weak eigenstates and the masses are
$$\begin{array}{cccccc}m_T:\hfill & |\mathrm{\Delta }L|=2,\hfill & & \mathrm{\Delta }I_W=1& & (\text{Majorana}),\hfill \\ m_D:\hfill & |\mathrm{\Delta }L|=0,\hfill & & \mathrm{\Delta }I_W=\frac{1}{2}& & (\text{Dirac}),\hfill \\ m_S:\hfill & |\mathrm{\Delta }L|=2,\hfill & & \mathrm{\Delta }I_W=0& & (\text{Majorana}).\hfill \end{array}$$
Diagonalizing the matrix in (28) yields two Majorana mass eigenvalues and two Majorana mass eigenstates, $`\nu _i=\nu _{iL}+\nu _{iR}^c=\nu _i^c`$, with $`i=1,2`$. The weak and mass bases are related by the unitary transformations
$$\begin{array}{ccccc}\left(\begin{array}{c}\nu _{1L}\\ \nu _{2L}\end{array}\right)=U_L^\nu \left(\begin{array}{c}\nu _L^0\\ N_L^{0c}\end{array}\right),\hfill & & \left(\begin{array}{c}\nu _{1R}^c\\ \nu _{2R}^c\end{array}\right)=U_R^\nu \left(\begin{array}{c}\nu _R^{0c}\\ N_R^0\end{array}\right)& & \end{array}.$$
(29)
$`U_L`$ and $`U_R`$ are generally different for Dirac mass matrices, which need not be Hermitian. However, the general $`2\times 2`$ neutrino mass matrix in (28) is symmetric because of (21). In our phase convention, in which $`\nu _{iR}^c=C\overline{\nu _{iL}}^T`$, this implies $`U_L^\nu =U_R^\nu `$.
### 5.5 Special Cases
$`m_D=0`$ in (28) corresponds to the pure Majorana case: the mass matrix is diagonal, with $`m_1=m_T`$, $`m_2=m_S`$, and
$$\begin{array}{ccc}\nu _{1L}=\nu _L^0,\hfill & & \hfill \nu _{1R}^c=\nu _R^{0c},\\ \nu _{2L}=N_L^{0c},\hfill & & \hfill \nu _{2R}^c=N_R^0.\end{array}$$
$`m_T=m_S=0`$ is the Dirac limit. There are formally two Majorana mass eigenstates, with eigenvalues $`m_1=m_D`$ and $`m_2=m_D`$ and eigenstates
$$\begin{array}{ccc}\nu _{1L}=\frac{1}{\sqrt{2}}(\nu _L^0+N_L^{0c}),\hfill & & \hfill \nu _{2L}=\frac{1}{\sqrt{2}}(\nu _L^0N_L^{0c}),\\ \nu _{1R}^c=\frac{1}{\sqrt{2}}(\nu _R^{0c}+N_R^0),\hfill & & \hfill \nu _{2R}^c=\frac{1}{\sqrt{2}}(\nu _R^{0c}N_R^0).\end{array}$$
Note that $`\nu _{1,2}`$ are degenerate in the sense that $`|m_1|=|m_2|`$. To recover the Dirac limit, expand the mass term
$$\begin{array}{ccc}\hfill & =& \frac{m_D}{2}(\overline{\nu }_{1L}\nu _{1R}^c\overline{\nu }_{2L}\nu _{2R}^c)+h.c.\hfill \\ & =& \frac{m_D}{2}(\overline{\nu }_L^0N_R^0+\overline{N}_R^0\nu _L^0),\hfill \end{array}$$
(30)
which clearly conserves lepton number (i.e., there is no $`\nu _L^0N_L^{0c}`$ or $`\nu _R^{0c}N_R^0`$ mixing). Thus, a Dirac neutrino can be thought of as two Majorana neutrinos, with maximal ($`45^{}`$) mixing and with equal and opposite masses. This interpretation is useful in considering the Dirac limit of general models.
$`m_T=0`$ and $`m_Sm_D`$ (e.g., $`m_D=O(m_u,m_e,m_d)`$ and $`m_S=O(M_X)`$, where $`M_X10^{14}`$ GeV) is the seesaw limit , with eigenstates and eigenvalues
$$\begin{array}{ccc}\nu _{1L}\nu _L^0\frac{m_D}{m_S}N_L^{0c}\nu _L^0,\hfill & m_1\frac{m_D^2}{m_S}m_D,& \\ \nu _{2L}\frac{m_D}{m_S}\nu _L^0+N_L^{0c}N_L^{0c},\hfill & m_2m_S.& \end{array}$$
The seesaw mechanism is illustrated in Figure 13.
this is a perturbation on the Dirac case, with $`m_T,m_Sm_D`$. There is a small lepton number violation, and a small splitting between the magnitudes of the mass eigenvalues. For example, $`m_T=ϵ,m_S=0`$ leads to $`|m_{1,2}|=m_D\pm ϵ/2`$.
The general case in which $`m_D`$ and $`m_S`$ (and/or $`m_T`$) are both small and comparable leads to non-degenerate Majorana mass eigenvalues and significant ordinary – sterile ($`\nu _L^0N_L^{0c}`$) mixing (as would be suggested if the LSND results are confirmed). Only this and the pseudo-Dirac cases allow such mixings.
As we have seen, the fermion mass eigenvalues can be negative or even complex. However, only the magnitude is relevant for most purposes, and in fact fermion masses can always be made real and positive by chiral transformations (e.g., redefining the phases of the $`\nu _{iR}^c`$). However, the relative signs (or phases) may reappear in the $`\beta \beta _{0\nu }`$ amplitude or new (beyond the SM) interactions.
### 5.6 Extension to Three Families
It is straightforward to generalize (28) to three or more families, or even to the case of different numbers of active and sterile neutrinos. For $`F=3`$ families, the Lagrangian is
$$=\frac{1}{2}\left(\begin{array}{cc}\overline{\nu }_L^0& \overline{N}_L^{0c}\end{array}\right)\left(\begin{array}{cc}m_T& m_D\\ m_D^T& m_S\end{array}\right)\left(\begin{array}{c}\nu _R^{0c}\\ N_R^0\end{array}\right)+h.c.,$$
(31)
where $`\nu _L^0`$ and $`N_L^{0c}`$ are three-component vectors
$$\begin{array}{ccc}\nu _L^0=\left(\begin{array}{c}\nu _{1L}^0\\ \nu _{2L}^0\\ \nu _{3L}^0\end{array}\right),\hfill & & \hfill N_L^{0c}=\left(\begin{array}{c}N_{1L}^{0c}\\ N_{2L}^{0c}\\ N_{3L}^{0c}\end{array}\right),\end{array}$$
(32)
and similarly for $`\nu _R^{0c}`$ and $`N_R^0`$. $`m_S`$, $`m_D`$ and $`m_T`$ are $`3\times 3`$ matrices, where $`m_S=m_S^T`$ and $`m_T=m_T^T`$ because of (21). There are six Majorana mass eigenvalues and eigenvectors. The transformation to go from the weak to the mass basis is given by
$$\nu _L=U_L^\nu \left(\begin{array}{c}\nu _L^0\\ N_L^{0c}\end{array}\right),$$
(33)
where $`U_L^\nu `$ is a $`6\times 6`$ unitary matrix and $`\nu _L`$ is a six-component vector. The analogous transformation for the $`R`$ fields involves $`U_R=U_L^{}`$ because the $`6\times 6`$ Majorana mass matrix is necessarily symmetric.
First consider the case in which there are only Dirac masses, which is analogous to the quarks and charged leptons. The mass Lagrangian is
$$=\overline{\nu }_L^0m_DN_R^0+h.c.,$$
(34)
where $`m_D`$ need not be Hermitian. We may change from the weak to the mass basis with a unitary transformation,
$$=\overline{\nu }_L\widehat{m}_DN_R,$$
(35)
where $`\widehat{m}_D`$ is the diagonal matrix of eigenvalues, $`\widehat{m}_D=U_L^\nu m_DU_R^\nu =\mathrm{diag}(m_1,\mathrm{}m_3)`$, and the eigenvectors are
$$\begin{array}{ccc}\nu _L=U_L^\nu \nu _L^0,\hfill & & \hfill N_R=U_R^\nu N_R^0\end{array}.$$
(36)
If $`m_D`$ is not Hermitian, $`U_{L,R}^\nu `$ can be found by diagonalizing the Hermitian matrices $`m_Dm_D^{}`$ and $`m_D^{}m_D`$,
$$U_L^\nu m_Dm_D^{}U_L^\nu =U_R^\nu m_D^{}m_DU_R^\nu =|\widehat{m}_D|^2=\mathrm{diag}(|m_1|^2\mathrm{}|m_3|^2).$$
(37)
The $`U_{L,R}^\nu `$ are not unique. If a given $`U_L^\nu `$ satisfies (37), then so does
$$\widehat{U_L^\nu }=U_L^\nu K_L^\nu ,$$
(38)
where $`K_L^\nu `$ is a diagonal phase matrix<sup>11</sup><sup>11</sup>11There is more freedom if there are degenerate eigenvalues., $`K_L^\nu =\mathrm{diag}\left(\begin{array}{ccc}e^{i\alpha _1}& e^{i\alpha _2}& e^{i\alpha _3}\end{array}\right)`$. This of course corresponds to the freedom to redefine the phases of the mass eigenstate fields. Similar statements apply to $`U_R^\nu `$ and to the unitary transformations for the quark and charged lepton fields.
The weak charged current is given by
$$J_W^\mu =2\overline{e}_L^0\gamma ^\mu \nu _L^0=2\overline{e}_L\gamma ^\mu U_L^eU_L^\nu \nu _L,$$
(39)
and $`U_L^eU_L^\nu `$ is the lepton mixing matrix $`V_{PMNS}`$<sup>12</sup><sup>12</sup>12This is for Maki, Nakagawa, and Sakata in 1962 and Pontecorvo in 1968.. A general unitary $`3\times 3`$ matrix involves 9 parameters: 3 mixing angles and 6 phases. However, five of the phases in $`V_{PMNS}`$ (and similarly in $`V_{CKM}`$) are unobservable in the Dirac case because they depend on the relative phases of the left-handed neutrino and charged lepton fields, i.e., they can be removed by an appropriate choice of $`K_L^\nu `$ and the analogous $`K_L^e`$, so there is only one physical CP-violating phase. The analogous $`K_R^{\nu ,e}`$ in $`U_R^{\nu ,e}`$ can then be chosen to make the mass eigenvalues real and positive. The $`K_R`$ phases are unobservable unless there are new BSM interactions involving the $`R`$ fields.
The diagonalization of a Majorana mass matrix is similar, except for the constraint $`U_L^\nu =U_R^\nu `$. This implies that $`K_L^\nu `$ and $`K_R^\nu `$ cannot be chosen independently (because the $`L`$ and $`R`$ mass eigenstates are adjoints of each other). One must then use the freedom in $`K_L^\nu `$ to make the mass eigenvalues real and positive, so that one cannot remove as many phases from the leptonic mixing matrix. For three light mass eigenstates (e.g., for three active neutrinos with $`m_T0,m_D=0`$; or for the seesaw limit of the $`6\times 6`$ case), this implies that there are two additional “Majorana” phases, i.e., $`V_{PMNS}=\widehat{V}_{PMNS}\times \mathrm{diag}\left(\begin{array}{ccc}e^{i\beta _1}& e^{i\beta _2}& 1\end{array}\right)`$, where $`\widehat{V}_{PMNS}`$ has the canonical form with one phase, and $`\beta _{1,2}`$ are CP-violating relative phases associated with the Majorana neutrinos. $`\beta _{1,2}`$ do not affect neutrino oscillations, but can in principle affect $`\beta \beta _{0\nu }`$ and new interactions.
Let us conclude this section with a few comments.
* The LSND oscillation results would, if confirmed, strongly suggest the existence of very light sterile neutrinos which mix with active neutrinos of the same chirality . The pure Majorana and pure Dirac cases do not allow any such mixing (sterile neutrinos are not even required in the Majorana case). The seesaw limit only has very heavy sterile neutrinos and negligible mixing. Only the general and pseudo-Dirac limits allow significant ordinary-sterile mixing, but in these cases one must find an explanation for two very small types of masses.
* There is no distinction between Dirac and Majorana neutrinos except by their masses (or by new interactions). As the masses go to zero, the active components reduce to standard active Weyl spinors in both case. There are additional sterile Weyl spinors in the massless limit of the Dirac case, but these decouple from the other particles.
* One can ignore $`V_{PMNS}`$ in processes for which the neutrino masses are too small to be relevant. In that limit, the neutrino masses are effectively degenerate (with vanishing mass) and one can simply work in the weak basis.
### 5.7 Models of Neutrino Mass
There are an enormous number of models of neutrino mass .
Models to generate Majorana masses are most popular, because no Standard Model gauge symmetry forbids them. (However, models descended from underlying string constructions may be more restrictive because of the underlying symmetries and selection rules .) Models constructed to yield small Majorana masses include: the ordinary (type I) seesaw , often combined with additional family and grand unification symmetries; models with heavy Higgs triplets (type II seesaw) ; TeV (extended) seesaws , with $`m_\nu m^{p+1}/M^p,p>1`$, e.g., with $`M`$ in the TeV range; radiative masses (i.e., generated by loops) ; supersymmetry with $`R`$-parity violation, which may include loop effects; supersymmetry with mass generation by terms in the Kähler potential; anarchy (random entries in the mass matrices); and large extra dimensions (LED), possibly combined with one of the above.
Small Dirac masses may be due to: higher dimensional operators (HDO) in intermediate scale models (e.g., associated with an extended $`U(1)^{}`$ gauge symmetry or supersymmetry breaking); large intersection areas in intersecting brane models or large extra dimensions, from volume suppression if $`N_R`$ propagates in the bulk .
Simultaneous small Dirac and Majorana masses, as motivated by LSND, may be due, e.g., to HDO or to sterile neutrinos from a “mirror world” .
There are also many “texture” models , involving specific guesses about the form of the $`3\times 3`$ neutrino mass matrix or the Dirac and Majorana matrices entering seesaw models. These are often studied in connection with models also involving quark and charged lepton mass matrices, such as grand unification (GUTs), family symmetries, or left-right symmetry.
## 6 Laboratory and Astrophysical Constraints on Neutrino Counting and Mass
### 6.1 Laboratory Limits
The most precise measurement of the number of light ($`m_\nu <M_Z/2`$) active neutrino types and therefore the number of associated fermion families comes from the invisible $`Z`$ width $`\mathrm{\Gamma }_{inv}`$, obtained by subtracting the observed width into quarks and charged leptons from the total width from the lineshape. The number of effective neutrinos $`N_\nu `$ is given by
$$N_\nu =\frac{\mathrm{\Gamma }_{inv}}{\mathrm{\Gamma }_l}\left(\frac{\mathrm{\Gamma }_l}{\mathrm{\Gamma }_\nu }\right)_{SM},$$
(40)
where $`(\mathrm{\Gamma }_l/\mathrm{\Gamma }_\nu )_{SM}`$ is the SM expression for the ratio of widths into a charged lepton and a single active neutrino, introduced to reduce the model dependence. The experimental value is $`N_\nu =2.984\pm 0.008`$ , excluding the possibility of a fourth family unless the neutrino is very heavy. Other unobserved particles from $`Z`$ decay would also give a positive contribution to $`N_\nu `$. For example, the decay $`ZMS`$ in models with spontaneous lepton number violation and a Higgs triplet , where $`M`$ is a Goldstone boson (Majoron) and $`S`$ is a light scalar, would yield a contribution of 2 to $`N_\nu `$ and is therefore excluded.
Kinematic laboratory measurements on neutrino masses are relatively weak . For the $`\tau `$-neutrino the most stringent limit is $`m_{\nu _\tau }<18.2`$ MeV, which comes from the decay channel $`\tau 5\pi +\nu _\tau `$. For the $`\mu `$-neutrino the most sensitive measurement gives $`m_{\nu _\mu }<0.19`$ MeV from $`\pi \mu \nu _\mu `$ decay<sup>13</sup><sup>13</sup>13The limits on $`m_{\nu _\tau }`$ and $`m_{\nu _\mu }`$ are now obsolete because of oscillation and cosmological constraints.. The analysis of tritium beta decay at low energies gives a limit on the mass of the $`e`$-neutrino of $`m_{\nu _e}2.8`$ eV. Including mixing, this should be interpreted as a limit on $`m_\beta \sqrt{\mathrm{\Sigma }_i|V_{ei}|^2|m_i|^2}`$, where $`V=V_{PMNS}`$. The latter should be improved in the future to a sensitivity of around 0.2 eV by the KATRIN experiment.
For Majorana masses, the amplitude for neutrinoless double beta decay is $`AA_{nuc}m_{\beta \beta }`$, where $`A_{nuc}`$ contains the nuclear matrix element and $`m_{\beta \beta }_i(V_{ei})^2m_i`$ is the effective Majorana mass in the presence of mixing between light Majorana neutrinos<sup>14</sup><sup>14</sup>14There can also be other BSM contributions to the decay, such as heavy Majorana neutrinos or $`R`$-parity violating effects in supersymmetry.. $`m_{\beta \beta }`$ is just the $`(1,1)`$ element of $`m_T`$ or of the effective Majorana mass matrix in a seesaw model. It involves the square of $`V_{ei}`$ rather than the absolute square, and is therefore sensitive in principle to the Majorana phases (though this is difficult in practice). The expression allows for the possibility of cancellations between different mass eigenstates<sup>15</sup><sup>15</sup>15In the convention for the field phases that the masses are real and positive the cancellations are due to the phases in $`V_{ei}`$. Alternatively, one can use the convention that the mass eigenvalues may be negative or complex., and in fact shows why $`\beta \beta _{0\nu }`$ vanishes for a Dirac neutrino, which can be viewed as two canceling Majorana neutrinos. At present, $`|m_{\beta \beta }|<(0.351)`$ eV, where the range is from the theoretical uncertainty in $`A_{nuc}`$. Members of one experiment claimed an observation of $`\beta \beta _{0\nu }`$, which corresponds to $`m_{\beta \beta }0.39`$ eV, but this has not been confirmed. Future experiments should be sensitive to $`m_{\beta \beta }(0.10.2\text{ eV})`$.
### 6.2 Cosmological Constraints
The light nuclides $`{}_{}{}^{4}He`$, $`D`$, $`{}_{}{}^{3}He`$, $`{}_{}{}^{7}Li`$ were synthesized in the first thousand seconds in the early evolution the universe , corresponding to temperatures from 1 MeV to 50 keV. At temperatures above the freezeout temperature $`T_f\mathrm{few}`$ MeV, the neutron to proton ratio was kept in equilibrium by the reactions
$$\begin{array}{ccc}n+\nu _e& & p+e^{},\\ & & \\ n+e^+& & p+\overline{\nu }_e.\end{array}$$
(41)
For $`T<T_f`$ their rates became slow compared to the expansion rate of the universe, and the neutron to proton ratio, $`n/p`$, froze at a constant (except for neutron decay) value $`\frac{n}{p}=\mathrm{exp}(\frac{m_nm_p}{T_f})`$. Most of the neutrons were incorporated into $`{}_{}{}^{4}He`$, so that the primordial abundance (relative to hydrogen) can be predicted<sup>16</sup><sup>16</sup>16There is also a weak dependence on the baryon density relative to photons, which is determined independently by the $`D`$ abundance and by the cosmic microwave background (CMB) anisotropies. in terms of $`T_f`$. $`T_f`$ is predicted by comparing the reaction rate for the processes in Eq. (41), $`\mathrm{\Gamma }G_F^2T^5`$, and the Hubble expansion rate, $`H=1.66\sqrt{g_{}}\frac{T^2}{M_{\mathrm{Pl}}}\sqrt{g_{}}T^2`$, where $`M_{\mathrm{Pl}}`$ is the Planck scale. Therefore, $`T_fg_{}^{1/6}`$, where $`g_{}`$ counts the number of relativistic particle species, determining the energy density in radiation. It is given by $`g_{}=g_B+\frac{7}{8}g_F`$, where $`g_F=10+2\mathrm{\Delta }N_\nu ^{}`$, and where the $`10`$ is due to $`3\nu +3\overline{\nu }+2\text{ helicities each of }e^\pm `$. $`\mathrm{\Delta }N_\nu ^{}`$ is the effective number of additional neutrinos present at $`T\begin{array}{c}>\hfill \\ \hfill \end{array}T_f`$. It includes new active neutrinos with masses<sup>17</sup><sup>17</sup>17There would be an enhanced contribution from a $`\nu _\tau `$ in the 1-20 MeV range, which was once important in constraining its mass, or a reduced contribution from $`\nu _\tau `$ decay. $`\begin{array}{c}<\hfill \\ \hfill \end{array}1`$ MeV, and also light sterile neutrinos, which could be produced by mixing with active neutrinos for a wide range of mixing angles<sup>18</sup><sup>18</sup>18The effects of light sterile neutrinos can be avoided in variant scenarios or compensated by a large $`\nu _e\overline{\nu }_e`$ asymmetry .. It does not include the right-handed sterile components of light Dirac neutrinos, which would not have been produced in significant numbers unless they couple to BSM interactions. The prediction of the primordial mass abundance of $`{}_{}{}^{4}He`$ relative to $`H`$ is $`24`$% for $`\mathrm{\Delta }N_\nu ^{}=0`$. There is considerable uncertainty in the observational value, but most estimates yield $`\mathrm{\Delta }N_\nu ^{}<0.11`$.
Neutrinos with masses in the eV range would contribute hot dark matter to the universe . They would close the universe, $`\mathrm{\Omega }_\nu =1`$, for the sum of masses of the light neutrinos (including sterile neutrinos, weighted by their abundance relative to active neutrinos) $`\mathrm{\Sigma }\mathrm{\Sigma }_i|m_i|35`$ eV. However, even though some 30% of the energy density is believed to be in the form dark matter, most of it should be cold (such as weakly interacting massive particles) rather than hot (neutrinos), because the latter cannot explain the formation of smaller scale structures during the lifetime of the universe. The Wilkinson Microwave Anisotropy Probe (WMAP), together with the Sloan Digital Sky Survey (SDSS), the Lyman alpha forest (Ly$`\alpha `$), and other observations, leads to strong constraints on $`\mathrm{\Sigma }\begin{array}{c}<\hfill \\ \hfill \end{array}1`$ eV , with the most stringent claimed limit of 0.42 eV . Using future Planck data, it may be possible to extend the sensitivity down to $`0.050.1`$ eV, close to the minimum value $`0.05\text{ eV}\sqrt{|\mathrm{\Delta }m_{\mathrm{atm}}^2|}`$ allowed by the neutrino oscillation data. However, there are significant theoretical uncertainties.
### 6.3 Neutrino Oscillations
Neutrino oscillations can occur due to the mismatch between weak and mass eigenstates, and are analogous to the time evolution of a quantum system which is not in an energy eigenstate, or a classical coupled oscillator in which one starts with an excitation that is not a normal mode.
Consider a system in which there are only two neutrino flavors, e.g., $`\nu _e`$ and $`\nu _\mu `$. Then
$$\begin{array}{c}|\nu _e=|\nu _1\mathrm{cos}\theta +|\nu _2\mathrm{sin}\theta ,\hfill \\ |\nu _\mu =|\nu _1\mathrm{sin}\theta +|\nu _2\mathrm{cos}\theta ,\hfill \end{array}$$
(42)
where $`\theta `$ is the neutrino mixing angle. Suppose that at initial time, $`t=0`$, we create a pure weak eigenstate (as is typically the case), such as $`\nu _\mu `$ from the decay $`\pi ^+\mu ^+\nu _\mu `$, i.e., $`|\nu (0)=|\nu _\mu `$. Since $`|\nu _\mu `$ is a superposition of mass eigenstates, each of which propagates with its own time dependence, then at a later time the state may have oscillated into $`|\nu _e`$, which can be identified by its interaction, e.g., $`\nu _ene^{}p`$. To quantify this, after a time $`t`$ the state $`|\nu (0)`$ will have evolved into
$$|\nu (t)=|\nu _1\mathrm{sin}\theta e^{iE_1t}+|\nu _2\mathrm{cos}\theta e^{iE_2t},$$
(43)
where $`E_1=\sqrt{p^2+m_1^2}p+\frac{m_1^2}{2p}`$ and $`E_2=\sqrt{p^2+m_2^2}p+\frac{m_2^2}{2p}`$, and we have assumed that the neutrino is highly relativistic, $`pm_{1.2}`$, as is typically the case. After traveling a distance $`L`$ the oscillation probability, i.e., the probability for the neutrino to interact as a $`\nu _e`$, is
$$\begin{array}{ccc}P_{\nu _\mu \nu _e}(L)\hfill & =& |\nu _e|\nu (t)|^2\hfill \\ & =& \mathrm{sin}^2\theta \mathrm{cos}^2\theta |e^{iE_1t}+e^{iE_2t}|^2\hfill \\ & =& \mathrm{sin}^22\theta \mathrm{sin}^2\left(\frac{\mathrm{\Delta }m^2L}{4E}\right)=\mathrm{sin}^22\theta \mathrm{sin}^2\left(\frac{1.27\mathrm{\Delta }m^2(\mathrm{eV}^2)L(\mathrm{km})}{E(\mathrm{GeV})}\right),\hfill \end{array}$$
(44)
where $`Lt`$, $`E=p`$ and $`\mathrm{\Delta }m^2=m_2^2m_1^2`$. The probability for the state to remain a $`\nu _\mu `$ is $`P_{\nu _\mu \nu _\mu }(L)=1P_{\nu _\mu \nu _e}(L)`$.
There are two types of oscillation searches:
1. Appearance experiments, in which one looks for the appearance of a new flavor, e.g., of $`\nu _e`$ or $`\nu _\tau `$ (i.e., of $`e^{}`$ or $`\tau ^{}`$) in an initially pure $`\nu _\mu `$ beam.
2. Disappearance experiments, in which one looks for a change in, e.g., the $`\nu _\mu `$ flux as a function of $`L`$ and $`E`$.
Even with more than two types of neutrino, it is often a good approximation to use the two neutrino formalism in the analysis of a given experiment. However, a more precise or general analysis should take all three neutrinos into account<sup>19</sup><sup>19</sup>19If there are light sterile neutrinos which mix with the active neutrinos one should generalize to include their effects as well.. The general lepton mixing for three families contains three angles $`\theta _{12},\theta _{13},`$ and $`\theta _{23}`$; one Dirac CP-violating phase $`\delta `$; and two Majorana phases $`\beta _{1,2}`$. (The latter do not enter the oscillation formulae.) In the basis in which the charged leptons are mass eigenstates, the neutrino mass eigenstates $`\nu _i`$ and the weak eigenstates $`\nu _a=(\nu _e,\nu _\mu ,\nu _\tau )`$ are related by a unitary transformation $`\nu _a=V_{ai}\nu _i`$, where $`V`$is the $`3\times 3`$ lepton mixing matrix $`V_{PMNS}`$, which is parametrized as
$$\begin{array}{ccc}V_{PMNS}\hfill & =& \left[\begin{array}{ccc}c_{12}c_{13}& s_{12}c_{13}& s_{13}e^{i\delta }\\ s_{12}c_{23}c_{12}s_{23}s_{13}e^{i\delta }& c_{12}c_{23}s_{12}s_{23}s_{13}e^{i\delta }& s_{23}c_{13}\\ s_{12}s_{23}c_{12}c_{23}s_{13}e^{i\delta }& c_{12}s_{23}s_{12}c_{23}s_{13}e^{i\delta }& c_{23}c_{13}\end{array}\right]\hfill \\ & & \\ & & \times \mathrm{diag}(e^{i\beta _1/2},e^{i\beta _2/2},1),\hfill \end{array}$$
(45)
where $`c_{ij}=\mathrm{cos}\theta _{ij}`$ and $`s_{ij}=\mathrm{sin}\theta _{ij}`$. The oscillation probability for $`\nu _a\nu _b`$ after a distance $`L`$ is then
$$P_{\nu _a\nu _b}(L)\underset{ab}{\underset{}{=}}\mathrm{\Sigma }_j\left|V_{aj}V_{bj}\right|^2+\mathrm{Re}\mathrm{\Sigma }_{ij}V_{ai}V_{bi}^{}V_{aj}^{}V_{bj}e^{i\mathrm{\Delta }_{ij}L/2E},$$
(46)
where $`\mathrm{\Delta }_{ij}=m_i^2m_j^2`$.
## 7 Atmospheric Neutrinos
Many experiments have searched for neutrino oscillations at reactors, accelerators, and from astrophysical sources . Although the first indications of an effect involved the Solar neutrinos, the first unambiguous evidence came from the oscillations of atmospheric neutrinos. Atmospheric neutrinos are the product of pion and muon decays, which are produced in the upper layers of the atmosphere due to the interaction of primary cosmic rays. The data from the Kamiokande and Super-Kamiokande water $`\stackrel{ˇ}{\mathrm{C}}`$herenkov detectors indicated the disappearance of $`\mu `$-neutrinos. This was first seen in the ratio of the $`\nu _\mu /\nu _e`$ fluxes, and later confirmed by the zenith angle distribution of $`\nu _\mu `$ events. Results from other experiments such as MACRO and Soudan confirm the results, as does the recent long-baseline K2K experiment involving neutrinos produced at the KEK lab and observed in the Super-Kamiokande detector (Figure 14). The details of the observations have now established that the dominant effect is indeed oscillations of the $`\nu _\mu `$ into<sup>20</sup><sup>20</sup>20Oscillations into $`\nu _e`$ or a sterile neutrino as the dominant effect are excluded by the Super-Kamiokande data and reactor experiments. $`\nu _\tau `$, with near-maximal mixing ($`\mathrm{sin}^22\theta _{23}>0.92`$), and the difference between the mass squares of the two mass eigenstates of order $`|\mathrm{\Delta }m_{atm}^2|2\times 10^3`$ eV.
## 8 Solar Neutrinos
The first indications of neutrino oscillations were from the low number of $`\nu _e`$ events observed in the Davis Solar neutrino chlorine experiment, compared with the predictions of the Standard Solar Model (SSM) , shown in Figure 15.
The deficit was later confirmed (and that the observed events do indeed come from the Sun) in the Kamiokande and Super-Kamiokande water $`\stackrel{ˇ}{\mathrm{C}}`$herenkov experiments. These observed the high energy $`{}_{}{}^{8}B`$ neutrinos by the reaction $`\nu e^{}\nu e^{}`$, for which the cross section for oscillated $`\nu _{\mu ,\tau }`$ by the neutral current is about 1/6 that for $`\nu _e`$. The gallium experiments GALLEX, SAGE, and (later) GNO were sensitive to the entire Solar spectrum and also showed a deficit. Comparing the depletions in the various types of experiments, which constituted a crude measurement of the spectrum, showed that neutrino oscillations of $`\nu _e`$ into $`\nu _{\mu ,\tau }`$ or possibly a sterile neutrino were strongly favored over any plausible astrophysical uncertainty in explaining the results . Combining with information on the spectrum and time dependence of the signal (which could vary between day and night due to matter effects in the Earth) zeroed in on two favored regions for the oscillation parameters, one with small mixing angles (SMA) similar to the quark mixing, and one with large mixing angles (LMA), and with $`\mathrm{\Delta }m^210^510^4\text{ eV}^2`$. The LMA solution allowed oscillations to $`\nu _{\mu ,\tau }`$ but not to sterile neutrinos as the primary process (they were differentiated by the sensitivity to $`\nu _{\mu ,\tau }`$ in the water $`\stackrel{ˇ}{\mathrm{C}}`$herenkov experiments), while the SMA solution allowed both<sup>21</sup><sup>21</sup>21Big bang nucleosynthesis also disfavored the LMA solution for sterile neutrinos..
The situation was clarified by the Sudbury Neutrino Observatory (SNO) experiment . SNO uses heavy water $`D_2O`$, and can measure three reactions,
$$\begin{array}{c}\nu _e+De^{}+p+p,\\ \nu _{e,\mu ,\tau }+D\nu _{e,\mu ,\tau }+p+n,\\ \nu _{e,\mu ,\tau }+e^{}\nu _{e,\mu ,\tau }+e^{}.\end{array}$$
In the first (charged current) reaction, the deuteron breakup can be initiated only by the $`\nu _e`$, while the second (neutral current) reaction can be initiated by neutrinos of all the active flavors with equal cross section. SNO could therefore determine the $`\nu _e`$ flux arriving at the detector and the total flux into active neutrinos (which would equal the initially produced $`\nu _e`$ flux if there are no oscillations into sterile neutrinos) separately, as in Figure 16. SNO observed that the total flux is around three times that of $`\nu _e`$, establishing that oscillations indeed take place, and that the total flux agrees well with the predictions of the SSM. They also verified the LMA solution and that the Solar mixing angle $`\theta _{12}`$ is large but not maximal, i.e., $`\mathrm{sin}^22\theta _{12}0.8`$.
The Solar neutrino conversions are actually not due just to the vacuum oscillations as the neutrinos propagate from the Sun to the Earth. One must also take into account the coherent forward scattering of the neutrinos from matter in the Sun (and in the Earth at night), which introduces the analog of an index of refraction. This is of order $`G_F`$ rather than $`G_F^2`$, and it distinguishes $`\nu _e`$, $`\nu _{\mu ,\tau }`$, $`\overline{\nu }_e`$, $`\overline{\nu }_{\mu ,\tau }`$, and sterile neutrinos $`\nu _s`$ because of their different weak interactions. The propagation equation in electrically neutral matter is
$$i\frac{d}{dt}\left(\begin{array}{c}\nu _e\\ \nu _{\mu ,\tau ,s}\end{array}\right)=\left(\begin{array}{cc}\frac{\mathrm{\Delta }m^2}{2E}\mathrm{cos}2\theta +\sqrt{2}G_Fn& \frac{\mathrm{\Delta }m^2}{2E}\mathrm{sin}2\theta \\ \frac{\mathrm{\Delta }m^2}{2E}\mathrm{sin}2\theta & \frac{\mathrm{\Delta }m^2}{2E}\mathrm{cos}2\theta \sqrt{2}G_Fn\end{array}\right),$$
where
$$n=\{\begin{array}{ccc}n_e\hfill & \text{ for}& \hfill \nu _e\nu _{\mu ,\tau },\\ n_e\frac{1}{2}n_n\hfill & \text{ for}& \hfill \nu _e\nu _s,\end{array}$$
and where $`n_e(n_n)`$ is the number density of electrons (neutrons). This reduces to the vacuum oscillation case (with parameters $`\theta `$ and $`\mathrm{\Delta }m^2`$) for $`n=0`$. Under the right conditions, the matter effect can greatly enhance the transitions. In particular, if the Mikheyev-Smirnov-Wolfenstein (MSW) resonance condition $`\frac{\mathrm{\Delta }m^2}{2E}\mathrm{cos}2\theta =\sqrt{2}G_Fn`$ is satisfied, the diagonal elements vanish and even small vacuum mixing angles lead to a maximal effective mixing angle. In practice $`n`$ decreases as the neutrinos propagate from the core of the Sun where they are produced to the outside. For the LMA parameters, the higher energy Solar neutrinos all encounter a resonance layer. It is hoped that a future precise experiment sensitive to the lower energy $`pp`$ neutrinos will observe the transition between the vacuum and MSW dominated regimes.
Recently, the KamLAND long-baseline reactor experiment in Japan has beautifully confirmed the LMA solution, free of any astrophysical uncertainties, by observing a depletion and spectral distortion of $`\overline{\nu }_e`$ produced by power reactors located $`O(100)`$ km away. The combination of KamLAND and the Solar experiments also limits the amount of oscillations into sterile neutrinos that could accompany the dominant oscillations into $`\nu _{\mu ,\tau }`$.
## 9 Neutrino Oscillation Patterns
The neutrino oscillation data is sensitive only to the mass-squared differences. The atmospheric and Solar neutrino mass-squares are given by
$$\begin{array}{c}|\mathrm{\Delta }m_{32}^2||\mathrm{\Delta }m_{atm}^2|2\times 10^3\text{ eV}^2,\\ \mathrm{\Delta }m_{21}^2\mathrm{\Delta }m_{solar}^28\times 10^5\text{ eV}^2.\end{array}$$
Vacuum oscillations depend only on the magnitude of $`\mathrm{\Delta }m^2`$, so the sign of $`\mathrm{\Delta }m_{atm}^2`$ is not known, while MSW (matter) effects establish that $`\mathrm{\Delta }m_{solar}^2>0`$. The respective mixing angles are $`\mathrm{sin}^22\theta _{23}>0.92`$ (90%), consistent with maximal; and $`\mathrm{sin}^22\theta _{12}0.8`$, i.e., $`\mathrm{tan}^2\theta _{12},=0.40_{0.07}^{+0.09}`$, which is large but not maximal. On the other hand, the third angle is small, $`\mathrm{sin}^2\theta _{13}<0.03`$ (90%) from short-baseline ($`1`$ km) reactor disappearance limits, especially CHOOZ .
The LSND experiment at Los Alamos has claimed evidence for oscillations, especially $`\overline{\nu }_\mu \overline{\nu }_e`$, with $`|\mathrm{\Delta }m_{\mathrm{LSND}}^2|\begin{array}{c}>\hfill \\ \hfill \end{array}1\text{ eV}^2`$ and small mixing. This has not been confirmed by the KARMEN experiment, but there is a small parameter region allowed by both. The Fermilab MiniBooNE experiment is currently running, and should be able to confirm or exclude the LSND results. If LSND is confirmed, the most likely explanation would be that there is mixing of a fourth neutrino with the three known ones. This would have to be sterile because the invisible $`Z`$ width excludes a fourth light active neutrino.
### 9.1 Three Neutrino Patterns
If the LSND results are not confirmed, then the remaining data can be described by oscillations amongst the three light active neutrinos. Let us choose a phase convention in which the mass eigenvalues are real and positive, and label the states such that $`m_1<m_2`$ is responsible for the Solar oscillations, while the atmospheric oscillations are due to the $`32`$ mass-squared difference. Since neither the sign of $`\mathrm{\Delta }m_{32}^2`$ nor the absolute scale of the masses is known, there are several possible patterns for the masses. These include:
The normal (or ordinary) hierarchy, i.e., $`m_1m_2m_3`$, as in the left diagram in Figure 17. In this case, $`m_3\sqrt{\mathrm{\Delta }m_{atm}^2}0.040.05`$ eV, $`m_2\sqrt{\mathrm{\Delta }m_{sol}^2}0.009`$ eV, and $`m_10`$.
The inverted (or quasi-degenerate) hierarchy, i.e., $`m_1m_2\sqrt{\mathrm{\Delta }m_{atm}^2}0.040.05\text{ eV}m_3`$, as in the right diagram in Figure 17.
The degenerate case, i.e., $`m_1m_2m_3`$, with small splittings responsible for the oscillations.
Of course, these are only limiting cases. One can interpolate smoothly from cases (1) and (2) to the degenerate case by increasing the mass of the lightest neutrino.
### 9.2 Four Neutrino Patterns
If the LSND data is confirmed, then there is most likely at least one light sterile neutrino. In the case of four neutrinos the mass-squared difference for LSND is $`\mathrm{\Delta }m_{LSND}^21`$ eV<sup>2</sup>. In the $`2+2`$ pattern in Figure 18 the sterile component must be mixed in significantly with either the Solar or atmospheric pair. However, it is well established that neither the Solar nor the atmospheric oscillations are predominantly into sterile states, and this scheme is excluded. In the $`3+1`$ schemes a predominantly sterile state is separated from the three predominantly active states. This is consistent with the Solar and atmospheric data, but excluded when reactor and accelerator disappearance limits are incorporated . However, some 5 $`\nu `$ (i.e., $`3+2`$) patterns involving mass splittings around 1 eV<sup>2</sup> and 20 eV<sup>2</sup> are more successful . All of these run into cosmological difficulties, although there are some (highly speculative/creative) loopholes .
## 10 Conclusions
Non-zero neutrino mass is the first necessary extension of the Standard Model. Most extensions of the SM predict non-zero neutrino masses at some level, so it is difficult to determine their origin. Many of the promising mechanisms involve very short distance scales, e.g., associated with grand unification or string theories. There are many unanswered questions. These include:
* Are the neutrinos Dirac or Majorana? Majorana masses, especially if associated with a seesaw or Higgs triplet, would allow the possibility of leptogenesis, i.e., that the observed baryon asymmetry was initially generated as a lepton asymmetry associated with neutrinos, and later converted to a baryon (and lepton) asymmetry by nonperturbative electroweak effects. The observation of neutrinoless double beta decay would establish Majorana masses (or at least $`L`$ violation), but foreseeable experiments will only be sensitive to the inverted or degenerate spectra. If the neutrinos are Dirac, this would suggest that additional TeV scale symmetries or string symmetries/selection rules are forbidding Majorana mass terms.
* What is the absolute mass scale (with implications for cosmology)? This is very difficult, but ordinary and double beta decay experiments, as well as future CMB experiments, may be able to establish the scale.
* Is the hierarchy ordinary or inverted? Possibilities include matter effects in future long-baseline experiments or in the observation of a future supernova, the observation of $`\beta \beta _{0\nu }`$ if the neutrinos are Majorana, and possibly effects in future Solar neutrino experiments.
* What is $`\theta _{13}`$? This is especially important because the observation of leptonic CP violation requires a non-zero $`\theta _{13}`$. This will be addressed in a program of future reactor and long-baseline experiments, and possibly at a dedicated neutrino factory (from a muon storage ring).
* Why are the mixings large, while those of the quarks are small? For example, the simplest grand unified theory seesaws would lead to comparable mixings, although this can be evaded in more complicated constructions.
* If the LSND result is confirmed, it will suggest mixing between ordinary and sterile neutrinos, presenting a serious challenge both to particle physics and cosmology, or imply something even more bizarre, such as CPT violation.
* Are there any new $`\nu `$ interactions or anomalous properties such as large magnetic moments? Most such ideas are excluded as the dominant effect for the Solar and atmospheric neutrinos, but could still appear as subleading effects.
Answering these questions and unraveling the origin of the masses is therefore an important and exciting probe of new particle physics. Future Solar neutrino experiments, observations of the neutrinos from a future core-collapse supernova, and observation of high energy neutrinos in large underground/underwater experiments would also be significant probes not only of the neutrino properties but also the underlying astrophysics.
## Acknowledgements
This work was supported in part by CONACYT (México) contract 42026–F, by DGAPA–UNAM contract PAPIIT IN112902, and by the U.S. Department of Energy under Grant No. DOE-EY-76-02-3071. |
warning/0506/math0506226.html | ar5iv | text | # Thinness for Scalar-Negative Singular Yamabe Metrics
## 1 Introduction
### 1.1 Singular Yamabe problem
Yamabe problem was to prove that for any compact Riemannian manifold $`(M,g)`$ of dimension $`n3`$ we can find a metric conformal to $`g`$ with a constant scalar curvature by solving a certain variational problem. This was proved in three subsequent contributions by Trudinger , Aubin , and Schoen , see also , for an exposition. Later other proofs were given, see for a survey. To a large extent, it is the Yamabe problem that stimulated the development of the modern geometric analysis. An intensive work was done on Yamabe problem for manifolds with boundary. In this case seeks a conformal metric with constant curvatures in the interrior and on the boundary, cf. e.g. , , . There were also generalisations to non-Riemannian settings , .
In 1988 Schoen and Yau arrived at a different (non-variational) problem , . Namely, is it possible to characterise domains on the unit sphere admitting a conformal deformation of the standard metric to a complete metric with a constant scalar curvature? Schoen and Yau were led to this problem by their research on geometry and topology of locally conformally flat manifolds and were mainly interested in the case of non-negative curvature. The case of negative scalar curvature goes back to an early paper by Loewner and Nirenberg . Further motivations for the problem can be found in #36 from , .
The goal of this paper is to introduce methods of nonlinear potential theory to this problem. They allow, in particular, to solve the problem in the negative curvature case. Our Theorem 1.1 states that the conformal deformation to a complete scalar-negative metric is possible in $`\mathrm{\Omega }`$ if and only if its complement is not thin. Thinness, see sec.1.2, is a basic concept in potential theory first introduced by Wiener in his works on the classical Dirichlet problem. Developments in nonlinear potential theory easily allow to relate thinness with geometric properties. Let us now describe the previous work and our results on the problem in more details.
Under the conformal change of metric $`g=u^{4/(n2)}\stackrel{\text{ }_{}}{g}`$, $`n3`$, the scalar curvature changes according to the formula
(1.1)
$$R(g)=u^{(n+2)/(n2)}\left(\frac{4(n1)}{n2}\mathrm{\Delta }u+R(\stackrel{\text{ }_{}}{g})u\right).$$
Here $`\mathrm{\Delta }u=\mathrm{div}\left(\mathrm{grad}u\right)`$ is the Laplace-Beltrami operator on $`(𝐒^n,\stackrel{\text{ }_{}}{g})`$ and $`R(\stackrel{\text{ }_{}}{g})=n(n1)`$ is the scalar curvature of the standard metric $`\stackrel{\text{ }_{}}{g}`$ induced by the embedding $`𝐒^n𝐑^{n+1}`$. Thus analytically for given $`\mathrm{\Omega }𝐒^n`$ and $`R\{1,0,1\}`$ one seeks a smooth solution to the following problem:
$`{\displaystyle \frac{4(n1)}{n2}}\mathrm{\Delta }uR(\stackrel{\text{ }_{}}{g})u+Ru^{(n+2)/(n2)}=0`$ $`\mathrm{in}`$ $`\mathrm{\Omega }`$
(1.2) $`u>0`$ $`\mathrm{in}`$ $`\mathrm{\Omega }`$
$`u^{4/(n2)}\stackrel{\text{ }_{}}{g}\mathrm{is}\mathrm{complete}\mathrm{metric}`$ $`\mathrm{in}`$ $`\mathrm{\Omega }.`$
The case $`R=1`$ is regarded as the hardest among the three. The present work focuses on the negative curvature case
(1.3)
$$R=1.$$
Survey by McOwen describes the progress on the problem and open problems. Let us explain how our results fit in the general picture. We set $`K=𝐒^n\mathrm{\Omega }`$.
Investigations of the negative curvature case were started in 1974 by Loewner and Nirenberg . They proved that if the problem (1.1), (1.3) admits a solution then the complement of $`\mathrm{\Omega }`$ must satisfy
$$^{(n2)/2}(K)=\mathrm{}.$$
Here $`^\alpha `$ denotes the Hausdorff $`\alpha `$-measure. Their work together with Aviles and Veron showed that if $`K`$ is a smooth submanifold of $`𝐒^n`$ of the corresponding dimension $`k>(n2)/2`$ then problem (1.1), (1.3) has a solution. Mazzeo showed, in particular, that for such $`K`$ the solution is unique. Finn , , established the solvability under weaker conditions on $`K`$. Namely he required that it has a structure similar to (actually, more general than) Lipschitz submanifold of the corresponding dimension, see also . The gap between such requirents and the sufficient condition of Loewner and Nirenberg still remained broad. In section 1.3 we show how all these results follow from Theorem 1.1.
This paper studies (1.1), (1.3) in dimensions $`n3`$. In the case of $`𝐒^2`$ the complete metric conformal to $`\stackrel{\text{ }_{}}{g}`$ and having the constant negative curvature is called the Poincare metric. The equation for Poincare metric is slightly different from (1.1). Mazzeo and Taylor proved that the Poincare metric in $`\mathrm{\Omega }𝐒^2`$ always exists provided the complement $`K`$ has at least two distinct points.
In 1988 Schoen and Yau were led by their research on locally conformally flat maniofolds to the case $`R0`$. They found that a necessary condition for solvability of (1.1) with $`R0`$, as oposed to the case (1.3), is smallness of $`K`$. For example, they proved that if the solution exists then the Newtonian capacity of $`K`$ must vanish. They also established that the solvabilty in this case implies that
$$^{\epsilon +(n2)/2}(K)=0\mathrm{for}\mathrm{all}\epsilon >0.$$
Thus the Hausdorff dimension $`(n2)/2`$ separates the cases of the negative and non-negative curvature. Similarly to the negative curvature case, the existence of a solution is known at the moment only in cases when $`K`$ has much more structure that vanishing Hausdorff measure or capacity. Despite the similarity in statements, the results in the case $`R=1`$ are much more difficult to prove. In a seminal paper Schoen established the existence of (1.1) with $`R=1`$ when $`K`$ is a finite number (at least two) of points. Mazzeo and Pacard generalising earlier results , , extended Schoen’s result to the case when $`K`$ is a finite number of disjoint smooth submanifolds of the dimension $`k(n2)/2`$. There is also a construction of a solution using Kleinian groups in the case when $`K`$ is a certain Cantor-type set
The case $`R=0`$ is easier becuase equation (1.1) becomes linear. In this case it is known that the solution exists provided that $`K`$ is essentially a finite union of Lipschitz submanifolds of dimension $`k(n2)/2`$ , , .
In the paper we are interested only in the basic problem of existence for (1.1). However, other questions about solutions of (1.1) can be asked as well. For example problems of uniqueness, asymptotic behaviour of $`u`$ near $`\mathrm{\Omega }`$, structure of moduli space of solutions, gluing different solutions, are investigated in , , , . Some of the results mentioned above hold for more general manifolds than $`𝐒^n`$. The result directly related to the present paper was proved by Aviles and McOwen , , . They established that an open subset $`\mathrm{\Omega }`$ of any closed Riemannian manifold $`(M,g)`$ admits a complete metric with the constant negative scalar curvature conformal to $`g`$ provided $`K`$ is a finite union of closed smooth submanifolds of dimensons $`k>(n2)/2`$. In a future publication we introduce a suitable capacity and extend our Theorem 1.1 to more general manifolds.
### 1.2 Main theorem
We investigate the solvability of (1.1) by attracting ideas from the nonlinear potential theory. Let us recall a fundamental result form the classical potential theory for the Laplace equation. This is the Wiener test for the classical Dirichlet problem for harmonic functions . Wiener theorem states that the Dirichlet problem
$$\{\begin{array}{ccc}\hfill \mathrm{\Delta }w=0& \mathrm{in}& D\hfill \\ \hfill w=f& \mathrm{on}& D\hfill \end{array}$$
in a bounded domain $`D𝐑^n`$, $`n3`$, is solvable for all boundary data $`fC(D)`$ if and only if $`𝐑^nD`$ is not thin. Explicitly the latter means that
$$_0^1\frac{cap(B(x,r)D)}{cap(B(x,r))}\frac{dr}{r}=+\mathrm{}\mathrm{for}\mathrm{any}xD.$$
Here $`1`$ can be replaced by any $`\delta >0`$, and $`cap`$ is the classical (electrostatic) capacity. Our main Theorem 1.1 states that problem (1.1), (1.3) admits a solution if and only if a Wiener-type test with a certain capacity holds.
Let us scetch the definition of the capacity apropriate for problem (1.1), see section 2.2 for more details. Take a compact set $`E𝐒^n`$, $`n3`$, with
$$\mathrm{diam}_{\stackrel{\text{ }_{}}{g}}(E)\pi /3.$$
After a rotation we can assume that such $`E`$ lies in the southern hemisphere. We set
$`𝒞(E)`$ $`=`$ $`inf\left\{{\displaystyle _{𝐒^n}}\left|^2\phi \right|^{(n+2)/4}𝑑vol_{\stackrel{\text{ }_{}}{g}}\right\},`$
$`{\displaystyle \frac{1}{(n+2)/(n2)}}+{\displaystyle \frac{1}{(n+2)/4}}=1.`$
Here symbols $`dvol_{\stackrel{\text{ }_{}}{g}}`$, $``$, and $`||`$, stand respectively for the volume element, connection, and norm with respect to the metric $`\stackrel{\text{ }_{}}{g}`$. The infimum is taken over all $`\phi C^{\mathrm{}}(𝐒^n)`$ such that $`\phi |_E1`$ and $`\phi 0`$ on the northern hemisphere. Essentially, $`𝒞`$ is the Bessel capacity for the Sobolev space $`W^{2,(n+2)/4}(𝐑^n)`$. Bessel capacities have been intensively investigated in the nonlinear potential theory. Nonlinear potential theory originates in early works of Maz’ya and Serrin in the 1960s and was extensively developed later in 1970s and 1980s by many authors. Our paper heavily relies on it. The main references will be monographs by Adams and Hedberg , Maz’ya , and Ziemer . There the reader can also find a rich bibliography and historical notes Now we state the main theorem.
###### Theorem 1.1
Let $`\mathrm{\Omega }𝐒^n`$, $`n3`$, be an open set and $`K=𝐒^n\mathrm{\Omega }`$. Then the following properties are equivalent:
(i) In $`\mathrm{\Omega }`$ there exists a complete metric with constant negative scalar curvature conformal to $`\stackrel{\text{ }_{}}{g}`$.
(ii) The compactum $`K`$ is not thin, that is for any $`pK`$
(1.4)
$$_0^{1/2}\left(\frac{𝒞(B(p,r)K)}{𝒞(B(p,r))}\right)^{2/(n2)}\frac{dr}{r}=+\mathrm{}.$$
Wiener test (1.4) is a capacitary condition on $`K`$. Geometric properties of the capacity $`𝒞`$ are well understood due to investigations in nonlinear potential theory. Using the information avialable there, we show in section 1.3 that more transparent geometric results easily follow from Theorem 1.1.
In view of Theorem 1.1 it would be interesting to clarify how the condition
$$𝒞(𝐒^n\mathrm{\Omega })=0$$
relates to the conformal deformation to nonnegative scalar curvature $`R0`$. In we apply potential theory ideas to the scalar flat case $`R=0`$. In this situation as opposed to Theorem 1.1 the set $`K`$ should be small.
We mention that ideas from potential theory have been used in conformal geometry before. For example Schoen and Yau , used capacity related to Sobolev space $`W^{1,q}(𝐑^n)`$. Capacity $`𝒞`$ was implicitely used at some stage in to prove distribution removability of isolated singularities for the equation in (1.1) with $`R=1`$.
###### Remark 1.2
Intuitevely, the completeness condition forces solutions of (1.1) to blow up in some sense near $`\mathrm{\Omega }`$. Dhersin and LeGall considered the problem
(1.5)
$$\{\begin{array}{ccc}\hfill \mathrm{\Delta }uu^2=0& \mathrm{in}& D\hfill \\ \hfill u(x)+\mathrm{}& \mathrm{when}& xD\hfill \end{array}$$
in general domains $`D𝐑^n`$. They proved that a Wiener test characterises domains $`D`$ for which (1.5) is solvable. Their paper provided an important inspiration for our work, although the consideration in is based on probabilistic methods. In fact, there is a strong connection between $`u`$ in (1.5) and a certain branching random process (so-called Brownian snake) , . To find an adequate probabilistic interpretation for $`p>2`$ is an important open problem in the area , . However, in we established a solvability criterion for
$$\mathrm{\Delta }uu^p=0$$
blowing up at the boundary in the full range of $`p>1`$. Relying there on entirely analytic ideas we establish estimates for solutions in terms of the capacity associated with the variational integral
$$_{𝐑^n}|D^2\phi |^p^{},\phi C_0^{\mathrm{}}(𝐑^n),\frac{1}{p}+\frac{1}{p^{}}=1.$$
Comparison of Theorem 1.1 with the condition from implies that any conformal factor in (1.1), (1.3) must blow up pointwisely near $`K`$. Iterplay between Brownian motion on a Riemannian manifold and geometric properties of the manifold is an important subject. We refer to Grigor’yan’s survey for an exposition of the classical and new results. It would be interesting to understand the relation between super-Brownian motion, Brownian snakes and the geometric problems from this paper (or properties of Riemannian manifolds in general). Some results in this direction can be found in .
### 1.3 Examples
We illustrate how Theorem 1.1 allows to establish the existence for (1.1), (1.3) in concrete situations. In particular, apparently all necessary or sufficient conditions from previous papers can be easily derived from (1.4). The reason for this is that capacity $`𝒞`$ had appeared before in different problems related to the interaction between nonlinear potentials and the Littlewood-Paley theory. As a result, it was intensively studied in the 1970s-1980s and its geometric properties are well known. They are carefully documented e.g. in , .
###### Example 1.3
The necessity of Loewner-Nirenberg condition
$$^{(n2)/2}(K)=+\mathrm{}$$
for solvability of (1.1), (1.3) follows immediately from Theorem 1.1 and the implication
$$^{(n2)/2}(K)<+\mathrm{}𝒞(K)=0,$$
valid for our capacity $`𝒞`$, see and (2.15) in section 2.2 below.
###### Example 1.4
Assume that $`K`$ is a smooth immersed submanifold in $`𝐒^n`$ of dimension $`k`$,
(1.6)
$$k>(n2)/2.$$
Then (1.1), (1.3) has a solution. In fact, according to Theorem 1.1 we need to show that (1.4) holds. Fix any $`pK`$. By the definition of immersion there exists an open smooth submanifold $`E`$ of dimension $`k`$ embedded in $`𝐒^n`$ such that $`pE`$ and $`EK`$. Therefore
$$𝒞(KB(p,r))𝒞(EB(p,r)).$$
The exponential map $`\mathrm{exp}_p`$ is a diffeomorphism of a neighbourhood of the origin $`0`$ in $`T_p𝐒^n`$. In the sufficiently small neighbourhood of $`p`$ the smooth submanifold $`E`$ is well approximated by the image under $`\mathrm{exp}_p`$ of a neighbourhood of $`0`$ in $`T_pE`$, $`T_pET_p𝐒^n`$. Capacities of a set and its image under a diffeomorphism are equivalent . Hence utilising the scaling property (2.13) we find a small number $`r_0>0`$, such that
$$𝒞(EB(p,r))C(E)\left(\frac{r}{r_0}\right)^{(n2)/2}𝒞(EB(p,r_0))\mathrm{for}r(0,r_0).$$
Here the constant $`C(E)>0`$ depends on the smoothness of $`E`$. Capacity of the ball can be estimated by
$$𝒞(B(p,r))r^{(n2)/2}\mathrm{for}r(0,r_0),$$
see (2.16). Now (1.4) follows.
###### Example 1.5
Let $`d=\epsilon +(n2)/2`$, $`\epsilon >0`$. Assume that for any $`xK`$ there is a positive constant $`C`$ such that
(1.7)
$$_{\mathrm{}}^d(KB(x,r))Cr^d$$
for all $`r`$ near $`0`$. That is, the Hausdorff $`d`$-content of $`K`$ is big at all small scales. Then the conformal metric from Theorem1.1 exists. This follows at once from (2.14), (2.16). Density condition (1.7) allows to recover the results of Finn , , about sets $`K`$ with stratified cone-type tangent structure. Also (1.7) allows to establish existence of the complete metrics in the cases when $`K`$ satisfies different geometric conditions invariant with respect to quasiconformal maps, see e.g. .
###### Example 1.6
Let $`K`$ be the Lebesgue cusp. That is for a fixed $`\rho >0`$ and for a continuous positive nondecreasing function $`h`$ on the real line, $`h(r)=O(r)`$, $`r0`$, we set
$$T_h=\{x𝐑^n:0x_n\rho ,\left(x_1^2+\mathrm{}+x_{n1}^2\right)^{1/2}h(x_n)\}.$$
Then define $`K`$ to be the preimage of $`T_h`$ under the stereogrphic projection,
$$K=\sigma ^1(T_h).$$
The existence of the singular conformal metric $`g`$ from Theorem 1.1 in $`𝐒^nK`$ depends on the dimension $`n`$. If $`n=3`$ then $`g`$ always exists. For higher dimensions $`g`$ exists if and only if
$`{\displaystyle _0^1}{\displaystyle \frac{dr}{r\mathrm{log}\left(r/h(r)\right)}}=+\mathrm{}\mathrm{for}n=4`$
$`{\displaystyle _0^1}\left({\displaystyle \frac{h(r)}{r}}\right)^{(n4)/(n2)}{\displaystyle \frac{dr}{r}}=+\mathrm{}\mathrm{for}n>4.`$
Indeed, we only need to check that (1.4) holds for the South pole $`S=\sigma ^1(0)`$. To verify (1.4), first recall that for the cyllinder
$$\mathrm{\Pi }=\sigma ^1\left((\delta ,\delta )\times \mathrm{}\times (\delta ,\delta )\times (r/2,r)\right),4\delta <r,$$
with $`r>0`$ small enough, the capacity is given by the following formulae , Ch. 9:
$`𝒞(\mathrm{\Pi })`$ $``$ $`𝒞(B(S,r))\mathrm{for}n=3,`$
$`𝒞(\mathrm{\Pi })`$ $``$ $`{\displaystyle \frac{1}{\mathrm{log}\left(r/\delta \right)}}𝒞(B(S,r))\mathrm{for}n=4,`$
$`𝒞(\mathrm{\Pi })`$ $``$ $`\left({\displaystyle \frac{\delta }{r}}\right)^{(n4)/2}𝒞(B(S,r))\mathrm{for}n>4.`$
Now just apply elementary estimate (2.28).
### 1.4 Organisation of the paper
In section 2 we introduce the capacity and use it to prove some preliminary estimates for solutions of the equation. We also describe there the unique feature of equation (1.1), (1.3). Namely, the existence of a finite maximal solution $`u_\mathrm{\Omega }`$ dominating all other solutions pointwisely.
In section 3 we prove the crucial estimates for solutions of (1.1), (1.3) in terms of the capacity. The principal difficulty in the proof of the main Theorem 1.1 is the analysis of the completeness condition in (1.1). This involves understanding the behaviour of the conformal factor $`u`$ near $`\mathrm{\Omega }`$ under no assumptions (say, when proving (i)$``$(ii) in Theorem 1.1) on the structure of $`\mathrm{\Omega }`$. The first estimate in section 3, Theorem 3.1, controlls the solution pointwisely away from $`\mathrm{\Omega }`$. The second estimate, Theorem 3.2, provides the integral control when we stay arbitrarily close to $`\mathrm{\Omega }`$.
With all this background in place we proceed to prove our main result, Theorem1.1, in section 4. The sufficiency of the Wiener test (1.4) will follow rather straighforwardly from the pointwise estimate from section 3. To prove the necessity we supose that the negation of (1.4) holds. In other words, suppose that the complement of $`\mathrm{\Omega }`$ is thin at some point. We will find a curve in $`\mathrm{\Omega }`$ approaching this point, such that its length with respect to $`u_\mathrm{\Omega }`$ (and hence with respect to any other solution of (1.1), (1.3)) is finite. How to construct such a curve without any assumptions on $`\mathrm{\Omega }`$? The key idea here is to reduce this issue to an integral estimate. To achieve this we bring in the estimate from Theorem 3.2.
Throughout this paper, we will use the notation
$$q=\frac{n+2}{n2},q^{}=\frac{n+2}{4},\frac{1}{q}+\frac{1}{q^{}}=1.$$
By $`g_E`$ we denote the Euclidean metric in $`𝐑^n`$ and by $`\stackrel{\text{ }_{}}{g}`$ the standard meric on the sphere induced by $`g_E`$. By $`B(p,r)`$ we denote the ball of radius $`r`$ centered at $`p`$ for $`\stackrel{\text{ }_{}}{g}`$ or $`g_E`$. It will be clear from the context which metric is taken. For an integer $`j`$ we put $`r_j=2^j`$. By $`B_j`$ we denote the dyadic ball in $`𝐑^n`$, $`B_j=B(0,r_j)`$. We denote the Green’s function for the Laplacian in $`B(0,R)𝐑^n`$ by $`G_R`$. By $`C`$, $`\stackrel{~}{C}`$, $`C_1`$, $`\mathrm{}`$, we denote positive constants depending only on the dimension. The value of $`C`$, $`\stackrel{~}{C}`$, $`C_1`$, $`\mathrm{}`$, may vary even within the same line. We write
$$AB(AB)$$
if
$$ACB(ACB)$$
for some $`C`$. We write
$$AB$$
if
$$ABA.$$
### 1.5 Acknowledgements.
I would like to thank Rick Schoen for many intersting and stimulating discussions on the problem. I also would like to thank Michael Struwe, Tom Ilmanen, Reiner Schätzle, Neil Trudinger, and Xu-Jia Wang for their interest and support.
## 2 Preliminaries on the equation and capacity
### 2.1 Equation on sphere and in space
Let $`g`$ be a metric on a manifold $`M`$ of dimension $`n`$, $`n3`$. The operator
$$_g=4\frac{n1}{n2}\mathrm{\Delta }_g+R(g)$$
from (1.1) is called conformal Laplacian . If we change the metric conformally
$$\widehat{g}=\phi ^{4/(n2)}g,$$
then
$`R(\widehat{g})=\phi ^{(n+2)/(n2)}_g\phi ,`$
$`_{\widehat{g}}v=\phi ^{(n+2)/(n2)}_g\left(\phi v\right).`$
More generally, let $`\stackrel{~}{M}`$ be another manifold with the metric $`\stackrel{~}{g}`$, and let $`f:M\stackrel{~}{M}`$ be a diffeomorphism. Assume that $`f`$ changes the metric conformally
$$f^{}\stackrel{~}{g}=\phi ^{4/(n2)}g.$$
Then
$`f^{}(R(\stackrel{~}{g}))=\phi ^{(n+2)/(n2)}_g\phi ,`$
$`f^{}(_{\stackrel{~}{g}}v)=\phi ^{(n+2)/(n2)}_g\left(\phi f^{}v\right).`$
Now, the stereographic projection $`\sigma :𝐒^n\{N\}𝐑^n`$ is a conformal diffeomorphism between $`(𝐒^n\{N\},\stackrel{\text{ }_{}}{g})`$ and $`(𝐑^n,g_E)`$ because
$`\left(\sigma ^1\right)^{}\stackrel{\text{ }_{}}{g}`$ $`=`$ $`\left({\displaystyle \frac{2}{1+|x|^2}}\right)^2g_E`$
$`=`$ $`\mathrm{{\rm Y}}^{4/(n2)}g_E`$
with
$$\mathrm{{\rm Y}}(x)=\left(\frac{2}{1+|x|^2}\right)^{(n2)/2}x𝐑^n.$$
According to the above formulae for the conformal changes we have the following correspondence.
Let $`\mathrm{\Omega }𝐒^n`$, $`N\mathrm{\Omega }`$, and let the function $`v`$ satisfy
(2.1)
$$v^{(n+2)/(n2)}_{\stackrel{\text{ }_{}}{g}}v=1,v>0\mathrm{in}\mathrm{\Omega }.$$
Then the function
(2.2) $`u(x)`$ $`=`$ $`\mathrm{{\rm Y}}(x)\left(\sigma ^1\right)^{}v(x)`$
$`=`$ $`\mathrm{{\rm Y}}(x)v(\sigma ^1x),x𝐑^n,`$
satisfies
$$u^{(n+2)/(n2)}_{g_E}u=1,u>0\mathrm{in}\sigma (\mathrm{\Omega }).$$
Thus after multiplication by a constant, $`u`$ satisfies
(2.3)
$$u>0,\mathrm{\Delta }uu^{(n+2)/(n2)}=0.$$
Conversely, for any solution $`u`$ of (2.3) defined in $`\sigma (\mathrm{\Omega })`$ set
$$v=\sigma ^{}\left(u/\mathrm{{\rm Y}}\right).$$
Then after multiplication by a constant, $`v`$ satisfies (2.1) in $`\mathrm{\Omega }\{N\}`$. Moreover we know that for any $`u`$ solving (2.3) in a neighbourhoud of infinity in $`𝐑^n`$, there exists a constant $`A>0`$ such that
$$u(x)=\frac{A}{|x|^{n2}}+o\left(\frac{1}{|x|^{n2}}\right),x\mathrm{}.$$
Hence we can extend $`v`$ to $`N`$ by continuity, remove the isolated singularity, and conclude that (2.1) holds.
Clearly the metric $`v^{4/(n2)}\stackrel{\text{ }_{}}{g}`$ is complete in $`\mathrm{\Omega }`$ for $`v`$ from (2.1) if and only if $`u^{4/(n2)}g_E`$ is complete in $`\sigma (\mathrm{\Omega })\{\mathrm{}\}`$ for the corresponding $`u`$ from (2.2), (2.3).
The previous discussion shows that the existence of the singular Yamabe metric in a domain on the sphere is equivalent to finding a complete solution of (2.3) in the exterior domain in $`𝐑^n`$. Let us describe the main features of equation (2.3) in $`𝐑^n`$. Omited proofs can be found for example in .
The crucial fact about solutions of (2.3) that will be used constantly in this paper is the elliptic comparison principle. As a consequence of this principle, local regularity estimates hold for $`u`$. In particular, if $`uL_{loc}^q`$ is a distributional solution of (2.3) then, in fact, $`uC_{loc}^{\mathrm{}}`$ and $`u`$ is the classical solution. Moreover, let $`u`$ be any solution of (2.3) in an open set $`O𝐑^n`$. Then
(2.4)
$$u(x)\frac{1}{\mathrm{dist}(x,O)^{(n2)/2}}\mathrm{for}\mathrm{all}xO.$$
This is an estimate uniform in $`u`$. It was first discovered by Keller and Osserman , and also follows from the comparison principle.
Estimate (2.4) combined with the elliptic Perron argument implies the existence of the finite solution $`u_O`$ which is maximal in $`O`$. It means that the inequality
$$uu_O\mathrm{in}O$$
holds for any other $`u`$ solving (2.3) in $`O`$. Clearly
$$u_{O_1}u_{O_2}\mathrm{in}O_2\mathrm{when}O_1O_2.$$
Let $`K_1`$, $`\mathrm{}`$, $`K_m`$ be compact sets in $`𝐑^n`$, let
$$K=K_1\mathrm{}K_m,$$
and let $`u`$, $`u_1`$, $`\mathrm{}`$, $`u_m`$ be the maximal solutions of (2.3) in $`K^c`$, $`K_1^c`$, $`\mathrm{}`$, $`K_m^c`$ respectively. Then the Hölder inequality and the comparison ensure that
(2.5)
$$m^{(n2)/(n+2)}\underset{i=1}{\overset{m}{}}u_iu\underset{i=1}{\overset{m}{}}u_i\mathrm{in}K^c.$$
If $`x_0O`$ and $`(\mathrm{\Omega })B(x_0,r)`$ is a smooth hupersurface for some $`r>0`$, then
(2.6) $`u_O(x)\mathrm{dist}(x,O)^{(n2)/2}`$ $``$ $`\left({\displaystyle \frac{n(n2)}{4}}\right)^{(n2)/4},`$
$`\mathrm{when}x(\mathrm{\Omega })B(x_0,r).`$
Asymptotic behaviour (2.6) holds in fact for any solution of (2.3) blowing up at $`(\mathrm{\Omega })B(x_0,r)`$.
Solutions to (2.3) exhibit the following dilation invariance: for all $`a>0`$ and $`r>0`$,
(2.7) $`u\mathrm{solves}(\text{2.3})\mathrm{in}B(0,r)`$ $``$ $`a^{(n2)/2}u(a)\mathrm{solves}(\text{2.3})`$
$`\mathrm{in}B(0,r/a).`$
Finally consider equation (2.1) on the sphere. As a direct consequence of the properties of the stereographic projection which we discussed above, there exists the maximal solution of (2.1) and the estimates analogous to (2.4)–(2.6) hold.
### 2.2 Capacity
In this paragraph we define the capacity $`𝒞`$ for subsets of the unit sphere $`𝐒^n`$. Esentially it is a particular Bessel capacity $`𝐂`$ in $`𝐑^n`$. Omited proofs of the statements about $`𝐂`$ can be found in monographs , , and . By $`𝐒_S^n`$ and $`𝐒_N^n`$ we denote southern and northern hemispheres.
Fix the spherical cup around the south pole $`S`$ by writing
(2.8)
$$𝐔=\{x𝐒^n:d_{\stackrel{\text{ }_{}}{g}}(S,x)\pi /3\},𝐔𝐒_S^n.$$
Take any compact set $`K𝐒^n`$ with
$$\mathrm{diam}_{\stackrel{\text{ }_{}}{g}}(K)\pi /3.$$
The rotation group $`SO(n+1)`$ acts transitively on $`𝐒^n𝐑^{n+1}`$. Map $`K`$ by a rotation $`\mathrm{\Phi }SO(n+1)`$ in a way that $`\mathrm{\Phi }(K)𝐔`$. Define
(2.9) $`𝒞(K)=inf\left\{{\displaystyle _{𝐒^n}}\right|^2\phi |^{(n+2)/4}dvol_{\stackrel{\text{ }_{}}{g}}:`$ $`\phi C^{\mathrm{}}(𝐒^n),\phi |_{𝐒_N^n}=0,`$
$`\phi |_{\mathrm{\Phi }(K)}1\}.`$
We will prove that different choices of $`\mathrm{\Phi }SO(n+1)`$ lead to equivalent capacities. First we give an alternative description of the capacity. Stereographic projection $`\sigma `$ is a smooth quasiisometry between $`𝐒_S^n`$ and $`B(0,1)𝐑^n`$. Hence
(2.10) $`𝒞(K)inf\left\{{\displaystyle _{𝐑^n}}\right|D^2\psi |^{(n+2)/4}dx:`$ $`\psi C_0^{\mathrm{}}(B(0,1)),`$
$`\psi |_{\sigma \mathrm{\Phi }(K)}1\}.`$
Let us introduce the corresponding capacity for sets in $`𝐑^n`$. For a compact set $`EB((0,1))𝐑^n`$ its Bessel capacity is defined as
(2.11) $`𝐂(E)=inf\left\{{\displaystyle _{𝐑^n}}\right|D^2\psi |^{(n+2)/4}dx:`$ $`\psi C_0^{\mathrm{}}(B(0,2)),`$
$`\psi |_E1\}.`$
Notice that the set $`\sigma \mathrm{\Phi }(K)`$ stays away from the boundary of the unit ball:
$$\sigma \mathrm{\Phi }(K)B(0,99/100).$$
Hence properties of Bessel capacities imply that for $`E=\sigma \mathrm{\Phi }(K)`$ the right hand sides of (2.11) and (2.10) are equivalent. Thus
(2.12)
$$𝒞(K)𝐂(\sigma \mathrm{\Phi }(K)).$$
Now take another $`\stackrel{~}{\mathrm{\Phi }}SO(n+1)`$, $`\stackrel{~}{\mathrm{\Phi }}(K)𝐔`$. The same variational procedure as (2.9) gives the new capacity $`\stackrel{~}{𝒞}(K)`$. Bessel capacity (2.11) of a compactum and of its image under a bi-Lipschitz homeomorphism are equivalent. Apply this to the locally bi-Lipschitz map
$$\sigma \stackrel{~}{\mathrm{\Phi }}\mathrm{\Phi }^1\sigma ^1:𝐑^n𝐑^n$$
which sends $`\sigma \mathrm{\Phi }(K)`$ to $`\sigma \stackrel{~}{\mathrm{\Phi }}(K)`$, and utilise (2.12) to derive that
$$𝒞(K)\stackrel{~}{𝒞}(K).$$
Clearly property (1.4) of the set to be not thin does not change when we pass to an equivalent capacity. Set functions $`𝒞`$ and $`𝐂`$ enjoy subadditvity and monotonicity properties. Standard scheme of axiomatic potential theory extends them to arbitrary sets as the outer measure.
Now we list some well-known metric estimates for the capacity. The following important scaling holds:
(2.13)
$$𝐂(tE)t^{(n2)/2}𝐂(E),t(0,1),EB(0,1).$$
Next, for $`\alpha >0`$ the Hausdorff $`\alpha `$-content of $`E𝐑^n`$ (or $`E𝐒^n`$) is defined as
$$_{\mathrm{}}^\alpha (E)=inf\underset{j}{}r_j^\alpha ,$$
where the infimum is taken over all coverings of $`E`$ by countable unions of euclidean balls $`\{B(x_j,r_j)\}`$ in $`𝐑^n`$ (or $`d_{\stackrel{\text{ }_{}}{g}}`$-balls in $`𝐒^n`$). The set function $`_{\mathrm{}}^\alpha `$ is subadditive and monotone. For the Hausdorff measure we have
$$^\alpha (E)=0_{\mathrm{}}^\alpha (E)=0.$$
There is a strong connection between the capacity and the Hausdorf content and measure. For any
$$\alpha >(n2)/2$$
there is a constant $`C(n,\alpha )>0`$ such that
(2.14)
$$\left(_{\mathrm{}}^\alpha (E)\right)^{(n2)/2}C(n,\alpha )𝐂(E)^\alpha ,EB(0,1).$$
Hence sets of the capacity $`0`$ have the Hausdorf dimension at most $`(n2)/2`$. In the converse direction the following implication holds for the Hausdorf measure:
(2.15)
$$^{(n2)/2}(E)<+\mathrm{}𝐂(E)=0$$
for $`EB(0,1)`$. According to (2.12) statements (2.14) and (2.15) also hold for $`𝒞`$. From (2.12), (2.13), and (2.14) we also derive that
(2.16)
$$𝒞(B(p,r))r^{(n2)/2},p𝐒^n,0r\pi /6.$$
### 2.3 An estimate
In this paragraph we provide an integral estimate for any solution $`u`$ of (2.3) outside a compact set $`K`$ in terms of the capacity of $`K`$. It will be frequently used in the sequel. More precisely, the following lemma produces a cut-off function $`\eta `$ which vanishes in a neighbourhoud of $`K`$, equals $`1`$ away from $`K`$, and bounds the rate of a possible blow-up of $`u`$ via estimates (2.18) and (2.19).
###### Lemma 2.1
Let $`KB(0,1)`$ be a compact set in $`𝐑^n`$, $`n3`$, and
$$m\frac{n+2}{2}.$$
Let $`u`$ solve (2.3) in $`K^c`$. Then there exists a function $`\phi C_0^{\mathrm{}}(B(0,2))`$ such that $`0\phi 1`$ in $`B(0,2)`$, $`\phi =1`$ in an open neighbourhood of $`K`$,
(2.17)
$$_{B(0,2)}|D^2\phi |^{(n+2)/4}𝐂(K),$$
and such that for $`\eta =(1\phi )^m`$ the inequalities
(2.18)
$$_{𝐑^n}u(|D\eta |+|\mathrm{\Delta }\eta |)C(m,n)𝐂(K),$$
(2.19)
$$_{𝐑^n}u^{(n+2)/(n2)}\eta C(m,n)𝐂(K)$$
hold.
Proof. 1. The open set $`K^c`$ can be approximated from the interior by domains with smooth boundaries. Consequently, by standard continuity properties of capacity, we can assume in the proof that $`K`$ is a disjoint union of a finite number of closed domains with smooth boundaries. We set $`\widehat{B}=B(0,2)`$.
We claim that there exists a function $`\phi C_0^{\mathrm{}}(\widehat{B})`$ with $`0\phi 1`$ in $`\widehat{B}`$ and $`\phi =1`$ in an open neighbourhood of $`K`$ such that (2.17) holds. To prove this, we first recall a well-known result in nonlinear potential theory Chapter 2, Chapter 9, that states that there exists a function $`\stackrel{~}{\phi }C_0^{\mathrm{}}(\widehat{B})`$ such that
$$\stackrel{~}{\phi }|_K1,_{\widehat{B}}|D^2\stackrel{~}{\phi }|^q^{}𝐂(K),\mathrm{and}\stackrel{~}{\phi }_{L^{\mathrm{}}(\widehat{B})}1.$$
Next, take a function $`HC^{\mathrm{}}(𝐑^1)`$ such that
$$H(t)=0\mathrm{for}t<1/3,H(t)=1\mathrm{for}t>1/2.$$
Now we take $`\phi `$ to be the smooth truncation of $`\stackrel{~}{\phi }`$, $`\phi =H(\stackrel{~}{\phi })`$. Then
$$_{\widehat{B}}|D^2\phi |^q^{}_{\widehat{B}}|H^{\prime \prime }(\stackrel{~}{\phi })|^q^{}|D\stackrel{~}{\phi }|^{2q^{}}+_{\widehat{B}}|H^{}(\stackrel{~}{\phi })|^q^{}|D^2\stackrel{~}{\phi }|^q^{}.$$
To obtain (2.17), we just apply the Gagliardo-Nirenberg interpolation inequality , Chapter 9, to the first term: if $`1<r<\mathrm{}`$, then for any $`fC_0^{\mathrm{}}(\widehat{B})`$
(2.20)
$$Df_{L^{2r}(\widehat{B})}D^2f_{L^r(\widehat{B})}^{1/2}f_{L^{\mathrm{}}(\widehat{B})}^{1/2}.$$
We remark that arguments of this type are well known, cf. Chapter 9, Chapter 3.
2. Let $`u`$ be a solution of (2.3). Take any $`\epsilon >0`$. Appealing to decay (2.4), we choose $`R=R(\epsilon )`$, $`R>4`$, such that
$$u\epsilon \mathrm{on}B(0,R).$$
Set $`B=B(0,R)`$, $`\widehat{B}B`$. Let $`v`$ solve the problem
$$\{\begin{array}{ccc}\hfill \mathrm{\Delta }vv^q=0& \mathrm{in}& BK\hfill \\ \hfill v(x)+\mathrm{}& \mathrm{when}& xK\hfill \\ \hfill v=0& \mathrm{on}& B.\hfill \end{array}$$
Then
$$\mathrm{\Delta }(v+\epsilon )(v+\epsilon )^q0\mathrm{in}BK.$$
Hence by asymptotic condition (2.6) and the comparison principle
(2.21)
$$uv+\epsilon \mathrm{in}BK.$$
In what follows we first prove (2.18) (2.19) for $`v`$ and then let $`\epsilon `$ vanish.
3. Let $`\psi =1\phi `$. We claim that
(2.22)
$$_Bv^q\psi ^mC(m,n)𝐂(K)\mathrm{for}m2q^{}.$$
In fact, by Green’s formula
$`{\displaystyle _B}v^q\psi ^m`$ $`=`$ $`{\displaystyle _B}(\mathrm{\Delta }v)\psi ^m`$
$`=`$ $`{\displaystyle _B}v\mathrm{\Delta }(\psi ^m)+{\displaystyle _B}\left(\psi ^m{\displaystyle \frac{v}{\nu }}v{\displaystyle \frac{\psi ^m}{\nu }}\right),`$
where $`\nu `$ is the outer normal on $`B`$. Since $`\psi |_{\{|x|2\}}=1`$ we conclude that
$$\frac{\psi ^m}{\nu }=0\mathrm{on}B.$$
By the comparison principle, $`v|_{BK}>0`$. Hence
$$_B\psi ^m\frac{v}{\nu }0.$$
Using the Hölder inequality, we compute:
(2.24) $`{\displaystyle _B}v^q\psi ^m`$ $``$ $`{\displaystyle _B}v\mathrm{\Delta }(\psi ^m)`$
$``$ $`{\displaystyle _B}v|\mathrm{\Delta }(\psi ^m)|`$
$``$ $`m{\displaystyle _B}\left(v\psi ^{m1}|\mathrm{\Delta }\psi |\right)+m(m1){\displaystyle _B}\left(v\psi ^{m2}|D\psi |^2\right)`$
$``$ $`m\left({\displaystyle _B}v^q\psi ^m\right)^{1/q}\left({\displaystyle _{\widehat{B}}}\psi ^X|\mathrm{\Delta }\phi |^q^{}\right)^{1/q^{}}`$
$`+m(m1)\left({\displaystyle _B}v^q\psi ^m\right)^{1/q}\left({\displaystyle _{\widehat{B}}}\psi ^Y|D\phi |^{2q^{}}\right)^{1/q^{}},`$
where
$$X=mq^{},Y=m2q^{}.$$
We can assume that the left-hand side in (2.22) is positive. From (2.24) it then follows that
$$_Bv^q\psi ^mm^{2q^{}}_{\widehat{B}}\left(|\mathrm{\Delta }\phi |^q^{}+|D\phi |^{2q^{}}\right).$$
Applying inequality (2.20), we obtain
$$_Bv^q\psi ^mC(m,n)_{\widehat{B}}|D^2\phi |^q^{},$$
and (2.22) follows from (2.17).
4. We claim that
(2.25)
$$_Bv(|\mathrm{\Delta }\eta |+|D\eta |)C(m,n)𝐂(K).$$
In fact, we have by the same calculations as in (2.24):
(2.26) $`{\displaystyle _B}v|\mathrm{\Delta }\eta |`$ $``$ $`m\left({\displaystyle _B}v^q\psi ^{(m1)q}\right)^{1/q}\left({\displaystyle _{\widehat{B}}}|\mathrm{\Delta }\phi |^q^{}\right)^{1/q^{}}`$
$`+m(m1)\left({\displaystyle _B}v^q\psi ^{(m2)q}\right)^{1/q}`$
$`\times \left({\displaystyle _{\widehat{B}}}|D\phi |^{2q^{}}\right)^{1/q^{}},`$
(2.27) $`{\displaystyle _B}v|D\eta |`$ $``$ $`m\left({\displaystyle _B}v^q\psi ^{(m1)q}\right)^{1/q}\left({\displaystyle _{\widehat{B}}}|D\phi |^q^{}\right)^{1/q^{}}.`$
For $`m2q^{}`$ we have
$$(m2)q2q^{},(m1)q2q^{}+q.$$
Thus we can use (2.22) to estimate the integrals containing $`v^q`$ in (2.26) and (2.27). Applying interpolation inequality (2.20) to the last term in (2.26), we conclude on the basis of (2.17) that
$`{\displaystyle _B}v|\mathrm{\Delta }\eta |`$ $``$ $`C(m,n)𝐂(K)^{1/q}\left({\displaystyle _{\widehat{B}}}|D^2\phi |^q^{}\right)^{1/q^{}}`$
$``$ $`C(m,n)𝐂(K).`$
Similarly, applying the Poincaré inequality to the last integral in (2.27) gives
$$_Bv|D\eta |C(m,n)𝐂(K).$$
We conclude that (2.25) indeed holds.
5. From (2.21) and (2.25) we obtain
$`{\displaystyle _{𝐑^n}}u(|D\eta |+|\mathrm{\Delta }\eta |)`$ $`=`$ $`{\displaystyle _{\widehat{B}}}u(|D\eta |+|\mathrm{\Delta }\eta |)`$
$``$ $`C(m,n)\left(𝐂(K)+\epsilon {\displaystyle _{\widehat{B}}}(|D\eta |+|\mathrm{\Delta }\eta |)\right).`$
To establish (2.18) we let $`\epsilon 0`$ both in (2.21) and in the last inequality. A similar limit argument applied to (2.22) gives us (2.19).
Finally, we record a useful elementary inequality, (see for example or ). Let $`J𝐙`$, and let the function $`\zeta :(0,r_J)𝐑^1`$ be either nondecreasing or nonincreasing. Then for any $`\kappa 𝐑`$
(2.28)
$$\underset{j=J+1}{\overset{\mathrm{}}{}}\zeta (r_j)r_j^\kappa _0^{r_J}\zeta (r)r^\kappa \frac{dr}{r}\underset{j=J}{\overset{\mathrm{}}{}}\zeta (r_j)r_j^\kappa .$$
## 3 Capacitary estimates
In this section we prove first estimates on $`u`$ near $`\mathrm{\Omega }`$. We will work in $`𝐑^n`$ instead of $`𝐒^n`$. According to sections 2.1, 2.2 transition to the sphere is immmediate.
Theorems from this section will play the following role in the proof of the main result. Let $`u`$ solve
(3.1)
$$u>0,\mathrm{\Delta }uu^q=0$$
outside a compact set $`K𝐑^n`$. When estimating the length of a curve $`\gamma `$ in the metric $`u^{4/(n2)}g_E`$ we will distinguish two regions. In the first region $`\gamma `$ is far enough from $`K`$. Then pointwise estimate (3.2) from Theorem 3.1 will be applied. In the second region $`\gamma `$ is arbitrarily close to $`K`$. Then we will use integral estimate (3.16) from Theorem 3.2.
###### Theorem 3.1
Let $`KB(0,r)`$ be a compact set in $`𝐑^n`$, $`0<r<1`$, $`n3`$, Let $`u`$ be the maximal solution of (3.1) in $`K^c`$. Then
(3.2)
$$u(x)\frac{𝐂(K)}{|x|^{n2}},|x|2r.$$
Proof. \[of the upper estimate in (3.2)\] 1. According to the scalings (2.7), (2.13) we need to prove that for a compact set $`K`$, $`KB(0,1)`$, the following estimate holds:
(3.3)
$$u(x)𝐂(K)\mathrm{for}\mathrm{all}x\mathrm{such}\mathrm{that}2|x|3.$$
Fix any such $`x`$. Let $`\eta `$ be the function for our set $`K`$ from Lemma 2.1 with some fixed $`m`$.
2. Utilising decay (2.4) we can choose $`R>0`$ so big that we have
(3.4) $`u(x)`$ $`=`$ $`(u\eta )(x)`$
$``$ $`{\displaystyle _{B(0,R)}}G_R(x,y)\mathrm{\Delta }(u\eta )(y)𝑑y.`$
Denote further $`B=B(0,R)`$, $`G=G_R`$. Equation (3.1) gives us
$`\mathrm{\Delta }(u\eta )`$ $`=`$ $`(\mathrm{\Delta }u)\eta +2DuD\eta +u(\mathrm{\Delta }\eta )`$
$``$ $`2DuD\eta +u\mathrm{\Delta }\eta .`$
Substitute this into (3.4) and integrate by parts to deduce that
$`u(x)`$ $``$ $`2{\displaystyle _B}D_yG(x,y)D\eta (y)u(y)𝑑y`$
$`{\displaystyle _B}G(x,y)u(y)\mathrm{\Delta }\eta (y)𝑑y.`$
Next, the choice of $`x`$ and elementary bounds for $`G`$ give that
$$u(x)_{𝐑^n}u(|D\eta |+|\mathrm{\Delta }\eta |).$$
Now estimate (2.18) from Lemma 2.1 leads us to (3.3).
Proof. \[of the lower estimate in (3.2)\] 1. According to the scalings (2.7), (2.13) we need to prove that for a compact set $`K`$, $`KB(0,1)`$, the following estimate holds:
(3.5)
$$u(x)𝐂(K)\mathrm{for}\mathrm{all}x\mathrm{such}\mathrm{that}2|x|3.$$
Taking a suitable approximation we can assume that $`K`$ in (3.5) is the closure of a finite number of domains with smooth boundaries.
Now we recall the fundamental result in potential theory, Ch. 2. The Bessel kernel $`𝒥_2C_{loc}^{\mathrm{}}(𝐑^n\{0\})`$ is defined via the formula
$$(1\mathrm{\Delta })^1f=𝒥_2f\mathrm{for}\mathrm{all}f𝒮.$$
It satisfies the estimates (see, for instance, Chapter 1):
(3.6) $`𝒥_2(x)`$ $``$ $`|x|^{n+2}\mathrm{for}xB(0,1),`$
$`𝒥_2(x)`$ $``$ $`e^{|x|}|x|^{(n+1)/2}\mathrm{for}xB(0,1)^c.`$
The theorem from nonlinear potential theory states that there exists a Radon measure $`\mu ^K`$, $`\mu ^K0`$, such that
$$\mathrm{supp}(\mu ^K)K,$$
and
$$𝐂(K)\mu ^K(K)_{𝐑^n}\left(𝒥_2\mu ^K\right)^q.$$
Hence, after the regularisation of $`\mu ^K`$ and a possible additional smooth approximation of $`K`$ we obtain a function $`gC_0^{\mathrm{}}(𝐑^n)`$, $`g0`$, such that
(3.7)
$$\mathrm{supp}(g)K,$$
and
(3.8)
$$𝐂(K)_{𝐑^n}g_{𝐑^n}\left(𝒥_2g\right)^q.$$
2. Set $`R=10`$ and $`B=B(0,R)`$. For a fixed $`\epsilon >0`$ consider the Dirichlet problem
$$\{\begin{array}{ccccc}\hfill \mathrm{\Delta }v& =& v^q\epsilon g\hfill & \mathrm{in}\hfill & B\hfill \\ \hfill v& =& 0\hfill & \mathrm{on}\hfill & B.\hfill \end{array}$$
As a simple consequence of the comparison principle , it has the unique smooth solution $`v=v_\epsilon `$, $`v>0`$ in $`B`$. Our goal will be to show that there exists $`\epsilon >0`$, $`\epsilon =\epsilon (n)`$, such that
(3.9)
$$v(x)𝐂(K)\mathrm{for}\mathrm{all}x\mathrm{such}\mathrm{that}2|x|3.$$
To prove this we set $`G(x,y)=G_R(x,y)`$, and note that by comparison principle
$$v(x)\epsilon _BG(x,y)g(y)𝑑y\mathrm{for}\mathrm{all}xB.$$
Consequently
(3.10) $`v(x)`$ $`=`$ $`\epsilon {\displaystyle _B}G(x,y)g(y)𝑑y+{\displaystyle _B}G(x,y)v(y)^q𝑑y`$
$``$ $`\epsilon {\displaystyle _B}\left|G(x,y)\right|g(y)𝑑y`$
$`\epsilon ^q{\displaystyle _B}\left|G(x,y)\right|\left({\displaystyle _B}\left|G(y,z)\right|g(z)𝑑z\right)^q𝑑y`$
$`=`$ $`\epsilon I(x)\epsilon ^qII(x)\mathrm{for}\mathrm{all}xB.`$
Hence, to obtain (3.9) we need to estimate $`I`$ from below and $`II`$ from above.
3. Define
$$S=\{x𝐑^n:2|x|3\}.$$
The sets $`S`$, $`\mathrm{supp}(g)`$, and $`B`$ are located at a distance at least $`1`$ from each other. Consequently applying (3.7), (3.8), and invoking the elementary properties of $`G`$, we derive
(3.11) $`I(x)`$ $`=`$ $`{\displaystyle _{\mathrm{supp}(g)}}\left|G(x,y)\right|g(y)𝑑y`$
$``$ $`{\displaystyle _{𝐑^n}}g`$
$``$ $`𝐂(K)\mathrm{for}\mathrm{all}xS.`$
4. We claim that
(3.12)
$$II(x)𝐂(K)\mathrm{for}\mathrm{all}xS.$$
Indeed, fix $`x_0S`$. Introduce the shell
$$\stackrel{~}{S}=\{x𝐑^n:21/100x3+1/100\},$$
and utilise estimate (3.6) for $`𝒥_2`$ to write
(3.13) $`II(x_0)`$ $``$ $`{\displaystyle _{\stackrel{~}{S}}}{\displaystyle \frac{1}{|x_0y|^{n2}}}\left({\displaystyle _B}G(y,z)g(z)𝑑z\right)^q𝑑y`$
$`+{\displaystyle _{B\stackrel{~}{S}}}{\displaystyle \frac{1}{|x_0y|^{n2}}}\left(𝒥_2g\right)^q(y)𝑑y`$
$`=`$ $`X+Y.`$
We estimate $`X`$ and $`Y`$ separately.
To estimate $`X`$ define the function $`H:B𝐑^1`$ by writing
$$H(y)=_BG(y,z)g(z)𝑑z,yB.$$
Notice that according to (3.7), $`H`$ is positive in $`B`$ and harmonic in $`BK`$. Consequently $`H^q`$ is subharmonic in $`BK`$. Hence by the mean value property
$`X`$ $``$ $`\left(\underset{\stackrel{~}{S}}{\mathrm{max}}H\right)^q{\displaystyle _{\stackrel{~}{S}}}{\displaystyle \frac{dy}{|x_0y|^{n2}}}`$
$``$ $`\underset{\stackrel{~}{S}}{\mathrm{max}}H^q`$
$``$ $`{\displaystyle _B}H(y)^q𝑑y.`$
Now (3.6) and (3.8) allow us to conclude that
(3.14) $`X`$ $``$ $`{\displaystyle _{𝐑^n}}\left(𝒥_2g\right)^q(y)𝑑y`$
$``$ $`𝐂(K).`$
To estimate $`Y`$ notice that
$$|x_0y|1/100\mathrm{for}\mathrm{all}yB\stackrel{~}{S}.$$
Therefore utilising (3.8) we derive
(3.15) $`Y`$ $``$ $`{\displaystyle _{B\stackrel{~}{S}}}\left(𝒥_2g\right)^q(y)𝑑y`$
$``$ $`{\displaystyle _{𝐑^n}}\left(𝒥_2g\right)^q`$
$``$ $`𝐂(K).`$
Substituting (3.14) and (3.15) into (3.13) we deduce (3.12).
5. Now we conclude the proof of the theorem. First we establish (3.9). Substitute (3.11) and (3.12) into (3.10):
$$v(x)\left(\epsilon C_1(n)\epsilon ^qC_2(n)\right)𝐂(K)\mathrm{for}\mathrm{all}xS.$$
Choosing the suitable $`\epsilon >0`$ derive (3.9).
Finally, the regularity of $`K`$ implies that our maximal solution $`u`$ blows up near $`K`$ as in (2.6). Therefore
$$uv\mathrm{on}(BK).$$
Owing to (3.7) and the comparison principle,
$$uv\mathrm{in}BK.$$
This inequality and (3.9) complete the proof of (3.5).
Next we establish an integral estimate for any solution of (3.1). This is Theorem 3.2 below. It has a particularly simple proof when $`n4`$ and hence
$$\frac{2}{n2}1.$$
In this case it follows more or less directly from the representation formula for solution of the linear Poisson equation. However, such approach does not work for $`n=3`$ because the singularity of the Green function is too strong then. Proof of Theorem 3.2 given below does not use representation formula. Instead we rely on techniques common in quasilinear elliptic regularity theory. Such arguments were first used by Moser , for linear equations, and by Trudinger for nonlinear equations.
###### Theorem 3.2
Let $`KB(0,1)`$ be a compact set, $`\phi `$ be the function from Lemma 2.1, and let
$$m=\frac{n+2}{2}+100n.$$
Then for any $`u`$ solving (3.1) in $`K^c`$ the estimate
(3.16)
$$_{B(0,10)}u^{2/(n2)}(1\phi )^m𝐂(K)^{2/(n2)}$$
holds.
Proof. 1. We claim that for any number $`\epsilon `$, $`0<\epsilon <1`$, the inequality
(3.17) $`\left({\displaystyle _{𝐑^n}}u^{(1\epsilon )n/(n2)}|\zeta |^{2n/(n2)}\right)^{(n2)/n}`$ $``$ $`C(\epsilon )({\displaystyle _{𝐑^n}}u^{1\epsilon }|D\zeta |^2`$
$`+{\displaystyle _{𝐑^n}}u^{q\epsilon }\zeta ^2)`$
holds for all functions $`\zeta C_0^{\mathrm{}}(𝐑^n)`$, such that $`\zeta =0`$ in an open neighbourhood of $`K`$. In fact, multiplying the equation
$$\mathrm{\Delta }uu^q=0\mathrm{in}𝐑^nK$$
by $`u^\epsilon \zeta ^2`$, integrating by parts, and invoking the formula
$$D(u^\epsilon \zeta ^2)=\epsilon u^{\epsilon 1}\zeta ^2Du+2\zeta u^\epsilon D\zeta ,$$
we deduce that
$`\epsilon {\displaystyle _{𝐑^n}}|Du|^2u^{\epsilon 1}\zeta ^2`$ $``$ $`2{\displaystyle _{𝐑^n}}|Du||D\zeta |u^\epsilon |\zeta |`$
$`+{\displaystyle _{𝐑^n}}u^{q\epsilon }\zeta ^2.`$
Since
$$|Du||D\zeta |u^\epsilon |\zeta |\delta |Du|^2u^{\epsilon 1}\zeta ^2+\frac{1}{4\delta }|D\zeta |^2u^{\epsilon +1}$$
for each $`\delta >0`$, we derive that
$$_{𝐑^n}|Du|^2u^{\epsilon 1}\zeta ^2C(\epsilon )\left(_{𝐑^n}u^{1\epsilon }|D\zeta |^2+_{𝐑^n}u^{q\epsilon }\zeta ^2\right).$$
After some calculations we find
$$_{𝐑^n}\left|D\left(u^{(1\epsilon )/2}\zeta \right)\right|^2C(\epsilon )\left(_{𝐑^n}u^{1\epsilon }|D\zeta |^2+_{𝐑^n}u^{q\epsilon }\zeta ^2\right).$$
Now (3.17) follows from the Sobolev inequality applied to the left hand side.
2. Set $`B=B(0,10)`$ and $`\widehat{B}=B(0,20)`$. In (3.17) choose $`\epsilon (0,1)`$ such that
$$(1\epsilon )\frac{n}{n2}=\frac{2}{n2}.$$
Then select a smooth cutoff function $`\theta C_0^{\mathrm{}}(\widehat{B})`$, such that $`\theta =1`$ on $`B`$. Take the function $`\eta =(1\phi )^m`$ from Lemma 2.1. Now set $`\zeta =\eta \theta `$, in estimate (3.17) to discover that
$`\left({\displaystyle _B}u^{2/(n2)}(1\phi )^m\right)^{(n2)/2}`$ $`(`$ $`{\displaystyle _{𝐑^n}}u^{1\epsilon }|D\eta |^2`$
$`+{\displaystyle _{𝐑^n}}u^{1\epsilon }|D\theta |^2`$
$`+{\displaystyle _{𝐑^n}}u^{q\epsilon }(\eta \theta )^2)^{1/(1\epsilon )}.`$
We estimate three integrals in the right hand side of (3) as follows. Applying Holder inequality and estimates (2.17), (2.19) from Lemma 2.1 we deduce that
$`{\displaystyle _{𝐑^n}}u^{1\epsilon }|D\eta |^2`$ $``$ $`\left({\displaystyle _B}u^{q(1\epsilon )}(1\phi )^{(n+2)/2}\right)^{1/q}\left({\displaystyle _B}|D\phi |^{2q^{}}\right)^{1/q^{}}`$
$``$ $`𝐂(K)^{(q\epsilon )/q}.`$
Estimate (3.2) from Theorem 3.1 implies
$`{\displaystyle _{𝐑^n}}u^{1\epsilon }|D\theta |^2`$ $``$ $`u_{L^{\mathrm{}}(\widehat{B}B)}^{1\epsilon }`$
$``$ $`𝐂(K)^{1\epsilon }.`$
Finally, owing to Holder inequality and (2.19) we have
$$_{𝐑^n}u^{q\epsilon }(\eta \theta )^2𝐂(K)^{(q\epsilon )/q}.$$
Substituting these estimates in (3) we arrive at
$`\left({\displaystyle _B}u^{2/(n2)}(1\phi )^m\right)^{(n2)/2}`$ $``$ $`𝐂(K)+𝐂(K)^{(q\epsilon )/(q\epsilon q)}`$
$``$ $`𝐂(K).`$
This is assertion (3.16).
## 4 Proof of the Wiener test for conformal metrics
### 4.1 Sufficiency
We prove the implication $`(ii)(i)`$ in Theorem 1.1.
1. In the proof we will work on $`𝐒^n`$. A curve $`\gamma :[0,+\mathrm{})\mathrm{\Omega }`$ is said to converge to infinity if for every compact set $`M\mathrm{\Omega }`$, there is a time $`T`$, $`0<T<+\mathrm{}`$, such that $`\gamma (t)M`$ for all $`t>T`$. By a version of the Hopf-Rinow theorem, a metric is complete in $`\mathrm{\Omega }`$ if and only if every smooth curve converging to infinity has the infinite length.
Let $`u_\mathrm{\Omega }`$ be the maximal solution of the conformal scalar curvature equation (2.1) in $`\mathrm{\Omega }`$. We set
$$g=u_\mathrm{\Omega }^{4/(n2)}\stackrel{\text{ }_{}}{g}$$
Fix any smooth curve $`\gamma :[0,+\mathrm{})\mathrm{\Omega }`$ converging to infinity. To prove statement (i) in Theorem 1.1 we need to show that
(4.1)
$$L_g(\gamma )=_0^{\mathrm{}}u_\mathrm{\Omega }(\gamma )^{2/(n2)}\stackrel{\text{ }_{}}{g}(\dot{\gamma },\dot{\gamma })^{1/2}𝑑t=+\mathrm{}.$$
In the rest of the proof we establish (4.1).
2. Compactness of $`𝐒^n`$ and convergence of $`\gamma `$ to infinity imply the existence of a point $`pK`$ such that
(4.2)
$$d_{\stackrel{\text{ }_{}}{g}}(\gamma (T_k),p)0\mathrm{for}\mathrm{a}\mathrm{sequence}\{T_k\},T_k+\mathrm{}.$$
For $`j=1`$, $`2`$, $`\mathrm{}`$ we define $`\mathrm{\Gamma }_j`$ to be that part of $`\gamma `$ whose image is contained in the shell $`S_j`$,
$$S_j=\{x𝐒^n:r_j<d_{\stackrel{\text{ }_{}}{g}}(x,p)<r_{j1}\}.$$
The smoothness of $`\gamma `$ implies that for any $`jj_0`$ the set $`\mathrm{\Gamma }_j`$ is at most a countable union of open smooth curves. Utilising condition (4.2) we deduce that $`\mathrm{\Gamma }_j\mathrm{}`$, and moreover
$$L_{\stackrel{\text{ }_{}}{g}}(\mathrm{\Gamma }_j)\frac{r_j}{100}\mathrm{for}\mathrm{all}jj_0.$$
After a rotation we can assume that
$$K\overline{B}(p,r_j)𝐔\mathrm{for}\mathrm{all}jj_010,$$
where $`𝐔`$ is cup (2.8) around the south pole from the definition of $`𝒞`$. We claim that for all $`jj_0`$ the inequality
(4.3)
$$L_g(\mathrm{\Gamma }_j)r_j\left(\frac{𝒞\left(K\overline{B}(p,r_{j+2})\right)}{r_j^{n2}}\right)^{2/(n2)}$$
holds. In fact, define the open set $`\mathrm{\Omega }_j`$, $`\mathrm{\Omega }_j\mathrm{\Omega }`$, by writing
$$\mathrm{\Omega }_j=𝐒^n(K\overline{B}(p,r_{j+2})).$$
Let $`u_j`$ be the maximal solution to our equation (2.1) in $`\mathrm{\Omega }_j`$. Pull estimate (3.2) from Theorem 3.1 back to the sphere via the stereographic projection, keeping in mind that the conformal factor in (2.2) satisfies
$$\mathrm{{\rm Y}}(x)1\mathrm{for}\mathrm{all}x,|x|10.$$
We discover that
$`u_\mathrm{\Omega }(x)`$ $``$ $`u_j(x)`$
$``$ $`{\displaystyle \frac{𝒞\left(K\overline{B}(p,r_{j+2})\right)}{r_j^{n2}}}\mathrm{for}\mathrm{all}xS_j\mathrm{\Omega }.`$
Let $`I_j`$, $`I_j(0,+\mathrm{})`$, be the open set such that
$$\mathrm{\Gamma }_j:I_j\mathrm{\Omega }S_j.$$
Then we derive that
$`L_g(\mathrm{\Gamma }_j)`$ $`=`$ $`{\displaystyle _{I_j}}u_\mathrm{\Omega }(\gamma )^{2/(n2)}\stackrel{\text{ }_{}}{g}(\dot{\gamma },\dot{\gamma })^{1/2}𝑑t`$
$``$ $`\left(\underset{S_j}{inf}u_\mathrm{\Omega }\right)^{2/(n2)}L_{\stackrel{\text{ }_{}}{g}}(\mathrm{\Gamma }_j)`$
$``$ $`\left({\displaystyle \frac{𝒞\left(K\overline{B}(p,r_{j+2})\right)}{r_j^{n2}}}\right)^{2/(n2)}r_j,`$
thereby obtaining (4.3).
3. We claim that (4.1) holds. Indeed, the sets $`S_j`$ are disjoint, and thus
$$L_g(\gamma )\underset{j1}{}L_g(\mathrm{\Gamma }_j).$$
To each term with sufficiently large number in this sum we apply estimate (4.3) and recall (2.16) to derive that
$`L_g(\gamma )`$ $``$ $`{\displaystyle \underset{jj_0}{}}r_{j+2}\left({\displaystyle \frac{𝒞\left(K\overline{B}(p,r_{j+2})\right)}{r_{j+2}^{n2}}}\right)^{2/(n2)}`$
$``$ $`{\displaystyle \underset{jj_0+100}{}}\left({\displaystyle \frac{𝒞(B(p,r_j)K)}{𝒞(B(p,r_j))}}\right)^{2/(n2)}.`$
Finally utilise (2.28) and (1.4) to establish (4.1). This completes the proof of implication $`(ii)(i)`$ in Theorem 1.1.
### 4.2 Necessity
Now we prove the implication $`(i)(ii)`$ in Theorem 1.1.
1. Seeking a contradiction assume that $`(ii)`$ does not hold. Hence
(4.4)
$$_0^{1/2}\left(\frac{𝒞(B(P,r)K)}{𝒞(B(P,r))}\right)^{2/(n2)}\frac{dr}{r}<+\mathrm{}$$
for some $`P\mathrm{\Omega }`$. The desired contradiction will follow if the maximal solution of (2.1) does not give the metric complete in $`\mathrm{\Omega }`$. Let $`U`$ be this maximal solution, and let
$$g=U^{2/(n2)}\stackrel{\text{ }_{}}{g}.$$
According to the Hopf-Rinow theorem, to prove the non-completeness of $`g`$ we must show that there exists a smooth curve $`c`$,
(4.5) $`c:[0,1)\mathrm{\Omega },`$ $`\mathrm{such}\mathrm{that}d_{\stackrel{\text{ }_{}}{g}}(c(t),P)0\mathrm{as}t1,`$
$`\mathrm{and}L_g(c)<+\mathrm{}.`$
2. First we reformulate claim (4.5). Fix a parameter $`\rho >0`$, which we will later choose small. Set
$$\stackrel{~}{K}=K\overline{B}(P,\rho ).$$
Let $`U_1`$ be the maximal solution of (2.1) in $`\stackrel{~}{\mathrm{\Omega }}`$,
$$\stackrel{~}{\mathrm{\Omega }}=𝐒^n\stackrel{~}{K},$$
and let $`U_2`$ be the maximal solution of (2.1) in $`\mathrm{\Omega }B(P,\rho )`$. From (2.5) we deduce that
$$U^{2/(n2)}U_1^{2/(n2)}+U_2^{2/(n2)}\mathrm{in}\mathrm{\Omega }.$$
At the same time (2.4) implies
$$U_2(x)C(\rho )\mathrm{for}\mathrm{all}xB(P,\rho /2).$$
Therefore (4.5) is equivalent to the same statement with $`\mathrm{\Omega }`$ replaced by $`\stackrel{~}{\mathrm{\Omega }}`$, and $`g`$ replaced by
$$\stackrel{~}{g}=U_1^{2/(n2)}\stackrel{\text{ }_{}}{g}.$$
To prove this statement it will be convinient to transform the problem to $`𝐑^n`$.
Applying a suitable rotation and stereographic projection we can achieve that $`P`$ is mapped to $`0`$. We denote the image of $`K`$ under such map by the same letter $`K`$. In $`𝐑^n`$ we set
$$K_j=K\overline{B}_j\mathrm{and}B=B(0,1).$$
By $`u`$ we denote the conformal pullback (2.2) of $`U_1`$,
$$u(x)=\mathrm{{\rm Y}}(x)U_1(\sigma ^1x),x\sigma (\stackrel{~}{\mathrm{\Omega }}).$$
As it is shown in section 2.1, $`u`$ is the maximal solution of (2.3) in $`𝐑^nK_J`$ for some $`J`$, $`J=J(\rho )`$. From (4.4) and (2.12) we deduce that
(4.6)
$$_0^1\left(\frac{𝐂(B(0,r)K_J)}{𝐂(B(0,r))}\right)^{2/(n2)}\frac{dr}{r}<+\mathrm{}.$$
Finally, to establish (4.5) we must prove that $`u^{2/(n2)}g_E`$ is not complete, that is
(4.7) $`{\displaystyle _\gamma }u^{2/(n2)}𝑑s<+\mathrm{}`$ $`\mathrm{for}\mathrm{a}\mathrm{smooth}\mathrm{curve}\gamma :[0,1)BK_J,`$
$`\mathrm{such}\mathrm{that}\gamma (t)0\mathrm{as}t1.`$
3. We intend to establish (4.7). The construction of $`\gamma `$ in (4.7) will be indirect. More precisely, let us first reduce the proof of (4.7) to an integral estimate for our maximal solution $`u`$.
We assert that it is possible to choose large enough $`J`$ in (4.6) (equivalently, to choose small enough $`\rho >0`$) such that there exists a compact set $`\mathrm{\Sigma }`$,
$$K_J\mathrm{\Sigma }B,$$
with the following two properties:
(4.8)
$$_{B\mathrm{\Sigma }}u(x)^{2/(n2)}\frac{1}{|x|^{n1}}𝑑x<+\mathrm{},$$
and
(4.9)
$$^{n1}\left(\pi (\mathrm{\Sigma }\{0\})\right)<^{n1}(B),$$
where $`\pi `$ is the radial projection on $`B`$,
$$\pi :B\{0\}B,x\frac{x}{|x|}.$$
This assertion is the core of the proof. Before passing to its verification we conclude the current step by showing that (4.8), (4.9) immediately imply (4.7) and hence the theorem.
Indeed, for $`\omega B`$ we define the interval $`\mathrm{}(\omega )`$ by writing
$$\mathrm{}(\omega )=\{x𝐑^n:x=s\omega ,0<s1\}.$$
Set
$$\mathrm{\Xi }=B\pi (\mathrm{\Sigma }\{0\}).$$
First notice that
$$\pi ^1(\mathrm{\Xi })B\mathrm{\Sigma }.$$
Hence, using the polar coordinates $`(r,\omega )`$, $`r>0`$, $`\omega B`$, we deduce at once from (4.8) that
$`+\mathrm{}`$ $`>`$ $`{\displaystyle _{\pi ^1(\mathrm{\Xi })}}u(x)^{2/(n2)}{\displaystyle \frac{1}{|x|^{n1}}}𝑑x`$
$`=`$ $`{\displaystyle _\mathrm{\Xi }}{\displaystyle _0^1}u(x(r,\omega ))^{2/(n2)}{\displaystyle \frac{1}{r^{n1}}}r^{n1}𝑑r𝑑^{n1}(\omega )`$
$`=`$ $`{\displaystyle _\mathrm{\Xi }}\left({\displaystyle _{\mathrm{}(\omega )}}u^{2/(n2)}𝑑s\right)𝑑^{n1}(\omega ).`$
Next, apply (4.9) to discover that
$$^{n1}(\mathrm{\Xi })=^{n1}(B)^{n1}\left(\pi \left(\mathrm{\Sigma }\{0\}\right)\right)>0.$$
Consequently
$$_{\mathrm{}(\omega _0)}u^{2/(n2)}𝑑s<+\mathrm{}\mathrm{for}\mathrm{some}\omega _0\mathrm{\Xi }.$$
By our definitions
$$\mathrm{}(\omega _0)K_J=\mathrm{},$$
and we conclude that (4.7) holds for the curve $`\gamma =\mathrm{}(\omega _0)`$.
Thus, to establish the theorem it is left to construct $`\mathrm{\Sigma }`$ satisfying (4.8) and (4.9). The rest of the proof is devoted entirely to this construction.
4. For any $`jJ`$ let
$$v=v_j$$
be the maximal solution for $`K_{j2}`$. We now apply Theorem 3.2 and scalings (2.7), (2.13) to bound $`v`$. Applying the scaling
$$xr_{j3}x,x𝐑^n,$$
to estimate (3.16) we deduce that there exists a function
$$\phi =\phi _j,$$
such that:
$`\phi C_0^{\mathrm{}}(B(0,2)),0\phi 1\mathrm{in}B(0,2),`$
$`\phi =1\mathrm{in}\mathrm{a}\mathrm{neighbourhood}\mathrm{of}K_{j2},`$
(4.10) $`{\displaystyle _{B(0,2)}}\left|D^2\phi \right|^{(n+2)/4}𝐂(K_{j2}),`$
and
(4.11) $`{\displaystyle \frac{1}{r_j^{n1}}}{\displaystyle _{B_{j2}}}v^{2/(n2)}(1\phi )^m`$ $``$ $`\left(𝐂(K_{j2}/r_{j3})\right)^{2/(n2)}`$
$``$ $`\left({\displaystyle \frac{𝐂(K_{j2})}{𝐂(B_{j2})}}\right)^{2/(n2)}.`$
5. Now we construct the compactum $`\mathrm{\Sigma }`$ for (4.8), (4.9). Set
$$S_j=\{x:r_j|x|r_{j1}\}.$$
First for a fixed $`jJ`$, define the compact set $`E_j`$ by writing
$$E_j=\{xS_j:\phi _j(x)\frac{99}{100}\},$$
where the function $`\phi _j`$ is taken from (4.10), (4.11). Then define
$$\mathrm{\Sigma }=\left(\underset{j=J}{\overset{\mathrm{}}{}}E_j\right)\{0\}.$$
According to the construction, the set $`\mathrm{\Sigma }`$ is compact and
$$K_J\mathrm{\Sigma }B.$$
We claim that (4.9) holds. Indeed, the definition of the capacity and (4.10) imply that
$`𝐂(E_j)`$ $``$ $`{\displaystyle _{B(0,2)}}\left|D^2\left({\displaystyle \frac{100}{99}}\phi _j\right)\right|^{(n+2)/4}`$
$``$ $`𝐂(K_{j2}).`$
The metric estimate (2.14) therefore ensures
$$_{\mathrm{}}^{n1}(E_j)^{(n2)/2}𝐂(K_{j2})^{n1}.$$
The projection $`\pi `$ restricted to $`S_j`$ distorts the distances at most $`1/r_j`$ times. Consequently
$`_{\mathrm{}}^{n1}(\pi (\mathrm{\Sigma }\{0\}))`$ $``$ $`{\displaystyle \underset{j=J}{\overset{\mathrm{}}{}}}_{\mathrm{}}^{n1}(\pi (E_j))`$
$``$ $`{\displaystyle \underset{j=J}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{r_j^{n1}}}_{\mathrm{}}^{n1}(E_j)`$
$``$ $`{\displaystyle \underset{j=J}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{r_j^{n1}}}𝐂(K_{j2})^{(n1)2/(n2)}`$
$``$ $`\left({\displaystyle \underset{j=J2}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{𝐂(K_j)}{𝐂(B_j)}}\right)^{2/(n2)}\right)^{n1}.`$
According to (2.28) and (4.6) we can make the last series as small as we wish by choosing $`J`$ large enough. Thus for any $`\epsilon >0`$ we may fix $`J`$ in (4.6) so that
$$_{\mathrm{}}^{n1}\left(\pi (\mathrm{\Sigma }\{0\})\right)<\epsilon .$$
For sets lying on an $`s`$-dimensional smooth submanifold, the Hausdorff $`s`$-measure is equivalent to the Lebesgue $`s`$-measure. Consequently
$$_{\mathrm{}}^{n1}(F)^{n1}(F)\mathrm{for}\mathrm{any}FB.$$
This gives (4.9).
6. It is left to prove (4.8). Splitting the integral there we find that
(4.12)
$$_{B\mathrm{\Sigma }}u(x)^{2/(n2)}\frac{1}{|x|^{n1}}𝑑x\underset{j=J}{\overset{\mathrm{}}{}}\frac{1}{r_j^{n1}}_{S_j\mathrm{\Sigma }}u^{2/(n2)}.$$
Thus our task is to estimate $`u`$ in $`S_j\mathrm{\Sigma }`$. Fix $`jJ`$. Let as before $`v`$ be the maximal solution for $`K_{j2}`$. For $`l=1`$, $`2`$, $`\mathrm{}`$, $`j2`$ let $`w_l`$ be the maximal solution for $`KS_l`$. From (2.5) we deduce that
$$uv+\underset{l=1}{\overset{j2}{}}w_l\mathrm{in}S_j.$$
Next observe that
$$|xy|r_l\mathrm{for}\mathrm{all}xS_j,yS_l,lj2.$$
Hence applying the scaled estimate (3.2) from Theorem 3.1 to $`w_l`$, we derive that
(4.13) $`u(x)`$ $``$ $`v(x)+{\displaystyle \underset{l=1}{\overset{j2}{}}}{\displaystyle \frac{𝐂(KS_l)}{r_l^{n2}}}`$
$``$ $`v(x)+{\displaystyle \underset{l=1}{\overset{j2}{}}}{\displaystyle \frac{𝐂(K_l)}{r_l^{n2}}}\mathrm{for}\mathrm{all}xS_j.`$
To estimate $`v`$ in $`S_j`$ notice that (4.11) and the definition of $`\mathrm{\Sigma }`$ ensure that
$$\frac{1}{r_j^{n1}}_{S_j\mathrm{\Sigma }}v^{2/(n2)}\left(\frac{𝐂(K_{j2})}{𝐂(B_{j2})}\right)^{2/(n2)}.$$
Utilising (4.13) we thereupon conclude that
$`{\displaystyle \frac{1}{r_j^{n1}}}{\displaystyle _{S_j\mathrm{\Sigma }}}u^{2/(n2)}`$ $``$ $`{\displaystyle \frac{1}{r_j^{n1}}}{\displaystyle _{S_j\mathrm{\Sigma }}}v^{2/(n2)}`$
$`+{\displaystyle \frac{1}{r_j^{n1}}}\left({\displaystyle \underset{l=1}{\overset{j2}{}}}{\displaystyle \frac{𝐂(K_l)}{r_l^{n2}}}\right)^{2/(n2)}|S_j|`$
$``$ $`\left({\displaystyle \frac{𝐂(K_{j2})}{𝐂(B_{j2})}}\right)^{2/(n2)}`$
$`+r_j\left({\displaystyle \underset{l=1}{\overset{j2}{}}}{\displaystyle \frac{𝐂(K_l)}{r_l^{n2}}}\right)^{2/(n2)}.`$
Now continue (4.12) to find
(4.14) $`{\displaystyle _{B\mathrm{\Sigma }}}u(x)^{2/(n2)}{\displaystyle \frac{1}{|x|^{n1}}}𝑑x`$ $``$ $`{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{𝐂(K_j)}{𝐂(B_j)}}\right)^{2/(n2)}`$
$`+{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}r_j\left({\displaystyle \underset{l=1}{\overset{j}{}}}{\displaystyle \frac{𝐂(K_l)}{r_l^{n2}}}\right)^{2/(n2)}`$
$`=`$ $`I+II.`$
Thus to prove (4.8) we need to bound $`I`$ and $`II`$.
7. Utilising (2.28) we deduce at once that
$$I_0^1\left(\frac{𝐂(KB(0,r))}{𝐂(B(0,r))}\right)^{2/(n2)}\frac{dr}{r}.$$
To estimate $`II`$ we define the function $`\mathrm{\Phi }:(0,1)𝐑^1`$ by writing
$$\mathrm{\Phi }(r)=𝐂(KB(0,r)),0<r<1.$$
First assume that $`n4`$ and hence
$$\frac{2}{n2}1.$$
In this case by the simple change of the summation order we discover that
$`II`$ $``$ $`{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}r_j{\displaystyle \underset{l=1}{\overset{j}{}}}\left({\displaystyle \frac{𝐂(K_l)}{r_l^{n2}}}\right)^{2/(n2)}`$
$``$ $`{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{𝐂(K_l)}{r_l^{n2}}}\right)^{2/(n2)}r_l`$
$``$ $`{\displaystyle _0^1}\left({\displaystyle \frac{\mathrm{\Phi }(r)}{r^{n2}}}\right)^{2/(n2)}𝑑r.`$
Assume next that $`n=3`$, and hence
$$\frac{2}{n2}=2$$
Then Hardy’s inequality implies that
$`II`$ $``$ $`{\displaystyle _0^1}\left({\displaystyle _t^1}{\displaystyle \frac{\mathrm{\Phi }(r)}{r}}{\displaystyle \frac{dr}{r}}\right)^2𝑑t`$
$``$ $`{\displaystyle _0^1}\left({\displaystyle \frac{\mathrm{\Phi }(r)}{r}}\right)^2𝑑r.`$
Thus for any $`n3`$ we have
$$II_0^1\left(\frac{\mathrm{\Phi }(r)}{r^{(n2)/2}}\right)^{2/(n2)}\frac{dr}{r}.$$
Returning to (4.14) and recalling (2.16) we derive
$$_{B\mathrm{\Sigma }}u(x)^{2/(n2)}\frac{1}{|x|^{n1}}𝑑x_0^1\left(\frac{𝐂(KB(0,r))}{𝐂(B(0,r))}\right)^{2/(n2)}\frac{dr}{r}.$$
Employing (4.6) we establish (4.8). This completes the proof of the implication $`(i)(ii)`$ in Theorem 1.1. |
warning/0506/quant-ph0506055.html | ar5iv | text | # Universality of quantum critical dynamics in a planar OPO
## Abstract
We analyze the critical quantum fluctuations in a coherently driven planar optical parametric oscillator. We show that the presence of transverse modes combined with quantum fluctuations changes the behavior of the ‘quantum image’ critical point. This zero-temperature non-equilibrium quantum system has the same universality class as a finite-temperature magnetic Lifshitz transition.
The non-equilibrium system consisting of a nonlinear crystal inside a laser driven Fabry-Perot interferometer, that couples sub-harmonic intra-cavity modes to a harmonic pump Peter , is known as an optical parametric oscillator (OPO). As well as demonstrating quantum squeezing Wu and EPR entanglement in numerous quantum information experiments, this system is widely used in frequency conversion applications. In single-mode experiments, there is a critical point in the phase diagram. This is caused by an increase in the pump intensity, which results in a transition from a disordered (but quantum squeezed) phase below threshold, to an ordered phase with a coherent output above threshold.
When extended to an interferometer with multiple transverse modes, more complex dynamical effects occur due to diffraction of the down-converted light, which are governed by the Swift-Hohenberg equation near thresholdClassical . The theory can be quantizedDrummond , and hence includes quantum fluctuationsLugiato . Experimentally observableFabre ; Pointer ‘quantum images’ are evidence for quantum pattern formation with spatio-temporal correlations in the output quadratures. Thus, the system can show both spatial critical fluctuations and non-equilibrium spontaneous pattern formation, which occurs in many fields of physics and other sciencesGollub .
In this Letter we apply the theory of finite size scaling to solve for the critical quantum dynamical properties, and obtain the universality class this phase-transition corresponds to. We provide a full quantum description of this non-equilibrium system using the positive P-representation, focusing on the nature of the critical point and critical fluctuations. This allows us to obtain an analytic solution for the functional distribution of the large critical fluctuations caused by quantum noise in the down-conversion process. There is an unexpected universality property in the solutions. Even though this is a non-equilibrium quantum system of coupled boson fields, we find the in one quadrature, the quantum fluctuations have exactly the same behavior as a classical thermal system of fields at a two-dimensional Lifshitz point, which is a model commonly used to describe the phase transition to a modulated magnetic phase; in the complementary quadrature there is strong entanglement.
The unitary evolution of the OPO system can be described by the HamiltonianLugiato
$`\widehat{}`$ $`=`$ $`{\displaystyle \underset{n=0,1}{}}{\displaystyle d^2\stackrel{}{x}\left\{\widehat{A}_n^{}\left[\omega _n\frac{v^2}{2\omega _n}^2\right]\widehat{A}_n\right\}}`$ (1)
$`+i\mathrm{}{\displaystyle d^2\stackrel{}{x}\widehat{A}_0\left\{^{}e^{2i\omega _Lt}\chi \widehat{A}_1^2\right\}}.`$
The term $`\chi `$ is a coupling parameter that depends on the nonlinear crystal, the frequencies of the field modes are $`\omega _1`$, $`\omega _0=2\omega _1`$, $`v`$ is the intracavity group velocity, and $`\widehat{A}_n`$ is the $`nth`$ photon field. The pump is described by the amplitude $``$ that could carry a spatial structure - but here we will assume a constant plane wave input. In addition, there are damping effects due to output couplings from the cavity mirrors, which can be well approximated using as a Markovian master equation for the density matrix $`\widehat{\rho }`$, so that:
$$\frac{\widehat{\rho }}{t}=\frac{1}{i\mathrm{}}[\widehat{},\widehat{\rho }]+\underset{n=0,1}{}\gamma _n_n\left[\widehat{\rho }\right],$$
(2)
where $`_n\left[\widehat{\rho }\right]=d^2\stackrel{}{x}\left[2\widehat{A}_n\widehat{\rho }\widehat{A}_n^{}\widehat{\rho }\widehat{A}_n^{}\widehat{A}_n\widehat{A}_n^{}\widehat{A}_n\widehat{\rho }\right]`$ describes the output coupling from the $`n`$th intra-cavity mode, with damping rate $`\gamma _n`$. This leads to a set of Fokker-Planck equations, mapped from the the quantum density matrix, using operator representation theory. These are valid provided boundary terms vanish in the mapping transformation, which we have checked numerically. Using the positive P-representation, we derive the following stochastic equationsPeter ; Lugiato ; Fabre in a rotating frame at frequency $`\omega _L`$:
$`{\displaystyle \frac{A_1}{t}}`$ $`=`$ $`\stackrel{~}{\gamma }_1A_1+\chi A_1^+A_0+i\gamma _1D^2A_1+\sqrt{\chi A_0}\xi _1(t,\stackrel{}{x}).`$
$`{\displaystyle \frac{A_0}{t}}`$ $`=`$ $`\stackrel{~}{\gamma }_0A_0+{\displaystyle \frac{\chi ^{}}{2}}A_1^2+{\displaystyle \frac{i\gamma _1}{2}}D^2A_0.`$ (3)
Here we write the two dimensional Laplacian causing diffraction, as $`^2=^2/x^2+^2/y^2`$. The complex relaxation rates are $`\stackrel{~}{\gamma }_i=\gamma _i(1+i\mathrm{\Delta }_i)`$. The relative detunings between the pump laser at $`2\omega _L`$, and the modes supported by the cavity are $`\mathrm{\Delta }_0=(\omega _02\omega _L)/\gamma _0`$, and $`\mathrm{\Delta }_1=(\omega _1\omega _L)/\gamma _1`$ . The diffraction rate is defined as $`D=v^2/(2\gamma _1\omega _1)`$. The stochastic field $`\xi _1`$ which describes quantum noise is real and Gaussian, with correlations of $`\xi _1(t)=0`$ and $`\xi _1(\stackrel{}{x},t)\xi _1(\stackrel{}{x}^{},t^{})=\delta ^2(\stackrel{}{x}\stackrel{}{x}^{})\delta (tt^{})`$.
In addition, there are equations that correspond to the hermitian conjugate fields. As elsewhere in this Letter, we obtain these by conjugating the constant terms and replacing the stochastic and noise fields according to: $`A_iA_i^+`$, $`\xi _1\xi _1^+`$, where $`\xi _1`$, $`\xi _1^+`$ are independent real Gaussian noises. These two noise fields are sufficient to generate all quantum effects, and are physically caused by the discrete nature of the photon pairs produced in down-conversion. The c-number fields $`A_i(t,\stackrel{}{x}),A_i^+(t,\stackrel{}{x})`$ are therefore not complex conjugate, although they are stochastically equivalent in terms of normally-ordered operator moments to photon operator fields $`\widehat{A}_i(t,\stackrel{}{x}),\widehat{A}_i^{}(t,\stackrel{}{x})`$. Thus for example, the photon number density is$`\widehat{A}_i^{}(t,\stackrel{}{x})\widehat{A}_i(t,\stackrel{}{x}^{})=A_i^+(t,\stackrel{}{x})A_i(t,\stackrel{}{x^{}}).`$
If we remove the transverse modes from the above equations we return to the well known single mode OPO theory. This system has a quantum critical point - a phase transition in the infinite volume limit, where the quantum fluctuations are reduced below the vacuum level for the squeezed quadrature and become huge for the unsqueezed quadraturePlimak . In a recent analysis CDD of this problem near the critical point, going beyond the linear theory, we obtained a scaling law for the squeezing quadrature spectrum near threshold, and the parameters for the optimum squeezing.
The introduction of transverse modes generates a spatial structure in the sub-harmonic field with an intensity correlation function known as a “quantum image” Lugiato , since it is supported by quantum fluctuations. To treat this problem analytically, we can perform an adiabatic elimination of the stable pump mode in the limit of $`\gamma _0\gamma _1`$ and $`\mathrm{\Delta }_00`$. That is, we assume that the pump mode has a short relaxation time.
Neglecting pump diffraction - which is negligible in the critical regime - we obtain an adiabatic solution for the pumped field: $`\overline{A}_0=\left(\chi ^{}A_1^2/2\right)/\gamma _0`$, together with a similar equation for the conjugate term. This solution takes into account the depletion of the pumping mode that supplies energy for down-converted light, and leads to an adiabatic equation for the down-converted field:
$`{\displaystyle \frac{A_1}{t}}`$ $`=`$ $`\stackrel{~}{\gamma }_1A_1+{\displaystyle \frac{\chi }{\gamma _0}}\left({\displaystyle \frac{\chi ^{}}{2}}A_1^2\right)A_1^+`$ (4)
$`+`$ $`i\gamma _1D^2A_1+\sqrt{\chi \overline{A}_0}\xi _1(t,\stackrel{}{x}).`$
Next, we introduce dimensionless variables $`\tau =t/t_0`$ and $`\stackrel{}{r}=\stackrel{}{x}/x_0`$, with a corresponding down-converted field $`\alpha =x_0A_1`$. Due to critical slowing down, the characteristic length $`x_0`$ and time $`t_0`$ scale as $`t_0=1/(g\gamma _1)`$ and $`x_0^2=D/\sqrt{g_c}`$, where the effective nonlinear coefficient is $`g_c=|\chi |^{4/3}/[8D\gamma _0\gamma _1]^{2/3}`$ , and we will assume that $`g_c1`$. The dimensionless driving field is $`\stackrel{~}{\mu }=\chi /[\gamma _1\gamma _0]=\mu +i\theta `$. We also introduce appropriately scaled noise fields with $`\xi (\tau ,\stackrel{}{r})=x_0\sqrt{t_0}\xi _1(t,\stackrel{}{x})`$, and the corresponding hermitian conjugate terms.
With these definitions, we find that:
$`{\displaystyle \frac{\alpha }{\tau }}`$ $`=`$ $`{\displaystyle \frac{1}{g_c}}\left[(1+i\mathrm{\Delta }_1)\alpha +\left(\stackrel{~}{\mu }4g_c^2\alpha ^2\right)\alpha ^+\right]`$ (5)
$`+`$ $`{\displaystyle \frac{1}{\sqrt{g_c}}}\left[i_r^2\alpha +\xi \sqrt{\left(\stackrel{~}{\mu }4g_c^2\alpha ^2\right)}\right].`$
This equation includes both the down-conversion term proportional to $`\alpha _1^+`$, which generates a squeezed signal — together with a nonlinear saturation term proportional to $`\alpha _1^2\alpha _1^+`$, which limits the down-converted amplitude, and leads to finite size critical fluctuations.
Critical fluctuations are most usefully analyzed with scaled quadratures that correspond to experimentally accessible homodyne detection. These are defined as
$`X(\tau ,\stackrel{}{r})`$ $`=`$ $`\sqrt{g_c}\left[\alpha _1(\tau ,\stackrel{}{r})+\alpha _1^+(\tau ,\stackrel{}{r})\right]`$
$`Y(\tau ,\stackrel{}{r})`$ $`=`$ $`i\left[\alpha _1^+(\tau ,\stackrel{}{r})\alpha _1(\tau ,\stackrel{}{r})\right].`$ (6)
Similarly, there are quadrature noise fields defined as $`\xi _x(\tau ,\stackrel{}{r})=\xi (\tau ,\stackrel{}{r})+\xi ^+(\tau ,\stackrel{}{r})`$, and $`\xi _y(\tau ,\stackrel{}{r})=i\sqrt{g_c}\left[\xi ^+(\tau ,\stackrel{}{r})\xi (\tau ,\stackrel{}{r})\right]`$. The resulting quantum dynamical equations for these signal field quadratures are:
$`{\displaystyle \frac{X}{\tau }}`$ $`=`$ $`\left[\gamma _x+X^2+g_cY^2\right]X\left[\gamma _{xy}+^2\right]Y+\xi _x`$
$`g_c{\displaystyle \frac{Y}{\tau }}`$ $`=`$ $`\left[\gamma _y+g_cX^2+g_{c}^{}{}_{}{}^{2}Y^2\right]Y\left[\gamma _{yx}^2\right]X+\xi _y.`$
The linear decay matrix that couples the $`X`$ and $`Y`$ quadratures is given by:
$$\left[\begin{array}{cc}\gamma _x& \gamma _{xy}\\ \gamma _{yx}& \gamma _y\end{array}\right]=\left[\begin{array}{cc}(1\mu )/g_c& (\theta +\mathrm{\Delta }_1)/\sqrt{g_c}\\ (\mathrm{\Delta }_1\theta )/\sqrt{g_c}& (1+\mu )\end{array}\right].$$
(8)
We assume also that close to threshold, and for small enough detunings, $`\gamma _x=O(1)`$ and $`\gamma _{xy}=(\theta +\mathrm{\Delta }_1)/\sqrt{g_c}=O(1)`$. We can always choose quadrature phases so that $`\theta =\mathrm{\Delta }_1+O(g_c)`$. With this choice, $`Y`$ is mainly coupled to the $`X`$ quadrature via the diffraction term, which couples noise from the critical fluctuations back into the squeezed quadrature. This implies that $`\stackrel{~}{\mu }=1+O(g_c)`$ and $`\gamma _y=2+O(\sqrt{g_c})`$, so that the noise correlations are given by;
$$\xi _x(\tau ,\stackrel{}{r})\xi _x(\tau ^{},\stackrel{}{r}^{})=2\delta (\tau \tau ^{})\delta ^2(\stackrel{}{r}\stackrel{}{r}^{})+O(g_c).$$
(9)
We can now perform a second type of adiabatic elimination, which is valid in a neighbourhood of the critical point. This takes into account the fact that the fluctuations in the $`X`$ quadrature become very slow near threshold, while the $`Y`$ quadrature still responds on fast time-scales of order $`1/\gamma _1`$. To leading order we can drop terms of $`O(\sqrt{g_c})`$ where $`g_c1`$, and approximate the above equations as follows:
$`{\displaystyle \frac{X}{\tau }}`$ $`=`$ $`\gamma _xX\gamma _{xy}YX^3^2Y+\xi _x`$
$`0`$ $`=`$ $`2Y+^2X.`$ (10)
We can therefore eliminate the fast or non-critical quadrature variable $`Y`$, by writing the steady state solution of the $`Y`$ quadrature as $`Y^2X/2`$ . This produces a reduced equation for the critical quadrature variable $`X`$, which is valid near threshold:
$$\frac{X}{\tau }=\gamma _xXX^3\frac{\gamma _{xy}}{2}^2X\frac{1}{2}^4X+\xi _x.$$
(11)
The above Langevin equation is a Ginzburg-Landau equation describing the critical quadrature dynamics. Unlike the usual application of this equation, we note that system is a non-equilibrium one. The noise term $`\xi _x`$ is of quantum origin rather than thermal origin, and is present at zero temperature. It is possible to write an equivalent functional Fokker Planck equation for the probability density $`P[X]`$,
$$\frac{P}{\tau }=\frac{\delta }{\delta X}\left[\left(\gamma _x+X^2+\frac{\gamma _{xy}}{2}^2+\frac{1}{2}^4\right)X+\frac{\delta }{\delta X}\right]P,$$
(12)
and look for the equilibrium distribution in the form $`P[X]=Nexp(V[X])`$, where $`V(X)`$ is a potential functional. Making this substitution, the solution for the distribution $`P[X]`$ is given by :
$$Pe^{\left[{\scriptscriptstyle d^2\stackrel{}{r}\left(2\gamma _xX^2+X^4\gamma _{xy}[X]^2+[^2X]^2\right)/4}\right]}.$$
(13)
This expression is exactly the same as the Ginzburg-Landau free energy of a next nearest neighbor magnetic interaction, where $`X`$ plays the role of an order parameter. That is, we have been able to map this problem into a soluble magnetic phase-transition equation with a Lifshitz pointHornreich . The phase diagram of this optical system should therefore have two ordered phases, one of them a spatially modulated phase associated with a pattern formation. This generic behavior is known to occur in an OPO, from previous analysisClassical .
In this analogy, the “optical paramagnetic phase” corresponds to a random photon emission from the OPO operating below threshold, the “optical ferromagnetic phase” to a continuum emission uniformly distributed in the transverse plane parallel to the cavity operating above threshold. In the “optical ferromagnetic modulated phase”, we have a continuum emission but with modulated quadrature in this plane. At the Lifshitz point, all three phases co-exist.
The line $`\gamma _x=0`$ is the line of the second-order phase transition between order-disorder (coherent-incoherent) states. In the incoherent phase below threshold, $`\gamma _x>0`$, and in the uniform coherent phase above threshold $`\gamma _x<0`$, as expected in the single mode case. If $`\gamma _x`$ vanishes we have a Lifshitz point over the line $`\gamma _x=0`$, and thus a triple point characterizing the coexistence of the three phases. This holds in the case of perfect tuning of the signal field inside the cavity, so that $`\mathrm{\Delta }_1=0`$.
In condensed matter physics, the nature of the Lifshitz pointKaplan ; Hornreich79 is crucially dependent on the order parameter and spatial dimension. Depending on the dimensionality of the order parameter, the system may or may not have a true phase transition in two dimensions. Our system has a one-dimensional real order parameter and two spatial dimensions with transverse modes, so we expect that this system should have a true phase transition in the infinite volume limit at finite temperature. According to the Mermin-Wagner theoremMerminWagner , increasing the order parameter dimension (as in type II down-conversion) would result in phase fluctuations that completely destroy any long range order. Similarly, reducing the spatial dimension would result in a continuous transition without a threshold.
A numerical simulation of the Ginzburg-Landau equations (11) with a continuously scanned input shows that the sub-harmonic quadrature correlations appear to have a true critical point in two transverse dimensions. This result is shown in Figure (1), which graphs $`|X(k)|^2`$ around $`k=0`$, as a function of the driving field $`\gamma _x`$ near threshold.
While these large fluctuations are occurring, we note that there are still strong non-classical correlations in the squeezed quadrature. This can be seen by analysing the relevant equations to the next order in $`g_c`$, which we also simplify by using a Gaussian factorizationStochdiagram :
$$g_c\frac{Y}{\tau }\gamma _yY+^2Xg_cYX^2\stackrel{~}{\gamma }_{yx}X+\xi _y,$$
(14)
where $`\stackrel{~}{\gamma }_{yx}=\gamma _{yx}+2g_cXY`$. In general, one can always choose an optimum local oscillator phase so that $`\theta =k^2+\mathrm{\Delta }_1+2g_cXY`$, in order to minimise the feedback of critical fluctuations into the squeezed quadrature at a given transverse momentum $`k`$. This leads to Fourier solutions which showing that entanglementDrummondFicek between the modes of momentum $`k`$ and $`k`$ can still occur at small enough wave-vectors, resulting in a universal squeezing spectrum as a function of frequency:
$$V(\mathrm{\Omega })=1\frac{1g_c(X^2+\gamma _x)}{(g_c\mathrm{\Omega }/2)^2+\left(1+g_c(X^2\gamma _x)/2\right)^2}.$$
(15)
This result differs from the linearised predictions of earlier treatmentsLugiato . A graph of the resulting spectrum in the Gaussian approximation is shown in Fig (2), compared to the linearized squeezing spectrum, showing large differences near threshold.
In summary, we have shown that the planar non-equilibrium OPO with quantum noise can be mapped to a magnetic phase-transition in two dimensions. Since the present case has a scalar, real order parameter it is analogous to the uni-axial ($`m=1`$) magnetic order parameter case, which is known to have a thermal equilibrium Lifshitz-point phase-transition at finite temperatureHornreich79 . This demonstrates a striking resemblance between known thermal equilibrium phase-transitions, and a quantum non-equilibrium system in which quantum noise replaces thermal noise. In this system there are also quantum correlations of the emitted photons, causing quantum squeezing and entanglement. Nevertheless, this highly non-classical behavior is found only in the squeezed ($`Y)`$ quadrature which has no critical slowing down \- and co-exists with a rather classical and universal critical fluctuation field in the conjugate ($`X`$) quadrature.
PDD acknowledges support from the Australian Research Council. |
warning/0506/math-ph0506046.html | ar5iv | text | # 1 Introduction
## 1 Introduction
In certain physical problems there may exist extra hidden symmetries which are not apparent, unless searched for. Some of the examples from classical considerations are, a particle in a $`\frac{1}{r^2}`$ potential in one dimension , the conserved Runge-Lenz vector of the Kepler problem and the extra symmetries of a charge moving in the field of a magnetic monopole, and the generators beyond the Poincare invariance that give rise to conformal invariance in electrodynamics as well as in Yang-Mills theory. The existence of symmetries help in classifying and obtaining energy levels and eigenstates in quantum mechanical problems, generating new solutions and also formulating conservation laws. As we know manifestation of scale invariance in deep inelastic scattering had deep significance in the development gauge theories. The invariance under scale and conformal transformations also motivated the construction of a simple classical model which leads to conformal quantum mechanics. Recently there have been a revival of interest in this model. This is due to the observation in string theory dynamics that a particle near a black hole possesses SO(2,1) symmetries as in conformal quantum mechanics..
The symmetries of charged particle - monopole system and the conformal quantum mechanics were obtained from physical reasonings and scale invariance. However, there exists a general programme to obtain the symmetries of the equations of motion of any such system by using the group theoretic methods of Lie. In this paper we use this method to find the Lie point symmetries of the monopole system as well as some other physically motivated systems.
A knowledge of the symmetry group of a system of differential equations leads to several types of applications. For example, the symmetry group is helpful in finding solutions to the set of equations as well as constructing new solutions to the systems from the known ones. There have been efforts to generalise and extend the Lie approach by considering non-standard symmetries, as well. This helps in getting a wider class of solutions. Non-classical symmetries result by the weakening of invariance requirements of the differential equation under the symmetry transformations. Another approach is to enlarge the space of independent variables by adding auxiliary variables and finding the symmetries. Then the symmetries related to the original system are figured out and the corresponding symmetries are called non-local. This way, it also provides a method to classify different classes of solutions corresponding to different symmetries. One can use, on the other hand, symmetry groups to classify families of differential equations depending on arbitrary parameters or functions. As an extension of this idea, Krause has introduced the important idea of complete symmetry of a differential equation so as to expand the symmetry group such that the manifold of solutions is an homogeneous space of the group and the group is specific to the system, i.e., no other system admits that symmetry group. The complete symmetry group of the system is the group represented by the set of symmetries required to specify completely a system and its point symmetries.
Further, just the enumeration of the symmetry generators sometimes provide much physical insight and quantitative physical results for which the full solutions are not required. The derivation of Kepler’s third law of planetary motion and Runge-Lenz vectors, calculation of energy levels for hydrogen like atoms and generalized Kepler’s problems, harmonic oscillators, Morse potentials, electron in a specific nonuniform magnetic field, being some such examples. In these the energy eigenvalues are obtained through the method of spectrum generating algebras which gives the Casimir invariants directly without explicit recourse to the solutions. Of course, the solutions are also obtained from the representation theory. For finding the continuous symmetries, Lie’s method of group analysis seems to be the most powerful technique available. For example, Witten had considered an example of the equation of motion of a particle in three dimensions constrained to move on the surface of a sphere in the presence of a magnetic monopole. This is the classical analogue of the Wess-Zumino model. The equations of motion of such a system in the presence of a magnetic monopole cannot be obtained from the usual Lagrangian formulation unless one goes to a higher dimension. Hence, the usual method of finding the symmetries through Noether’s theorem would have difficulties. So, to look for the continuous symmetries associated with such classical systems one has to analyze directly the equations of motion. Similar is the case for Korteweg-de Vries equation which is not amenable to a direct Lagrangian formulation when expressed as a lowest order equation. Another example is the Lorenz system of equations which have been dealt in the papers by Sen and Tabor, and Nucci. For the classical systems, this procedure of finding directly the symmetries from equation of motion is, in some sense, more fundamental. This is because in certain cases many different Lagrangians may give rise to the same equations of motion. The group analysis of the equations of motion gives all the Lie point group symmetry generators. In the cases where a Lagrangian formulation is possible, the usual Noether symmetries are a subset of the above generators. This subset of generators acting on the Lagrangian gives zero. However, besides these there may be other generators obtained through group analysis which have direct physical significance, but not explicitly available from the consideration of usual Noether symmetries alone. The reproduction of Kepler’s third law in the planetary motion problem is such an example. The extension of this idea to the notion of Lie dynamical symmetries contains similarly a subclass known as Cartan symmetries. The Runge-Lenz vector can be obtained from such considerations. These symmetries are, further, related to the Lie-Bäcklund symmetries.
The application of these types of analysis to nonlocal cases have been widely studied through Bäcklund transformations and related techniques in the context of integrable systems containing infinite number of conservation laws. The Thirring model has been analyzed by Morris . The differential geometric forms developed earlier are used in the above analysis to obtain the prolongation structure. Some other applications of these ideas to important problems from physics is comprehensibly covered by Gaeta.
In this paper, as a first step towards a complete group analysis, we find the Lie point symmetries of the coupled set of differential equations representing the motion of a charged particle in three dimensions:
(i) in the presence of a magnetic field proportional to the coordinate vector,
(ii) in the presence of a magnetic monopole stationed at the centre of a sphere, with the constraint that the particle moves on the surface of this sphere of unit radius,
(iii) in the presence of a $`\frac{1}{r^2}`$ potential
(iv) in the presence of a magnetic monopole,
(v) in the presence of a dyon,
(vi) in a model magnetic field of the form $`𝐁=\{0,0,\frac{}{x^2}\}`$. Here $``$ is a constant,
and
(vii) a particle in a type of velocity dependent potential,
We also obtain the generators of the complete symmetry group of Krause for most of the above examples in this paper.
It should be noted that in the context of the symmetries of Wess-Zumino-Witten models, the symmetries in the higher dimension play by far the most important role and these have been fruitfully exploited .
Section 2 provides the outline of the method for group analysis of the equations of motion. In section 3 we use this method to find the generators of the point symmetries for the above examples. The symmetries of the equation of a scalar quantum particle near the horizon of a massive blackhole is considered in section 4. Next we follow the method of reduction of order introduced by Nucci to obtain the complete Krause symmetries for four of the above examples which is the content of section 5. Finally, Section 6 is devoted to physical interpretations and conclusions.
## 2 Symmetry Conditions
Typically we are interested in the coupled nonlinear set of equations representing the equations of motion of a particle in three dimensions. These are of the form
$`\ddot{x}_a=\beta \omega _a(x_i,\dot{x}_i,t)`$ (1)
where a dot represents derivative with respect to time, $`a,i=1,2,`$ and $`3`$, and $`\beta `$ is a constant involving mass, coupling constant etc.. The expressions for the function $`\omega _a`$ will be given explicitly for each example.
These set of equations can be analyzed by means of one parameter groups by infinitesimal transformations. We demand the equation to be invariant under infinitesimal changes of the explicit variable $`t`$, as well as simultaneous infinitesimal changes of the dependent functions $`x_a`$ in the following way,
$`t\stackrel{~}{t}`$ $`=`$ $`t+ϵ\tau (t,x_1,x_2,x_3)+O(ϵ^2),`$
$`x_a\stackrel{~}{x}_a`$ $`=`$ $`x_a+ϵ\eta _a(t,x_i)+O(ϵ^2).`$ (2)
Under $`t\stackrel{~}{t}`$ and $`x_a\stackrel{~}{x}`$, the equation changes to,
$`\ddot{\stackrel{~}{x}}_a`$ $`=`$ $`\beta \stackrel{~}{\omega }_a(\stackrel{~}{t},\stackrel{~}{x}_i,\dot{\stackrel{~}{x}}_i)`$ (3)
To illustrate the procedure consider the simple case in one space dimension. We express the above equation in terms of $`t`$ and $`q`$ by using the transformation (2). Then the invariance condition implies that an expression containing various partial derivatives of $`\tau `$ and $`\eta `$ is obtained which equates to zero. For example, we get
$`{\displaystyle \frac{d\stackrel{~}{x}}{d\stackrel{~}{t}}}`$ $`=`$ $`{\displaystyle \frac{dx+ϵ(\frac{\eta }{t}dt+\frac{\eta }{x}dx)}{dt+ϵ(\frac{\tau }{t}dt+\frac{\tau }{x}dx)}}+O(ϵ^2)`$ (4)
and now relate the left hand side with $`\frac{dx}{dt}`$ by using binomial theorem for the denominator to obtain
$`{\displaystyle \frac{d\stackrel{~}{x}}{d\stackrel{~}{t}}}`$ $`=`$ $`{\displaystyle \frac{dx}{dt}}+ϵ\left[\left({\displaystyle \frac{\eta }{x}}{\displaystyle \frac{\tau }{t}}\right){\displaystyle \frac{dx}{dt}}{\displaystyle \frac{\tau }{x}}\left({\displaystyle \frac{x}{t}}\right)^2\right]+O(ϵ^2).`$ (5)
A similar procedure is followed to express $`\frac{d^2\stackrel{~}{x}}{d\stackrel{~}{t}^2}`$ likewise. By substituting equations (2 ) - (5 ) for a given explicit expression for $`\omega _a`$ and remembering that $`\frac{d^2x_a}{dt^2}\omega _a`$ is zero, we obtain the desired partial differential equation whose solution would determine $`\tau (t,x)`$ and $`\eta (t,x)`$. In our case, of course, we have to find $`\tau (t,x_1,x_2,x_3)`$ and $`\eta _a(t,x_1,x_2,x_3)`$’s.
To relate these to the generators of the infinitesimal transformations we write
$`\stackrel{~}{t}(t,x_i;ϵ)=t+ϵ\tau (t,x_i)+\mathrm{}=t+ϵ𝐗t+\mathrm{}`$ (6)
$`\stackrel{~}{x}_a(t,x_i;ϵ)=x_a+ϵ\eta (t,x_i)+\mathrm{}=x_a+ϵ𝐗x_a+\mathrm{}`$ (7)
where the functions $`\tau `$ and $`\eta _a`$ are components of tangent vectors at the points $`\stackrel{~}{t}`$ and $`\stackrel{~}{x}_a`$ defined by
$`\tau (t,x_i)`$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{t}}{ϵ}}_{ϵ=0},`$ (8)
$`\eta _a(t,x_i)`$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{x_a}}{ϵ}}_{ϵ=0}`$ (9)
and the operator $`𝐗`$ is given by
$`𝐗`$ $`=`$ $`\tau (t,x_i){\displaystyle \frac{}{t}}+\eta _a(t,x_i){\displaystyle \frac{}{x_a}}.`$ (10)
where repeated indices are summed. Following Stephani , we will find the infinitesimal generators of the symmetry under which the system of differential equations do not change. The symmetry is generated by $`𝐗`$ and its extension
$`\dot{𝐗}`$ $`=`$ $`\tau {\displaystyle \frac{}{t}}+\eta _a{\displaystyle \frac{}{x_a}}+\dot{\eta }_a{\displaystyle \frac{}{\dot{x}_a}}`$ (11)
and the symmetry condition under transformations represented by equation (2) determines $`\dot{\eta }_a`$. In the expanded form the symmetry condition becomes
$`\eta _b\omega _{a}^{}{}_{,b}{}^{}+(\eta _{b}^{}{}_{,t}{}^{}+\dot{x}_c\eta _{b}^{}{}_{,c}{}^{}\dot{x}_b\tau _{,t}\dot{x}_b\dot{x}_c\tau _{,c}){\displaystyle \frac{\omega _a}{\dot{x}_b}}+\tau \omega _{a}^{}{}_{,t}{}^{}+2\omega _a(\tau _{,t}+\dot{x}_b\tau _{,b})`$
$`+\omega _b(\dot{x}_a\tau _{,b}\eta _{a}^{}{}_{,b}{}^{})+\dot{x}_a\dot{x}_b\dot{x}_c\tau _{,bc}+\dot{x}_a\tau _{,tt}+2\dot{x}_a\dot{x}_c\tau _{,tc}`$
$`\dot{x}_c\dot{x}_b\eta _{a}^{}{}_{,bc}{}^{}2\dot{x}_b\eta _{a}^{}{}_{,tb}{}^{}\eta _{a}^{}{}_{,tt}{}^{}=0`$ (12)
where $`f_{,t}=\frac{f}{t}`$ and $`f_{,c}=\frac{f}{x_c}`$. By herding together coefficients of the terms that are cubic, quartic, and linear in $`\dot{x_a}`$ , and the ones independent of $`\dot{x_a}`$ separately, and equating each of these to zero we obtain an over determined set of partial differential equations and solve for $`\tau `$ and $`\eta _a`$.
## 3 Symmetry generators, classical particle
The solutions of the symmetry conditions provide us the generators of the group. In this section, we explicitly obtain the generators for the cases mentioned in the introduction.
For a charged particle moving in the absence of electromagnetic field the equation of motion is given by,
$`m\ddot{x}_k=0`$ (13)
and it is well known that the symmetry condition, which is equation(12), gives rise to the eight parameter symmetry generator of the general projective transformation as given by
$`𝐗=[a_1+a_2t+a_3x_a+a_4tx_a+a_5t^2]{\displaystyle \frac{}{t}}`$
$`+[a_6+a_7t+a_8x_a+a_5tx_a+a_4(x_a)^2]{\displaystyle \frac{}{x_a}}`$ (14)
for each of the equations(13). We mention this result, so that the generators can be compared with the results we would obtain later for our examples.
We first give the complete analysis of a simpler case, which is the case(i) as mentiones in the introduction.
Case (i):
For the motion of a charged particle in the presence of a magnetic field proportional to the coordinate vectors, the equations of motion are
$`\ddot{x}_k`$ $`=`$ $`\beta \epsilon _{kbc}\dot{x}_bx_c=\omega _k`$ (15)
where the above expression in the middle is the Lorentz force acting on a charged particle. The magnetic field for this case is proportional to $`x_c`$.
Substituting
$`\omega _k`$ $`=`$ $`\beta \epsilon _{kbc}\dot{x}_bx_c`$ (16)
into equation (12) we obtain, in general, coupled partial differential equations for $`\tau `$ and $`\eta `$ by equating to zero the terms corresponding to various powers of $`\dot{x}_l`$.
Consideration of the term with $`\dot{x}_a\dot{x}_b\dot{x}_c`$ in Eq.(12) tells us
$`\tau _{,bc}`$ $`=`$ $`0.`$ (17)
Hence we may have
$`\tau `$ $`=`$ $`A_l(t)x_l+B(t)+C.`$ (18)
The terms quadratic in $`\dot{x}`$ give
$`\beta \dot{x}_b\dot{x}_cx_m\tau _{,c}\epsilon _{abm}+2\beta \dot{x}_l\dot{x}_bq_m\tau _{,b}\epsilon _{alm}`$
$`+\beta \dot{x}_r\dot{x}_ax_s\tau _{,b}\epsilon _{brs}+2\dot{x}_a\dot{x}_c\tau _{,tc}\dot{x}_c\dot{x}_b\eta _{a}^{}{}_{,bc}{}^{}`$ $`=`$ $`0`$ (19)
This shows that $`\tau `$ has to be independent of $`x_l`$. Hence
$`\tau `$ $`=`$ $`B(t)+C,`$ (20)
and $`\eta `$ may have the form
$`\eta _a=D(t)x_a+E_l(t)\epsilon _{lam}x_m+F(t)+G.`$ (21)
The terms linear in $`\dot{x}_l`$ provide
$`\beta \dot{x}_l\epsilon _{alb}\eta _b+\beta \dot{x}_cx_m\epsilon _{abm}\eta _{b}^{}{}_{,c}{}^{}\beta \dot{x}_bx_m\epsilon _{abm}\tau _{,t}+2\beta \dot{x}_lx_m\epsilon _{alm}\tau _{,t}`$
$`\beta \dot{x}_rx_s\epsilon _{brs}\eta _{a}^{}{}_{,b}{}^{}2\dot{x}_b\eta _{a}^{}{}_{,tb}{}^{}+\dot{x}_a\tau _{,tt}=0.`$ (22)
This demands
$`B(t)`$ $`=`$ $`tH+C,`$ (23)
and also $`\eta _a`$ has to be independent of $`t`$, giving
$`\eta _a`$ $`=`$ $`Hx_a+E_l\epsilon _{lam}x_m.`$ (24)
Thus we obtain five generators
$`𝐗_a=\epsilon _{akb}x_b{\displaystyle \frac{}{x_k}},𝐗_4={\displaystyle \frac{}{t}},𝐗_5=t{\displaystyle \frac{}{t}}x_a{\displaystyle \frac{}{x_a}}`$ (25)
and their Lie algebra
$`[𝐗_a,𝐗_b]=\epsilon _{abc}𝐗_c,[𝐗_a,𝐗_4]=0,[𝐗_a,𝐗_5]=0,[𝐗_4,𝐗_5]=𝐗_4.`$ (26)
A comparison with the results of similar analysis for the Kepler problem shows that the first four generators are identical, the first three corresponding to the generators of the three dimensional rotation group and $`𝐗_4`$ is the generator for time translation. However in this case the law corresponding to Kepler’s third law goes instead like
$`\stackrel{~}{t}\stackrel{~}{r}`$ $`=`$ $`tr=constant.`$ (27)
The Lie algebra represented by equation (26) corresponds, in the notation of Stephani , to the group $`SO(3)\times G_2IIa`$, where the $`G_2IIa`$ is a group with the two generators
$`𝐗_4={\displaystyle \frac{}{t}},𝐗_5=t{\displaystyle \frac{}{t}}x_a{\displaystyle \frac{}{x_a}}.`$ (28)
We can find their extension from the formula
$`\dot{\eta }_a={\displaystyle \frac{d\eta _a}{dt}}\dot{x}_a{\displaystyle \frac{d\tau }{dt}}`$ (29)
and obtain, denoting the extensions by $`\dot{𝐗}`$,
$`\dot{𝐗}_a`$ $`=`$ $`\epsilon _{akb}(x_b{\displaystyle \frac{}{x_k}}+\dot{x}_b{\displaystyle \frac{}{\dot{x}_k}}),`$ (30)
$`\dot{𝐗}_4`$ $`=`$ $`{\displaystyle \frac{}{t}},`$ (31)
$`\dot{𝐗}_5`$ $`=`$ $`t{\displaystyle \frac{}{t}}x_a{\displaystyle \frac{}{x_a}}2\dot{x}_a{\displaystyle \frac{}{\dot{x}_a}},scaling.`$ (32)
Henceforth we scale $`\beta `$, which is a function of the coupling constant, mass, etc. to one.
Through an analysis that is similar to the above procedure, we find the the generators of the symmetry groups for the following examples.
Case (ii):
If the particle is further constrained to move on the surface of a sphere of unit radius, the equation of motion becomes
$`\ddot{x}_a`$ $`=`$ $`\epsilon _{abc}\dot{x}_bx_cx_a\dot{x}_k\dot{x}_k.`$ (33)
This is equivalent to the case of a particle moving in the presence of a magnetic monopole centered at the origin of the sphere. Witten has generalized this idea to arbitrary dimensions for field theoretic considerations. In the above and henceforth we have scaled $`\beta `$ to one. As has been pointed out by Witten , one faces trouble in attempting to derive these equations of motion by using the usual procedure of variation of a Lagrangian since no obvious term can be included in the Lagrangian whose variation would give the equation of motion (15). Hence it would be more appropriate here to consider the group analysis of the equations of motion directly to obtain all the Lie point symmetries. Usually the Noether symmetries are a subclass of these. However, the present analysis cannot give any of the non-Lie symmetries.
With $`\omega _a`$ being equal to the right hand side of equation (33), the group analysis shows that there is only one trivial time translation besides the generators of the angular momentum for this problem,
$`𝐗_a=\epsilon _{akb}x_b{\displaystyle \frac{}{x_k}},𝐗_\mathrm{𝟒}`$ $`=`$ $`{\displaystyle \frac{}{t}}.`$ (34)
Same is the case if we ignore the term containing $`\epsilon _{abc}`$in equation (33). So also for a magnetic field $`𝐁=bx_1𝐞_1+(B_0+bx_3)𝐞_3`$, which is an idealized version of the Stern-Gerlach magnetic field .
Case (iii):
If a potential like $`\frac{1}{r^2}`$ is only present we find the generators with extensions to be
$`\dot{𝐗}_a=\epsilon _{abc}\left(x_c{\displaystyle \frac{}{x_b}}+\dot{x}_c{\displaystyle \frac{}{\dot{x}_b}}\right),spacerotations,`$
$`\dot{𝐗_4}={\displaystyle \frac{}{t}},timetranslation,`$
$`\dot{𝐗}_5=2t{\displaystyle \frac{}{t}}+x_a{\displaystyle \frac{}{x_a}}{\displaystyle \frac{1}{2}}\dot{x}_a{\displaystyle \frac{}{\dot{x}_a}},`$
$`Keplerlikescalinglaw{\displaystyle \frac{t}{r^2}}=constant,`$
$`\dot{𝐗}_6=t^2{\displaystyle \frac{}{t}}+tx_a{\displaystyle \frac{}{x_a}}+x_a{\displaystyle \frac{}{\dot{x}_a}}\dot{x}_a{\displaystyle \frac{}{\dot{x}_a}}`$ (35)
The vector fields have the commutation relations
$`[𝐗_a,𝐗_b]=\epsilon _{abc}𝐗_b,`$
$`[𝐗_a,𝐗_4]=[𝐗_a,𝐗_5]=[𝐗_a,𝐗_6]=0`$
$`[𝐗_4,𝐗_5]=2𝐗_4,[𝐗_4,𝐗_6]=𝐗_5,[𝐗_5,𝐗_6]=2𝐗_6`$ (36)
The classical Kepler problem with $`\frac{1}{r}`$ potential has a different scaling law of $`\frac{t^2}{r^3}`$ and also does not possess the symmetry corresponding to generator $`𝐗_\mathrm{𝟔}`$. However, it possesses a Runge-Lenz vector. Stephani has given a general method to obtain such conserved vectors in the Lagrangian formulation. In the quantum mechanical case, if the eigenvalues are taken instead the Hamiltonian operator, an enhanced symmetry with closed Lie algebra occurs for $`\frac{1}{r}`$ potential. For $`\frac{1}{r^2}`$ potential we could not find a classical Runge-Lenz vector by Stephani’s method. This appears to be related to orbits being not closed in such a potential. However, as has already been noted, in this case new vector fields result leading to the extra symmetries.
Case (iv):
Jackiw has considered the symmetries of equation of motion, Lagrangian, and Hamiltonian for a charged particle in the field of a magnetic monopole . He had discovered an extra $`SO(2,1)`$ hidden symmetry by scaling and physical considerations. Leonhardt and Piwnicki have explored the theoretical possibility of obtaining the field of quantized monopoles when a classical dielectric moves in a charged capacitor . Since the magnetic field due to a magnetic monopole is $`B_a=\frac{x_a}{r^3}`$, the equation of motion is
$`\ddot{x}_a=\epsilon _{abc}{\displaystyle \frac{\dot{x}_bx_c}{r^3}}=\omega _a`$ (37)
We have taken the coupling constant, mass etc. to be unity. Lie point symmetries of these equations were obtained in. We get the same generators with extension as in the case (iii) for the $`\frac{1}{r^2}`$ potential. Zwanzinger had considered the motion of a charged particle in the presence of monopole along with a $`\frac{1}{r^2}`$ potential. Hence the equation of motion
$`\ddot{x}_a=\epsilon _{abc}{\displaystyle \frac{\dot{x}_bx_c}{r^3}}+{\displaystyle \frac{\mu ^2x_a}{mr^4}}`$ (38)
also possesses the same above symmetry. The last terms of the equation reminds of an electric dipole potential at large distances in the one dimensional case.
Case(v):
We have obtained for the case of the field due to both a monopole and a charge, i.e. a dyon, only the first four generators of equation (35) for case(iii).
Case (vi):
But for a velocity dependent potential with equations of motions of the form
$`\ddot{x}_a`$ $`=`$ $`\dot{x}_ax_kx_k`$ (39)
we again find five symmetry generators, the first four being the same as $`\dot{𝐗}_a`$ and $`\dot{𝐗}_4`$ while
$`\widehat{𝐗}_5`$ $`=`$ $`2t{\displaystyle \frac{}{t}}x_a{\displaystyle \frac{}{x_a}},`$ (40)
and with its extension,
$`\dot{\widehat{𝐗}}_5`$ $`=`$ $`2t{\displaystyle \frac{}{t}}x_a{\displaystyle \frac{}{x_a}}3\dot{x}_a{\displaystyle \frac{}{\dot{x}_a}}.`$ (41)
The length and time scale in this case as
$`\stackrel{~}{t}\stackrel{~}{r}^2`$ $`=`$ $`tr^2.`$ (42)
Case (vii):
For a charged particle moving in a model magnetic field of the form
$`𝐁=\{B_x=0,B_y=0,B_z={\displaystyle \frac{}{x^2}}\}`$ (43)
the equations of motion are
$`\ddot{x}_1={\displaystyle \frac{\dot{x}_2}{x_{1}^{}{}_{}{}^{2}}},\ddot{x}_2={\displaystyle \frac{\dot{x}_1}{x_{1}^{}{}_{}{}^{2}}},\ddot{x}_3=0.`$ (44)
Here $``$ is a constant. This magnetic field may be obtained from a current density $`𝐉=(0,\frac{}{x^3},0)`$, which is singular. It is interesting to note, however, that the Schrödinger equation can be exactly solved and corresponding energy levels be obtained in the manner of Landau . The condition (12) for $`a=1`$, and $`2`$ gives
$`\tau =\lambda ,\eta _1=0,\eta _2=\sigma ,`$ (45)
and for $`a=3`$ we obtain
$`\eta _3=\rho +x_3`$ (46)
where $`\lambda `$, $`\sigma `$, and $`\rho `$ are constants. This gives rise to the vector fields
$`𝐗_\tau =\lambda {\displaystyle \frac{}{t}},𝐗_{\eta _2}=\sigma {\displaystyle \frac{}{x_2}},𝐗_3=\rho {\displaystyle \frac{}{x_3}},𝐗_{\eta _3}=x_3{\displaystyle \frac{}{x_3}},`$ (47)
that forms a solvable Lie algebra. The commutation relations are given by
$`[𝐗_\tau ,𝐗_{\eta _2}]=0,[𝐗_\tau ,𝐗_3]=0,[𝐗_\tau ,𝐗_{\eta _3}]=0,`$
$`[𝐗_{\eta _2},𝐗_3]=0,[𝐗_{\eta _2},𝐗_{\eta _3}]=0,[𝐗_3,𝐗_{\eta _3}]=𝐗_3.`$ (48)
and correspond to direct products of two abelian groups and and the group $`G_2IIb`$ which has the generators
$`𝐗_3=\rho {\displaystyle \frac{}{x_3}}`$
$`𝐗_{\eta _3}=x_3{\displaystyle \frac{}{x_3}}.`$ (49)
We also analyse the Landau problem with a constant magnetic field in the $`x_3`$ direction. The equation of motion is given by,
$`\ddot{x}_k`$ $`=`$ $`\beta \epsilon _{kbc}\dot{x}_bB_c=\omega _k`$ (50)
where $`𝐁=(0,0,B)`$, with $`B`$ a constant. The vector field obtained is
$`𝐗_1={\displaystyle \frac{}{x_1}},𝐗_2={\displaystyle \frac{}{x_2}},𝐗_3=x_1{\displaystyle \frac{}{x_2}}x_2{\displaystyle \frac{}{x_1}},𝐗_4={\displaystyle \frac{}{t}}.`$ (51)
with commutation relations
$`[𝐗_1,𝐗_2]=0,[𝐗_1,𝐗_3]=𝐗_2,[𝐗_2,𝐗_3]=𝐗_1,[𝐗_a,𝐗_4]=0.`$ (52)
Thus we have found the generators that specifies the corresponding symmetry groups for all our examples considered in the introduction.
## 4 Symmetry generators, quantum particle
The simplest one dimensional version of the equation
$`\ddot{x}_a={\displaystyle \frac{\mu ^2x_a}{mr^4}}`$ (53)
possesses remarkable symmetries which were exploited by de Alfaro, Fubini, and Furlan to construct conformal quantum mechanics. Here $`x`$ is considered as a field in zero space and one time dimension. The quantum mechanical equation for the wave function $`u`$ becomes, in our notation,
$`\left({\displaystyle \frac{d^2}{dx^2}}+{\displaystyle \frac{g}{x^2}}+{\displaystyle \frac{x^2}{a^2}}\right)u={\displaystyle \frac{4r}{a}}u`$ (54)
where $`\frac{\mu ^2}{m}`$ is replaced by $`g`$. Here $`a`$ is a constant which plays a fundamental role in the theory and $`r`$ is related to appropriate raising and lowering operators. This equation can be expressed in terms of differential operator realization of $`su(1,1)`$ algebra and was studied in detail in.
There has been earlier works, where it has been shown that the dynamics of a scalar particle approaching the event horizon of a blackhole is governed by an Hamiltonian with an inverse square potential. The scalar field can be used as a probe to study the geometry in the vicinity of the horizon and its dynamics is expected to provide clues to the inherent symmetry properties of the system. The Hamiltonian of conformal quantum mechanics fits into this. This Hamiltonian also arise as a limiting case of the brickwall model describing the low energy quantum dynamics of a field in the background of a massive Schwarzschild blackhole of mass $`M`$ . By factorizing such a Hamiltonian, Birmingham, Gupta and Sen have found the Virasoro symmetry of the system and have studied the representation of the algebra as well as the scaling properties of the time independent modes. They had obtained the full Virasoro algebra by the requirement of unitarity of the representation. The Hamiltonian operator is in the enveloping algebra. However, here we aim to find the underlying Lie point symmetry of the equation of motion of the scalar particle, viewed as a differential equation. Of course, mathematically any two linear homogeneous ordinary differential equations can be transformed to the form,where a prime denotes a differentiation with respect to $`x`$,
$`u^{\prime \prime }=0.`$ (55)
This equation has the eight dimensional symmetry of projective transformations. But the two equations could be different from the physics point of view having different eigenvalues and eigenfunctions. Hence we would like to see explicitly what are the Lie point symmetries of the particular equation.
For the equation
$`u^{\prime \prime }=\omega (x,u,u^{})`$ (56)
where
$`\omega (x,u,u^{})=({\displaystyle \frac{C}{x^2}}+{\displaystyle \frac{D}{x}}+\widehat{E})u(x).`$ (57)
the symmetry generators are obtained from the conditions given by the equation (12) which reduces in the one dimensional case to
$`\omega (\eta _{,u}2\tau _{,x}3u^{}\tau _{,u})\omega _{,x}\tau \omega _{,u}\eta \omega _{,u^{}}[\eta _{,x}+u^{}(\eta _{,u}\tau _{,x})u_{}^{}{}_{}{}^{2}\tau _{,u}]`$
$`+\eta _{,xx}+u^{}(2\eta _{,xu}\tau _{,xx})+u_{}^{}{}_{}{}^{2}(\eta _{,uu}2\tau _{,xu})u_{}^{}{}_{}{}^{3}\tau _{,uu}=0`$ (58)
For the case $`C=g`$, $`D=\widehat{E}=0`$, equating to zero the coefficients of $`u_{}^{}{}_{}{}^{3}`$ and $`u_{}^{}{}_{}{}^{2}`$ in (58) we get
$`\tau _{,uu}=0,\eta _{,uu}=2\tau _{,xu}`$ (59)
which are satisfied for
$`\tau =u\alpha (x)+\beta (x),\eta =u^2\alpha ^{}(x)+u\gamma (x)+\delta (x).`$ (60)
Using these and equating to zero the coefficient of $`u^{}`$ and then considering the the terms not involving $`u^{}`$, we find that an interesting symmetry exists only when the coupling constant $`g`$ is equal to 2. For this case we obtain
$`\tau ={\displaystyle \frac{1}{x}}Au+Fx`$
$`\eta ={\displaystyle \frac{1}{x^2}}Au^2+Bu+\delta (x)`$ (61)
where $`A,F`$, and $`B`$ are constants and $`\delta (x)`$ satisfies the same equation as $`u`$ does. The vector fields are
$`𝐗_1=x{\displaystyle \frac{d}{dx}},𝐗_\mathrm{𝟐}=u{\displaystyle \frac{d}{du}},𝐗_3={\displaystyle \frac{1}{x}}u{\displaystyle \frac{d}{dx}}{\displaystyle \frac{1}{x^2}}u^2{\displaystyle \frac{d}{du}}`$ (62)
with commutation relations
$`[𝐗_1,𝐗_2]=0,[𝐗_1,𝐗_3]=2𝐗_3,[𝐗_2,𝐗_3]=𝐗_3`$ (63)
For the case, considered in , the relevant equation is
$`{\displaystyle \frac{d^2u}{dx^2}}+{\displaystyle \frac{1}{x^2}}\left[{\displaystyle \frac{1}{4}}+R^2E^2\right]u=0`$ (64)
where $`E`$ is a generic eigenvalue and $`R=2M`$. Hence $`g`$ corresponds to $`\left[\frac{1}{4}+R^2E^2\right]`$ in this case. However, our result shows that only when $`E`$ is imaginary with $`\left[\frac{1}{4}+R^2E^2\right]=2`$, the symmetry will show up. Further, the $`\frac{1}{x}`$ and $`\frac{1}{x^2}`$ factors in $`𝐗_3`$ makes it ill defined as $`x0`$ similar to the $`L_n`$ operators of conformal field theory or the $`P_m`$ operators considered by Birmingham, Gupta, and Sen.
## 5 Complete Krause symmetry
Krause has introduced the concept of the complete symmetry group of a system by specifying two extra properties in the definition of a Lie symmetry group. This requires the manifold of solutions to be a homogeneous space on which the group action takes place and the group is specific to the system with no other system admitting it.
Besides the Lie point symmetries and the contact symmetries, new types of symmetries are to be included in order to obtain the complete symmetry group. For an N dimensional system, the generators of the new symmetry was defined to be
$`Y=\left[{\displaystyle \xi (t,x_1,x_2,\mathrm{},x_N)𝑑t}\right]_t+{\displaystyle \underset{k=1}{\overset{N}{}}}\eta _k(t,x_1,x_2,\mathrm{},x_N)_{x_k}`$ (65)
which is different from the generators of a Lie point transformation because of the appearance of the integral of $`\xi `$. This makes it a nonlocal operator.
Nucci has developed a method based on the reduction of order to derive all Lie symmetries. Later Nucci and Leach have found the existence of more nonlocal symmetries. These symmetries become local on reduction of order. In this technique, for an autonomous system, one of the unknown function is taken as the new independent variable and the system is written in the modified form. Then the standard Lie group analysis of this transformed system yields the extra symmetries leading to a complete attainment of Krause symmetries. The method can be extended to include non-autonomous systems. Using a different technique, it has also again being found that the three dimensional Kepler problem is completely specified by six symmetries.
Besides the original Kepler problem, this method has been used to analyze many problems of physics, space science, meteorology etc. which include the Kepler problem with a drag, motion in an angle dependent force, Lorenz equation, Euler-Poinsot system and Kowalevsky top, as well as relativistically spherically symmetric systems. Nucci’s interactive code for determination of Lie symmetries has been used to arrive at the above results.
We follow the method and notation of references to determine the Krause symmetries for our examples. In our case the system of differential equations are given by
$`\ddot{x}_k=F_k(x_1,x_2,x_3,\dot{x}_1,\dot{x}_2,\dot{x}_3)`$ (66)
where $`k=1,2,`$ and $`3`$ and the $`\omega _k`$ of equation(1) equals $`F_k`$. By standard techniques one obtains the generators of the Lie point group for this system and a generator is written in the form
$`X=\tau (t,x_1,x_2,x_3)_t+{\displaystyle \underset{k=1}{\overset{3}{}}}\eta _k(t,x_1,x_2,x_3)_{x_k}`$ (67)
To treat the velocities in the same footing as the coordinates and for reduction of order, equation(66) is next made into a set of six ordinary differential equations
$`\dot{u}_k=u_{3+k}`$ (68)
$`\dot{u}_{3+k}=F_k(u_1,u_2,u_3,u_4,u_5,u_6)`$ (69)
where $`u_k`$s are related to $`x_k`$s and $`\dot{x}_k`$s. Then one of the dependent variables, $`u_i`$s, is chosen as the new independent variable $`y`$. We take $`u_3=y`$. The system (68)-(69) is now converted to a set of five ordinary differential equations depending on the variable $`y`$, with
$`{\displaystyle \frac{du_j}{dy}}={\displaystyle \frac{u_{3+j}}{u_6}},`$ (70)
$`{\displaystyle \frac{du_{3+j}}{dy}}={\displaystyle \frac{F_j(u_1,u_2,y,u_4,u_5,u_6)}{u_6}},`$ (71)
$`{\displaystyle \frac{du_6}{dy}}={\displaystyle \frac{F_3(u_1,u_2,y,u_4,u_5,u_6)}{u_6}},`$ (72)
where $`j=1,2`$. Using equation (70) we obtain
$`u_{3+j}=u_6{\displaystyle \frac{du_j}{dy}}.`$ (73)
This is put back in equations (71) and (72) to give the two ordinary second order equations and one first order equation for the unknowns $`u_j=u_j(y)`$, and $`u_6=u_6(y)`$
$`u_{j}^{}{}_{}{}^{\prime \prime }={\displaystyle \frac{\left[F_j(u_1,u_2,y,u_{}^{}{}_{1}{}^{},u_{}^{}{}_{2}{}^{},u_6)F_3(u_1,u_2,y,u_{}^{}{}_{1}{}^{},u_{}^{}{}_{2}{}^{},u_6)u_{}^{}{}_{j}{}^{}\right]}{u_{6}^{}{}_{}{}^{2}}},`$ (74)
$`u_6^{}={\displaystyle \frac{1}{u_6}}F_3(u_1,u_2,y,u_{}^{}{}_{1}{}^{},u_{}^{}{}_{2}{}^{},u_6),`$ (75)
where a prime denotes differentiation with respect to $`y`$. For the above system we write a generator for the Lie symmetry group as
$`Z=V(y,u_1,u_2,u_6)_y+{\displaystyle \underset{j=1}{\overset{2}{}}}G_j(y,u_1,u_2,u_6)_{u_j}+G_6(y,u_1,u_2,u_6)_{u_6}.`$ (76)
These can be transformed to the old form, $`Y`$, of the operators of equation (65) by replacing $`u_j,y,u_6`$ with $`x_j,x_3,\dot{x_3}`$, respectively, and solving the following system of equations for $`\xi `$ and $`\eta _k`$
$`Y(x_j)\eta _j=G_j,`$ (77)
$`Y(x_3)\eta _3=V,`$ (78)
$`Y^{(1)}{\displaystyle \frac{d\eta _3}{dt}}\xi \dot{x_3}=G_6,`$ (79)
with $`Y^{(1)}`$ being the first prolongation of Y.
On application of the above technique to our example corresponding to the equation of motion (15),
$`\ddot{x}_k`$ $`=`$ $`\epsilon _{kbc}\dot{x}_bx_c=\omega _k`$
we obtain the following set of equations.
$`\mathrm{\Omega }_1u_{1}^{}{}_{}{}^{\prime \prime }={\displaystyle \frac{1}{u_{6}^{}{}_{}{}^{2}}}\left[(u_{2}^{}{}_{}{}^{}yu_2)u_{1}^{}{}_{}{}^{}(u_{1}^{}{}_{}{}^{}u_2u_{2}^{}{}_{}{}^{}u_1)\right]`$ (81)
$`\mathrm{\Omega }_2u_{2}^{}{}_{}{}^{\prime \prime }={\displaystyle \frac{1}{u_{6}^{}{}_{}{}^{2}}}\left[(u_1u_{1}^{}{}_{}{}^{}y)u_{2}^{}{}_{}{}^{}(u_{1}^{}{}_{}{}^{}u_2u_{2}^{}{}_{}{}^{}u_1)\right]`$ (82)
$`\mathrm{\Omega }_6u_{6}^{}{}_{}{}^{}=1\left(u_{1}^{}{}_{}{}^{}u_2u_{2}^{}{}_{}{}^{}u_1\right)`$ (83)
For the determination of $`\xi `$ and the $`\eta `$’s we note that, first denoting $`\eta _3=\zeta `$,
$`{\displaystyle \frac{\stackrel{~}{u_l}}{\stackrel{~}{y}}}={\displaystyle \frac{u_l}{y}}+ϵ\left[\eta _{l,y}\zeta _{,y}u_{l}^{}{}_{}{}^{}+\left(\eta _{l,u_i}u_{l}^{}{}_{}{}^{}\zeta _{,u_i}\right)u_{i}^{}{}_{}{}^{}\right]`$ (84)
and, for example,we find
$`\eta _{6}^{}{}_{}{}^{}=\eta _{6,y}+\left(\eta _{6,u_i}\zeta _{,u_i}u_{6}^{}{}_{}{}^{}\right)u_{i}^{}{}_{}{}^{}\zeta _{,y}u_{6}^{}{}_{}{}^{}`$ (85)
with $`u_{6}^{}{}_{}{}^{}=\mathrm{\Omega }_6`$. Denoting the vector-fields by $`\overline{𝐗}`$, the symmetry condition for our first order equation is given by
$`\overline{𝐗}\mathrm{\Omega }_6=\eta _{6,y}+\left(\eta _{6,u_i}\zeta _{,u_i}\mathrm{\Omega }_6\right)u_{i}^{}{}_{}{}^{}\zeta _yu_{6}^{}{}_{}{}^{}`$ (86)
where
$`\overline{𝐗}=\zeta (y,u_i){\displaystyle \frac{}{y}}+\eta _j(y,u_i){\displaystyle \frac{}{u_j}}`$ (87)
and for the second order equations the symmetry conditions take the form
$`\zeta \mathrm{\Omega }_{1,y}+\eta _1\mathrm{\Omega }_{1,u_1}+\eta _2\mathrm{\Omega }_{1,u_2}+\eta _6\mathrm{\Omega }_{1,u_6}`$
$`+\left[\eta _{1,y}+\left(\eta _{1,u_i}\zeta _{,u_i}u_{1}^{}{}_{}{}^{}\right)u_{i}^{}{}_{}{}^{}\zeta _{,y}u_{1}^{}{}_{}{}^{}\right]\mathrm{\Omega }_{1,u_{1}^{}{}_{}{}^{}}`$
$`+\left[\eta _{2,y}+\left(\eta _{2,u_i}\zeta _{,u_i}u_{2}^{}{}_{}{}^{}\right)u_{i}^{}{}_{}{}^{}\zeta _{,y}u_{}^{}{}_{2}{}^{}\right]\mathrm{\Omega }_{1,u_{2}^{}{}_{}{}^{}}`$
$`\left[\eta _{1,u_1}2\zeta _{,y}3\zeta _{,u_i}u_{i}^{}{}_{}{}^{}\right]\mathrm{\Omega }_1\eta _{1,yy}\left(2\eta _{1,yu_i}\zeta _{,yu_i}\right)u_{i}^{}{}_{}{}^{}+\zeta _{,yy}u_{1}^{}{}_{}{}^{}`$
$`+\zeta _{,u_i}u_{1}^{}{}_{}{}^{}\mathrm{\Omega }_i\eta _{1,u_ju_i}u_{i}^{}{}_{}{}^{}u_{j}^{}{}_{}{}^{}+\zeta _{,yu_i}u_{1}^{}{}_{}{}^{}u_{i}^{}{}_{}{}^{}+\zeta _{,u_iu_j}u_{1}^{}{}_{}{}^{}u_{i}^{}{}_{}{}^{}u_{j}^{}{}_{}{}^{}=0.`$ (88)
In the above the appropriate equations (81),(82), and (83) are to be substituted. For equations(86) and (88) to be compatible we have
$`\zeta =y,\eta _1=u_1,\eta _2=u_2`$ (89)
and we obtain
$`\xi =1`$ (90)
For our example of equation (39) which is
$`\ddot{x}_a=\dot{x}_ax_kx_k`$
the equations for $`u_i`$’s become
$`\mathrm{\Omega }_1u_{1}^{}{}_{}{}^{\prime \prime }=0,\mathrm{\Omega }_2u_{2}^{}{}_{}{}^{\prime \prime }=0,\mathrm{\Omega }_6u_{6}^{}{}_{}{}^{}=u_{1}^{}{}_{}{}^{2}+u_{2}^{}{}_{}{}^{2}+y^2.`$ (92)
The compatibility of the above three equations forces $`\xi `$ and $`\eta `$’s to be
$`\zeta =y,\eta _1=u_1,\eta _2=u_2,\eta _6=3u_6.`$ (93)
Thus the vector field is given by,
$`\overline{𝐗}=y{\displaystyle \frac{}{y}}+u_1{\displaystyle \frac{}{u_1}}+u_2{\displaystyle \frac{}{u_2}}+3u_6{\displaystyle \frac{}{u_6}}`$ (94)
which has a corresponding part in $`\dot{𝐗}_5`$ of equation(41). From equation (79) we obtain $`\xi `$ to be
$`\xi =2.`$ (95)
In the case of a charged particle in the magnetic field of a monopole, with the equation of motion given by (equation (37)),
$`\ddot{x}_a=\epsilon _{abc}{\displaystyle \frac{\dot{x}_bx_c}{r^3}}=\omega _a`$
we get
$`\mathrm{\Omega }_1u_{1}^{}{}_{}{}^{\prime \prime }={\displaystyle \frac{1}{u_6}}\left(u_2u_{2}^{}{}_{}{}^{}yu_{1}^{}{}_{}{}^{}{}_{}{}^{2}u_2u_1u_{1}^{}{}_{}{}^{}u_{2}^{}{}_{}{}^{}\right)`$ (96)
$`\mathrm{\Omega }_2u_{2}^{}{}_{}{}^{\prime \prime }={\displaystyle \frac{1}{u_6}}\left(u_{1}^{}{}_{}{}^{}yu_1+u_{1}^{}{}_{}{}^{}u_2u_{2}^{}{}_{}{}^{}u_1\right)`$ (97)
$`\mathrm{\Omega }_6u_{6}^{}{}_{}{}^{}=u_1u_{2}^{}{}_{}{}^{}u_2u_{1}^{}{}_{}{}^{}`$ (98)
The symmetry conditions are satisfied for
$`\zeta =y,\eta _1=u_1,\eta _2=u_2,\eta _6=u_6`$ (99)
and this gives rise to the vector field
$`\overline{𝐗}=y{\displaystyle \frac{}{y}}+u_1{\displaystyle \frac{}{u_1}}+u_2{\displaystyle \frac{}{u_2}}u_6{\displaystyle \frac{}{u_6}}.`$ (100)
Consequently one obtains,
$`\xi =2.`$ (101)
For the case of the last example,i.e, the Landau problem with a constant magnetic field, the equations of motion are
$`\ddot{x}_k=\epsilon _{klm}\dot{x}_lB_m`$ (102)
with $`𝐁=(0,0,\frac{1}{x^2})`$. After reduction of order the equations are
$`\mathrm{\Omega }_1u_{}^{\prime \prime }{}_{1}{}^{}={\displaystyle \frac{u_{}^{\prime \prime }{}_{2}{}^{}}{u_6u_1}},\mathrm{\Omega }_2u_{}^{\prime \prime }{}_{2}{}^{}={\displaystyle \frac{u_{1}^{}{}_{}{}^{}}{u_6u_1^2}},\mathrm{\Omega }_6u_{6}^{}{}_{}{}^{}=0.`$ (103)
The analysis results in
$`\zeta =y,\eta _1=u_1,\eta _6=u_6.`$ (104)
So we get, also for this case,
$`\xi =2.`$ (105)
This gives us the vector fields
$`𝐯_1=y{\displaystyle \frac{}{y}},𝐯_2=u_1{\displaystyle \frac{}{u_1}},𝐯_3=u_6{\displaystyle \frac{}{u_6}},`$ (106)
So we get back the vector field $`𝐗_{\eta _3}`$ in $`𝐯_1`$, but there are now two extra vector fields $`𝐯_2`$ and $`𝐯_3`$, which are to be included in the complete symmetry group.
Thus we have found the generators that specifies the corresponding symmetry groups for all our examples as well as the generators of the complete symmetry of Krause. In our examples, we find that $`\eta _k`$’s do not depend on $`\dot{x}_N`$ and $`\xi `$ is a constant, and hence $`Z`$ can be transformed into a generator of a Lie point symmetry group.
## 6 Conclusion
As has already been explicitly mentioned , the equations of motion of a free particle admit eight symmetries for each of the $`x_a`$s. This is the maximum number of symmetries for an ordinary second order differential equation. By including different $`x_a`$, $`\dot{x_a}`$ dependent terms in the equations we do explicitly see which generators survive as symmetries and we have found corresponding complete Lie algebras. We have chosen some cases motivated by problems from physics. The original motivation of including Wess-Zumino terms in the Lagrangian has been to reduce some of its symmetries, and here we find that the equations of motion now support the three dimensional rotations and a time translation symmetry instead of the six vector fields as for the monopole problem without any constraint. For the other examples considered here, we get some interesting result in the form of Kepler’s scaling law and the full structure of the inherent symmetry group. This we get without solving the equations of motion. In the cases where a Lagrangian could be set up, those generators operating on the Lagrangian giving zero include all the Noether symmetries.
It is also expected that related group analysis may provide useful information when terms are modified in the Lagrangian, due to quantum corrections, for example. We have shown explicitly how many and which generators remain as symmetries. These symmetries correspond to the transformations of the solutions. Similar analysis for specific cases of some nonlinear equations arising out of linear equations have also been carried out in . We expect that these will ultimately lead to a better understanding of the spontaneous symmetry breaking, as well.
We have also carried out the analysis using Nucci’s method of reduction of order to find the complete Krause symmetries in four of our examples. It is found that in our last two examples the generators of the complete symmetry group can be transformed into a generator of a Lie point symmetry group.
For the case of a scalar particle probing the near horizon structure of a blackhole, under certain limits the Hamiltonian contains $`\frac{1}{r}`$ and $`\frac{1}{r^2}`$ potentials. It is found that there exist a symmetry with three generators only for specific constant value of the coefficient of the $`\frac{1}{r^2}`$ term and in the absence of $`\frac{1}{r}`$ term.
The quantum mechanical problem of a charged particle in the presence of even a constant magnetic field has many interesting mathematical structures and under certain limits can make space coordinates noncommutative . Klishevich and Plyuschay have found a universal algebraic structure at the quantum level for the two dimensional case in the presence of certain magnetic fields. Nonlinear superconformal symmetry of the fermion-monopole system has been extensively studied in. It would be a motivation to seek the existence of analogous structures for our cases. It would be interesting to study these quantum aspects as well as the non-Abelian quantum kinematics in the framework of group theoretic quantization programme of Krause for above nonuniform magnetic fields.
Acknowledgement I would like to thank Professor Dieter Lüst for the warm hospitality and for providing the academic facilities at the Arnold Sommerfeld Centre, Ludwig Maximilians University, Munich, where this work has been done. I am grateful to Professor F. Haas for pointing out the reference to me. |
warning/0506/cond-mat0506037.html | ar5iv | text | # Diagnosis of weaknesses in modern error correction codes: a physics approach
## Appendix A Appendix A: Supplementary Figures
Figure S1. Parity check matrix $`\widehat{H}`$ for the $`(155,64,20)`$ LDPC code. The matrix consists of $`3\times 5`$ blocks. Each block is a square $`31\times 31`$ matrix. Empty/filled elements of the matrix stand for $`0/1`$. Bits are numbered from “0” to “154”. The girth of this code is eight.
Figure S2. Frame-Error-Rate (FER) vs SNR<sup>2</sup> for the $`(155,64,20)`$ code and Belief Propagation decoding. The filled/empty circle-marks correspond to result of Monte Carlo evaluation of FER for $`4/1024`$ iterations of BP. The straight/dashed line corresponds to the (a)-instanton asymptotic, $`\mathrm{exp}[(46^2/210)s^2/2]`$ and the ML asymptotic, $`\mathrm{exp}[20s^2/2]`$.
Figure S3. Min-sum decoding for the (a) instanton on the computational tree. All numbers are in units of $`1/105`$. The numbers shown on the bits/circles are the values of the magnetic fields, $`h`$. The number shown next to each check/square is the message, $`\eta `$, arriving at the check on the $`q`$-th step of the iteration procedure, where $`4q`$ corresponds to the number of bits separating this check from the tree center. Other notations/marks are in accordance with the caption of Fig. 1 of the main text.
Figure S4. Min-sum decoding for the (b) instanton on the computational tree. All numbers are in units of $`1/79`$. Other notations/marks and explanations are in accordance with captions of Fig. 1 of the main text and Fig. S3.
Figure S5. Min-sum decoding for the (c) instanton on the computational tree. All numbers are in the units of $`1/47`$. Other notations/marks and explanations are in accordance with captions of Fig. 1 of the main text and Fig. S3.
Figure S6. Geometrical interpretation of the degeneracy observed in case (b). The plot shows distance $`l^2`$ for the noise configuration in case (b), minimized with respect to all the magnetic fields except $`h_0`$ and $`h_{77}`$, vs the two remaining fields. The relevant part of the $`(h_0,h_{77})`$ plane is the quadrant $`h_0<0`$, $`h_{77}>0`$. Looking for the instanton in the $`|h_0|<|h_{77}|`$ domain and thus assuming that the “0” bit connected to the red check in Fig. 1B of the main text is colored while bit “77” is pale, one gets the paraboloid, shown in pink in the Figure, that achieves its minimum outside of the domain, i.e. in the $`|h_0|>|h_{77}|`$ semi-quadrant. Similarly, if one looks for the instanton in the $`|h_0|>|h_{77}|`$ domain one again finds an inconsistency as the respective paraboloid, shown in yellow in the Figure achieves its minimum in the $`|h_0|<|h_{77}|`$ semi-quadrant. Therefore, the point of actual minimum, shown by the green dot on the Figure, lies exactly at the minimum of the angled join of the two domains, $`|h_0|=|h_{77}|`$.
## Appendix B Appendix B: Instantons for the min-sum decoding
These notes consist of three parts. The first part is devoted to explaining how the entire instanton family for an arbitrary LDPC code decoded by the min-sum algorithm can be fully characterized using the computational tree approach. The second part describes an alternative exposition which allows one to represent an instanton as a configuration of magnetic fields equidistant from some set of codewords on the computational tree. The third part formulates a relation between the theoretical and numerical approaches and suggests challenges and questions that need to be addressed in the future.
### B.1 Colored/signature structure and constraint minimization
The basic object for our construction is the computational tree and its colored/pale/uncolored parts (as briefly introduced in the main text). The computational tree is a tree constructed by a simple unwrapping of the Tanner graph of the code into a tree starting from the bit where the BER is calculated. The number of generations on the tree is equal to the number of iterations of the decoder, $`n_{\mathrm{it}}`$. If $`n_{\mathrm{it}}`$ is larger than, or equal to, a quarter of the code girth (defined as the length of the shortest loop of the original Tanner graph, measured by the number of edges within the loop), then the computational tree contains more than one replica of some of the bits of the original code.
Consider an arbitrary configuration of magnetic field, $`𝒉`$ (or noise field, $`𝝋`$) on the Tanner graph, and on the computational tree respectively. Calculating the magnetization (or switching from physics jargon to communication theory jargon, a-posteriori log-likelihood) at the $`n_{\mathrm{it}}`$-th iteration in the center of the tree one derives, $`m_{\mathrm{center}}^{(n_{\mathrm{it}})}=𝒉𝒏`$. $`n_i`$, defined on a bit of the original Tanner graph is an integer. It is a sum of contributions, each originating from the respective bit/replica on the computational tree. Colored are bits on the tree that contribute the integer $`+1`$ or $`1`$. (In the Figures of the main text and Supporting Figures we use different colors to identify bits in the computational tree originating from different bits of the original Tanner graph.) A colored bit on the computational tree has the signature $`+1/1`$ if it contributes $`+1/1`$ to the integer associated with the respective bit on the Tanner graph. Uncolored bits (i.e. bits not shown in the Figures) or pale bits (i.e. bits shown pale in the Figures) on the tree do not contribute to the respective integer (one may also say that the respective contribution is just zero) according to the min-part of the min-sum rules described in Eq. (1) of the main text. We draw a bit on the computational tree in pale if it does not contribute to respective integer, however at least one of its siblings, i.e. bits on the tree originating from the same bit on the computational tree, does contribute to the integer.)
Let us now describe how an individual contribution of a colored bit on the tree to the respective integer (that is $`+1`$ or $`1`$) is calculated. We aim to calculate the contribution to magnetization at the tree center counting integers according to the min-sum rule. We assign $`+1`$ signatures to the colored leaves of the tree and start an iterative procedure which assigns signatures moving from the tree leaves towards the tree center. Consider the case when at a certain step of the iteration procedure a check receives messages from some number of bits among which only one is colored. Then one calculates the product of the signatures associated with the messages this check receives from the remaining bits. If the resulting product is $`+1/1`$ the signature of the colored bit is $`+1/1`$ and the signatures of the colored bits, laying on the tree branch grown from the given colored bit, do not/do change. The other possibility (that will be called degenerate) is that a check receives two (or more) messages that all have the same absolute value, which is also the minimal of all the messages received. Then, one has the freedom to color only one of the degenerate bits with a colored branch grown from it with the signatures assigned as described above. The iterative procedure of the signature assignment is terminated once the three center is reached. One calls a check marked if it lies in between two colored bits of different signatures.
The union of all colored bits, i.e. bits contributing to $`m_{\mathrm{center}}^{(n_{\mathrm{it}})}`$, is called the colored structure. Any check connected to the colored structure is actually connected to two bits of the structure. Another important characteristic complementing the notion of the colored structure on the computational tree is the set of aforementioned signatures $`\pm 1`$ associated with any bit of the colored structure. In a degenerate case, one finds multiple colored/signature structures associated with a given configuration of magnetic fields. Each degenerate structure will actually correspond to a distinct linear combination of magnetic fields equal to $`m_{\mathrm{center}}^{(n_{\mathrm{it}})}`$. Therefore, of the whole variety of possible degenerate colored/signature structures corresponding to the same magnetic field one can always select a set of linearly independent ones.
So far, this description was generic, i.e. not restricted to a specific configuration of the magnetic field. Let us now fix the family of linear independent colored/signature structures just explained and allow variation in the value of magnetic fields. Our goal here is to find an instanton conditioned to the specific form of the family of the linear independent colored/signature structures. Finding the instanton means minimizing $`l^2=(\mathrm{𝟏}𝒉)^2`$ with respect to $`𝒉`$ under the additional set of linearly independent conditions, $`m^{(n_{\mathrm{it}})}=𝒉𝒏^{(\mu )}=0`$, where $`\mu `$ is an index enumerating these conditions each corresponding to a certain colored structure. (The expression presented above for the length $`l`$ applies to the white Gaussian channel, however generalization for any other channel is straightforward.) The resulting expression for the optimal configuration of the noise, $`𝝋=\mathrm{𝟏}𝒉`$, is
$$𝝋=\underset{\mu ,\nu }{}𝒏^{(\mu )}(\widehat{G}^1)_{\mu \nu }\underset{i}{}n_i^{(\nu )},G_{\mu \nu }𝒏^{(\mu )}𝒏^{(\nu )}.$$
(S1)
This expression (also generalizing Eq. (3) of the main text for an arbitrary type of instanton degeneracy) should be checked for consistency with the family of the colored/signature structures assumed for the instanton. If the consistency check is met, the instanton construction is completed.
Let us now demonstrate how this formal description works for the three instanton examples (a), (b) and (c) described in Figs. 1, 2 of the main text and also illustrated in Figs. S3–S6. Instantons (a) and (c) are both explained by a single colored structure. For the (a) instanton each bit contributes $`+1`$ to the respective component of $`𝒏`$. For the (c) example all contributions are $`+1`$ except of the one coming from the “75” bit connected to the marked check. Since $`h_0<0`$, the message originating from this bit contributes with the opposite sign to the magnetization. There exist $`6`$ replicas of the “75” bit on the computational tree, however taking into account that one replica of the bit contributes $`1`$, one finds that the actual value of the noise is $`4`$ rather than $`6`$. Since $`|h_0|>|h_{75}|`$, the “0” bit is pale thus it does not contribute to the magnetization. Considering the (b) instanton one finds that this is a degenerate example, with $`|h_0|=|h_{77}|`$. There are two linearly independent colored/signature structures describing the (b) instanton. The two structures are different only at the two bits on the tree leaves shown in Fig. 1B of the main text adjusted to the marked check. The first colored structure does not contain bit “77” (zero contribution to the magnetization) while bit “0” contributes $`+1`$ to the respective component of $`𝒏`$. The second structure does not contain bit “0” while bit “77” contributes $`1`$ to the respective component of $`𝒏`$ (simply because $`h_0<0`$, thus forcing the respective message to contribute to the magnetization with the opposite sign).
One natural question to ask about the degenerate case (b) is the following: can one of the two colored structures describing the instanton be forming its own non-degenerate instanton? The answer is negative. Indeed, the colored/signature structure generates (through the minimization procedure described above) such a configuration of magnetic fields that will not be consistent with the colored/signature structure one started from. Considering the structure with the “0” bit connected to the marked check being pale (and thus the signature field associated with the colored bit “77” connected to the red check being $`1`$) one finds that the magnetic field minimizing $`l`$ will actually be inconsistent with the colored/signature structure. Considering the other configuration (the “0” bit is colored with $`+1`$ signature while the “77” bit is pale) one finds again that the resulting magnetic field is inconsistent with the colored structure. To resolve this inconsistency one needs to account for the two configurations simultaneously, thus introducing two constraints, not one. The degeneracy of the (b) instanton is illustrated in Fig. S6.
Let us also notice that the discrete nature of consistency check (yes/no answer as a result) puts degenerate configurations on equal footing (in the sense of counting all possibilities) with the non-degenerate configurations: considering the family of all possible instantons for the given computational tree one finds that the number of degenerate instantons is comparable with the number of nondegenerate instantons.
### B.2 Instantons as medians between pseudo-codewords
We consider an instanton with the set of linearly independent structures already established. Each colored structure, indexed by $`\mu `$, corresponds to the constraint, $`𝒉𝒏^{(\mu )}=0`$, imposed on the magnetic field, $`𝒉`$. Each of the constraints can actually be reformulated in terms of a pair of pseudo-codewords on the computational tree:
$$\underset{i\mathrm{tree}}{}h_i\sigma _i^{(\mu ;+)}=\underset{i\mathrm{tree}}{}h_i\sigma _i^{(\mu ;)},$$
(S2)
where $`i`$ stands for index assigned to a bit on the computational tree; the magnetic field on the computational tree bit is equal to the magnetic field defined on the respective bit of the original graph; and the pseudo-codewords, $`𝝈^{(\mu ;\pm )}`$, are the two distinct configurations of the binary field, $`\sigma _i=\pm 1`$, defined on each bit $`i`$ of the computational tree that satisfy all the checks on the computational tree.
Let us now discuss how the pseudo-codewords can be constructed if the respective constraint $`\mu `$ described by Eq. (S2), is already established. If the signature field, described in the previous Section of the Notes, corresponding to the structure $`\mu `$, does not contain a single $`1`$ element, then $`𝝈^{(\mu ;+)}`$ is the all unity codeword ($`+1`$ on all bits of the computational tree) and $`𝝈^{(\mu ;)}`$ is the pseudo-codeword containing $`1`$ on all the colored bits of the structure and $`+1`$ on all other bits. If, however, the colored structure does contain some $`1`$ signatures the situation is more elaborate as both pseudo-codewords are nontrivial. The algorithm that allows restoration of the pseudo-code words starts by determining the values of the colored bits for $`𝝈^{(\mu ;+)}`$ that are set equal to the values of the signatures. The uncolored bits are assigned values $`+1`$. Although it is possible to determine the bit values of the pale substructures the procedure is elaborate and we are not presenting it here. Finally, the pseudo-codeword $`𝝈^{(\mu ;)}`$ is obtained from $`𝝈^{(\mu ;+)}`$ by changing the signs of colored bits with the uncolored and pale bits remaining the same.
Fig. 2 of the main text shows three examples of the pseudo-codeword construction for the three instantons discussed in the manuscript. Examples (a), (c) contain one pair of competing pseudo-codewords. However, the two cases are different. In the case (a) the colored/signature structure does not contain $`1`$ bits, thus one pseudo-codeword is just the all unity codeword and another pseudo-codeword contains $`1`$ at all the bits of the colored structure and $`+1`$ at all other bits. In the (c) case the colored/signature structure does contain $`1`$ bit thus resulting in two distinct pseudo-codewords shown in Fig. 2C of the main text. Example (b) corresponds to the degenerate case with the two pairs of pseudo-codewords involved in the conditions (S2). However, one pseudo-codeword enters both conditions (that is the one shown on the top diagram of Fig. 2B of the main text) therefore the total count for the case (b) gives three pseudo-codewords being equidistant from the instanton configuration of the magnetic fields.
### B.3 General Remarks
Let us note that the analysis presented above, in addition to its theoretical significance, may be helpful for accelerating the instanton-amoeba numerical procedure, e.g. through guiding selection on the final stage of the minimization. We also expect that this theoretical analysis will be instrumental for formulating the right questions to address by the instanton-amoeba method, or by other minimization methods aiming at finding the instanton numerically. In what follows, we conclude by posing some questions that we did not yet study but plan to address in the future.
* More detailed exploration of the phase space, especially in the context of describing not only the minimal distance $`l_{\mathrm{min}}`$ contribution but also the family of other “low laying” instantons. The particular question of interest here is to estimate the “density of states/instantons”, that is to answer the question: how many instantons are found within the $`\delta l`$ vicinity of the one correspondent to $`l_{\mathrm{min}}`$?
* Dependence of BER on the number of iterations. As we already indicated in the main text our preliminary tests show that instantons and thus asymptotic estimates for BER do change with the number of iterations. We will be interested to explain this dependence. We will also be testing with our instanton-amoeba approach, the validity of the graph covers method suggested recently 03KV .
* Dependence on the code length. It is important to analyze the family of LDPC codes with varying code length, $`N`$. Of a special interest are the regular LDPC codes where the Hamming distance grows with $`N`$, e.g. Margulis codes 82Mar . Then, the relevant question is: how does $`l_{\mathrm{min}}`$ (and other characteristics of the error floor) change with $`N`$ for a given family of codes? This study will essentially lead to analysis of the finite-size effects, already discussed in the water-fall domain 04ARMR , but not yet explored in the asymptotic regime of the error-floor.
* Does BP/min-sum decoding perform better than other suboptimal algorithms (that can possibly exist) of the same complexity, e.g. linear in $`N`$? Even if the answer is yes (that is by no means guaranteed), what would be the best decoding for a higher level of complexity, e.g. $`N^a`$, where $`a>1`$? Once an idea of better decoding is formulated, our instanton-amoeba toolbox will be indispensable in answering the aforementioned questions and also testing in depth the performance of the new decoding.
* Other types of codes, e.g. turbo codes. Turbo codes show remarkable performance at moderate SNR but they are also infamous for demonstrating much higher (than comparable in size LDPC codes) error floors. Even though some important similarities between the LDPC codes and turbo codes are established 98MMC , the decoders of these two types of codes are different and it becomes important to analyze the performance of the turbo scheme, especially in light of the turbo-codes popularity.
* Other, application specific, channels. The instanton-amoeba approach is not limited to the white Gaussian channel, which we choose primarily for the purpose of demonstration, but can be applied straightforwardly to other types of channels, e.g. with correlations among received samples. Of special interest will be to analyze the performance of fiber-optic communication channels where the effects of fiber dispersion 03CCDGKL , birefringence and amplifier noise 04CCGKL will be accounted for. Another two interesting channel types are magnetic and optical recording channels exhibiting high level of nonlinearity and correlations among received samples 04BV\_book .
* There are many problems in the information and computer sciences that are different from standard coding problem but are also dependent or sensitive to rare errors. Therefore, estimating performance/BER in these problems is a major step required for their comprehensive analysis. Two interesting examples here are (i) inter-symbol interference, that is especially challenging in the context of two-dimensional 03WOSI and three-dimensional information storage, and (ii) estimating algorithmic errors in the domain of typically good performance within a combinatorial optimization K-SAT setting 02MPZ . |
warning/0506/nlin0506035.html | ar5iv | text | # Multiple permanent-wave trains in nonlinear systems
## 1 Introduction
Nonlinear wave systems have been studied for a few decades. Much progress has been made on integrable equations where the inverse scattering transform method can be applied . For non-integrable equations, the general analytical treatment has been elusive so far and will likely remain so in the near future. A less ambitious goal, then, would be to generally study the permanent waves in non-integrable systems. Such waves often contain valuable information on the system’s general solution behaviors. An interesting fact is that, in many non-integrable systems, simple permanent waves can be matched together and form multiple permanent-wave trains (, , , , , , etc.). If solitary waves exist, multiple solitary-wave trains can be constructed by a perturbation method proposed by Karpman and Solov’ev and Gorshkov and Ostrovsky (see ). If permanent waves with exponentially decaying and oscillating tails are present, the existence of countably infinite multiple permanent-wave trains has been proved for certain types of nonlinear systems by variational methods (, ). In this paper, if a nonlinear wave system allows permanent waves which exponentially approach a constant at infinity, we will construct widely-separated multiple permanent-wave trains by a new and general method, namely, the asymptotic tail-matching method. This method is stimulated in part by another matching method for non-local solitary waves (see and ). Under some general assumptions, we will show that an arbitrary number of permanent waves can be matched together and form multiple permanent-wave trains if and only if the exponential tails of these permanent waves satisfy certain simple algebraic conditions. These conditions will also determine the spacings between adjacent permanent waves if such matching takes place. This asymptotic tail-matching method differs from Karpman et al’s perturbation method in two major aspects. First, it can be applied directly to the matching of kink and anti-kink type permanent waves. Second, its results are explicit, simple and insightful. As applications of these general results, we will discuss fourth-order systems and the coupled nonlinear Schrödinger equations. For fourth-order systems which allow permanent waves exponentially and oscillatorily approaching a constant at infinity, we will show that countably infinite multiple permanent-wave trains exist and can be readily constructed. Thus the results in and are reproduced. For the coupled nonlinear Schrödinger equations, we will show that countably infinite multiple solitary-wave trains can be constructed in a large portion of the parameter space. Numerical results will also be presented and compared with the theoretical predictions when appropriate.
## 2 Construction of multiple permanent-wave trains
We consider a general nonlinear wave system
$$F(U,D_x,D_t)=0,$$
(2.1)
where $`U`$ is the unknown vector variable, and $`F`$ is a nonlinear vector function. Suppose it allows permanent waves of certain form which, when substituted into Eq. (2.1), reduces it into an autonomous complex system of first-order nonlinear ordinary differential equations
$$d\mathrm{\Phi }/dx=G(\mathrm{\Phi }),$$
(2.2)
where $`\mathrm{\Phi }(x)`$ is a $`n`$-component vector variable. If Eq. (2.2) has permanent wave solutions which exponentially approach a constant at infinity, then we next will develop a new method to determine if those permanent waves can be matched together and form widely-separated multiple permanent-wave trains or not. The idea is to perturb each permanent wave such that the exponential tails of each perturbed wave match those of the adjacent permanent waves. We first discuss the matching of solitary waves, followed by that of general permanent waves.
### 2.1 Solitary-wave trains
Suppose Eq. (2.2) allows solitary waves $`\mathrm{\Phi }(x)`$ which exponentially decay to zero as $`|x|\mathrm{}`$. We make the following general assumptions:
1. the eigenvalues of the constant (Jacobian) matrix $`G(0)`$ all have non-zero real parts;
2. for any solitary wave $`\mathrm{\Phi }(x)`$, the linear behavior dominates at infinity, i.e., as $`x\mathrm{}`$ or $`\mathrm{}`$, $`\mathrm{\Phi }(x)`$ approaches a solution of the linear equation
$$d\stackrel{~}{\mathrm{\Phi }}/dx=G(0)\stackrel{~}{\mathrm{\Phi }};$$
(2.3)
3. for any solitary wave $`\mathrm{\Phi }(x)`$, the linearized equation of (2.2) around $`\mathrm{\Phi }(x)`$
$$d\stackrel{~}{\mathrm{\Phi }}/dx=G(\mathrm{\Phi })\stackrel{~}{\mathrm{\Phi }}$$
(2.4)
and its adjoint equation
$$d\mathrm{\Psi }/dx=G^\text{T}(\mathrm{\Phi })\mathrm{\Psi }$$
(2.5)
each have a single linearly independent localized solution. Here “T” represents the transpose and “\*” the complex conjugate.
Remark: Since Eq. (2.2) is autonomous, any spatial translation of $`\mathrm{\Phi }(x)`$ is still (2.2)’s solution. Therefore Eq. (2.4) always has a nontrivial localized solution $`d\mathrm{\Phi }(x)/dx`$. The requirement for Eq. (2.4) is just that $`d\mathrm{\Phi }(x)/dx`$ is its single linearly independent localized solution. This can be guaranteed if $`\mathrm{\Phi }(x)`$ is isolated in $`H^1(R,R^2)`$ up to spatial translations (the so-called non-degenerency condition in some of the literature. See ).
We also introduce the following notations. In view of assumption A1, let us denote $`G(0)`$’s eigenvalues as $`\lambda _1,\lambda _2,\mathrm{},\lambda _n`$ where
$$\text{Re}(\lambda _1)\text{Re}(\lambda _2)\mathrm{}\text{Re}(\lambda _s)<0<\text{Re}(\lambda _{s+1})\text{Re}(\lambda _{s+2})\mathrm{}\text{Re}(\lambda _n).$$
(2.6)
Suppose for an eigenvalue $`\lambda `$, $`G(0)`$ has a chain of eigenvector and generalized eigenvectors $`v_i(i=1,\mathrm{},l)`$ such that
$$(G(0)\lambda I)v_1=0,$$
(2.7)
$$(G(0)\lambda I)v_{i+1}=v_i,i=1,\mathrm{},l1.$$
(2.8)
Define the polynomial functions $`\xi _i(x)(i=1,\mathrm{},l)`$ as
$$\xi _i(x)=v_i+xv_{i1}+\mathrm{}+\frac{x^{i1}}{(i1)!}v_1,i=1,\mathrm{},l,$$
(2.9)
then
$$\xi _{i+1}^{}(x)=\xi _i(x),$$
(2.10)
and $`\{\xi _i(x)e^{\lambda x},i=1,\mathrm{},l\}`$ form a chain of linearly independent solutions of Eq. (2.3). According to the theory of linear differential equations with constant coefficients, we can find such chains of solutions which together form a fundamental set of solutions of Eq. (2.3). Thus according to assumption A2, we have
$$\mathrm{\Phi }(x)\{\begin{array}{c}_{i=1}^sc_i\xi _i(x)e^{\lambda _ix},x\mathrm{},\hfill \\ _{i=s+1}^nc_i\xi _i(x)e^{\lambda _ix},x\mathrm{},\hfill \end{array}$$
(2.11)
where $`c_i(i=1,\mathrm{},n)`$ are complex constants. We point out that in the special case where $`G(0)`$ has $`n`$ linearly independent eigenvectors, $`\{\xi _i,i=1,\mathrm{},n\}`$ are just those constant eigenvectors. The fundamental matrix of the adjoint equation (2.5) at infinity is
$$[\eta _1e^{\lambda _1^{}x}\eta _2e^{\lambda _2^{}x}\mathrm{}\eta _ne^{\lambda _n^{}x}],$$
(2.12)
where
$$[\eta _1\eta _2\mathrm{}\eta _n]=\{[\xi _1\xi _2\mathrm{}\xi _n]^1\}^\text{T}.$$
(2.13)
Note that for $`1i,jn`$,
$$\xi _i(x)\eta _j^{}(x)=\{\begin{array}{c}1,i=j,\hfill \\ 0,ij.\hfill \end{array}$$
(2.14)
Thus for the single linearly independent localized solution $`\mathrm{\Psi }(x)`$ of Eq. (2.5), we have
$$\mathrm{\Psi }(x)\{\begin{array}{c}_{i=1}^sd_i\eta _i(x)e^{\lambda _i^{}x},x\mathrm{},\hfill \\ _{i=s+1}^nd_i\eta _i(x)e^{\lambda _i^{}x},x\mathrm{},\hfill \end{array}$$
(2.15)
where $`d_i(i=1,\mathrm{},n)`$ are complex constants.
Now suppose $`\{\mathrm{\Phi }^{(1)},\mathrm{\Phi }^{(2)},\mathrm{},\mathrm{\Phi }^{(N)}\}`$ are $`N`$ solitary waves of Eq. (2.2) with
$$\mathrm{\Phi }^{(k)}(x)\{\begin{array}{c}_{i=1}^sc_i^{(k)}\xi _i(x)e^{\lambda _ix},x\mathrm{},\hfill \\ _{i=s+1}^nc_i^{(k)}\xi _i(x)e^{\lambda _ix},x\mathrm{}.\hfill \end{array}$$
(2.16)
For each $`\mathrm{\Phi }^{(k)}`$, the single linearly independent localized solution $`\mathrm{\Psi }^{(k)}`$ of the adjoint equation
$$d\mathrm{\Psi }^{(k)}/dx=G^\text{T}(\mathrm{\Phi }^{(k)})\mathrm{\Psi }^{(k)}$$
(2.17)
has the following asymptotic behavior at infinity:
$$\mathrm{\Psi }^{(k)}(x)\{\begin{array}{c}_{i=1}^sd_i^{(k)}\eta _i(x)e^{\lambda _i^{}x},x\mathrm{},\hfill \\ _{i=s+1}^nd_i^{(k)}\eta _i(x)e^{\lambda _i^{}x},x\mathrm{}.\hfill \end{array}$$
(2.18)
Consider a new solitary wave which looks like a superposition of the above $`N`$ solitary waves $`\{\mathrm{\Phi }^{(k)}\}`$ widely separated, with the $`k`$-th wave $`\mathrm{\Phi }^{(k)}`$ located at $`x=x_k`$ $`(k=1,2,\mathrm{},N)`$. Let
$$x_1<x_2<\mathrm{}<x_N,$$
(2.19)
and denote
$$\mathrm{}_k=x_{k+1}x_k(1),k=1,2,\mathrm{},N1.$$
(2.20)
We will call this new solitary wave as a $`N`$-pulse wavetrain. It can be constructed explicitly by the following theorem.
###### Theorem 1
Under the assumptions A1, A2, A3 and the above notations, the $`N`$ solitary waves $`\{\mathrm{\Phi }^{(1)},\mathrm{},\mathrm{\Phi }^{(N)}\}`$ can match each other and form a widely-separated $`N`$-pulse wavetrain if and only if the spacings $`\mathrm{}_k(1)(k=1,\mathrm{},N1)`$ asymptotically satisfy the following $`N`$ conditions
$$\underset{j=s+1}{\overset{n}{}}c_j^{(2)}d_j^{(1)}e^{\lambda _j\mathrm{}_1}=0,$$
(2.21a)
$$\underset{j=1}{\overset{s}{}}c_j^{(k1)}d_j^{(k)}e^{\lambda _j\mathrm{}_{k1}}=\underset{j=s+1}{\overset{n}{}}c_j^{(k+1)}d_j^{(k)}e^{\lambda _j\mathrm{}_k},(2kN1),$$
(2.21b)
$$\underset{j=1}{\overset{s}{}}c_j^{(N1)}d_j^{(N)}e^{\lambda _j\mathrm{}_{N1}}=0.$$
(2.21c)
The relative errors in Eqs. (1) are exponentially small with the spacings.
We will prove this theorem by the asymptotic tail-matching method to be developed next.
Proof: Suppose such a $`N`$-pulse wavetrain $`\mathrm{\Phi }(x)`$ exists. Then around the $`k`$-th wave $`(2kN1)`$, the solution is
$$\mathrm{\Phi }(x)=\mathrm{\Phi }^{(k)}(xx_k)+\stackrel{~}{\mathrm{\Phi }}^{(k)}(xx_k),$$
(2.22)
where $`\stackrel{~}{\mathrm{\Phi }}^{(k)}1`$. The linearized equation for $`\stackrel{~}{\mathrm{\Phi }}^{(k)}`$ is
$$d\stackrel{~}{\mathrm{\Phi }}^{(k)}(x)/dx=G(\mathrm{\Phi }^{(k)}(x))\stackrel{~}{\mathrm{\Phi }}^{(k)}(x).$$
(2.23)
According to assumption A3, Eq. (2.23) has a single linearly independent localized solution which is $`d\mathrm{\Phi }^{(k)}(x)/dx`$. From Eq. (2.16) we get
$$d\mathrm{\Phi }^{(k)}(x)/dx\underset{i=1}{\overset{s}{}}c_i^{(k)}(\xi _i^{}(x)+\lambda _i\xi _i(x))e^{\lambda _ix},x\mathrm{}.$$
(2.24)
Clearly not all the $`c_i^{(k)}`$’s $`(i=1,\mathrm{},s)`$ are equal to zero. Without loss of generality, we assume that $`c_1^{(k)}0`$. Then we denote the other $`n1`$ solutions of Eq. (2.23) as $`\stackrel{~}{\mathrm{\Phi }}_j^{(k)}(j=2,\mathrm{},n)`$. We require that
$$\stackrel{~}{\mathrm{\Phi }}_j^{(k)}(x)\xi _j(x)e^{\lambda _jx},x\mathrm{}.$$
(2.25)
As $`x\mathrm{}`$, we generally have
$$\stackrel{~}{\mathrm{\Phi }}_j^{(k)}(x)\underset{i=1}{\overset{n}{}}a_{ji}^{(k)}\xi _i(x)e^{\lambda _ix},j=2,\mathrm{},n,$$
(2.26)
where $`a_{ji}^{(k)}`$ are constants. Since $`c_1^{(k)}0`$, these $`n`$ solutions $`\{d\mathrm{\Phi }^{(k)}/dx,\stackrel{~}{\mathrm{\Phi }}_2^{(k)},\mathrm{},\stackrel{~}{\mathrm{\Phi }}_n^{(k)}\}`$ are linearly independent at $`x`$ equal to infinity, and they form a fundamental set of solutions of Eq. (2.23). Thus the general solution for $`\stackrel{~}{\mathrm{\Phi }}^{(k)}(xx_k)`$ is
$$\stackrel{~}{\mathrm{\Phi }}^{(k)}(xx_k)=h_1\frac{d\mathrm{\Phi }^{(k)}}{dx}(xx_k)+\underset{j=2}{\overset{n}{}}h_j\stackrel{~}{\mathrm{\Phi }}_j^{(k)}(xx_k),$$
(2.27)
where $`h_j(j=1,\mathrm{},n)`$ are constants. The first term in (2.27) can be absorbed into $`\mathrm{\Phi }^{(k)}(xx_k)`$ and cause a position shift to it. By normalization we make $`h_1=0`$. When $`x_kxx_{k+1}`$, dropping the exponentially small terms, we get
$$\stackrel{~}{\mathrm{\Phi }}^{(k)}(xx_k)\underset{j=s+1}{\overset{n}{}}h_j\xi _j(xx_k)e^{\lambda _j(xx_k)}.$$
(2.28)
Similarly, when $`x_{k1}xx_k`$, we have
$$\stackrel{~}{\mathrm{\Phi }}^{(k)}(xx_k)\underset{j=1}{\overset{s}{}}(\underset{i=2}{\overset{n}{}}a_{ij}^{(k)}h_i)\xi _j(xx_k)e^{\lambda _j(xx_k)}.$$
(2.29)
The key idea in the asymptotic tail-matching method is that, in order for the matching to occur, we need to require that in the region $`x_kxx_{k+1}`$, $`\stackrel{~}{\mathrm{\Phi }}^{(k)}(xx_k)`$’s exponentially growing terms match the left tail of the right-hand wave $`\mathrm{\Phi }^{(k+1)}(xx_{k+1})`$; in the region $`x_{k1}xx_k`$, $`\stackrel{~}{\mathrm{\Phi }}^{(k)}(xx_k)`$’s exponentially decaying terms match the right tail of the left-hand wave $`\mathrm{\Phi }^{(k1)}(xx_{k1})`$. In the region $`x_kxx_{k+1}`$, this requirement is
$$\underset{j=s+1}{\overset{n}{}}h_j\xi _j(xx_k)e^{\lambda _j(xx_k)}=\underset{j=s+1}{\overset{n}{}}c_j^{(k+1)}\xi _j(xx_{k+1})e^{\lambda _j(xx_{k+1})};$$
(2.30)
and in the region $`x_{k1}xx_k`$, it is
$$\underset{j=1}{\overset{s}{}}(\underset{i=2}{\overset{n}{}}a_{ij}^{(k)}h_i)\xi _j(xx_k)e^{\lambda _j(xx_k)}=\underset{j=1}{\overset{s}{}}c_j^{(k1)}\xi _j(xx_{k1})e^{\lambda _j(xx_{k1})}.$$
(2.31)
We now need to select the constants $`h_j(j=2,\mathrm{},n)`$ and spacings $`\mathrm{}_k(k=1,\mathrm{},N1)`$ so that the above two conditions are satisfied.
First consider condition (2.30). Recall that functions $`\{\xi _j(x)\}`$ are of the form (2.9). If for $`\lambda =\lambda _m(s+1mn)`$, the chain of such functions is $`\{\xi _m,\mathrm{},\xi _{m+l1}\}`$, where $`\xi _m`$ is a constant vector and
$$\xi _{j+1}^{}(x)=\xi _j(x),j=m,\mathrm{},m+l2.$$
(2.32)
Then we select $`\{h_m,\mathrm{},h_{m+j1}\}`$ from the equation
$$\underset{j=m}{\overset{m+l1}{}}h_j\xi _j(xx_k)e^{\lambda (xx_k)}=\underset{j=m}{\overset{m+l1}{}}c_j^{(k+1)}\xi _j(xx_{k+1})e^{\lambda (xx_{k+1})}.$$
(2.33)
The right hand side of this equation is
$$\underset{j=m}{\overset{m+l1}{}}c_j^{(k+1)}\xi _j(xx_k\mathrm{}_k)e^{\lambda (xx_k)}e^{\lambda \mathrm{}_k}$$
$$=\underset{j=m}{\overset{m+l1}{}}c_j^{(k+1)}(\underset{i=0}{\overset{jm}{}}\frac{(\mathrm{}_k)^i}{i!}\frac{d^i\xi _j}{dx^i}(xx_k))e^{\lambda (xx_k)}e^{\lambda \mathrm{}_k}$$
$$=\underset{j=m}{\overset{m+l1}{}}\underset{i=m}{\overset{j}{}}c_j^{(k+1)}\frac{(\mathrm{}_k)^{ji}}{(ji)!}\xi _i(xx_k)e^{\lambda (xx_k)}e^{\lambda \mathrm{}_k}$$
$$=\underset{i=m}{\overset{m+l1}{}}(\underset{j=i}{\overset{m+l1}{}}\frac{(\mathrm{}_k)^{ji}}{(ji)!}c_j^{(k+1)})e^{\lambda \mathrm{}_k}\xi _i(xx_k)e^{\lambda (xx_k)}.$$
(2.34)
Now we choose $`h_i(i=m,\mathrm{},m+l1)`$ to be
$$h_i=(\underset{j=i}{\overset{m+l1}{}}\frac{(\mathrm{}_k)^{ji}}{(ji)!}c_j^{(k+1)})e^{\lambda \mathrm{}_k},i=m,\mathrm{},m+l1,$$
(2.35)
then Eq. (2.33) is valid. Repeating this procedure for the other chains of {$`\xi _j(x)`$} functions in the form (2.9), we can successfully select $`h_i(i=s+1,\mathrm{},n)`$ so that condition (2.30) is satisfied.
Next consider condition (2.31). Similar analysis shows that we can reduce its right hand side to
$$\underset{j=1}{\overset{s}{}}c_j^{(k1)}\xi _j(xx_{k1})e^{\lambda _j(xx_{k1})}=\underset{j=1}{\overset{s}{}}\alpha _je^{\lambda _j\mathrm{}_{k1}}\xi _j(xx_k)e^{\lambda _j(xx_k)},$$
(2.36)
where $`\alpha _j(j=1,\mathrm{},s)`$ are constants and determined by $`c_j^{(k1)}(j=1,\mathrm{},s)`$ and $`\mathrm{}_{k1}`$. Then condition (2.31) becomes
$$\underset{i=2}{\overset{s}{}}a_{ij}^{(k)}h_i=\alpha _je^{\lambda _j\mathrm{}_{k1}}\underset{i=s+1}{\overset{n}{}}a_{ij}^{(k)}h_i,j=1,\mathrm{},s.$$
(2.37)
This is a linear system of $`s`$ equations for $`s1`$ unknowns $`h_i(i=2,\mathrm{},s)`$. We now show that the matrix $`(a_{ij}^{(k)})_{s\times (s1)}`$ on the left side of Eq. (2.37) has rank $`s1`$. Consider the solution of Eq. (2.23)
$$T(x)=\underset{j=2}{\overset{s}{}}p_j\stackrel{~}{\mathrm{\Phi }}_j^{(k)}(x),$$
(2.38)
where $`p_j(j=2,\mathrm{},s)`$ are constants. Dropping exponentially small terms we get
$$T(x)\{\begin{array}{c}0,x\mathrm{},\hfill \\ _{j=1}^s(_{i=2}^sa_{ij}^{(k)}p_i)\xi _j(x)e^{\lambda _jx},x\mathrm{}.\hfill \end{array}$$
(2.39)
According to assumption A3, the only localized solution of Eq. (2.23) is $`d\mathrm{\Phi }^{(k)}(x)/dx`$. Moreover, $`c_1^{(k)}`$ in (2.24) is non-zero. Thus $`T(x)`$ can not be a localized solution. In other words, the linear system of equations
$$\underset{i=2}{\overset{s}{}}a_{ij}^{(k)}p_i=0,j=1,\mathrm{},s,$$
(2.40)
has no non-trivial solutions for $`p_i(i=1,\mathrm{},s)`$. Therefore the matrix $`(a_{ij}^{(k)})_{s\times (s1)}`$ has rank $`s1`$. Without loss of generality, we assume that the last $`(s1)`$ rows of the matrix are linearly independent. Then the linear system
$$\underset{i=2}{\overset{s}{}}a_{ij}^{(k)}h_i=\alpha _je^{\lambda _j\mathrm{}_{k1}}\underset{i=s+1}{\overset{n}{}}a_{ij}^{(k)}h_i,j=2,\mathrm{},s,$$
(2.41)
has a unique solution for $`h_i(i=2,\mathrm{},s)`$. With $`h_i(i=2,\mathrm{},n)`$ given by (2.35) and (2.41), the only matching condition left to be satisfied now is
$$\underset{i=2}{\overset{s}{}}a_{i1}^{(k)}h_i=\alpha _1e^{\lambda _1\mathrm{}_{k1}}\underset{i=s+1}{\overset{n}{}}a_{i1}^{(k)}h_i,$$
(2.42)
which will determine the spacings of this $`N`$-pulse wavetrain. Since the matrix $`(a_{ij}^{(k)})_{s\times (n1)}`$ is not readily available, to determine the spacings from Eq. (2.42) is difficult. But this can be easily done with the aid of the solution $`\mathrm{\Psi }^{(k)}`$ of the adjoint equation (2.17). With $`h_i(i=2,\mathrm{},n)`$ given by (2.35) and (2.41), it is easy to show that Eqs. (2.28) and (2.29) become
$$\stackrel{~}{\mathrm{\Phi }}^{(k)}(xx_k)\underset{j=s+1}{\overset{n}{}}c_j^{(k+1)}\xi _j(xx_{k+1})e^{\lambda _j(xx_{k+1})},x_kxx_{k+1},$$
(2.43)
and
$$\stackrel{~}{\mathrm{\Phi }}^{(k)}(xx_k)w\xi _1(xx_{k1})e^{\lambda _1(xx_{k1})}+\underset{j=2}{\overset{s}{}}c_j^{(k1)}\xi _j(xx_{k1})e^{\lambda _j(xx_{k1})},x_{k1}xx_k,$$
(2.44)
where $`w`$ is a constant. Condition (2.42) is equivalent to
$$w=c_1^{(k1)}.$$
(2.45)
For $`x_{k1}y_1x_k`$ and $`x_ky_2x_{k+1}`$, we have
$$\begin{array}{ccc}\hfill 0& =& _{y_1}^{y_2}\{d\stackrel{~}{\mathrm{\Phi }}^{(k)}(xx_k)/dxG(\mathrm{\Phi }^{(k)}(xx_k))\stackrel{~}{\mathrm{\Phi }}^{(k)}(xx_k)\}\mathrm{\Psi }^{(k)}(xx_k)𝑑x\hfill \\ & =& \stackrel{~}{\mathrm{\Phi }}^{(k)}(xx_k))\mathrm{\Psi }^{(k)}(xx_k)|^{y_2}_{y_1}\hfill \\ & & +_{y_1}^{y_2}\stackrel{~}{\mathrm{\Phi }}^{(k)}(xx_k)\{d\mathrm{\Psi }^{(k)}(xx_k)/dxG^\text{T}(\mathrm{\Phi }^{(k)}(xx_k))\mathrm{\Psi }^{(k)}(xx_k)\}^{}𝑑x\hfill \\ & =& \stackrel{~}{\mathrm{\Phi }}^{(k)}(xx_k))\mathrm{\Psi }^{(k)}(xx_k)|^{y_2}_{y_1}.\hfill \end{array}$$
(2.46)
For $`\mathrm{}_{k1}1`$ and $`\mathrm{}_k1`$, asymptotically we get
$$wd_1^{(k)}e^{\lambda _1\mathrm{}_{k1}}+\underset{j=2}{\overset{s}{}}c_j^{(k1)}d_j^{(k)}e^{\lambda _j\mathrm{}_{k1}}=\underset{j=s+1}{\overset{n}{}}c_j^{(k+1)}d_j^{(k)}e^{\lambda _j\mathrm{}_k}.$$
(2.47)
Condition (2.45) is satisfied if and only if Eq. (1b) is valid. For the first and last waves in this $`N`$-pulse wavetrain, the analysis is simpler, and we get Eqs. (1a,c) for matching. In summary, the $`N`$ pulses $`\{\mathrm{\Phi }^{(1)},\mathrm{},\mathrm{\Phi }^{(N)}\}`$ can be matched and form a $`N`$-pulse wavetrain if and only if the spacings $`\mathrm{}_k(1)(k=1,\mathrm{},N1)`$ asymptotically satisfy the $`N`$ conditions (1).
Now we discuss the accuracy of the above results. Error is created mainly by the matching requirements (see Eqs. (2.30) and (2.31)) and the negligence of nonlinear terms in Eq. (2.23). First we discuss the error in the matching requirements. Let us reconsider the solution (2.22) around the $`k`$-th wave. When $`x_kxx_{k+1}`$, beside the exponentially decaying terms in $`\mathrm{\Phi }^{(k)}(xx_k)`$, there are also such terms in $`\stackrel{~}{\mathrm{\Phi }}^{(k)}(xx_k)`$ (see Eq. (2.27)). The combined exponentially decaying tails are
$$c_1^{(k)}\xi _1(xx_k)e^{\lambda _1(xx_k)}+\underset{j=2}{\overset{s}{}}(c_j^{(k)}+h_j)\xi _j(xx_k)e^{\lambda _j(xx_k)}.$$
(2.48)
Thus Eqs. (1) would be more accurate if the $`c_j^{(k)}`$ values are replaced by $`c_j^{(k)}+h_j`$. Recall that $`h_j(j=2,\mathrm{},s)`$ are determined from Eq. (2.41), so they are exponentially small for large $`\mathrm{}_{k1}`$ and $`\mathrm{}_k`$. As a result, the negligence of tail contribution from $`\stackrel{~}{\mathrm{\Phi }}^{(k)}(xx_k)`$ only causes exponentially small relative errors in Eqs. (1). Simple reasoning also shows that the exclusion of nonlinear terms in Eq. (2.4) also causes only exponentially small relative errors in (1). The proof of theorem 1 is now completed. It should be pointed out that, if the eigenvalues {$`\lambda _i,i=1,\mathrm{},s`$} or {$`\lambda _i,i=s+1,\mathrm{},n`$} are real-valued and close to each other, those exponentially small relative errors in Eqs. (1) may become significant. In such cases, caution is needed in interpreting the results from (1)
### 2.2 General permanent-wave trains
The results in the previous section can be readily extended to the matching of permanent waves which exponentially approach a complex constant at infinity. Suppose such permanent waves exist in Eq. (2.2), then we make the following general assumptions: for any permanent wave $`\mathrm{\Phi }(x)`$ where
$$\mathrm{\Phi }(x)\{\begin{array}{c}b_2,x\mathrm{},\hfill \\ b_1,x\mathrm{},\hfill \end{array}$$
(2.49)
1. the eigenvalues of the constant (Jacobian) matrices $`G(b_1)`$ and $`G(b_2)`$ all have non-zero real parts, and the number of $`G(b_1)`$’s eigenvalues with negative real parts is equal to that of $`G(b_2)`$’s eigenvalues with negative real parts;
2. the linear behavior dominates at infinity, i.e., as $`x\mathrm{}`$ and $`\mathrm{}`$, $`\mathrm{\Phi }(x)`$ approaches a solution of the linear equations
$$d\stackrel{~}{\mathrm{\Phi }}/dx=G(b_1)\stackrel{~}{\mathrm{\Phi }},$$
(2.50)
and
$$d\stackrel{~}{\mathrm{\Phi }}/dx=G(b_2)\stackrel{~}{\mathrm{\Phi }},$$
(2.51)
respectively;
3. the linearized equation of (2.2) around $`\mathrm{\Phi }(x)`$
$$d\stackrel{~}{\mathrm{\Phi }}/dx=G(\mathrm{\Phi })\stackrel{~}{\mathrm{\Phi }}$$
(2.52)
and its adjoint equation
$$d\mathrm{\Psi }/dx=G^\text{T}(\mathrm{\Phi })\mathrm{\Psi }$$
(2.53)
each have a single linearly independent localized solution.
Now suppose $`\{\mathrm{\Phi }^{(1)},\mathrm{\Phi }^{(2)},\mathrm{},\mathrm{\Phi }^{(N)}\}`$ are $`N`$ permanent waves with
$$\mathrm{\Phi }^{(k)}(x)\{\begin{array}{c}b_2^{(k)},x\mathrm{},\hfill \\ b_1^{(k)},x\mathrm{},\hfill \end{array}$$
(2.54)
where $`1kN`$. If they are to be matched and form a widely-separated $`N`$-permanent-wave train, we need to require that
$$b_1^{(k)}=b_2^{(k1)},2kN.$$
(2.55)
One consequence is that all the matrices $`G(b_1^{(k)})`$ and $`G(b_2^{(k)})`$ $`(1kN)`$ have the same number of eigenvalues with negative real parts, which we denote as $`s`$. We introduce the following notations. Denote $`G(b_1^{(k)})`$’s $`n`$ eigenvalues as $`\lambda _1^{(k)},\lambda _2^{(k)},\mathrm{},\lambda _n^{(k)}`$ with
$$\text{Re}(\lambda _1^{(k)})\text{Re}(\lambda _2^{(k)})\mathrm{}\text{Re}(\lambda _s^{(k)})<0<\text{Re}(\lambda _{s+1}^{(k)})\text{Re}(\lambda _{s+2}^{(k)})\mathrm{}\text{Re}(\lambda _n^{(k)}),$$
(2.56)
and $`G(b_2^{(k)})`$’s as $`\mathrm{\Lambda }_1^{(k)},\mathrm{\Lambda }_2^{(k)},\mathrm{},\mathrm{\Lambda }_n^{(k)}`$ with
$$\text{Re}(\mathrm{\Lambda }_1^{(k)})\text{Re}(\mathrm{\Lambda }_2^{(k)})\mathrm{}\text{Re}(\mathrm{\Lambda }_s^{(k)})<0<\text{Re}(\mathrm{\Lambda }_{s+1}^{(k)})\text{Re}(\mathrm{\Lambda }_{s+2}^{(k)})\mathrm{}\text{Re}(\mathrm{\Lambda }_n^{(k)}).$$
(2.57)
The fundamental sets of solutions of the linear equations
$$d\stackrel{~}{\mathrm{\Phi }}/dx=G(b_1^{(k)})\stackrel{~}{\mathrm{\Phi }}$$
(2.58)
and
$$d\stackrel{~}{\mathrm{\Phi }}/dx=G(b_2^{(k)})\stackrel{~}{\mathrm{\Phi }}$$
(2.59)
are respectively $`\{\xi _i^{(k)}(x)e^{\lambda _i^{(k)}x},i=1,\mathrm{},n\}`$ and $`\{\theta _i^{(k)}(x)e^{\mathrm{\Lambda }_i^{(k)}x},i=1,\mathrm{},n\}`$ which consist of chains of linearly independent solutions of Eqs. (2.58) and (2.59) as defined before (see Eq. (2.9)). The fundamental matrices of the adjoint equations of (2.58) and (2.59) are then
$$[\eta _1^{(k)}e^{\lambda _1^{(k)}x}\eta _2^{(k)}e^{\lambda _2^{(k)}x}\mathrm{}\eta _n^{(k)}e^{\lambda _n^{(k)}x}]$$
(2.60)
and
$$[\zeta _1^{(k)}e^{\mathrm{\Lambda }_1^{(k)}x}\zeta _2^{(k)}e^{\mathrm{\Lambda }_2^{(k)}x}\mathrm{}\zeta _n^{(k)}e^{\mathrm{\Lambda }_n^{(k)}x}],$$
(2.61)
with
$$[\eta _1^{(k)}\eta _2^{(k)}\mathrm{}\eta _n^{(k)}]=\{[\xi _1^{(k)}\xi _2^{(k)}\mathrm{}\xi _n^{(k)}]^1\}^\text{T}$$
(2.62)
and
$$[\zeta _1^{(k)}\zeta _2^{(k)}\mathrm{}\zeta _n^{(k)}]=\{[\theta _1^{(k)}\theta _2^{(k)}\mathrm{}\theta _n^{(k)}]^1\}^\text{T}.$$
(2.63)
Note that
$$\lambda _i^{(k)}=\mathrm{\Lambda }_i^{(k1)},\xi _i^{(k)}=\theta _i^{(k1)},\eta _i^{(k)}=\zeta _i^{(k1)},i=1,\mathrm{},n$$
(2.64)
in view of (2.55). In those notations, we have
$$\mathrm{\Phi }^{(k)}(x)\{\begin{array}{c}b_2^{(k)}+_{i=1}^sc_i^{(k)}\theta _i^{(k)}(x)e^{\mathrm{\Lambda }_i^{(k)}x},x\mathrm{},\hfill \\ b_1^{(k)}+_{i=s+1}^nc_i^{(k)}\xi _i^{(k)}(x)e^{\lambda _i^{(k)}x},x\mathrm{},\hfill \end{array}$$
(2.65)
according to assumption B2. For the single linearly independent localized solution $`\mathrm{\Psi }^{(k)}`$ of the adjoint equation
$$d\mathrm{\Psi }^{(k)}/dx=G^\text{T}(\mathrm{\Phi }^{(k)})\mathrm{\Psi }^{(k)},$$
(2.66)
we have
$$\mathrm{\Psi }^{(k)}(x)\{\begin{array}{c}_{i=1}^sd_i^{(k)}\eta _i^{(k)}(x)e^{\lambda _i^{(k)}x},x\mathrm{},\hfill \\ _{i=s+1}^nd_i^{(k)}\zeta _i^{(k)}(x)e^{\mathrm{\Lambda }_i^{(k)}x},x\mathrm{}.\hfill \end{array}$$
(2.67)
Here $`c_i^{(k)}`$ and $`d_i^{(k)}`$ $`(1in)`$ are complex constants.
Now consider a widely-separated permanent-wave train matched by the above $`N`$ permanent waves $`\{\mathrm{\Phi }^{(1)}`$, $`\mathrm{\Phi }^{(2)}`$, …, $`\mathrm{\Phi }^{(N)}\}`$. Assume that the $`k`$-th wave $`\mathrm{\Phi }^{(k)}`$ is located at $`x=x_k`$ $`(k=1,2,\mathrm{},N)`$, and $`\mathrm{}_k`$ is as defined in Eq. (2.20), then we have the following result.
###### Theorem 2
Under the assumptions B1, B2, B3, (2.55) and the above notations, the $`N`$ permanent waves $`\{\mathrm{\Phi }^{(1)},\mathrm{},\mathrm{\Phi }^{(N)}\}`$ can match each other and form a widely-separated $`N`$-permanent-wave train if and only if the spacings $`\mathrm{}_k(1)(k=1,\mathrm{},N1)`$ asymptotically satisfy the following $`N`$ conditions
$$\underset{j=s+1}{\overset{n}{}}c_j^{(2)}d_j^{(1)}e^{\mathrm{\Lambda }_j^{(1)}\mathrm{}_1}=0,$$
(2.68a)
$$\underset{j=1}{\overset{s}{}}c_j^{(k1)}d_j^{(k)}e^{\lambda _j^{(k)}\mathrm{}_{k1}}=\underset{j=s+1}{\overset{n}{}}c_j^{(k+1)}d_j^{(k)}e^{\mathrm{\Lambda }_j^{(k)}\mathrm{}_k},(2kN1),$$
(2.68b)
$$\underset{j=1}{\overset{s}{}}c_j^{(N1)}d_j^{(N)}e^{\lambda _j^{(N)}\mathrm{}_{N1}}=0.$$
(2.68c)
The relative errors in Eqs. (2) are exponentially small with the spacings.
The proof for this theorem is similar to that for theorem 1, and is thus omitted here.
Remark: In applying theorem 2 to a given nonlinear wave system, the major difficulty is the determination of the coefficients $`\{d_j^{(k)}\}`$ in the localized solution $`\mathrm{\Psi }^{(k)}(x)`$ of the adjoint equation (2.66). In general, this has to be done numerically. But in many cases, Eq. (2.52) can be cast into a self-adjoint system (see , and ). Then $`\mathrm{\Psi }^{(k)}(x)`$ and its coefficients $`\{d_j^{(k)}\}`$ can be readily obtained from $`d\mathrm{\Phi }^{(k)}/dx`$, and the verification of conditions (2) can proceed.
## 3 Applications
### 3.1 Fourth-order systems
The permanent waves in many nonlinear wave problems are governed by fourth-order systems (2.2) (see , and ). In this section, we apply theorems 1 and 2 to certain classes of such systems. In particular, we will establish the existence of countably infinite multiple permanent-wave trains under some general assumptions.
We first consider the matching of identical permanent waves in a fourth-order system (2.2). Suppose $`\mathrm{\Phi }(x)`$ is a permanent wave in (2.2) where $`\mathrm{\Phi }(x)b`$ as $`|x|\mathrm{}`$, and the assumptions B1, B2 and B3 are satisfied. Moreover, we suppose the eigenvalues of $`G(b)`$ are $`\pm \lambda _1`$ and $`\pm \lambda _2`$, where $`\lambda _1\lambda _2`$ and $`\text{Re}(\lambda _i)>0(i=1,2)`$. Then corresponding to the four distinct eigenvalues $`\lambda _1,\lambda _2,\lambda _1`$ and $`\lambda _2`$, $`G(b)`$ has four linearly independent eigenvectors $`\xi _i(i=1,\mathrm{},4)`$. If we denote
$$[\eta _1\eta _2\eta _3\eta _4]=\{[\xi _1\xi _2\xi _3\xi _4]^1\}^\text{T},$$
(3.1)
then
$$\mathrm{\Phi }(x)\{\begin{array}{c}b+c_1\xi _1e^{\lambda _1x}+c_2\xi _2e^{\lambda _2x},x\mathrm{},\hfill \\ b+c_3\xi _3e^{\lambda _1x}+c_4\xi _4e^{\lambda _2x},x\mathrm{},\hfill \end{array}$$
(3.2)
and
$$\mathrm{\Psi }(x)\{\begin{array}{c}d_1\eta _1e^{\lambda _1^{}x}+d_2\eta _2e^{\lambda _2^{}x},x\mathrm{},\hfill \\ d_3\eta _3e^{\lambda _1^{}x}+d_4\eta _4e^{\lambda _2^{}x},x\mathrm{}.\hfill \end{array}$$
(3.3)
For some fourth-order problems, Eq. (2.52) can be cast into a self-adjoint system, and one has either $`\lambda _1`$ and $`\lambda _2`$ real-valued with $`(d_1d_3)(c_3c_1)`$ and $`(d_2d_4)(c_4c_2)`$, or $`\lambda _1`$ and $`\lambda _2`$ complex-valued with $`\lambda _2=\lambda _1^{}`$, $`(d_1d_3)(c_4c_2)`$ and $`(d_2d_4)(c_3c_1)`$. In such cases, conditions (2) for the matching of $`N`$ identical permanent waves $`\{\mathrm{\Phi }(x),\mathrm{},\mathrm{\Phi }(x)\}`$ simply become
$$c_1d_1^{}e^{\lambda _1\mathrm{}_k}+c_2d_2^{}e^{\lambda _2\mathrm{}_k}=0,k=1,\mathrm{},N1.$$
(3.4)
In the second case, if furthermore (2.2) is a real system, then $`\lambda _2=\lambda _1^{}`$, $`c_2=c_1^{},d_2=d_1^{}`$, and Eq. (3.4) becomes
$$c_1d_1^{}e^{i\text{Im}(\lambda _1)\mathrm{}_k}+c_1^{}d_1e^{i\text{Im}(\lambda _1)\mathrm{}_k}=0,k=1,\mathrm{},N1.$$
(3.5)
The spacings $`\mathrm{}_k`$ can then be easily obtained from (3.5) as
$$\mathrm{}_k=(\text{arg}(c_1d_1^{})+\frac{\pi }{2}+m_k\pi )/\text{Im}(\lambda _1),k=1,\mathrm{},N1,$$
(3.6)
where $`m_k`$ is any non-negative integer. Note that in this case, the exponentially small relative errors in (2) make little difference, especially when $`m_k`$ is large. Thus we conclude that an arbitrary number of identical permanent waves $`\mathrm{\Phi }(x)`$ can be matched together and form multiple permanent-wave trains, whose spacings are given asymptotically by Eq. (3.6). Clearly a countably infinite number of such wavetrains can be formed. In the paper by Buffoni and Sere , they proved the existence of countably infinite multi-pulse permanent wave solutions for a class of coupled-nonlinear-Schrödinger-type equations. When those equations are cast into a fourth-order system of the two variables and their first derivatives, it is easy to check that they fall into the above category. Thus their result is a special case of ours. But differences also exist between their result and ours. In their result, $`m_k`$ in Eq. (3.6) is an even integer; while in ours, it is any integer. This means that we identified twice as many solitary-wave trains as they did. For fourth-order systems where $`\mathrm{\Phi }(x)b`$ and $`\mathrm{\Psi }(x)`$ element-wise are both even or odd in $`x`$, or one of them is even (odd) and the other one odd (even), then $`(c_3c_4)=\pm (c_1c_2)`$, and $`(\xi _3\xi _4)`$ is row-wise equal to or opposite of $`(\xi _1\xi _2)`$. It is easy to show from (3.1) that $`(\eta _3\eta _4)`$ is also row-wise equal to or opposite of $`(\eta _1\eta _2)`$ and $`(d_3d_4)=\pm (d_1d_2)`$. Thus the $`N`$ matching conditions (2) also reduce to (3.4). If further more, $`\lambda _2=\lambda _1^{}`$, then we will find countably infinite multiple permanent-wave trains whose spacings are given by (3.6).
Next we consider the matching of different permanent waves in a fourth-order system (2.2). Suppose $`G`$ is an odd function of $`\mathrm{\Phi }`$, i.e.,
$$G(\mathrm{\Phi })=G(\mathrm{\Phi }),$$
(3.7)
and $`\mathrm{\Phi }(x)`$ is a permanent wave in Eq. (2.2) with
$$\mathrm{\Phi }(x)\{\begin{array}{c}b,x\mathrm{},\hfill \\ b,x\mathrm{},\hfill \end{array}$$
(3.8)
then $`\mathrm{\Phi }(x)`$ is also a permanent wave in (2.2). It is easy to show from (3.7) that $`G(\mathrm{\Phi })=G(\mathrm{\Phi })`$, thus $`G(b)=G(b)`$. Beside the assumptions B1, B2 and B3, we also assume that $`G(b)`$’s four eigenvalues are $`\pm \lambda _1`$ and $`\pm \lambda _2`$ with $`\lambda _1\lambda _2`$ and $`\text{Re}(\lambda _i)>0(i=1,2)`$. Suppose the eigenvectors corresponding to $`\lambda _1,\lambda _2,\lambda _1`$ and $`\lambda _2`$ are denoted as $`\xi _i(i=1,\mathrm{},4)`$, then we have
$$\mathrm{\Phi }(x)\{\begin{array}{c}b+c_1\xi _1e^{\lambda _1x}+c_2\xi _2e^{\lambda _2x},x\mathrm{},\hfill \\ b+c_3\xi _3e^{\lambda _1x}+c_4\xi _4e^{\lambda _2x},x\mathrm{}.\hfill \end{array}$$
(3.9)
For the localized solution $`\mathrm{\Psi }(x)`$ of the adjoint equation (2.53), we have
$$\mathrm{\Psi }(x)\{\begin{array}{c}d_1\eta _1e^{\lambda _1^{}x}+d_2\eta _2e^{\lambda _2^{}x},x\mathrm{},\hfill \\ d_3\eta _3e^{\lambda _1^{}x}+d_4\eta _4e^{\lambda _2^{}x},x\mathrm{},\hfill \end{array}$$
(3.10)
where $`\eta _i(i=1,\mathrm{},4)`$ are given by (3.1). For those equations (2.2) where Eq. (2.52) can be cast into a self-adjoint system and one has either $`(d_1d_3)(c_3c_1)`$ and $`(d_2d_4)(c_4c_2)`$ with $`\lambda _1`$ and $`\lambda _2`$ real, or $`(d_1d_3)(c_4c_2)`$ and $`(d_2d_4)(c_3c_1)`$ with $`\lambda _2=\lambda _1^{}`$, conditions (2) for the matching of permanent waves $`\{\mathrm{\Phi },\mathrm{\Phi },\mathrm{\Phi },\mathrm{\Phi },\mathrm{}\}`$ or $`\{\mathrm{\Phi },\mathrm{\Phi },\mathrm{\Phi },\mathrm{\Phi },\mathrm{}\}`$ will also reduce to (3.4). When $`\lambda _2=\lambda _1^{}`$, if furthermore (2.2) is a real system, then we can show as before that such matchings are always possible and the spacings are given by Eq. (3.6). An infinite number of such wavetrains will be obtained. We point out that the fourth-order systems studied by Kalies and VanderVorst falls into this category and is thus a special case of the above results. Here again we identified twice as many permanent-wave trains as they did since $`m_k`$ in Eq. (3.6) needs to be an even integer in their result.
### 3.2 Coupled nonlinear Schrödinger equations
The coupled nonlinear Schrödinger equations govern the evolution of two interacting wave packets in nonlinear and dispersive physical systems . These equations are particularly important in nonlinear optics as they govern the pulse propagation in birefringent nonlinear optical fibers . In recent years, the experimental design of high-speed optical-soliton-based telecommunication systems stimulated great interest in these equations, and much work has been done on them (see and the references therein). In particular, simple and multi-pulse solitary waves in these equations have been found and classified in . In this section, we study the multiple permanent-wave trains in these equations. We primarily discuss the focusing case where solitary waves exist. In the end of this section, we comment on the defocusing case where dark solitons arise.
The solitary waves in coupled nonlinear Schrödinger equations (focusing case) are governed by the following set of equations
$$r_{1xx}r_1+(r_1^2+\beta r_2^2)r_1=0,$$
(3.11a)
$$r_{2xx}\omega ^2r_2+(r_2^2+\beta r_1^2)r_2=0,$$
(3.11b)
where $`r_1`$ and $`r_2`$ approach zero as $`x`$ goes to infinity, and $`\beta `$ and $`\omega `$ are positive parameters. To apply theorem 1 to these equations, we first rewrite them as the following first order system
$$dU/dx=G(U),$$
(3.12)
where
$$U=(u_1,u_2,u_3,u_4)^T=(r_1,r_{1x},r_2,r_{2x})^T,$$
(3.13)
and
$$G(U)=\left(\begin{array}{c}u_2\\ u_1(u_1^2+\beta u_3^2)u_1\\ u_4\\ \omega ^2u_3(u_3^2+\beta u_1^2)u_3\end{array}\right).$$
(3.14)
It is easy to check that the above system satisfies the assumptions A1, A2 and A3 when $`\beta 1`$. Thus in the following we assume that $`\beta 1`$. The eigenvalues of the matrix $`G(0)`$ are $`1,\omega ,1`$ and $`\omega `$, and the corresponding eigenvectors are
$$\xi _1=\left(\begin{array}{c}1\\ 1\\ 0\\ 0\end{array}\right),\xi _2=\left(\begin{array}{c}0\\ 0\\ 1\\ \omega \end{array}\right),\xi _3=\left(\begin{array}{c}1\\ 1\\ 0\\ 0\end{array}\right),\xi _2=\left(\begin{array}{c}0\\ 0\\ 1\\ \omega \end{array}\right).$$
(3.15)
For a solitary wave $`(r_1,r_2)`$ of Eqs. (3.2) with
$$r_1(x)\{\begin{array}{c}c_1e^x,x\mathrm{},\hfill \\ c_3e^x,x\mathrm{},\hfill \end{array}$$
(3.16)
and
$$r_2(x)\{\begin{array}{c}c_2e^{\omega x},x\mathrm{},\hfill \\ c_4e^{\omega x},x\mathrm{},\hfill \end{array}$$
(3.17)
we have
$$U(x)\{\begin{array}{c}c_1\xi _1e^x+c_2\xi _2e^{\omega x},x\mathrm{},\hfill \\ c_3\xi _3e^x+c_4\xi _4e^{\omega x},x\mathrm{}.\hfill \end{array}$$
(3.18)
The linearized equation of (3.12) around a solitary wave $`U(x)`$ is
$$d\stackrel{~}{U}/dx=G(U)\stackrel{~}{U},$$
(3.19)
where $`\stackrel{~}{U}=(\stackrel{~}{u}_1,\stackrel{~}{u}_2,\stackrel{~}{u}_3,\stackrel{~}{u}_4)^T`$, and
$$G(U)=\left(\begin{array}{cccc}0& 1& 0& 0\\ 13u_1^2\beta u_3^2& 0& 2\beta u_1u_3& 0\\ 0& 0& 0& 1\\ 2\beta u_1u_3& 0& 13u_3^2\beta u_1^2& 0\end{array}\right).$$
(3.20)
The single localized solution of the above equations is $`dU/dx`$. If $`\stackrel{~}{u}_2`$ and $`\stackrel{~}{u}_4`$ in (3.19) are eliminated in favor of $`\stackrel{~}{u}_1`$ and $`\stackrel{~}{u}_3`$, then the linear system for $`\stackrel{~}{u}_1`$ and $`\stackrel{~}{u}_3`$ are self-adjoint. The adjoint equation of (3.19) is
$$dV/dx=G^\text{T}(U)V$$
(3.21)
with $`V=(v_1,v_2,v_3,v_4)`$. It is easy to see that when $`v_1`$ and $`v_3`$ are eliminated from (3.21), the equations for $`v_2`$ and $`v_4`$ are the same as those for $`\stackrel{~}{u}_1`$ and $`\stackrel{~}{u}_3`$. Thus the single localized solution of Eqs. (3.21) is
$$V=(u_{1xx},u_{1x},u_{3xx},u_{3x})^T.$$
(3.22)
At infinity,
$$V(x)\{\begin{array}{c}d_1\eta _1e^x+d_2\eta _2e^{\omega x},x\mathrm{},\hfill \\ d_3\eta _3e^x+d_4\eta _4e^{\omega x},x\mathrm{},\hfill \end{array}$$
(3.23)
where $`\{\eta _i,i=1,\mathrm{},4\}`$ are obtained from Eq. (3.1) as
$$\eta _1=\left(\begin{array}{c}\frac{1}{2}\\ \frac{1}{2}\\ 0\\ 0\end{array}\right),\eta _2=\left(\begin{array}{c}0\\ 0\\ \frac{1}{2}\\ \frac{1}{2\omega }\end{array}\right),\eta _3=\left(\begin{array}{c}\frac{1}{2}\\ \frac{1}{2}\\ 0\\ 0\end{array}\right),\eta _4=\left(\begin{array}{c}0\\ 0\\ \frac{1}{2}\\ \frac{1}{2\omega }\end{array}\right),$$
(3.24)
and
$$d_1=2c_3,d_3=2c_1,d_2=2\omega ^2c_4,d_4=2\omega ^2c_2.$$
(3.25)
Now we consider the matching of $`N`$ solitary waves $`\{(r_1^{(k)},r_2^{(k)}),k=1,\mathrm{},N\}`$ where
$$r_1^{(k)}(x)\{\begin{array}{c}c_1^{(k)}e^x,x\mathrm{},\hfill \\ c_3^{(k)}e^x,x\mathrm{},\hfill \end{array}$$
(3.26)
and
$$r_2^{(k)}(x)\{\begin{array}{c}c_2^{(k)}e^{\omega x},x\mathrm{},\hfill \\ c_4^{(k)}e^{\omega x},x\mathrm{}.\hfill \end{array}$$
(3.27)
In view of (3.25), the matching conditions (1) become
$$c_1^{(k)}c_3^{(k+1)}e^\mathrm{}_k+\omega ^2c_2^{(k)}c_4^{(k+1)}e^{\omega \mathrm{}_k}=0,k=1,\mathrm{},N1.$$
(3.28)
Interestingly Eq. (3.28) indicates that the $`N`$ solitary waves $`\{(r_1^{(k)},r_2^{(k)}),k=1,\mathrm{},N\}`$ can be matched if and only if all the adjacent solitary waves can. Thus the matching of solitary waves in Eqs. (3.2) is a “local” phenomenon. This fact would make the construction of those multiple solitary-wave trains much easier. In what follows, we discuss the matching of some special types of solitary waves.
First we consider the matching of wave and daughter wave solutions. In such solutions, either $`r_2r_1`$ or $`r_1r_2`$. Without loss of generality, we assume that $`r_2r_1`$. These solutions exist near the curves
$$\omega =(\sqrt{1+8\beta }1)/2m$$
(3.29)
in the $`(\omega ,\beta )`$ parameter plane . Here $`m`$ is a non-negative integer and $`m<(\sqrt{1+8\beta }1)/2`$. In these solutions, $`r_1`$ is symmetric; $`r_2`$ is symmetric for even values of $`m`$ and anti-symmetric for odd values of $`m`$. Suppose $`(\widehat{r}_1,\widehat{r}_2)`$ is such a solution, then
$$\widehat{r}_1c_1e^{|x|},|x|\mathrm{},$$
(3.30a)
and
$$\widehat{r}_2\{\begin{array}{c}c_2e^{\omega x},x\mathrm{},\hfill \\ (1)^mc_2e^{\omega x},x\mathrm{}.\hfill \end{array}$$
(3.30b)
Here $`c_21`$. Notice that if $`r_i(x)`$ ($`i`$=1 or 2) is a solution of Eqs. (3.2), so is $`r_i(x)`$. Without loss of generality, we require that $`c_i>0(i=1,2)`$. Now we consider $`N`$ wave and daughter wave solutions $`\{(r_1^{(k)},r_2^{(k)}),k=1,\mathrm{},N\}`$ where
$$r_1^{(k)}(x)=q_1^{(k)}\widehat{r}_1(x),r_2^{(k)}(x)=q_2^{(k)}\widehat{r}_2(x),k=1,\mathrm{},N,$$
(3.31)
and $`q_i^{(k)}=\pm 1(i=1,2)`$. The matching condition (3.28) for these solitary waves are simply
$$q_1^{(k)}q_1^{(k+1)}c_1^2e^\mathrm{}_k+(1)^mq_2^{(k)}q_2^{(k+1)}\omega ^2c_2^2e^{\omega \mathrm{}_k}=0,k=1,\mathrm{},N1,$$
(3.32)
i.e.
$$e^{(1\omega )\mathrm{}_k}=(1)^{m+1}\frac{q_2^{(k)}q_2^{(k+1)}}{q_1^{(k)}q_1^{(k+1)}}\frac{\omega ^2c_2^2}{c_1^2},k=1,\mathrm{},N1.$$
(3.33)
For these conditions to be satisfied, we need to require that $`\omega <1`$ and
$$(1)^{m+1}q_1^{(k)}q_1^{(k+1)}q_2^{(k)}q_2^{(k+1)}=1,k=1,\mathrm{},N1.$$
(3.34)
Suppose $`(q_1^{(k)},q_2^{(k)})`$ is fixed, then condition (3.34) shows that $`(q_1^{(k+1)},q_2^{(k+1)})`$ can take two sets of values. In other words, there are two possible types of matching. Thus these $`N`$ wave and daughter wave solutions can form $`2^N`$ topologically distinct solitary-wave trains. Since $`N`$ is arbitrary, countably infinite multiple-pulse solitary waves will be formed. The spacings between adjacent waves in those wavetrains are
$$\mathrm{}_k=\frac{\mathrm{ln}(\omega ^2c_2^2/c_1^2)}{\omega 1},k=1,\mathrm{},N1,$$
(3.35)
which are the same throughout an entire wavetrain. As $`\omega `$ approaches the wave and daughter wave boundary $`(\sqrt{1+8\beta }1)/2m`$, $`c_1`$ approaches $`2\sqrt{2}`$, $`c_2`$ approaches 0, and thus $`\mathrm{}_k`$ approaches infinity. The above theoretical results can be checked numerically. We first select $`(\beta ,\omega )`$ to be (2/3, 0.85) which is close to the curve (3.29) with $`m`$ equal to zero. With these parameter values, it is easy to find numerically that $`c_1`$ and $`c_2`$ as in Eq. (3.29) are equal to 2.6592 and 1.1744 respectively. Eqs. (3.34) and (3.35) then predict that the two wave and daughter waves $`(\widehat{r}_1,\widehat{r}_2)`$ and $`(\widehat{r}_1,\widehat{r}_2)`$ can be matched with the spacing approximately equal to 13.0635. This is indeed the case. Numerically we found this exact two-pulse solitary wave and plotted it in Fig. 1a. The exact spacing (measured as the distance between the two extrema in $`r_1`$) is 13.064, which is very close to the theoretical prediction. Next we select $`(\beta ,\omega )`$ to be (2, 0.6) which is close to the curve (3.29) with $`m=1`$. In this case, we numerically found that $`c_1`$ and $`c_2`$ in (3.29) are equal to 3.0386 and 0.6041. Then we predict from (3.34) and (3.35) that $`(\widehat{r}_1,\widehat{r}_2)`$ and itself can be matched with the spacing approximately equal to 10.6308. Indeed, that exact two-pulse solution was numerically found and plotted in Fig. 1b. The exact spacing is 10.40, close to the predicted value. The predictions on other types of matchings were also verified with good accuracy. We point out that each multiple-pulse solitary wave will generate a family of solitary waves as the parameter pair $`(\omega ,\beta )`$ moves away from the curves (3.29). Therefore countably infinite families of solitary waves will be generated near those curves.
Next we discuss mixed matchings between wave and daughter wave solutions and other types of solitary waves. When $`(\omega ,\beta )`$ is near the curve (3.29) with $`m=0`$, beside the wave and daughter wave solutions, another type of solitary waves (belonging to family $`D_2`$) also exist . Suppose $`(\widehat{r}_1,\widehat{r}_2)`$ is a wave and daughter wave solution whose large $`x`$ behavior is given by (3.29) (with $`m=0`$), and $`(\overline{r}_1,\overline{r}_2)`$ is a solitary wave with
$$\overline{r}_1\alpha _1e^{|x|},|x|\mathrm{},$$
(3.36a)
$$\overline{r}_2\alpha _2\text{sgn}(x)e^{\omega |x|},|x|\mathrm{},$$
(3.36b)
and $`\alpha _i>0(i=1,2)`$. Consider the mixed matching of the solitary waves $`(q_1\widehat{r}_1,q_2\widehat{r}_2)`$ and $`(q_3\overline{r}_1,q_4\overline{r}_2)`$ where $`q_i(i=1,\mathrm{},4)`$ are either 1 or $`1`$. The matching condition is
$$q_1q_3c_1\alpha _1e^{\mathrm{}}q_2q_4\omega ^2c_2\alpha _2e^\omega \mathrm{}=0,$$
(3.37)
or
$$e^{(1\omega )\mathrm{}}=\frac{q_2q_4}{q_1q_3}\frac{\omega ^2c_2\alpha _2}{c_1\alpha _1}$$
(3.38)
where $`\mathrm{}`$ is the spacing. This condition can be satisfied if and only if $`\omega <1`$ and the sign of $`q_1q_2q_3q_4`$ is equal to 1. As an example, we choose $`(\beta ,\omega )`$ as (2/3, 0.78). Then it is easy to find that $`c_1`$, $`c_2`$, $`\alpha _1`$ and $`\alpha _2`$ are 2.7967, 0.5210, 7.8105 and 8.4171 respectively. The above results predict that $`(\widehat{r}_1,\widehat{r}_2)`$ and $`(\overline{r}_1,\overline{r}_2)`$ can match each other and form a new two-pulse solitary wave. This was verified numerically. The exact matched solution is plotted in Fig. 2 with the spacing 10.26, while the predicted value for the spacing is 9.5571. Mixed matching between many copies of $`(\widehat{r}_1,\widehat{r}_2)`$ and $`(\overline{r}_1,\overline{r}_2)`$ can be similarly analysed. Once again, countably infinite multiple-pulse solitary waves will be formed by these mixed matchings.
Lastly we discuss the matching of solitary waves near $`\omega =1`$. In this case, single-hump solitary waves with $`r_1r_2`$ are present. Suppose $`(\widehat{r}_1,\widehat{r}_2)`$ is such a solution with
$$\widehat{r}_1c_1e^{|x|},|x|\mathrm{},$$
(3.39a)
$$\widehat{r}_2c_2e^{\omega |x|},|x|\mathrm{},$$
(3.39b)
then $`c_1c_2`$. If we consider the matching of these solitary waves $`\{(q_1^{(k)}\widehat{r}_1,q_2^{(k)}\widehat{r}_2)\}`$ where $`q_i^{(k)}=\pm 1(i=1,2)`$, the matching condition would again be Eq. (3.32) (with $`m=0`$). But here since $`G(0)`$’s eigenvalues 1 and $`\omega `$ are close, the exponentially small relative errors in (1) and (3.32) may become important. Thus condition (3.32) should be treated with caution. For instance, when $`(\beta ,\omega )`$ is (2, 0.99), we found $`c_1`$ and $`c_2`$ to be 1.6142 and 1.6355. In this case, $`\omega ^2c_2^2/c_1^2>1`$. Thus according to (3.32), $`(\widehat{r}_1,\widehat{r}_2)`$ and $`(\widehat{r}_1,\widehat{r}_2)`$ can not be matched. But our numerical results show otherwise .
Theorems 1 and 2 can also be used to study the matching of dark solitons which exist in coupled nonlinear Schrödinger equations (defocusing case) . In this case, our results on the matching of some classes of dark solitons indicate that such matchings are impossible since conditions (2) can not be satisfied. We suspect that any dark solitons can not match each other to form widely-separated dark-soliton trains.
## 4 Discussion
The results in this paper can be readily applied to general nonlinear wave systems for the construction of widely-separated multiple permanent-wave trains. Such wavetrains geometrically look like a superposition of individual permanent waves. This is somewhat analogous to the superposition principle of solutions in a linear system. But the difference here is that, due to the nonlinear nature of Eq. (2.2), those individual permanent waves have to be properly spaced (according to Eq. (1) or (2)) in order to form a wavetrain. When such wavetrains exist, one important question is their stability. For the coupled nonlinear Schrödinger equations, we indicated in that they are all unstable. For certain Ginzburg-Landau and coupled-nonlinear-Schrödinger type systems, Malomed argued that multi-pulse trains exist and are stable by an approximate method based on the variational principle and effective potential (, ). Such results need to be viewed with caution due to the approximations involved. The clear evidence that some multi-pulse waves are stable can be found in the experimental results on binary fluid convection () and the numerical results on subcritical Ginzburg-Landau equations (). We will investigate those systems in the near future.
## Acknowledgment
This work was supported in part by the National Science Foundation under the grant DMS-9622802. |
Subsets and Splits
No community queries yet
The top public SQL queries from the community will appear here once available.